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# 1 Introduction
## 1 Introduction
A strong field which is resonant for one of the allowed transitions in a medium can change the emission and absorption spectra corresponding to adjacent transitions \[1-3\]. This happens because the photons of a weak field may be emitted in a manner correlated with that in which photons of the strong field are emitted . In addition, a strong field changes the populations and splits the energy levels corresponding to the resonant transition .
Each of the three effects has a different dependence on the unsaturated-level populations and relaxation characteristics of the medium. These features are displayed most clearly for uniformly broadened optical transitions, because optical transitions are characterized by a larger number of relaxation constants than are microwave transitions, and the population differences corresponding to the various transitions may differ markedly. By choosing transitions appropriately one can suppress some effects while intensifying others.
Below we analyze the conditions for the appearance of each of these effects separately, and we analyze the associated changes in the spectral properties of the amplification-absorption coefficient for conditions typical of the optical range.
## 2 Equation for the Amplification Coefficient
We consider the amplification (or absorption) of a weak field at frequency $`\omega _\mu `$, approximately equal to the frequency $`\omega _{gn}`$ of the g-n transition, in the presence of a strong field $`𝐄`$ the frequency of which $`\omega `$ is approximately equal to the transition frequency $`\omega _{mn}`$($`E_m>E_g>E_n`$). The weak-field amplification coefficient $`\alpha _\mu (\mathrm{\Omega }_\mu )`$ is related in a simple manner to the emission power per unit volume $`w_{gn}(\mathrm{\Omega }_\mu )`$:
$$\alpha _\mu (\mathrm{\Omega }_\mu )=\left\{\frac{c}{8\pi }|E_\mu |^2\right\}^1w_{gn}(\mathrm{\Omega }_\mu ),$$
where $`E_\mu `$ is the amplitude of the ”weak field”, $`\mathrm{\Omega }_\mu =\omega _\mu \omega _{gn}`$, and $`w_{gn}`$ is calculated as in . We thus have:
$`\alpha _\mu `$ $`=`$ $`\alpha _\mu ^0\mathrm{\Gamma }_{gn}\{[\mathrm{\Gamma }_{gn}+i\mathrm{\Omega }_\mu ^{}+|G|^2(\mathrm{\Gamma }_{gm}+i(\mathrm{\Omega }_\mu ^{}\mathrm{\Omega }^{}))^1]^1`$ (2.1)
$`\times [1|G|^2{\displaystyle \frac{\mathrm{\Delta }n_{mn}/\mathrm{\Delta }n_{gn}}{\mathrm{\Gamma }^2(1+\varkappa )+\mathrm{\Omega }^2}}((1{\displaystyle \frac{\gamma _{mn}}{\mathrm{\Gamma }_m}}){\displaystyle \frac{2\mathrm{\Gamma }}{\mathrm{\Gamma }_n}}+{\displaystyle \frac{\mathrm{\Gamma }+i\mathrm{\Omega }^{}}{\mathrm{\Gamma }_{gm}+i(\mathrm{\Omega }_\mu ^{}\mathrm{\Omega }^{})}})]\}.`$
Here $`G=d_{mn}E/2\mathrm{}`$; $`\alpha _\mu ^0`$ is the coefficient at the line center in the absence of an external field ($`|G|^2=0`$); $`\mathrm{\Gamma }_{ik}`$, $`\mathrm{\Gamma }_i`$ are the Lorentz broadenings of the lines and levels ($`\mathrm{\Gamma }_{mn}\mathrm{\Gamma }`$); $`\gamma _m`$ is the probability for a transition from level $`m`$ to level $`n`$, $`\mathrm{\Omega }_\mu ^{}=\omega _\mu k_\mu v`$, $`\mathrm{\Omega }^{}=\omega kv`$, $`v`$ is the atomic velocity, $`(n_mn_n)/(n_gn_n)\mathrm{\Delta }n_{nm}/\mathrm{\Delta }n_{gn}`$ is the ratio of the unsaturated population differences corresponding to the transitions $`mn`$ and $`gn`$; $`\varkappa =(\mathrm{\Gamma }_m+\mathrm{\Gamma }_n\gamma _{mn})(\mathrm{\Gamma }_m\mathrm{\Gamma }_n\mathrm{\Gamma })^12|G|^2\tau ^22|G|^2`$. The population differences $`\rho _{mm}\rho _{nn}`$ and $`n_g\rho _{nn}`$ depend on the field in the following manner:
$`\rho _{mm}\rho _{nn}`$ $`=`$ $`(\mathrm{\Gamma }_2+\mathrm{\Omega }^2)\mathrm{\Delta }n_{mn}[\mathrm{\Gamma }^2(1+\varkappa )+\mathrm{\Omega }^2]^1,`$
$`n_g\rho _{nn}`$ $`=`$ $`\mathrm{\Delta }n_{gn}\mathrm{\Delta }n_{mn}\left(1{\displaystyle \frac{\gamma _{mn}}{\mathrm{\Gamma }_m}}\right){\displaystyle \frac{2\mathrm{\Gamma }}{\mathrm{\Gamma }_n}}|G|^2[\mathrm{\Gamma }^2(1+\varkappa )+\mathrm{\Omega }^2]^1.`$ (2.2)
The term proportional to $`|G|^2`$ in the common denominator in Eq. (2.1) reflects the broadening and splitting of the line by the strong field; expanding Eq. (2.1) in simple fractions, we can write the expression for $`\alpha _\mu `$ as:
$`\alpha _\mu `$ $`=`$ $`\alpha _\mu ^0Re\{{\displaystyle \frac{\mathrm{\Gamma }_{gn}(\alpha _1\alpha _2)^1}{\mathrm{\Gamma }_{gn}\alpha _2^{}+i(\mathrm{\Omega }_\mu ^{}\alpha _2^{\prime \prime })}}(\alpha _1|G|^2{\displaystyle \frac{\mathrm{\Delta }n_{mn}/\mathrm{\Delta }n_{gn}}{\mathrm{\Gamma }^2(1+\varkappa )+\mathrm{\Omega }^2}}`$ (2.3)
$`\times [(1{\displaystyle \frac{\gamma _{mn}}{\mathrm{\Gamma }_m}}){\displaystyle \frac{2\mathrm{\Gamma }}{\mathrm{\Gamma }_n}}\alpha _1(\mathrm{\Gamma }+i\mathrm{\Omega }^{})]){\displaystyle \frac{\mathrm{\Gamma }_{gn}(\alpha _1\alpha _2)^1}{\mathrm{\Gamma }_{gn}\alpha _1^{}+i(\mathrm{\Omega }_\mu ^{}\alpha _1^{\prime \prime })}}`$
$`\times (\alpha _2|G|^2{\displaystyle \frac{\mathrm{\Delta }n_{mn}/\mathrm{\Delta }n_{gn}}{\mathrm{\Gamma }^2(1+\varkappa )+\mathrm{\Omega }^2}}[(1{\displaystyle \frac{\gamma _{mn}}{\mathrm{\Gamma }_m}}){\displaystyle \frac{2\mathrm{\Gamma }}{\mathrm{\Gamma }_n}}\alpha _2(\mathrm{\Gamma }+i\mathrm{\Omega }^{})])\},`$
where
$$\alpha _{1,2}^{}+\alpha _{1,2}^{\prime \prime }=\frac{1}{2}\left(\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm}+i\mathrm{\Omega }^{}\pm \sqrt{(\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm}+i\mathrm{\Omega }^{})^24|G|^2}\right).$$
(2.4)
It follows from Eqs.(2.3) and (2.4) that it is easiest to achieve level splitting with $`\mathrm{\Gamma }_{gn}=\mathrm{\Gamma }_{gm}`$, and this splitting can be observed most easily in its pure form in the case $`n_m=n_n`$.
The terms proportional to $`1(\gamma _{mn}/\mathrm{\Gamma }_{mn})`$ describe the changes caused in the populations of levels $`m`$ and $`n`$ by the strong field, and the proportional quantities $`\mathrm{\Gamma }+i\mathrm{\Omega }^{}`$ corresponds to nonlinear interference effects . The relative weights of these effects depend on several factors: the relaxation properties of the system, the ratio $`\mathrm{\Delta }n_{mn}/\mathrm{\Delta }n_{gn}`$, and the atomic velocity distribution. The case of a Maxwell velocity distribution was analyzed in . Below we take up the case of monoenergetic atoms (of which a particular case is that of atoms at rest), since all the effects are displayed most clearly with a uniform broadening of spectral lines. Below we will omit the primes from $`\mathrm{\Omega }^{}`$ and $`\mathrm{\Omega }_\mu ^{}`$ and use $`\mathrm{\Omega }`$ and $`\mathrm{\Omega }_\mu `$ to signify the deviation from resonance in the inertial reference system. A real or effective monoenergetic beam can be produced artificially; an effective beam can be produced, for example, by exciting atoms with a coherent field from the ground level to one of the higher-lying levels and by the subsequent relaxation of the atoms to the $`m`$, $`n`$, and $`g`$ levels. Conditions can be arranged such that for the levels of interest the projections of the atomic velocity on the direction of $`𝐤_0`$ will lie in a very narrow velocity range $`\mathrm{\Delta }v=\gamma _0/k_0\mathrm{\Gamma }/k`$, $`\mathrm{\Gamma }_{gn}/k_\mu `$ near the velocity $`v_0=\mathrm{\Omega }_0/k_0`$. Here $`\gamma _0`$ is the line half-width, and $`k_0`$ and $`\mathrm{\Omega }_0`$ are the modulus of the wave vector and the deviation from resonance for the exciting transition. Then we can neglect the atomic velocity distributions at the m, n, and g levels. For a gas with nonuniform broadening these results can be used to show how the individual atoms interact with the field.
## 3 Spectral Properties of the Amplification Coefficient
We first take up the case $`\mathrm{\Delta }n_{mn}=0`$, in which the only effect of the field is to split the levels. Equation (2.3) becomes:
$$\frac{\alpha _\mu }{\alpha _\mu ^0}=Re\left\{\frac{\mathrm{\Gamma }_{gn}}{\alpha _1\alpha _2}\left[\frac{\alpha _1}{\mathrm{\Gamma }_{gn}\alpha _2^{}+i(\mathrm{\Omega }_\mu \alpha _2^{\prime \prime })}\frac{\alpha _2}{\mathrm{\Gamma }_{gn}\alpha _1^{}+i(\mathrm{\Omega }_\mu \alpha _1^{\prime \prime })}\right]\right\}.$$
(3.1)
To determine how the line shape $`\alpha _\mu (\mathrm{\Omega }_\mu )`$ depends on $`|G|^2`$ it is sufficient to analyze the following limiting cases.
For the case $`\mathrm{\Omega }=0`$ the roots $`\alpha _1`$ and $`\alpha _2`$ can be written as:
$`\alpha _1\{\begin{array}{cc}(\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm})\left(1{\displaystyle \frac{|G|^2}{(\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm})^2}}\right),\hfill & 4|G|^2(\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm})^2,\hfill \\ (\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm})/2+i|G|,\hfill & 4|G|^2(\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm})^2;\hfill \end{array}`$
$`\alpha _2\{\begin{array}{cc}|G|^2/(\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm}),\hfill & 4|G|^2(\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm})^2,\hfill \\ (\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm})/2i|G|,\hfill & 4|G|^2(\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm})^2.\hfill \end{array}`$ (3.2)
From Eqs. (3.1) and (3) we see that with $`4|G|^2(\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm})^2`$ the spectrum is a set of two components having half-widths $`\mathrm{\Gamma }_{gn}|G|^2/(\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm})`$ and $`\mathrm{\Gamma }_{gm}+|G|^2/(\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm})`$ The maxima of the two components occur at the same frequency, but the maximum of the component having a half-width $`\mathrm{\Gamma }_{gn}|G|^2/(\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm})`$ is higher than that of the second component by a factor of about $`\mathrm{\Gamma }_{gm}(\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm})^2/\mathrm{\Gamma }_{gn}|G|^2`$. Accordingly, the overall effect of a low-intensity external field in the case $`\mathrm{\Omega }=0`$ is a slight broadening of the spectral line corresponding to the $`gn`$ transition in the case $`\mathrm{\Gamma }_{gm}>\mathrm{\Gamma }_{gn}`$, or there is a narrowing in the case $`\mathrm{\Gamma }_{gm}<\mathrm{\Gamma }_{gn}`$ For a high-intensity external field $`|G|^2(\mathrm{\Gamma }_{gm}\mathrm{\Gamma }_{gn})^2`$, the spectrum consists of two components having the same intensity and the same half-widths $`(\mathrm{\Gamma }_{gm}+\mathrm{\Gamma }_{gn}/2`$. The components are centered at positions symmetric with respect to the frequency $`\mathrm{\Omega }_\mu =0`$ and are separated by $`\mathrm{\Delta }\omega =2|G|`$.
In the other limiting case of $`\mathrm{\Gamma }_{gm}=\mathrm{\Gamma }_{gn}`$, $`\mathrm{\Omega }0`$ we find:
$`\alpha _1i\mathrm{\Omega }(1+|G|^2/\mathrm{\Omega }^2),\alpha _2i|G|^2/\mathrm{\Omega },4|G|^2\mathrm{\Omega }^2;`$
$`\alpha _1i(\mathrm{\Omega }+2|G|)/2,\alpha _2i(\mathrm{\Omega }2|G|)/2,4|G|^2\mathrm{\Omega }^2.`$ (3.3)
In this case the two spectral components have the same half-width, $`\mathrm{\Gamma }_{gn}=\mathrm{\Gamma }_{gm}`$. For weak fields ($`|G|^2\mathrm{\Omega }^2`$) one line has a maximum at $`\mathrm{\Omega }_\mu =|G|^2/\mathrm{\Omega }`$, while the other has a maximum at $`\mathrm{\Omega }_\mu =\mathrm{\Omega }+|G|^2/\mathrm{\Omega }`$; the intensity at the maximum of the first line is higher than that at the maximum of the center of the second by a factor of $`\mathrm{\Omega }^2/|G|^2`$. In intense fields ($`|G|^2\mathrm{\Omega }^2`$) the two lines have the same intensity and lie at symmetric positions with respect to frequency $`\mathrm{\Omega }_\mu =\mathrm{\Omega }/2`$, separated by $`2|G|`$.
We can draw the following conclusion regarding the change in the spectrum accompanying an increase of the intensity of the external field on the basis of these arguments. When the external field is applied, we find, in addition to the fundamental component, having a width of approximately $`2\mathrm{\Gamma }_{gn}`$, an additional component, having a width approximately equal to that of the line corresponding to the Raman transition ($`2\mathrm{\Gamma }_{gm}`$), the center of which is near $`\mathrm{\Omega }_\mu =\mathrm{\Omega }`$. As the external field is intensified, the additional component becomes relatively more important. The width of each line changes in such a manner that in the limit of high external field intensities the widths of both components become the same, equal to $`\mathrm{\Gamma }_{gn}+\mathrm{\Gamma }_{gm}`$. The width change is accompanied by an increase in the separation between the centers of the spectral components; this separation depends on both $`|G|^2`$ and $`\mathrm{\Omega }`$, so that in the limit these components lie symmetrically about the frequency $`\mathrm{\Omega }_\mu =\mathrm{\Omega }/2`$ separated by $`2|G|`$. These spectral changes reflect a modification of the properties of stepped and multiphoton processes caused by an intensification of the external field . The integral radiation intensity, on the other hand, is governed only by the change in the quantity $`n_g\rho _{nn}(\varkappa )`$ caused by the field and is independent of the external field with $`\mathrm{\Delta }n_{mn}=0`$. This behavior is illustrated in Fig.1 where the average values of $`\mathrm{\Omega }`$ and $`\varkappa `$ for the case of a model having the relaxation properties of the neon $`3s_22p_4`$ and $`2s_22p_4`$ transitions.
We turn now to an analysis of the spectral properties of $`\alpha _\mu `$, taking into account the interference term proportional to $`\mathrm{\Gamma }+i\mathrm{\Omega }`$. We restrict the discussion to the case in which the numerator in Eq. (2.1) is governed primarily by this term, i.e., to the case in which we have
$$\frac{\mathrm{\Gamma }}{\mathrm{\Gamma }_{gm}}\left(1\frac{\gamma _{mn}}{\mathrm{\Gamma }_m}\right)\frac{2}{\mathrm{\Gamma }_n},\frac{\mathrm{\Gamma }}{\mathrm{\Gamma }_{gm}}\frac{|G|^2\mathrm{\Delta }n_{mn}/\mathrm{\Delta }n_{gn}}{\mathrm{\Gamma }^2(1+\varkappa )+\mathrm{\Omega }^2}1.$$
In this case amplification is possible even at $`n_g\rho _{nn}(\varkappa )<0`$ and $`\alpha _\mu `$ may change sign as a function of $`𝛀_\mu `$. It follows from Eq. (2.1) that with $`\mathrm{\Delta }n_{mn}/\mathrm{\Delta }n_{gn}>0`$ amplification occurs in the frequency band between $`(\mathrm{\Omega }_\mu )_1`$ and $`(\mathrm{\Omega }_\mu )_2`$, given by
$`(\mathrm{\Omega }_\mu )_{1,2}={\displaystyle \frac{1}{2\mathrm{\Gamma }}}\left\{(\mathrm{\Gamma }_{gn}+\mathrm{\Gamma }+\mathrm{\Gamma }_{gm})\mathrm{\Omega }\pm \sqrt{(\mathrm{\Gamma }_{gn}+\mathrm{\Gamma }+\mathrm{\Gamma }_{gm})^2\mathrm{\Omega }^24\mathrm{\Gamma }_{gn}\mathrm{\Gamma }\mathrm{\Omega }^2+4\mathrm{\Gamma }^2(\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm}+|G|^2)}\right\}.`$ (3.4)
It follows from this equation that the band width increases with increasing $`\mathrm{\Omega }^2`$ and $`|G|^2`$. In the limiting cases in which the intense field is far from and close to resonance, we find from Eq. (3.4)
$`(\mathrm{\Omega }_\mu )_{1,2}{\displaystyle \frac{1}{2}}\left(1+{\displaystyle \frac{\mathrm{\Gamma }_{gn}}{\mathrm{\Gamma }}}\right)\mathrm{\Omega }\pm \sqrt{\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm}+|G|^2},\text{if}\mathrm{\Gamma }_{gm}\mathrm{\Gamma },\mathrm{\Gamma }_{gn};\mathrm{\Omega }^2{\displaystyle \frac{4\mathrm{\Gamma }^2(\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm}+|G|^2)}{(\mathrm{\Gamma }_{gn}\mathrm{\Gamma })^2}},`$
$`(\mathrm{\Omega }_\mu )_1{\displaystyle \frac{\mathrm{\Gamma }_{gn}}{\mathrm{\Gamma }}}\mathrm{\Omega },(\mathrm{\Omega }_\mu )_2\mathrm{\Omega },\text{if}\mathrm{\Gamma }_{gm}\mathrm{\Gamma },\mathrm{\Gamma }_{gn};\mathrm{\Omega }^2{\displaystyle \frac{4\mathrm{\Gamma }^2(\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm}+|G|^2)}{(\mathrm{\Gamma }_{gn}\mathrm{\Gamma })^2}}.`$ (3.5)
With $`\mathrm{\Omega }=0`$ we have $`(\mathrm{\Omega }_\mu )_{1,2}=\pm \sqrt{\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm}+|G|^2}`$ for any values of $`\mathrm{\Gamma }_{gn}`$, $`\mathrm{\Gamma }_{gm}`$ or $`|G|^2`$. Here the half-width at half-height of the amplification line is
$`(\mathrm{\Delta }\mathrm{\Omega }_\mu )_{1,2}^2=`$ $`{\displaystyle \frac{1}{2}}\left\{\sqrt{(\mathrm{\Gamma }_{gn}+\mathrm{\Gamma }_{gm})^4+4(\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm}+|G|^2)^2}(\mathrm{\Gamma }_{gn}+\mathrm{\Gamma }_{gm})^2\right\},`$
$`(\mathrm{\Delta }\mathrm{\Omega }_\mu )_{1,2}^2`$ $`{\displaystyle \frac{\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm}+|G|^2}{\mathrm{\Gamma }_{gn}+\mathrm{\Gamma }_{gm}}},{\displaystyle \frac{\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm}+|G|^2}{(\mathrm{\Gamma }_{gn}+\mathrm{\Gamma }_{gm})^2}}1.`$ (3.6)
It follows from Eqs. (3) and (3) that with $`|G|^2\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm}`$ and $`\mathrm{\Gamma }_{gm}\mathrm{\Gamma }_{gn}`$ the width of the amplification band is governed by the geometric average of $`\mathrm{\Gamma }_{gn}`$ and $`\mathrm{\Gamma }_{gm}`$, while the width of the amplification line is $`2\mathrm{\Gamma }_{gm}`$ and may be much narrower than the natural line width corresponding to the $`gn`$ transition. When the frequency of the strong field is scanned, the amplification band of the weak field also shifts; the band width depends on both $`\mathrm{\Omega }^2`$ and $`|G|^2`$.
We see from Eq. (2.1) that as field $`E`$ increases there are increases in the population and interference contributions to $`\alpha _\mu `$. On the other hand, there is a tendency for the amplification coefficient at the center the line to fall off with increasing $`|G|^2`$ because of the level splitting. Analysis of Eq. (2.1) for $`\mathrm{\Omega }_\mu =\mathrm{\Omega }=0`$ shows that the optimum value at $`\varkappa `$, corresponding to the maximum value of $`\alpha _\mu `$ at the line center, is given by $`\varkappa _{opt}=\varkappa _1(x)\{1+\sqrt{1+[x_1\varkappa _1(x)]^1(2\tau ^2\mathrm{\Gamma }_{gm}\mathrm{\Gamma }_{gn}x+x_1)}\},x>x_1,`$ where $`x=\mathrm{\Delta }n_{mn}/\mathrm{\Delta }n_{gn}`$, $`x_1=[(\mathrm{\Gamma }_m\gamma _{mn}+\mathrm{\Gamma }_n\mathrm{\Gamma }_m)(2\mathrm{\Gamma }_{gm})^1]^1(\mathrm{\Gamma }_m\gamma _{mn}+\mathrm{\Gamma }_n)`$, $`\varkappa _1(x)=(xx_1^11)^1`$.
The spectral properties of the function $`\alpha _\mu (\mathrm{\Omega }_\mu )/\alpha _\mu ^0`$ for the model discussed above are illustrated in Figs. 2-4, where $`x`$ is set equal to 4.14, corresponding to the optimum field $`\varkappa =2`$. Figure 2 corresponds to the case $`\mathrm{\Omega }=0`$ for values $`\varkappa >\varkappa _{opt}`$. The change in $`\alpha _\mu `$ at the maximum is very rapid while $`\varkappa `$ varies near the optimum. (These cases are not illustrated in Fig. 2.) For example, with $`\varkappa =\varkappa _{opt}=2`$ we have $`\alpha _\mu (0)/\alpha _\mu ^0=32`$. As $`\varkappa `$ falls off to half its value, $`\alpha _\mu (0)/\alpha _\mu ^0`$ falls off to about one-third its value. As $`\varkappa `$ increases to a value 50% above $`\varkappa _{opt}`$, the value of $`\alpha _\mu (0)/\alpha _\mu ^0`$ falls off by a factor of 60. As $`\varkappa `$ changes from 2 to 3, the line half-width at half-height changes from $`0.2\mathrm{\Gamma }_{gn}`$ to $`\mathrm{\Gamma }_{gn}`$. With $`\varkappa =\varkappa _{opt}`$ the $`\alpha _\mu `$ profile is symmetric and changes sign at $`\mathrm{\Omega }_\mu ^2\mathrm{\Gamma }_{gn}^2`$. The absorption in the wings changes very slowly with increasing $`|\mathrm{\Omega }_\mu |`$. The maximum absorption is roughly $`1/30`$ the value of $`|\alpha _\mu (0)|`$.
Figure 2 shows the change in $`\alpha _\mu (\mathrm{\Omega }_\mu )`$ corresponding to a further increase in $`\varkappa `$. The case $`\varkappa =3`$ corresponds to the vanishing of $`n_g\rho _{nn}(\varkappa )`$. The integral value of the coefficient $`\alpha _\mu `$ corresponding to $`gn`$ transition also vanishes. As the external field is intensified, the line splits, so that with $`\varkappa =40`$ a plateau appears on the amplification curve, about $`6\mathrm{\Gamma }_{gn}`$ in width.
Figures 3 and 4 show the changes in the spectral properties of the amplification coefficient as the frequency and intensity of the strong field are changed. Figure 3 shoves the line profile for $`\mathrm{\Omega }=\mathrm{\Gamma }`$, while figure 4 shows this profile for $`\mathrm{\Omega }=10\mathrm{\Gamma }`$. These cases correspond to quasiresonant Raman scattering through a common lower level. We see from Figs. 3 and 4 that the frequency separation between the amplification and absorption maxima increases as the deviation of the strong field from resonance increases. For fixed $`\mathrm{\Delta }n_{gn}<0`$, an increase in $`|\mathrm{\Omega }|`$ requires a more intense external field for appreciable amplification. This effect is accompanied by a change in $`n_g\rho _{nn}(\varkappa )`$ which is significant in comparison with $`\mathrm{\Delta }n_{gn}`$.
## 4 Conclusion
In conclusion we will compare the results for the cases of a beam having a Maxwell atomic velocity distribution and a monoenergetic beam. With a Maxwell velocity distribution, peaks or troughs appear against the background of the $`\alpha _\mu (\mathrm{\Omega }_\mu )`$ Doppler profile under the influence of the traveling wave of the intense field; these peaks and dips are described by
$`{\displaystyle \frac{k_\mu }{k}}{\displaystyle \frac{1\pm \sqrt{1+\varkappa }}{\sqrt{1+\varkappa }}}Re\{(N_mN_n)|G|^2[\mathrm{\Gamma }_0+i(\mathrm{\Omega }_\mu {\displaystyle \frac{k_\mu }{k}}\mathrm{\Omega })+`$
$`+{\displaystyle \frac{|G|^2}{\mathrm{\Gamma }_\pm +i\left(\mathrm{\Omega }_\mu {\displaystyle \frac{k_\mu }{k}}\mathrm{\Omega }\right)}}]^1[(1{\displaystyle \frac{\gamma _{mn}}{\mathrm{\Gamma }_m}}){\displaystyle \frac{2}{\stackrel{~}{\mathrm{\Gamma }}_n}}+{\displaystyle \frac{1}{\mathrm{\Gamma }_\pm +i\left(\mathrm{\Omega }_\mu {\displaystyle \frac{k_\mu }{k}}\mathrm{\Omega }\right)}}]\},`$ (4.1)
where
$`k_\mu >k,\mathrm{\Gamma }_0=\mathrm{\Gamma }_{gn}+{\displaystyle \frac{k_\mu }{k}}\mathrm{\Gamma }\sqrt{1+\varkappa },\mathrm{\Gamma }_\pm =\mathrm{\Gamma }_{gm}+\left(1{\displaystyle \frac{k_\mu }{k}}\right)\mathrm{\Gamma }\sqrt{1+\varkappa },\stackrel{~}{\mathrm{\Gamma }}_n=\mathrm{\Gamma }_n(1\pm \sqrt{1+\varkappa }).`$ (4.2)
The upper sign corresponds to the case $`k_\mu k>0`$, and the lower sign corresponds to the case $`k_\mu k<0`$ . Comparing Eq. (4) with Eq. (2.1), we conclude that the line profile for the peaks (or dips) for a gas with a Doppler velocity distribution is the same as the line shape for an effective monoenergetic (atomic) beam, for which we have $`|N_mN_n||N_gN_n|,\mathrm{\Omega }^{}=0,\mathrm{\Gamma }_{gn}=\mathrm{\Gamma }_0,\mathrm{\Gamma }_n=\stackrel{~}{\mathrm{\Gamma }}_n,\mathrm{\Gamma }_{gm}=\mathrm{\Gamma }_\pm `$ and for which the resonant frequency is $`\omega _{gn}\pm (k_\mu /k)\mathrm{\Omega }`$. However, in the case of a Maxwell distribution the line shape has a completely different dependence on the magnitude of the external, field, and there is no profile-asymmetry effect, as may be displayed in the case of a monoenergetic beam.
The authors thank S. G. Rautian for useful discussions and S. I. Mortseva for assistance in the numerical calculations.
Novosibirk State University. Institute of Semiconductor Physics, Siberian Branch, Academy of Sciences of the USSR.
## LITERATURE CITED
1. V. M. Fain, Ya. I. Khanin, and E. G. Yashchin, lzv. VUZ. Radiofiz., 5, 697 (1962); V. M. Fain and Ya. I. Khanin, Quantum Radio Physics \[in Russian\], lzd. Sov. Radio (1965).
2. G. E. Notkin, S. G. Rautian, and A. A. Feoktistov, Zh. Eksp. Teor. Fiz., 52, 1673 (1967).
3. T. Ya. Popova, A. K. Popov, S. G. Rautian, and R. I. Sokolovskii, Zh. Eksp. Teor. Fiz., 57, 85 (1969).
4. H.K. Holt, Phys. Rev. Lett D 19, 1275 (1967); ibid 20, 410(1968).
5. T. Ya. Popova, A. K. Popov, S. G. Rautian, and A. A. Feoktistov, Zh. Eksp. Teor. Fiz., 57, 444 (1969).
6. A. K. Popov, Zh. Fksp. Teor. Fiz., 58, 1623 (1970).
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# Parametrization and distillability of three-qubit entanglement
## I Introduction
The importance of quantum entanglement, both as a resource for quantum information processing and as a ubiquitous feature of quantum systems, has become increasingly apparent over the last few years . Recent developments in quantum information theory, in particular, have stimulated interest in the quantification and manipulation of entanglement.
For bipartite pure states an essentially complete theory of entanglement now exists , though the situation for mixed states is less definite . All descriptions of bipartite pure state entanglement start with the Schmidt decomposition. It is possible to find orthonormal bases $`\{|i_A\}`$ and $`\{|i_B\}`$ for systems A and B such that we can write the joint state of the system in the form
$$|\mathrm{\Psi }_{AB}=\underset{i}{}\sqrt{p_i}|i_A|i_B,p_i>0,\underset{i}{}p_i=1.$$
(1)
These Schmidt coefficients $`\{p_i\}`$ are uniquely defined by the state $`|\mathrm{\Psi }_{AB}`$, and are equal to the eigenvalues of the reduced density matrix $`\rho _A`$ (or equivalently, of $`\rho _B`$); the bases $`\{|i_A\}`$ and $`\{|i_B\}`$ are eigenbases of $`\rho _A`$ and $`\rho _B`$, respectively, so the local density matrices are diagonal in this choice of bases. This choice of bases also minimizes the number of terms needed to represent $`|\mathrm{\Psi }_{AB}`$.
For tripartite or multipartite states, there is no equivalent to the Schmidt decomposition (1) ). Three main approaches to parametrizing tripartite or multipartite entanglement have been followed so far. First, one may choose the local bases to put the joint state into a standard form. Often these standard forms are intended to generalize some aspect of the Schmidt decomposition in the bipartite case . Second, one may try to identify a complete set of locally invariant quantities, functions of the state which are invariant under local unitary transformations , and which uniquely characterize equivalent states. The coefficients of a standard form are obviously such quantities, but they may not have readily meaningful physical interpretations. Third, one may identify operational quantities, such as the number of Greenberger-Horne-Zeilinger (GHZ) triplets or EPR pairs that can be distilled from the state by some procedure .
In section II we consider a number of proposals for standard forms of three-qubit pure states, concentrating especially on those which generalize some aspect of the bipartite Schmidt decomposition: the minimal form , the two-term form , and the Schmidt form. This form was given briefly in and independently in . In section III we examine it in greater detail. We give an explicit parametrization of the coefficients in terms of five locally invariant quantities, and discuss their physical significance.
We make use of the Schmidt form to prove analytically the reliability of a proposed distillation technique for GHZ triplets from general three-qubit pure states ; we present this proof in section IV.
In Linden and Popescu proposed characterizing the entanglement properties of three-qubit states by examining the “orbits” of the states under general local unitary transformation. This was carried a step further by Carteret and Sudbery , who proved that most states have a certain generic behavior under such transformations, but identified classes of ‘special’ states which they speculated to have unusual entanglement properties. In section V we briefly review these ‘special’ states, then analytically evaluate the yield of GHZs under the distillation protocol of and section IV. By examining the yield of these states as a function of the invariant parameters from section III, and also of the locally invariant “residual tangle” $`\tau _{ABC}`$ of Coffman, Kundu and Wootters , we verify that these classes of states are indeed exceptional by this operational measure, representing extremes of distillability or undistillability. We briefly compare the results using this protocol to the recently discovered optimal distillation method of , and find that they are entirely consistent. Our conclusions are summarized in section VI.
## II Review of standard forms for three-qubit states
Two qubits can always be represented in their Schmidt decomposition (1)
$$|\psi =\sqrt{p}|00+\sqrt{1p}|11,$$
(2)
characterized by a single Schmidt coefficient $`p`$ (or equivalently $`1p`$). Without loss of generality, we adopt the convention that $`p1/2`$ and the corresponding eigenvector is $`|0`$.
Most attempts to define a standard form for a three-qubit state attempt either to generalize some property of the bipartite Schmidt decomposition, or make use of the Schmidt decomposition between one of the bits and the other two, or both. For instance, we can make a Schmidt decomposition between qubit A and qubits B and C, writing the three-qubit state in the form
$$|\psi =\sqrt{p}|0_A|\psi _0_{BC}+\sqrt{1p}|1_A|\psi _1_{BC}.$$
(3)
Choosing the Schmidt basis for qubit A guarantees that the correlated states of qubits B and C must be orthogonal: $`\psi _0|\psi _1=0.`$
The sixteen real parameters to describe a generic pure state of three qubits can be reduced to fifteen by normalization, and to five which are invariant under the ten-dimensional group of local unitary transformations ; unfortunately, no single choice of five quantities has proven completely satisfactory. One simple parametrization that has been proposed is the Linden-Popescu-Schlienz (LPS) standard form. One begins with a state in form (3). One can then choose one of the two correlated states, say $`|\psi _0_{BC}`$, and find its corresponding Schmidt bases. The resulting state for the three qubits has the form
$`|\psi `$ $`=`$ $`\sqrt{p}|0_A\left(a|00_{BC}+\sqrt{1a^2}|11_{BC}\right)`$ (5)
$`+\sqrt{1p}|1_A\left(\gamma (\sqrt{1a^2}|00_{BC}a|11_{BC})+f|01_{BC}+g|10_{BC}\right),`$
where $`p`$, $`a`$ and $`f`$ are real positive numbers, $`g`$ is complex, and $`\gamma =(1f^2|g|^2)^{1/2}`$. Together these give five independent real parameters.
The vectors $`|\psi _0_{23}`$ and $`|\psi _1_{23}`$ span a two-dimensional subspace of the Hilbert space for qubits 2 and 3. It’s possible to make an interesting variation on the LPS idea using a result of Niu and Griffiths, who showed that any such two-dimensional subspace can be given basis vectors of the form
$`|\chi _0`$ $`=`$ $`\sqrt{q}|00_{23}+\sqrt{1q}|11_{23},`$ (6)
$`|\chi _1`$ $`=`$ $`\sqrt{r}|01_{23}+\sqrt{1r}|10_{23},`$ (7)
for some choice of a product basis for the 4-D Hilbert space of the two bits, where $`q`$ and $`r`$ are real numbers between 0 and 1. Using this basis leads to a unique standard form
$`|\psi `$ $`=`$ $`\sqrt{p}|0\left(a\sqrt{q}|00+a\sqrt{1q}|11+b\sqrt{r}|01+b\sqrt{1r}|10\right)`$ (9)
$`+\sqrt{1p}|1\left(b^{}\sqrt{q}|00b^{}\sqrt{1q}|11+a\sqrt{r}|01+a\sqrt{1r}|10\right),`$
where $`a`$ is real and $`a^2+|b|^2=1`$. This then gives five independent real parameters: $`p`$, $`q`$, $`r`$, $`a`$, and the phase of $`b`$. This form treats the $`|0`$ and $`|1`$ terms more symmetrically than LPS; however, there is still a lack of symmetry under interchange of the bits.
More interesting from a fundamental point of view are attempts to generalize some aspect of the Schmidt decomposition. Three such properties suggest themselves. First, the Schmidt decomposition is the choice of orthonormal bases for the local Hilbert spaces which minimizes the number of terms needed to represent the state. Second, any two qubit state can be written as the sum of only two product vectors. (For $`N`$-dimensional systems, $`N`$ product vectors are needed.) Third, the Schmidt decomposition diagonalizes the reduced density matrices of the local subsystems. No single representation for tripartite systems has all three properties, but they can be generalized individually.
Acín et al. have shown that all three-qubit states can be written in the form
$$|\psi =\lambda _0|000+\lambda _1\mathrm{e}^{i\varphi }|100+\lambda _2|101+\lambda _3|110+\lambda _4|111$$
(10)
by a suitable choice of basis, where the $`\lambda _i`$ are all real and positive and $`\varphi `$ is a phase between $`0`$ and $`\pi `$. With only five terms, this is a minimal description, and in that sense a generalization of the bipartite Schmidt decomposition. A similar form has been described by Higuchi and Sudbery ,
$$|\psi =\lambda _0\mathrm{e}^{i\varphi }|000+\lambda _1|100+\lambda _2|010+\lambda _3|001+\lambda _4|111$$
(11)
which has the added benefit of being symmetric under interchange of the qubits. Carteret, Higuchi and Sudbery have shown how to generalize this construction to give a unique minimal representation for systems of any dimension. These minimal forms have practical benefits: with a small number of terms, they can simplify the calculation of locally invariant quantities . However, the $`\lambda _i`$ and $`\varphi `$ themselves have no obvious physical interpretation.
This minimal property can be generalized in another way, by relaxing the requirement that the product vectors be orthogonal. Acín et al. and Dür, Vidal and Cirac have also shown that almost all three qubit states can be written in the form
$$|\mathrm{\Psi }_{ABC}=\mu _1|a_1b_1c_1+\mu _2\mathrm{e}^{i\varphi }|a_2b_2c_2,$$
(12)
where the vectors are normalized but not orthogonal. There are six real parameters, $`\mu _1`$, $`\mu _2`$, $`a_1|a_2`$, $`b_1|b_2`$, $`c_1|c_2`$ and $`\varphi `$; imposing normalization reduces this to five.
Interestingly, not all three qubit states can be written in the form (12); a small subclass of states require a minimum of three product terms . Dür, Vidal and Cirac made use of this result to prove that there are two classes of three-qubit pure states which cannot be interconverted with nonzero probability . The class that requires three terms is a three-parameter family, and is characterized by vanishing residual tangle $`\tau _{ABC}=0`$ (see section V). Acín, Dür and Vidal also used this form to demonstrate a method of converting a single copy of a three qubit state into a GHZ triplet with maximum probability .
The third generalization is to find bases for all three qubits which diagonalize their reduced density matrices. That is, one can simultaneously put each bit in its Schmidt decomposition with respect to the other two. This form was proposed in and independently in . The state has the form
$`|\psi `$ $`=`$ $`a|000+b|001+c|010+d|011`$ (14)
$`+e|100+f|101+g|110+h|111,`$
which looks just like a generic three-qubit state with 16 parameters. However, using each of the three qubits in turn we can write $`|\psi `$ in a form similar to (3), with orthogonality conditions which impose restrictions on the possible values of the coefficients in (14). We can use these relationships to reduce these coefficients to five independent parameters, as we show in the next section.
## III Parametrizing the Schmidt form
By redefining the relative phases of the basis vectors
$$|0_j,|1_j\mathrm{exp}(i\varphi _j)|0,\mathrm{exp}(i\theta _j)|1,$$
(15)
we can choose to make four of the coefficients real. A convenient choice is to make $`a,d,f,g`$ real, while $`b,c,e,h`$ remain complex. The state must also be normalized, which imposes the condition
$$a^2+|b|^2+|c|^2+d^2+|e^2|+f^2+g^2+|h|^2=1.$$
(16)
This leaves 11 undetermined parameters.
We can now express the larger eigenvalues $`p_{A,B,C}`$ of the reduced density matrices $`\rho _{A,B,C}`$ in terms of the coefficients:
$$p_A=a^2+|b|^2+|c|^2+d^2,\mathrm{etc}.,$$
(17)
(the smaller eigenvalues obviously being $`1p_{A,B,C}`$). Finally, the states $`|\psi _{0,1}_{kl}`$ correlated with basis vectors $`|0_j`$ and $`|1_j`$ must be orthogonal to each other. This gives three more equations:
$$ae^{}+bf+cg+dh^{}=0,\mathrm{etc}.$$
(18)
Because these equations are complex, they are equivalent to six real equations.
Combining these restrictions, we now have fourteen equations in sixteen unknowns. Thus, in addition to the eigenvalues $`p_{A,B,C}`$ we would expect there to be two more free parameters. Can we identify reasonable candidates for these parameters? It turns out that natural choices are the two probabilities $`a^2`$ and $`|h|^2`$. These parameters are symmetric under interchanges of the three qubits, and have a fairly simple physical interpretation: they are the probabilities of all three qubits giving the same result (0 or 1, respectively) when measured in their Schmidt bases. Moreover, the coefficients of the other state vectors can all be calculated in terms of the five probabilities $`a^2,|h|^2`$, and $`p_{A,B,C}`$, up to a sign.
Define $`p_{\mathrm{sum}}=p_A+p_B+p_C`$. The expressions for the norms of the coefficients are then simple:
$`|b|^2,|c|^2,|e|^2`$ $`=`$ $`{\displaystyle \frac{(2p_{C,B,A}1)|h|^2(p_{\mathrm{sum}}p_{C,B,A}1)(2a^2p_{\mathrm{sum}}+1)}{2p_{\mathrm{sum}}3}}`$ (19)
$`d^2,f^2,g^2`$ $`=`$ $`{\displaystyle \frac{(2p_{A,B,C}1)a^2(p_{\mathrm{sum}}p_{A,B,C}1)(2|h|^2+p_{\mathrm{sum}}2)}{2p_{\mathrm{sum}}3}}.`$ (20)
The phases of $`b,c,e`$ are more complicated. If we define the variables $`\varphi _{b,c,e}`$ by $`b=|b|\mathrm{exp}(i\varphi _b)`$, $`c=|c|\mathrm{exp}(i\varphi _c)`$, and $`e=|e|\mathrm{exp}(i\varphi _e)`$, the constraint equations (18) imply after a bit of algebra that
$`\mathrm{cos}(\varphi _b),\mathrm{sin}(\varphi _b)`$ $`=`$ $`(Q_{1,2}/|b|)(2adf+g(a^2+d^2+f^2g^2)),`$ (21)
$`\mathrm{cos}(\varphi _c),\mathrm{sin}(\varphi _c)`$ $`=`$ $`(Q_{1,2}/|c|)(2adg+f(a^2+d^2f^2+g^2)),`$ (22)
$`\mathrm{cos}(\varphi _e),\mathrm{sin}(\varphi _e)`$ $`=`$ $`(Q_{1,2}/|e|)(2afg+d(a^2d^2+f^2+g^2)),`$ (23)
$`\mathrm{cos}(\varphi _h),\mathrm{sin}(\varphi _h)`$ $`=`$ $`(Q_{1,2}/|h|)(2dfg+a(a^2+d^2+f^2+g^2)),`$ (24)
where $`Q_1`$ and $`Q_2`$ are two constants. We can solve for the values of $`Q_1`$ and $`Q_2`$ by using the identity $`\mathrm{sin}^2(\varphi )+\mathrm{cos}^2(\varphi )=1`$ and substituting (20) for $`|b|,\mathrm{},g`$.
In the Schmidt form for three-qubit pure states, each of the five parameters has a reasonably straightforward physical interpretation. The three parameters $`p_A,p_B,p_C`$ are the larger (i.e., $`p1/2`$) eigenvalues of the reduced density operators for each of the three qubits, and correspond to the probabilities of obtaining the more likely of the two possible outcomes (which by convention we label $`|0`$) when we measure each of the qubits in its Schmidt basis. These parameters are closely related to the minimum absolutely selective information for each qubit, which is given by the entropy function
$$\mathrm{min}S_i=(p_i\mathrm{log}_2p_i+(1p_i)\mathrm{log}_2(1p_i)).$$
(25)
This quantity is the minimum amount of fundamentally unpredictable classical information generated by carrying out a measurement on qubit $`i`$, given a free choice of measurement basis . By using the Schmidt form to choose measurement bases we can simultaneously minimize the absolutely selective information for all three qubits.
The parameters $`p_A,p_B,p_C`$ range from $`1/2`$ to $`1`$ (since they are defined to be the larger eigenvalues of their corresponding local density matrices). Similarly, $`a^2`$ ranges from $`0`$ to $`1`$, and $`|h|^2`$ from $`0`$ to $`1/2`$. However, this does not mean that these parameters can take arbitrary values within these ranges. Some choices of parameter values correspond to no physical state, and give nonsensical values for (20) and (24).
In particular, the local probabilities must obey the triangle inequalities
$`p_A(1p_A)+p_B(1p_B)`$ $``$ $`p_C(1p_C),`$ (26)
$`p_B(1p_B)+p_C(1p_C)`$ $``$ $`p_A(1p_A),`$ (27)
$`p_C(1p_C)+p_A(1p_A)`$ $``$ $`p_B(1p_B);`$ (28)
these imply, for instance, that if $`p_A=1`$ then $`p_B=p_C`$. The restrictions on $`a^2`$ and $`|h|^2`$ are more complicated, but they too display an interdependency in their range of values. In particular, as $`p_{\mathrm{sum}}3`$ we must have $`a^21`$ and $`|h|^20`$.
## IV Proof of distillability
The Schmidt form can provide analytical insight when addressing specific problems. For example, the efficacy of a recently proposed tripartite distillation protocol can be demonstrated with its help.
Consider a state of three qubits in an arbitrary product basis, which can be written in the form (14). We can straightforwardly calculate the quantity $`p_A(1p_A)`$
$`p_A(1p_A)`$ $`=`$ $`|afbe|^2+|agce|^2+|ahde|^2`$ (30)
$`+|bgcf|^2+|bhdf|^2+|chdg|^2,`$
This expression is a polynomial in the coefficients and their complex conjugates, and is correct in any basis. If the state is in the Schmidt form, this simplifies to
$$p_A(1p_A)=(a^2+|b|^2+|c|^2+d^2)(|e|^2+f^2+g^2+|h|^2).$$
(31)
Let us assume that we have written the state in Schmidt form, such that the states $`\{|0,|1\}`$ for each qubit $`j`$ are eigenstates of the local density matrix with eigenvalues $`p_j`$ and $`1p_j`$, respectively. Suppose we now perform a weak measurement on each of the three qubits. First, allow each qubit to interact with a separate ancilla bit initially in state $`|0`$, such that
$`|0|0_{\mathrm{anc}}`$ $``$ $`\sqrt{1ϵ}|0|0_{\mathrm{anc}}+\sqrt{ϵ}|0|1_{\mathrm{anc}},`$ (32)
$`|1|0_{\mathrm{anc}}`$ $``$ $`|1|0_{\mathrm{anc}},`$ (33)
where $`ϵ1`$. Then measure the three ancilla bits. With a probability of $`ϵ(p_{\mathrm{sum}})`$ one will find one or more of the ancilla bits in state $`|1_{\mathrm{anc}}`$, in which case the procedure has failed. Otherwise, this step has succeeded and the three qubits are now in a new state with slightly different coefficients $`a^{},b^{},\mathrm{},h^{}`$. The changes in the coefficients are
$`\mathrm{\Delta }a`$ $`=`$ $`(ϵ/2)(3p_{\mathrm{sum}})a,`$ (34)
$`\mathrm{\Delta }(b,c,e)`$ $`=`$ $`(ϵ/2)(2p_{\mathrm{sum}})(b,c,e),`$ (35)
$`\mathrm{\Delta }(d,f,g)`$ $`=`$ $`(ϵ/2)(1p_{\mathrm{sum}})(d,f,g),`$ (36)
$`\mathrm{\Delta }h`$ $`=`$ $`(ϵ/2)p_{\mathrm{sum}}h.`$ (37)
This very simple form results because the state is in Schmidt form. After this procedure the bases for the three bits will generally no longer be the correct Schmidt basis (though it will be close to it), so the expression (31) cannot be used; but (30) is always correct. Thus we get a change in $`p_A(1p_A)`$
$`\mathrm{\Delta }[p_A(1p_A)]`$ $`=`$ $`ϵ(42(p_A+p_B+p_C))(|afbe|^2+|agce|^2)`$ (40)
$`ϵ(32(p_A+p_B+p_C))(|ahde|^2+|bgcf|^2)`$
$`ϵ(22(p_A+p_B+p_C))(|bhdf|^2+|chdg|^2)`$
$`=`$ $`ϵ(32(p_A+p_B+p_C))p_A(1p_A)`$ (42)
$`(ϵ/2)(|afbe|^2+|agce|^2|bhdf|^2|chdg|^2)`$
By making use of equations (17) and (18), this expression simplifies to
$`\mathrm{\Delta }[p_A(1p_A)]`$ $`=`$ $`ϵ[(2(p_A+p_B+p_C)3)p_A(1p_A)`$ (44)
$`+p_A(a^2|e|^2+|h|^2d^2)+d^2a^2],`$
which using (20) further simplifies to
$`\mathrm{\Delta }[p_A(1p_A)]`$ $`=`$ $`{\displaystyle \frac{ϵ(2p_A1)}{2p_A+2p_B+2p_C3}}[2(a^2+|h^2|)(p_B+p_C1)`$ (46)
$`(2p_A1)(p_A+p_B+p_C1)(p_A+p_B+p_C2)].`$
The prefactor to (46) is strictly positive, as is the first term inside the brackets. The second term is positive if $`p_A+p_B+p_C<2`$; any state that satisfies this criterion will evolve towards the GHZ state and have a nonzero yield.
For $`p_A+p_B+p_C2`$, the sign of (46) depends on the relative sizes of the first and second terms inside the brackets. The last two equations of (20) show that for $`p_A+p_B+p_C2`$, the fact that $`f^2+g^2>0`$ implies
$$2a^2(p_B+p_C1)(2p_A+p_B+p_C2)(p_A+p_B+p_C2),$$
(47)
which yields the inequalities
$`2a^2(p_B+p_C1)(2p_A1)(p_A+p_B+p_C1)(p_A+p_B+p_C2)`$ (48)
$``$ $`(2p_A+p_B+p_C2)(p_A+p_B+p_C2)`$ (50)
$`(2p_A1)(p_A+p_B+p_C1)(p_A+p_B+p_C2)`$
$`=`$ $`(1p_A)(p_A+p_B+p_C2)(2p_A+2p_B+2p_C3)0.`$ (51)
This straightforwardly implies
$$\mathrm{\Delta }[p_A(1p_A)]ϵ(2p_A1)(1p_A)(p_A+p_B+p_C2)0.$$
(52)
Because of the symmetry of the protocol, $`p_B(1p_B)`$ and $`p_C(1p_C)`$ must also increase. So one step of this protocol must move the state towards the GHZ with nonvanishing probability, and will (in general) produce a nonzero yield of GHZ triplets.
There are three circumstances in which this result can fail. First, no product state can ever be distilled to a GHZ by this method. At least one of $`p_A,p_B,p_C`$ must equal 1 in this case, which causes the rate (46) corresponding to it to vanish. This is not immediately obvious from the form of (46), but it is easily checked using (20) and (28)—if $`p_A=1`$, then $`p_B=p_C=a^2`$, and (46) is equal to zero.
Second, there are states with $`p_A+p_B+p_C=2`$ for which $`a^2=|h|^2=0`$, again making (46) vanish. These are a subset of the triple states discussed in section V below, which are equivalent to states of the form (55); these states have vanishing residual tangle. Finally, it is possible for a state with $`p_A+p_B+p_C>2`$ to evolve to one of these triple states. All such states will also have vanishing residual tangle , and conversely all states with vanishing residual tangle will evolve under this distillation protocol to a triple state with $`p_A+p_B+p_C=2`$, and hence have zero yield of GHZs. This can be clearly seen in Fig. 1.
## V Entanglement and distillability
Linden and Popescu proposed characterizing three-qubit states by the dimensions of their orbits under the action of the local unitary group. Generically, tripartite pure states of qubits have ten-dimensional orbits, equal to the dimension of the local unitary group. The very interesting results of Carteret and Sudbery give a complete classification of all states for three qubits which behave nongenerically under local unitary transformations; these ‘special’ states have stabilizers of nonzero dimension, and hence orbits of dimension $`<10`$ (see ). This behavior suggests that these ‘special’ classes have unusual entanglement properties, which might be evident in other measures of entanglement.
We have numerically simulated the distillation of generic states by the algorithm described above in section IV, in order to determine the yield of GHZ triplets as a function of various parameters, especially the parameters used to describe the Schmidt form. We have also calculated analytical expressions for the yield of states in the exceptional classes enumerated by Carteret and Sudbery. We find that these states are indeed exceptional by this operational criterion, as we describe below.
Most important in calculating the yield of GHZs is the sum of the local eigenvalues $`p_{\mathrm{sum}}p_A+p_B+p_C`$. This quantity determines the probability of failure in one step of the infinitesimal distillation procedure of section IV, with the probability of failure being $`ϵp_{\mathrm{sum}}`$. If it takes $`N`$ steps to become sufficiently close to a GHZ triplet, the expected yield is
$$Y=\underset{n=1}{\overset{N}{}}(1ϵp_{\mathrm{sum}}(n))\mathrm{exp}\left\{\underset{n}{}ϵp_{\mathrm{sum}}(n)\right\},$$
(53)
where $`p_{\mathrm{sum}}(n)`$ is the value of $`p_{\mathrm{sum}}`$ at the $`n`$th step. In the limit of infinitesimal steps the sum inside the exponent becomes an integral. Calculating $`p_{\mathrm{sum}}(n)`$ analytically is no simple matter for a general state; the equations (37) for the change in the coefficients become differential equations in the limit, but must be supplemented by an additional change of basis between steps, since in general the bases will no longer be the Schmidt basis for the new state. While this is simple to do numerically, analytically it is challenging.
Fortunately, the classes of exceptional states are generally expressible in simple forms which make it possible to integrate the equations (37) in closed form, and derive simple expressions for the yield of GHZs. Interestingly, the steps of the GHZ distillation technique commute with local unitary transformations. Because of this, the distillation procedure preserves the stabilizer of the initial state, and hence must take ‘special’ states to other ‘special’ states of the same type. This gives another way of understanding why certain special states are not distillable to GHZs.
In addition to $`a^2`$, $`|h|^2`$, and $`p_{\mathrm{sum}}`$, we looked at the dependence of the GHZ yield on one other locally invariant quantity. This is the residual tangle of Coffman et al. , which can be written
$$\tau _{ABC}=2(\lambda _1^{AB}\lambda _2^{AB}+\lambda _1^{AC}\lambda _2^{AC}),$$
(54)
where $`\lambda _1^{ij}`$ and $`\lambda _2^{ij}`$ are the (positive) eigenvalues of the matrix $`\sqrt{\rho _{ij}\stackrel{~}{\rho }_{ij}}`$. Here $`\rho _{ij}`$ is the density operator for the two-party $`ij`$ system, and $`\stackrel{~}{\rho }_{ij}`$ is the “spin-flipped” density operator: $`\stackrel{~}{\rho }_{ij}=(\sigma _y\sigma _y)\rho _{ij}^{}(\sigma _y\sigma _y)`$. It has been suggested that the residual tangle is a measure of the irreducible three-way (“GHZ-type”) entanglement of a tripartite state, beyond any two-party (“EPR-type”) entanglement that may be contained in such a state. As such, it is of particular interest in discussing distillability below. Also, its square $`\tau _{ABC}^{}{}_{}{}^{2}`$ is a polynomial quantity, which makes it analytically tractable.
Triple States. For this set of states the residual tangle vanishes . We previously described states in this set as “triple” states , because they are equivalent under local unitary transformations to states with just three components:
$$|\psi _{\mathrm{tr}}=b|001+c|010+e|100.$$
(55)
Carteret and Sudbery refer to these as “beechnut” states; they all have $`p_{\mathrm{sum}}2`$. For triple states with $`p_{\mathrm{sum}}>2`$ each step of the infinitesimal distillation protocol reduces $`p_{\mathrm{sum}}`$, but leaves the state a triple state. If $`p_{\mathrm{sum}}=2`$, the actions on the three qubits cancel out, leaving the state unchanged. States of this type have vanishing primary yield for the tripartite distillation protocols described in section IV and in ; indeed, Dür, Vidal and Cirac have shown that no procedure can transform one copy of a state with zero residual tangle into a GHZ with nonzero probability .
Because the distillation procedures of section IV and preserve the classes of ‘special’ states, it is easy to see why they cannot produce GHZs from triple states; because all triple states have $`p_{\mathrm{sum}}2`$, they cannot include the GHZ state ($`p_{\mathrm{sum}}=3/2`$) as a limit. The set of product states (or “bystander states” in the terminology of Carteret and Sudbery) is similarly undistillable. The result of Dür, Vidal and Cirac, however, goes beyond this, since it assumes nothing about the symmetry of the procedure.
The symmetric version of state (55) (with $`b=c=e=1/\sqrt{3}`$) is termed by Dür, Vidal and Cirac the “W” state, and seems to fill a role for the zero residual tangle states similar to the role filled by the GHZ for all other states: it is, in some sense, maximally entangled. We will say a bit more about this below. All other ‘special’ classes include the GHZ as a limit, and therefore are distillable.
Generalized GHZ states. These states can be written in Schmidt form
$$|\psi =a|000+h|111.$$
(56)
They have $`p_A=p_B=p_C=a^2`$, residual tangle $`\tau _{ABC}=4a^2h^2`$. A single step of the infinitesimal distillation procedure gives a new generalized GHZ with coefficients $`a^{}=a+\mathrm{\Delta }a`$, $`h^{}=h+\mathrm{\Delta }h`$:
$$\mathrm{\Delta }a=(ϵ/2)(3p_{\mathrm{sum}})a,\mathrm{\Delta }h=+(ϵ/2)p_{\mathrm{sum}}h,$$
(57)
so (53) can readily be evaluated to give the yield of GHZs
$$Y=1\sqrt{1\tau _{ABC}}=(2/3)(3p_{\mathrm{sum}}).$$
(58)
These states are the most distillable three qubit states as a function of both $`\tau _{ABC}`$ and $`p_{\mathrm{sum}}`$; we can see this in Figures 1 and 2 below.
Slice states. In Schmidt form these are
$$|\psi =a|000+d|011e|100+h|111,ae=dh,$$
(59)
plus similar states derived by permuting the order of the bits. These states have $`p_B=p_C=a^2+e^2`$, $`p_A=a^2+d^2`$, $`\tau _{ABC}=4(ah+de)^2=4a^2(h+e^2/h)^2`$. Imposing normalization and the orthogonality condition on (59) we see that this is a two-parameter family of states. For these two parameters we may choose $`a^2`$ and $`h^2`$, or equivalently $`p_A`$ and $`p_B`$. One step of the infinitesimal distillation protocol applied to state (59) leaves qubits B and C in their Schmidt bases, but not qubit A; a change of basis must be applied to A to put the new state in Schmidt form. This new state is still a slice state, and has new parameters $`p_A^{}=p_A+\mathrm{\Delta }p_A`$, $`p_B^{}=p_B+\mathrm{\Delta }p_B`$,
$`\mathrm{\Delta }p_A`$ $`=`$ $`2ϵ(p_A+2p_B1)p_A,`$ (60)
$`\mathrm{\Delta }p_B`$ $`=`$ $`2ϵ(p_A+2p_B2)p_Bϵp_A(p_A+p_B1)/(2p_A1).`$ (61)
The yield is difficult to evaluate analytically, but numerical evidence shows that generic slice states are not extremes of distillability. With each step of the distillation protocol, the parameters $`p_B=p_C`$ approach $`1/2`$, but $`p_A`$ actually moves away. However, when $`p_B=p_C=1/2`$, this subclass of slice states does have extremal behavior. Carteret and Sudbery term this subclass the maximal slice states.
Maximal slice or Slice-ridge states are of form (59) with $`a^2+e^2=1/2=p_B=p_C`$. This subclass is parametrized by a single number, which can be taken to be $`p_A`$. Because only $`p_A`$ is larger than $`1/2`$, there is no need to perform the GHZ distillation procedure on qubits B and C; performing it on A alone preserves the form of the state, with
$$\mathrm{\Delta }p_A=(ϵ/2)p_A(1p_A),$$
(62)
giving a yield of GHZs
$$Y=1\sqrt{1\tau _{ABC}}=2(2p_{\mathrm{sum}}).$$
(63)
The expression for the primary yield in terms of the residual tangle is identical to that for the GHZ-type states, while in terms of $`p_{\mathrm{sum}}`$ it is not. In terms of $`\tau _{ABC}`$ it is one of the most distillable types of state (see Fig. 1). In terms of $`p_{\mathrm{sum}}`$ (Fig. 2) it appears to be one of the least distillable types of states; this is because maximal slice states have the minimum $`\tau _{ABC}`$ of all states with a given $`p_{\mathrm{sum}}`$.
In addition to these ‘special’ states, there are two classes of states that deserve additional attention. While these states have stabilizers of zero dimension like generic states, these classes are also preserved by the above distillation protocols. Like the ‘special’ states, they extremize distillability as a function of $`\tau _{ABC}`$ and $`p_{\mathrm{sum}}`$.
Generalized triple or Tetrahedral states. These states can be written
$$|\psi =b|001+c|010+e|100+h|111.$$
(64)
We are mainly interested here in the symmetric state $`b=c=e`$; for this case $`p_A=p_B=p_C=2b^2`$, $`p_{\mathrm{sum}}2`$. This form is preserved by the steps of the infinitesimal distillation protocol, which make the coefficients evolve according to (37); the yield is easily integrated according to (53) to give $`Y=2(2p_{\mathrm{sum}})=4(13b^2)`$; the residual tangle is $`\tau _{ABC}=16b^3\sqrt{13b^2}=\sqrt{(4Y)^3Y/27}`$. This yield is identical to that of the maximal slice states as a function of $`p_{\mathrm{sum}}`$, but not as a function of $`\tau _{ABC}`$; from Figs. 1 and 2, we see that they are states of minimal distillability in terms of both $`p_{\mathrm{sum}}`$ and $`\tau _{ABC}`$.
Zero residual tangle (ZRT) states. Dür, Vidal and Cirac have shown that no states with $`\tau _{ABC}=0`$ can be converted to GHZ triplets with nonzero probability, so $`Y=0`$. They also showed that all such states can be written in the form
$$|\psi =a|000+b|001+c|010+e|100.$$
(65)
This is in general not in the Schmidt form of section III. These states include the triple states $`a=0`$ as a subclass (for which (65) is in Schmidt form). The triple states form a boundary of this set, and any ZRT state will evolve under the distillation protocol to a triple state. These states have $`p_{\mathrm{sum}}2`$.
All these ‘special’ states have symmetries which account for both their enlarged stabilizers and their extremal distillability. One way of seeing this is to note that the various standard forms given in section II, which for generic states all require distinct bases, often coincide for these special states. For instance, the generalized GHZ states (56) are simultaneously in Schmidt, two-term, minimal, LPS and Griffiths-Niu form. ZRT and triple states cannot be written in two-term form, but can be written with three terms; the triple states (55) are simultaneously in Schmidt, minimal, three-term and Griffiths-Niu form. Slice states written in the form (59) are simultaneously in both Schmidt and LPS standard forms.
It is easiest to see how the the distillability of these states compares to that of generic states by plotting their yields $`Y`$ as a function of $`p_{\mathrm{sum}}`$ and $`\tau _{ABC}`$ along with the numerical results for a large sample of randomly generated states. We have plotted these quantities in Figs. 1 and 2, with the families of ‘special’ states indicated. We see that most of these states are indeed special as far as distillation is concerned: they form the boundaries of the plotted regions. The quantity $`\tau _{ABC}`$ does seem to be closely related to distillability, as conjectured, though this relationship is not exact; for a given value of $`\tau _{ABC}`$ states with a range of $`Y`$ values exist, but the range is not very wide. This range is bounded at the top by the generalized GHZ and maximal slice states, and at the bottom by the symmetric generalized triple state. All ZRT states have $`Y=\tau _{ABC}=0`$. There is also a relationship between $`p_{\mathrm{sum}}`$ and $`Y`$, though again for a given $`p_{\mathrm{sum}}`$ there is a range of $`Y`$ values. This range too is bounded above by the generalized GHZs, and below by the ZRT states, generalized triples and maximal slice states. These upper and lower bounds are both linear; the upper bound is exactly the same as that for Bernstein and Bennett’s Procrustean technique , reflecting the fact that generalized GHZ states can be distilled by exactly the same techniques which work in the bipartite case.
A reasonable question is to what extent these yields are artifacts of the particular distillation protocol we use. After all, this technique is only one possible way of producing GHZs, in general not the optimal method even for a single copy of a three-qubit state.
Fortunately, we can actually answer this question. Recently, Acín, Jané, Dür and Vidal have discovered the optimal algorithm for transforming a single copy of a three-qubit state into a GHZ. This involves performing a POVM on each of the three bits, designed to project the states in the two-term representation onto tri-orthogonal vectors. Finding the correct POVM for an arbitrary state involves maximizing a somewhat involved function, but is easily done numerically. We have done so for a large sample of random states, as well as for the members of the ‘special’ classes enumerated by Carteret and Sudbery.
The optimal yield is higher, in general, than that of the infinitesimal algorithm of section IV, though they are surprisingly close for most states. However, for the ‘special’ states, the yields are identical. In other words, the infinitesimal distillation technique gives the optimal yield for these classes of states. Quite remarkably, if we plot Figures 1 and 2 for the optimal GHZ distillation protocol, the figures look completely unchanged.
Thus we can see that by both the optimal and the infinitesimal techniques, these classes of special states extremize the yield of GHZs as a function of both $`\tau _{ABC}`$ and $`p_{\mathrm{sum}}`$. This strongly supports the conclusion that these states do indeed have unusual entanglement properties, and are worthy of further study.
## VI Conclusions
We have examined tripartite entanglement from both an analytical and an operational point of view. In the bipartite case, which is well understood and to which we have turned for clues, the analytical and operational aspects of entanglement are closely related: the entanglement properties of a single copy are given by the locally invariant parameters, the Schmidt coefficients, which also determine their operational characteristics. We have looked for similar connections in the three-qubit case. Here at least five locally invariant parameters are required, as opposed to just one in the two-qubit case. We have examined several ways of choosing these five parameters, looking in particular at generalizations of the bipartite Schmidt decomposition. One representation in particular, the “Schmidt form,” has useful properties which made it simple to prove the efficacy of the infinitesimal GHZ distillation protocol of ; it can also be parametrized in terms of five physically meaningful quantities.
We have looked for connections between these parameters and yields in distilling GHZ triplets, as well as connections with the residual tangle of Coffman et al. We have shown that the ‘special’ classes of states enumerated by the theorem of Carteret and Sudbery extremize the distillation yield as functions of the residual tangle $`\tau _{ABC}`$ and $`p_{\mathrm{sum}}=p_A+p_B+p_C`$.
Although a certain amount amount of progress towards understanding tripartite entanglement has been made, at least for qubits, many important questions remain unanswered. For example, the number of states in the asymptotic minimum reversible entanglement generating set (MREGS) for three-qubit states, and for tripartite states in general, is still unknown. No asymptotically reversible (or optimal but irreversible) distillation technique for GHZ states is known. The search for solutions to these and related problems is ongoing.
## Acknowledgments
We would like to thank H.A. Carteret, W. Dür, R.B. Griffiths, A. Sudbery and G. Vidal for many useful conversations. This work was supported by NSF Grant No. PHY-9900755.
Figure 1. Here we plot the primary yield of GHZ triplets from the infinitesimal distillation algorithm of section III vs. the square of the residual tangle $`\tau _{ABC}^{}{}_{}{}^{2}`$ for various ‘special states’ as well as a random sample of generic states. We see that all states lie between two curved boundaries; the generalized GHZ and maximal slice states lie on the upper boundary, while the generalized triple states lie on the lower boundary. The triple states all have both $`\tau _{ABC}`$ and the yield equal to zero. Interestingly, the maximal slice states appear to be high-yield states when plotted against $`\tau _{ABC}`$, but low-yield when plotted against $`p_{\mathrm{sum}}=p_A+p_B+p_C`$; for a given value of $`\tau _{ABC}`$ these states minimize $`p_{\mathrm{sum}}`$.
Figure 2. Here we plot the primary yield of GHZ triplets from the infinitesimal distillation algorithm of section III vs. $`p_{\mathrm{sum}}=p_A+p_B+p_C`$ for various ‘special states’ as well as a random sample of generic states. We see that all states lie between two linear boundaries; the generalized GHZ states lie on the upper boundary, while the maximal slice and generalized Triple states lie on the lower boundary, and the triple states are the zero-yield states between $`p_{\mathrm{sum}}=2`$ and $`p_{\mathrm{sum}}=3`$. The upper linear boundary corresponds to the yield of Bernstein and Bennett’s Procrustean method of EPR distillation in the bipartite case.
Figure 1.
Figure 2.
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# ORFEUS II echelle spectra : H2 measurements in the Magellanic Clouds Data partly obtained under the DARA guest observing program in the ORFEUS II Mission
## 1 Introduction
The molecular hydrogen (H<sub>2</sub>) is by far the most abundant molecule in the interstellar medium (ISM) and thus plays a key role for our understanding of the molecular gas in the ISM of the Milky Way and other galaxies. Despite its large abundance, H<sub>2</sub> in the ISM is difficult to measure because it is not seen in radio emission, in striking contrast to the second most abundant interstellar molecule, carbon monoxide (CO). H<sub>2</sub> emission lines are seen in the near infrared (NIR), but unfortunately they are weak (quadropole transitions) and thus can not be used to study the overall interstellar abundance of H<sub>2</sub>. Molecular hydrogen in the diffuse ISM can only be studied by way of absorption spectrosopy in the far ultraviolet (FUV) toward stars or other bright UV background sources. During the seventies, considerable effort was put into the investigation of H<sub>2</sub> absorption lines with the Copernicus satellite. Savage et al. (1977; hereafter S77) summarized Copernicus H<sub>2</sub> measurements of 102 lines of sight toward nearby stars in the Milky Way. One of the most striking results was the correlation between the H<sub>2</sub> column density $`N`$(H<sub>2</sub>) and colour excess $`E(BV)`$, representative of the dust amount along a sight line. This relation has been interpreted in terms of the self-shielding effect of H<sub>2</sub> (Federman et al. 1979). S77 showed that the transition from low to high molecular fractions in the local Galactic gas is found at total hydrogen column densities near $`5.0\times 10^{20}`$ cm<sup>-2</sup>.
Measurements with Copernicus, however, were limited to the very local interstellar gas of the Milky Way. For more distant background sources, Copernicus was not sensitive enough. Later UV satellites, such as IUE and HST, had better sensitivity, but these instruments do not cover the wavelength range below 1150 Å where the transitions of H<sub>2</sub> are seen. The analysis of extragalactic H<sub>2</sub> gas is of great importance since the abundance of H<sub>2</sub> can be studied in environments very different from those of the Milky Way. The Magellanic Clouds, the most nearby satellite galaxies of the Milky Way, are ideal hunting grounds for extragalactic H<sub>2</sub> measurements, because they provide many bright stars suitable as UV background sources for absorption spectroscopy. Moreover, H<sub>2</sub> has been detected from warm regions in both galaxies in the near-IR emission lines (Koornneef & Israel 1985, Israel & Koornneef 1988, 1991). It has been suggested that, due to the lower metallicity and the lower dust content, the amount of H<sub>2</sub> in the diffuse interstellar medium of the Magellanic Clouds is significantly lower than in the Milky Way (Clayton et al. 1995).
The ORFEUS telescope, launched for its second mission in 1996, was the first instrument able to measure H<sub>2</sub> absorption lines in the LMC (de Boer et al. 1998) and SMC (Richter et al. 1998). In addition, the spectrum of LH 10:3120 was used to detemine an upper limit for the H<sub>2</sub>/CO ratio in the LMC gas along this line of sight (Richter et al. 1999a). Together with the observations presented here, these ORFEUS spectra provide the first opportunity to investigate the relations between $`N(`$H i$`)`$, $`N(`$H$`{}_{2}{}^{})`$ and $`E(BV)`$ in diffuse interstellar gas of the Magellanic Clouds in comparison to the Milky Way.
## 2 Observations
The observations have been carried out during the second mission of ORFEUS on the ASTRO-SPAS space shuttle mission in Nov./Dec. 1996. ORFEUS is equipped with two alternatively operating spectrometers, the echelle spectrometer (Krämer et al. 1990) and the Berkely spectrometer (Hurwitz & Bowyer 1996). The spectroscopic data presented here were obtained with the Heidelberg/Tübingen echelle spectrometer. This instrument has a resolution of somewhat better than $`\lambda /\mathrm{\Delta }\lambda =10^4`$ (Barnstedt et al. 1999), working in the spectral range between $`912`$ and $`1410`$ Å. A detailed description of the instrument is given by Barnstedt et al. (1999).
Here we study the ORFEUS spectra of 4 LMC stars and one SMC star. Basic information about the targets is given in Table 1. The primary data reduction was performed by the ORFEUS team in Tübingen (Barnstedt et al. 1999). In order to improve signal-to-noise ratios (S/N), all spectra have been filtered by a wavelet algorithm (Fligge & Solanki 1997). The resolution after filtering is $`30`$ km s<sup>-1</sup>. Heliocentric velocities have been transformed for each line of sight into the LSR (Local Standard of Rest) system.
## 3 Data analysis
The complex line-of-sight structure in direction of the Magellanic Clouds, with contributions from local Galactic gas (0 km s<sup>-1</sup>), Galactic halo gas (near $`+60`$ and $`+120`$ km s<sup>-1</sup> in front of the LMC) and Magellanic Cloud gas (near $`+250`$ km s<sup>-1</sup> for the LMC and $`+150`$ km s<sup>-1</sup> for the SMC; see Savage & de Boer 1979, 1981; de Boer et al. 1980; Bomans et al. 1996) makes the thorough analysis of H<sub>2</sub> absorption lines at LMC velocities and $`30`$ km s<sup>-1</sup> resolution a difficult task. The main problem is that for the vast majority of the H<sub>2</sub> transitions line blends from atomic or molecular species can not be excluded, even when many of these blendings might are unlikely. As a consequence, the number of unambigously identified H<sub>2</sub> features at high radial velocities is strongly limited to only a few wavelength regions. Typical line strengths for low H<sub>2</sub> column densities (as observed in the spectra presented in the following) have values $`100`$ mÅ, which is (at low S/N) comparable with the strength of noise peaks. For most of these lines, only upper limits for the equivalent widths can be obtained. For diffuse H<sub>2</sub> clouds it is known that the process of UV pumping (Spitzer & Zweibel 1974) often leads to an excitation of the higher rotational states, particulary if the total H<sub>2</sub> column density is below the limit for the self-shielding effect. Thus, the rotational excitation of H<sub>2</sub> in the most diffuse medium often does not reflect the actual kinetic temperature of the gas. H<sub>2</sub> line strengths in diffuse clouds for excited rotational states might be significantly higher than for the ground states (see the Copernicus spectrum of $`\zeta `$ Pup, as published by Morton & Dinerstein 1976, where the strongest H<sub>2</sub> lines occur for $`J=3`$), even if the gas is cold. In the most complex case, the equivalent widths for the ground state lines are below the detection limit, while in the same spectrum, absorption from higher rotational levels is visible.
Velocity information from metal lines as well as model spectra for the excited rotational states were used to interpret the complex H<sub>2</sub> absorption pattern found in the ORFEUS spectra of stars in the Magellanic Clouds.
## 4 ORFEUS H<sub>2</sub> measurements
For the lines of sight toward HD 269698, HD 269546 and HD 36402 the analysis of H<sub>2</sub> line strengths is presented in the following. For 2 other lines of sight, ORFEUS H<sub>2</sub> measurements of Magellanic-Cloud gas have recently been published by de Boer et al. (1998; LH 10:3120) and Richter et al. (1998, 1999a; HD 5980, LH 10:3120).
Wavelengths and oscillator strengths for the H<sub>2</sub> lines have been taken from the list of Morton & Dinerstein (1976) <sup>1</sup><sup>1</sup>1We note that the Lyman P(1), 10-0 line is located at $`982.839`$ Å and not at $`982.383`$ Å, as given by Morton & Dinerstein (1976). See the wavelength list of Abgrall & Roueff (1989).. We measured equivalent widths ($`W_\lambda `$) by using either trapezoidal or gaussian fits. For the error determination we used the algorithm of Jenkins et al. (1973), taking into account photon statistics and the number of pixels involved for each line. In order to estimate the uncertainty for the choice of the continuum, we fitted a maximum and a minimum continuum to the data in the vicinity of each line and derived a mean deviation. The error for $`W_\lambda `$ given in Table 2 represents the total uncertainty calculated from all contributions discussed above. Column densities were derived by using a standard curve-of-growth technique.
### 4.1 HD 269698
The ORFEUS spectrum of HD 269698 in the Large Magellanic Cloud shows weak H<sub>2</sub> absoption at LMC velocities near $`+220`$ km s<sup>-1</sup>. Six lines from the two rotational ground states ($`J=0,1`$) with high oscillator strengths are clearly seen in the spectrum and are not blended by other transitions. For additional five lines from higher rotational states we find upper limits for the equivalent widths of $`W_\lambda 82`$ mÅ (Table 2). Fig. 1 shows three of the detected H<sub>2</sub> absorption lines plotted on a velocity scale. The lack of absorption in higher rotational states indicates that the H<sub>2</sub> gas is not strongly excited. Constructing curves of growth for each rotational state we obtain column densities of $`4.0_{2.0}^{+4.0}\times 10^{14}`$ cm<sup>-2</sup> for $`J=0`$ and $`5.0_{3.0}^{+6.0}\times 10^{15}`$ cm<sup>-2</sup> for the $`J=1`$ state, using a $`b`$ value of 8 km s<sup>-1</sup> (best fit). The total H<sub>2</sub> column density in the LMC gas toward HD 269698, derived by summing over $`N(0)`$ and $`N(1)`$, is $`N`$(H<sub>2</sub>)$`{}_{\mathrm{total}}{}^{}=5.4_{3.1}^{+7.3}\times 10^{15}`$ cm<sup>-2</sup>. The error is derived from the uncertainty for the fit to the curve of growth and includes the error for the individual equivalent widths and the uncertainty for the $`b`$ value. From the detection limits for the lines from $`J2`$ we can exclude the possibilty that the higher rotational states will significantly contribute to the total H<sub>2</sub> column density.
HD 269698 is located in the OB association N 57 at the rim of the supergiant shell LMC 4 where the H i emission (Rohlfs et al. 1984) has a minimum. The IUE spectrum of HD 269698 (Domgörgen et al. 1994) reveals three S ii components at LMC velocities in front of the star: near $`+220`$ km s<sup>-1</sup>, near $`+245`$ km s<sup>-1</sup> and near $`+290`$ km s<sup>-1</sup>. The detected H<sub>2</sub> lines obviously belong to the first component.
### 4.2 HD 269546
In the ORFEUS spectrum of the LMC star HD 269546, no clear H<sub>2</sub> absorption is visible at LMC velocities. The presence of H<sub>2</sub> at LMC velocities (near $`+200`$ km s<sup>-1</sup>) in ORFEUS data was indicated by Widmann et al. (1998), using a coaddition of 25 Lyman- and Werner lines. However, the Werner R(0), R(1) line-pair (in Fig. 1 plotted on a velocity scale) gives no hint for the presence of H<sub>2</sub> absorption at LMC velocities. Marginal H<sub>2</sub> absorption might be present in the Lyman P(1), 2-0 line ($`\lambda =1078.925`$ Å) and in the Lyman R(3), 4-0 line ($`\lambda =1053.976`$ Å) near $`+265`$ km s<sup>-1</sup> (Fig. 1), but these absorption features are not clearly distinguishable from noise peaks and no other H<sub>2</sub> profiles from $`J=1,3`$ show similar features at $`+265`$ km s<sup>-1</sup>. Metal lines in the LMC gas near $`+250`$ km s<sup>-1</sup> have been found in the IUE spectrum of HD 269546 (Grewing & Schultz-Luepertz 1980). Moreover, the IUE data reveal absorption over the whole velocity range between $`0`$ and $`+290`$ km s<sup>-1</sup>, most likely related to Galactic halo gas and weaker LMC components. H<sub>2</sub> absorption at $`+120`$ km s<sup>-1</sup> is seen in some of the stronger lines, indicating that the Galactic high-velocity gas in front of HD 269546 contains molecular gas and dust (Richter et al. 1999b). The H i emission line data from Rohlfs et al. (1984) show the LMC gas at $`+250`$ km s<sup>-1</sup>. HD 269546 is member of the OB association LH 58 in the N 144 superbubble complex northwest of 30 Doradus. The H i gas seen in 21 cm emission at $`+250`$ km s<sup>-1</sup> is most likely in front of N 144.
Detection limits for eight H<sub>2</sub> absorption lines at LMC velocities near $`+265`$ km s<sup>-1</sup> are used to obtain upper limits for the column densities of $`N(J)`$ for $`J4`$ by fitting the values of log ($`W_\lambda /\lambda `$) to the linear part of the curve of growth. We calculate an upper limit for the total H<sub>2</sub> column density by modeling the population of the rotational states for T$`{}_{\mathrm{ex}}{}^{}300`$ K. From that we derive $`N`$(H<sub>2</sub>)$`{}_{\mathrm{total}}{}^{}2.3\times 10^{15}`$ cm<sup>-2</sup> for the LMC gas toward HD 269546.
### 4.3 HD 36402
No H<sub>2</sub> absorption is seen in the spectrum of HD 36402 at LMC velocities in the range $`+220`$ to $`+320`$ km s<sup>-1</sup>. In this velocity range, atomic absorption has been found by de Boer & Nash (1982). We place an upper limit on the H<sub>2</sub> column density in the LMC gas after inspecting some of the strongest of the H<sub>2</sub> transitions in the rotational states $`J=0`$ to $`J=4`$. HD 36402 is located in the superbubble N 51D and shows hydrogen emission and metal absorption from LMC foreground gas near $`+300`$ km s<sup>-1</sup> (de Boer & Nash 1982). Therefore we expect H<sub>2</sub> absorption from LMC gas roughly at the same velocity. Inspecting the R(0), R(1) line-pair plotted on the velocity scale (Fig. 1, lowest panel) there is very weak absorption near $`+290`$ km s<sup>-1</sup>, but this feature has no significance with respect to the noise.
In the same way as for HD 269546, we determine upper limits for the individual column densities $`N(J)`$ for $`J4`$ from the detection limits for some of the stronger H<sub>2</sub> transitions near $`+300`$ km s<sup>-1</sup>. We find an upper limit for the total H<sub>2</sub> column density in the LMC gas toward HD 36402 of $`N`$(H<sub>2</sub>)$`{}_{\mathrm{total}}{}^{}1.0\times 10^{15}`$ cm<sup>-2</sup>. Again, this upper limit is calculated under the assumption that the excitation temperature of possibly existing H<sub>2</sub> gas in the LMC would not exceed a value of 300 K.
## 5 H i measurements
For two of our sight lines (HD 269698 and HD 269546) we present the determination of H i column densities from the analysis of the Ly $`\alpha `$ absorption near 1215 Å. The values for the H i column densities along the other three lines of sight have been adopted from the literature. All H i column densities are summarized in Table 3. For HD 5980 and HD 36402, the column density of H i has been determined by Fitzpatrick & Savage (1983) and de Boer & Nash (1982), respectively, using H i emission line data in combination with the Ly $`\alpha `$ absorption near 1215 Å. They derive H i column densities of $`N`$(H i$`)=1.0\times 10^{21}`$ cm<sup>-2</sup> for the SMC gas toward HD 5980 and $`1.6\times 10^{20}`$ cm<sup>-2</sup> for the LMC gas toward HD 36402. For LH 10:3120, the LMC H i column density is $`2.0\times 10^{21}`$ cm<sup>-2</sup>, obtained by a multi-component fit of the Ly $`\alpha `$ profile (Richter et al. 1999a).
We use a similar technique for the determination of H i column densities in the LMC gas toward HD 269698 and HD 269546. For HD 269546, we fix the LMC component at a velocity of $`+265`$ km s<sup>-1</sup>, similar to the velocity for which we had determined the upper limit for the H<sub>2</sub> column density in Sect. 4.2 . The velocity structure seen in metal lines (Grewing & Schulz-Luepertz 1980), however, indicates that there are definitely additional (weaker) absorption components in front of HD 269546. We thus might slightly overestimate the H i column density in the LMC gas at $`+265`$ km s<sup>-1</sup> by fitting one single LMC component to the Ly $`\alpha `$ absorption structure. The situation is even more difficult for the Ly $`\alpha `$ profile in the spectrum of HD 269698. IUE data of HD 269698 show the presence of three velocity components in this sight line (Domgörgen et al. 1994), near $`+220,+245`$ and $`+290`$ km s<sup>-1</sup>. The H i emission (Rohlfs et al. 1984) shows a weak component near $`+256`$ km s<sup>-1</sup> which could be associated with the absorption component near $`+245`$ km s<sup>-1</sup> (see Domgörgen et al. 1994). The S ii abundances found by Domgörgen et al. indicate similar total gas quantities for the two main components at $`+245`$ km s<sup>-1</sup> and at $`+220`$ km s<sup>-1</sup>. The H<sub>2</sub> absorption was found in the latter component (see Sect. 4.1). For the Ly $`\alpha `$ fit, we fix the two LMC components at $`+220`$ (cloud I) and $`+245`$ (cloud II) km s<sup>-1</sup>, assuming equal H i column densities. For the fitting procedure we use a multi-component Voigt profile. Galactic foreground absorption by H i is taken into account by a fit component at $`0`$ km s<sup>-1</sup>. The multiple velocity components are not resolved in the Ly $`\alpha `$ profile and we do not take into account additional absorption from Galactic intermediate and high-velocity gas. Thus, it is clear that our results derived by this method represent only rough estimates for the distribution of the H i gas in front of the stars. However, for the comparison between $`N(`$H i$`)`$, $`N(`$H$`{}_{2}{}^{})`$ and $`E(BV)`$, as presented in Sect. 7, the determined H i column densities are sufficiently accurate.
The Ly $`\alpha `$ fit for HD 269546 provides the best agreement with the data with a Galactic foreground absorption of $`N(`$H i$`)_{\mathrm{MW}}=3.5\pm 0.6\times 10^{20}`$ cm<sup>-2</sup> and an additional LMC component at $`+265`$ km s<sup>-1</sup> of $`N(`$H i$`)_{\mathrm{LMC}}=2.0\pm 0.5\times 10^{20}`$ cm<sup>-2</sup>. For HD 269698, the best fit is found for $`N(`$H i$`)_{\mathrm{MW}}=3.5\pm 0.7\times 10^{20}`$ cm<sup>-2</sup> and $`N(`$H i$`)_{\mathrm{LMC},\mathrm{comp}.\mathrm{I}}N(`$H i$`)_{\mathrm{LMC},\mathrm{comp}.\mathrm{II}}=1.0\pm 0.4\times 10^{20}`$ cm<sup>-2</sup> for the two LMC clouds at $`+220`$ and $`+245`$ km s<sup>-1</sup>. The H i column densities in the LMC gas in these two lines of sight are significantly lower than found for HD 5980 and LH 10:3120, but comparable with the H i column density found in the LMC gas toward HD 36402.
## 6 Colour excess $`E(BV)`$
The colour excess $`E(BV)`$ for each line of sight, adopted from different publications, is given in Table 1. All lines of sight show total values of $`E(BV)`$ lower than 0.20 mag. The main problem is to separate the contributions of the Galactic foreground from the colour excess within the Magellanic Clouds. For LH 10:3120 (LMC) and HD 5980 (SMC), we adopt values for $`E(BV)_{\mathrm{LMC}}`$ from previous publications (de Boer et al. 1998; Richter et al. 1998). For HD 269698 and HD 269546, we have calculated the foreground extinction from the H i column densities of the Galactic foreground gas by using the mean gas-to-colour excess relation, as given by Bohlin et al. (1978). The values for $`E(BV)_{\mathrm{LMC}}`$ are presented in Table 3. Their errors are based on the uncertainty for $`E(BV)_{\mathrm{total}}`$ (as cited in the literature; see Table 1) in addition to the uncertainty for the Galactic foreground reddening, which shows variations in the range of $`0.05`$ mag in direction of the LMC and $`0.02`$ mag toward the SMC (Bessel 1991).
## 7 Discussion
Values of $`N(`$H i$`)`$, $`N(`$H$`{}_{2}{}^{})`$ and $`E(BV)`$ for all five lines of sight measured with ORFEUS (summarized in Table 3) are used to investigate the diffuse molecular ISM of the Magellanic Clouds. In order to extend our data, we include recent results from Gunderson et al. (1998) for two lines of sight toward Sk $``$66 19 and Sk $``$69 270 in the LMC. They estimate, on the basis of low dispersion spectra with the Hopkins Ultraviolet Telescope (HUT), column densities of molecular hydrogen in the LMC gas by fitting H<sub>2</sub> line profiles (see Table 3).
Fig. 2 presents correlations between $`N(`$H$`{}_{2}{}^{})`$, $`E(BV)`$, $`f=2N(`$H<sub>2</sub>$`)/[N(`$H i)$`+2N(`$H$`{}_{2}{}^{})]`$, and $`N(`$H i$`+`$H$`{}_{2}{}^{})=N(`$H i$`)+2N(`$H$`{}_{2}{}^{})`$. The left panel shows log $`N(`$H$`{}_{2}{}^{})`$ plotted versus $`E(BV)`$. In principle, we find the typical relation known from the Copernicus H<sub>2</sub> survey from S77 for Galactic gas. In both Milky Way and the Magellanic Clouds the logarithmic H<sub>2</sub> column density (log $`N(`$H$`{}_{2}{}^{})`$) undergoes a transition from low to to high values at $`E(BV)0.08`$ (dashed line) due to the self-shielding effect of H<sub>2</sub> (Federman et al. 1979).
It is known that the Magellanic Clouds have a significantly lower dust content than the Milky Way. Typical gas-to-dust ratios $`N(`$H i$`+`$H$`{}_{2}{}^{})`$/$`E(BV)`$ in the Magellanic Clouds are 4 times (LMC) and 8 times (SMC) higher than in Milky Way gas (Koornneef 1982; Bouchet et al. 1985, respectively). In our sample, we find gas-to-dust ratios as high as $`3.0\times 10^{22}`$ cm<sup>-2</sup> for the gas in the LMC and SMC (see Table 3), consistent with these results. For HD 36402, HD 269698 and HD 269546, the ratios are significantly lower, but note that these values most likely represent lower limits due to the large uncertainty for the gas-to-dust ratio near the zero point of the $`E(BV)`$ scale. With respect to the generally lower dust content and the relation between H<sub>2</sub> column density and $`E(BV)`$ (Fig. 2, left panel) one should expect that the fraction of gas in molecular form is significantly lower in the Magellanic Cloud than in the Milky Way. From more theoretical considerations, Elmegreen (1989) concluded that interstellar clouds in Magellanic type irregular galaxies should be mostly atomic, since their lower metallicity directly influences the shielding function $`S`$ for the cloud layers. This author also showed that the H to H<sub>2</sub> conversion also depends sensitively on the pressure and radiation field in the ISM (Elmegreen 1993). Accordingly, typical sight lines through the Magellanic Clouds might not contain any measureable column density of H<sub>2</sub>, except for those, whose column density in H i is high enough to allow a significant fraction of the gas to convert into molecular form.
As the right panel of Fig. 2 shows, the discussed effects are slightly visible in the FUV absorption line data. The figure shows the molecular fraction $`f`$ plotted against the total hydrogen column density $`N(`$H i$`+`$H$`{}_{2}{}^{})`$. The Copernicus sample (S77) shows that the transition from low ($`f10^2`$) to high ($`f>10^2`$) molecular fractions in the local Galactic gas is found at a total hydrogen column density (‘transition column density’ $`N_\mathrm{T}`$(H i)) near $`5.0\times 10^{20}`$ cm<sup>-2</sup> (right panel, dashed line). We find high total hydrogen column densities ($`10^{21}`$ cm<sup>-2</sup>) but low molecular hydrogen fractions ($`f10^2`$) for the Magellanic Clouds gas along two of seven lines of sight. The data points of these two lines of sight toward LH 10:3120 and HD 5980 indicate that the transition column density $`N_\mathrm{T}`$ from low to high molecular fractions could be indeed higher in the Magellanic Clouds than in the Milky Way. Only for sight lines with a very high total hydrogen column density (Sk $`66`$$`19`$ and Sk $`69`$$`270`$), the molecular fractions exceeds values above 1 percent. For sight lines with $`N(`$H$`{}_{\mathrm{total}}{}^{})10^{21}`$ cm<sup>-2</sup> the molecular fractions in the Magellanic Cloud gas seem to be negligible.
Additional sight line measurements toward the Magellanic Clouds, however, are required to investigate these relations on a statistically more significant level. With a larger data set it might be possible to determine the transition column density from low to high molecular fractions as a function of the overall metallicity. Since it is known that the SMC is even more metal-poor than the LMC, it is of special interest to also investigate differences in the molecular gas fractions between LMC and SMC. For that, the FUSE satellite, launched in June 1999, holds the prospect for fresh H<sub>2</sub> absorption line data in the near future.
###### Acknowledgements.
I thank K.S. de Boer and the Heidelberg-Tübingen team for permission to use the G.O. and P.I. data on Magellanic Cloud star spectra for this study and for their great support. I thank K.S. de Boer for helpful comments on this work. When this paper was prepared, PR was supported by a grant from the DARA (now DLR) under code 50 QV 9701 3.
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# The 𝑐-Functions of Noncommutative Yang-Mills Theory from Holography
## I Introduction
Yang-Mills theory on a non-commutative space , or simply noncommutative Yang-Mills theory (NCYM), has recently received increasing attention in string theory community. A few years ago, coordinates for coincident D-branes were shown to naturally promote to matrices , signaling the relevance of noncommutative gauge theory . Later, NCYM was shown to actually appear in the D-string solution to the IIB matrix model and in various string/M(atrix) theory compactification with constant NS-NS B-background . This is not too surprising, because for a single D-brane in a background with constant gauge field-strength or rank-two anti-symmetric B-tensor, some appropriate limit should lead to a situation similar to that of a particle in the lowest Landau level, where the guiding-center coordinates are known to be non-commuting. Indeed in a recent seminal paper, among other things, Seiberg and Witten have explicitly identified the precise limit in string theory for NCYM to work, which is similar to the limit in M theory for discrete light cone quantization of Matrix theory to work . In this way, NCYM arises as a new limit in string theory, providing a new probe to non-perturbative effects in string/M(atrix) theory.
In this paper we study NCYM by exploring its supergravity dual. By now it is widely believed that gauge theory is dual to a certain limit of string theory ; in particular, type IIB supergravity on an anti-de Sitter background, say of five dimensions, can be used to describe a large-$`N`$ supersymmetric Yang-Mills (SYM) theory on the four dimensional boundary, which is known to be a conformal field theory (CFT). One important test of this AdS/CFT correspondence is the holographic derivation of the quantum Weyl anomaly in the $`D=4`$, $`𝒩=4`$, $`SU(N)`$ SYM from the generally covariant boundary counter-terms in the classical action of its bulk AdS gravity dual , with central charges reproducing the expected large-$`N`$ behavior. It seems natural to extend this correspondence between gauge theory and gravity to NCYM. The supergravity backgrounds with non-vanishing B-fields that are supposed to be dual to NCYM have been recently suggested in refs. and . Furthermore, it was observed in ref. that these supergravity duals can be derived from the Seiberg-Witten relations between closed and open string moduli, by assuming the running string tension is a simple power function of energy. This observation suggests that one should try to use the supergravity duals to explore the running behavior of NCYM.
It is known that NCYM is no longer conformally invariant, because of the length scale associated with a non-vanishing B-background. Thus, one expects that the ”central charges” of NCYM, defined as the coefficients in its Weyl anomaly, should run as a function of the energy scale. It is interesting to calculate these functions, the so-called $`c`$-functions, and to see whether they obey a generalization of the famous $`C`$-theorem in two dimensions, that asserts the $`c`$-function is always monotonically increasing with the energy scale. The consistent coupling of NCYM to a curved background is not known yet, so it is not possible at this moment to directly calculate the Weyl anomaly of NCYM on the field theory side. The goal of the present paper is to study the $`c`$-functions in a holographic manner through the supergravity dual, thus providing constraints and shedding light on the problem of consistently coupling NCYM to a curved background.
The method we are going to use to calculate the holographic Weyl anomaly is the Hamilton-Jacobi approach to the 5-dimensional bulk gravity developed by de Boer, Verlinde and Verlinde. In this approach an analogue of the first-order Callan-Symanzik equations for the 4-dimensional dual field theory on the boundary can be derived from the bulk Hamilton-Jacobi equations. The key point is to interpret the Hamilton-Jacobi functional as the quantum effective action in the dual field theory resulting from integrating out the matter degrees of freedom coupled to the boundary background gravity. Then from it one can derive the holographic Weyl anomaly and $`c`$-functions. Moreover, in the de Boer-Verlinde-Verlinde formalism there are dilaton-like scalar fields in 5-dimensional gravity. The radial profile of these fields in the bulk represents the renormalization group (RG) running of certain coupling constants in the dual field theory on the boundary. To apply this formalism, one needs to show that the NCYM’s gravity dual given in for the full 10-dimensional IIB background really has a 5-dimensional dilatonic gravity Lagrangian after dimensional reduction. In this paper we will show that this is indeed the case for an NCYM with self-dual $`\theta `$-parameters, corresponding to isotropic non-commutativity, whose gravity dual has a self-dual B-background, such that the B-field does not explicitly appear in the action for the dilaton-gravity sector after dimensional reduction to 5 dimensions.
The paper is organized as follows. In Sec. II we show that after dimensional reduction from 10-dimensional IIB supergravity, the gravity dual of an NCYM with self-dual non-commutativity parameters has a 5-dimensional Lagrangian in the form of a dilatonic gravity. In Sec. III the de Boer-Verlinde-Verlinde formalism for holographic renormalization group flows is adopted to calculate the $`c`$-functions of the NCYM at low energies under the assumption of potential dominance. In Sec. IV, we show that the $`c`$-functions defined in Sec. II transform as vectors on the dilaton-space as the beta function does. The final section, Sec. V, is devoted to conclusions and discussions. In the appendix we show how to generalize the de Boer-Verlinde-Verlinde formalism to the non-canonical form of dilatonic gravity.
## II Dilaton Gravity Dual of Noncommutative Yang-Mills
In contrast to the gravity dual of ordinary Yang-Mills theory, the supergravity dual of NCYM involves turning on nontrivial scalar and various $`r`$-form (anti-symmetric $`r`$-tensor) backgrounds with a radial profile. One may wonder if there exists a 5-dimensional gravity action from which the equations of motion dimensionally reduced from 10-dimensional IIB supergravity can be derived. We will show that there is indeed such a 5-dimensional dilatonic gravity action, at least for the case with self-dual B-backgrounds.
The bosonic action for type IIB supergravity in ten dimensions in the Einstein frame is
$$I_{10}=\frac{1}{2\kappa _{10}^2}d^{10}z\sqrt{detG}[R\frac{1}{2}(\varphi )^2\frac{1}{2}e^{2\varphi }(\chi )^2\frac{1}{23!}e^\varphi H_3^2\frac{1}{23!}e^\varphi F_3^2\frac{1}{45!}F_5^2],$$
(1)
where $`\varphi `$ is the dilaton, $`\chi `$ the RR scalar and the form strengths are defined as
$`H_3`$ $`=`$ $`dB_2,`$ (2)
$`F_3`$ $`=`$ $`dA_2\chi H_3,`$ (3)
$`F_5`$ $`=`$ $`dA_4{\displaystyle \frac{1}{2}}A_2H_3+{\displaystyle \frac{1}{2}}B_2F_3.`$ (4)
Here $`B_2`$ and $`A_2`$ are respectively the NS-NS and RR 2-form potentials, $`A_4`$ are the RR 4-form potential and the 5-form strength $`F_5`$ is self-dual, that is
$$F_5=iF_5,$$
(5)
where $``$ is the 10-dimensional Hodge dual.
In this paper, we only consider self-dual B-backgrounds, together with the following conditions motivated by supersymmetry:
$`\chi `$ $``$ $`ie^\varphi =ic,`$ (6)
$`F_3`$ $`=`$ $`icH_3,`$ (7)
$`B_{01}`$ $`=`$ $`B_{23},A_{01}=A_{23}.`$ (8)
where $`c`$ is a real constant.
These conditions are consistent with the equations of motion for scalars and two-form potentials. Moreover, they make the contributions to the energy-momentum tensor from the NS-NS and RR sector cancel each other, except the one from self-dual five-form strength. The Einstein equations and the Gauss law for the five-form strength thus form a closed set of equations:
$`R_{MN}=T_{MN}^F,`$ (9)
$`_M(\sqrt{detG}F^{MNPQR})=0,`$ (10)
where $`T_{MN}^F=\frac{1}{44!}F_{MPQRS}F_N^{PQRS}`$ is the energy-momentum tensor of the five-form strength. Once the solution to this set is known, the other fields can be solved from their equations of motion.
We now perform dimensional reduction by using the following ansatz for the 10-dimensional metric:
$$ds^2=G_{MN}(z)dz^Mdz^N=\widehat{g}_{mn}(x)dx^mdx^n+\mathrm{}^2\mathrm{\Phi }(x)\overline{g}_{ab}(y)dy^ady^b,$$
(11)
with the indices {$`m,n`$} running on the reduced 5-dimensional manifold $``$, and indices {$`a,b`$} on a prescribed 5-sphere with $`x`$-dependent radius $`\sqrt{\mathrm{}^2\mathrm{\Phi }(x)}`$, where $`\mathrm{}^2`$ is the typical length scale of $``$ related to the D3-brane charge (or the 5-form flux), and the Jordan-Brans-Dicke scalar $`\mathrm{\Phi }`$ turns out to be the dilaton in the reduced theory.
Using this ansatz we can solve the Gauss law equation (10) for the five-form strength, and the solution is
$$F_{mnopq}=i4\mathrm{}^1\mathrm{\Phi }^{5/2}\sqrt{det\widehat{g}}ϵ_{mnopq},$$
(12)
where the $`ϵ`$symbol is equal to $`1`$$`(1)`$ for even (odd) permutations of $`0,1,2,3,r`$, and to $`0`$ otherwise.
The components of the corresponding energy-momentum tensor are
$$T_{mn}^F=\frac{4}{\mathrm{}^2}\mathrm{\Phi }^5\widehat{g}_{mn},T_{ab}^F=\frac{4}{\mathrm{}^2}\mathrm{\Phi }^4\overline{g}_{ab},T_{ma}^F=0.$$
(13)
The 10-dimensional Einstein equation (9) then decomposes into
$`R_{mn}`$ $`=`$ $`\widehat{R}_{mn}{\displaystyle \frac{5}{2}}\mathrm{\Phi }^1\widehat{}_m\widehat{}_n\mathrm{\Phi }+{\displaystyle \frac{5}{4}}\mathrm{\Phi }^2\widehat{}_m\mathrm{\Phi }\widehat{}_n\mathrm{\Phi }={\displaystyle \frac{4}{\mathrm{}^2}}\mathrm{\Phi }^5\widehat{g}_{mn},`$ (14)
$`R_{ab}`$ $`=`$ $`({\displaystyle \frac{4}{\mathrm{}^2}}{\displaystyle \frac{1}{2}}\widehat{}^2\mathrm{\Phi }{\displaystyle \frac{3}{4}}\mathrm{\Phi }^1(\widehat{}\mathrm{\Phi })^2)\overline{g}_{ab}={\displaystyle \frac{4}{\mathrm{}^2}}\mathrm{\Phi }^4\overline{g}_{ab},`$ (15)
$`R_{ma}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Phi }^1_m\mathrm{\Phi }\overline{g}^{bc}(\overline{}_c\overline{g}_{ab}\overline{}_a\overline{g}_{bc})=0,`$ (16)
where $`\widehat{R}_{mn}`$ is the Ricci tensor of the metric $`\widehat{g}_{mn}`$, while $`\widehat{}`$ and $`\overline{}`$ the covariant derivative with respect to $`\widehat{g}_{mn}`$ and $`\overline{g}_{ab}`$ respectively.
Given the prescribed 5-sphere metric $`\overline{g}_{ab}`$, eq. (15) reduces to a single equation of motion for the dilaton $`\mathrm{\Phi }`$, and eq. (16) is just an identity because of the metricity condition.
It turns out that this reduced set of equations of motion (14) and (15) for $`\widehat{g}_{mn}`$ and $`\mathrm{\Phi }`$ can be derived from the following 5-dimensional action for a dilatonic gravity
$$I_5=\frac{V_5}{2\kappa _{10}^2}_{}d^5x\sqrt{det\widehat{g}}\mathrm{\Phi }^{5/2}[\widehat{R}+5\mathrm{\Phi }^2(\widehat{}\mathrm{\Phi })^2)+\mathrm{}^2(20\mathrm{\Phi }^18\mathrm{\Phi }^5)],$$
(17)
where $`V_5`$ is the volume of the unit 5-sphere. This action reduces to the familiar action for AdS gravity if we set $`\mathrm{\Phi }=1`$.
To bring the gravity action to the canonical Einstein-Hilbert form we need to do the following Weyl transformation
$$\widehat{g}_{mn}=\mathrm{\Phi }^{5/3}g_{mn},$$
(18)
and the corresponding new 5-dimensional action is
$$I_5^{EH}=\frac{V_5}{2\kappa _{10}^2}_{}d^5x\sqrt{detg}[R\frac{10}{3}\mathrm{\Phi }^2(\mathrm{\Phi })^2+\mathrm{}^2\mathrm{\Phi }^{8/3}(208\mathrm{\Phi }^4)],$$
(19)
where the unhatted quantities are with respect to the new 5-dimensional metric $`g_{mn}`$.
This action was derived before from the gauged supergravity point of view in a different context. Our above discussions establish that the proposed gravity dual of an NCYM with isotropic non-commutativity has a 5-dimensional dilatonic gravity description given by the action (17) or (19).
One may feel odd at first sight that in the 5-dimensional reduced dilatonic gravity dual, the 2-form field that specifies the non-commutativity in the original boundary Yang-Mills theory does not show up explicitly. The puzzle is resolved by noticing that the self-duality conditions (8) place strong restrictions on the holographic profile of the dilaton $`\mathrm{\Phi }`$ to make it dependent on the asymptotic value of the 2-form field. This is most easily seen from the full expression of Maldacena and Russo’s solution (in the near horizon limit):
$`ds_E^2=\mathrm{}^2r^2\mathrm{\Phi }\{\mathrm{\Phi }^2(dx_0^2+dx_1^2+dx_2^2+dx_3^2)+r^4dr^2+r^2d\mathrm{\Omega }_5^2\},`$ (20)
$`\mathrm{\Phi }`$ $`=`$ $`(1+a^4r^4)^{1/2},F_{0123r}=i4\mathrm{}^1\mathrm{\Phi }^{5/2}\sqrt{det\widehat{g}}=i4\mathrm{}^4r^3\mathrm{\Phi }^4`$ (21)
$`B_{01}`$ $`=`$ $`B_{23}=\sqrt{g_s}a^2r^4\mathrm{\Phi }^2,A_{01}=A_{23}=i{\displaystyle \frac{a^2\mathrm{}^2}{\sqrt{g_s}}}r^4\mathrm{\Phi }^2,`$ (22)
$`e^{2\varphi }`$ $`=`$ $`g_s^2\mathrm{\Phi }^2,\chi =i{\displaystyle \frac{a^4r^4}{g_s}},`$ (23)
where $`g_s`$ is the string coupling in the IR limit $`r=0`$, and $`a^2=\sqrt{4\pi \alpha ^{}_{}{}^{}2g_sN}/B^{\mathrm{}}`$, $`\mathrm{}^2=\alpha ^{^{}}\sqrt{4\pi N}`$ with $`\alpha ^{^{}}`$ the string scale, $`N`$ the D3-brane charge and $`B^{\mathrm{}}`$ the boundary value of B-field at $`r=\mathrm{}`$. Note that the world-volume (specified by the directions 0,1,2,3) quantities have been properly re-scaled as in , and also that with a self-dual 2-form background (or non-commutativity parameters) the world volume metric for the NCYM remains isotropic.
As one can see, the profile of $`\mathrm{\Phi }`$ is chosen so that the solutions of NSNS and RR scalars in (23) satisfy the first equation of the self-duality conditions (8), which would not hold for an arbitrary dilaton profile. Following the proposal of , we then conclude that the dilatonic gravity (19) with the specific dilaton radial profile given in (23) is the holographic dual of the NCYM with isotropic non-commutativity. Moreover, we can read the 5-dimensional background metrics from (23), that is
$$\widehat{g}_{mn}dx^mdx^n=\mathrm{}^2r^2\mathrm{\Phi }dr^2+\mathrm{}^2r^2\mathrm{\Phi }^1dx_{}^2$$
(24)
for the action $`I_5`$, and
$$g_{mn}dx^mdx^n=\mathrm{}^2r^2\mathrm{\Phi }^{8/3}dr^2+\mathrm{}^2r^2\mathrm{\Phi }^{2/3}dx_{}^2$$
(25)
for the action $`I_5^{EH}`$, where $`x_{}`$ represents the longitudinal coordinates and $`r`$ is called the holographic coordinate which is the energy scale from the field theory point of view.
## III Holographic RG flow of NCYM in Self-Dual B-background
The existence of a consistent effective 5-dimensional dilatonic gravity allows us to generalize a counter-term generating algorithm in the AdS/CFT correspondence, known as the holographic renormalization group (RG) flow that is determined by the dilaton profile. The dilaton is interpreted as an effective coupling running with the energy scale in the dual field theory. If the dilaton has a constant radial profile, the theory reduces to pure AdS gravity with the holographic dual a CFT having a vanishing beta function. When the dilaton has a nontrivial radial profile, the holographic Callan-Symanzik RG equations have been constructed by de Boer, Verlinde and Verlinde in an elegant formalism using the standard Hamilton-Jacobi theory. The $`c`$-functions in the Weyl anomaly can then be calculated.
The essence of the de Boer-Verlinde-Verlinde formalism is the observation that though the equations of motion of the 5-dimensional supergravity is of second order, the evolution equation of its on-shell action $`S`$, derived from the standard Hamilton-Jacobi theory, is of first order and takes the usual form of the Callan-Symanzik equations, therefore $`S`$ can be interpreted as the 4-dimensional quantum effective action after integrating out all the matter degrees of freedom coupled to the background gravity.
According to the holographic interpretation of the gravity dual, a preferred radial coordinate in the bulk gravity can be selected out as representing the energy scale of the dual field theory. For simplicity, we choose the ”temporal” gauge for 5-dimensional metric
$$g_{mn}dx^mdx^n=d\rho ^2+\gamma _{\sigma \nu }(\rho ,x)dx^\sigma dx^\nu ,$$
(26)
where $`\rho `$ is the holographic radial coordinate.
For the metric on a boundary screen located at the radial position $`\rho `$, we can further separate out the holographic coordinate dependence:
$$\gamma _{\sigma \nu }(\rho ,x)=\mu ^2(\rho )\overline{\gamma }_{\sigma \nu }(x).$$
(27)
where $`\overline{\gamma }_{\sigma \nu }`$ is the background geometry seen by the dual field theory at some fundamental scale, and the warp factor $`\mu ^2`$ is the overall length scale on the boundary screen. According to the holographic UV/IR relation, $`\mu `$ stands for the energy scale of the boundary QFT, and we define the beta function for the dilaton $`\mathrm{\Phi }`$ by
$$\beta \mu \frac{d\mathrm{\Phi }}{d\mu },$$
(28)
which can be easily calculated once the 5-dimensional metric and the dilaton profile are given.
For example, the proposed gravity dual (23) of NCYM with isotropic commutativity has the dilaton profile
$$\mathrm{\Phi }=(1+a^4r^4)^{1/2},$$
(29)
and the energy scale can be read from the defining metric (25), (26) and (27):
$$\mu =\mathrm{}r\mathrm{\Phi }^{1/3}=\frac{\mathrm{}}{a}\mathrm{\Phi }^{1/3}(\mathrm{\Phi }^21)^{1/4},$$
(30)
and the resulting beta function from (28) and (30) is
$$\beta =\frac{6\mathrm{\Phi }(\mathrm{\Phi }^21)}{5\mathrm{\Phi }^22}.$$
(31)
Note that $`\mu `$ is a monotonically increasing function of $`\mathrm{\Phi }`$ and $`r`$, so the UV limit $`r\mathrm{}`$ corresponds to $`\mu \mathrm{}`$ and $`\mathrm{\Phi }\mathrm{}`$, and the IR limit $`r0`$ to $`\mu 0`$ and $`\mathrm{\Phi }1`$.
To develop the Hamilton-Jacobi theory, we shall cast the 5-dimensional dilatonic gravity action into the canonical formalism using the above metric:
$`I`$ $`=`$ $`{\displaystyle \frac{1}{2\kappa _5^2}}{\displaystyle _{}}d^5x\sqrt{detg}[R+{\displaystyle \frac{1}{2}}G(\mathrm{\Phi })(\mathrm{\Phi })^2+V(\mathrm{\Phi })]`$ (32)
$``$ $`{\displaystyle \frac{1}{2\kappa _5^2}}{\displaystyle 𝑑\rho L},`$ (33)
$`L`$ $`=`$ $`{\displaystyle d^4x\sqrt{det\gamma }[\pi _{\sigma \nu }\dot{\gamma }^{\sigma \nu }+\mathrm{\Pi }\dot{\mathrm{\Phi }}]}.`$ (34)
Here $`\dot{}`$ denotes the derivative with respect to $`\rho `$, and the canonical momenta and the Hamiltonian density are defined by
$`\pi _{\sigma \nu }{\displaystyle \frac{1}{\sqrt{det\gamma }}}{\displaystyle \frac{\delta L}{\delta \dot{\gamma }^{\sigma \nu }}},\mathrm{\Pi }{\displaystyle \frac{1}{\sqrt{det\gamma }}}{\displaystyle \frac{\delta L}{\delta \dot{\mathrm{\Phi }}}},`$ (35)
$``$ $``$ $`{\displaystyle \frac{1}{3}}\pi ^2\pi _{\sigma \nu }\pi ^{\sigma \nu }+{\displaystyle \frac{\mathrm{\Pi }^2}{2G}},`$ (36)
$``$ $``$ $`+{\displaystyle \frac{1}{2}}G\gamma ^{\sigma \nu }_\sigma \mathrm{\Phi }_\nu \mathrm{\Phi }+V,`$ (37)
with $``$ the Ricci scalar of the boundary metric $`\gamma _{\mu \nu }`$. Note that $`d^4x`$ is the action dimensionally reduced from (32).
The defining equations of canonical momenta (35) can be inverted to obtain the first order flow equations
$`\dot{\gamma }_{\sigma \nu }`$ $`=`$ $`2\pi _{\sigma \nu }{\displaystyle \frac{2}{3}}\gamma _{\sigma \nu }\pi _\sigma ^\sigma ,`$ (38)
$`\dot{\mathrm{\Phi }}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Pi }}{G}}.`$ (39)
These equations will be helpful in solving the resulting Hamilton-Jacobi equation.
In the canonical formulation of the gravity theory, the Hamiltonian density $``$ gives rise to a constraint $`=0`$, imposed upon the canonical variables:
$$\frac{1}{3}\pi ^2\pi _{\sigma \nu }\pi ^{\sigma \nu }+\frac{\mathrm{\Pi }^2}{2G}=+\frac{1}{2}G\gamma ^{\sigma \nu }_\sigma \mathrm{\Phi }_\nu \mathrm{\Phi }+V.$$
(40)
We then introduce the Hamilton-Jacobi functional $`S`$ with a properly assumed form, and see if we can derive first-order evolution equations for terms in $`S`$. With the hint of the AdS/CFT correspondence, one interprets $`S`$ as the quantum effective action of the dual field theory after integrating out the matter degrees of freedom coupled to the background gravity, which is assumed of the usual form on a curved space:
$`S[\gamma ,\mathrm{\Phi }]`$ $`=`$ $`S_{EH}[\gamma ,\mathrm{\Phi }]+\mathrm{\Gamma }[\gamma ,\mathrm{\Phi }],`$ (41)
$`S_{EH}[\gamma ,\mathrm{\Phi }]`$ $`=`$ $`{\displaystyle d^4x\sqrt{det\gamma }[Z(\mathrm{\Phi })+\frac{1}{2}M(\mathrm{\Phi })\gamma ^{\sigma \nu }_\sigma \mathrm{\Phi }_\nu \mathrm{\Phi }+U(\mathrm{\Phi })]}.`$ (42)
$`S_{EH}`$ is the tree level renormalized action which is similar in structure to the Lagrangian density $``$, and $`\mathrm{\Gamma }`$ contains the higher-derivative and non-local terms.
In the Hamilton-Jacobi theory, the canonical momenta are related to the Hamilton-Jacobi functional $`S`$ by
$$\pi _{\sigma \nu }=\frac{1}{\sqrt{det\gamma }}\frac{\delta S}{\delta \gamma ^{\sigma \nu }},\mathrm{\Pi }=\frac{1}{\sqrt{det\gamma }}\frac{\delta S}{\delta \mathrm{\Phi }}.$$
(43)
With these relations and the interpretation of $`S`$ as the effective quantum action, the quantum average of the boundary stress tensor $`<T_{\sigma \nu }>`$ and that of the gauge invariant operator $`<O_\mathrm{\Phi }>`$ to which $`\mathrm{\Phi }`$ couples can be related to $`\mathrm{\Gamma }`$ by
$$<T_{\sigma \nu }>=\frac{2}{\sqrt{det\gamma }}\frac{\delta \mathrm{\Gamma }}{\delta \gamma ^{\sigma \nu }},<O_\mathrm{\Phi }>=\frac{1}{\sqrt{det\gamma }}\frac{\delta \mathrm{\Gamma }}{\delta \mathrm{\Phi }}.$$
(44)
The factor of two is determined from the Hamilton-Jacobi equation by requiring the correct proportionality to the beta-function term in the Weyl anomaly:
$$<T_\sigma ^\sigma >=\beta <O_\mathrm{\Phi }>c_{\sigma \nu }^{\sigma \nu }+d^2,$$
(45)
where $`\beta `$ is the beta function defined in (28), and $`c`$ and $`d`$ are the $`c`$-functions.
Substituting (42) into (43) we obtain the explicit form of the canonical momenta, which will be helpful in solving the Hamilton-Jacobi equation (40),
$`\pi _{\sigma \nu }`$ $`=`$ $`{\displaystyle \frac{1}{2}}<T_{\sigma \nu }>+Z_{\sigma \nu }+({\displaystyle \frac{M}{2}}Z^{^{\prime \prime }})_\sigma \mathrm{\Phi }_\nu \mathrm{\Phi }Z^{^{}}_\sigma _\nu \mathrm{\Phi }`$ (46)
$``$ $`{\displaystyle \frac{1}{2}}\gamma _{\sigma \nu }[Z+({\displaystyle \frac{M}{2}}2Z^{^{\prime \prime }})(\mathrm{\Phi })^22Z^{^{}}^2\mathrm{\Phi }+U],`$ (47)
$`\mathrm{\Pi }`$ $`=`$ $`<O_\mathrm{\Phi }>+Z^{^{}}{\displaystyle \frac{M^{^{}}}{2}}(\mathrm{\Phi })^2M^2\mathrm{\Phi }+U^{^{}},`$ (48)
where $`^{}`$ denotes the derivative with respect to $`\mathrm{\Phi }`$ and the covariant derivatives here are with respect to $`\gamma _{\sigma \nu }`$.
To derive the evolution equations for terms in $`S_{EH}`$, we insert the expansion (47) and (48) into (40), and solve the resulting Hamilton-Jacobi equation by equating terms on both sides with the same functional form. With this procedure we get from the potential term,
$$\frac{U^2}{3}+\frac{U^{}_{}{}^{}2}{2G}=V,$$
(49)
and from the curvature term,
$$\frac{U}{3}Z+\frac{U^{^{}}}{G}Z^{^{}}=1.$$
(50)
Note both are first-order evolution equations.
Moreover, combining the second-order curvature terms and the first-order terms in the quantum average $`<T_\sigma ^\sigma >`$ and $`<O_\mathrm{\Phi }>`$, we can obtain the expression for the Weyl anomaly which is of the form of (45), with the $`c`$-functions given by
$$c=\frac{6Z^2}{U},d=\frac{2}{U}(Z^2+\frac{3Z^{}_{}{}^{}2}{2G}).$$
(51)
We can rewrite the curvature part of (45) in terms of the Euler density $``$, the Weyl density $`𝒲`$ and the Ricci scalar squared as follows
$$c_{\sigma \nu }^{\sigma \nu }+d^2=\frac{c}{2}(𝒲)+(d\frac{c}{3})^2,$$
(52)
where
$$=^24_{\sigma \nu }^{\sigma \nu }+_{\sigma \nu \lambda \delta }^{\sigma \nu \lambda \delta },𝒲=\frac{1}{3}^22_{\sigma \nu }^{\sigma \nu }+_{\sigma \nu \lambda \delta }^{\sigma \nu \lambda \delta }.$$
(53)
Note that $``$ is a topological density, and $`𝒲`$ is an invariant under Weyl transformations, so that the combination $`𝒲`$ is invariant up to a total derivative under Weyl transformations. However, the $`^2`$ term is not a Weyl invariant, whose presence signals the non-conformal nature of NCYM when $`c3d`$, as shown later.
Because of the nonlinearity, it is difficult to solve $`U`$ from (49). We, however, can solve it from the flow equations by substituting (47) and (48) into (38) and (39). Assuming that the theory is at sufficiently low energy scale compared to the cutoff so that the potential term dominates, it then yields
$`U`$ $`=`$ $`{\displaystyle \frac{6\dot{\mu }}{\mu }},`$ (54)
$`\beta `$ $``$ $`\mu {\displaystyle \frac{d\mathrm{\Phi }}{d\mu }}={\displaystyle \frac{6U^{^{}}}{GU}}.`$ (55)
Clearly the effective cosmological constant $`U`$ is over-determined by three equations (49), (54) and (55), the consistency of the solutions among them will imply the validity of the formalism and the assumption of potential dominance, which reminds us that the theory is at sufficiently low energy scale.
On the other hand, the effective inverse Newton constant $`Z`$ will be determined by (50) up to an integration constant given by the initial conditions. There are also equations determining $`M`$ in the kinetic term from the input of $`U`$ and $`Z`$; however, we omit them since our interest is the $`c`$-functions which are independent of $`M`$, and it is easy to show that $`M`$ can be consistently solved from the Hamilton-Jacobi equation.
Having the formalism at hand, we are ready to calculate the running behavior of the quantum effective action $`S`$ for the NCYM from its dilatonic gravity dual defined by (19) and (25). The beta function for $`\mathrm{\Phi }`$ has been given in (31). Compare (19) and (32), we have
$$G(\mathrm{\Phi })=20\mathrm{\Phi }^2/3,V(\mathrm{\Phi })=\mathrm{}^2\mathrm{\Phi }^{8/3}(208\mathrm{\Phi }^4).$$
(56)
With these data and eq. (31) for the beta function, we find the solutions for the effective cosmological constant $`U`$ from the three equations mentioned above agree with each other, all giving
$$U=\frac{2\mathrm{\Phi }^{10/3}}{\mathrm{}}(5\mathrm{\Phi }^22).$$
(57)
The running behavior of the effective inverse Newton constant is then determined from (50) and is given by
$$Z=\frac{\mathrm{}}{6}\mathrm{\Phi }^{2/3}(\mathrm{\Phi }^2+2)+Z_0\mathrm{\Phi }^{2/3}(\mathrm{\Phi }^21)^{1/2}.$$
(58)
Note that the second term blows up in the IR limit $`\mathrm{\Phi }1`$ if $`Z_00`$, which will violate the assumption of potential dominance at low energy scale; and thus we are forced to set $`Z_0=0`$.
Finally, from (51) the resulting $`c`$-functions are (for $`Z_0=0`$)
$$c=\frac{\mathrm{}^3}{12}\frac{\mathrm{\Phi }^2(\mathrm{\Phi }^2+2)^2}{5\mathrm{\Phi }^22},d=\frac{\mathrm{}^3}{60}\frac{\mathrm{\Phi }^2(\mathrm{\Phi }^4+8\mathrm{\Phi }^2+6)}{5\mathrm{\Phi }^22}.$$
(59)
In the above equations (56) to (59), the profile of the dilaton $`\mathrm{\Phi }`$ is given by eq. (29). Like the beta function, the $`c`$-functions are monotonically increasing with $`\mathrm{\Phi }`$ (and thus with $`\mu `$) for $`\mathrm{\Phi }1`$. This is a generalization of Zamolodchikov’s C-theorem in two dimensions that the c functions are always monotonically increasing with the energy scale; so one may say that the C-theorem holds true in the present case. Away from the IR limit, $`\mathrm{\Phi }>1`$ and the ratio $`c/d3`$, differing from the one ($`c/d=3`$) for ordinary Yang-Mills theory in the usual AdS/CFT correspondence with $`\mathrm{\Phi }1`$, which is the IR limit of the NCYM.
## IV $`c`$-functions as vectors on the $`\mathrm{\Phi }`$-Space
In the section II we have seen that the form of the 5-dimensional dilatonic gravity action, dimensionally reduced from 10-dimensional supergravity as the dual of NCYM, is not unique. We have obtained two such actions, one given by (19) with a canonical Einstein-Hilbert term for gravity and the other (17) of a non-canonical form; they are related to each other by a Weyl transformation (18). In the section III, we have chosen to work with the canonical form of the action (19). One may wonder what are the resulting beta and $`c`$-functions if we work with the non-canonical action (17). The answer for the beta function is straightforward: from its definition (28), it should transform as a vector on the $`\mathrm{\Phi }`$-space which can be thought as the coupling constant space of NCYM. To be explicit, let us call the energy scale parameter $`\mu _q`$<sup>*</sup><sup>*</sup>*The subscript $`q`$ will be used to specify the non-canonical counterparts of the quantities defined in section III; the same convention will be also adopted in the Appendix. for the non-canonical gravity in contrast to the parameter $`\mu `$ defined for the canonical one. These two quantities are related to each other by the Weyl transformation (18), which through (27) leads to
$$\mu _q=\mathrm{\Phi }^{5/6}\mu .$$
(60)
From this, the beta functions in the two cases are related by
$`\beta `$ $``$ $`\mu {\displaystyle \frac{d\mathrm{\Phi }}{d\mu }}=\mathrm{\Omega }\mu _q{\displaystyle \frac{d\mathrm{\Phi }}{d\mu _q}}\mathrm{\Omega }\beta _q,`$ (61)
$`\mathrm{\Omega }`$ $``$ $`{\displaystyle \frac{\mu }{\mu _q}}{\displaystyle \frac{d\mu _q}{d\mu }}={\displaystyle \frac{3}{5\mathrm{\Phi }^22}}.`$ (62)
Though the beta function has clear geometric meaning by its definition, it is not clear if the $`c`$-functions have also the geometric meaning as vectors on the $`\mathrm{\Phi }`$-space. To answer this question, we need to generalize the de-Boer-Verlinde-Verlinde formalism to the non-canonical action. The generalization is straightforward but tedious, we will leave the details to Appendix A. The resulting $`c`$\- functions turn out to be
$$c_q=\mathrm{\Omega }^1c,d_q=\mathrm{\Omega }^1d,$$
(63)
and are thus vectors on the $`\mathrm{\Phi }`$-space. Note that (63) is true as long as the integration constants for $`Z_q`$ and $`Z`$ are set to equal, that is
$$Z_q=\mathrm{\Phi }^{5/3}Z=\frac{\mathrm{}}{6}\mathrm{\Phi }(\mathrm{\Phi }^2+2)+\frac{Z_0\mathrm{\Phi }}{\sqrt{\mathrm{\Phi }^21}}.$$
(64)
Now that the $`c`$-functions have a geometric interpretation, it would be interesting to see if the C-theorem may have a generic geometric origin. This issue has been explored in the recent works on the gravity side. We leave this problem for NCYM for future study.
## V Conclusions and Discussions
Since Maldacena and Russo proposed the supergravity dual of NCYM, not much has been done along this line. In this paper, we first pointed out that the gravity dual of NCYM with isotropic non-commutativity has a consistent 5-dimensional action in the form of a dilatonic gravity, which enables us to adopt the holographic RG flow approach to investigate the physics on the dual field theory side, generalizing the usual AdS/CFT correspondence.
We adopted the de Boer-Verlinde-Verlinde formalism to evaluate the $`c`$-functions at low energies, under the assumption of potential dominance, and found that the $`C`$-theorem holds true in the present case. The ratio of the two coefficient fucntions in the Weyl anomaly away from the IR limit is different from that in ordinary Yang-Mills theory, indicating the non-conformal nature of the NCYM. All of these were seen from the dual gravity side. To examine these phenomena directly inside the NCYM itself is worthwhile, especially because the perturbative techniques of non-commutative field theory seem to have become matured in a series of recent works.
The calculations of Weyl anomaly and boundary counter-terms for the boundary conformal theory from the AdS gravity have been performed in many different ways, they all agree to each other. Not much similar efforts have been spent for the non-commutative cases. Besides the method adpoted in this paper, there is an alternative approach by generalizing the method of Henningson and Skenderis to dilaton gravity which applies only to asymptotically AdS spacetime. However, as pointed out in the second paper of ref. , the NCYM dual at hand, corresponding to eq. (57) there, has not asymptotic AdS region in UV. It would be interesting to see whether an improvement of their approach can be applied to the NCYM dual.
Although we have defined the $`c`$-functions from its gravity dual by calculating the Weyl anomaly, we still lack a general understanding from the field theory side. It has been an issue of defining sensible $`c`$-functions in 4 dimensions and there is an on-going debate about the validity of a general 4-dimensional C-theorem. In section IV, we clarified the nature of $`c`$-functions on the coupling constant space, and showed they are indeed vectors on the coupling constant space as the beta function. We hope this geometric understanding will help in constructing a geometric realization of the C-theorem in 4 dimensions.
Up to now, we have only considered the supergravity background with self-dual B-field configurations. It would be interesting to consider more general B-backgrounds, which will correspond to NCYM with anisotropic non-commutativity. The 5-dimensional gravity dual will then be a dilaton gravity coupled to the 2-form potentials, and we need to generalize the de Boer-Verlinde-Verlinde formalism to include the dynamics of 2-form potentials, which may help us to understand more about the physics of NCYM from its gravity dual.
## A Generalization of de Boer-Verlinde-Verlinde Formalism For Non-Canonical Gravity
We start with the non-canonical action (17) and cast it into the ADM form as done for the canonical one:
$`I_5`$ $`=`$ $`{\displaystyle \frac{V_5}{2\kappa _{10}^2}}{\displaystyle d^5x\sqrt{det\widehat{g}}[X_q(\mathrm{\Phi })\widehat{R}+\frac{1}{2}G_q(\mathrm{\Phi })(\widehat{}\mathrm{\Phi })^2+V_q(\mathrm{\Phi })]}`$ (A.1)
$``$ $`{\displaystyle \frac{V_5}{2\kappa _{10}^2}}{\displaystyle 𝑑rL_q},`$ (A.2)
where $`X_q`$, $`G_q`$ and $`V_q`$ can be read from (17).
Decompose the metric into the warped form
$$\widehat{g}_{mn}dx^mdx^n=N^2d\rho ^2+\gamma _{\sigma \nu }(\rho ,x)dx^\sigma dx^\nu ,N=\pm 1,$$
(A.3)
(with $`N`$ the lapse function). Using the identity
$$\widehat{R}=2_m(n^m_nn^n)+(𝒦_{\sigma \nu }𝒦^{\sigma \nu }𝒦^2)$$
(A.4)
where $`n^m`$ is the boundary unit normal, $``$ and $`𝒦`$ are the intrinsic and extrinsic boundary curvature respectively, we then have
$$L_q=d^4x\sqrt{det\gamma }[\pi _{\sigma \nu }\dot{\gamma }^{\sigma \nu }+\mathrm{\Pi }\dot{\mathrm{\Phi }}N],$$
(A.5)
with
$``$ $``$ $`N[{\displaystyle \frac{1}{X_q}}({\displaystyle \frac{1}{3}}\pi ^2\pi _{\sigma \nu }\pi ^{\sigma \nu })+{\displaystyle \frac{1}{2G_q}}\mathrm{\Pi }^2+{\displaystyle \frac{X_q^{^{}}\pi \dot{\mathrm{\Phi }}}{3X_q}}{\displaystyle \frac{X_q^{^{}}𝒦\mathrm{\Pi }}{G_qN}}],`$ (A.6)
$``$ $``$ $`X_q+{\displaystyle \frac{1}{2}}G_q\gamma ^{\sigma \nu }_\sigma \mathrm{\Phi }_\nu \mathrm{\Phi }+V_q,`$ (A.7)
where $`^{}`$ denotes derivative with respect to $`\mathrm{\Phi }`$, and $`\dot{}`$ with respect to $`\rho `$.
The canonical momenta are defined by
$`\pi _{\sigma \nu }`$ $``$ $`{\displaystyle \frac{1}{\sqrt{det\gamma }}}{\displaystyle \frac{\delta L_q}{\delta \dot{\gamma }^{\sigma \nu }}}={\displaystyle \frac{X_q}{N}}(𝒦_{\sigma \nu }\gamma _{\sigma \nu }𝒦)\gamma _{\sigma \nu }X_q^{^{}}\dot{\mathrm{\Phi }},`$ (A.8)
$`\mathrm{\Pi }`$ $``$ $`{\displaystyle \frac{1}{\sqrt{det\gamma }}}{\displaystyle \frac{\delta L_q}{\delta \dot{\mathrm{\Phi }}}}=G_q\dot{\mathrm{\Phi }}+{\displaystyle \frac{2X_q^{^{}}𝒦}{N}}.`$ (A.9)
By inverting these equations, we obtain the flow equations
$`𝒦_{\sigma \nu }`$ $``$ $`{\displaystyle \frac{1}{2N}}\dot{\gamma }_{\sigma \nu }={\displaystyle \frac{N}{X_q}}[\pi _{\sigma \nu }{\displaystyle \frac{1}{3}}\gamma _{\sigma \nu }(\pi +X_q^{^{}}\dot{\mathrm{\Phi }})],`$ (A.10)
$`\dot{\mathrm{\Phi }}`$ $`=`$ $`F(\mathrm{\Pi }+{\displaystyle \frac{2X_q^{^{}}}{3X_q}}\pi ),F{\displaystyle \frac{1}{G_q\frac{8X_q^{}_{}{}^{}2}{3X_q}}},`$ (A.11)
and then substituting these two equation into (A.6), the Hamiltonian density $``$ can be expressed completely in terms of the canonical momenta.
Define the Hamiltonian-Jacobi functional as before
$$S[\gamma ,\mathrm{\Phi }]=\mathrm{\Gamma }[\gamma ,\mathrm{\Phi }]+𝑑x^4\sqrt{det\gamma }[Z_q(\mathrm{\Phi })+\frac{1}{2}M_q(\mathrm{\Phi })(\mathrm{\Phi })^2+U_q(\mathrm{\Phi })],$$
(A.12)
and solve the Hamiltonian-Jacobi equation and the flow equations by adopting the new energy scale parameter $`\mu _q`$ defined in (60). We obtain the beta function
$$\beta _q=\mathrm{\Omega }^1\beta =2\mathrm{\Phi }(\mathrm{\Phi }^21),\mathrm{\Omega }\frac{3}{5\mathrm{\Phi }^22},$$
(A.13)
and the renormalized dilatonic potential and coefficient of the scalar curvature
$$U_q=\mathrm{\Phi }^{10/3}U=\frac{10\mathrm{\Phi }^24}{\mathrm{}},Z_q=\mathrm{\Phi }(\mathrm{\Phi }^2+2)+\frac{Z_{q0}\mathrm{\Phi }}{\sqrt{\mathrm{\Phi }^21}}.$$
(A.14)
As mentioned in section IV, if we take $`Z_{q0}=Z_0`$, then $`Z_q=\mathrm{\Phi }^{5/3}Z`$, and the resulting $`c`$-functions transform as vectors on the $`\mathrm{\Phi }`$-space.
The formal expressions for the $`c`$-functions are somewhat involved:
$`c`$ $`=`$ $`{\displaystyle \frac{1}{T}}{\displaystyle \frac{Z_q^2}{X_q}},`$ (A.15)
$`d`$ $`=`$ $`{\displaystyle \frac{1}{T}}[{\displaystyle \frac{Z_q^2}{3X_q}}+{\displaystyle \frac{Z_q^{}_{}{}^{}2}{2G_q}}Z_qH{\displaystyle \frac{X_q^{^{}}Z_q^{^{}}Z_q}{3G_qX_q}}+{\displaystyle \frac{4X_q^{^{}}Z_q^{^{}}H}{G_q}}],`$ (A.16)
with
$`T`$ $``$ $`{\displaystyle \frac{U_q}{3X_q}}+{\displaystyle \frac{FX_q^{^{}}}{3X_q}}(U_q^{^{}}{\displaystyle \frac{8X_q^{^{}}U_q}{3X_q}})+{\displaystyle \frac{X_q^{^{}}U_q^{^{}}}{3G_qX_q}}+{\displaystyle \frac{8FX_q^{}_{}{}^{}3U_q^{^{}}}{9G_qX_q^2}},`$ (A.17)
$`H`$ $``$ $`{\displaystyle \frac{FX_q^{^{}}(Z_q^{^{}}\frac{2X_q^{^{}}Z_q}{3X_q})}{3X_q}}.`$ (A.18)
However, the final expressions are very simple
$$c_q=\mathrm{\Omega }^1c,d_q=\mathrm{\Omega }^1d,$$
(A.19)
as long as $`Z_{q0}=Z_0`$. Indeed, by continuity at $`\mathrm{\Phi }=1`$, we are forced to take $`Z_{q0}=Z_0=0`$.
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# A Nearly Minimum Redundant Correlator Interpolation Formula for Gravitational Wave Chirp Detection
## Introduction
The detection of gravitational wave (henceforth GW) chirps from unknown inspiraling compact binary sources (henceforth CBS) is a primary goal for the early operation of broadband interferometric detectors, including TAMA300 , GEO600 , the two LIGOs and VIRGO , in view of the sizeable expected rate of observable events .
For additive gaussian stationary noise, the correlator-bank threshold-detector is the optimal one, yielding the smallest false-dismissal probability, at any fixed false-alarm probability and signal to noise ratio .
The issues of optimum template parametrization and placement, and the related computational burden have been discussed by several Authors -. A lucid account of the main relevant landmarks is given in .
Curiously, the question of possible efficient interpolation among the correlators has been left yet unsolved .
In this paper we set up and test an efficient interpolated representation of the (reduced, noncoherent) correlator for the simplest paradigm case of newtonian chirps. The proposed representation is proven to get close to the absolute minimum template density required by a prescribed minimal-match condition, which follows from the theory of quasi-bandlimited (henceforth q-BL) functions in the $`L^{\mathrm{}}`$ norm. The statistical performance of the proposed representation are shown to be essentially equivalent to those of the plain lattice.
This paper is accordingly organized as follows. In Section I we recall a number of relevant concepts and results. In Section II we review the design of the plain template-bank for the simplest (newtonian) case, and discuss its statistical detection properties. In Section III we briefly introduce q-BL functions and cardinal expansions, and derive the proposed representations. In Section IV we compare the computational burden and the statistical detection/estimation properties of the cardinal-interpolated (reduced, noncoherent) correlator lattice to those of the plain lattice. Conclusions follow under Section V.
## I Background
In this section we resume a number of well known concepts relevant to CBS chirp detection, and introduce the notation.
### A Noncoherent Correlator. Deflection and SNR
Detecting GW chirps from unknown inspiraling CBS requires the computation of a suitable set of non-coherent correlators (henceforth NCC) :
$$c[\overline{T}]=2\left|_{f_{inf}}^{f_{sup}}\frac{A(f)\overline{T}^{}(f)}{\mathrm{\Pi }(f)}𝑑f\right|,$$
(1)
where $`(f_{inf},f_{sup})`$ is the useful antenna spectral window, $`A(f)=S(f)+N(f)`$ are the noise corrupted (spectral) data, resulting from the superposition of a (possibly null) signal $`S(f)`$ and a realization $`N(f)`$ of the antenna noise (assumed gaussian and stationary), $`\overline{T}(f)`$ is an element of a suitable set of unit-norm chirp-templates such that
$$\overline{T}=\left|2_{f_{inf}}^{f_{sup}}\frac{\overline{T}(f)\overline{T}^{}(f)}{\mathrm{\Pi }(f)}𝑑f\right|^{1/2}=1,$$
(2)
and $`\mathrm{\Pi }(f)`$ is the (one-sided) antenna noise power spectral density (henceforth PSD). The random variables $`c`$ have Ricean probability densities ,
$$w(c)=c\mathrm{exp}\left(\frac{c^2+d^2}{2}\right)I_0\left(cd\right),$$
(3)
where $`I_0()`$ is the modified Bessel function of first kind and zero order, and
$$d=\left|2_{f_{inf}}^{f_{sup}}\frac{S(f)\overline{T}^{}(f)}{\mathrm{\Pi }(f)}𝑑f\right|$$
(4)
is the deflection obtained using $`\overline{T}`$. The moments of (3) can be written in terms of Kummer’s confluent hypergeometric function :
$$c^n=2^{n/2}\mathrm{\Gamma }\left(1+\frac{n}{2}\right){}_{1}{}^{}F_{1}^{}(\frac{n}{2};1;\frac{d^2}{2}).$$
(5)
For $`d\stackrel{>}{}5`$ the riceans (3) merge into gaussians, with:
$$E\left[c\right]d,\text{var}\left[c\right]1.$$
(6)
The deflection attains its maximum value iff the template $`\overline{T}`$ is matched to the signal, viz.:
$$\overline{T}(f)=\frac{S(f)}{S(f)},$$
(7)
yielding:
$$d=d_{max}=|2_{f_{inf}}^{f_{sup}}\frac{S(f)S^{}(f)}{\mathrm{\Pi }(f)}df|^{1/2}=:SNR,$$
(8)
where SNR is the intrinsic signal-to-noise ratio.
### B Chirp Templates
The stationary phase principle (see for a thorough discussion of its validity) can be used to show that the asymptotic principal part of a general , , reduced post newtonian (henceforth PN) chirp can be written :
$$S(f;\varphi _c,T_c,\stackrel{}{\xi })=Af^{7/6}\mathrm{exp}\left\{j\left[2\pi fT_c\varphi _c+\psi (f,\stackrel{}{\xi })\right]\right\},$$
(9)
where $`A`$ is a constant (real, unknown) amplitude factor, $`T_c`$ is the (fiducial) coalescency time , $`\varphi _c`$ is the template phase at $`t=T_c`$, and $`\stackrel{}{\xi }`$ represents the remaining intrinsic source parameters .
Equation (9) is used to construct the needed chirp templates. A further suffix $`T`$ will be used to label the template parameters $`A_T`$, $`\varphi _{c_T}`$, $`T_{c_T}`$, and $`\stackrel{}{\xi }_T`$.
All template amplitudes $`A_T`$ will be chosen so as to comply with the normalization condition (2), viz.:
$$A_T=\left[2_{f_{inf}}^{f_{sup}}\frac{f^{7/3}}{\mathrm{\Pi }(f)}𝑑f\right]^{1/2}.$$
(10)
### C Maximum Likelihood Criterion. Fitting Factor
Equations (6) imply that (under the assumption of a uniform distribution of the unknown source parameters) the largest correlator will most likely correspond to the special template yielding the largest deflection (maximum likelihood (ML) estimation criterion ).
Data analysis for detecting chirps reduces thus to the following. Given the (spectral) noisy data, and a set (lattice) of templates, suitably covering the chirp parameter space, the corresponding (noncoherent) correlators $`\{c_k|k=1,2,\mathrm{},N\}`$ are computed. The largest among these correlators is used as a detection statistic , viz., whenever this latter exceeds a suitable threshold, set by the prescribed false-alarm probability (surveillance strategy ), a signal is declared to have been observed , and the corresponding template is taken as the most likely estimate of the observed signal .
It is convenient to measure the goodness of fit between a given signal $`S(f)`$ and the best available template in the set in terms of the so called fitting factor :
$$FF=\underset{k}{\mathrm{max}}\frac{d_k}{SNR}.$$
(11)
The set of templates should be constructed in such a way that for any admissible signal,
$$FF\mathrm{\Gamma },$$
(12)
where $`1\mathrm{\Gamma }^3`$ gauges the fraction of potentially observable sources which could be lost as an effect of template mismatch , .
### D The Reduced Correlator
In view of eq. (9), the noncoherent correlator (1) can be written:
$$c=\frac{\left|2{\displaystyle _{f_{inf}}^{f_{sup}}}{\displaystyle \frac{A(f)f^{7/6}e^{j\psi _T(f,\stackrel{}{\xi })]}}{\mathrm{\Pi }(f)}}\mathrm{exp}(j2\pi fT_{c_T})𝑑f\right|}{\left[2{\displaystyle _{f_{inf}}^{f_{sup}}}{\displaystyle \frac{f^{7/3}}{\mathrm{\Pi }(f)}}𝑑f\right]^{1/2}}.$$
(13)
Equation (13) is formally (but for inessential factors) the absolute value of the ($`fT_{c_T}`$) Fourier transform of the (complex-valued) function:
$$K(f)=\{\begin{array}{c}\frac{A(f)\overline{T}^{}(f;0,0,\stackrel{}{\xi })}{\mathrm{\Pi }(f)},f_{inf}ff_{sup}\hfill \\ 0,f<f_{inf},f>f_{sup}.\hfill \end{array}$$
(14)
Maximizing the noncoherent correlator (13) w.r.t. $`T_{c_T}`$, is thus equivalent to taking the largest absolute value of the ($`fT_{c_T}`$) Fourier transform of (14). The resulting reduced correlator will be denoted with a capital letter, viz.:
$$C=\underset{T_{c_T}}{sup}c.$$
(15)
### E The Newtonian Deflection and the Match
For illustrative purposes, in this paper we shall restrict to the simplest Newtonian ($`0PN`$) signals and templates. The $`0PN`$ function $`\psi _T(f;\stackrel{}{\xi })`$ in (9) reads:
$$\psi _T(f)=\frac{3}{128}\left(\frac{\pi G}{c^3}\right)^{5/3}_T^{5/3}f^{5/3},$$
(16)
where $`_T`$ is the template chirp-mass. It is convenient to introduce the following dimensionless variables and parameters:
$`\overline{f}={\displaystyle \frac{f}{f_{inf}}},\mathrm{\Theta }=f_{inf}(T_cT_{c_T}),\overline{}={\displaystyle \frac{}{M_{}}},`$
$$\mathrm{\Delta }=\overline{}_s^{5/3}\overline{}_T^{5/3},\mathrm{\Lambda }=\frac{3}{128}\left(\frac{\pi Gf_{inf}M_{}}{c^3}\right)^{5/3},$$
(17)
where $`M_{}`$ is the solar mass, so as to recast the deflection (4) into the form:
$$d(\mathrm{\Delta },\mathrm{\Theta })=SNR\frac{|{\displaystyle _1^{\overline{f}_{sup}}}d\overline{f}{\displaystyle \frac{\overline{f}^{7/3}}{\mathrm{\Pi }(\overline{f})}}\mathrm{exp}\left[j(2\pi \mathrm{\Theta }\overline{f}+\mathrm{\Lambda }\overline{f}^{5/3}\mathrm{\Delta }]\right|}{{\displaystyle _1^{\overline{f}_{sup}}}𝑑\overline{f}{\displaystyle \frac{\overline{f}^{7/3}}{\mathrm{\Pi }(\overline{f})}}}.$$
(18)
It is also useful to introduce the reduced deflection:
$$D(\mathrm{\Delta })=\underset{\mathrm{\Theta }}{\mathrm{max}}d(\mathrm{\Delta },\mathrm{\Theta }),$$
(19)
and the related normalized functions:
$$\overline{d}()=\frac{d()}{SNR},\overline{D}()=\frac{D()}{SNR}.$$
(20)
The function $`\overline{D}(\mathrm{\Delta })`$ is known as the (newtonian) match. The quantity $`\mathrm{\Gamma }`$ in (12) is accordingly also named the minimal match .
The functions $`\overline{d}`$ and $`\overline{D}`$ are displayed in Fig.s 1 and 2, respectively, for the special case of a LIGO-like noise PSD,
$$\mathrm{\Pi }(f)=\frac{\mathrm{\Pi }_0}{5}\left\{\left(\frac{f_0}{f}\right)^4+2\left[1+\left(\frac{f}{f_0}\right)^2\right]\right\},$$
(21)
with $`f_0=300Hz`$ ($`\mathrm{\Pi }_0`$ is a constant of no concern to us here), and for a spectral window with:
$$f_{inf}=40Hz,f_{sup}=400Hz.$$
(22)
The value $`\mathrm{\Theta }_{max}`$ of $`\mathrm{\Theta }`$ which maximizes $`\overline{d}(\mathrm{\Delta },\mathrm{\Theta })`$ in a neighbourhood of $`\mathrm{\Delta }=0`$ is shown in Fig. 3 as a function of $`\mathrm{\Delta }`$.
## II The Plain (Newtonian) Lattice
Given a range $`[_{min},_{max}]`$ of allowed source chirp masses, let the set of template chirp-masses be:
$$\overline{}_1^{5/3}=\overline{}_{max}^{5/3},\overline{}_{n+1}^{5/3}=\overline{}_n^{5/3}+\delta _L,n=1,2,\mathrm{},N_L1,$$
(23)
where $`\delta _L`$ is the lattice-spacing, and
$$N_L=\frac{\overline{}_{min}^{5/3}\overline{}_{max}^{5/3}}{\delta _L}$$
(24)
is the total number of templates. Obviously $`\mathrm{\Delta }`$ can take only the discrete values:
$$\mathrm{\Delta }_n=\overline{}_s^{5/3}\overline{}_n^{5/3},n=1,2,\mathrm{},N_L.$$
(25)
### A Lattice Design
Given $`\overline{}_s`$, the fitting factor of the lattice is:
$$FF=\underset{n}{\mathrm{max}}\overline{D}(\mathrm{\Delta }_n)=\overline{D}\left(\underset{n}{\mathrm{min}}|\mathrm{\Delta }_n|\right).$$
(26)
The minimal-match condition (12) should be obviously enforced in the worst case, where:
$$\overline{}_s^{5/3}=\overline{}_q^{5/3}+\frac{\delta _L}{2},$$
(27)
yielding:
$$\overline{D}\left(\frac{\delta _L}{2}\right)=\mathrm{\Gamma }.$$
(28)
Equation (28) uniquely determines the lattice spacing $`\delta _L`$, and hence via (24) also the total number of templates.
### B False Alarm and False Dismissal Probabilities
The statistical distribution of the lattice detection statistic
$$\underset{k}{\mathrm{max}}C_k=\underset{h}{\mathrm{max}}c_h$$
(29)
is not known in exact form, and one should resort to numerical simulations aided by intuition to compute the lattice false-alarm and false-dismissal probabilities.
The detection threshold $`\gamma `$ is determined from the prescribed (tolerated) false alarm probability $`\alpha `$, by solving the equation:
$$\alpha =\text{prob}[k:C_k>\gamma |SNR=0]=1\text{prob}[h,c_h<\gamma |SNR=0].$$
(30)
The joint probability in (30) is difficult to compute, since the $`c_h`$ are not statistically independent in general. For most practical purposes, a decent approximation is:
$$\text{prob}[h,c_h<\gamma |SNR=0]=(1e^{\gamma ^2/2})^M,$$
(31)
which would be appropriate if the $`c_h`$ were a collection of $`M`$ independent ricean random variables. Numerical experiments suggest an almost linear dependence of $`M`$ on the total number of NCC used .
The probability of false dismissal of a signal with $`SNR0`$ is:
$$\beta (\gamma ,SNR)=\text{prob}[\underset{k}{\mathrm{max}}C_k<\gamma |SNR0]=\text{prob}[h,c_h<\gamma |SNR0].$$
(32)
A simple (conservative) approximation of (32) is , :
$$\beta (\gamma ,SNR)\text{prob}[C_{}<\gamma ,C_+<\gamma |SNR0],$$
(33)
where $`C_{}`$, $`C_+`$ denote the reduced correlators corresponding to the nearest-neighbouring templates, with chirp masses $`\overline{}_\pm `$ such that $`\overline{}_s[\overline{}_{},\overline{}_+]`$. Under the same (reasonably large SNR) assumptions leading to (33), the involved joint probability density can be approximated by a gaussian bivariate :
$$w(C_{},C_+)\frac{\mathrm{exp}\left[{\displaystyle \frac{(C_{}d_{})^2+(C_+d_+)^22R(C_{}d_{})(C_+d_+)}{2(1R^2)}}\right]}{2\pi (1R^2)^{1/2}},$$
(34)
where:
$$d_\pm =d[\mathrm{\Delta }_\pm ,\mathrm{\Theta }_{max}(\mathrm{\Delta }_\pm )],R\overline{d}[\delta _L,\mathrm{\Theta }_{max}(\mathrm{\Delta }_+)\mathrm{\Theta }_{max}(\mathrm{\Delta }_{})],$$
(35)
and the function $`\mathrm{\Theta }_{max}()`$ has been defined in Sect. II.F and shown in Fig. 3.
Letting:
$$\overline{}_s^{5/3}=\overline{}_q^{5/3}+\eta \delta _L,\eta [0,1[,$$
(36)
the false dismissal probability (33) is obviously a function of $`\eta `$. Within the limits of validity of (30) to (35), for a fixed $`SNR`$ and a prescribed $`\alpha `$, the following qualitative dependence of $`\beta `$ on the lattice spacing $`\delta _L`$ is observed. In a neighbourhood of $`\eta =0.5`$, the false dismissal probability is reduced by reducing the spacing $`\delta _L`$ among the templates. On the other hand, in a neighbourhood of $`\eta =0`$, reducing the template spacing produces an increase of $`\beta `$. This is due to the dominant effect of the parallel increase of $`\gamma `$, needed to keep $`\alpha `$ unchanged.
A judicious tradeoff should be obviously sought, to choice a value of $`\mathrm{\Gamma }`$, and hence of $`\delta _L`$, via (28), which minimizes, e.g., the average value of $`\beta `$ w.r.t $`\eta `$, under the assumption of a uniform distribution of the source chirp mass.
## III A Nearly Minimum Redundant Interpolated Lattice
This section contains the main new results. A short introduction to the theory of q-BL functions is included to make the paper self-contained.
### A q-BL Functions and Cardinal Expansions
A function $`f:x[a,b]`$ is q-BL in the $`L^{\mathrm{}}`$ norm iff :
$$\gamma ,B_c^+:\underset{x[a,b]}{sup}\left|f(x)f_B(x)\right|=\mathrm{exp}[\gamma (BB_c)],$$
(37)
$`f_B(x)`$ being obtained by taking the inverse Fourier transform of the spectrum of $`f(x)`$ chopped at $`|y|B`$, viz.:
$$f_B(x)=_{yx}^1\left\{W\left(\frac{y}{B}\right)_{xy}\left[f(x)\right]\right\},$$
(38)
where:
$$W(x)=\{\begin{array}{c}1,|x|1\hfill \\ 0,|x|>1.\hfill \end{array}$$
(39)
For a strictly bandlimited function $`f(x)`$, whose spectrum vanishes identically outside $`[B,B]`$, eq. (38) provides an exact interpolating representation known as cardinal expansion , :
$$f_B(x)=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}f(x_n)\text{sinc}\left[\frac{\pi }{\delta }(xx_n)\right],x_{n+1}x_n=\delta ,\delta =\frac{1}{2B},$$
(40)
where $`\text{sinc}(x)=\mathrm{sin}(x)/x`$. For a q-BL function on the other hand, one can prove that eq. (40), while reproducing exactly $`f(x)`$ at $`x=x_k,k𝒩`$, satisfies eq. (37), i.e. that
$$ϵ>0,B:\underset{x[a,b]}{sup}\left|f(x)f_B(x)\right|<ϵ.$$
(41)
Equation (40) is an approximate sample-interpolating representation, where the sample density $`\delta ^1=2B`$ depends on the prescribed $`L^{\mathrm{}}`$ approximation error $`ϵ`$. It is important to note that the exponential decay of the error in (37) implies that reducing $`ϵ`$ in (41) by orders of magnitude does not change the order of magnitude of $`B`$.
Usually one needs to compute $`f(x)`$ in a finite interval $`[a,b]`$ including only
$$N=\frac{(ba)}{\delta }$$
(42)
samples . However, using eq. (40) to compute $`f(x)`$ in $`[a,b]`$ requires, in principle, knowledge of an infinite number of samples outside the interval of interest. This limitation can be circumvented by using generalized (economized) cardinal expansions. These expansions have the general form :
$`f(x)={\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}f(x_n)\text{sinc}\left[{\displaystyle \frac{\pi }{\delta ^{}}}(xx_n)\right]\theta (xx_n),`$
$$x_{n+1}x_n=\delta ^{},\delta ^{}=(2\chi B)^1,\chi >1$$
(43)
where $`\theta (x)`$ is a suitable windowing function such that:
$$\theta (0)=1,_{x\xi }\left[\theta (x)\right]=0,\xi >(\chi 1)B,\chi >1.$$
(44)
The expansion (43) is nothing but the std. cardinal expansion of the function $`f(x)\theta (ux)`$, whose bandwidth under the assumptions (44) is $`B^{}=\chi B`$, viz.:
$`f(x)\theta (ux)={\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}f(x_n)\theta (ux_n)\text{sinc}\left[{\displaystyle \frac{\pi }{\delta ^{}}}\left(xx_n\right)\right],`$
$$x_{n+1}x_n=\delta ^{},\delta ^{}=\frac{1}{2\chi B},$$
(45)
evaluated at $`u=x`$. The Fourier spectrum of $`f(x)\theta (ux)`$ is a smoothed version of the plain spectrum of $`f(x)`$; a judicious choice of $`\theta (x)`$ can thus make in principle, the decay rate of $`f(x_n)\theta (ux_n)`$ as $`|ux_n|\mathrm{}`$ as fast as desired .
The Knab window function :
$$\theta (x)=K_P(x):=\frac{\text{sinh}\left\{\pi P(1\chi ^1)\left[1\left({\displaystyle \frac{x}{P\delta ^{}}}\right)^2\right]^{1/2}\right\}}{\left[1\left({\displaystyle \frac{x}{P\delta ^{}}}\right)^2\right]^{1/2}\text{sinh}[\pi P(1\chi ^1)]}$$
(46)
satisfies all constraints (44) and is essentially confined in $`|x|P\delta ^{}`$. This allows to truncate (43) at $`|xx_n|P\delta ^{}`$, so that for any given $`x`$, only $`2P`$ samples symmetrically placed around $`x`$ are essentially needed to reconstruct $`f(x)`$.
The error resulting from truncation of (43) with (46) at $`|xx_n|=P\delta ^{}`$ has been discussed in . A simple (and conservative) upper bound is given by:
$`\left|{\displaystyle \underset{|xx_n|>P\delta ^{}}{}}f(x_n)\text{sinc}\left[{\displaystyle \frac{\pi }{\delta ^{}}}(xx_n)\right]K_P(xx_n)\right|<`$
$$<\frac{M}{\text{sinh}[\pi P(1\chi ^1)]},M=\underset{x[a,b]}{sup}f(x).$$
(47)
Usually, one enforces the condition:
$$\frac{M}{\text{sinh}[\pi P(1\chi ^1)]}=ϵ^{}ϵ,$$
(48)
where $`ϵ`$ is the prescribed $`L^{\mathrm{}}`$ error in (41). Equation (48) can be solved to express $`P`$ as a function of $`\chi `$,
$$P=\frac{\text{sinh}^1\left(M/ϵ^{}\right)}{\pi (1\chi ^1)}.$$
(49)
The total number of samples
$$N_T=\frac{(ba)}{\delta ^{}}+2P=\chi \frac{(ba)}{\delta }+2P$$
(50)
needed to represent $`f(x)`$ in $`[a,b]`$ using (45) and (46), within a prescribed $`L^{\mathrm{}}`$ error and under the constraint (48), can be accordingly minimized by letting:
$$\chi =1+\left[\frac{2\delta \text{sinh}^1(M/ϵ^{})}{\pi (ba)}\right]^{1/2}.$$
(51)
### B The Cardinal-Interpolated Newtonian Match
The match $`\overline{D}(\mathrm{\Delta })`$ is a q-BL function in the $`L^{\mathrm{}}`$ norm. This can be seen from Fig. 4, where the exponential decay of the $`L^{\mathrm{}}`$ error in $`[0,\mathrm{}[`$ between $`\overline{D}(\mathrm{\Delta })`$ and the cardinal expansion:
$$\overline{D}_B(\mathrm{\Delta })=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\overline{D}(\mathrm{\Delta }_n)\text{sinc}\left[\frac{\pi }{\delta _C}\left(\mathrm{\Delta }\mathrm{\Delta }_n\right)\right],\mathrm{\Delta }_{n+1}\mathrm{\Delta }_n=\delta _C,$$
(52)
is displayed as a function of $`\delta _C^1`$ on a Log-Lin plot . Switching back to the original variables, eq. (52) reads:
$$\overline{D}_B(\overline{}_s^{5/3}\overline{}_T^{5/3}):=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\overline{D}(\overline{}_s^{5/3}\overline{}_n^{5/3})\text{sinc}\left[\frac{\pi }{\delta _C}\left(\overline{}_T^{5/3}\overline{}_n^{5/3}\right)\right],$$
(53)
where :
$$\overline{}_{n+1}^{5/3}\overline{}_n^{5/3}=\delta _C.$$
(54)
Given $`\overline{}_s`$, the fitting factor obtained using (53) is given by:
$$FF=\underset{\overline{}_T}{\mathrm{max}}\overline{D}_B(\overline{}_s^{5/3}\overline{}_T^{5/3})=:\overline{D}_B(\overline{}_s^{5/3}\overline{}_{}^{5/3}).$$
(55)
It is convenient to let:
$$\overline{}_s^{5/3}=\overline{}_q^{5/3}+\eta \delta _C,\eta [0,1[,$$
(56)
$$\overline{}_{}^{5/3}=\overline{}_q^{5/3}+\eta _{}\delta _C,\eta _{}[0,1[,$$
(57)
so as to rewrite (55) as:
$$FF=\overline{D}_B(\eta \eta _{}).$$
(58)
The difference $`\eta \eta _{}`$ turns out to depend on $`\eta `$ as shown in Fig. 5, and hence the fitting factor (58) depends on $`\eta `$ as shown in Fig. 6. The minimal-match condition (12) should again be enforced in the worst case(s), i.e., as seen from Fig.s 5, 6, for $`\eta =0.5`$ and $`\eta _{}=\eta `$, yielding:
$$\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\overline{D}\left[\left(n+\frac{1}{2}\right)\delta _C\right]\text{sinc}\left[\left(n+\frac{1}{2}\right)\pi \right]=\mathrm{\Gamma }.$$
(59)
This condition is notably indipendent on $`q`$, and fixes the sample spacing $`\delta _C`$.
In practice, as discussed in the previous section, it is convenient to use an economized cardinal expansion, viz:
$$\overline{D}_B(\overline{}_s^{5/3}\overline{}_T^{5/3}):=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\overline{D}(\overline{}_s^{5/3}\overline{}_n^{5/3})\mathrm{\Psi }_n\left(\overline{}_T^{5/3}\overline{}_n^{5/3}\right),$$
(60)
where (see eq.s (43) and (46)):
$$\overline{}_{n+1}^{5/3}\overline{}_n^{5/3}=\delta _C^{}=\chi ^1\delta _C,$$
(61)
and:
$$\mathrm{\Psi }_n(x)=\text{sinc}\left(\frac{\pi x}{\delta _C^{}}\right)\frac{\text{sinh}\left\{\pi P(1\chi ^1)\left[1\left({\displaystyle \frac{x}{P\delta _C^{}}}\right)^2\right]^{1/2}\right\}}{\left[1\left({\displaystyle \frac{x}{P\delta _C^{}}}\right)^2\right]^{1/2}\text{sinh}[\pi P(1\chi ^1)]}.$$
(62)
As shown in the previous section, the function (62) is essentially contained in the interval: $`|x|<P\delta _C^{}`$, and hence the infinite sum in (60) is essentially restricted to:
$$\frac{\overline{}_T^{5/3}}{\delta _C^{}}Pn\frac{\overline{}_T^{5/3}}{\delta _C^{}}+P.$$
(63)
Capitalizing on eq. (47) we shall enforce the condition:
$$\frac{1}{\text{sinh}[\pi P(1\chi ^1)]}=\frac{1\mathrm{\Gamma }}{10}$$
(64)
to guarantee that the minimal-match condition will not be affected within the last significant figure of $`\mathrm{\Gamma }`$, when using using eq. (60) truncated according to (63) in place of eq. (53). Further, in view of eq.s (49) and (51), we shall take:
$$\chi =\chi _{opt}=:1+\left[\frac{2\delta _C\text{sinh}^1\left({\displaystyle \frac{10}{1\mathrm{\Gamma }}}\right)}{\pi \left(\overline{}_{min}^{5/3}\overline{}_{max}^{5/3}\right)}\right]^{1/2}$$
(65)
and:
$$P=P_{opt}=:\frac{\text{sinh}^1\left({\displaystyle \frac{10}{1\mathrm{\Gamma }}}\right)}{\pi (1\chi _{opt}^1)},$$
(66)
so as to minimize the total number of correlators
$$N_C=\chi \frac{\overline{}_{min}^{5/3}\overline{}_{max}^{5/3}}{\delta _C}+2P.$$
(67)
needed to evaluate (60) throughout the range $`[_{min},_{max}]`$ of $`_T`$.
### C The Cardinal-Interpolated Reduced Correlator
As a next step, we make the ansatz that an approximate representation of the reduced (newtonian) noncoherent correlator in terms of a (generalized) cardinal expansion also holds, viz. :
$$C_B:=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}C_n\mathrm{\Psi }_n(\overline{}_T^{5/3}\overline{}_n^{5/3}),$$
(68)
where:
$$C_n=\underset{T_{c_T}}{\mathrm{max}}\left|2_{f_{inf}}^{f_{sup}}\frac{A(f)\overline{T}_n^{}(f)}{\mathrm{\Pi }(f)}𝑑f\right|,$$
(69)
and the infinite sum is truncated according to (63).
In (69) the templates $`\overline{T}_n`$ are defined by eq.s (9), (10) and (16), where the (scaled) chirp masses take the values
$$\overline{}_k^{5/3}=\overline{}_{max}^{5/3}+k\delta _C^{},k=P,P+1,\mathrm{},N_C1+P,$$
(70)
the interpolating functions $`\mathrm{\Psi }_n(x)`$ are given by (62), and the parameters $`\delta _C`$, $`\chi `$ and $`P`$ are computed from the prescribed minimal match $`\mathrm{\Gamma }`$ as explained In Sect. III.A.
## IV Plain vs. Cardinal Interpolated Lattice
In this section we shall compare the (newtonian) cardinal-interpolated (reduced, noncoherent) correlator to the plain-lattice of (reduced noncoherent) correlators in terms of computational cost and statistical features (detection and estimation performance). The assumed noise PSD and spectral window are given by (21) and (22), respectively.
### A Computational Burden
The plain lattice template spacings $`\delta _L`$ and total number of correlators $`N_L`$ needed to cover the range $`(0.2M_{},10M_{})`$ for some values of the minimal match $`\mathrm{\Gamma }`$ are compared in Table I to the corresponding quantities $`\delta _C`$ and $`N_C`$ of the cardinal interpolated correlator.
It is seen that at any value of the minimal match $`\mathrm{\Gamma }`$, the cardinal-interpolated representation requires some $`30\%`$ less many templates than the plain lattice.
On the other hand, evaluating (68) at any value $`_T_k`$ is substantially cheaper than computing the corresponding (reduced, noncoherent) correlator.
Indeed, to use (68) one really needs to evaluate the interpolating functions $`\mathrm{\Psi }_n(x)`$ only at a finite number of (equispaced) values of $`\overline{}_T^{5/3}`$ between the samples $`\overline{}_q^{5/3}`$. The corresponding values of the interpolating functions can be computed once for all, and stored in a look-up table. As a result, only $`2P`$ floating point operations are needed to compute (68) at each of the above values of $`\overline{}_T^{5/3}`$, with a typical $`P10^2`$.
### B Statistical Features
The statistical properties of the cardinal-interpolated correlator have been compared to those of the plain-lattice via extensive Monte Carlo simulations. The number of different realizations used to derive the statistics was $`10^4`$.
Simulated data were sampled in time at twice the Nyquist rate. To limit running times, the minimal-match was set at $`\mathrm{\Gamma }=0.9`$, and the chirp-mass range was chosen in such a way that the longest observable waveform spanned $`2^{15}`$ time bins. In order to avoid circular-correlation artifacts, and to have equal statistics for all reduced correlators, all templates were zero-padded up to a total length of $`2^{16}`$ bins. Gaussian uncorrelated noise samples were generated using a feedback-shift-register routine from the IMSL package, featuring an extremely large period , followed by a Box-Müller transformation . The noise samples were added to the whitened data in the spectral domain.
In the following we shall denote the cardinal-interpolated and plain lattice test-statistics as :
$$C_{max}^{(C)}=\underset{_𝒯}{\mathrm{max}}\underset{k}{}C_k\mathrm{\Psi }_k\left(\overline{}_T\overline{}_k\right)$$
(71)
and:
$$C_{max}^{(L)}=\underset{k}{\mathrm{max}}C_k.$$
(72)
respectively. We shall denote the estimated mass, i.e., the value of $`\overline{}_T`$ which yields the maximum in (71) or (72) as $`_{est}`$, and let:
$$\overline{}_{est}^{5/3}=\overline{}_q^{5/3}+\eta _{est}\delta ,\eta _{est}[0,1[,$$
(73)
$$\overline{}_s^{5/3}=\overline{}_q^{5/3}+\eta \delta ,\eta [0,1[,$$
(74)
Whenever needed a suffix/superfix $`L,C`$ will be used to identify the plain-lattice and cardinal-interpolated cases in (73), (74).
The CDFs of $`C_{max}^{(C)}`$ (dashed lines) and $`C_{max}^{(L)}`$ (full lines) in the absence of signal ($`SNR=0`$), are compared in Fig. 7-a, for template spacings corresponding to $`\mathrm{\Gamma }=0.9`$. The corresponding PDFs are displayed in Fig. 7-b. The observed difference falls within the $`3\sigma `$ uncertainty interval related to the finite number of realizations.
The CDFs of $`C_{max}^{(C)}`$ (dashed lines) and $`C_{max}^{(L)}`$ (full lines) in the presence of a signal with $`SNR=6,8,10`$ are shown in Fig. 8-a, for the (worst) case where $`\eta =0.5`$ in (74). The corresponding PDFs are displayed in Fig. 8-b. Again, the observed differences fall within the $`3\sigma `$ uncertainty interval related to the finite number of realizations. Note that the expected value always exceeds the design value $`\mathrm{\Gamma }SNR`$, as might be expected as an effect of the supremum-taking operations in (15), (71) and (72).
As $`\eta `$ in (74) changes between $`0`$ and $`0.5`$, the PDFs of $`C_{max}^{(C)}`$ and $`C_{max}^{(L)}`$ change in turn. The limiting PDFs corresponding to $`\eta =0`$ and $`\eta =0.5`$ are shown in Fig.s 9 and 10, for the special case $`SNR=8`$, for of $`C_{max}^{(C)}`$ and $`C_{max}^{(L)}`$.
It can be concluded that the detection performance of the cardinal interpolated (reduced, noncoherent) correlator is essentially equivalent to that of the computationally more expensive plain lattice of (reduced, noncoherent) correlators.
We turn now to a comparison of the pertinent estimation features. To this end we let:
$$\xi =\frac{\overline{}_{est}^{5/3}\overline{}_s^{5/3}}{\delta }=\eta \eta _{est}.$$
(75)
The PDFs of $`\xi `$ for the cardinal-interpolated correlator at $`SNR=8`$ and $`\mathrm{\Gamma }=0.9`$ is shown in Fig. 11 for $`\eta =0`$ and $`\eta =0.5`$.
The corresponding probabilities $`P(\xi )`$ for the plain lattice of correlators are shown in Fig. 12. Note that for the plain lattice of correlators, $`\eta _{est}`$ can take only values which are integer multiples of $`\delta _L`$.
Both the cardinal-interpolated correlator and the plain lattice of correlators provide biased estimates. The bias $`E[\xi ]`$ becomes for both essentially independent of the $`SNR`$ at sufficiently high $`SNR`$ levels ($`SNR\stackrel{>}{}8`$). For the cardinal-interpolated correlator the asymptotic large-SNR bias is shown in Fig. 13, which is a close akin of Fig. 5. For the plain lattice, it is displayed in Fig. 14. The cardinal-interpolated correlator is seen to exhibit a smaller bias.
The std. deviations of the cardinal-interpolated and plain lattice estimators are nearly the same. For instance, for $`\eta =0.5`$ (worst case), at $`SNR=8`$ one has $`\sigma (\xi _L)=1.12`$ and $`\sigma (\xi _C)=.92`$; at $`SNR=10`$ $`\sigma (\xi _L)=0.84`$ and $`\sigma (\xi _C)=0.63`$.
### C Extension to PN Models
The cardinal-interpolated approach can be extended in principle, to higher order PN models, provided the structure and q-BL properties of the (reduced) correlator are preserved. In the geometrical language of , this is equivalent to requiring that the chosen parameter space be (globally) flat and Euclidean . This is surely the case for 1PN models , and almost the case for the new $`0`$spin 2PN coordinates proposed in . One should expect an even more substantial computational saving, in view of the higher dimension of the parameter space.
## V Conclusions and Recommendations
Quasi-bandlimited function approximation theory can be used to build a (nearly) minimum redundant cardinal-interpolated representations of the noncoherent correlator for detecting gravitational wave chirps. An explicit expression has been provided and tested, for the simplest case of newtonian waveforms.
The number of correlators to be computed and interpolated in order to maintain the match above a given minimal value $`\mathrm{\Gamma }`$ has been shown to be substantially less then required by the std. (lattice) approach, and the computational gain goes up with $`\mathrm{\Gamma }`$. On the other hand, evaluating the cardinal-interpolated representation at any value of $`_T`$ is substantially cheaper than computing the corresponding correlator.
We suggest that cardinal-interpolated expansions could be used to improve the efficiency of hierarchical searches, at all hierarchical levels.
Extension to PN templates should be straighforward, in principle, insofar as the structure and q-BL property of the correlators is preserved, and lead to an even increased computational gain. Work in this direction is in progress.
## VI Acknowledgements
This work has been sponsored in part by the European Community through a 1998 Senior Visiting Scientist Grant to I.M. Pinto at NAO - Spacetime Astronomy Division, Tokyo, Japan, in connection with the TAMA project. I.M. Pinto wishes to thank all the TAMA staff at NAO, and in particular prof. Fujimoto Masa-Katsu and prof. Kawamura Seiji for kindly hospitality and stimulating discussions.
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# Search for the Identification of 3EG J1835+5918: Evidence for a New Type of High-Energy Gamma-ray Source
## 1 Introduction
One of the most important advances in high-energy astrophysics in recent years is the discovery of 271 persistent high energy $`\gamma `$-ray sources by the EGRET instrument aboard the Compton Gamma-ray Observatory (CGRO, Hartman et al. 1999). While the detection of these sources is a major success, identification of their nature and origin has turned out to be a more challenging task. The principal method of identification, which relies on statistical evidence that blazars are the dominant population, is to find positional coincidences between EGRET sources and flat-spectrum radio/millimeter sources (Thompson et al. 1995, 1996; Mattox et al. 1997; Bloom et al. 1997). By definition blazars are flat-spectrum, radio-loud AGNs with polarized and variable optical emission. Although numerous efforts have been made at various wavelengths, only about one third of all EGRET sources have been identified with any degree of confidence. On the latest count these identifications include 66 blazars, i.e., flat-spectrum radio quasars or BL Lac objects (Hartman et al. 1999), seven rotation-powered pulsars (Hartman et al. 1999, Kaspi et al. 2000, Ramanamurthy et al. 1995), the nearby radio galaxy Cen A, and the Large Magellanic Cloud. Therefore approximately 196 EGRET sources remain unidentified with roughly half of these located at high Galactic latitude, $`b>10^{}`$.
Many difficulties attend the identification of EGRET sources close to the Galactic plane, but even at high Galactic latitude, the size of the typical error circle and the lack of a tight relation between gamma-ray flux and other properties such as X-ray flux and core radio flux prevent all but the brightest counterparts from being identified securely on the basis of position alone. The absence of obvious counterparts also admits the possibility that there is another population with characteristics unlike the identified EGRET sources. We have decided to explore the latter possibility by means of detailed work at other wavelengths, while in the long term the situation should improve considerably with the next generation high-energy $`\gamma `$-ray mission GLAST, which will produce more precise source locations.
We have chosen for a case study the unidentified EGRET source 3EG J1835+5918. This object may be the best candidate for the prototype of a new population different from blazars or pulsars. It is the brightest of the as-yet unidentified EGRET sources at high Galactic latitude ($`\mathrm{},b=89^{},25^{}`$), and the one with the smallest error circle. Because it is strongly detected and well away from the confusing diffuse emission in the Galactic plane, 3EG J1835+5918 is localized to within a radius of only $`12^{}`$ at 99% confidence, which makes a deep multiwavelength search for a counterpart feasible. The latest analysis of the EGRET observations of 3EG J1835+5918 leads to the conclusion that it shows no strong evidence for variability (Reimer et al. 2000). Its spectrum can be fitted by a power law of photon index –1.7 from 70 MeV to 4 GeV, with a turndown above 4 GeV. Such temporal and spectral behavior is more consistent with a rotation-powered pulsar than a blazar. Unlike 3EG J1835+5918, blazars are highly variable, and exhibit steeper spectra.
Prior to the observations reported herein, there were no known active galactic nuclei (AGNs) or pulsars in the error circle of 3EG J1835+5918. Examination of existing catalogs finds no flat-spectrum radio source (Mattox et al. 1997), no 1.4 GHz radio source of any type brighter than 4 mJy in the NRAO-VLA Sky Survey catalog (NVSS, Condon et al. 1998), and no 4.85 GHz source brighter than 20 mJy (Becker, White, & Edwards 1991). Observations by Nice & Sayer (1997) find and no radio pulsar to an upper limit of 1 mJy at 770 MHz. Furthermore, all of the known gamma-ray blazars and pulsars appear brighter in X-rays than the upper limit that we shall present for 3EG J1835+5918. In light of these facts, 3EG J1835+5918 cannot be a blazar unless it is a radio-quiet one (requiring a redefinition of this concept), nor a pulsar unless, as we shall show, it is one with unprecedented characteristics.
In this paper we present the results of radio, X-ray, and optical observations of the location of 3EG J1835+5918. The outline of the paper is as follows: §2 describes our multiwavelength data acquisition and selection techniques. §3 describes the optical spectroscopy of candidates and the overall results. §4 details notable properties of individual objects and assesses their prospects as the identification of 3EG J1835+5918. Multiwavelength comparisons with known $`\gamma `$-ray sources are addressed in §5, and the implications and conclusions of our work are discussed in §§6 and 7.
## 2 Observations
### 2.1 Optical Photometry and QSO Candidate Selection
The principal body of optical data for this study is a series of standard $`UBV`$ and Cousins $`R`$ CCD images of the error circle of 3EG J1835+5918 which we obtained using the MDM Observatory 1.3m telescope during a photometric run in 1998 June and July. A thinned, back-illuminated $`2048\times 2048`$ SITe CCD was used to cover a $`17^{}\times 17^{}`$ field with multiple exposures. A mosaic of four such overlapping fields enabled us to observe a $`32^{}\times 32^{}`$ region centered on the most likely EGRET source position (B. Dingus, private communication). Our images thus cover the entire 99% confidence region specified in the Third EGRET Catalog (Hartman et al. 1999), which can be approximated as an ellipse of major axis $`24^{}`$. In 1997 July we had covered the same field in the $`V`$ and $`I`$ bands only, and all of the $`V`$-band images were used to search for variability on long (year) and short (hours to days) time scales. The images were processed using standard IRAF/DAOPHOT procedures. Approximately 5000 objects were measured inside a $`15^{}`$ radius circle. The photometry described here was calibrated using Landolt standard stars (Landolt 1992). Typical limiting detections achieved were $`U=22.1,B=23.4,V=22.5`$, and $`R=22.5`$. Galactic extinction in this field is small but not negligible; Schlegel, Finkbeiner, & Davis (1998) give $`E(BV)=0.045`$, corresponding to $`A_U=0.25`$, $`A_B=0.20`$, $`A_V=0.15`$, and $`A_R=0.12`$. Magnitudes quoted in this paper are observed, i.e., not corrected for extinction.
We derived a list of QSO candidates from this photometry using the standard ultraviolet excess selection technique. Following Hall et al. (1996), we required plausible quasar candidates to have either $`(BV)<0.4`$ and $`(UB)>0.3`$, or $`(BV)<0.6`$ and $`(UB)<0.3`$. This selection is effective in separating QSOs from the stellar locus, and is efficient in detecting them out to $`z=2.2`$ (Hall et al. 1996; Fan 1999). We note that of the current identifications in the 3EG catalog, which are unbiased by optical selection, the largest redshift is only $`z=2.286`$, and all of their optical counterparts are brighter than $`V=22.1`$. Our color selection should also permit the discovery of any object that has a power-law continuum, which produces a UV excess, and especially a synchrotron spectrum which peaks above the optical band, e.g., those blazars commonly referred to as high-energy peaked. Thus, our technique is sensitive to most of the known EGRET blazars, and useful to search for a UV excess counterpart that might be expected on the basis of the absence of strong radio emission.
The major complication in this search comes in separating quasars from white dwarfs, blue field stars, and compact emission-line galaxies which often have similar blue colors and are known major contaminants of quasar color surveys. Further criteria can be applied using additional colors, but we decided to allow maximum freedom in the selection criteria in order to avoid excluding possibly interesting candidates. A total of 40 such candidates to a limiting magnitude of $`B=21`$ were selected for follow-up spectroscopy. In subsequent sections of this paper we discuss the eight QSOs that were discovered in our spectroscopic observations.
### 2.2 X-ray Observations
A total of three X-ray observations were made that cover the entire 99% error ellipse of 3EG J1835+5918, two by the ROSAT High Resolution Imager (HRI) and one by ASCA. The first ROSAT observation took place on 1995 February 2–4, with a total exposure time of 9,186 s. Five point-like X-ray sources were detected in this image, which reached a minimum detectable intrinsic flux of $`7.4\times 10^{14}`$ erg cm<sup>-2</sup> s<sup>-1</sup> in the 0.1–2.4 keV band, assuming a power-law spectrum with photon index 2.0 and Galactic $`N_\mathrm{H}=4.6\times 10^{20}`$ cm<sup>-2</sup>. A longer HRI observation of the same field was obtained between 1997 December 15 and 1998 January 20, with a total exposure time of 61,269 s. This deeper observation detected a number of fainter X-ray sources above a limiting unabsorbed flux of $`2\times 10^{14}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, including four of the five previous sources, as well as 10 new ones. Nine sources fall within the 99% confidence ellipse of 3EG J1835+5918. All of these sources are listed in Table 1, together with information about their optical identifications, which are radio-quiet QSOs or coronal emitting stars. The HRI astrometry was recalibrated using the optical counterparts of five well-localized X-ray sources, for which an average translation of $`2.^{\prime \prime }3`$ was required. After this shift, the five fiducial X-ray sources have a dispersion of only $`0.^{\prime \prime }8`$ from their optical positions. In Table 1 we list optical position, or recalibrated X-ray position in the case that no firm optical identification has been made. X-ray fluxes are calculated assuming a power law of photon index –2.0 and the full Galactic $`N_\mathrm{H}`$ for QSOs and unidentified sources, and a Raymond-Smith thermal plasma of $`T=3\times 10^6`$ K and $`N_\mathrm{H}=1\times 10^{20}`$ cm<sup>-2</sup> for stars.
An ASCA observation took place from 1998 April 20–22 for a total clean exposure time of $`68,900`$ s in each of the two Gas Imaging Spectrometers (GIS). Figure 1 shows the combined GIS image. The detection threshold for this ASCA observation was $`1.1\times 10^{13}`$ erg cm<sup>-2</sup> s<sup>-1</sup> (1–10 keV) assuming a photon index of –1.7. Several sources are detected far from the EGRET error ellipse, and only one faint source falls within it, a radio-quiet QSO at $`z=0.973`$ that was also detected by ROSAT. In Table 1 we give information about this and four additional ASCA sources outside the EGRET error ellipse that we were able to identify.
Diffuse X-ray emission at the western edge of the ASCA GIS image appears to be coming from an uncatalogued cluster of galaxies that is evident on our CCD images. We have not attempted to measure the X-ray flux of this source as it is too close to the edge of the detector and may extend outside it. The brightest galaxies in this vicinity are members of the cluster at $`z=0.102`$ and have $`R14`$ and $`R15`$, at J2000 coordinates $`18^\mathrm{h}32^\mathrm{m}38.^\mathrm{s}01,+59^{}23^{}43.^{\prime \prime }8`$, and $`18^\mathrm{h}32^\mathrm{m}49.^\mathrm{s}52,+59^{}21^{}49.^{\prime \prime }4`$, respectively. This X-ray source is well outside the 3EG J1835+5918 error ellipse, and we have no reason to suspect that they are related. In particular, there is no evidence of an AGN in this cluster.
The field of view of the ASCA Solid-state Imaging Spectrometer (SIS) detectors, even when operated in 4-CCD mode during this observation, is too small to cover the EGRET error ellipse. No X-ray sources were detected in the SIS images, so we do not discuss them further here.
### 2.3 Radio Observations
We reduced an archival VLA observation of this field which was taken at a frequency of 1.4 GHz on 1995 February 21 in the D configuration. We found 14 sources stronger than 2.5 mJy in the neighborhood of 3EG J1835+5918. They have a positional accuracy of approximately $`7^{\prime \prime }`$ for the fainter sources, and $`1^{\prime \prime }`$ for sources stronger than 15 mJy. For completeness, we examined the NVSS catalog at the same frequency to confirm six more faint sources that were marginally detected in the 1995 pointing. To incorporate information at other radio frequencies, we searched the Westerbork Northern Sky Survey (WENSS), which covered this field to a limiting flux of 18mJy at 326 MHz (Rengelink et al. 1997), and the NRAO 4.85 GHz catalog of Becker et al. (1991), which has a flux limit of 20 mJy at this location. A combined total of 20 radio sources were found inside and outside the error ellipse. Their properties are listed in Table 2, and their positions are shown in Figure 2. Most notably, there are no flat-spectrum sources in this field, and there are only three sources within the 99% confidence error ellipse of 3EG J1835+5918, all fainter than 4 mJy at 1.4 GHz.
## 3 Optical Spectroscopy and Results
We used a number of spectrographs to obtain moderate-resolution spectra of candidate X-ray and radio counterparts as well as UV excess objects selected from our optical imaging survey. These instruments include the Goldcam spectrograph on the KPNO 2.1m telescope, the Mark III spectrograph on the MDM 1.3m McGraw-Hill and 2.4m Hiltner telescopes, the Kast double spectrograph on the 3m Shane reflector at Lick Observatory, the Low Resolution Spectrograph (LRS) on the Hobby-Eberly telescope, and the Low Resolution Imaging Spectrograph (LRIS) on the Keck II telescope. Most spectra were analyzed independently by two authors and an agreement on classification was reached after comparing separate findings. The spectra were analyzed for emission and absorption lines and classified as either as star, galaxy, white dwarf, AGN, or uncertain. We have completed spectroscopy to a limiting magnitude of $`B=20.3`$, which includes 43 out of 53 optical candidates. In addition we have spectra of two objects fainter than $`B=20.3`$. Finding charts for the classified objects are given in Figures 3 and 4, and their spectra are shown in Figures 5 and 6.
Thus far we have found eight QSOs by the UV excess technique in the magnitude range $`18.5<B<21.3`$. Their redshifts range from 0.504 to 2.21. These are listed in Table 3 and their positions are shown in Figure 2. By design, they all fall within or very close to the 3EG J1835+5918 error ellipse. The efficiency of our color selection agrees fairly well with the number counts reported by Koo & Kron (1998) and Hall et al. (1996), which would predict that six QSOs with $`B<20.3`$ and $`z<2.3`$ would be found within a region of this size. Several additional candidates were found to have featureless blue spectra that we cannot securely classify. Since their colors are consistent with those of white dwarfs, we suspect that they are of the weak-lined (DC) variety.
Of the X-ray sources, six have been identified with radio-quiet QSOs, including five that were independently selected by UV excess colors. A seventh X-ray quasar is an ASCA and radio source at $`z=0.668`$ that lies well outside the EGRET error circle. Four more X-ray sources are identified with coronal emitting stars of types G, K, and dMe whose X-ray fluxes are normal for their optical magnitudes.
Two radio sources outside the EGRET error ellipse are identified with bright, early type galaxies at redshifts of 0.106 and 0.156, respectively, that lack any emission lines or evidence of non-stellar continuum in their optical spectra (Figure 6). Neither of these are promising $`\gamma `$-ray source candidates. The lower-redshift galaxy is close to the X-ray emitting galaxy cluster that is west of the EGRET error ellipse and it is apparently a member of the cluster.
We have had less success in identifying the faint radio sources within the EGRET error ellipse. Bright optical objects near their positions have proven to be ordinary stars, indicating that their true optical counterparts are likely to be fainter than our limiting magnitude for spectroscopy. Finding charts for both of the radio galaxies, as well as for several unidentified radio sources, are displayed in Figures 3 and 4.
## 4 Notes on Individual Interesting Objects
RX J1834.1+5913: This is the brightest quasar in the EGRET error ellipse ($`V=18.8,z=0.973`$) and it is detected by both ASCA and ROSAT. Its X-ray flux decreased between the two ROSAT observations, from $`1.9\times 10^{13}`$ erg cm<sup>-2</sup> s<sup>-1</sup> in 1995 to $`4.76\times 10^{14}`$ erg cm<sup>-2</sup> s<sup>-1</sup> in 1997–98. However, we are cautious about this variability since the source was near the edge of the detector in the later observation. We have several optical measurements of it in 1997, 1998, 1999 which also show modest variability. The largest change of $`0.39`$ magnitudes occurred between 1998 June and 1999 September, but there is no evidence for rapid variability on time scales of days. In addition, the equivalent width of its Mg II emission line did not vary in spectra taken at two different epochs. Thus, the spectral and variability properties of RX J1834.1+5913 offer no strong reason to argue that it is a candidate identification for 3EG J1835+5918. However, as the brightest QSO in the EGRET error ellipse, it does warrant continued scrutiny. In §5, we compare the properties of this source to those of the identified EGRET blazars in order to illustrate how unusual any AGN counterpart of 3EG J1835+5918 must be.
UVQ J1834.3+5926: At $`z=2.21`$, this is the highest redshift QSO that we have found near 3EG J1835+5918. Its optical spectrum is somewhat unusual in that it is the reddest of all the QSOs in this field, and its emission lines are broad but weak. We suspect that its Ly $`\alpha `$ line, which falls just blueward of our Keck spectrum, is responsible for boosting its $`U`$-band flux and helping it to meet the UV excess criterion.
RX J1834.4+5920: This relatively bright ROSAT source ($`5.3\times 10^{14}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, assuming a $`T=3\times 10^6`$ K thermal plasma spectrum) remains unidentified, although it lies near the edge of the HRI detector where the point-spread function is very poor. An M star of magnitude $`R=17.8`$ has been suggested as a possible identification even though it lies $`15^{\prime \prime }`$ from the X-ray position (Carramiñana et al. 2000).
VLA J1834.7+5918: This faint radio source of 3.7 mJy remains without spectroscopy, yet a blue optical object with $`V=21.4`$ falls just inside the western boundary of its error circle (see Figure 3). Although lacking X-ray emission, it is still a possible quasar or BL Lac object and worth further study, especially spectroscopy of the optical candidate. Since this is the brightest and most promising radio source of those within the EGRET error ellipse, we adopt its radio flux as an upper limit for 3EG J1835+5918 in subsequent discussion.
VLA J1835.1+5906: This is the brightest radio galaxy ($`R=15.1,z=0.156`$) at the edge of the EGRET error ellipse. Its optical spectrum was examined for any evidence of a BL Lac object in its nucleus, the principal indicator of which would be a shallower than normal break at 4000 Å. However, no such evidence is seen. This plus its steep radio spectrum, $`\alpha =0.53`$ between 1.4 and 4.85 GHz and absence of X-ray emission argue against VLA J1835.1+5906 being a BL Lac identification of 3EG J1835+5918.
VLA J1835.6+5939 (=AX J1835.7+5939): This is a quasar at $`z=0.668`$ and the brightest radio source near 3EG J1835+5918, with a 1.4 GHz flux of 359 mJy. However, it is outside of the 99% error ellipse by $`8^{}`$, and this plus its steep radio spectrum, $`\alpha =0.84`$ between 1.4 and 4.85 GHz, argue against considering it as a strong $`\gamma `$-ray candidate.
RX J1836.2+5925: This is perhaps the most intriguing object found in all of our searches. It was the brightest X-ray source within the error ellipse ($`1.6\times 10^{13}`$ erg cm<sup>-2</sup> s<sup>-1</sup>), at least during the second ROSAT observation, but it was undetected in the first ROSAT pointing or in the ASCA observation. Thus, it must have varied by at least a factor of 2 in the long term, although it emitted steadily over the one-month span which comprises the second ROSAT observation. RX J1836.2+5925 of interest here primarily because it does not have an optical counterpart in any color (Figure 7) to limits of $`U>22.3,B>23.4,V>23.3`$, and $`R>22.5`$. A red stellar object of $`R19.7`$ is located $`11.^{\prime \prime }4`$ west of the X-ray centroid, but it is not a viable candidate given the precision with which several other X-ray sources in this field line up with their established optical counterparts. As described above, the HRI astrometry in this figure was recalibrated using the optical counterparts of five well-localized X-ray sources, for which an average translation of $`2.^{\prime \prime }3`$ was required. After this shift, the five X-ray sources have a dispersion of only $`0.^{\prime \prime }8`$ from their optical positions. Thus, the illustrated error box which is $`8^{\prime \prime }`$ on a side must include the true position beyond a reasonable doubt.
None of the optical objects near the error box of RX J1836.2+5925 show any proper motion which could account for their positional discrepancy with the X-ray source. We have not obtained spectroscopy for any of these faint neighbors, but a deeper and more exhaustive optical study of this X-ray source would be important to evaluate its qualifications as a possible new type of $`\gamma `$-ray source counterpart. By the definition of Stocke et al. (1991), this X-ray source has an X-ray to optical flux ratio $`f_X/f_V>78`$. Such a high ratio is found only among low-mass X-ray binaries and isolated neutron stars. As we argue below, neither of these object classifications would make 3EG J1835+5918 compatible with the broad-band spectra of the well-identified EGRET sources.
If RX J1836.2+5925 is not the counterpart of 3EG J1835+5918, then it might be similar to the newly discovered class of luminous soft X-ray transients that have been found by ROSAT in the nuclei of non-active galaxies. These are as luminous as $`10^{44}`$ erg s<sup>-1</sup> and last for several months (Grupe, Thomas, & Leighly 1999; Komossa & Greiner 1999; Komossa & Bade 1999). A promising interpretation of these events is tidal disruption and accretion of stellar debris by a central black hole. If RX J1836.2+5925 is such an event, then it could reside in a host galaxy at $`z0.5`$ which deeper optical imaging could detect.
## 5 Multiwavelength Comparisons to Known Classes of EGRET Sources
### 5.1 Blazars
Our radio, optical, and X-ray data on active objects in the field of 3EG J1835+5918 can be compared with other identified EGRET sources to evaluate whether 3EG J1835+5918 can still fall within the multiwavelength parameters of any of the known classes of $`\gamma `$-ray emitters. Beginning with blazars, Figure 8 shows radio, optical, X-ray, and $`\gamma `$-ray fluxes of the sample of well-identified EGRET blazars defined by Mattox et al. (1997, and personal communication). ROSAT and Einstein fluxes are taken from Fossati et al. (1998), and $`V`$ magnitudes and total 4.85 GHz radio fluxes from Mattox (personal communication). The EGRET spectral points from 3EG J1835+5918 are taken from Reimer et al. (2000). Of the numerous candidate identifications which we could superpose, we chose two, namely, the brightest QSO within the error ellipse (RX J1834.1+5913, $`z=0.973`$), and the brightest radio source within the error ellipse (VLA J1834.7+5918). For the latter, we hypothesize that the suggestive $`V=21.4`$ optical identification is correct, and we graph an X-ray upper limit from the deeper ROSAT observation. For the QSO, we assign an upper limit of 0.5 mJy at 1.4 GHz, from the VLA image.
The smooth curves fitted to these two candidates correspond to the sum of two empirical third-order polynomials as applied by Comastri et al. (1995). This is not a model of blazar emission, but only a guide to the eye in making empirical estimates of the peak fluxes at low and high energy. In doing so we assume the presence of two emission mechanisms, a low-energy synchrotron component and a high-energy component peaking in the $`\gamma `$-ray band, possibly due to inverse Compton scattering.
While the optical and X-ray properties of our brightest candidates are not unprecedented, they lie at the faint end of the distributions. In particular, the X-ray upper limit for VLA J1834.7+5918 (or any of the other radio sources in the error ellipse) falls below the faintest blazars by at least an order of magnitude. More significant are their faint radio fluxes which, in the case of VLA J1834.7+5918 is two orders of magnitude fainter than the faintest radio counterpart of any well-identified EGRET blazar. RX 1834+5913 is nearly three orders of magnitude fainter in the radio band. Figure 9, in which the ratio of 4.85 GHz flux density to the peak $`\gamma `$-ray flux in the range $`E>100`$ MeV is graphed as a function of $`\gamma `$-ray flux for the Mattox blazars, confirms the highly discrepant positions of any of the QSOs or radio sources which are positionally coincident with 3EG J1835+5918 and candidates for identification with it. (We assume in this Figure a flat radio spectrum for 3EG J1835+5918, since none of its faint candidates were actually detected at 4.85 GHz.)
Another property of the majority of EGRET blazars is their rapid and large-amplitude flux variations. The absence of such obvious $`\gamma `$-ray variability from 3EG J1835+5918 already argues against a blazar nature for it (Reimer et al. 2000). We have also searched our $`V`$-band images obtained in 1997 and 1998 for objects with rapid or extreme optical variability, looking for variations of $`\mathrm{\Delta }V>0.3`$. Apart from the modest variability of the $`z=0.973`$ QSO RX J1834.1+5913 described above, no optical candidates for blazar activity were discovered in this manner.
### 5.2 Rotation-Powered Pulsars
Similar to the comparison with known blazars, we can examine how 3EG J1835+5918 compares to the EGRET pulsars. In Figure 10, we compare the 0.1–2.4 keV X-ray flux (Becker & Trümper 1997) and average flux $`E>100`$ MeV for EGRET pulsars (Fierro 1995; Kaspi et al. 2000; Ramanamurthy et al. 1995). Any possible pulsar counterpart of 3EG J1835+5918 should be assigned an X-ray flux upper limit equal to the flux of the brightest unidentified ROSAT source in the error ellipse. This role is therefore properly assigned to RX J1836.2+5925, although the fact that it is variable in X-rays already places some doubt upon its credentials as a pulsar candidate. Most of the soft X-ray flux observed from intermediate-age neutron stars is surface thermal emission, which should not vary from year to year. However, the additional nonthermal X-ray component which is present in Geminga and other $`\gamma `$-ray pulsars could in principle vary, and Halpern & Wang (1997a) suggested that it does in Geminga. Therefore, we use the quiescent flux upper limit of this source (from the 1995 ROSAT observation) for comparison in Figure 10. Such a comparison strains the analogy with Geminga. While the latter is a cooling neutron star with $`T5.6\times 10^5`$ K at $`d150`$ pc, 3EG J1835+5918 is about 50 times fainter in X-rays, thus either $`d>1`$ kpc, or if it is to be located at a similar distance as Geminga, its surface temperature should be less than $`3\times 10^5`$ K. The larger distance is problematic, since it implies a $`\gamma `$-ray luminosity of $`1.7\times 10^{35}(d/1\mathrm{kpc})^2`$ erg s<sup>-1</sup> if isotropic, which is at least 5 times larger than the spin-down power of Geminga, $`3.3\times 10^{34}`$ erg s<sup>-1</sup>. Alternatively, if it is closer than 1 kpc, then its surface must be cooler and it is likely to be older than $`3\times 10^5`$ yr, which would also strain its $`\gamma `$-ray efficiency. If 3EG J1835+5918 is a pulsar but RX J1836.2+5925 is not its counterpart, then its X-ray flux upper limit is reduced to $`5\times 10^{14}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, or 80 times fainter than Geminga.
Younger pulsars such as Vela, PSR B1951+32, and PSR B1706-44 are EGRET sources with luminosities in the range $`(12)\times 10^{35}`$ erg s<sup>-1</sup>, but 3EG J1835+5918 lacks the nonthermal X-ray emission and synchrotron nebulae that accompany those more luminous pulsars. Furthermore, it would be highly unexpected to find a pulsar of characteristic age $`\tau <1\times 10^5`$ yr at $`d400`$ pc from the Galactic plane. since this would require a kick velocity $`v>5000(\tau /10^5\mathrm{yr})^1`$ km s<sup>-1</sup>. An interesting possibility would be a recycled millisecond pulsar, which could be old yet energetic. But even such pulsars manage to channel at least $`5\times 10^4`$ of their spin-down power into either thermal (Halpern & Wang 1997b) or nonthermal X-rays (Becker & Trümper 1999; Mineo et al. 2000). In the case of 3EG J1835+5918, any pulsar counterpart would have $`L_X(0.12.4\mathrm{keV})/L_\gamma (>100\mathrm{MeV})<6\times 10^5`$, which places a uniquely low limit on the ratio of X-ray to spin-down power.
### 5.3 Other Possible $`\gamma `$-ray Sources
In addition to the well established classes of $`\gamma `$-ray blazars and pulsars, several associations have been suggested which are highly plausible even while not conclusively proven. Most notable is the radio star and Be/X-ray binary LSI $`+61^{}303`$ (Strickman et al. 1998), long associated with the $`\gamma `$-ray source 2CG 135+01. Similar objects might be the 47 ms pulsar B1259–63 with a Be star companion, detected up to 200 keV (Grove et al. 1995), and the Be/X-ray binary SAX J0635+0533 in the error circle of 2EG J0635+0521 (Kaaret et al. 1999). Since these systems all have neutron stars with Be star companions, their $`\gamma `$-ray emission is not necessarily confined to the pulsar magnetospheric mechanism, but may arise in the interaction of the relativistic pulsar wind with the wind of the companion, or with its radiation. However, all of these systems have bright optical companions, as well as strong X-ray emission at least at some of the time. It is estimated that there are only 200 Be/X-ray binaries within 5 kpc (Rappaport & van den Heuvel 1982); these are young systems which are confined to the Galactic disk. If a Be star binary, the location of 3EG J1835+5918 well away from the Galactic plane would probably make it the nearest such system, and virtually impossible to miss since its $`V`$ magnitude would be brighter than 9 if at $`d<1`$ kpc. Since no such Be star is present in this region, this scenario for 3EG J1835+5918 can safely be ruled out.
## 6 Implications for EGRET Source Identifications
The statistical issues concerning the identification of EGRET sources with flat-spectrum radio sources were rigorously addressed by Mattox et al. (1997), and it is hardly possible to improve upon that analysis at this time. To summarize, flat-spectrum radio sources are the only AGNs that have been detected by EGRET with any degree of confidence. Unfortunately, while the EGRET survey is flux limited, the radio identifications of EGRET sources in Figure 9 are not flux limited, but rather are plagued by source confusion due to the large size of the EGRET error circles and the large surface density of radio sources. Thus, the statistical reliability of EGRET source identifications is lower than that in any other branch of astronomy. As Mattox et al. calculate, the radio sources that are reasonably secure (i.e., $`>95\%`$ confidence) identifications of EGRET sources have 5 GHz flux densities $`>500`$ mJy. That is why the correlation between radio flux and $`\gamma `$-ray flux in Figure 9 is weak and less than linear. Below 50 mJy, it is not even possible to make a meaningful argument for identification because the mean separation of such radio sources on the sky is comparable to the size of the EGRET error circles. Accordingly, there are three radio sources within the error circle of 3EG J1835+5918, and they are all fainter than 4 mJy. None is an X-ray source. If any one of these radio sources were the true counterpart of the EGRET source, its ratio of radio to $`\gamma `$-ray flux would be two orders of magnitude smaller than that of any known blazar. We are not claiming that 3EG J1835+5918 is unique in this regard. Other unidentified EGRET sources may eventually prove to be similar.
Even though we cannot yet point to a likely identification of 3EG J1835+5918, it is apparent from our multiwavelength observations that the true counterpart must be physically different or extreme in its properties relative to the classes of EGRET sources that have been identified so far. This is true whether the counterpart is one of the candidates studied here, or an undetected fainter object. Furthermore, it is unlikely that a systematic error in $`\gamma `$-ray position has caused us to overlook a more conventional identification. In radio and X-ray we have explored a region approximately 4 times the size of the 99% confidence location, and even within this ample area there are no blazar or pulsar candidates. For example, even if the counterpart were the brightest radio quasar in Figure 1, which is $`8^{}`$ from the edge of the 99% confidence region, that object is a steep-spectrum radio source, as are all of the other bright radio sources outside the EGRET error ellipse. Thus, an error of this magnitude in the location of 3EG J1835+5918 will not change the basic conclusion that a new or extreme type of counterpart is responsible.
One possible implication of this result is that radio-steep or radio-quiet quasars could be counterparts of some of the unidentified EGRET sources, despite the analysis of Mattox et al. (1997) which argues that such a new population is not needed. Instead of interpreting the hard $`\gamma `$-ray spectrum and lack of variability as pulsar-like, it might be that these properties are also characteristic of the less violently variable AGNs. The obstacle to identifying a potential radio-weak or radio-quiet EGRET source population is not sensitivity, but source overlap. There are simply too many such AGNs in any EGRET error circle. While it is almost certainly the case that weaker radio blazars will be identified with high-energy $`\gamma `$-ray sources once their error circles are reduced by GLAST, it remains to be seen whether or not qualitatively different types of AGN will be also be represented.
An interesting scenario for a new type of $`\gamma `$-ray AGN has been suggested by Ghisellini (1999), who posits the existence of blazars whose synchrotron spectrum peaks in the MeV band, and an inverse-Compton component that peaks in the TeV. A variation of such a model could fit the multiwavelength spectrum of the $`z=0.973`$ QSO RX J1834.1+5913 or any of the fainter QSOs in the field provided that the proper index for the power-law electron energy distribution can be accommodated, and only if the observed optical emission is dominated by the usual thermal accretion-disk emission so that it can represent an upper limit to the underlying synchrotron power law. In such a model the hard X-ray emission is due entirely to the synchrotron component. The absence of a radio counterpart is naturally explained by the form of the power law, which in this case requires a flat spectral index $`\alpha 0.45`$ where $`F_\nu \nu ^\alpha `$, thus the power-law index of the electron energy distribution is $`p1.9`$ . Such a prediction can easily be tested by more sensitive hard X-ray spectra of the QSO RX J1834.1+5913.
Radio-quiet blazars have been hypothesized theoretically (Ghisellini 1999; Mannheim 1993; Schlickeiser 1984) but so far none have been identified (Stocke et al. 1990; Jannuzi et al. 1993), and it is not even clear what such a phenomenon would mean. Could the multiwavelength properties of 3EG J1835+5918 be evidence of the hadronic model, the so called proton blazar? Such a theory proposes to explain $`\gamma `$-ray emission in blazars, relying on protons accelerated by shocks moving through the jet. The accelerated protons then interact with soft-photons which lead to the creation of pions that further decay and cascade into electron-positron pairs, $`\gamma `$-rays and neutrinos. Such a model (Mannheim 1993) could fit the observations of 3EG J1835+5918 if the energy density ratio of protons to electrons is greater than 10.
If 3EG J1835+5918 is a pulsar, it implies that highly efficient (or highly beamed) $`\gamma `$-ray pulsars can avoid producing soft X-rays at a level below $`10^4`$ of their apparent $`\gamma `$-ray luminosity. At least two mechanisms of X-ray emission have been observed to accompany all $`\gamma `$-ray pulsars at such levels or higher (Wang et al. 1998). In the outer-gap model, synchrotron emission from secondary pairs that are produced by conversion of $`\gamma `$-rays in the inner magnetosphere where $`B>2\times 10^{10}`$ G can explain the nonthermal X-ray component from pulsars like Geminga and PSR B1055–52. The second mechanism is thermal emission arising from the heated polar caps that are impacted by the inward-going accelerated particles from the outer-gap accelerator. There is good evidence that polar-cap heating occurs even in recycled pulsars which are not detectable EGRET sources (Zavlin & Pavlov 1998; Halpern & Wang 1997b). Therefore, it is difficult to reconcile such a theory, as well as the observational fact that pulsars are X-ray sources of $`L_X>10^4I\mathrm{\Omega }\dot{\mathrm{\Omega }}`$, with a pulsar origin for 3EG J1835+5918. If many of the unidentified EGRET sources are similar radio-quiet pulsars in the Galactic plane, X-ray absorption makes them exceedingly difficult to identify, and perhaps they will be revealed only when $`\gamma `$-ray observations are sensitive enough to detect their pulsations independently.
## 7 Conclusions and Further Work
We identified all but one of the X-ray sources in the field of 3EG J1835+5918 to a flux limit of approximately $`5\times 10^{14}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. These are radio-quiet QSOs \[$`F(1.4\mathrm{GHz})<0.5`$ mJy\], coronal emitting stars, and a cluster of galaxies. There are no flat-spectrum radio sources in the vicinity to a flux limit of $`20`$ mJy, and no radio sources in the EGRET error ellipse brighter than 4 mJy at 1.4 GHz. In addition, we find no evidence of a BL Lac object hosted in any low-redshift galaxy. We also found several QSOs, as one would expect, using purely optical color selection. Multiple-epoch optical imaging of the entire EGRET error ellipse has not revealed any notable variability. The discovery of only radio-quiet quasars in the error circle of 3EG J1835+5918 is a sobering development in the search for its identification. Although the $`\gamma `$-ray properties of 3EG J1835+5918 are more similar to those of Geminga and other EGRET pulsars, no other indirect evidence for a pulsar, apart from one unidentified X-ray source (RX J1836.2+5925) whose optical counterpart is probably fainter than $`B=23.4,V=23.3`$, and $`R=22.5`$, has been found. Yet, the fact that this X-ray source is variable by at least a factor of 2 would make it unique among rotation-powered pulsars. Taken together, these findings point to the possibility of a truly remarkable object, one that cannot be matched by any known class of $`\gamma `$-ray source.
Even in the absence of a definite identification, it is clear that 3EG J1835+5918 is lacking in one or more of the physically essential attributes of any known class of $`\gamma `$-ray emitter. Its radio flux is at least two orders of magnitude fainter than any of the securely identified EGRET blazars, and its soft X-ray flux is at least 50 times fainter than that of Geminga and similar EGRET pulsars. If it is an AGN it lacks the beamed radio emission of blazars. If it is an isolated neutron star, it lacks the steady thermal X-rays from a cooling surface and the magnetospheric non-thermal X-ray emission that is characteristic of all EGRET pulsars. If a pulsar, 3EG J1835+5918 must be either older or more distant than Geminga, and probably an even more efficient or highly beamed $`\gamma `$-ray engine.
We have plans to complete the optical spectroscopy of fainter candidates in this field to $`B21.5`$ and we will also study fundamental properties such as polarization and optical variability of the newly discovered AGNs. Perhaps the most important technique which we have not yet applied is polarimetry. Polarimetry provides a definitive test for synchrotron emission in an ordered magnetic field, and polarization is one of the essential properties of blazars. Perhaps the blazar nature of a radio-quiet beam in an AGN can only be demonstrated in this way. A deeper radio pulsar search would also be warranted. Finally, we will pursue the optical identification of the ROSAT source RX J1836.2+5925 to the faintest magnitudes that are necessary in order to find our whether or not it is a neutron star. In combination, these observations may result in the identification of an important EGRET source, and possibly the prototype of a new class of $`\gamma `$-ray emitter.
We thank Eric Gotthelf for his assistance with the reduction of ASCA data, Karen Leighly and John Tomsick for assistance with the optical imaging, and John Mattox, Greg Madejski, Brenda Dingus, and Reshmi Mukherjee for helpful discussions. This work was supported by NASA grants NAG 5-3229 and NAG 5-7814.
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# Multiplets of representations and Kostant’s Dirac operator for equal rank loop groups
## 0. Introduction
Although this paper is chiefly concerned with representations of Lie groups and loop groups, the motivation for these results originally comes from M-Theory. In physics, the Lie group $`\mathrm{Spin}(9)`$ arises as the little group for massive particles in 10 dimensional superstring theories and as the little group for massless particles in 11 dimensional supergravity. Recently, Pengpan and Ramond noticed that the irreducible representations of $`\mathrm{Spin}(9)`$ come in triples, with the Casimir operator taking the same value on all three representations, and where the dimensions of two such representations sum to the dimension of the third. Ramond brought this curious fact to the attention of Sternberg, who in collaboration with Gross and Kostant then showed that these triples of representations of $`B_4=\mathrm{Spin}(9)`$ actually correspond to representations of the exceptional Lie group $`F_4`$, which contains $`B_4`$ as an equal rank subgroup.
In fact, this is not an isolated phenomenon. In , Gross, Kostant, Ramond, and Sternberg consider the general case where $`𝔥`$ is a reductive Lie algebra which is a maximal rank subalgebra of some semi-simple Lie algebra $`𝔤`$. Letting $`G`$ and $`H`$ denote the compact, simply-connected Lie groups with Lie algebras $`𝔤`$ and $`𝔥`$ respectively, associated to any irreducible representation of $`G`$ is a set of $`\chi (G/H)`$ irreducible representations of $`H`$, where $`\chi (G/H)`$ is the Euler number of the homogeneous space $`G/H`$. We shall refer to such a set of $`H`$-representations as a *multiplet*. As in the case of $`B_4F_4`$, all of the representations in a multiplet share the same value of the Casimir operator, and the alternating sum of the dimensions of these representations vanishes. The relation between a representation of $`G`$ and the $`H`$-representations in the corresponding multiplet is given by the following homogeneous generalization of the Weyl character formula, viewed as an identity in the representation ring $`R(H)`$:
(1)
$$V_\lambda 𝕊^+V_\lambda 𝕊^{}=\underset{cC}{}(1)^cU_{c(\lambda +\rho _G)\rho _H},$$
where $`V_\lambda `$ and $`U_\mu `$ denote the irreducible representations of $`G`$ and $`H`$ with highest weight $`\lambda `$ and $`\mu `$ respectively, $`𝕊=𝕊^+𝕊^{}`$ is the spin representation associated to the complement of $`𝔥`$ in $`𝔤`$, the subset $`CW_G`$ of the Weyl group of $`G`$ has one representative from each coset of $`W_H`$, and $`(1)^c`$ is the sign of the element $`c`$.
In representation theory, the Casimir operator of a Lie algebra is analogous to the Laplacian. Using the spin representation, we can also consider operators analogous to the Dirac operator. Furthermore, we can choose a particular Dirac operator such that its square is the Casimir operator shifted by a constant, giving a representation theory version of the Weitzenböck formula. Such a Dirac operator was introduced in a more formal setting by Alekseev and Meinrenken in , and the geometric version of this Dirac operator is examined in . Since the Casimir operator takes the same value on all of the representations in a multiplet, it follows that this Dirac operator likewise takes a constant value, up to sign, on each multiplet.
In the homogeneous case, for any linear operator
$$\overline{)}:V_\lambda 𝕊^+V_\lambda 𝕊^{},$$
since both the domain and range are finite dimensional, the index of $`\overline{)}`$ must be given by (1). This prompted Kostant to seach for a Dirac operator whose kernel and cokernel are precisely those representations on the right hand side of (1). In , Kostant constructs a Dirac operator $`\overline{)}_{𝔤/𝔥}`$ on $`V_\lambda 𝕊`$ with a cubic term associated to the fundamental 3-form on $`𝔤`$. The kernel of Kostant’s Dirac operator is
(2)
$$\mathrm{Ker}\overline{)}_{𝔤/𝔥}=\underset{cC}{}U_{c(\lambda +\rho _G)\rho _H},$$
and the signs $`(1)^c`$ on the right side of (1) can be recovered by decomposing the operator $`\overline{)}_{𝔤/𝔥}`$ according to the positive and negative half-spin representations. Taking the kernel of Kostant’s Dirac operator therefore gives an explicit construction of the multiplet of $`H`$-representations corresponding to a given representation of $`G`$.
This paper takes the results discussed above and reformulates them in the Kac-Moody setting, replacing the equal rank Lie groups $`HG`$ with their corresponding loop groups $`LHLG`$. After briefly reviewing the representation theory of loop groups in §1, we introduce the positive energy spin representation $`𝒮_{L𝔤}`$ associated to a loop group in §2, using it to reformulate the Weyl-Kac character formula. In §3, we prove the following homogeneous version of the Weyl-Kac character formula:
$$_𝝀𝒮_{L𝔤/L𝔥}^+_𝝀𝒮_{L𝔤/L𝔥}^{}=\underset{c𝒞}{}(1)^c𝒰_{c(𝝀𝝆_𝔤)+𝝆_𝔥},$$
where $`_𝝀`$ and $`𝒰_𝝁`$ denote the positive energy representations of the central extensions $`\stackrel{~}{L}G`$ and $`\stackrel{~}{L}H`$ with lowest weights $`𝝀`$ and $`𝝁`$ respectively, the subset $`𝒞𝒲_G`$ now lives in the affine Weyl group of $`G`$, and $`𝝆_𝔤`$ and $`𝝆_𝔥`$ are the lowest weights of the spin representations $`𝒮_{L𝔤}`$ and $`𝒮_{L𝔥}`$.
In §§46, we return to the case of compact Lie groups, reviewing various results of and . There we construct Kostant’s Dirac operator, compute its square, and prove that its kernel has the form given by (2). Our approach here differs slightly from Kostant’s, which views the Lie algebra $`𝔤`$ as an orthogonal extension of $`𝔥`$. Instead, we first consider the Dirac operator on $`𝔤`$ and a twisted Dirac operator on $`𝔥`$ and then construct Kostant’s Dirac operator as their difference, an idea borrowed from . In addition, we avoid working with a basis for $`𝔤`$ wherever possible, which greatly simplifies the computations and hopefully elucidates their meanings. These sections can stand alone as an alternative exposition on Kostant’s Dirac operator, and they provide a outline of the more advanced material in the subsequent sections.
The remaining sections reprise these results for the loop group case. In §7 we examine the Clifford algebra associated to a loop group, which builds on the treatment of infinite dimensional Clifford algebras given in . We then introduce the Dirac and Casimir operators associated to a loop group in §8, and we construct the loop group analogue $`\overline{)}_{L𝔤/L𝔥}`$ of Kostant’s Dirac operator in §9. These Dirac and Casimir operators appear in the physics literature in and as the odd and even zero-mode generators for the $`N=1`$ superconformal algebras associated to current (Lie group) and coset space (homogeneous space) models. In contrast, our treatment builds these operators on a mathematical foundation, viewing them as canonical objects rather than working in terms of a basis. Finally, we compute the square of the Dirac operator $`\overline{)}_{L𝔤/L𝔥}`$, and in §10 we prove that its kernel is
$$\mathrm{Ker}\overline{)}_{L𝔤/L𝔥}=\underset{c𝒞}{}𝒰_{c(𝝀𝝆_𝔤)+𝝆_𝔥},$$
just as for compact Lie groups. So once again, taking the kernel of this Dirac operator provides an explicit construction for the multiplet of representations of $`\stackrel{~}{L}H`$ corresponding to any given representation of $`\stackrel{~}{L}G`$.
*Note.* Anthony Wassermann, who has independently obtained results similar to those in this paper, pointed out to me that with only minor modifications, the arguments presented here provide a quick proof of the Weyl-Kac character formula.
## 1. Loop groups and their representations
### 1.1. Loop Groups
Let $`G`$ be a compact connected Lie group, and let $`LG`$ denote the group of free loops on $`G`$, i.e., the space of smooth maps from $`S^1`$ to $`G`$, where the product of two loops is taken pointwise. The Lie algebra of the loop group $`LG`$ is simply the vector space $`L𝔤`$ of loops on the Lie algebra $`𝔤`$ of $`G`$, with brackets again taken pointwise. The group $`\mathrm{Diff}(S^1)`$ of diffeomorphisms of the circle acts on loop spaces by reparameterizing the loops, and in particular the subgroup $`S^1`$ of rigid rotations of the circle acts on $`LG`$ and $`L𝔤`$. This circle action induces a $``$-grading on the complexified Lie algebra $`L𝔤_{}`$, which is the closure of the direct sum of the Fourier components $`_k𝔤_{}z^k`$, where $`𝔤_{}z^k`$ denotes loops of the form $`zXz^k`$ for $`X𝔤_{}`$. We are interested in those representations of $`LG`$ which likewise admit a $``$-grading intertwining with the $`S^1`$-action on $`LG`$, or in other words representations of the semi-direct product $`S^1LG`$. Such a representation $``$ then decomposes into eigenspaces $`_k(k)`$ according to the $`S^1`$-weight $`k`$, called the *energy*. (This terminology comes from an analogy with quantum mechanics, where the energies are eigenvalues of the Hamiltonian operator, which generates time translation.)
### 1.2. The central extension
The representations that we will consider are actually projective representations of $`LG`$. To realize them as true representations, we must introduce a central extension $`\stackrel{~}{L}G`$ of $`LG`$ by $`S^1`$. This is analogous to taking the universal cover of a compact Lie group, except that here we lift to a circle bundle rather than a finite cover.
The corresponding central extension of the Lie algebra, which is called a *Kac-Moody algebra*, is $`\stackrel{~}{L}𝔤=L𝔤I`$, where $`I`$ is the infinitesimal generator of the central $`S^1`$ subgroup. The Lie bracket on the central extension $`\stackrel{~}{L}𝔤`$ is determined by a choice of $`\mathrm{ad}`$-invariant inner product on $`L𝔤`$. Any $`\mathrm{ad}`$-invariant inner product on the Lie algebra $`𝔤`$ induces an inner product on the $`L𝔤`$ by averaging the pointwise inner products. For loops $`\xi ,\eta L𝔤`$, this gives
(3)
$$\xi ,\eta =\frac{1}{2\pi }_0^{2\pi }\xi (\theta ),\eta (\theta )d\theta ,$$
which is $`\mathrm{ad}`$-invariant on $`L𝔤`$. To extend this inner product to $`\stackrel{~}{L}𝔤`$, we must actually go one step further and extend it to the semi-direct sum $`\stackrel{~}{}\stackrel{~}{L}𝔤`$, where $``$ is generated by the infinitesimal rotation $`_\theta `$, and we define the inner product by
$$a_\theta +\xi +xI,b_\theta +\eta +yI=\xi ,\eta aybx$$
for $`a,b,x,y`$ and $`\xi ,\eta L𝔤`$. This inner product is $`\mathrm{ad}`$-invariant on the extended Lie algebra $`\stackrel{~}{}\stackrel{~}{L}𝔤`$ provided that the Lie bracket on the central extension $`\stackrel{~}{L}𝔤`$ is
(4)
$$[\xi ,\eta ]_{\stackrel{~}{L}𝔤}=[\xi ,\eta ]_{L𝔤}+\xi ,_\theta \eta I.$$
Although this central extension depends on the original choice of inner product on $`𝔤`$, there is a unique $`\mathrm{ad}`$-invariant inner product on $`𝔤`$ (up to scaling) if $`𝔤`$ is *simple*. In this case, the *universal central extension* corresponds to the smallest possible scaling for which the Lie algebra $`\stackrel{~}{L}𝔤`$ exponentiates to give a central extension $`\stackrel{~}{L}G`$ of the loop group $`LG`$. This smallest inner product on $`𝔤`$ is the *basic inner product*, which is scaled so that the highest root $`\alpha _{\mathrm{max}}`$ of $`𝔤`$ satisfies $`\alpha _{\mathrm{max}}^2=2`$.
If $`G`$ is not simple but only semi-simple, then a given projective representation of $`LG`$ can still be lifted to a true representation of some $`S^1`$ extension of $`LG`$. However, the universal central extension of $`LG`$ is no longer a circle bundle, but rather an extension by the torus $`T^d`$, where $`d`$ counts the number of simple components of $`G`$. At the Lie algebra level, an $`\mathrm{ad}`$-invariant inner product on $`𝔤`$ can be scaled separately on each of the simple components, and the central term in the Lie bracket (4) now becomes $`d`$ separate terms corresponding to the basic inner products for each of these components.
###### Remark.
Let $`G`$ be simply connected. Topologically, the invariant inner products on $`𝔤`$ correspond to elements of the Lie algebra cohomology $`H^3(𝔤)H^3(G;)`$ by associating to any inner product its fundamental 3-form $`\omega \mathrm{\Lambda }^3(𝔤^{})`$ given by $`\omega (X,Y,Z)=X,[Y,Z]`$ for $`X,Y,Z𝔤`$. The possible central extensions of the Lie algebra $`L𝔤`$ by a circle thus correspond to elements of the real cohomology $`H^3(G;)`$, and the universal central extension of $`L𝔤`$ is then an extension by the dual space $`K=H_3(G;)`$. On the other hand, the central extensions of the loop group $`LG`$ correspond to circle bundles, which are classified by their Chern classes $`c_1H^2(LG;)H^3(G;)`$ in the integral lattice of $`H^3(G;)`$. Writing $`L=H_3(G;)`$ for the dual lattice in $`K`$, the universal central extension $`\stackrel{~}{L}G`$ is an extension of $`LG`$ by the torus $`K/L`$. Using the cohomology spectral sequence for this extension and noting that $`H^1(LG;)=H^2(G;)=0`$, we obtain the exact sequence
$$0H^1(\stackrel{~}{L}G;)H^1(K/L;)\stackrel{d_2}{}H^2(LG;)H^2(\stackrel{~}{L}G;)0.$$
Now, by our construction of the torus $`K/L`$, we have a canonical isomorphism $`H^1(K/L;)H^3(G;)`$, and we also have a canonical isomorphism $`H^2(LG;)H^3(G;)`$. The map $`d_2`$ is therefore a homomorphism $`d_2:H^3(G;)H^3(G;)`$, and the universality condition becomes the assertion that $`d_2`$ be the identity map. In particular, if $`\stackrel{~}{L}G`$ is the universal central extension, then $`d_2`$ must be an isomorphism, and it follows that $`H^1(\stackrel{~}{L}G;)=H^2(\stackrel{~}{L}G;)=0`$, which in terms of homotopy implies that $`\stackrel{~}{L}G`$ is 2-connected. So, whereas taking the universal cover of a compact Lie group $`G`$ kills the obstruction $`\pi _1(G)`$, the loop group $`LG`$ is already simply connected, but taking its universal central extension kills the obstruction $`\pi _2(LG)`$.
The semi-direct sum $`\stackrel{~}{}\stackrel{~}{L}𝔤`$ which we introduced above is the Lie algebra of the semi-direct product $`S^1\stackrel{~}{L}G`$, and from here on we refer to representations of $`S^1\stackrel{~}{L}G`$ as representations of $`LG`$. Given such a representation, we call the weight of the central $`S^1`$ in $`\stackrel{~}{L}G`$ the *level* or *central charge*, and since this circle by definition commutes with the rest of the loop group, it follows that the level is constant on each irreducible representation. Unless stated otherwise, from here on we assume that $`G`$ is simply connected and simple, we use the basic inner product on $`𝔤`$, and we let $`\stackrel{~}{L}G`$ denote the universal central extension. However, the following discussion can be generalized to the semi-simple case by treating the $`d`$ simple components separately and viewing the level as a $`d`$-vector.
### 1.3. Affine roots and the affine Weyl group
Let $`T`$ be a maximal torus of $`G`$. When considering the representation theory of loop groups, rather than taking the abelian subgroup $`LT`$ as the maximal torus of $`LG`$, we instead use the maximal torus $`S^1\times T\times S^1`$ of $`S^1\stackrel{~}{L}G`$. Here the first $`S^1`$ factor corresponds to rotation of loops, while the second comes from the central extension. The Cartan subalgebra is then $`𝔱`$, and the weights of $`LG`$ are of the form $`𝝀=(m,\lambda ,h)`$, where $`m`$ is the energy, $`\lambda `$ is a weight of $`G`$, and $`h`$ is the level. In this notation, the roots of $`LG`$, also called the *affine roots* of $`G`$, consist of the weights $`(m,\alpha ,0)`$ with $`m`$ and $`\alpha `$ a root of $`G`$, as well as the weights $`(m,0,0)`$ for nonzero $`m`$, counted with multiplicity $`\mathrm{rank}G=dim𝔱`$. Given a system of positive roots for $`G`$, we take the positive roots of $`LG`$ to be the roots $`(0,\alpha ,0)`$ for $`\alpha >0`$, as well as all roots $`(m,\alpha ,0)`$ with $`m>0`$, including roots of the form $`(m,0,0)`$. If $`\{\alpha _i\}`$ is a set of simple roots for $`G`$, then the corresponding simple affine roots for $`LG`$ are $`(0,\alpha _i,0)`$, as well as the root $`(1,\alpha _{\text{max}},0)`$, where $`\alpha _{\text{max}}`$ is the highest root of $`G`$.
The affine Weyl group $`𝒲_G`$ of $`G`$ is the group generated by the reflections through the hyperplanes corresponding to the affine roots of $`G`$. In terms of loop groups, given any root $`𝜶=(k,\alpha ,0)`$ of $`LG`$, there is a corresponding $`𝔰𝔲(2)`$ subalgebra of $`\stackrel{~}{L}𝔤`$ generated by the loops $`E_\alpha z^k`$ and $`E_\alpha z^k`$ and the coroot
$$H_{k,\alpha }=[E_\alpha z^k,E_\alpha z^k]_{\stackrel{~}{L}𝔤}=H_\alpha +\frac{1}{2}ikH_\alpha ^2I,$$
where $`\{E_\alpha ,E_\alpha ,H_\alpha \}`$ span the $`𝔰𝔲(2)`$ subalgebra of $`𝔤`$ associated to the root $`\alpha `$. Note that these elements are normalized so that $`E_\alpha ,E_\alpha =\frac{1}{2}H_\alpha ^2=2\alpha ,\alpha ^1`$. The reflection of a weight $`𝝀=(m,\lambda ,h)`$ through the hyperplane orthogonal to $`𝜶`$ is then
(5)
$$\begin{array}{cc}\hfill s_{k,\alpha }(𝝀)& =𝝀𝝀(H_{k,\alpha })𝜶\hfill \\ & =(m\lambda (H_\alpha )k+\frac{1}{2}hH_\alpha ^2k^2,\lambda \lambda (H_\alpha )\alpha +\frac{1}{2}hH_\alpha ^2k\alpha ,h).\hfill \end{array}$$
Furthermore, these $`s_{k,\alpha }`$ are generated by the reflections $`s_{0,\alpha }`$, which act solely on the $`𝔱^{}`$ component and generate the usual Weyl group $`W_G`$, as well as the transformations
$$t_\alpha (𝝀)=s_{1,\alpha }s_{0,\alpha }(𝝀)=(m+\lambda (H_\alpha )+\frac{1}{2}hH_\alpha ^2,\lambda +hH_\alpha ,h),$$
where we use the inner product to identify the coroot $`H_\alpha 𝔱`$ with the weight $`\frac{1}{2}H_\alpha ^2\alpha `$ in $`𝔱^{}`$. Restricting to $`𝔱^{}`$, the $`t_\alpha `$ are simply translations by the coroots, which generate the coweight lattice $`L𝔱`$. We therefore have $`𝒲_GW_GL`$.
Note that under the action of the affine Weyl group, the level $`h`$ is fixed, while the energy $`m`$ is shifted so as to preserve the inner product
(6)
$$(m_1,\lambda _1,h_1)(m_2,\lambda _2,h_2)=\lambda _1,\lambda _2m_1h_2m_2h_1$$
on $`𝔱^{}`$. Thus, at any given level $`h`$, the affine Weyl action is completely determined by its restriction to $`𝔱^{}`$. In particular, the element $`s_{k,\alpha }`$ corresponds to the reflection through the hyperplane given by the equation $`\lambda ,\alpha =hk`$. These hyperplanes divide $`𝔱^{}`$ into connected components called *alcoves*, and the affine Weyl group acts simply transitively on these alcoves. Given a positive root system for $`LG`$, the corresponding *fundamental alcove* is the unique alcove satisfying $`𝝀𝜶0`$ for all $`𝜶>0`$. This alcove is bounded by the hyperplanes corresponding to the negatives of the simple affine roots, or in other words, a weight $`𝝀=(m,\lambda ,h)`$ lies in the fundamental alcove if and only if $`\lambda `$ is in the positive Weyl chamber for $`G`$ and $`\lambda ,\alpha _{\text{max}}h`$.
### 1.4. Positive energy representations
A representation $``$ of $`LG`$ is a *positive energy representation* if $`(k)=0`$ for all $`k<m`$ for some fixed integer $`m`$, or in other words, there is a minimum energy when $``$ is decomposed into its constant energy eigenspaces. In the literature, positive energy representations are often normalized so that this minimum energy is 0. However, we will consider positive energy representations with the full spectrum of minimum energies. When restricted to the positive energy representations, the representation theory of loop groups behaves quite analogously to the representation theory of compact Lie groups. In particular, the positive energy representations satisfy the following fundamental properties (for a complete discussion, see ):
1. A positive energy representation is *completely reducible* into a direct sum of (possibly infinitely many) irreducible positive energy representations.
2. An irreducible positive energy representation $``$ is of *finite type*: each of the constant energy subspaces $`(k)`$ is a finite dimensional representation of $`G`$.
3. Every irreducible positive energy representation $``$ has a unique *lowest weight* $`𝝀=(m,\lambda ,h)`$, in the sense that $`𝝀𝜶`$ is not a weight of $``$ for any positive root $`𝜶`$ of $`LG`$. The lowest weight space is one dimensional and generates $``$.
4. A weight $`𝝀=(m,\lambda ,h)`$ is *anti-dominant* for $`LG`$ if it lies in the fundamental Weyl alcove described at the end of §1.3 above. The lowest weight of a positive energy representation is anti-dominant, and every anti-dominant weight is realized as the lowest weight of some positive energy representation.
As a consequence of (iii), an irreducible positive energy representation $``$ is completely characterized by its minimum energy $`m`$, its minimum energy subspace $`(m)V_\lambda `$, and its level $`h`$. Property (iv) implies that for a positive energy representation, the level $`h`$ is always non-negative and is zero only for the trivial representation. Also, for a fixed minimum energy $`m`$, there are only finitely many positive energy representations at each level $`h`$, but as the level tends to infinity, the representation theory of $`LG`$ resembles that of $`G`$.
If $`_𝝀`$ is the irreducible positive energy representation with lowest weight $`𝝀=(0,\lambda ,h)`$, then $`_𝝀`$ also contains all the weights in the orbit of $`𝝀`$ under the affine Weyl group $`𝒲_G`$. Recalling that the affine Weyl group action preserves the inner product (6), it turns out that the orbit of $`𝝀`$ consists of all weights $`𝝁=(m,\mu ,h)`$ at level $`h`$ satisfying $`𝝀𝝀=𝝁𝝁`$, or equivalently $`\mu ^22mh=\lambda ^2`$. This equation sweeps out a paraboloid, and the weights of $`_𝝀`$ all lie in its interior. (As the level $`h`$ tends to infinity, this paraboloid flattens into a cone.) For an example, see Figure 1 at the end of §3, which gives the weights of the irreducible representation of $`L\mathrm{SU}(2)`$ with lowest weight $`(0,1,2)`$.
## 2. The spin representation
If $`V`$ is a finite dimensional vector space with an inner product, and $`V=WW^{}`$ is a polarization of $`V`$ into a maximal isotropic subspace $`W`$ and its dual, then the spin representation of the Clifford algebra $`\mathrm{Cl}(V)`$ can be written in the form
(7)
$$𝕊_V=\mathrm{\Lambda }^{}(W)(detW)^{\frac{1}{2}},$$
where $`detW`$ denotes the top exterior power of $`W`$. The resulting spin representation $`𝕊_V`$ is independent of the choice of polarization, which is accounted for by the factor of $`(detW)^{1/2}`$. On the other hand, if $`V`$ is infinite dimensional, then this determinant factor does not make sense, and so we can no longer use (7) to define the spin representation. Without this determinant factor to correct for the choice of polarization, different polarizations give rise to distinct spin representations. For a general discussion of infinite dimensional Clifford algebras and their spin representations, see .
For our purposes, consider the Lie algebra $`L𝔤`$ with the inner product (3) induced by the basic inner product on $`𝔤`$. If we complexify $`L𝔤`$, then the orthogonal complement of the Cartan subalgebra $`𝔱_{}`$ in $`L𝔤_{}`$ decomposes into the sum of the positive and negative root spaces, each of which is isotropic with respect to the inner product on $`L𝔤_{}`$. We can therefore use this polarization to define a positive energy spin representation associated to the complement of $`𝔱`$ in $`L𝔤`$:
(8)
$$𝒮_{L𝔤/𝔱}:=𝕊_{𝔤/𝔱}\mathrm{\Lambda }^{}\left(\underset{k>0}{}𝔤_{}z^k\right)=𝕊_{𝔤/𝔱}\underset{k>0}{}\mathrm{\Lambda }^{}\left(𝔤_{}z^k\right),$$
where we have explicitly factored out the contribution $`𝕊_{𝔤/𝔱}`$ coming from the constant loops (or zero modes). Here, we have used the expression (7) for the spin representation, except that we have dropped the portion of the $`(detW)^{1/2}`$ factor coming from the positive energy modes. If we were to include that factor, it would contribute an overall anomalous energy shift of
(9)
$$\left(\underset{k>0}{}z^{kdim𝔤}\right)^{\frac{1}{2}}=z^{\frac{1}{2}_{k>0}kdim𝔤}=z^{\frac{1}{24}dim𝔤},$$
where in the last equality we use the Riemann zeta function trick to write the infinite sum as $`_{k>0}k=\zeta (1)=\frac{1}{12}`$. Fortunately, by normalizing the spin representation to have minimum energy 0, we can safely ignore this factor.
For the moment, we are interested only in the character of the spin representation. The restriction of the character of $`𝒮_{L𝔤/𝔱}`$ to $`S^1\times T`$ is completely determined by the description (8) of the spin representation. However, in correcting for the infinite determinant factor, the spin representation acquires a nonzero central charge.
###### Theorem 1.
If $`G`$ is simple, then the central charge of the spin representation $`𝒮_{L𝔤/𝔱}`$ is the value of the quadratic Casimir operator of $`𝔤`$ in the adjoint representation:
$$c_G=\mathrm{\Delta }_{\mathrm{ad}}^𝔤=\frac{1}{2}\underset{i}{}(\mathrm{ad}X_i)^2=\rho _G,\alpha _{\mathrm{max}}+1,$$
where $`\rho _G`$ is half the sum of the positive roots, $`\alpha _{\mathrm{max}}`$ is the highest root of $`G`$, and $`\{X_i\}`$ is an orthonormal basis for $`𝔤`$.
###### Proof.
To compute the central charge of the spin representation $`𝒮_{L𝔤/𝔱}`$, we extend it to obtain the spin representation associated to the entire Lie algebra $`L𝔤`$. Since the construction of spin representations is multiplicative, we have
$$𝒮_{L𝔤}𝕊_𝔱𝒮_{L𝔤/𝔱}.$$
These two spin representations have the same central charge since they differ only by the finite dimensional factor $`𝕊_𝔱`$. However, the extended spin representation $`𝒮_{L𝔤}`$ admits an action of the full Lie algebra $`\stackrel{~}{L}𝔤`$, and in fact, $`𝒮_{L𝔤}`$ is the direct sum of $`dim𝕊_𝔱`$ copies of an irreducible positive energy representation of $`L𝔤`$. Examining the structure of this representation, the first three energy levels of $`𝒮_{L𝔤}`$ are as follows:
$`𝒮_{L𝔤}(0)`$ $`=𝕊_𝔤,`$
$`𝒮_{L𝔤}(1)`$ $`=𝕊_𝔤𝔤_{},`$
$`𝒮_{L𝔤}(2)`$ $`=𝕊_𝔤𝔤_{}𝕊_𝔤\mathrm{\Lambda }^2(𝔤_{}).`$
Letting $`\alpha `$ denote the highest root of $`𝔤`$, and $`c`$ the central charge of $`𝒮_{L𝔤}`$, the highest weights of $`𝒮_{L𝔤}(0)`$ and $`𝒮_{L𝔤}(1)`$ are then $`(0,\rho ,c)`$ and $`(1,\rho +\alpha ,c)`$ respectively, while the weight $`(2,\rho +2\alpha ,c)`$ is *not* present in $`𝒮_{L𝔤}(2)`$. The weights $`(0,\rho ,c)`$ and $`(1,\rho +\alpha ,c)`$ thus form a complete string of weights for the root $`𝜶=(1,\alpha ,0)`$, and so they must be related to each other by the affine Weyl element $`s_{1,\alpha }`$, the reflection through the hyperplane orthogonal to $`𝜶`$. By (5), the difference of these weights is $`(1,\rho +\alpha ,c)(0,\rho ,c)=𝜶=(0,\rho ,c)(H_{1,\alpha })𝜶`$, so we obtain
$$1=(0,\rho ,c)(H_{1,\alpha })=\rho (H_\alpha )\frac{1}{2}H_\alpha ^2c=\rho ,\alpha c,$$
where $`\frac{1}{2}H_\alpha ^2=1`$ and $`\rho (H_\alpha )=\rho ,\alpha `$ in the basic inner product since $`\alpha `$ is the highest root. The central charge of the spin representation is thus $`c=\rho ,\alpha +1`$.
The quadratic Casimir operator of a Lie algebra does not depend on the choice of orthonormal basis, and it commutes with the action of the Lie algebra. It therefore acts by a constant times the identity on each irreducible representation. On the irreducible representation of highest weight $`\alpha `$, the value of the Casimir operator is $`\frac{1}{2}\alpha ^2+\alpha ,\rho `$. In particular, if $`G`$ is simple, then the adjoint representation is irreducible, and taking $`\alpha `$ to be the highest root of $`G`$, which satisfies $`\alpha ^2=2`$ in the basic inner product, we again obtain the value $`\rho ,\alpha +1`$ as desired. ∎
We can now compute the character of $`𝒮_{L𝔤/𝔱}`$ directly from the decomposition (8) and Theorem 1. Written in terms of the affine roots $`𝜶=(k,\alpha ,0)`$, the character is
$$\chi (𝒮_{L𝔤/𝔱})=u^{c_G}\underset{\alpha >0}{}\left(e^{\frac{i\alpha }{2}}+e^{\frac{i\alpha }{2}}\right)\underset{k>0,\alpha }{}\left(1+e^{i\alpha }z^k\right)=e^{i𝝆_G}\underset{𝜶>0}{}(1+e^{i𝜶}),$$
where $`u`$ is a parameter on the central $`S^1`$ extension in $`\stackrel{~}{L}G`$, and $`𝝆_G=(0,\rho _G,c_G)`$. Here, $`𝝆_G`$ is the lowest weight of $`𝒮_{L𝔤/𝔱}`$, which corresponds to the square root of the determinant in (7). This weight is the loop group version of $`\rho _G`$, half the sum of the positive roots of $`G`$, which is also characterized by the identity $`\rho _G(H_\alpha )=1`$ for each of the simple roots $`\alpha `$ of $`G`$. In the loop group case, the identity $`𝝆_G(H_𝜶)=1`$ must hold as $`𝜶`$ ranges over the simple *affine* roots, including the additional root $`(1,\alpha _{\text{max}},0)`$. However, in our proof of Theorem 1, the condition $`𝝆_G(H_{1,\alpha _{\text{max}}})=1`$ is the same equation (up to sign) that we used to compute the central charge $`c_G`$.
The spin representation decomposes as $`𝒮_{L𝔤/𝔱}=𝒮_{L𝔤/𝔱}^+𝒮_{L𝔤/𝔱}^{}`$ into the sum of two half-spin representations. In particular, since the complement of $`𝔱`$ in $`𝔤`$ is even dimensional, the zero mode factor $`𝕊_{𝔤/𝔱}`$ of $`𝒮_{L𝔤/𝔱}`$ splits into half-spin representations, and the exterior algebra in (8) splits into its even and odd degree components. The difference of the characters of these half-spin representations is
(10)
$$\chi \left(𝒮_{L𝔤/𝔱}^+\right)\chi \left(𝒮_{L𝔤/𝔱}^{}\right)=e^{i𝝆_G}\underset{𝜶>0}{}\left(1e^{i𝜶}\right),$$
which can be viewed either as a supertrace on $`𝒮_{L𝔤/𝔱}`$ or as the character of the virtual representation $`𝒮_{L𝔤/𝔱}^+𝒮_{L𝔤/𝔱}^{}`$. Using the notation of spin representations, the Weyl-Kac character formula becomes
###### Theorem 2 (Weyl-Kac Character Formula).
If $`G`$ is simply connected and simple, then the character of the irreducible positive energy representation $`_𝛌`$ of $`\stackrel{~}{L}G`$ with lowest weight $`𝛌`$ is given by the quotient
(11)
$$\chi (_𝝀)=\frac{_{w𝒲_G}(1)^we^{iw(𝝀𝝆_G)}}{\chi \left(𝒮_{L𝔤/𝔱}^+\right)\chi \left(𝒮_{L𝔤/𝔱}^{}\right)},$$
where $`𝒲_G`$ is the affine Weyl group of $`G`$ and $`𝛒_G=(0,\rho _G,c_G)`$.
Note that as an immediate consequence of the Weyl-Kac character formula, if we consider the trivial representation with $`𝝀=0`$, we obtain the identity
$$\chi \left(𝒮_{L𝔤/𝔱}^+\right)\chi \left(𝒮_{L𝔤/𝔱}^{}\right)=\underset{w𝒲_G}{}(1)^we^{iw(𝝆_G)},$$
which gives an alternative expression for the signed character (10) of the spin representation appearing in the denominator of (11).
###### Remark.
If $`G`$ is semi-simple, then we recall that the universal central extension of $`LG`$ is an extension not by a circle but rather by the torus $`T^d`$, where $`d`$ counts the number of simple components. In this case the central charge of the spin representation is the $`d`$-vector $`𝐜_G=(c_{G_1},\mathrm{},c_{G_d})`$, where $`G_1,\mathrm{},G_d`$ are the simple components of $`G`$. If we work with the universal central extension of $`LG`$ and define $`𝝆_G=(0,\rho _G,𝐜_G)`$, then the Weyl-Kac character formula still holds as written. In fact, using the appropriate universal central extension, the Weyl-Kac character formula continues to hold for an arbitrary compact Lie group $`G`$.
## 3. The homogeneous Weyl-Kac formula
Let $`𝔤`$ be a compact, semi-simple Lie algebra, and let $`𝔥`$ be a reductive subalgebra of maximal rank in $`𝔤`$. In , Gross, Kostant, Ramond, and Sternberg prove a homogeneous generalization of the Weyl character formula, associating to each $`𝔤`$-representation a set of $`𝔥`$-representations with similar properties, called a multiplet.
###### Theorem 3 (Homogeneous Weyl Formula).
Let $`V_\lambda `$ and $`U_\mu `$ denote the irreducible representations of $`𝔤`$ and $`𝔥`$ with highest weights $`\lambda `$ and $`\mu `$ respectively. The following identity holds in the representation ring $`R(𝔥)`$:
(12)
$$V_\lambda 𝕊_{𝔤/𝔥}^+V_\lambda 𝕊_{𝔤/𝔥}^{}=\underset{cC}{}(1)^cU_{c(\lambda +\rho _𝔤)\rho _𝔥},$$
where the sum is taken over the subset $`C`$ of elements $`cW_𝔤`$ of the Weyl group of $`𝔤`$ for which $`c(\lambda +\rho _𝔤)\rho _𝔥`$ are dominant weights of $`𝔥`$.
Note that if $`𝔥=𝔱`$ is a Cartan subalgebra of $`𝔤`$, then $`C`$ is the full Weyl group $`W_𝔤`$, and (12) becomes the Weyl character formula. Also note that by stating this result in terms of the Lie algebras $`𝔥𝔤`$ rather than their corresponding Lie groups $`HG`$, we bypass the issue of whether the spin representation $`𝕊_{𝔤/𝔥}`$ exponentiates to give a true representation of $`H`$. Geometrically, this is equivalent to the condition that $`G/H`$ be a spin manifold.
Theorem 3 has an immediate analogue for loop groups. The only complication is that simply working at the level of Lie algebras is no longer sufficient to avoid the geometric obstruction, which in this case is the condition that $`G/H`$ admit a *string structure* (see ). Rather, we must work with the universal central extensions. Given $`𝔤`$ and $`𝔥`$ as described above, let $`\stackrel{~}{L}𝔤`$ be the universal central extension of $`L𝔤`$, and let $`\stackrel{~}{L}𝔥`$ be the restriction of $`\stackrel{~}{L}𝔤`$ to $`L𝔥`$. Note that $`\stackrel{~}{L}𝔥`$ is not in general the universal central extension of $`L𝔥`$, which we denote by $`\widehat{L}𝔥`$. Rather, $`\stackrel{~}{L}𝔥`$ is a quotient of $`\widehat{L}𝔥`$. Since $`𝔥`$ has the same rank as $`𝔤`$, if $`𝔱`$ is a Cartan subalgebra of $`𝔥`$, then it is likewise a Cartan subalgebra of $`𝔤`$. The Cartan subalgebras of $`\stackrel{~}{}\widehat{L}𝔥`$ and $`\stackrel{~}{}\stackrel{~}{L}𝔤`$ are then $`𝔱^{d_𝔥}`$ and $`𝔱^{d_𝔤}`$ respectively, where $`d_𝔤`$ is the number of simple components of $`𝔤`$ and $`d_𝔥d_𝔤`$. In other words, we have the commutative diagram
$$\begin{array}{ccccc}𝔱^{d_𝔥}& \stackrel{\text{quotient}}{}& 𝔱^{d_𝔤}& =& 𝔱^{d_𝔤}\\ & & & & & & \\ \stackrel{~}{}\widehat{L}𝔥& \stackrel{\text{quotient}}{}& \stackrel{~}{}\stackrel{~}{L}𝔥& \stackrel{\text{inclusion}}{}& \stackrel{~}{}\stackrel{~}{L}𝔤\end{array}$$
where the vertical arrows are inclusions of Cartan subalgebras.
The weights of $`L𝔥`$ and $`L𝔤`$ live in the dual spaces to their Cartan subalgebras, and dual to the quotient map we have an inclusion
$$𝔱^{}^{d_𝔤}𝔱^{}^{d_𝔥}.$$
We may therefore view the weight lattice of $`L𝔤`$ as a subset of the weight lattice of $`L𝔥`$. On the other hand, if we ignore the central extension (i.e., restrict to weights of level 0), then the weight lattices are identical, and the roots of $`L𝔥`$ are a subset of the roots of $`L𝔤`$. Consequently, the affine Weyl group $`𝒲_𝔥`$ of $`𝔥`$, which is generated by the reflections through the hyperplanes orthogonal to the roots of $`L𝔥`$, is a subgroup of the affine Weyl group $`𝒲_𝔤`$ of $`𝔤`$.
###### Theorem 4 (Homogeneous Weyl-Kac Formula).
Let $`_𝛌`$ and $`𝒰_𝛍`$ denote the irreducible positive energy representations of $`\stackrel{~}{L}𝔤`$ and $`\widehat{L}𝔥`$ with lowest weights $`𝛌`$ and $`𝛍`$ respectively. We then have the following identity for virtual representations of $`\widehat{L}𝔥`$:
(13)
$$_𝝀𝒮_{L𝔤/L𝔥}^+_𝝀𝒮_{L𝔤/L𝔥}^{}=\underset{c𝒞}{}(1)^c𝒰_{c(𝝀𝝆_𝔤)+𝝆_𝔥},$$
where the sum is taken over the subset $`𝒞`$ of elements $`c𝒲_𝔤`$ of the affine Weyl group of $`𝔤`$ for which $`c(𝛌𝛒_𝔤)+𝛒_𝔥`$ are anti-dominant weights of $`\widehat{L}𝔥`$.
###### Proof.
We first note that the construction of the spin representation is multiplicative, provided that the underlying vector spaces are even dimensional. In our case, the positive and negative energy subspaces pair off, while for the zero modes, the maximal rank condition implies that the complement of $`𝔥`$ in $`𝔤`$ and the complement of $`𝔱`$ in $`𝔥`$ are even dimensional, so we have
(14)
$$𝒮_{L𝔤/𝔱}^+𝒮_{L𝔤/𝔱}^{}=\left(𝒮_{L𝔤/L𝔥}^+𝒮_{L𝔤/L𝔥}^{}\right)\left(𝒮_{L𝔥/𝔱}^+𝒮_{L𝔥/𝔱}^{}\right).$$
Applying the Weyl-Kac character formula (11) to the left side of (13), and factoring the Weyl-Kac denominator using (14), we obtain
(15)
$$\chi \left(_𝝀𝒮_{L𝔤/L𝔥}^+\right)\chi \left(_𝝀𝒮_{L𝔤/L𝔥}^{}\right)=\frac{_{w𝒲_𝔤}(1)^we^{iw(𝝀𝝆_𝔤)}}{\chi \left(𝒮_{L𝔥/𝔱}^+\right)\chi \left(𝒮_{L𝔥/𝔱}^{}\right)}.$$
Recall that the affine Weyl group acts simply transitively on the Weyl alcoves. Due to the $`𝝆_𝔤`$ shift, the weight $`𝝀𝝆_𝔤`$ lies in the interior of the fundamental Weyl alcove for $`𝔤`$, and thus for any $`w𝒲_𝔤`$, the weight $`w(𝝀𝝆_𝔤)`$ likewise lies in the interior of some Weyl alcove. Furthermore, the Weyl alcoves for $`𝔤`$ are completely contained inside the Weyl alcoves for $`𝔥`$, and so there exists a unique element $`w^{}𝒲_𝔥`$ such that $`w^{}w(𝝀𝝆_𝔤)`$ lies in the interior of the fundamental Weyl alcove for $`𝔥`$. Shifting back by $`𝝆_𝔥`$, we see that the weight $`w^{}w(𝝀𝝆_𝔤)+𝝆_𝔥`$ is anti-dominant for $`\widehat{L}𝔥`$. Putting $`c=w^{}w`$, we can write $`w=(w^{})^1c`$, and more generally we have $`𝒲_𝔤=𝒲_𝔥𝒞`$. Using this decomposition to rewrite the numerator on the right side of (15), we have
$$\begin{array}{cc}\hfill \frac{_{w𝒲_𝔤}(1)^we^{iw(𝝀𝝆_𝔤)}}{\chi \left(𝒮_{L𝔥/𝔱}^+\right)\chi \left(𝒮_{L𝔥/𝔱}^{}\right)}& =\underset{c𝒞}{}(1)^c\frac{_{w𝒲_𝔥}(1)^we^{iwc(𝝀𝝆_𝔤)}}{\chi \left(𝒮_{L𝔥/𝔱}^+\right)\chi \left(𝒮_{L𝔥/𝔱}^{}\right)}\hfill \\ & =\underset{c𝒞}{}(1)^c\chi \left(𝒰_{c(𝝀𝝆_𝔤)+𝝆_𝔥}\right),\hfill \end{array}$$
where the second line follows by applying the Weyl-Kac character formula (11) for $`\widehat{L}𝔥`$. This proves the character form of the identity (13). ∎
The subset $`𝒞𝒲_𝔤`$ appearing in Theorem 4 does not depend on the weight $`𝝀`$. Rather, it consists of all elements of the affine Weyl group of $`𝔤`$ that map the fundamental Weyl alcove for $`𝔤`$ into the fundamental Weyl alcove for $`𝔥`$. Since the affine Weyl group acts simply transitively on the Weyl alcoves, it follows that the cardinality of $`𝒞`$ is the ratio of the volumes of the fundamental alcoves for $`𝔤`$ and $`𝔥`$. Equivalently, the elements of $`𝒞`$ are representatives of the cosets of $`𝒲_𝔥`$ in $`𝒲_𝔤`$, so the cardinality of $`𝒞`$ is the index of $`𝒲_𝔥`$ in $`𝒲_𝔤`$. In particular, the sum appearing in (13) is finite if and only if $`𝔥`$ is semi-simple. In such cases, $`|𝒞|`$ is the index of $`W_𝔥`$ in $`W_𝔤`$, which is also the Euler number of the corresponding homogeneous space $`G/H`$. Examples of pairs $`𝔥𝔤`$ with both $`𝔥`$ and $`𝔤`$ semi-simple include $`D_nB_n`$ with $`|𝒞|=2`$, as well as the case $`B_4F_4`$ with $`|𝒞|=3`$ that prompted . On the other hand, for pairs $`𝔥𝔤`$ corresponding to *complex* homogeneous spaces $`G/H`$, the group $`H`$ must contain a $`\mathrm{U}(1)`$ component, and so (13) is an infinite sum. We note that in the physics literature (see ), the $`N=1`$ superconformal coset models on $`G/H`$ possess an additional $`N=2`$ symmetry precisely when $`𝒞`$ is infinite.
At the other extreme, if $`𝔥=𝔱`$ is a Cartan subalgebra of $`𝔤`$, then $`𝝆_𝔥`$ vanishes, $`𝒞`$ is the full affine Weyl group $`𝒲_𝔤`$, and the homogeneous Weyl-Kac formula becomes
(16)
$$_𝝀𝒮_{L𝔤/L𝔱}^+_𝝀𝒮_{L𝔤/L𝔱}^{}=\underset{w𝒲_𝔤}{}(1)^w𝒰_{w(𝝀𝝆_𝔤)}.$$
This identity is equivalent to the Weyl-Kac character formula (11), but it is expressed slightly differently. Since $`𝔱`$ is abelian, the irreducible positive energy representation $`𝒰_𝝁`$ of $`\stackrel{~}{L}𝔱`$ takes the particularly simple form
$$𝒰_𝝁=\mathrm{Sym}^{}\left(\underset{k>0}{}𝔱_{}z^k\right)=\underset{k>0}{}\mathrm{Sym}^{}\left(𝔱_{}z^k\right),$$
where $`\mathrm{Sym}^{}`$ is the symmetric algebra, and the character of this representation is
(17)
$$\chi (𝒰_𝝁)=e^{i𝝁}\underset{k>0}{}(1+z^k+\mathrm{})^{dim𝔱}=e^{i𝝁}\underset{k>0}{}(1z^k)^{dim𝔱}.$$
On the other hand, the signed character of the spin representation on $`L𝔱/𝔱`$ is
(18)
$$\chi \left(𝒮_{L𝔱/𝔱}^+\right)\chi \left(𝒮_{L𝔱/𝔱}^{}\right)=\underset{k>0}{}(1z^k)^{dim𝔱},$$
since the product in (10) is taken over the positive roots $`(k,0,0)`$, each counted with multiplicity $`dim𝔱`$. In particular, the products in the characters (17) and (18) cancel each other, yielding the Weyl-Kac character formula for $`LT`$:
$$\chi \left(𝒰_𝝁𝒮_{L𝔱/𝔱}^+\right)\chi \left(𝒰_𝝁𝒮_{L𝔱/𝔱}^{}\right)=e^{i𝝁}.$$
So, multiplying the formula (16) by the character (18), we recover the usual form of the Weyl-Kac character formula (11) for $`LG`$.
###### Example.
Take $`𝔤=𝔰𝔲(2)`$ and let $`𝔥=𝔲(1)`$ be the Cartan subalgebra of diagonal elements. In this particular case, we can use the homogeneous Weyl-Kac formula to explicitly compute the character of the entire spin representation $`𝒮_{L𝔤/L𝔥}`$, not just the difference of the two half-spin representations. Here we have $`𝝆_𝔤=(0,\rho _𝔤,c_𝔤)=(0,1,2)`$, and so the lowest weight of the spin representation is $`𝝆_𝔤=(0,1,2)`$. The half-spin representation $`𝒮^+`$ (resp. $`𝒮^{}`$) is obtained by acting on a lowest weight vector by an even (resp. odd) number of Clifford multiplications by the positive generators $`E_+`$ and $`E_\pm z^n`$ for $`n>0`$ of $`L𝔤/L𝔥`$. Since each of these generators shifts the $`𝔰𝔲(2)`$ weight by $`\pm 2`$, the $`𝔰𝔲(2)`$ weights of all elements in $`𝒮^+`$ must be of the form $`4n1`$, while the weights for $`𝒮^{}`$ are all of the form $`4n+1`$. Hence the weights of $`𝒮^+`$ and $`𝒮^{}`$ are distinct, and thus there is no cancellation when we take their difference.
Applying (16) for the case of the trivial representation with $`𝝀=0`$, we obtain
$`𝒮_{L𝔤/L𝔥}^+`$ $`={\displaystyle \underset{w𝒲_𝔤^+}{}}𝒰_{w(0,1,2)}={\displaystyle \underset{n}{}}𝒰_{(2n^2n,4n1,2)},`$
$`𝒮_{L𝔤/L𝔥}^{}`$ $`={\displaystyle \underset{w𝒲_𝔤^{}}{}}𝒰_{w(0,1,2)}={\displaystyle \underset{n}{}}𝒰_{(2n^2+n,4n+1,2)},`$
where we have explicitly written out the action of $`𝒲_{𝔰𝔲(2)}_2`$:
$$w_n^\pm (m,\lambda ,h)=(m\pm \lambda n+hn^2,\pm \lambda +2hn,h).$$
Using (17) for $`\chi (𝒰_𝝁)`$, the characters of the half-spin representations are
$`\chi \left(𝒮_{L𝔤/L𝔥}^+\right)(z,w,u)`$ $`=u^2{\displaystyle \underset{n}{}}w^{4n1}z^{2n^2n}{\displaystyle \underset{k>0}{}}(1z^k)^1,`$
$`\chi \left(𝒮_{L𝔤/L𝔥}^{}\right)(z,w,u)`$ $`=u^2{\displaystyle \underset{n}{}}w^{4n+1}z^{2n^2+n}{\displaystyle \underset{k>0}{}}(1z^k)^1,`$
where the powers of $`z`$, $`w`$, and $`u`$ correspond to the energy, $`𝔰𝔲(2)`$ weight, and level respectively. Combining these half-spin representations, the total spin representation has character
$$\chi \left(𝒮_{L𝔤/L𝔥}\right)(z,w,u)=u^2\underset{n}{}w^{2n1}z^{\frac{1}{2}n(n1)}\underset{k>0}{}(1z^k)^1.$$
The orbit of the lowest weight $`𝝆_𝔤=(0,1,2)`$ under the affine Weyl group consists of all weights $`(m,\lambda ,2)`$ with $`\lambda `$ odd and $`m=\frac{1}{8}(\lambda ^21)`$. This equation sweeps out a parabola, and the remaining weights live inside this parabola, satisfying $`m>\frac{1}{8}(\lambda ^21)`$. The weights of $`𝒮_{L𝔤/L𝔥}`$ are shown in Figure 1, with the orbit of $`𝝆_𝔤`$ drawn as open circles. The multiplicity of any such weight can be derived from (17) and is given by the number of partitions of $`m\frac{1}{8}(\lambda ^21)`$ into positive integers.
## 4. The Clifford algebra $`\mathrm{Cl}(𝔤)`$
Let $`𝔤`$ be a finite dimensional Lie algebra with an $`\mathrm{ad}`$-invariant inner product. Recall that the Clifford algebra $`\mathrm{Cl}(𝔤)`$ is generated by the elements of $`𝔤`$ subject to the anti-commutator relation $`\{X,Y\}=XY+YX=2X,Y`$ for all $`X,Y𝔤`$. There is a natural Clifford action on the exterior algebra $`\mathrm{\Lambda }^{}(𝔤^{})`$, which is given on the generators $`X𝔤`$ by $`c(X)=\iota _X+\epsilon _X^{}`$, where $`\iota _X`$ is interior contraction by $`X𝔤`$ and $`\epsilon _X^{}`$ is exterior multiplication by the dual element $`X^{}𝔤^{}`$ satisfying $`X^{}(Y)=X,Y`$. Using the distinguished element $`1`$ of the exterior algebra, the map $`xc(x)1`$ gives an isomorphism $`\mathrm{Cl}(𝔤)\mathrm{\Lambda }^{}(𝔤^{})`$ of left $`\mathrm{Cl}(𝔤)`$-modules, called the Chevalley identification. We may therefore view the Clifford algebra as the exterior algebra $`\mathrm{\Lambda }^{}(𝔤^{})`$ with the alternative multiplication
(19)
$$X^{}\eta =X^{}\eta +\iota _X\eta $$
for $`X𝔤`$ and $`\eta \mathrm{\Lambda }^{}(𝔤^{})`$.
Consider the graded Lie superalgebra $`\widehat{𝔤}=𝔤_1𝔤_0_1`$, where the subscript denotes the integer grading. The exterior algebra $`\mathrm{\Lambda }^{}(𝔤^{})`$ is a representation of this Lie superalgebra $`\widehat{𝔤}`$, with $`𝔤_1`$ acting by interior contraction, $`𝔤_0`$ acting by the coadjoint action, and the generator $`d_1`$ acting as the exterior derivative. On the generators $`\xi 𝔤^{}`$, these operators are given by
$`\iota _X\xi `$ $`=\xi (X)`$ $`\iota _X`$ $`:\mathrm{\Lambda }^k(𝔤^{})\mathrm{\Lambda }^{k1}(𝔤^{})`$
$`(\mathrm{ad}_X^{}\xi )(Y)`$ $`=\xi (\mathrm{ad}_XY)`$ $`\mathrm{ad}_X^{}`$ $`:\mathrm{\Lambda }^k(𝔤^{})\mathrm{\Lambda }^{k+0}(𝔤^{})`$
$`(d\xi )(X,Y)`$ $`=\frac{1}{2}\xi ([X,Y])`$ $`d`$ $`:\mathrm{\Lambda }^k(𝔤^{})\mathrm{\Lambda }^{k+1}(𝔤^{})`$
for $`X,Y𝔤`$. These operators then extend as super-derivations to the full exterior algebra, and they satisfy the identities $`[\mathrm{ad}_X^{},\iota _Y]=\iota _{[X,Y]}`$ and $`\{d,\iota _X\}=\mathrm{ad}_X^{}`$. If we perturb this action by taking $`d^{}=d\iota _\mathrm{\Omega }^{}`$, where $`\mathrm{\Omega }`$ is a closed $`𝔤`$-invariant form of odd degree, then the commutation relations on $`\widehat{𝔤}`$ are unchanged.
###### Theorem 5.
Using the Chevalley identification, the action of $`\widehat{𝔤}=𝔤_1𝔤_0_1`$ on $`\mathrm{\Lambda }^{}(𝔤^{})`$ can be expressed in terms of the adjoint action of the Clifford algebra as
(20) $`\mathrm{ad}X^{}`$ $`=2\iota _X`$ $`X^{}`$ $`\mathrm{\Lambda }^1(𝔤^{})`$
(21) $`\mathrm{ad}dX^{}`$ $`=2\mathrm{ad}_X^{}`$ $`dX^{}`$ $`\mathrm{\Lambda }^2(𝔤^{})`$
(22) $`\mathrm{ad}\mathrm{\Omega }`$ $`=2d2\iota _\mathrm{\Omega }^{}`$ $`\mathrm{\Omega }`$ $`\mathrm{\Lambda }^3(𝔤^{})`$
where $`\mathrm{\Omega }`$ is the fundamental 3-form given by $`\mathrm{\Omega }(X,Y,Z)=\frac{1}{6}X,[Y,Z]`$.
###### Proof.
First, we show that the operators $`\iota _X`$ and $`\mathrm{ad}_X^{}`$ are super-derivations with respect to the Clifford multiplication (19). For $`X,Y𝔤`$ and $`\eta \mathrm{\Lambda }^{}(𝔤^{})`$, we have
$`\iota _X(Y^{}\eta )`$ $`=\iota _X(Y^{}\eta )+\iota _X\iota _Y\eta `$
$`=(\iota _XY^{})\eta Y^{}\iota _X\eta \iota _Y\iota _X\eta `$
$`=(\iota _XY^{})\eta Y^{}\iota _X\eta ,`$
$`\mathrm{ad}_X^{}(Y^{}\eta )`$ $`=\mathrm{ad}_X^{}(Y^{}\eta )+\mathrm{ad}_X^{}\iota _Y\eta `$
$`=(\mathrm{ad}_X^{}Y^{})\eta +Y^{}\mathrm{ad}_X^{}\eta +\iota _Y\mathrm{ad}_X^{}\eta +\iota _{[X,Y]}\eta `$
$`=(\mathrm{ad}_X^{}Y^{})\eta +Y^{}\mathrm{ad}_X^{}\eta .`$
Now, to prove the identities (20) and (21), we need only verify them for the generators $`𝔤^{}=\mathrm{\Lambda }^1(𝔤^{})`$, but it follows from the definition of the Clifford algebra that
$$\{X^{},Y^{}\}=2X,Y=2\iota _XY^{},$$
and by applying (20) and the identity $`\{d,\iota _X\}=\mathrm{ad}_X^{}`$, we obtain
$$[dX^{},Y^{}]=2\iota _YdX^{}=2\mathrm{ad}_Y^{}X^{}=2\mathrm{ad}_X^{}Y^{}.$$
To prove (22), we first verify that it holds when acting on a generator $`X^{}𝔤^{}`$:
$$\{\mathrm{\Omega },X^{}\}(Y,Z)=(2\iota _X\mathrm{\Omega })(Y,Z)=X,[Y,Z]=(2dX^{})(Y,Z).$$
Finally we show that $`d^{}=d\iota _\mathrm{\Omega }^{}`$ is a super-derivation for Clifford multiplication:
$$\begin{array}{cc}\hfill d^{}(X^{}\eta )& =d(X^{}\eta )\iota _\mathrm{\Omega }^{}(X^{}\eta )+d\iota _X\eta \iota _\mathrm{\Omega }^{}\iota _X\eta \hfill \\ & =(dX^{})\eta X^{}d\eta \iota _{(dX^{})^{}}\eta +X^{}\iota _\mathrm{\Omega }^{}\eta \hfill \\ & \iota _Xd\eta +\mathrm{ad}_X^{}\eta +\iota _X\iota _\mathrm{\Omega }^{}\eta \hfill \\ & =(d^{}X^{})\eta X^{}d^{}\eta ,\hfill \end{array}$$
where we use the expansion $`(d^{}X^{})\eta =(dX^{})\eta =(dX^{})\eta +\mathrm{ad}_X^{}\eta \iota _{(dX^{})^{}}\eta .`$
Although the Clifford algebra $`\mathrm{Cl}(𝔤)`$ does not admit an integer grading, it does have the distinguished subspaces $`𝔤`$ and $`𝔰𝔭𝔦𝔫(𝔤)`$, which correspond via the Chevalley identification to the first two degrees of the exterior algebra:
$$\mathrm{\Lambda }^1(𝔤^{})𝔤\mathrm{Cl}(𝔤),\mathrm{\Lambda }^2(𝔤^{})𝔰𝔭𝔦𝔫(𝔤)\mathrm{Cl}(𝔤).$$
Since $`\mathrm{Spin}(𝔤)`$ is the double cover of $`\mathrm{SO}(𝔤)`$, there is a Lie algebra isomorphism $`𝔰𝔭𝔦𝔫(𝔤)𝔰𝔬(𝔤)`$, and given any element $`a𝔰𝔬(𝔤)`$, the corresponding element of $`\stackrel{~}{a}𝔰𝔭𝔦𝔫(𝔤)`$ is uniquely determined by the identity $`[\stackrel{~}{a},X^{}]=(aX)^{}`$ for all $`X𝔤`$. In particular, the adjoint action $`\mathrm{ad}:𝔤𝔰𝔬(𝔤)`$ lifts to a Lie algebra homomorphism $`\stackrel{~}{\mathrm{a}}\mathrm{d}:𝔤𝔰𝔭𝔦𝔫(𝔤)`$ satisfying
$$[\stackrel{~}{\mathrm{a}}\mathrm{d}X,Y^{}]=(\mathrm{ad}_XY)^{}=\mathrm{ad}_X^{}Y^{}.$$
However, from the identity (21), we see that the spin lift of the adjoint action must be $`\stackrel{~}{\mathrm{a}}\mathrm{d}X=\frac{1}{2}dX^{}`$.
Let $`\{X_i\}`$ be a basis of $`𝔤`$, and let $`\{X_i^{}\}`$ denote the corresponding dual basis of $`𝔤`$ satisfying $`X_i^{},X_j=\delta _{ij}`$. In terms of this basis, the map $`\stackrel{~}{\mathrm{a}}\mathrm{d}:X\frac{1}{2}dX^{}`$ is
(23)
$$\stackrel{~}{\mathrm{a}}\mathrm{d}X=\frac{1}{4}\underset{i}{}X_i^{}[X,X_i],$$
while the element $`\gamma =\frac{1}{4}\mathrm{\Omega }`$ corresponding to the fundamental 3-form is given by
(24)
$$\gamma =\frac{1}{24}\underset{i,j}{}X_i^{}X_j^{}[X_i,X_j]=\frac{1}{6}\underset{i}{}X_i^{}\stackrel{~}{\mathrm{a}}\mathrm{d}X_i.$$
Rewriting Theorem 5 in terms of this new notation, we obtain the following:
###### Corollary 6.
The elements $`1`$, $`X`$, $`\stackrel{~}{\mathrm{a}}\mathrm{d}X`$, $`\gamma `$ for $`X𝔤`$ span a Lie superalgebra $`_+𝔤_{}𝔤_+_{}`$ in the Clifford algebra $`\mathrm{Cl}(𝔤)`$ with the commutation relations
$`[\stackrel{~}{\mathrm{a}}\mathrm{d}X,Y]`$ $`=[X,Y],`$ $`[\stackrel{~}{\mathrm{a}}\mathrm{d}X,\stackrel{~}{\mathrm{a}}\mathrm{d}Y]`$ $`=\stackrel{~}{\mathrm{a}}\mathrm{d}[X,Y],`$ $`[\stackrel{~}{\mathrm{a}}\mathrm{d}X,\gamma ]`$ $`=0,`$
$`\{X,Y\}`$ $`=2X,Y,`$ $`\{\gamma ,X\}`$ $`=\stackrel{~}{\mathrm{a}}\mathrm{d}X,`$ $`\{\gamma ,\gamma \}`$ $`=\frac{1}{24}\mathrm{tr}_𝔤\mathrm{\Delta }_{\mathrm{ad}}^𝔤,`$
where $`\mathrm{\Delta }_{\mathrm{ad}}^𝔤=\frac{1}{2}_i\mathrm{ad}_{X_i^{}}\mathrm{ad}_{X_i}`$ is the quadratic Casimir operator.
###### Proof.
All of the commutation relations follow immediately from Theorem 5 and the above discussion with the exception of that for $`\{\gamma ,\gamma \}`$. For example, we derive
$`[\stackrel{~}{\mathrm{a}}\mathrm{d}X,\stackrel{~}{\mathrm{a}}\mathrm{d}Y]`$ $`=\frac{1}{2}\mathrm{ad}_X^{}dY^{}=\frac{1}{2}d\mathrm{ad}_X^{}Y^{}=\stackrel{~}{\mathrm{a}}\mathrm{d}[X,Y],`$
$`[\stackrel{~}{\mathrm{a}}\mathrm{d}X,\gamma ]`$ $`=[\frac{1}{4}\mathrm{\Omega },\frac{1}{2}dX^{}]=\frac{1}{4}ddX^{}=0.`$
To compute $`\{\gamma ,\gamma \}`$, we note that the fundamental 3-form is closed, so we have
$$\begin{array}{c}\hfill \{\gamma ,\gamma \}=\{\frac{1}{4}\mathrm{\Omega },\frac{1}{4}\mathrm{\Omega }\}=\frac{1}{8}d\mathrm{\Omega }\frac{1}{8}\iota _\mathrm{\Omega }^{}\mathrm{\Omega }=\frac{1}{8}\mathrm{\Omega },\mathrm{\Omega }.\end{array}$$
Written in terms of an orthonormal basis $`\{X_i\}`$ for $`𝔤`$, the fundamental 3-form is
$$\mathrm{\Omega }=\underset{i<j<k}{}X_i,[X_j,X_k]X_i^{}X_j^{}X_k^{},$$
and so its norm is given by
$$\begin{array}{cc}\hfill \mathrm{\Omega },\mathrm{\Omega }& =\frac{1}{6}\underset{i,j,k}{}X_i,[X_j,X_k]^2=\frac{1}{6}\underset{j,k}{}[X_j,X_k],[X_j,X_k]\hfill \\ & =\frac{1}{6}\underset{j,k}{}X_k,[X_j,[X_j,X_k]]=\frac{1}{3}\mathrm{tr}_𝔤\mathrm{\Delta }_{\mathrm{ad}}^𝔤,\hfill \end{array}$$
which yields the desired anti-commutator $`\{\gamma ,\gamma \}=\frac{1}{24}\mathrm{tr}_𝔤\mathrm{\Delta }_{\mathrm{ad}}^𝔤`$. ∎
Note that the map $`\stackrel{~}{\mathrm{a}}\mathrm{d}:𝔤\mathrm{Cl}(𝔤)`$ is not necessarily injective; rather its kernel is the center of $`𝔤`$. So, Corollary 6 actually gives us an inclusion of the superalgebra
(25)
$$\stackrel{~}{𝔤}:=𝔤[𝔤,𝔤]𝔤,[𝔤,𝔤]\mathrm{\Lambda }^{}(𝔤^{})\mathrm{Cl}(𝔤)$$
into the Clifford algebra of $`𝔤`$. Also note that this Lie superalgebra $`\stackrel{~}{𝔤}`$, with the commutation relations given by Corollary 6, is the quantized form of the graded Lie superalgebra $`\widehat{𝔤}=𝔤_1𝔤_0_1`$ discussed above.
## 5. The Dirac operator on $`𝔤`$
Let $`𝔤`$ be a Lie algebra with an $`\mathrm{ad}`$-invariant inner product, and let $`𝕊_𝔤`$ be the complex spin representation of the Clifford algebra $`\mathrm{Cl}(𝔤)`$. If $`𝔤`$ is even dimensional then we have $`\mathrm{Cl}(𝔤)\mathrm{End}(𝕊_𝔤)`$, and in general the spin representation $`\mathrm{Cl}(𝔤)\mathrm{End}(𝕊_𝔤)`$ is faithful. To simplify our notation, in the following we implicitly identify $`\mathrm{Cl}(𝔤)`$ with its image in $`\mathrm{End}(𝕊_𝔤)`$ under the spin representation. We recall from the previous section that the adjoint action $`\mathrm{ad}`$ of $`𝔤`$ on itself lifts to the representation $`\stackrel{~}{\mathrm{a}}\mathrm{d}:𝔤\mathrm{Cl}(𝔤)`$ given by (23).
Let $`V`$ be an arbitrary $`𝔤`$-module where the $`𝔤`$-action is the map $`r:𝔤\mathrm{End}(V)`$. Alternatively, this representation $`r`$ may be viewed as the $`\mathrm{End}(V)`$-valued 1-form $`\widehat{r}\mathrm{End}(V)\mathrm{\Lambda }^{}(𝔤^{})`$ given tautologically by $`\widehat{r}(X)=r(X)`$ for all $`X𝔤`$. Identifying $`\mathrm{\Lambda }^{}(𝔤^{})`$ with $`\mathrm{Cl}(𝔤)`$ via the Chevalley map, the element $`\widehat{r}\mathrm{End}(V)\mathrm{Cl}(𝔤)`$ becomes an operator on the tensor product $`V𝕊_𝔤`$. Perturbing this operator slightly we define the Dirac operator $`\overline{)}_r`$ on $`V𝕊_𝔤`$ to be the element
$$\overline{)}_r:=\widehat{r}+1\frac{1}{2}\mathrm{\Omega }\mathrm{End}(V)\mathrm{Cl}(𝔤),$$
where $`\mathrm{\Omega }\mathrm{Cl}(𝔤)`$ is the cubic term given by (22). Written in terms of a basis $`\{X_i\}`$ for $`𝔤`$ and the dual basis $`\{X_i^{}\}`$ satisfying $`X_i^{},X_j=\delta _{i,j}`$, this Dirac operator is
$$\begin{array}{cc}\hfill \overline{)}_r& =\underset{i}{}X_i^{}r(X_i)\frac{1}{12}\underset{i,j}{}X_i^{}X_j^{}[X_i,X_j]\hfill \\ & =\underset{i}{}X_i^{}\left(r(X_i)+\frac{1}{3}\stackrel{~}{\mathrm{a}}\mathrm{d}X_i\right).\hfill \end{array}$$
Note that the second form of this operator resembles a geometric Dirac operator for the connection $`_X=r(X)+\frac{1}{3}\stackrel{~}{\mathrm{a}}\mathrm{d}X`$. Indeed, if $`r`$ is the right action of $`𝔤`$ on functions, then this is the *reductive connection* on the spin bundle over $`G`$ (see ).
Rather than choosing a particular representation $`V`$, we can instead take $`r`$ to be the canonical inclusion $`r:𝔤U(𝔤)`$ of $`𝔤`$ into its universal enveloping algebra $`U(𝔤)`$. This gives us a universal Dirac operator $`\overline{)}`$, which is an element of the non-abelian Weil algebra $`U(𝔤)\mathrm{Cl}(𝔤)`$, introduced by Alekseev and Meinrenken in . Again identifying $`\mathrm{Cl}(𝔤)`$ with $`\mathrm{\Lambda }^{}(𝔤^{})`$, the element $`\overline{)}`$ is characterized by the identity
(26)
$$\iota _X\overline{)}=\varrho (X):=r(X)1+1\stackrel{~}{\mathrm{a}}\mathrm{d}X$$
for all $`X𝔤`$, where $`\varrho =r1+1\stackrel{~}{\mathrm{a}}\mathrm{d}`$ is the diagonal action of $`𝔤`$ on $`U(𝔤)\mathrm{Cl}(𝔤)`$. Squaring the Dirac operator, we obtain the Weitzenböck formula
(27)
$$\begin{array}{cc}\hfill \overline{)}^2& =\widehat{r}\widehat{r}+\{\widehat{r},\frac{1}{2}\mathrm{\Omega }\}+\frac{1}{2}\mathrm{\Omega }\frac{1}{2}\mathrm{\Omega }\hfill \\ & =\widehat{r},\widehat{r}+\widehat{r}\widehat{r}+d\widehat{r}+\frac{1}{8}\{\mathrm{\Omega },\mathrm{\Omega }\}=2\mathrm{\Delta }_r^𝔤\frac{1}{12}\mathrm{tr}_𝔤\mathrm{\Delta }_{\mathrm{ad}}^𝔤,\hfill \end{array}$$
where the “curvature” term $`d\widehat{r}+\widehat{r}\widehat{r}`$ vanishes since
$$(d\widehat{r}+\widehat{r}\widehat{r})(X,Y)=\frac{1}{2}\left(r([X,Y])+[r(X),r(Y)]\right)=0.$$
Note that the square of the Dirac operator has no Clifford algebra component and is thus an element $`\overline{)}^2U(𝔤)`$ of the universal enveloping algebra. In fact, since the Casimir operator commutes with the $`𝔤`$-action, the element $`\overline{)}^2`$ lies in the center of $`U(𝔤)`$. Considering the Dirac operator itself, given any $`X𝔤`$ we have
(28)
$$[\varrho (X),\overline{)}]=[\iota _X\overline{)},\overline{)}]=\iota _X(\overline{)}\overline{)})=0,$$
and thus $`\overline{)}`$ is invariant under the diagonal action $`\varrho `$ of $`𝔤`$ on $`U(𝔤)\mathrm{Cl}(𝔤)`$. We can summarize the above results by stating that the Lie superalgebra
$$(1𝔤)\varrho (𝔤)\overline{)}\mathrm{\Delta }^𝔤U(𝔤)\mathrm{Cl}(𝔤)$$
is a central extension of the Lie superalgebra $`\stackrel{~}{𝔤}`$ from Corollary 6 by the span of the quadratic Casimir operator $`\mathrm{\Delta }^𝔤`$. The commutation relations in this extension are the same as in Corollary 6, with the exception of the square of the Dirac operator which is given by (27). To obtain the corresponding “classical” algebra, we let this superalgebra act on the non-abelian Weil algebra via the adjoint action. Since the elements $`1`$ and $`\mathrm{\Delta }^𝔤`$ lie in the center of the universal enveloping algebra, we are left with the graded Lie superalgebra $`\widehat{𝔤}`$, as Alekseev and Meinrenken show in .
###### Theorem 7.
The non-abelian Weil algebra $`U(𝔤)\mathrm{Cl}(𝔤)`$ is a representation of the graded Lie superalgebra $`\widehat{𝔤}=𝔤_1𝔤_0_1`$ spanned by the operators $`\iota _X`$, $`_X`$, $`d`$ for $`X𝔤`$ given by
$$\iota _X=\mathrm{ad}\left(\frac{1}{2}X\right),_X=\mathrm{ad}\left(\varrho (X)\right),d=\mathrm{ad}\left(\overline{)}\right).$$
Now, suppose that $`𝔤`$ is reductive. If $`V_\lambda `$ is the irreducible representation of $`𝔤`$ with highest weight $`\lambda `$, then the value of the quadratic Casimir operator $`\mathrm{\Delta }_\lambda ^𝔤`$ on $`V_\lambda `$ is
$$\mathrm{\Delta }_\lambda ^𝔤=\frac{1}{2}\left(\lambda +\rho _𝔤^2\rho _𝔤^2\right).$$
In addition, for reductive Lie algebras we have the identity $`\frac{1}{12}\mathrm{tr}_𝔤\mathrm{\Delta }_{\mathrm{ad}}^𝔤=\rho _𝔤^2`$, and it follows that the square of the Dirac operator $`\overline{)}_\lambda `$ acting on $`V_\lambda 𝕊_𝔤`$ is simply the constant $`\overline{)}_\lambda ^2=\lambda +\rho _𝔤^2`$.
## 6. The Dirac operator on $`𝔤/𝔥`$
Let $`𝔥`$ be a Lie subalgebra of $`𝔤`$, and let $`𝔭`$ be the orthogonal complement of $`𝔥`$ with respect to the $`\mathrm{ad}`$-invariant inner product on $`𝔤`$. The adjoint action of $`𝔥`$ on $`𝔤=𝔥𝔭`$ respects this decomposition, so we obtain $`𝔥`$-representations $`\mathrm{ad}_𝔥`$ and $`\mathrm{ad}_𝔭`$ on $`𝔥`$ and $`𝔭`$ respectively. The Clifford algebra also decomposes into the product $`\mathrm{Cl}(𝔤)\mathrm{Cl}(𝔥)\mathrm{Cl}(𝔭)`$ of two Clifford algebras, and with it the spin lift of the adjoint action becomes the sum $`\stackrel{~}{\mathrm{a}}\mathrm{d}_𝔤=\stackrel{~}{\mathrm{a}}\mathrm{d}_𝔥1+1\stackrel{~}{\mathrm{a}}\mathrm{d}_𝔭`$ of separate spin actions $`\stackrel{~}{\mathrm{a}}\mathrm{d}_𝔥:𝔥\mathrm{Cl}(𝔥)`$ and $`\stackrel{~}{\mathrm{a}}\mathrm{d}_𝔭:𝔥\mathrm{Cl}(𝔭)`$. The spin representations $`𝕊_𝔥`$ and $`𝕊_𝔭`$ of $`\mathrm{Cl}(𝔥)`$ and $`\mathrm{Cl}(𝔭)`$ are therefore representations of $`𝔥`$, and if one or both of $`𝔥`$ or $`𝔭`$ is even dimensional, then we have $`𝕊_𝔤𝕊_𝔥𝕊_𝔭`$.
Let $`\overline{)}_𝔤U(𝔤)\mathrm{Cl}(𝔤)`$ denote the universal Dirac operator on $`𝔤`$ discussed in the previous section. Now consider the twisted Dirac operator $`\overline{)}_𝔥^{}`$ on $`𝔥`$ given by
$$\overline{)}_𝔥^{}:=\widehat{r}_𝔥^{}+\frac{1}{2}\mathrm{\Omega }_𝔥\left(U(𝔥)\mathrm{Cl}(𝔭)\right)\mathrm{Cl}(𝔥)U(𝔥)\mathrm{Cl}(𝔤)$$
where $`r^{}=r1+1\stackrel{~}{\mathrm{a}}\mathrm{d}_𝔭`$ is the diagonal action of $`𝔥`$ on $`U(𝔥)\mathrm{Cl}(𝔭)`$. In other words, given any representation $`U`$ of $`𝔥`$, this twisted Dirac operator $`\overline{)}_𝔥^{}`$ is the usual Dirac operator $`\overline{)}_𝔥`$ acting on the twisted space $`(U𝕊_𝔭)𝕊_𝔥U𝕊_𝔤`$. As we saw in (26), this Dirac operator $`\overline{)}_𝔥^{}`$ is characterized by the identity
$$\iota _Z\overline{)}_𝔥^{}=\varrho _𝔥^{}(Z)=\varrho _𝔤(Z)=\iota _Z\overline{)}_𝔤$$
for all $`Z𝔥`$, where $`\varrho _𝔥^{}`$ is the diagonal action of $`𝔥`$ on $`\left(U(𝔥)\mathrm{Cl}(𝔭)\right)\mathrm{Cl}(𝔥)`$. Note that $`\varrho _𝔥^{}`$ is just the restriction to $`𝔥`$ of the diagonal action $`\varrho _𝔤`$ of $`𝔤`$ on $`U(𝔤)\mathrm{Cl}(𝔤)`$. It then follows from (28) that the element $`\overline{)}_𝔥^{}`$ commutes with the diagonal action $`\varrho _𝔥^{}`$.
Define $`\overline{)}_{𝔤/𝔥}=\overline{)}_𝔤\overline{)}_𝔥^{}`$ to be the difference of these two operators. This element $`\overline{)}_{𝔤/𝔥}`$ is then *basic* with respect to $`\widehat{𝔥}`$, or in other words it satisfies the identities
$$\iota _Z\overline{)}_{𝔤/𝔥}=0,_Z\overline{)}_{𝔤/𝔥}=[\varrho _h^{}(Z),\overline{)}_{𝔤/𝔥}]=0,$$
for all $`Z𝔥`$. This Dirac operator can also be written as the element
$$\overline{)}_{𝔤/𝔥}=\widehat{r}_𝔭+1\frac{1}{2}\mathrm{\Omega }_𝔭\left(U(𝔤)\mathrm{Cl}(𝔭)\right)^𝔥,$$
where $`\widehat{r}_𝔭U(𝔤)\mathrm{\Lambda }^1(𝔭^{})`$ corresponds to the map $`r:𝔭U(𝔤)`$, and $`\mathrm{\Omega }_𝔭\mathrm{\Lambda }^3(𝔭^{})`$ is the fundamental 3-form given by $`\mathrm{\Omega }_𝔭(X,Y,Z)=\frac{1}{6}X,[Y,Z]`$ for all $`X,Y,Z𝔭`$. To keep track of the cubic terms, note that the fundamental 3-form decomposes as
$$\mathrm{\Omega }_𝔤=\mathrm{\Omega }_𝔥+\mathrm{\Omega }_𝔭+3\mathrm{\Omega }_{𝔥𝔭𝔭},$$
into its projections onto $`\mathrm{\Lambda }^3(𝔥^{})`$, $`\mathrm{\Lambda }^3(𝔭^{})`$, and $`𝔥^{}𝔭^{}𝔭^{}`$ respectively. This extra contribution corresponds to the twisted term $`\widehat{r}_𝔥^{}\widehat{r}_𝔥=\frac{3}{2}\mathrm{\Omega }_{𝔥𝔭𝔭}`$ appearing in $`\overline{)}_𝔥^{}`$.
If $`\{X_i\}`$ and $`\{X_i^{}\}`$ are dual bases for $`𝔭`$, then the Dirac operator $`\overline{)}_{𝔤/𝔥}`$ is
$$\begin{array}{cc}\hfill \overline{)}_{𝔤/𝔥}& =\underset{i}{}X_i^{}r(X_i)\frac{1}{12}\underset{i,j}{}X_i^{}X_j^{}[X_i,X_j]_𝔭\hfill \\ & =\underset{i}{}X_i^{}\left(r(X_i)+\frac{1}{3}\stackrel{~}{\mathrm{a}}\mathrm{d}_𝔭X_i\right),\hfill \end{array}$$
where $`[X,Y]_𝔭`$ for $`X,Y𝔭`$ denotes the projection of $`[X,Y]`$ onto $`𝔭`$, and
$$\stackrel{~}{\mathrm{a}}\mathrm{d}_𝔭X=\frac{1}{4}\underset{i}{}X_i^{}[X,X_i]_𝔭$$
for $`X𝔭`$. (Note that all of the sums here are taken over a basis of $`𝔭`$, not of $`𝔤`$.) The geometric version of this Dirac operator $`\overline{)}_{𝔤/𝔥}`$, viewed as an operator on twisted spinors on the homogeneous space $`G/H`$, is discussed in and .
To compute the square of $`\overline{)}_{𝔤/𝔥}`$, we first show that $`\overline{)}_𝔥^{}`$ and $`\overline{)}_{𝔤/𝔥}`$ decouple,
$$\{\overline{)}_𝔥^{},\overline{)}_{𝔤/𝔥}\}=\{\widehat{r}_𝔥^{},\overline{)}_{𝔤/𝔥}\}+\{\frac{1}{2}\mathrm{\Omega }_𝔥,\overline{)}_{𝔤/𝔥}\}=[r^{}(),\overline{)}_{𝔤/𝔥}]\stackrel{~}{}+d_𝔥\overline{)}_{𝔤/𝔥}=0.$$
We therefore have
$$\overline{)}_{𝔤/𝔥}^2=(\overline{)}_𝔤)^2(\overline{)}_𝔥^{})^2=2\left(\mathrm{\Delta }_r^𝔤\mathrm{\Delta }_r^{}^𝔥\right)\frac{1}{12}\left(\mathrm{tr}_𝔤\mathrm{\Delta }_{\mathrm{ad}}^𝔤\mathrm{tr}_𝔥\mathrm{\Delta }_{\mathrm{ad}}^𝔥\right).$$
Suppose that both $`𝔤`$ and $`𝔥`$ are reductive. If $`r:𝔤\mathrm{End}(V_\lambda )`$ is the irreducible representation of $`𝔤`$ with highest weight $`\lambda `$, then $`\overline{)}_{𝔤/𝔥}`$ is an $`𝔥`$-invariant operator on $`V_\lambda 𝕊_𝔭`$. Its square then takes the value
(29)
$$\overline{)}_{𝔤/𝔥}^2|_\mu =\lambda +\rho _𝔤^2+\mu +\rho _𝔥^2$$
on the $`𝔥`$-invariant subspace of $`V_\lambda 𝕊_𝔭`$ transforming like the irreducible representation $`U_\mu `$ of $`𝔥`$ with highest weight $`\mu `$. It follows that the kernel of $`\overline{)}_{𝔤/𝔥}^2`$, which is in turn the kernel of the Dirac operator $`\overline{)}_{𝔤/𝔥}`$ itself, consists of all $`𝔥`$-invariant subspaces of $`V_\lambda 𝕊_𝔭`$ transforming like $`U_\mu `$, where $`\mu `$ is a dominant weight of $`𝔥`$ satisfying $`\mu +\rho _𝔥^2=\lambda +\rho _𝔤^2`$. As we show in the following theorem, these subspaces are precisely the multiplet of $`𝔥`$-representations corresponding to the $`𝔤`$-representation $`V_\lambda `$, which we discussed in §3.
###### Theorem 8.
Let $`𝔤`$ be a semi-simple Lie algebra with a maximal rank reductive Lie subalgebra $`𝔥`$, and let $`V_\lambda `$ and $`U_\mu `$ denote the irreducible representations of $`𝔤`$ and $`𝔥`$ with highest weights $`\lambda `$ and $`\mu `$. The kernel of the Dirac operator $`\overline{)}_{𝔤/𝔥}`$ on $`V_\lambda 𝕊_𝔭`$ is
$$\mathrm{Ker}\overline{)}_{𝔤/𝔥}=\underset{cC}{}U_{c\lambda },$$
where $`c\lambda =c(\lambda +\rho _𝔤)\rho _𝔥`$, and $`CW_𝔤`$ is the subset of Weyl elements which map the positive Weyl chamber for $`𝔤`$ into the positive Weyl chamber for $`𝔥`$.
###### Proof.
Since the Weyl group acts by isometries, the weights $`c\lambda `$ satisfy the identity
$$(c\lambda )+\rho _𝔥^2=c(\lambda +\rho _𝔤)^2=\lambda +\rho _𝔤^2,$$
and it follows from (29) that the Dirac operator $`\overline{)}_{𝔤/𝔥}`$ vanishes on any $`𝔥`$-invariant subspace of $`V_\lambda 𝕊_𝔭`$ transforming like $`U_{c\lambda }`$. To complete the proof, it remains to show that each of these representations occurs exactly once in the domain of the Dirac operator and that no other $`𝔥`$-representations appear in its kernel. We establish these facts in the following two lemmas (see also ). ∎
###### Lemma 9.
For each $`cC`$, the irreducible representation $`U_{c\lambda }`$ of $`𝔥`$ with highest weight $`c\lambda =c(\lambda +\rho _𝔤)\rho _𝔥`$ occurs exactly once in the decomposition of $`V_\lambda 𝕊_𝔭`$.
###### Proof.
The highest weight space of an irreducible representation of $`𝔤`$ is always one dimensional, so the weight $`\lambda `$ appears with multiplicity 1 in $`V_\lambda `$. Now consider the complex spin representation $`𝕊_{𝔤/𝔱}`$ associated to the orthogonal complement of a Cartan subalgebra $`𝔱`$ in $`𝔤`$. Given a positive root system for $`𝔤`$, the character of this spin representation is $`\chi (𝕊_{𝔤/𝔱})=_{\alpha >0}(e^{i\alpha /2}+e^{i\alpha /2})`$, and so the highest weight $`\rho _𝔤=\frac{1}{2}_{\alpha >0}\alpha `$ of $`𝕊_{𝔤/𝔱}`$ appears with multiplicity 1. The highest weight of the tensor product $`V_\lambda 𝕊_{𝔤/𝔱}`$ is then $`\lambda +\rho _𝔤`$, appearing with multiplicity 1, and likewise the weights $`w(\lambda +\rho _𝔤)`$ for $`wW_𝔤`$ all have multiplicity 1 in $`V_\lambda 𝕊_{𝔤/𝔱}`$.
Choosing a common Cartan subalgebra $`𝔱𝔥𝔤`$, the spin representation factors as $`𝕊_{𝔤/𝔱}𝕊_𝔭𝕊_{𝔥/𝔱}`$. As we noted above, the weight $`\rho _𝔥`$ appears with multiplicity 1 in the second factor $`𝕊_{𝔥/𝔱}`$. It follows that the weights $`w\lambda `$ for $`wW_𝔤`$ can appear at most once in the tensor product $`V_\lambda 𝕊_𝔭`$, as each such weight contributes one weight of the form $`(w\lambda )+\rho _𝔥=w(\lambda +\rho _𝔤)`$ to the tensor product $`V_\lambda 𝕊_𝔭𝕊_{𝔥/𝔱}V_\lambda 𝕊_{𝔤/𝔱}`$. On the other hand, we see from the homogeneous Weyl formula (12) that the irreducible representations $`U_{c\lambda }`$ for $`cC`$ appear at least once in the decomposition of the tensor product $`V_\lambda 𝕊_{𝔤/𝔥}`$. We therefore conclude that the representations $`U_{c\lambda }`$ for $`cC`$ each occur exactly once in $`V_\lambda 𝕊_𝔭`$. ∎
###### Lemma 10.
If $`\mu `$ is a weight of $`V_\lambda 𝕊_𝔭`$ satisfying $`\mu +\rho _𝔥^2=\lambda +\rho _𝔤^2`$, then there exists a unique Weyl element $`wW_𝔤`$ such that $`\mu +\rho _𝔥=w(\lambda +\rho _𝔤)`$.
###### Proof.
If $`\mu `$ is a weight of $`V_\lambda 𝕊_𝔭`$, then $`\mu +\rho _𝔥`$ is a weight of the tensor product $`V_\lambda 𝕊_𝔭𝕊_{𝔥/𝔱}V_\lambda 𝕊_{𝔤/𝔱}`$. Since the Weyl group acts simply transitively on the Weyl chambers, there exists an element $`wW_𝔤`$ such that $`w^1(\mu +\rho _𝔥)`$ is dominant, where we recall that a weight $`\nu `$ is dominant if and only if $`\nu ,\alpha 0`$ for all positive roots $`\alpha `$. Note that every weight of $`𝕊_{𝔤/𝔱}`$ can be obtained from its highest weight $`\rho _𝔤`$ by subtracting a sum of positive roots. Likewise, for the tensor product $`V_\lambda 𝕊_{𝔤/𝔱}`$, the difference $`(\lambda +\rho _𝔤)w^1(\mu +\rho _𝔥)`$ is a sum of positive roots, and it follows that
$$\lambda +\rho _𝔤^2w^1(\mu +\rho _𝔥)^2,$$
with equality holding only when $`(\lambda +\rho _𝔤)w^1(\mu +\rho _𝔥)=0`$. As for the uniqueness of $`w`$, if $`\lambda `$ is dominant, then the weight $`\lambda +\rho _𝔤`$ lies in the interior of the positive Weyl chamber for $`𝔤`$, and thus the weights $`w(\lambda +\rho _𝔤)`$ for $`wW_𝔤`$ are distinct. ∎
Theorem 8 now follows immediately from the above two lemmas. We can actually be slightly more specific about the kernel of the Dirac operator, recovering the signs appearing in the homogeneous Weyl formula (12). Recall that the spin representation decomposes as $`𝕊_𝔭=𝕊_𝔭^+𝕊_𝔭^{}`$ into two half-spin representations. Since the Dirac operator is an odd element of the non-abelian Weil algebra $`U(𝔤)\mathrm{Cl}(𝔤)`$, it interchanges $`𝕊_𝔭^+`$ and $`𝕊_𝔭^{}`$. Restricting the domain of the Dirac operator to the positive half-spin representation, we obtain an operator
$$\overline{)}_{𝔤/𝔥}^+:V_\lambda 𝕊_𝔭^+V_\lambda 𝕊_𝔭^{}.$$
Furthermore, since the Dirac operator is formally self-adjoint, its adjoint is
$$\overline{)}_{𝔤/𝔥}^{}:V_\lambda 𝕊_𝔭^{}V_\lambda 𝕊_𝔭^+,$$
the restriction of the Dirac operator to the negative half-spin representation. Since these Dirac operator are acting on finite dimensional vector spaces, the index is the difference of the domain and range, so we have
(30)
$$\mathrm{Ker}\overline{)}_{𝔤/𝔥}^+\mathrm{Ker}\overline{)}_{𝔤/𝔥}^{}=V_\lambda 𝕊_𝔭^+V_\lambda 𝕊_𝔭^{},$$
which is given by the homogeneous Weyl formula (12). Comparing this with the kernel of $`\overline{)}_{𝔤/𝔥}=\overline{)}_{𝔤/𝔥}^+\overline{)}_{𝔤/𝔥}^{}`$ given in Theorem 8, we therefore obtain
$$\mathrm{Ker}\overline{)}_{𝔤/𝔥}^+=\underset{(1)^c=+1}{}U_{c\lambda },\mathrm{Ker}\overline{)}_{𝔤/𝔥}^{}=\underset{(1)^c=1}{}U_{c\lambda }.$$
In other words, there is no cancellation on the left hand side of equation (30), and the signed kernel of this Dirac operator picks out precisely those representations, with sign, appearing on the right hand side of the homogeneous Weyl formula (12).
## 7. The Clifford algebra $`\mathrm{Cl}(L𝔤)`$
In Section 4, we examined the Clifford algebra associated to a finite dimensional Lie algebra with an invariant inner product. The infinite dimensional case is more complicated, and the general theory of such infinite dimensional Clifford algebras and their spin representations is developed in the mathematical literature by Kostant and Sternberg in . Here, we consider the Clifford algebra associated to the Lie algebra $`L𝔤`$ of smooth maps from $`S^1`$ to a finite dimensional Lie algebra $`𝔤`$, where we restrict to the dense subspace of loops with finite Fourier expansions. This finiteness condition ensures that the Lie algebra $`L𝔤`$ has a countable basis, and the complexification of this loop space is then $`L𝔤_{}=_k𝔤_{}z^k=𝔤_{}[z,z^1]`$, the Lie algebra of finite Laurent series with values in $`𝔤_{}`$. Averaging the pointwise inner products over the loop, we obtain an invariant inner product on $`L𝔤`$ given by (3).
The Clifford algebra $`\mathrm{Cl}(L𝔤)`$ is spanned by *finite* sums of products of the form $`\xi _1\mathrm{}\xi _n`$ for loops $`\xi _iL𝔤`$, subject to the relation $`\{\xi ,\eta \}=2\xi ,\eta `$. However, the loop space analogues of the elements $`\stackrel{~}{\mathrm{a}}\mathrm{d}X`$ and $`\gamma `$ introduced in §4 are in fact infinite sums, so we must instead work with a formal completion of the Clifford algebra. Unfortunately, the product of two such infinite sums does not necessarily converge. On the other hand, given a spin representation $`𝒮_{L𝔤}`$ of the Clifford algebra $`\mathrm{Cl}(L𝔤)`$, we can view $`\mathrm{End}(𝒮_{L𝔤})`$ as a completion of $`\mathrm{Cl}(L𝔤)`$ with a well defined product given by the composition of endomorphisms. As we discussed in §2, to define the spin representation we must first choose a polarization. With respect to the action of the infinitesimal generator $`_\theta `$ of rotations, the complexified loop space $`L𝔤_{}`$ decomposes into its negative, zero, and positive energy subspaces $`L𝔤_{}=L𝔤_{}^{}𝔤_{}L𝔤_{}^+`$, where $`L𝔤_{}^+`$ and $`L𝔤_{}^{}`$ are isotropic subspaces which are dual to each other with respect to the inner product. The spin representation corresponding to this polarization is $`𝒮_{L𝔤}:=𝕊_𝔤\mathrm{\Lambda }^{}(L𝔤_{}^+)`$, and the Clifford action $`c:\mathrm{Cl}(L𝔤_{})\mathrm{End}(𝒮_{L𝔤})`$ is given by
$$c(\xi )=\{\begin{array}{cc}1\epsilon (\xi )\hfill & \text{for }\xi L𝔤_{}^+,\hfill \\ 1\iota (\xi )\hfill & \text{for }\xi L𝔤_{}^{},\hfill \\ c(\xi )(1)^F\hfill & \text{for }\xi 𝔤_{},\hfill \end{array}$$
where $`\epsilon `$ and $`\iota `$ are exterior multiplication and interior contraction respectively, and $`F`$ is the degree operator on the exterior algebra. If $`\{\eta _i\}`$ is a basis for $`\mathrm{Cl}(L𝔤_{}^{})`$, then when applied to a specific element of the spin representation $`𝒮_{L𝔤}`$, all but finitely many of the operators $`c(\eta _i)=\iota (\eta _i)`$ vanish. Formal infinite sums $`_ic(\omega _i)c(\eta _i)`$ with coefficients $`\omega _i\mathrm{Cl}(𝔤_{}L𝔤_{}^+)`$ therefore yield well defined operators on the spin representation, and in fact all elements of $`\mathrm{End}(𝒮_{L𝔤})`$ can be expressed in this form.
The exterior algebra that we shall consider here is not $`\mathrm{\Lambda }^{}(L𝔤^{})`$, but rather the algebra $`\mathrm{\Lambda }^{}(L𝔤)^{}`$ of skew-symmetric multilinear forms on $`L𝔤`$. Such forms can be expressed as formal infinite sums of basic products of the form $`\xi _1^{}\mathrm{}\xi _n^{}`$ for $`\xi _iL𝔤`$. In infinite dimensions, the Chevalley map $`\mathrm{ch}:\mathrm{Cl}(L𝔤)\mathrm{\Lambda }^{}(L𝔤)^{}`$ is no longer surjective; its image consists of all forms given by finite sums of the basic wedge products. Although the Chevalley map fails to converge if we attempt to extend it to the completion $`\mathrm{End}(𝒮_{L𝔤})`$ of $`\mathrm{Cl}(L𝔤)`$, we can perturb it by terms of lower degree to remove the infinite contributions. Separating the Clifford algebra into its positive and negative energy factors, we define the *normal ordering* map $`n:\mathrm{Cl}(L𝔤_{})\mathrm{\Lambda }^{}(L𝔤_{})^{}`$ by
$$n(\omega ^+\omega ^{})=\mathrm{ch}(\omega ^+)\mathrm{ch}(\omega ^{}),$$
where $`\omega ^+\mathrm{Cl}(𝔤_{}L𝔤_{}^+)`$ and $`\omega ^{}\mathrm{Cl}(L𝔤_{}^{})`$. The normal ordering map extends to the completion $`\mathrm{End}(𝒮_{L𝔤})`$ of the Clifford algebra, and its image is the subspace $`\mathrm{\Lambda }^{}(L𝔤_{})^+\mathrm{\Lambda }^{}(L𝔤_{})^{}`$ given by
$$\mathrm{\Lambda }^{}(L𝔤_{})^+=\left\{\omega \mathrm{\Lambda }^{}(L𝔤_{})^{}\right|(\iota _\eta \omega )^+\mathrm{\Lambda }^{}(L𝔤_{}^{})\text{ for all }\eta \mathrm{\Lambda }^{}(L𝔤_{}^{})\},$$
where $`()^+`$ denotes the projection of $`\mathrm{\Lambda }^{}(L𝔤_{})^{}`$ onto $`\mathrm{\Lambda }^{}(L𝔤_{}^+)^{}`$, and we identify $`\mathrm{\Lambda }^{}(L𝔤_{}^{})`$ with a subspace of $`\mathrm{\Lambda }^{}(L𝔤_{}^+)^{}`$ via the inner product. In terms of a basis $`\{\eta _i\}`$ for $`\mathrm{\Lambda }^{}(L𝔤^{})`$, we may write elements of $`\mathrm{\Lambda }^{}(L𝔤_{})^+`$ as formal infinite sums $`_i\omega _i^{}\eta _i^{}`$, with $`\omega _i\mathrm{\Lambda }^{}(𝔤_{}L𝔤_{}^+)`$ living in the zero and positive energy components.
###### Remark.
Decomposing $`L𝔤_{}=_k𝔤_{}z^k`$ in terms of its energy grading, we define a secondary degree on $`L𝔤_{}`$ counting only the negative energy contribution
$$\mathrm{sdeg}Xz^k=\{\begin{array}{cc}0\hfill & \text{for }k0,\hfill \\ k\hfill & \text{for }k<0,\hfill \end{array}$$
where $`X𝔤_{}`$ and $`Xz^k`$ is the loop $`zXz^k`$ for $`|z|=1`$. Let $`L𝔤_{}^{}=𝔤_{}^{}[z,z^1]`$ denote the reduced dual of $`L𝔤_{}`$. Extending $`\mathrm{sdeg}`$ to the exterior algebra $`\mathrm{\Lambda }^{}(L𝔤_{}^{})`$, we note that $`\mathrm{\Lambda }^{}(L𝔤_{})^+`$ is the completion of $`\mathrm{\Lambda }^{}(L𝔤_{}^{})`$ with respect to $`\mathrm{sdeg}`$. In other words, $`\mathrm{\Lambda }^{}(L𝔤_{})^+`$ consists of all formal infinite sums $`_i\omega _i`$ of $`\mathrm{sdeg}`$-homogeneous elements $`\omega _i\mathrm{\Lambda }^{}(L𝔤_{}^{})`$ with $`\mathrm{sdeg}\omega _i\mathrm{}`$.
We can now use the normal ordering identification $`n:\mathrm{End}(𝒮_{L𝔤})\mathrm{\Lambda }^{}(L𝔤_{})^+`$ to define product and bracket structures on $`\mathrm{\Lambda }^{}(L𝔤_{})^+`$. The normal ordered product $`\omega _1_n\omega _2=n(n^1\omega _1n^1\omega _2)`$ on the exterior algebra differs from the product induced by the Chevalley identification by terms of lower degree. However, many of the supercommutators remain unchanged. In particular, the normal ordered bracket with the dual $`\xi ^{}L𝔤_{}^{}\mathrm{\Lambda }^1(L𝔤_{})^+`$ of a loop $`\xi L𝔤_{}`$ is still given by
$$\begin{array}{cc}\hfill [\xi ^{},\omega ^+\omega ^{}]_n& =n[n^1(\xi ^{}),n^1(\omega ^+\omega ^{})]=n[\xi ,\mathrm{ch}^1\omega ^+\mathrm{ch}^1\omega ^{}]\hfill \\ & =n\left([\xi ,\mathrm{ch}^1\omega ^+]\mathrm{ch}^1\omega ^{}\pm \mathrm{ch}^1\omega ^+[\xi ,\mathrm{ch}^1\omega ^{}]\right)\hfill \\ & =2\iota _\xi \omega ^+\omega ^{}\pm \omega ^+2\iota _\xi \omega ^{}=2\iota _\xi (\omega ^+\omega ^{}),\hfill \end{array}$$
for $`\omega ^+\mathrm{\Lambda }^{}(𝔤_{}L𝔤_{}^+)^{}`$ and $`\omega ^{}\mathrm{\Lambda }^{}(L𝔤_{}^{})^{}`$ of homogeneous degree. Thus,
(31)
$$[\xi ^{},\omega ]_n=2\iota _\xi \omega \text{ for }\xi L𝔤_{}\text{ and }\omega \mathrm{\Lambda }^{}(L𝔤_{})^+\text{.}$$
Reprising the discussion of §4, for any $`\xi L𝔤`$, consider the 2-form $`d\xi ^{}`$ given by $`d\xi ^{}(\eta ,\zeta )=\frac{1}{2}\xi ,[\eta ,\zeta ]`$ for all $`\eta ,\zeta L𝔤`$. Although $`_\theta `$ is not an element of $`L𝔤`$, we can nevertheless define an analogous 2-form $`d_\theta ^{}`$ by $`d_\theta ^{}(\xi ,\eta ):=\frac{1}{2}\xi ,_\theta \eta `$. Note that $`d_\theta ^{}`$ is closed but not exact, so it defines a cohomology element in $`H^2(L𝔤)`$. Finally, the fundamental 3-form $`\mathrm{\Omega }`$ is given by $`\mathrm{\Omega }(\xi ,\eta ,\zeta )=\frac{1}{6}\xi ,[\eta ,\zeta ]`$ for $`\xi ,\eta ,\zeta L𝔤`$. These elements all lie in $`\mathrm{\Lambda }^{}(L𝔤_{})^+`$, and they satisfy the identities
(32)
$$\iota _\xi d\eta ^{}=[\xi ,\eta ]^{},\iota _\xi d_\theta ^{}=(_\theta \xi )^{},\iota _\xi \mathrm{\Omega }=d\xi ^{}.$$
Using the normal ordered product and bracket on $`\mathrm{\Lambda }^{}(L𝔤_{})^+`$ coming from $`\mathrm{End}(𝒮_{L𝔤})`$, we obtain the loop space version of Corollary 6.
###### Theorem 11.
If $`𝔤`$ is simple, then the elements $`1`$, $`\xi ^{}`$ for $`\xi L𝔤`$, $`\stackrel{~}{\mathrm{a}}\mathrm{d}\xi =\frac{1}{2}d\xi ^{}`$ for $`\xi \stackrel{~}{}L𝔤`$, and $`\gamma =\frac{1}{4}\mathrm{\Omega }`$ span a Lie superalgebra in $`\mathrm{\Lambda }^{}(L𝔤)^+\mathrm{End}(𝒮_{L𝔤})`$ satisfying
$`\{\xi ^{},\eta ^{}\}`$ $`=2\xi ,\eta ,`$
$`[\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ,\eta ^{}]`$ $`=[\xi ,\eta ]^{},`$ $`[\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ,\stackrel{~}{\mathrm{a}}\mathrm{d}\eta ]`$ $`=\stackrel{~}{\mathrm{a}}\mathrm{d}[\xi ,\eta ]+ic_𝔤\xi ,_\theta \eta ,`$
$`[\stackrel{~}{\mathrm{a}}\mathrm{d}_\theta ,\xi ^{}]`$ $`=(_\theta \xi )^{},`$ $`[\stackrel{~}{\mathrm{a}}\mathrm{d}_\theta ,\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ]`$ $`=\stackrel{~}{\mathrm{a}}\mathrm{d}(_\theta \xi ),`$
$`\{\gamma ,\xi ^{}\}`$ $`=\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ,`$ $`[\gamma ,\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ]`$ $`=\frac{1}{2}ic_𝔤(_\theta \xi )^{},`$
$`[\gamma ,\stackrel{~}{\mathrm{a}}\mathrm{d}_\theta ]`$ $`=0,`$ $`\{\gamma ,\gamma \}`$ $`=ic_𝔤\stackrel{~}{\mathrm{a}}\mathrm{d}_\theta \frac{1}{24}c_𝔤dim𝔤,`$
where $`c_𝔤`$ is the value of the Casimir operator of $`𝔤`$ in the adjoint representation.
###### Proof.
The bracket $`\{\xi ^{},\eta ^{}\}=2\xi ,\eta `$ is simply the definition of the Clifford algebra, while the brackets $`[\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ,\eta ^{}]=[\xi ,\eta ]^{}`$ and $`[\stackrel{~}{\mathrm{a}}\mathrm{d}_\theta ,\xi ^{}]=(_\theta \xi )^{}`$ and $`\{\gamma ,\xi ^{}\}=\stackrel{~}{\mathrm{a}}\mathrm{d}\xi `$ follow immediately from (31) and (32). By the Jacobi identity, for any $`\xi ,\eta \stackrel{~}{}L𝔤`$ and $`\zeta L𝔤`$ we have
$$\begin{array}{cc}\hfill [[\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ,\stackrel{~}{\mathrm{a}}\mathrm{d}\eta ],\zeta ^{}]& =[\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ,[\stackrel{~}{\mathrm{a}}\mathrm{d}\eta ,\zeta ^{}]][\stackrel{~}{\mathrm{a}}\mathrm{d}\eta ,[\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ,\zeta ^{}]]\hfill \\ & =[\xi ,[\eta ,\zeta ]]^{}[\eta ,[\xi ,\zeta ]]^{}=[[\xi ,\eta ],\zeta ]^{}=[\stackrel{~}{\mathrm{a}}\mathrm{d}[\xi ,\eta ],\zeta ^{}],\hfill \end{array}$$
which shows that $`\stackrel{~}{\mathrm{a}}\mathrm{d}`$ is a projective representation of $`\stackrel{~}{}L𝔤`$ on $`𝒮_{L𝔤}`$. In Theorem 1, we established that this spin representation has central charge $`c_𝔤`$, which gives us the brackets $`[\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ,\stackrel{~}{\mathrm{a}}\mathrm{d}\eta ]=\stackrel{~}{\mathrm{a}}\mathrm{d}[\xi ,\eta ]+ic_𝔤\xi ,_\theta \eta `$ and $`[\stackrel{~}{\mathrm{a}}\mathrm{d}_\theta ,\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ]=\stackrel{~}{\mathrm{a}}\mathrm{d}_\theta \xi `$.
To compute $`\gamma ^2`$, we write it as the sum $`\gamma ^2=(\gamma ^2)_0+(\gamma ^2)_2`$ of homogeneous forms of degrees 0 and 2. (We shall see that $`\gamma ^2`$ has no components of degrees 4 or 6.) Since $`\iota _\xi =\mathrm{ad}\xi ^{}`$ for $`\xi L𝔤`$ is a derivation with respect to the backet, we have
$$\iota _\xi \gamma ^2=[\iota _\xi \gamma ,\gamma ]=\frac{1}{2}[\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ,\gamma ].$$
Taking one further interior contraction, we obtain
$$\iota _\xi \iota _\eta \gamma ^2=\frac{1}{4}\left([\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ,\stackrel{~}{\mathrm{a}}\mathrm{d}\eta ]\stackrel{~}{\mathrm{a}}\mathrm{d}[\xi ,\eta ]\right)=\frac{1}{4}ic_𝔤\xi ,_\theta \eta ,$$
which is a constant. It follows that $`\gamma ^2`$ has no components of degree higher than 2, and that $`(\gamma ^2)_2`$ is the 2-cocycle determining the central extension of $`L𝔤`$ for the spin representation $`\stackrel{~}{\mathrm{a}}\mathrm{d}`$. In fact, this 2-cocycle is a multiple of $`\stackrel{~}{\mathrm{a}}\mathrm{d}_\theta `$, and we have
$$(\gamma ^2)_2(\xi ,\eta )=\frac{1}{2}\iota _\xi \iota _\eta \gamma ^2=\frac{1}{8}\left(ic_𝔤\xi ,_\theta \eta \right)=\frac{1}{2}ic_𝔤(\stackrel{~}{\mathrm{a}}\mathrm{d}_\theta )(\xi ,\eta ).$$
Going back up one level, we see that
$$[\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ,\gamma ]=2\iota _\xi \gamma ^2=ic_𝔤\iota _\xi \stackrel{~}{\mathrm{a}}\mathrm{d}_\theta =\frac{1}{2}ic_𝔤(_\theta \xi )^{}.$$
Finally, the value of the constant $`(\gamma ^2)_0`$ is the value of $`\gamma ^2`$ acting on the minimum energy subspace $`𝒮_{L𝔤}(0)`$ of the spin representation, since $`\stackrel{~}{\mathrm{a}}\mathrm{d}_\theta `$ vanishes there. However, all the terms in $`\gamma ^2`$ vanish on $`𝒮_{L𝔤}(0)`$ except the contribution from the constant loops, and thus $`(\gamma ^2)_0=\frac{1}{48}\mathrm{tr}_𝔤\mathrm{\Delta }_{\mathrm{ad}}^𝔤=\frac{1}{48}c_𝔤dim𝔤`$ as we proved in Corollary 6. ∎
Taking a slightly different view of this theorem, the commutation relations given in Theorem 11 determine a Lie superalgebra (with subscripts denoting the grading)
$$_{\text{even}}L𝔤_{\text{odd}}(\stackrel{~}{}L𝔤)_{\text{even}}_{\text{odd}},$$
and the identification of $`\mathrm{\Lambda }^{}(L𝔤)^+`$ with its image in $`\mathrm{End}(𝒮_{L𝔤})`$ gives a representation of this Lie superalgebra on the spin representation $`𝒮_{L𝔤}`$. Actually, we can extend this Lie superalgebra further. The component $`_{\text{even}}L𝔤_{\text{odd}}L𝔤_{\text{even}}`$ is called a *super Kac-Moody algebra*, and using superspace notation, its complexification is a central extension of the polynomial algebra $`𝔤[z,z^1,\mathrm{\Theta }]`$, where $`\mathrm{\Theta }`$ is an odd variable (i.e. $`\mathrm{\Theta }^2=0`$). The *super Virasoro algebra* $`\mathrm{SVir}`$ is the universal central extension of the Lie algebra of derivations of $`[z,z^1,\mathrm{\Theta }]`$. (Note that the even derivations are just the vector fields on the circle.) The super Virasoro algebra therefore acts on the super Kac-Moody algebra, and their semidirect sum is referred to as the *$`N=1`$ superconformal current algebra* (see ):
$$\mathrm{SVir}\stackrel{~}{}(_{\text{even}}L𝔤_{\text{odd}}L𝔤_{\text{even}}).$$
In our case, the elements $`\stackrel{~}{\mathrm{a}}\mathrm{d}_\theta `$ and $`\gamma `$ span the even and odd zero-mode subspaces of the super Virasoro algebra, with commutator $`\{\gamma ,\gamma \}=ic_𝔤\left(\stackrel{~}{\mathrm{a}}\mathrm{d}_\theta +\frac{1}{24}idim𝔤\right)`$. Here, the additional $`\frac{1}{24}dim𝔤`$ term, which is sometimes incorporated into the definition of $`\stackrel{~}{\mathrm{a}}\mathrm{d}_\theta `$, corresponds to the anomalous energy shift we encountered in (9).
Given an orthonormal basis $`\{X_i\}`$ for $`𝔤`$, the loops $`X_i^n=X_iz^n`$ for $`n`$ form a basis for $`L𝔤_{}`$ satisfying $`X_i^n,X_j^m=\delta _{i,j}\delta _{n,m}`$. In terms of this basis, we have
$`\stackrel{~}{\mathrm{a}}\mathrm{d}\xi `$ $`={\displaystyle \frac{1}{4}}{\displaystyle \underset{i,k}{}}X_i^k[\xi ,X_i^k],\stackrel{~}{\mathrm{a}}\mathrm{d}_\theta ={\displaystyle \frac{1}{2}}{\displaystyle \underset{j,k>0}{}}ikX_j^kX_j^k,`$
$`\gamma `$ $`={\displaystyle \frac{1}{24}}{\displaystyle \underset{i,j,k,l}{}}X_i^kX_j^l[X_i,X_j]^{k+l}={\displaystyle \frac{1}{6}}{\displaystyle \underset{i,k}{}}X_i^k\stackrel{~}{\mathrm{a}}\mathrm{d}X_i^k.`$
Note that in the expressions for $`\stackrel{~}{\mathrm{a}}\mathrm{d}\xi `$ and $`\gamma `$, the ordering of the factors does not matter (up to sign), since they are orthogonal and therefore anti-commute with each other. However, in the expression for $`\stackrel{~}{\mathrm{a}}\mathrm{d}_\theta `$, we have $`\{X_i^k,X_i^k\}=2`$, so changing the order of the factors shifts the operator by a constant. Here we see normal ordering in action, forcing us to write factors $`X_i^k`$ with $`k`$ positive on the left and factors $`X_i^k`$ with $`k`$ negative on the right. In physics notation, this would be written as $`\stackrel{~}{\mathrm{a}}\mathrm{d}_\theta =\frac{1}{4}_{j,k}ik:X_j^kX_j^k:,`$ where $`:\xi \eta :=n^1(n\xi n\eta )`$ denotes the normal ordered product in the Clifford algebra. (This colon notation is misleading as it is not a map on the Clifford algebra but rather an instruction to replace all Clifford products between the colons with normal ordered products.)
If the Lie algebra $`𝔤`$ is not simple, then Theorem 11 still holds, albeit with slightly modified commutation relations. For a general finite dimensional Lie algebra $`𝔤`$, the Casimir operator $`\mathrm{\Delta }_{\mathrm{ad}}^𝔤=\frac{1}{2}_i(\mathrm{ad}X_i)^2`$ no longer takes a constant value $`c_𝔤`$. In this case, the role of the quadratic element $`\stackrel{~}{\mathrm{a}}\mathrm{d}_\theta `$ is played by the 2-cocycle $`\omega _{\stackrel{~}{\mathrm{a}}\mathrm{d}}`$ for the projective spin representation $`\stackrel{~}{\mathrm{a}}\mathrm{d}`$, given on $`\xi ,\eta L𝔤`$ by
$$\omega _{\stackrel{~}{\mathrm{a}}\mathrm{d}}(\xi ,\eta ):=[\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ,\stackrel{~}{\mathrm{a}}\mathrm{d}\eta ]\stackrel{~}{\mathrm{a}}\mathrm{d}[\xi ,\eta ]=i\xi ,\mathrm{\Delta }_{\mathrm{ad}}^𝔤_\theta \eta ,$$
where the Casimir operator $`\mathrm{\Delta }_{\mathrm{ad}}^𝔤`$ acts pointwise on the loop space $`L𝔤`$. Viewing $`\omega _{\stackrel{~}{\mathrm{a}}\mathrm{d}}`$ as an element of the Clifford algebra, we have the commutator
$$[\omega _{\stackrel{~}{\mathrm{a}}\mathrm{d}},\xi ^{}]=2\iota _\xi \omega _{\stackrel{~}{\mathrm{a}}\mathrm{d}}=4i\left(\mathrm{\Delta }_{\mathrm{ad}}^𝔤_\theta \xi \right)^{},$$
so we may also view the projective cocycle as $`\omega _{\stackrel{~}{\mathrm{a}}\mathrm{d}}=4i\stackrel{~}{\mathrm{a}}\mathrm{d}(\mathrm{\Delta }^𝔤_\theta )`$, where $`\mathrm{\Delta }^𝔤`$ is the formal Casimir operator in the universal enveloping algebra of $`𝔤`$. We therefore have
$$[\omega _{\stackrel{~}{\mathrm{a}}\mathrm{d}},\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ]=4i[\stackrel{~}{\mathrm{a}}\mathrm{d}(\mathrm{\Delta }^𝔤_\theta ),\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ]=4i\stackrel{~}{\mathrm{a}}\mathrm{d}\left(\mathrm{\Delta }_{\mathrm{ad}}^𝔤_\theta \xi \right),$$
and the adjoint action of $`\gamma `$ in Theorem 11 then becomes
$`[\gamma ,\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ]`$ $`=\frac{1}{2}i\left(\mathrm{\Delta }_{\mathrm{ad}}^𝔤_\theta \xi \right)^{},`$
$`\{\gamma ,\gamma \}`$ $`=\frac{1}{4}\omega _{\stackrel{~}{\mathrm{a}}\mathrm{d}}\frac{1}{24}\mathrm{tr}_𝔤\mathrm{\Delta }_{\mathrm{ad}}^𝔤,`$
with the other commutation relations remaining unchanged. Alternatively, the projective cocycle $`\omega _{\stackrel{~}{\mathrm{a}}\mathrm{d}}`$ can be viewed as the 2-form component of the Casimir operator
$$\mathrm{\Delta }_{\stackrel{~}{\mathrm{a}}\mathrm{d}}^{L𝔤}=2i\stackrel{~}{\mathrm{a}}\mathrm{d}(\mathrm{\Delta }^𝔤_\theta )+\mathrm{\Delta }_{\stackrel{~}{\mathrm{a}}\mathrm{d}}^𝔤=\frac{1}{2}\omega _{\stackrel{~}{\mathrm{a}}\mathrm{d}}+\frac{1}{8}\mathrm{tr}_𝔤\mathrm{\Delta }_{\mathrm{ad}}^𝔤$$
for the spin representation $`\stackrel{~}{\mathrm{a}}\mathrm{d}`$ of $`L𝔤`$, which we discuss in Theorem 12 below.
## 8. The Dirac operator on $`L𝔤`$
Following our discussion in Section 5, given any positive energy representation $`r:\stackrel{~}{L}𝔤\mathrm{End}()`$, we construct a Dirac operator
$$\overline{)}_r:=\widehat{r}+1\frac{1}{2}\mathrm{\Omega }_{L𝔤}\mathrm{End}(𝒮_{L𝔤}),$$
where $`\widehat{r}`$ is the tautological $`\mathrm{End}()`$-valued 1-form on $`L𝔤`$ given by $`\widehat{r}(\xi )=r(\xi )`$ for all $`\xi L𝔤`$, and $`\mathrm{\Omega }_{L𝔤}`$ is the fundamental 3-form given by $`\mathrm{\Omega }_{L𝔤}(\xi ,\eta ,\zeta )=\frac{1}{6}\xi ,[\eta ,\zeta ]`$ for $`\xi ,\eta ,\zeta L𝔤`$. As in the previous section, we implicitly identify $`\mathrm{\Lambda }^{}(L𝔤)^+`$ with its image in $`\mathrm{End}(𝒮_{L𝔤})`$. Written in terms of a basis $`\{X_i^n\}`$ of $`L𝔤`$ satisfying $`X_i^n,X_j^m=\delta _{i,j}\delta _{n,m}`$, this Dirac operator is
$$\begin{array}{cc}\hfill \overline{)}_r& =\underset{i,n}{}X_i^nr(X_i^n)\frac{1}{12}\underset{i,j,m,n}{}X_i^nX_j^m[X_i,X_j]^{n+m}\hfill \\ & =\underset{i,n}{}X_i^n\left(r(X_i^n)+\frac{1}{3}\stackrel{~}{\mathrm{a}}\mathrm{d}X_i^n\right).\hfill \end{array}$$
Note that all of the individual factors in this expression (anti-)commute with each other, so $`\overline{)}_r`$ does indeed give a well-defined operator on the tensor product $`𝒮_{L𝔤}`$, without requiring normal ordering or dealing with any infinite constants.
In its most general form, if we take the representation $`r`$ to be the canonical inclusion $`r:L𝔤U(L𝔤)`$ of $`L𝔤`$ into its universal enveloping algebra $`U(L𝔤)`$, then the corresponding universal Dirac operator is an element of the formal completion of the non-abelian Weil algebra $`𝒜=U(\stackrel{~}{L}𝔤)\mathrm{Cl}(L𝔤)`$. (Alternatively, we may view $`𝒜`$ as the universal enveloping algebra of the super Kac-Moody algebra $`\stackrel{~}{L}𝔤_{\text{even}}L𝔤_{\text{odd}}`$.) As we saw in the previous section, the product of two such infinite formal sums does not necessarily converge. However, keeping in mind that we are really working with operators on Hilbert spaces, we can indeed extend multiplication to a suitable subspace $`𝒜^+`$ of the formal completion, which we define as the largest subspace for which the homomorphism $`𝒜\mathrm{End}(𝒮_{L𝔤})`$ extends to $`𝒜^+`$ for any positive energy representation $``$ of $`\stackrel{~}{L}𝔤`$. In particular, if $``$ is a faithful representation of $`U(\stackrel{~}{L}𝔤)`$—we can construct such a representation by taking the Hilbert space direct sum of countably many irreducible positive energy representations of $`\stackrel{~}{L}𝔤`$—then the homomorphism $`𝒜^+\mathrm{End}(𝒮_{L𝔤})`$ induces a product structure on $`𝒜^+`$. Fortunately, we can perform all of our computations here using the techniques of the previous section, working with $`U(\stackrel{~}{L}𝔤)`$-valued forms on $`\stackrel{~}{L}𝔤`$.
Using this extended multiplication, the square of the Dirac operator is
$$\overline{)}^2=\widehat{r}^2+\{\widehat{r},\frac{1}{2}\mathrm{\Omega }_{L𝔤}\}+\frac{1}{4}\mathrm{\Omega }_{L𝔤}^2.$$
Since $`\widehat{r}`$ is an $`\mathrm{End}()`$-valued 1-form on $`L𝔤`$, its square is a sum $`\widehat{r}^2=(\widehat{r}^2)_0+(\widehat{r}^2)_2`$ of forms of homogeneous degrees 0 and 2. For the degree 2 component, we have $`(\widehat{r}^2)_2=\widehat{r}\widehat{r}`$, and the “curvature” $`d\widehat{r}+\widehat{r}\widehat{r}`$ of the representation $`r`$ is given by
$$(d\widehat{r}+\widehat{r}\widehat{r})(\xi ,\eta )=\frac{1}{2}\left([r(\xi ),r(\eta )]r([\xi ,\eta ])\right)=\frac{1}{2}\omega _r(\xi ,\eta ),$$
where $`\omega _r\mathrm{\Lambda }^2(L𝔤)^+`$ is the 2-cocycle corresponding to the projective representation $`r`$. If $`𝔤`$ is simple and $`I`$ is the generator of the universal central extension of $`L𝔤`$, then $`\omega _r=4r(I)\stackrel{~}{\mathrm{a}}\mathrm{d}_\theta `$. The degree 0 component of $`\widehat{r}^2`$ is given by the following:
###### Theorem 12.
The operator $`\mathrm{\Delta }_r^{L𝔤}:=\frac{1}{2}\left(\widehat{r}^2\right)_0`$ is called the Casimir operator for the loop group $`L𝔤`$, and if $`𝔤`$ is simple then the Casimir operator acting on the irreducible positive energy representation $`_𝛌`$ with lowest weight $`𝛌=(m,\lambda ,h)`$ is given by
(33)
$$\begin{array}{cc}\hfill \mathrm{\Delta }_r^{L𝔤}& =i(h+c_𝔤)(r(_\theta )im)+\mathrm{\Delta }_\lambda ^𝔤\hfill \\ & =i(h+c_𝔤)r(_\theta )+\frac{1}{2}\left(𝝀𝝆_𝔤^2𝝆_𝔤^2\right),\hfill \end{array}$$
where $`𝛒_𝔤=(0,\rho _𝔤,c_𝔤)`$ and $`c_𝔤=\mathrm{\Delta }_{\mathrm{ad}}^𝔤`$ is the value of the quadratic Casimir operator of $`𝔤`$ acting on the adjoint representation, and the inner product is given by (6).
###### Proof.
In order to simply our calculations, we first note the following identites:
$$\begin{array}{cc}\hfill [\widehat{r},\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ](\eta )& =[r(\eta ),\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ]\{\widehat{r},\frac{1}{2}[\eta ,\xi ]^{}\}=r([\xi ,\eta ]),\hfill \\ \hfill [\widehat{r},r(\xi )](\eta )& =[r(\eta ),r(\xi )]=r([\eta ,\xi ])+\eta ,_\theta \xi r(I)\hfill \\ & =\left([\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ,\widehat{r}]+r(I)(_\theta \xi )^{}\right)(\eta ),\text{ and}\hfill \\ \hfill [\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ,\widehat{r}^2]_0& =[\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ,(\widehat{r}^2)_2]_0=[\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ,d\widehat{r}]_0\hfill \\ & =2(\stackrel{~}{\mathrm{a}}\mathrm{d}I)\widehat{r}(_\theta \xi )=2(\stackrel{~}{\mathrm{a}}\mathrm{d}I)r(_\theta \xi ).\hfill \end{array}$$
We now show that the commutator of $`\mathrm{\Delta }_r^{L𝔤}`$ with an element $`\xi L𝔤`$ is
$$\begin{array}{cc}\hfill [\mathrm{\Delta }_r^{L𝔤},r(\xi )]& =\frac{1}{2}[\widehat{r}^2,r(\xi )]_0=\frac{1}{2}\{\widehat{r},[\widehat{r},r(\xi )]\}_0\hfill \\ & =\frac{1}{2}\{\widehat{r},[\widehat{r},\stackrel{~}{\mathrm{a}}\mathrm{d}\xi ]\}_0\frac{1}{2}\{\widehat{r},r(I)(_\theta \xi )^{}\}_0\hfill \\ & =(\stackrel{~}{\mathrm{a}}\mathrm{d}I)r(_\theta \xi )r(I)r(_\theta \xi )\hfill \\ & =[(r(I)+\stackrel{~}{\mathrm{a}}\mathrm{d}I)r(_\theta ),r(\xi )].\hfill \end{array}$$
It follows that the operator $`\stackrel{~}{\mathrm{\Delta }}_r^{L𝔤}:=\mathrm{\Delta }_r^{L𝔤}+(r(I)+\stackrel{~}{\mathrm{a}}\mathrm{d}I)r(_\theta )`$ commutes with the action of $`L𝔤`$, and therefore takes a constant value on each irreducible representation. Acting on the minimum energy subspace $`_𝝀(m)`$ of $`_𝝀`$, the only terms contributing to $`\mathrm{\Delta }_r^{L𝔤}`$ are those coming from the constant loops, and thus this constant is
$$\stackrel{~}{\mathrm{\Delta }}_𝝀^{L𝔤}=\mathrm{\Delta }_r^{L𝔤}|_{_𝝀(m)}+i(h+c_g)r(_\theta )|_{_𝝀(m)}=\mathrm{\Delta }_\lambda ^𝔤(h+c_𝔤)m.$$
The desired result then follows immediately. ∎
By definition, the 0-form component of $`\widehat{r}^2`$ acts as the identity operator on $`𝒮_{L𝔤}`$. To compute the action of $`\mathrm{\Delta }_r^{L𝔤}`$, we can therefore restrict it to the minimum energy subspace $`𝒮_{L𝔤}(0)`$ of the spin representation. In terms of a basis $`\{X_i^n\}`$, we have
$$\begin{array}{cc}\hfill \mathrm{\Delta }_r^{L𝔤}& =\frac{1}{2}\underset{i,n}{}r(X_i^n)X_i^n\underset{j,m}{}r(X_j^m)X_j^m|_{𝒮_{L𝔤}(0)}\hfill \\ & =\frac{1}{2}\underset{i,j}{}\left(r(X_i)r(X_j)X_iX_j+\underset{n>0}{}r(X_i^n)r(X_j^n)X_i^nX_j^n\right)\hfill \\ & =\mathrm{\Delta }_r^𝔤\underset{i,n>0}{}r(X_i^n)r(X_i^n),\hfill \end{array}$$
which is the usual definition of the Casimir operator for a loop group. The Casimir operator can be used to define the energy operator $`r(_\theta )`$ in terms of the action of $`L𝔤`$. The constant term $`\mathrm{\Delta }_\lambda ^𝔤`$ is sometimes incorporated into $`r(_\theta )`$, in which case it is viewed as an anomalous energy shift due to the degeneracy of the vacuum.
Returning to our computation of the square of the Dirac operator, we note that the cross term is given by the anti-commutator $`\{\widehat{r},\frac{1}{2}\mathrm{\Omega }_{L𝔤}\}=d\widehat{r}`$, and we obtain
$$\begin{array}{c}\hfill \overline{)}^2=2\mathrm{\Delta }_r^{L𝔤}+\frac{1}{2}\omega _\varrho \frac{1}{12}\mathrm{tr}_𝔤\mathrm{\Delta }_{\mathrm{ad}}^𝔤,\end{array}$$
where $`\omega _\varrho `$ is the 2-cocycle corresponding to the diagonal action $`\varrho =r1+1\stackrel{~}{\mathrm{a}}\mathrm{d}`$ on the tensor product $`𝒮_{L𝔤}`$. If $`𝔤`$ is simple, then this 2-cocycle is $`\omega _\varrho =4\varrho (I)\stackrel{~}{\mathrm{a}}\mathrm{d}_\theta `$. Furthermore, if $`_𝝀`$ is the irreducible positive energy representation of $`L𝔤`$ with lowest weight $`𝝀=(m,\lambda ,h)`$, then using (33) for the Casimir operator, we have
(34)
$$\begin{array}{cc}\hfill \overline{)}_𝝀^2& =2i(h+c_𝔤)\left(\varrho (_\theta )im\right)\lambda +\rho _g^2\hfill \\ & =2\varrho (I)\varrho (_\theta )𝝀𝝆_𝔤^2.\hfill \end{array}$$
Note that unlike the finite dimensional case discussed in §5, the square of the Dirac operator for $`L𝔤`$ does not take a constant value on each irreducible representation. Here, the Dirac operator fails to commute with the diagonal action $`\varrho `$ of $`L𝔤`$ on $`𝒮_{L𝔤}`$. Since the Dirac operator satisfies the identity $`\iota _\xi \overline{)}=\varrho (\xi )`$, we have
$$[\varrho (\xi ),\overline{)}]=\iota _\xi \overline{)}^2=\frac{1}{2}\iota _\xi \omega _\varrho =2\varrho (I)\iota _\xi \stackrel{~}{\mathrm{a}}\mathrm{d}_\theta =\varrho (I)(_\theta \xi )^{},$$
and thus $`\overline{)}`$ commutes only with the subalgebra $`𝔤`$ of $`\stackrel{~}{}\stackrel{~}{L}𝔤`$.
If the Lie algebra $`𝔤`$ is reductive but not simple, then the expression (34) for the square of the Dirac operator still holds provided that the central extension satisfies $`\omega _\varrho =4\varrho (I)\stackrel{~}{\mathrm{a}}\mathrm{d}_\theta `$. In other words, the invariant inner product on $`𝔤`$ must satisfy
$$[\varrho (\xi ),\varrho (\eta )]\varrho ([\xi ,\eta ])=\varrho (I)\xi ,_\theta \eta $$
for some imaginary constant $`\varrho (I)`$. Given any irreducible positive energy projective representation $``$ of $`L𝔤`$, we can always choose an invariant inner product on $`𝔤`$ such that $`𝒮_{L𝔤}`$ is a true representation of the corresponding central extension $`\stackrel{~}{L}𝔤`$. However, this choice of inner product depends on the representation, so this approach does not give a universal expression for the Dirac operator.
## 9. The Dirac operator on $`L𝔤/L𝔥`$
As in Section 6, let $`𝔥`$ be a Lie subalgebra of $`𝔤`$, and let $`𝔭`$ denote the orthogonal complement of $`𝔥`$ with respect to the invariant inner product on $`𝔤`$. This orthogonal decomposition extends to the loop Lie algebra $`L𝔤=L𝔥L𝔭`$, and the Clifford algebra decomposes as $`\mathrm{Cl}(L𝔤)\mathrm{Cl}(L𝔥)\mathrm{Cl}(L𝔭)`$. If $`𝔭`$ is even dimensional, as is the case when $`𝔥`$ has the same rank as $`𝔤`$, then we can also factor the spin representation as $`𝒮_{L𝔤}𝒮_{L𝔥}𝒮_{L𝔭}`$, where $`𝒮_{L𝔥}`$ and $`𝒮_{L𝔭}`$ are representations of $`\stackrel{~}{L}𝔥`$ of levels $`c_𝔥`$ and $`c_𝔤c_𝔥`$ respectively, and the action of $`L𝔥`$ on $`𝒮_{L𝔭}`$ is
$`\stackrel{~}{\mathrm{a}}\mathrm{d}_{L𝔭}:L𝔥`$ $`\mathrm{\Lambda }^2(L𝔭)^+\mathrm{End}(𝒮_{L𝔭})`$
$`\zeta `$ $`(\stackrel{~}{\mathrm{a}}\mathrm{d}_{L𝔭}\zeta )(\xi ,\eta )=\frac{1}{4}\xi ,[\zeta ,\eta ]`$
for $`\zeta L𝔥`$ and $`\xi ,\eta L𝔭`$.
Given any positive energy representation $`r_{L𝔤}`$ of $`\stackrel{~}{L}𝔤`$ on a Hilbert space $``$, its restriction gives a representation $`r_{L𝔥}`$ of $`\stackrel{~}{L}𝔥`$ on $``$. Now consider the diagonal representation $`r_{L𝔥}^{}=r_{L𝔥}1+1\stackrel{~}{\mathrm{a}}\mathrm{d}_{L𝔭}`$ of $`\stackrel{~}{L}𝔥`$ on the tensor product $`𝒮_{L𝔭}`$. Using the construction of the previous section, we build the twisted Dirac operator
$$\overline{)}_{L𝔥}^{}=\widehat{r}_{L𝔥}^{}+\frac{1}{2}\mathrm{\Omega }_{L𝔥}\mathrm{End}(𝒮_{L𝔭}𝒮_{L𝔥})\mathrm{End}(𝒮_{L𝔤}).$$
Noting that the diagonal action $`\varrho _{L𝔥}^{}=r^{}1+1\stackrel{~}{\mathrm{a}}\mathrm{d}_{L𝔥}`$ on $`𝒮_{L𝔤}`$ is simply the restriction of the action $`\varrho _{L𝔤}=r1+1\stackrel{~}{\mathrm{a}}\mathrm{d}_{L𝔤}`$ to $`L𝔥`$, we obtain the identities
$`\iota _\zeta \overline{)}_{L𝔥}^{}`$ $`=\varrho _{L𝔥}^{}(\zeta )=\varrho _{L𝔤}(\zeta ),`$
$`[\varrho _{L𝔤}(\zeta ),\overline{)}_{L𝔥}^{}]`$ $`=\frac{1}{2}\iota _\zeta \omega _\varrho ^{}^{L𝔥}=\frac{1}{2}\iota _\zeta \omega _\varrho ^{L𝔤},`$
for $`\zeta L𝔥`$. The difference $`\overline{)}_{L𝔤/L𝔥}:=\overline{)}_{L𝔤}\overline{)}_{L𝔥}^{}`$ is basic with respect to $`L𝔥`$, i.e.
$$\iota _\zeta \overline{)}_{L𝔤/L𝔥}=0,[\varrho _{L𝔤}(\zeta ),\overline{)}_{L𝔤/L𝔥}]=0,$$
for all $`\zeta L𝔥`$, and thus it can be written as the $`L𝔥`$-equivariant operator
(35)
$$\overline{)}_{L𝔤/L𝔥}=\widehat{r}_{L𝔭}+\frac{1}{2}\mathrm{\Omega }_{L𝔭}\mathrm{End}(𝒮_{L𝔭})^{L𝔥},$$
where $`\widehat{r}_{L𝔭}`$ is the tautological $`\mathrm{End}()`$-valued 1-form on $`L𝔭`$ given by $`\widehat{r}(\xi )=r(\xi )`$ for $`\xi L𝔤`$, and $`\mathrm{\Omega }_{L𝔭}`$ is the fundamental 3-form given by $`\mathrm{\Omega }_{L𝔭}(\xi ,\eta ,\zeta )=\frac{1}{6}\xi ,[\eta ,\zeta ]`$ for $`\xi ,\eta ,\zeta L𝔭`$. Writing this Dirac operator in terms of a basis $`\{X_i^n\}`$ of $`L𝔭`$ satisfying $`X_i^n,X_j^m=\delta _{i,j}\delta _{n,m}`$, we have
$$\begin{array}{cc}\hfill \overline{)}_{L𝔤/L𝔥}& =\underset{i,n}{}X_i^nr(X_i^n)\frac{1}{12}\underset{i,j,m,n}{}X_i^nX_j^m[X_i,X_j]_𝔭^{n+m},\hfill \end{array}$$
where $`[X,Y]_𝔭`$ denotes the projection of $`[X,Y]`$ onto $`𝔭`$.
As we saw in the finite dimensional case, the two Dirac operators $`\overline{)}_{L𝔥}^{}`$ and $`\overline{)}_{L𝔤/L𝔥}`$ are decoupled, or in other words they anti-commute with each other:
$$\{\overline{)}_{L𝔥}^{},\overline{)}_{L𝔤/L𝔥}\}=\{\widehat{r}_{L𝔥}^{},\overline{)}_{L𝔤/L𝔥}\}+\{\frac{1}{2}\mathrm{\Omega }_{L𝔥},\overline{)}_{L𝔤/L𝔥}\}=0,$$
where the first summand vanishes since for all $`\zeta L𝔥`$ we have
$$\{\widehat{r}_{L𝔥}^{},\overline{)}_{L𝔤/L𝔥}\}(\zeta )=[r^{}(\zeta ),\overline{)}_{L𝔤/L𝔥}]=0,$$
and the second summand vanishes as the odd operators $`\frac{1}{2}\mathrm{\Omega }_{L𝔥}`$ and $`\overline{)}_{L𝔤/L𝔥}`$ act on distinct representations $`𝒮_{L𝔥}`$ and $`𝒮_{L𝔭}`$ and therefore anti-commute. Since these two operators are decoupled, the square of the Dirac operator on $`L𝔤/L𝔥`$ is
$$\begin{array}{cc}\hfill \overline{)}_{L𝔤/L𝔥}^2& =\left(\overline{)}_{L𝔤}\right)^2\left(\overline{)}_{L𝔥}^{}\right)^2\hfill \\ & =2\left(\mathrm{\Delta }_r^{L𝔤}\mathrm{\Delta }_r^{}^{L𝔥}\right)+\frac{1}{2}\left(\omega _\varrho ^{L𝔤}\omega _\varrho ^{}^{L𝔥}\right)+\frac{1}{12}\left(\mathrm{tr}_𝔤\mathrm{\Delta }_{\mathrm{ad}}^𝔤\mathrm{tr}_𝔥\mathrm{\Delta }_{\mathrm{ad}}^𝔥\right).\hfill \end{array}$$
Now consider the case where $`𝔤`$ is simple, $`𝔥`$ is reductive, and $`_𝝀`$ is the irreducible positive energy representation of $`L𝔤`$ with lowest weight $`𝝀`$. Since $`\overline{)}_{L𝔤/L𝔥}`$ is an $`L𝔥`$-equivariant operator on $`_𝝀𝒮_{L𝔭}`$, it is a constant on each of the irreducible subrepresentations of $`L𝔥`$. If $`𝒰_𝝁`$ is the irreducible positive energy representation of $`\stackrel{~}{L}𝔥`$ with lowest weight $`𝝁`$, then using (34), we see that the square of the Dirac operator takes the value
(36)
$$\begin{array}{cc}\hfill (\overline{)}_{L𝔤/L𝔥})^2|_𝝁& =2\varrho _{L𝔤}(I)\varrho _{L𝔤}(_\theta )𝝀𝝆_𝔤^2\hfill \\ & 2\varrho _{L𝔥}^{}(I)\varrho _{L𝔥}^{}(_\theta )+𝝁𝝆_𝔥^2=𝝀𝝆_𝔤^2+𝝁𝝆_𝔥^2,\hfill \end{array}$$
on the $`𝒰_𝝁`$ components of $`_𝝀𝒮_{L𝔭}`$. Note that the non-constant terms vanish since $`\varrho _{L𝔤}`$ and $`\varrho _{L𝔥}^{}`$ agree on $`\stackrel{~}{}\stackrel{~}{L}𝔥`$.
Note that in the above construction, we are using the invariant inner product on $`𝔥`$ obtained by restricting our invariant inner product on $`𝔤`$. When $`𝔤`$ is simple, we use the basic inner product on $`𝔤`$, which is normalized so that $`\alpha _{\text{max}}^2=2`$, where $`\alpha _{\text{max}}`$ is the highest root of $`𝔤`$. We recall that the basic inner product corresponds to the universal central extension $`\stackrel{~}{L}𝔤`$ of $`L𝔤`$, which in turn restricts to give a (not necessarily universal) central extension $`\stackrel{~}{L}𝔥`$ of $`L𝔥`$. Nevertheless, given any positive energy representation $``$ of $`\stackrel{~}{L}𝔤`$, the tensor product $`𝒮_{L𝔤}`$ is a true representation of this central extension $`\stackrel{~}{L}𝔥`$. So, if $`𝔥`$ is reductive, then the squares of the Dirac operators $`\overline{)}_{L𝔥}^{}`$ and $`\overline{)}_{L𝔤/L𝔥}`$ are indeed of the form given by (34) and (36).
On the other hand, if $`𝔤`$ is not simple but rather semi-simple, then the basic inner product on $`𝔤`$ is normalized so that $`\alpha _i^2=2`$, where the $`\alpha _i`$ are the highest roots of each of the simple components of $`𝔤`$. In this case, a projective positive energy representation of $`L𝔤`$ is not necessarily a true representation of the corresponding central extension $`\stackrel{~}{L}𝔤`$, so the expression (34) for the square of the Dirac operator on $`L𝔤`$ is not universal. However, if $`𝔥`$ is reductive, then the expression (36) for the square of the Dirac operator on $`L𝔤/L𝔥`$ does still hold, as the non-constant terms must vanish since the operator commutes with the action of $`L𝔥`$.
## 10. The kernel of the Dirac operator
Given a linear operator $`d:VW`$ between two finite dimensional vector spaces, the alternating sum of the dimensions in the exact sequence
$$0\mathrm{Ker}dV\stackrel{d}{}W\mathrm{Coker}d0$$
vanishes, and it follows that $`\mathrm{Index}d=dimVdimW`$. Furthermore, if $`V`$ and $`W`$ are $`G`$-modules and the operator $`d`$ is $`G`$-equivariant, then the analogous result $`\mathrm{Index}_Gd=VW`$ holds in the representation ring $`R(G)`$. In the infinite dimensional case, this result does not necessarily hold, but for representations of loop groups, it does hold provided that the representations are of finite type and that the operator commutes with rotating the loops.
###### Lemma 13.
If $`𝒱`$ and $`𝒲`$ are representations of $`LG`$ of finite type, and $`𝒟:𝒱𝒲`$ is an $`S^1LG`$-equivariant linear operator, then its $`LG`$-equivariant index is the virtual representation $`\mathrm{Index}_{LG}𝒟=𝒱𝒲`$.
###### Proof.
Since $`𝒟`$ is $`S^1`$-equivariant, it respects the decompositions of $`𝒱`$ and $`𝒲`$ into their constant energy subspaces, and it can be written in the block diagonal form $`𝒟=_k𝒟_k`$, with $`𝒟_k:𝒱(k)𝒲(k)`$. If both $`𝒱`$ and $`𝒲`$ are of finite type, then each of the subspaces $`𝒱(k)`$ and $`𝒲(k)`$ is a finite dimensional $`G`$-module, and so the $`S^1\times G`$-equivariant index of $`𝒟`$ is given by the $`R(G)`$-valued formal power series
$$\mathrm{Index}_{S^1\times G}𝒟=\underset{k}{}z^k\left(𝒱(k)𝒲(k)\right)=\underset{k}{}z^k𝒱(k)\underset{k}{}z^k𝒲(k).$$
Since a representation of the full loop group $`LG`$ is uniquely determined by its constant energy components, the $`LG`$-equivariant index must therefore be the difference of the domain and the range, hence $`\mathrm{Index}_{LG}𝒟=𝒱𝒲`$. ∎
Returning to the notation of the previous section, let $`𝔤`$ be semi-simple, and let $`𝔥`$ be a reductive subalgebra of $`𝔤`$ with maximal rank. If we decompose the spin representation as $`𝒮_{L𝔭}=𝒮_{L𝔭}^+𝒮_{L𝔭}^{}`$, the Dirac operator $`\overline{)}_{L𝔤/L𝔥}`$ interchanges the two half-spin representations and can thus be written as the sum of the operators
$`\overline{)}_{L𝔤/L𝔥}^+`$ $`:_𝝀𝒮_{L𝔭}^+_𝝀𝒮_{L𝔭}^{},`$
$`\overline{)}_{L𝔤/L𝔥}^{}`$ $`:_𝝀𝒮_{L𝔭}^{}_𝝀𝒮_{L𝔭}^+,`$
where $`\overline{)}_{L𝔤/L𝔥}^{}`$ is the adjoint of $`\overline{)}_{L𝔤/L𝔥}^+`$. When we introduced $`\overline{)}_{L𝔤/L𝔥}`$ in (35), we showed that it is $`L𝔥`$-equivariant, and all of our Dirac operators clearly commute with the generator $`_\theta `$ of rotations of the loops. The operator $`\overline{)}_{L𝔤/L𝔥}`$ is therefore $`S^1LH`$-equivariant, and since its domain and range are both of finite type, we may apply Lemma 13. The $`LH`$-equivariant index of $`\overline{)}_{L𝔤/L𝔥}^+`$ is thus the difference
$$\mathrm{Ker}\overline{)}_{L𝔤/L𝔥}^+\mathrm{Ker}\overline{)}_{L𝔤/L𝔥}^{}=_𝝀𝒮_{L𝔭}^+_𝝀𝒮_{L𝔭}^{},$$
which is given by the homogeneous Weyl-Kac formula (13).
On the other hand, to compute the kernel of $`\overline{)}_{L𝔤/L𝔥}=\overline{)}_{L𝔤/L𝔥}^+\overline{)}_{L𝔤/L𝔥}^{}`$ we proceed as in the computation of the kernel of the finite dimensional operator $`\overline{)}_{𝔤/𝔥}`$ in §6. In fact, the proofs of Lemmas 9 and 10 apply equally well in the Kac-Moody setting using the decomposition $`𝒮_{L𝔤/𝔱}𝒮_{L𝔭}𝒮_{L𝔥/𝔱},`$ and we obtain
###### Lemma 14.
For each $`c𝒞`$, the irreducible representation $`𝒰_{c𝛌}`$ of $`\stackrel{~}{L}𝔥`$ with lowest weight $`c𝛌=c(𝛌𝛒_𝔤)+𝛒_𝔥`$ occurs exactly once in the decomposition of $`_𝛌𝒮_{L𝔭}`$.
###### Lemma 15.
If $`𝛍`$ is a weight of $`_𝛌𝒮_{L𝔭}`$ satisfying $`𝛍𝛒_𝔥^2=𝛌𝛒_𝔤^2`$, then there exists a unique affine Weyl element $`w𝒲_𝔤`$ such that $`𝛍𝛒_𝔥=w(𝛌𝛒_𝔤)`$.
Then, in light of our formula (36) for the square of the Dirac operator $`\overline{)}_{L𝔤/L𝔥}`$, we immediately obtain the loop group analogue of Theorem 8.
###### Theorem 16.
Let $`𝔤`$ be a semi-simple Lie algebra with a maximal rank reductive Lie subalgebra $`𝔥`$. Let $`_𝛌`$ and $`𝒰_𝛍`$ be the irreducible representations of $`\stackrel{~}{L}𝔤`$ and $`\stackrel{~}{L}𝔥`$ with lowest weights $`𝛌`$ and $`𝛍`$. The kernel of the operator $`\overline{)}_{L𝔤/L𝔥}`$ on $`_𝛌𝒮_{L𝔭}`$ is
$$\mathrm{Ker}\overline{)}_{L𝔤/L𝔥}=\underset{c𝒞}{}𝒰_{c𝝀},$$
where $`c𝛌=c(𝛌𝛒_𝔤)+𝛒_𝔥`$, and $`𝒞𝒲_𝔤`$ is the subset of affine Weyl elements which map the fundamental Weyl alcove for $`𝔤`$ into the fundamental alcove for $`𝔥`$.
Comparing this result to the homogeneous Weyl-Kac formula (13), we obtain
$$\mathrm{Ker}\overline{)}_{L𝔤/L𝔥}^+=\underset{(1)^c=+1}{}𝒰_{c𝝀},\mathrm{Ker}\overline{)}_{L𝔤/L𝔥}^{}=\underset{(1)^c=1}{}𝒰_{c𝝀}.$$
Taking the kernels of these Dirac operators therefore gives an explicit construction for the multiplet of signed representations of $`\stackrel{~}{L}𝔥`$ corresponding to any given irreducible positive energy representation of $`\stackrel{~}{L}𝔤`$.
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# How to experimentally measure the number 5 of the 𝑆𝑂(5) theory?
## Abstract
According to Wilson’s theory of critical phenomena, critical exponents are universal functions of $`d`$, the dimension of space, and $`n`$, the dimension of the symmetry group. $`SO(5)`$ theory of antiferromagnetism and superconductivity predicts a bicritical point where $`T_N`$ and $`T_c`$ intersect. By measuring critical exponents close to the bicritical point, and knowing that $`d=3`$, one can experimentally measure the number $`5`$ of the $`SO(5)`$ theory.
In a system of many strongly interacting degrees of freedom, it is generally hard to make precise quantitative predictions which can be tested experimentally. Theories of high $`T_c`$ superconductivity generally have to resort to uncontrolled approximations, as a result, it is not possible for experiments to uniquely test the fundamental physical validity of the theory. However, at special values of physical parameters, the basic degrees of freedom may compete so strongly that a new critical point is reached. At this critical point, low energy properties depend only on universal quantities such as the number of space dimension $`d`$ and the dimension of the symmetry group $`n`$, and are independent of the microscopic details of the constituent materials. Close to such critical points, a new kind of simplicity and predictibility becomes possible, and the theoretical foundation can be tested unambiguously by experiments.
$`SO(5)`$ theory predicts a new bicritical point in the two dimensional phase diagram of temperature versus doping where the antiferromagnetic (AF) transition temperature $`T_N`$ intersects the superconducting (SC) transition temperature $`T_c`$. At this critical point, the three component AF order parameter is unified with the two component SC order parameter to form a five component superspin order parameter. At this point, the dimension of the symmetry group $`n`$ is enhanced to $`5`$, and as a result, the critical properties at this point are uniquely different from that of a lower symmetry world. By experimentally tuning into such a bicritical point, and by precisely measuring the critical exponents at this point, experiments can therefore determine the dimension of the symmetry group and test the fundamental validity of the $`SO(5)`$ theory.
The purpose of this paper is to summarize the various exponents predicted by the $`SO(5)`$ theory, and to encourage experiments to measure these exponents. Due to chemical complications, the bicritical point has not yet been clearly identified experimentally in the high $`T_c`$ cuprates. However, such a point does exist in two dimensional organic superconductors which share many physical properties with the high $`T_c`$ cuprates. Here we also review a insightful theoretical analysis by Murakami and Nagaosa, who showed that the the NMR experiments near such a bicritical point measure the dimension of the symmetry group $`n`$ to be very close to $`5`$.
The model: We start with a generic Ginzburg-Laudau form of the $`SO(5)`$ model,
$`H={\displaystyle \frac{1}{2}}{\displaystyle }d^d𝒓[r_c|\stackrel{}{n}|^2+|\stackrel{}{}\stackrel{}{n}|^2+r_s|\stackrel{}{m}|^2+|\stackrel{}{}\stackrel{}{m}|^2+`$
$`2\delta _c|\stackrel{}{n}|^4+4W|\stackrel{}{n}|^2|\stackrel{}{m}|^2+2\delta _s|\stackrel{}{m}|^4].`$ (1)
where $`\stackrel{}{n}`$ and $`\stackrel{}{m}`$ are the order parameters of the SC and the AF respectively. In this paper, we will fix the dimension to $`d=3`$ and the expansion parameter $`ϵ=4d=1`$. The mean field phase diagram and RG flows of above effective Hamiltonian have been derived in Ref.. Defining $`F=\delta _c\delta _sW^2`$, we summarize their results in following: (i) When $`F>0`$, the RG flow converges to the biconical fixed point $`(\delta _c,W,\delta _s)=2\pi ^2(0.0905,0.0847,0.0536)`$ which correponds to tetracritical phenomena in mean field phase diagram. In this case, the AF and SC orders can coexist in the low temperature phase; (ii) When $`F=0`$, the RG flow converges to a Heisenberg fixed point $`\delta _c=\delta _s=W=\frac{2\pi ^2}{13}`$ which corresponds to bicritical behavior. This fixed point has an exact $`SO(5)`$ rotational symmetry. In this case, there is a direct first order transition between AF and SC. (iii) When $`F<0`$, the RG flow goes to unstable region $`(\delta _c\delta _s)<0)`$. The first order transition line between SC and AF branches at a triple critical point and extends until two branches ends at tricritical points. All of above results were discussed and summarized in the schematic diagrams by the authors of Ref.. In the case of $`n=5`$, the bicritical and the tetracritical points are very close in parameter space. Starting from a generic point in the parameter space, there is a rapid RG flow towards the bicritical point, followed by a slow flow from the bicritical to the tetracritical point. Therefore, there is a large regime of parameters where the bicritical behavior dominates, and it is possible to observe the $`SO(5)`$ symmetry. Recent Monte Carlo simulations of the classical $`SO(5)`$ spin models are consistent with this interpretation of the bicritical point.
Static Exponents: We first discuss the static critical phenomena. At the bicritical point $`r_s=r_c=0`$, thermodynamic quantities obey scaling relations. However, $`g=r_cr_s`$ is a relevant parameter, and the scaling theory of a bicritical point requires a crossover critical exponent $`\varphi `$ (Ref.). The scaling postulate for the singular part of the free energy takes the form ($`t=|TT_c(g=0)|`$),
$`F(t,g,\stackrel{}{n},\stackrel{}{m})=t^{2\alpha }f({\displaystyle \frac{g}{t^\varphi }},{\displaystyle \frac{\stackrel{}{n}}{t^\beta }},{\displaystyle \frac{\stackrel{}{m}}{t^\beta }}).`$ (2)
The exponent $`\beta `$ and $`\alpha `$ take the same value as the isotropic vector model.
The critical exponents $`\alpha `$, $`\beta `$, $`\gamma `$, $`\delta `$, $`\nu `$ and $`\eta `$ satisfy the usual scaling relations:
$`\alpha `$ $`=`$ $`2d\nu ,\beta ={\displaystyle \frac{1}{2}}(d2+\eta ),`$
$`\gamma `$ $`=`$ $`\nu (2\eta ),\delta ={\displaystyle \frac{d+2\eta }{d2+\eta }}`$ (3)
Within second order $`ϵ`$ expansion, they are given by (see Ref. P.133 for $`\alpha `$,$`\beta `$,$`\eta `$ $`\nu `$ and $`\delta `$ Ref. P.611 for $`\varphi `$):
$`\alpha `$ $`=`$ $`{\displaystyle \frac{(n4)}{2(n+8)}}ϵ{\displaystyle \frac{(n+2)^2(n+28)}{4(n+8)^3}}ϵ^2`$
$`\beta `$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{3}{2(n+8)}}ϵ+{\displaystyle \frac{(n+2)(2n+1)}{2(n+8)^3}}ϵ^2`$
$`\gamma `$ $`=`$ $`1+{\displaystyle \frac{(n+2)}{2(n+8)}}ϵ+{\displaystyle \frac{(n+2)(n^2+22n+52)}{4(n+8)^3}}ϵ^2`$
$`\delta `$ $`=`$ $`3+ϵ+{\displaystyle \frac{(n^2+14n+60)}{2(n+8)^2}}ϵ^2`$
$`\nu `$ $`=`$ $`{\displaystyle \frac{1}{2}}+{\displaystyle \frac{n+2}{4(n+8)}}ϵ+{\displaystyle \frac{n+2}{8(n+8)^3}}(n^2+23n+60)ϵ^2`$
$`\eta `$ $`=`$ $`{\displaystyle \frac{n+2}{2(n+8)^2}}ϵ+{\displaystyle \frac{n+2}{8(n+8)^4}}(56n+272n^2)ϵ^2`$
$`\varphi `$ $`=`$ $`1+{\displaystyle \frac{n}{2(n+8)}}ϵ+{\displaystyle \frac{n^2+24n+68}{4(n+8)^3}}ϵ^2`$ (4)
Here we list explicitly the values of the critical exponents in the table 1. It is easy to check the scaling law is approximately satisfied.
Critical temperatures: As already discussed in Ref. , the behavior of the SC transition temperature $`T_c`$ and the AF transition temperature $`T_N`$ close to the bicritical point are governed by the $`SO(5)`$ bicritical exponent $`\varphi `$. In the neighborhood of the bi-critical point, divergent quantities generally behave like:
$`\chi (T,g)t^{\gamma _5}X(g/t^\varphi )`$ (5)
where $`t=TT_c(g=0)`$ and $`X(z)`$ is a scaling function, normalized such that $`X(0)=1`$. However, unlike the usual scaling functions, $`X(z)`$ diverges at two points $`z_2>0`$ and $`z_3<0`$:
$`X(z)(zz_2)^{\gamma _2};X(z)(zz_3)^{\gamma _3}`$ (6)
Therefore, sufficiently close to the $`SO(5)`$ bicritical point, $`g/t^\varphi 1`$, and the critical behavior is given by the new $`SO(5)`$ exponent $`\gamma _5`$. Away from the $`SO(5)`$ bicritical point, the divergence of physical quantities are determined by the divergence of $`X(z)`$. For $`g>0`$, the critical temperature is given by $`g/t^\varphi =z_2`$, or
$`T_c(g)=T_c(0)+Ag^{1/\varphi }`$ (7)
where $`A`$ is a constant. Similar arguments applies for the case of $`g<0`$. This way, by measuring the precise values of both $`T_c`$ and $`T_N`$ close to the bicritical point, one can determine the value of the crossover exponent $`\varphi `$ and compare it with the $`SO(5)`$ prediction of $`\varphi =1.314`$.
London penetration length: Near the superconducting to normal phase transition, this quantity scales like
$`\lambda \rho _s^{\frac{1}{2}}\xi ^{(2D)/2}t^{\nu /2}`$ (8)
This is a very interesting quantity, since its critical behavior has already been measured by Kamal et al. for YBCO superconductors, and the exponent was found to be consistent with the XY value of $`\nu _2=0.655`$. Here we suggest to measure the critical behavior of $`\lambda `$ for doping levels ranging from optimal to deeply underdoped regime. The $`SO(5)`$ theory predicts that data for all doping levels $`x`$ can be fit into a single scaling curve:
$`\chi (T,x)t^{\nu _5/2}Y(x/t^\varphi )`$ (9)
where $`Y(0)=1`$ and it diverges near $`z_2>0`$ as:
$`Y(z)(zz_2)^{\nu _2/2}`$ (10)
If one can get sufficiently close to the bicritical regime, one can determine both $`\nu _5`$ and $`\varphi `$ and compare with the $`SO(5)`$ predictions of $`\nu _5=0.714`$ and $`\varphi =1.314`$. Together with these precisely predicted values, the fitting into a single scaling curve for all doping levels provides a highly nontrivial quantitative test of the $`SO(5)`$ theory.
Other static quantities should all follow similar scaling relations close to the bicritical point.
Dynamic exponent: Under the dynamic scaling hypothesis, the typical frequency or the relaxation rate $`\omega `$ scales as $`\omega \xi ^z`$. Standard arguments in dynamical critical phenomena gives $`z=d/2`$ generally, and $`z=\varphi /\nu `$ near a bicritical point.
Nuclear magnetic relaxation rate: $`1/T_1`$ is given by the following response function:
$`1/T_1=lim_{\omega 0}{\displaystyle \frac{1}{\omega }}{\displaystyle \frac{d^dk}{(2\pi )^d}\chi (k,\omega )}`$ (11)
where $`\chi (k,\omega )`$ is the imaginary part of the spin response function, which near a critical point behaves like
$`\chi (k,\omega )=\xi ^{2\eta }Y(\overline{k},\overline{\omega })`$ (12)
Here $`Y(\overline{k},\overline{\omega })`$ is a scaled, dimensionless spin correlation function of the dimensionless variable $`\overline{k}=k\xi `$ and $`\overline{\omega }=\omega \xi ^z`$. Expressed in terms of the rescaled variables,
$`1/T_1=\xi ^{zd+2\eta }lim_{\overline{\omega }0}{\displaystyle \frac{1}{\overline{\omega }}}{\displaystyle \frac{d^d\overline{k}}{(2\pi )^d}Y(\overline{k},\overline{\omega })}`$ (13)
from which we can see that the scaling behavior of $`1/T_1`$ is given by:
$`1/T_1=\xi ^{zd+2\eta }=t^x`$ (14)
where $`x=\nu (z1\eta )`$. Applying the results of the $`ϵ`$ expansion listed in the previous table, we obtain $`x=\nu (z1\eta )=0.573`$ close to a $`SO(5)`$ bicritical point. For a regular antiferromagnetic transition of the $`O(3)`$ symmetry class, we obtain $`x=\nu (z1\eta )=0.67(1.510.039)=0.312.`$
Frequency-dependent conductivity: In the superconducting state $`T<T_c`$, and for low frequency, the complex conductivity takes the form
$`\sigma (\omega ){\displaystyle \frac{\rho _s}{i\omega }}`$ (15)
In the critical region, as we pointed out before, $`\rho _s\xi ^{2d}`$, therefore, the dynamic conductivity scales as
$`\sigma (\omega )\xi ^{2d}/\omega \omega ^{\frac{z+2d}{z}}`$ (16)
at $`T=T_c`$ for low frequency. At a $`SO(5)`$ biciritical point, the exponent is given by $`\frac{z+2d}{z}=0.46`$ compared with $`\frac{z+2d}{z}=0.33`$ for a ordinary superconductor to normal transition in the $`XY`$ universality class.
Experimental status: Possibly due to chemical complications, it is hard to reach a uniform state in the deeply underdoped regime of the high $`T_c`$ superconductors. For this reason, the existence of a bicritical in the high $`T_c`$ superconductors has neither been discovered nor refuted. One of the greatest experimental challenge in this field is to prepare better and more uniform materials in the deeply underdoped regime. Given such materials, it is most feasible to measure the critical properties of the London penetration length by microwave cavity experiment. As clearly demonstrated in this work, such system can provide definite and quantitative test of the $`SO(5)`$ theory of high $`T_c`$ superconductivity.
Encouraging experimental evidence for a $`SO(5)`$ bicritical point does exist in a class of 2D organic superconductors called $`bedt`$ salt. These material share most common physical properties with the cuprates, and a AF to SC transition can be induced by pressure. A bicritical point exists where $`T_c`$ and $`T_N`$ intersect each other. Kanoda and coworkers measured the $`1/T_1`$ rate both in the AF region and the bicritical region. Murakami and Nagaosa analyzed the experimental data. The $`1/T_1`$ exponent in the AF region was measured to be $`x_{AF}=0.30`$, compared with the theoretical prediction of $`x_3=0.312`$. In the bicritical region, the experimental fit gives $`x_{bi}=0.56`$, compared with the $`SO(5)`$ theoretical prediction of $`x_5=0.573`$. This is the first experiment which directly measures the dimension of the symmetry group close to a AF/SC bicritical point, and determines $`n`$ to be close to $`5`$.
Conclusions: $`SO(5)`$ theory makes precise and quantitative predictions on the critical exponents near a bicritical point. We strongly encourage experiments to be carried out in deeply underdoped regime of high $`T_c`$ superconductors and to look for the bicritical point. Measurement of the critical exponents associated with various physical quantities can uniquely test the fundamental validity of the $`SO(5)`$ theory, and can measure the number $`5`$ of the the $`SO(5)`$ theory in a direct and unambiguous fashion.
Acknowledgement: We would like to thank Dr. S. Murakami and N. Nagaosa for stimulating communications, their work directly inspired ours. We would like to thank the organizers of the m2s meeting for the invitation to contribute to this proceeding. This work is supported by the NSF under grant numbers DMR-9814289 and DMR-9400372. JP. Hu. is supported by the Stanford Graduate Fellowship Program.
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# Braggoriton–Excitation in Photonic Crystal Infiltrated with Polarizable Medium
## I Introduction
Photonic crystals and in particular photonic band gap (PBG) materials have recently attracted much attention due to their rich physics and possible applications. In these systems the dielectric function is periodically modulated and, as a result, their optical properties are dominated by light diffraction effects. When Bragg diffraction conditions are met then light scattering is very strong, so that within certain frequency intervals near the resonances light propagation is inhibited.
Since the subject of photonic crystals was introduced, one of the main goals of photonic band-structure calculations was to engineer structures with a complete band gap, i.e. with no propagating solutions of Maxwell’s equations within a certain forbidden gap. The pursuit of this goal has generated a stream of studies that are too numerous to be cited here; early works are reviewed in Refs.. Here we only mention that a complete band gap in two dimensions (2D) was theoretically predicted and experimentally demonstrated for an array of dielectric rods. In the quest for a structure having a complete PBG in three dimensions (3D), the diamond lattice was shown to be more promising than a simple face centered cubic (fcc) lattice.
The frequency gap in the photonic spectrum sets a stage for a number of physical effects. The prime effect, namely the inhibition of spontaneous emission for an emitter with transition frequency within the gap, was already suggested in the pioneering works. Furthermore, since light cannot leave the emitting atom, a coupled atom-field in-gap state is formed, in which the atomic level is “dressed” by its own exponentially localized radiation field. It was also demonstrated that although a single photon cannot propagate inside the gap, nevertheless a non-linear medium embedded inside the photonic crystal gives rise to multi-photon bound states, or gap solitons that result in self-induced transparency. Yet another consequence of PBG is the modification of cooperative emission with frequency close to the band edge. In particular PBG was shown to change the rate of superradiant emission from an ensemble of emitters. Lastly, PBG structures facilitate strong Anderson localization of photons because the sharp density of states within the gap spectral range necessitates a reinterpretation of the Ioffe-Regel criterion.
PBG structures with a defect constitute a separate area of study initiated by the classical works in Refs.. These structures are important since the defects cause localized intragap states. For these states, the PBG sample acts as a resonator with a very high quality factor. This property was recently used for designing a low-threshold PBG defect-mode laser.
Another class of materials with a forbidden gap for light propagation are spatially homogeneous, but frequency-dispersive media. The energy gap in these systems has a polaritonic origin, i.e. it is formed due to the interaction of light with the medium polarization. This energy gap can be viewed as the result of anticrossing between the photonic and excitonic dispersion relation branches. Some non-trivial manifestations of the polaritonic gap were recently explored in Refs.. In these papers a general model of two-level systems interacting with elementary electromagnetic excitations with a gap in the spectrum was solved by means of the Bethe ansatz technique. Within this model a very rich excitation spectrum was found, consisting of ordinary solitons, single-particle impurity bound states and massive pairs of confined gap excitations and their bound complexes — dissipationless quantum gap solitons.
Most of the available photonic crystals nowadays however have incomplete PBGs; this means that light propagation is forbidden only along certain directions inside the crystal. A prominent example are opals, representing self-assembled monodispersed silica balls arranged in a fcc type lattice. Although opals have only incomplete PBG, the voids between the balls can be infiltrated by various media, which brings about non-trivial physics. In particular, the medium may contain polarizable molecules. Infiltrated opal with polarizable molecules combines therefore polaritonic and Bragg-diffractive properties. Obviously, both effects coexist independently when the Bragg ($`\omega =\omega _\text{B}`$) and polaritonic ($`\omega =\omega _\text{T}`$) resonances are well separated in frequency. A completely different situation occurs when $`\omega _\text{B}\omega _\text{T}`$. This may be easily achieved in infiltrated opals that gives rise to a peculiar interplay between various frequency dispersions. This interplay is the subject of the present paper.
Our most important finding pertains to the case when the polaritonic gap of the polarized molecules infiltrating the opal lies within the opal PBG. We demonstrate that such an overlap gives rise to novel massive propagating excitations having frequencies inside the Bragg gap, that we dubbed here braggoritons. In other words, the Bragg gap splits into two sub-gaps, so that the braggoritonic branches are isolated from the rest of the spectrum. We found that the braggoriton dispersion relation is very sensitive to the frequency detuning between $`\omega _\text{B}`$ and $`\omega _\text{T}`$ and to the relative width of the polaritonic gap (or, alternatively Rabi frequency) and the Bragg gap.
The principal assumption we adopt here is that the Bragg gap, $`\mathrm{\Delta }\omega _\text{B}`$, is narrow compared to $`\omega _\text{B}`$; this is actually the case in opals. The small value of $`\mathrm{\Delta }\omega _\text{B}/\omega _\text{B}`$ ($`1`$) enables us then to obtain analytical results. In addition we also study the phase slip related intragap defect states for $`\omega _\text{B}\omega _\text{T}`$. In the absence of polaritonic effect, the underlying physics of the defect-induced intragap states was already discussed in the original PBG paper. An analogy was drawn between a defect state and a localized mode in a distributed feedback resonator, that originates from a phase slip. We extend this picture to incorporate polarizable medium and show that when the Bragg gap splits into two sub-gaps, then an existing phase slip gives rise to two localized states with frequencies within each of the sub-gaps.
Our paper is organized as follows: In order to introduce the notations, we separately review in Sec. II the derivation of the PBG and polaritonic spectra using the second quantization representation. In Sec. III we consider the combined Hamiltonian in the second quantization representation and diagonalize it by a unitary transformation. This yields the dispersion relations for the two excitations outside the gap, or Bloch-like waves, and the two intragap branches, or braggoritons. The properties of these novel excitations are analyzed in Sec. IV. We use them in Sec. V to determine the intragap frequencies of the defect-induced localized states. Concluding remarks are presented in Sec. VI.
## II Second quantized PBG and polaritonic Hamiltonians
The Hamiltonian $``$ of the system under study is the sum of three terms
$$=_{\text{ph}}+_\text{m}+_{\text{m-ph}}.$$
(1)
The first term, $`_{\text{ph}}`$, describes the photons in a photonic crystal. The second term, $`_\text{m}`$, is the Hamiltonian of the polarizable medium; $`_{\text{m-ph}}`$ describes the photon-medium coupling. In this Section we review two limiting cases: (i) no polarizable medium ($`_\text{m}0`$), and (ii) no modulation of the dielectric constant.
### A Incomplete PBG
The general form of the Hamiltonian $`_{\text{ph}}`$ is
$$_{\text{ph}}=\frac{1}{8\pi }𝑑𝒓\left[\epsilon (𝒓)𝑬^2+𝑯^2\right],$$
(2)
where $`𝑬`$ and $`𝑯`$ are respectively the electric and magnetic fields. For a constant dielectric function, $`\epsilon (𝒓)\epsilon _0`$, the second quantized form of the Hamiltonian (2) reduces to a sum over oscillators representing plane waves with frequencies $`\omega _k=ck/\sqrt{\epsilon _0}`$, where $`k`$ is the wave vector. Modulation of $`\epsilon (𝒓)`$ causes light diffraction, so that the plane wave solutions are no longer the correct eigenfunctions of the Hamiltonian (2). Below we consider a photonic crystal with an incomplete PBG along the $`z`$ axis. This situation can be approximated by a one-dimensional modulation of $`\epsilon (𝒓)`$ along the $`z`$ direction. If the Bragg gap is relatively narrow as assumed above, then it may be sufficient to consider only the first harmonics in $`\epsilon (z)`$:
$$\epsilon (z)=\epsilon _0+\delta \epsilon \mathrm{cos}(\sigma z+\varphi ).$$
(3)
Here $`\delta \epsilon `$ ($`\epsilon _0`$) is the modulation amplitude, $`\sigma =2\pi /d`$, where $`d`$ is the modulation period and $`\varphi `$ is the dielectric modulation phase. We assume for simplicity that the electromagnetic field propagates along the $`z`$-direction and is homogeneous in the $`xy`$ plane. In this case light polarization is irrelevant. Generalization to arbitrary propagation direction is straightforward. The Fourier components of $`\epsilon (z)`$ in (3) couple the original photon oscillators with momenta $`k`$ and $`k\pm \sigma `$. These coupled oscillators form an infinite series that is constructed by successive addition (subtraction) of $`\sigma `$. However, if $`\delta \epsilon \epsilon _0`$ and the wavevectors domain is restricted to the vicinity of the first Bragg resonance at $`k\sigma /2`$, then the Hamiltonian (2) can be truncated. In this case only the coupling to the near-resonance backscattered photons with momenta $`(\sigma k)\sigma /2`$ must be retained, so that the Hamiltonian (1) takes the form
$`_{\text{ph}}=`$ (4)
$`{\displaystyle \underset{q}{}}\{\omega (q)\widehat{a}_{}^+(q)\widehat{a}_{}(q)+\omega (q)\widehat{a}_{}^+(q)\widehat{a}_{}(q)`$ (5)
$`+\mathrm{\Omega }_\text{B}[e^{i\varphi }\widehat{a}_{}^+(q)\widehat{a}_{}(q)+e^{i\varphi }\widehat{a}_{}^+(q)\widehat{a}_{}(q)]\}.`$ (6)
Here, we introduced the notations: $`q=k\sigma /2`$, $`\widehat{a}_{}(q)=\widehat{a}_k`$ and $`\widehat{a}_{}(q)=\widehat{a}_{k\sigma }`$ for $`k\sigma /2`$, where $`\widehat{a}_k`$ is the usual photon annihilation operator. In the notations introduced in (6), the frequencies of the photonic branches are given by:
$$\omega (q)=\frac{c(q+\sigma /2)}{\sqrt{\epsilon _0}}=\omega _\text{B}\left(1+\frac{2q}{\sigma }\right),$$
(7)
where $`\omega _\text{B}=c\sigma /(2\sqrt{\epsilon _0})`$ is the Bragg frequency. We define the coupling constant, $`\mathrm{\Omega }_\text{B}`$, as the half-width of the Bragg gap, i.e. $`\mathrm{\Omega }_\text{B}=\frac{1}{2}\mathrm{\Delta }\omega _\text{B}`$. It can be shown that $`\mathrm{\Omega }_\text{B}`$ is related to the amplitude of the dielectric function modulation: $`\mathrm{\Omega }_\text{B}=\omega _\text{B}\delta \epsilon /(2\epsilon _0)`$. The summation in Eq. (6) is performed over the $`k`$ domain $`|q|\sigma /2`$.
It is straightforward to diagonalize the Hamiltonian in Eq. (6) with the use of the following unitary transformation
$$\begin{array}{c}\widehat{a}_{}(q)=\mathrm{cos}\theta \widehat{\beta }_1(q)+\mathrm{sin}\theta e^{i\varphi }\widehat{\beta }_2(q),\\ \widehat{a}_{}(q)=\mathrm{sin}\theta e^{i\varphi }\widehat{\beta }_1(q)+\mathrm{cos}\theta \widehat{\beta }_2(q),\end{array}$$
(8)
where
$$\mathrm{cos}2\theta =\frac{\omega (q)\omega (q)}{\sqrt{\left[\omega (q)\omega (q)\right]^2+4\mathrm{\Omega }_\text{B}^2}}.$$
(9)
The new operators $`\widehat{B}_1`$ and $`\widehat{B}_2`$ describe the creation (annihilation) of pairs of Bloch waves, that consist of forward and backscattered photons near the Bragg frequency. The diagonalized Hamiltonian (6) takes the form
$$_{\text{ph}}=\underset{q}{}\omega _\text{B}^{(1)}(q)\widehat{\beta }_1^+(q)\widehat{\beta }_1(q)+\omega _\text{B}^{(2)}(q)\widehat{\beta }_2^+(q)\widehat{\beta }_2(q),$$
(10)
where the dispersion relations of the two photonic branches are given by:
$`\omega _\text{B}^{(1,2)}(q)=`$ (11)
$`{\displaystyle \frac{1}{2}}\left[\omega (q)+\omega (q)\pm \sqrt{(\omega (q)\omega (q))^2+4\mathrm{\Omega }_\text{B}^2}\right]=`$ (12)
$`\omega _\text{B}\pm \sqrt{\left({\displaystyle \frac{2\omega _\text{B}}{\sigma }}\right)^2q^2+\mathrm{\Omega }_\text{B}^2}.`$ (13)
As mentioned above, the width of the PBG from Eq. (13) is $`2\mathrm{\Omega }_\text{B}`$.
### B Polarizable medium
The Hamiltonians $`_\text{m}`$ of polarizable medium and $`_{\text{m-ph}}`$ of light-polarization interaction in Eq. (1) can be written in the second quantization form as:
$`_\text{m}+_{\text{m-ph}}`$ $`=`$ $`\omega _\text{T}{\displaystyle \underset{k}{}}\widehat{b}_k^+\widehat{b}_k`$ (14)
$`+`$ $`\mathrm{\Omega }_\text{P}{\displaystyle \underset{k}{}}\left[\widehat{b}_k^+\widehat{a}_k+\widehat{a}_k^+\widehat{b}_k\right],`$ (15)
where $`\widehat{b}`$ is the annihilation operator of the medium excitations (e.g. excitons or optical phonons), which we assume here to be dispersionless having frequency $`\omega _\text{T}`$; $`\mathrm{\Omega }_\text{P}`$ denotes the light-medium coupling strength that is proportional to the Rabi frequency. In the absence of the Bragg scattering term ($`\mathrm{\Omega }_\text{B}=0`$) the complete Hamiltonian (1)reduces to the conventional polaritonic Hamiltonian $`_{\text{pol}}`$, given by:
$`_{\text{pol}}`$ $`=`$ $`{\displaystyle \underset{k}{}}\omega _k\widehat{a}_k^+\widehat{a}_k+\omega _\text{T}{\displaystyle \underset{k}{}}\widehat{b}_k^+\widehat{b}_k`$ (16)
$`+`$ $`\mathrm{\Omega }_\text{P}{\displaystyle \underset{k}{}}\left[\widehat{b}_k^+\widehat{a}_k+\widehat{a}_k^+\widehat{b}_k\right],`$ (17)
with eigenstates representing the mixture of light and medium excitations
$$\begin{array}{c}\widehat{a}_k=\mathrm{cos}\psi \widehat{\pi }_1(k)+\mathrm{sin}\psi \widehat{\pi }_2(k),\\ \widehat{b}_k=\mathrm{sin}\psi \widehat{\pi }_1(k)+\mathrm{cos}\psi \widehat{\pi }_2(k),\end{array}$$
(18)
where
$$\mathrm{cos}2\psi =\frac{\omega _k\omega _\text{T}}{\sqrt{(\omega _k\omega _\text{T})^2+4\mathrm{\Omega }_\text{P}^2}},\text{(}\omega _k=ck/\sqrt{\epsilon _0}\text{)}.$$
(19)
With the new operators $`\widehat{\pi }_1`$ and $`\widehat{\pi }_2`$, the Hamiltonian (17) is diagonalized:
$`_{\text{pol}}`$ $`=`$ $`{\displaystyle \underset{k}{}}\omega _\text{P}^{(1)}(k)\widehat{\pi }_1^+(k)\widehat{\pi }_1(k)`$ (20)
$`+`$ $`{\displaystyle \underset{k}{}}\omega _\text{P}^{(2)}(k)\widehat{\pi }_2^+(k)\widehat{\pi }_2(k),`$ (21)
where the frequencies of the polaritonic branches are given by
$$\omega _\text{P}^{(1,2)}(k)=\frac{1}{2}\left[\omega _k+\omega _\text{T}\pm \sqrt{(\omega _k\omega _\text{T})^2+4\mathrm{\Omega }_\text{P}^2}\right].$$
(22)
The Rabi splitting at resonance, i.e. at $`\omega _k=\omega _\text{T}`$ is $`2\mathrm{\Omega }_\text{P}`$. Eq. (22) allows to express the phenomenological parameter $`\mathrm{\Omega }_\text{P}`$ through the observables. Namely, $`\mathrm{\Omega }_\text{P}=\sqrt{\omega _\text{T}\omega _{\text{LT}}/2}`$, where $`\omega _{\text{LT}}\omega _\text{T}`$ is the transverse-longitudinal splitting.
We note that the above description is valid only for wavenumbers $`k`$ in the vicinity of “crossing” of the excitation branches, where $`\omega _k\omega _\text{T}`$. It does not capture, however, the correct behavior of the polaritonic branches for $`k0`$. In this limit an additional term of the type $`a_kb_k+\text{c.c.}`$ should be taken into account in Eq. (17).
## III Diagonalization of the full Hamiltonian
Now let us consider the full Hamiltonian (1) with both Bragg scattering and light-medium interaction included,
$``$ $`=`$ $`{\displaystyle \underset{q}{}}\left[\omega (q)\widehat{a}_{}^+(q)\widehat{a}_{}(q)+\omega (q)\widehat{a}_{}^+(q)\widehat{a}_{}(q)\right]`$ (23)
$`+`$ $`\mathrm{\Omega }_\text{B}{\displaystyle \underset{q}{}}\left[e^{i\varphi }\widehat{a}_{}^+(q)\widehat{a}_{}(q)+e^{i\varphi }\widehat{a}_{}^+(q)\widehat{a}_{}(q)\right]`$ (24)
$`+`$ $`\omega _\text{T}{\displaystyle \underset{q}{}}\left[\widehat{b}_{}^+(q)\widehat{b}_{}(q)+\widehat{b}_{}^+(q)\widehat{b}_{}(q)\right]`$ (25)
$`+`$ $`\mathrm{\Omega }_\text{P}{\displaystyle \underset{q}{}}[\widehat{b}_{}^+(q)\widehat{a}_{}(q)+\widehat{a}_{}^+(q)\widehat{b}_{}(q)`$ (26)
$`+`$ $`\widehat{b}_{}^+(q)\widehat{a}_{}(q)+\widehat{a}_{}^+(q)\widehat{b}_{}(q)],`$ (27)
where we have again truncated the “Bragg” Hamiltonian in (1) by including only near-resonance terms. If a column of operators $`\widehat{c}=\{\widehat{a}_{}(q),\widehat{a}_{}(q),\widehat{b}_{}(q),\widehat{b}_{}(q)\}`$ is introduced, then the Hamiltonian (23) can be formally rewritten in a matrix form $`=\widehat{c}^+H\widehat{c}`$, where
$$H=\left(\begin{array}{cccc}\omega (q)& \mathrm{\Omega }_\text{B}e^{i\varphi }& \mathrm{\Omega }_\text{P}& 0\\ \mathrm{\Omega }_\text{B}e^{i\varphi }& \omega (q)& 0& \mathrm{\Omega }_\text{P}\\ \mathrm{\Omega }_\text{P}& 0& \omega _\text{T}& 0\\ 0& \mathrm{\Omega }_\text{P}& 0& \omega _\text{T}\end{array}\right).$$
(28)
The four $`H`$ eigenvalues yield the dispersion relations of the four excitation branches, whereas the eigenvectors determine the unitary transformation diagonalizing $``$.
The characteristic equation for the eigenvalues $`\omega `$ of the matrix $`H`$ reads
$$\left[\left(\omega (q)\omega \right)\left(\omega (q)\omega \right)\mathrm{\Omega }_\text{B}^2\left(\frac{\omega (q)+\omega (q)2\omega }{\omega _\text{T}\omega }\right)\mathrm{\Omega }_\text{P}^2\right](\omega _\text{T}\omega )^2+\mathrm{\Omega }_\text{P}^4=0.$$
(29)
If the light-matter coupling is absent, i.e. $`\mathrm{\Omega }_\text{P}=0`$, then the roots of (29) reduce to two pure medium excitations with unperturbed frequency $`\omega _\text{T}`$ propagating in the forward and backward directions along $`z`$, and two purely photonic excitations with disperion relation given by Eq. (13) that results from the Bragg scattering. If, on the other hand, the Bragg scattering is absent, i.e. $`\mathrm{\Omega }_\text{B}=0`$, then the roots of Eq. (29) reduce to two pairs of polariton branches with dispersion relation given by Eq. (22).
It appears that the unitary transformation in four-dimensional space diagonalizing the Hamiltonian $``$ can be parameterized by two angles:
$$\left(\begin{array}{c}\widehat{a}_{}(q)\\ \widehat{a}_{}(q)\\ \widehat{b}_{}(q)\\ \widehat{b}_{}(q)\end{array}\right)=\left(\begin{array}{cccc}\mathrm{cos}\theta \mathrm{cos}\stackrel{~}{\psi }& \mathrm{sin}\theta \mathrm{sin}\stackrel{~}{\psi }& \mathrm{cos}\theta \mathrm{sin}\stackrel{~}{\psi }& \mathrm{sin}\theta \mathrm{cos}\stackrel{~}{\psi }\\ \mathrm{sin}\theta \mathrm{cos}\stackrel{~}{\psi }& \mathrm{cos}\theta \mathrm{sin}\stackrel{~}{\psi }& \mathrm{sin}\theta \mathrm{sin}\stackrel{~}{\psi }& \mathrm{cos}\theta \mathrm{cos}\stackrel{~}{\psi }\\ \mathrm{cos}\theta \mathrm{sin}\stackrel{~}{\psi }& \mathrm{sin}\theta \mathrm{cos}\stackrel{~}{\psi }& \mathrm{cos}\theta \mathrm{cos}\stackrel{~}{\psi }& \mathrm{sin}\theta \mathrm{sin}\stackrel{~}{\psi }\\ \mathrm{sin}\theta \mathrm{sin}\stackrel{~}{\psi }& \mathrm{cos}\theta \mathrm{cos}\stackrel{~}{\psi }& \mathrm{sin}\theta \mathrm{cos}\stackrel{~}{\psi }& \mathrm{cos}\theta \mathrm{sin}\stackrel{~}{\psi }\end{array}\right)\left(\begin{array}{c}\widehat{𝓑}_2\\ \widehat{\text{a}}_2\\ \widehat{\text{a}}_1\\ \widehat{𝓑}_1\end{array}\right),$$
(30)
where $`\widehat{𝓑}_1`$, $`\widehat{\text{a}}_1`$, $`\widehat{\text{a}}_2`$, and $`\widehat{𝓑}_2`$ are new operators that annihilate mixed light-matter states. The angle $`\theta `$ in Eq. (30) is precisely the “Bragg” rotation angle introduced in Eq. (9) \[for simplicity we set $`\varphi =0`$ for the modulation phase in Eq. (30)\]. The second angle, $`\stackrel{~}{\psi }`$ is defined by the following relation
$$\mathrm{cos}2\stackrel{~}{\psi }=\frac{\sqrt{\left[\omega (q)\omega (q)\right]^2+4\mathrm{\Omega }_\text{B}^2}2(\omega _\text{T}\omega _\text{B})}{\sqrt{\left\{2(\omega _\text{T}\omega _\text{B})\sqrt{\left[\omega (q)\omega (q)\right]^2+4\mathrm{\Omega }_\text{B}^2}\right\}^2+16\mathrm{\Omega }_\text{P}^2}}.$$
(31)
Naturally, for $`\mathrm{\Omega }_\text{B}=0`$ the angle $`\stackrel{~}{\psi }`$ reduces to the polaritonic rotation angle $`\psi `$ in Eq. (19). In the presence of the Bragg scattering, however, this rotation angle also depends on the “Bragg” parameters $`\omega _\text{B}`$ and $`\mathrm{\Omega }_\text{B}`$. Therefore it is the angle $`\stackrel{~}{\psi }`$ that characterizes the interplay between the polaritonic and diffraction effects.
## IV Braggoritonic excitations
In order to analyze the solutions of Eq. (29), it is convenient to introduce the following dimensionless variables. We measure frequencies, $`\mathrm{\Delta }`$, from the Bragg frequency, $`\omega _\text{B}`$, and express them in units of the Bragg gap $`2\mathrm{\Omega }_\text{B}`$:
$$\mathrm{\Delta }=\frac{\omega \omega _\text{B}}{2\mathrm{\Omega }_\text{B}}.$$
(32)
In analogy with Eq. (32) we introduce the dimensionless frequency detuning, $`\delta `$, of $`\omega _\text{T}`$ from the Bragg frequency $`\omega _\text{B}`$, where
$$\delta =\frac{\omega _\text{T}\omega _\text{B}}{2\mathrm{\Omega }_\text{B}}.$$
(33)
As seen from Eq. (13), the natural unit for the wavevector deviation, $`q`$, from the Bragg wavevector, $`\sigma /2`$, is $`\sigma \mathrm{\Omega }_\text{B}/\omega _\text{B}`$. Hence, we introduce the dimensionless parameter
$$Q=\left(\frac{\omega _\text{B}}{\sigma \mathrm{\Omega }_\text{B}}\right)q.$$
(34)
Fig. 1 Dispersion of mixed photonic-medium excitations for various coupling strengths: $`\alpha =0,0.35,0.57`$. (a) $`\delta =0`$; (b) $`\delta =0.3`$. The shaded area represents the two forbidden sub-gaps at $`\alpha =0.35`$.
With the new notations, the excitation spectrum determined by Eq. (29) can be rewritten in a more concise form
$$\left|Q\right|=\sqrt{\left(\mathrm{\Delta }\frac{\alpha ^2}{\mathrm{\Delta }\delta }\right)^2\frac{1}{4}},$$
(35)
where $`\alpha =\mathrm{\Omega }_\text{P}/(2\mathrm{\Omega }_\text{B})`$ characterizes the relative strength of the Bragg and polaritonic couplings. Expression (35) is our main result. It clearly demonstrates that the Bragg and polaritonic dispersion relations compete with each other.
Consider for simplicity the case of exact resonance, i.e. $`\delta =0`$. It is seen from Eq. (35) that in the absence of light-matter coupling ($`\alpha =0`$), the first term in the brackets gives rise to the conventional PBG. It is also seen that with increasing $`\alpha `$ (or, $`\mathrm{\Omega }_\text{P}`$) the decay length, $`\text{Im}Q^1`$, increases, and for sufficiently small $`\mathrm{\Delta }`$ we find Fig. 2 The effective mass for various excitations (in units of “free” Bragg mass $`M_𝓑=M_{𝓑_1}=M_{𝓑_2}`$ at $`\alpha =0`$) is plotted vs. coupling strength: solid lines are for $`\delta =0`$ where $`M_{𝓑1}=M_{𝓑2}`$, $`M_{\text{a}_1}=M_{\text{a}_2}`$; dashed lines are for $`\delta =0.3`$.
that $`Q`$ becomes real. This manifests the emergence of the novel allowed photonic states, or braggoritons, inside the PBG (see Fig. 1). The braggoriton branches in the excitation dispersion relations are described by the operators $`\widehat{\text{a}}_1`$ and $`\widehat{\text{a}}_2`$. They occupy the frequency ranges $`\mathrm{\Delta }=[0,\pm \frac{1}{4}\left(\sqrt{1+16\alpha ^2}1\right)]`$. For small $`\alpha `$ ($`\alpha 1`$), the braggoriton frequency interval reduces to $`(0,\pm 2\alpha ^2)`$. We note that due to the finite $`\alpha `$ value the Bragg gap broadens. Namely, the band edges of the branches described by the operators $`\widehat{𝓑}_1`$, $`\widehat{𝓑}_2`$ are respectively given for $`\delta =0`$ by $`\mathrm{\Delta }=\pm \frac{1}{4}\left(\sqrt{1+16\alpha ^2}+1\right)`$ (compare to $`\mathrm{\Delta }=\pm 1/2`$ for $`\alpha =0`$). The dispersion relations $`\mathrm{\Delta }(Q)`$ calculated using Eq. (35) are shown in Fig. 1(a) for different values of $`\alpha `$ in the case of exact resonance $`\omega _\text{B}=\omega _\text{T}`$, or $`\delta =0`$.
Moderate frequency detuning $`\delta 0`$ does not qualitatively change the above picture as seen in Fig. 1(b). The major effect of frequency detuning is that the braggoritonic branches $`\text{a}_1`$, $`\text{a}_2`$ acquire an asymmetry since they are “pinned” by $`\omega _\text{T}`$. Fig. 1 also shows that the Bragg-like photonic branches $`𝓑_1`$, $`𝓑_2`$ are affected by coupling or detuning only weakly.
To quantitatively describe the braggoriton dispersion relation, we consider two characteristics: (i) dimensionless effective mass, $`M`$ near the band edges that is defined from Eq. (35) by the relation,
$$\mathrm{\Delta }\mathrm{\Delta }_{Q=0}+\frac{Q^2}{2M},$$
(36)
and (ii) the density of states, $`N(\mathrm{\Delta })`$. Expanding Eq. (35) in $`Q`$ yields the following effective masses for the braggoritons
Fig. 3 Density of states for mixed photonic-medium excitations. Left panel is for the symmetric case ($`\delta =0`$), for $`\alpha =0.35`$ (solid line) and $`\alpha =0.57`$ (dashed line). Right panel is for $`\delta =0.3`$ and $`\alpha =0.35`$. The thin solid vertical lines indicate the band edges.
$`M_{\text{a}_1}`$ $`=`$ $`\left[1{\displaystyle \frac{12\delta }{\sqrt{\left(12\delta \right)^2+16\alpha ^2}}}\right]^1,`$ (37)
$`M_{\text{a}_2}`$ $`=`$ $`\left[1{\displaystyle \frac{1+2\delta }{\sqrt{\left(1+2\delta \right)^2+16\alpha ^2}}}\right]^1.`$ (38)
These masses are plotted in Fig. 2 versus the coupling strength $`\alpha `$. For $`\alpha 0`$, we have $`M_{\text{a}_1}`$, $`M_{\text{a}_2}\mathrm{}`$, reflecting the fact that at $`\alpha =0`$ the braggoritons reduce to dispersionless medium excitations that are not coupled to light. With the light-matter interaction switched on, the braggoriton effective mass rapidly decreases (the width of the in-gap branches increases).
The one-dimensional density of states $`N(\mathrm{\Delta })`$ is given from Eq. (35) by:
$`N(\mathrm{\Delta }){\displaystyle \frac{dQ}{d\mathrm{\Delta }}}=`$ (39)
$`{\displaystyle \frac{\mathrm{\Delta }{\displaystyle \frac{\alpha ^2}{\mathrm{\Delta }\delta }}}{\sqrt{\left(\mathrm{\Delta }{\displaystyle \frac{\alpha ^2}{\mathrm{\Delta }\delta }}\right)^2{\displaystyle \frac{1}{4}}}}}\left(1+{\displaystyle \frac{\alpha ^2}{(\mathrm{\Delta }\delta )^2}}\right),`$ (40)
and shown in Fig. 3 for different values of $`\alpha `$ and detuning, $`\delta `$. The density of states of the braggoritonic branches inside the gap exhibits conventional 1D square-root singularities at the band edges $`\mathrm{\Delta }=\delta `$ and $`\mathrm{\Delta }=\frac{1}{2}\left[\delta \frac{1}{2}\pm \sqrt{(\delta \pm \frac{1}{2})^2+4\alpha ^2}\right]`$.
As mentioned above, the upper and lower Bragg-like photonic branches, $`𝓑_1`$, $`𝓑_2`$, are only slightly affected by the coupling and/or detuning. In particular their effective masses,
$`M_{𝓑_1}`$ $`=`$ $`\left[1+{\displaystyle \frac{1+2\delta }{\sqrt{\left(1+2\delta \right)^2+16\alpha ^2}}}\right]^1,`$ (41)
$`M_{𝓑_2}`$ $`=`$ $`\left[1+{\displaystyle \frac{12\delta }{\sqrt{\left(12\delta \right)^2+16\alpha ^2}}}\right]^1,`$ (42)
change only by a factor of 2 as $`\alpha `$ varies from zero to infinity (see Fig. 2).
## V Intragap localized states
We now turn our attention to the localized photonic states caused by a phase-slip like defect. Note that in the absence of the polarizable medium, a structure with one-dimensional modulation (3) of the dielectric function can be viewed as a distributed feedback resonator first considered by Kogelnik and Shank in 1972. Later it was realized that a phase slip in the modulation
$$\epsilon (z)=\epsilon _0+\delta \epsilon \mathrm{cos}(\sigma z+\varphi (z)),$$
(43)
where
$$\varphi (z)=\{\begin{array}{cc}\varphi _1,\hfill & z<0\hfill \\ \varphi _2,\hfill & z>0\hfill \end{array},$$
(44)
results in a localized state inside the PBG. Within the second quantization formalism of Sec. II, the emergence of such a state can be established as follows. Consider the eigenstate annihilation operators of the Hamiltonian (6)
$$\left(\begin{array}{c}\widehat{\beta }_1\\ \widehat{\beta }_2\end{array}\right)=\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta e^{i\varphi }\\ \mathrm{sin}\theta e^{i\varphi }& \mathrm{cos}\theta \end{array}\right)\left(\begin{array}{c}\widehat{a}_{}\\ \widehat{a}_{}\end{array}\right).$$
(45)
It follows from (45) that the absolute value, $`\lambda `$, of the amplitude ratio of the left and right propagating waves constituting the eigenstates $`\beta _1`$, $`\beta _2`$ is either $`\lambda =\mathrm{tan}\theta `$, or $`\lambda =\mathrm{tan}^1\theta `$. These expressions are actually equivalent with appropriate choice of the sign of square root:
$$\lambda _\pm (\mathrm{\Delta })=2\left(\mathrm{\Delta }\pm \sqrt{\mathrm{\Delta }^2\frac{1}{4}}\right),$$
(46)
where we used the definition (9) of the rotation angle $`\theta `$.
In the presence of a phase slip (44) the continuity condition at $`z=0`$ reads
$$\lambda _{}(\mathrm{\Delta })e^{i\varphi _2}=\lambda _{}^{}(\mathrm{\Delta })e^{i\varphi _1}.$$
(47)
As is well known Eq. (47) has a unique in-gap solution, $`\mathrm{\Delta }^{}`$, for an arbitrary phase discontinuity $`\varphi _1\varphi _2`$,
$$\mathrm{\Delta }^{}=\mathrm{cos}\chi ,$$
(48)
where
$$\chi =\{\begin{array}{cc}\frac{\varphi _1\varphi _2}{2}+\pi ,\hfill & \pi <\varphi _1\varphi _2<0,\hfill \\ \frac{\varphi _1\varphi _2}{2},\hfill & 0<\varphi _1\varphi _2<\pi .\hfill \end{array}$$
(49)
Fig. 4 The frequencies of localized intragap states vs. the phase slip magnitude: D<sub>1</sub>, D<sub>2</sub> are defect levels inside the Bragg gap ($`\alpha =0`$, dashed line); $`D_1`$, $`\stackrel{~}{D}_1`$ are defect levels inside the lower subgap, and $`D_2`$, $`\stackrel{~}{D}_2`$ are defect levels inside the upper subgap ($`\delta =0`$, $`\alpha =0.35`$, solid line). Thin solid and dashed lines represent the band edges of the two forbidden gaps and of conventional Bragg gap (in the absence of coupling), respectively.
Generalization of the above consideration to include the polarizable medium is straightforward. It reduces to the following modification of the parameter $`\lambda `$ in Eq. (47):
$`\lambda _\pm (\mathrm{\Delta },\alpha ,\delta )=`$ (50)
$`2\left(\mathrm{\Delta }{\displaystyle \frac{\alpha ^2}{\mathrm{\Delta }\delta }}\pm \sqrt{\left(\mathrm{\Delta }{\displaystyle \frac{\alpha ^2}{\mathrm{\Delta }\delta }}\right)^2{\displaystyle \frac{1}{4}}}\right),`$ (51)
Then condition (47) yields the gap state solution, $`\mathrm{\Delta }^{}`$
$$\mathrm{\Delta }^{}=\frac{1}{2}\left[\delta +\frac{1}{2}\mathrm{cos}\chi \pm \sqrt{\left(\delta \frac{1}{2}\mathrm{cos}\chi \right)^2+4\alpha ^2}\right],$$
(52)
where $`\chi `$ is defined by Eq. (49). For $`\alpha =0`$ we return to the in-gap state (48). Remarkably, we note that for nonzero coupling parameter $`\alpha `$, when the Bragg gap is divided into two sub-gaps as in Fig. 1, the two values of $`\mathrm{\Delta }^{}`$ determined by Eq. (52) are located in each of the corresponding sub-gaps. In Fig. 4 the frequencies, $`\mathrm{\Delta }^{}`$ of the localized in-gap states are shown versus the magnitude $`\varphi _1\varphi _2`$ of the phase-slip for zero detuning $`\delta =0`$.
## VI Discussion
Here we discuss the connection of the present work to the earlier related studies.
The band structure of a photonic crystal with frequency dependent dielectric function was recently studied numerically in Ref., using the plane-wave method. In this work the photonic crystal was modeled as a two-dimensional array of GaAs rods. Frequency dispersion was introduced through the transverse-longitudinal splitting of the optical phonons. The authors (see also Ref.) observed that numerous branches of the band structure calculated for $`\omega _\text{B}\omega _\text{T}`$ become almost dispersionless at frequencies close to the frequency $`\omega _\text{T}`$.
In the context of the present work, this weakening of dispersion can be understood from Eqs. (37), (38) that describes the effective masses of the braggoritonic branches $`\text{a}_1`$, $`\text{a}_2`$. These masses rapidly increase as the coupling parameter $`\alpha `$ decreases.
Signatures of braggoritonic excitations studied in detail in the present paper can be also found in the numerical calculations of Ref.. In that work, transmission spectra of a photonic crystal identical to that of Ref. were calculated within the transfer-matrix formalism. Two minima in the transmission spectra were found instead of the usual single minimum that is caused by the Bragg diffraction in non-dispersive photonic crystal. In light of the theory developed in the present work, these two minima can be identified with the two forbidden subgaps in the excitation spectrum (Figs. 1 and 3). As we have demonstrated (Fig. 3), in the presence of light-matter coupling there are two spectral regions with zero density of states. Correspondingly, the transmission coefficient within these frequency regions must be low if the sample is sufficiently thick.
A very different realization of periodic polarizable structures was the subject of extensive theoretical studies during the last decade. The structures are multiple quantum wells separated by wide-gap semiconductor barriers. The width, $`d`$, of each barrier was assumed to be close to $`\lambda _0/2`$, where $`\lambda _0`$ is the wavelength corresponding to the intra-well exciton resonance frequency, $`\omega _\text{T}`$. The condition $`d\lambda _0/2`$ implies that the Bragg frequency, $`\omega _\text{B}`$ is close to $`\omega _\text{T}`$. Since the quantum wells with strong frequency dispersion at $`\omega \omega _\text{T}`$ had thickness much smaller than $`d`$, then a real Bragg gap in the structures was lacking. However, under the condition $`\omega _\text{T}\omega _\text{B}`$ the dispersion law of light propagating along the principal axis was shown to have a gap within a frequency range $`\left|\omega \omega _\text{T}\right|=\left(2\mathrm{\Gamma }_0\omega _\text{T}/\pi \right)^{1/2}=\mathrm{\Omega }^{\text{(eff)}}`$. Here $`\mathrm{\Gamma }_0`$ denotes the radiative rate for an exciton in a single well. In other words, $`\mathrm{\Omega }^{\text{(eff)}}`$ plays the role of the “effective” Bragg gap in the structures. Remarkably, a physical picture completely analogous to the multiple-quantum-well structures emerged from consideration of an optical lattice formed by laser-cooled atoms. Correspondingly, the light dispersion relation derived in Ref. has the same form as in Refs. . In conclusion we note that as was recently pointed out, a detuning $`(\omega _\text{B}\omega _\text{T})\mathrm{\Omega }^{\text{(eff)}}`$ gives rise to a band of propagating states within the “effective” Bragg gap of the multiple-quantum-well structures.
Acknowledgments: This work was supported by NSF grant DMR 9732820, the Petroleum Research Fund under grant ACS-PRF #34302-AC6, and the Army Research Office.
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# KUNS-1665hep-th/0005178 Scalar field theories in a Lorentz-invariant three-dimensional noncommutative space-time
## 1 Introduction
Quantum field theories in noncommutative space-times with the noncommutativity $`[x^\mu ,x^\nu ]=i\theta ^{\mu \nu }`$ are interesting in pursuing the new possibilities of quantum field theories in quantum space-times. Remarkably, this kind of noncommutativity between coordinates appears in string theory, for instance, in the toroidal compactification of Matrix Theory and in open string theory - with a B-field background. Thus it may be expected that field theories in this kind of noncommutative space-times are controllable in some way. On the other hand, however, perturbative analyses show that these noncommutative field theories have interesting but unusual behaviours -. Infrared singularities in the correlation functions were shown to appear even for massive theories -. In the case with space-time noncommutativity with $`\theta ^{0i}0`$ -, the causality is violated in interesting ways , and S-matrix does not satisfy unitarity constraints . Thus, before we become able to handle such noncommutative quantum field theories consistently, there still seems to remain a lot to learn about them.
An obvious direction to learn more would be to generalise the class of noncommutative space-time. A constant background of $`\theta ^{\mu \nu }`$ violates Lorentz symmetry in more than two-dimensional space-time. Since Lorentz symmetry is one of the fundamental components in the present theoretical physics, in this paper we consider a Lorentz-invariant three-dimensional space-time with the noncommutativity $`[x^\mu ,x^\nu ]=2ił_P\epsilon ^{\mu \nu \rho }x_\rho `$<sup>1</sup><sup>1</sup>1A noncommutative space-time with this commutation relation is often quoted a quantum sphere. For example, see -. In this paper, the noncommutative space-time is three-dimensional, rather than a two-dimensional subspace in it. See also .. Though this noncommutativity is different from that obtained from the constant background of the two-form field in string theory, a similar kind of noncommutativity between more than two coordinates appears on the boundary string of a membrane in M-theory with a non-vanishing background of the three-form field , and also on a D2-brane in a non-constant two-form field background in string theory . Another motivation comes from that this noncommutative space-time may be regarded as a fuzzy space-time with a Lorentz-invariant space-time uncertainty relation derived from a gedanken experiment in which only the general relativity and quantum mechanics are used -. Concerning the questions whether there is any relationship between the space-time uncertainty relation and quantum gravity or string theory , as well as whether the properties obtained so far for the specific noncommutative quantum field theories are general in other noncommutative space-times, it would be interesting to investigate quantum field theories in that noncommutative three-dimensional space-time. As for the latter question, since the time-coordinate is also noncommutative in the space-time we consider, it should be the most radical case with the above-mentioned violations of causality and unitarity -. In this paper, we shall discuss only scalar field theories in the noncommutative space-time.
The organisation of this paper is as follows. In section 2, we discuss the group theoretical structure of the one-particle Hilbert space of free scalar field theory. In section 3, we define the action and derive the Feynman rules. In section 4, we compute some one-loop diagrams. The non-planar diagrams are shown to be finite and have infrared singularities from the UV/IR mixing. We encounter the feature that the violation of the momentum conservation from the non-planar diagrams does not vanish even in the commutative limit $`l_P0`$. In section 5, we propose a translationally symmetric theory to remedy the defect. In section 6, we discuss a noncommutative analogue of the thermodynamics of free scalar field theory and compare the result with the qualitative estimation given previously in . Section 7 is devoted for the summary and discussions.
## 2 One-particle Hilbert space
In the paper , one of the present authors discussed the momentum space representation of the one-particle Hilbert space of free scalar field theory in a Lorentz-invariant three-dimensional noncommutative space-time. There the noncommutativity of the space-time is motivated by a space-time uncertainty relation derived from a certain gedanken experiment -. In this section we discuss the group theoretical structure of the one-particle Hilbert space.
### 2.1 Construction via $`ISO(2,2)`$ algebra
The three-dimensional noncommutative space-time in is represented by the following $`SO(1,2)`$ Lorentz-invariant commutation relations between the coordinates and momentum operators<sup>2</sup><sup>2</sup>2 We have rescaled the numerical constant of the algebra by a factor of 2 from that in the previous paper .:
$`[\widehat{x}^\mu ,\widehat{x}^\nu ]`$ $`=`$ $`2il_Pϵ^{\mu \nu \rho }\widehat{x}_\rho ,`$ (1)
$`[\widehat{p}^\mu ,\widehat{p}^\nu ]`$ $`=`$ $`0,`$ (2)
$`[\widehat{p}^\mu ,\widehat{x}^\nu ]`$ $`=`$ $`i\eta ^{\mu \nu }\sqrt{1+l_P^2\widehat{p}^2}+il_Pϵ^{\mu \nu \rho }\widehat{p}_\rho ,`$ (3)
where the greek indices run through 0 to 2, and the numerical constant $`l_P`$ should be in the order of Planck length.
Now let us consider the lie algebra of $`ISO(2,2)`$,
$`[\widehat{J}_{mn},\widehat{J}_{kl}]`$ $`=`$ $`i(\eta _{mk}\widehat{J}_{nl}\eta _{ml}\widehat{J}_{nk}\eta _{nk}\widehat{J}_{ml}+\eta _{nl}\widehat{J}_{mk}),`$ (4)
$`[\widehat{J}_{mn},\widehat{p}_k]`$ $`=`$ $`i(\eta _{mk}\widehat{p}_n\eta _{nk}\widehat{p}_m),`$ (5)
$`[\widehat{p}_m,\widehat{p}_n]`$ $`=`$ $`0,`$ (6)
where the roman indices run through $`1`$ to 2, and the signature is given by $`\eta _{mn}=(,,+,+)`$. By identifying
$`\widehat{x}_\mu `$ $`=`$ $`l_P(\widehat{J}_{1,\mu }{\displaystyle \frac{1}{2}}ϵ_\mu {}_{}{}^{\alpha \beta }\widehat{J}_{\alpha \beta }^{}),`$ (7)
$`\widehat{p}_\mu `$ $`=`$ $`\widehat{p}_{m=\mu },`$ (8)
and imposing a constraint
$$1+l_P^2\widehat{p}^m\widehat{p}_m=0,$$
(9)
we can easily show that the commutation relations (1) can be derived from (4). The momentum square $`\widehat{p}^m\widehat{p}_m`$ is one of the Casimir operators of the algebra (4), and hence we may consistently impose the constraint (9). The remaining three independent generators in (4) are now the $`SO(1,2)`$ Lorentz generators of the noncommutative space-time (1).
The representation of the algebra (1) is obtained from that of $`ISO(2,2)`$, which is given by $`\widehat{J}_{mn}=i(p_m\frac{}{p^n}p_n\frac{}{p^m})`$ in the momentum space representation. In order to impose the constraint (9), the following “polar” coordinate is convenient:
$`p_1`$ $`=`$ $`r\mathrm{cosh}\chi \mathrm{cos}\theta ,`$ (10)
$`p_0`$ $`=`$ $`r\mathrm{cosh}\chi \mathrm{sin}\theta ,`$ (11)
$`p_1`$ $`=`$ $`r\mathrm{sinh}\chi \mathrm{cos}\varphi ,`$ (12)
$`p_2`$ $`=`$ $`r\mathrm{sinh}\chi \mathrm{sin}\varphi .`$ (13)
This coordinate is only valid in the neighbourhood of the hyperboloid $`1+l_P^2p^mp_m=0`$, but this is enough for our purposes. Since $`\widehat{J}_{mn}p^lp_l=0`$ and $`\widehat{J}_{mn}`$ do not contain $`r`$-derivative, we can restrict the representation space to the functions on the hyperboloid, i.e. the functions depending only on $`\chi ,\theta ,\varphi `$. Thus the natural inner-product that makes $`\widehat{J}_{mn},\widehat{p}_m`$ hermite in the restricted representation space is given by
$`\mathrm{\Phi }_1|\mathrm{\Phi }_2=2{\displaystyle d^4p\delta (1+l_P^2p^mp_m)\mathrm{\Phi }_1^{}(p)\mathrm{\Phi }_2(p)}.`$ (14)
After integrating over $`p_1`$, this inner product is identical to that obtained previously in .
The constraint (9) shows that the mass square $`p^\mu p_\mu `$ in the noncommutative space-time should have an upper bound $`p^\mu p_\mu 1/l_P^2`$, and also that, coming from the choices of the sign of $`p_1`$, there exists two-fold degeneracy in the momentum space of the noncommutative space-time. These features agree with the results obtained in .
### 2.2 $`SL(2,R)`$ structure of the momentum space
In the previous subsection, we have shown that the momentum space is the hyperboloid, $`1+l_P^2p^mp_m=0`$. This hyperboloid can be mapped to the group manifold $`SL(2,R)`$ as follows. Let us define the matrices
$`𝕩^0`$ $`=`$ $`l_P\sigma ^2=l_P\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),`$ (15)
$`𝕩^1`$ $`=`$ $`l_Pi\sigma ^3=l_P\left(\begin{array}{cc}i& 0\\ 0& i\end{array}\right),`$ (16)
$`𝕩^2`$ $`=`$ $`l_Pi\sigma ^1=l_P\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),`$ (17)
which satisfy
$`𝕩^\mu 𝕩^\nu =l_P^2\eta ^{\mu \nu }+il_Pϵ^{\mu \nu \lambda }𝕩_\lambda .`$ (18)
Then the map between the hyperboloid and the group elements of $`SL(2,R)`$ is defined by
$`g=l_P\left(\begin{array}{cc}p_1p_1& p_0p_2\\ p_0p_2& p_1+p_1\end{array}\right)=l_Pp_1+i𝕩^\mu p_\mu =e^{ik𝕩}.`$ (19)
By a direct computation, we can verify $`p_\mu =k_\mu \mathrm{sinh}(l_P\sqrt{k^2})/l_P\sqrt{k^2}`$. As discussed in , this relation is identical to the relation between $`k_\mu `$ and the $`\widehat{p}_\mu `$-eigenvalue of the state $`e^{ik\widehat{x}}|0`$, where $`|0`$ denotes the momentum-zero eigen-state with $`p_1=1`$. Thus the one-particle Hilbert space can be parameterised by the group manifold $`SL(2,R)`$:
$`|p(g)=e^{ik(g)\widehat{x}}|0,`$ (20)
where $`p_m(g)`$ and $`k_\mu (g)`$ are defined by (19). This parameterisation with $`SL(2,R)`$ elements is superior to the parameterisation by the space-time momentum $`p_\mu `$, because distinct choices of the signs of $`p_1`$ correspond to distinct group elements and we do not need to worry about the two-fold degeneracy in the momentum space representation.
For later convenience, we collect some useful formulae in the followings. From (18) and (19), we can show that the group multiplication can be expressed in terms of $`p_m`$ by
$`p(gh)_1`$ $`=`$ $`l_P(p(g)_1p(h)_1+p(g)^\alpha p(h)_\alpha ),`$ (21)
$`p(gh)_\mu `$ $`=`$ $`l_P(p(g)_1p(h)_\mu +p(g)_\mu p(h)_1ϵ_\mu {}_{}{}^{\alpha \beta }p(g)_\alpha p(h)_\beta ).`$ (22)
Note that $`\mathrm{tr}(g)=2p(g)_1=\pm 2\sqrt{l_P^2+p_\mu (g)p^\mu (g)}`$. Since an adjoint action keeps the trace invariant, it keeps $`p_\mu p^\mu `$ invariant. Hence the adjoint action of $`SL(2,R)`$ corresponds to the $`SO(1,2)`$ Lorentz transformation of the space-time. Under the adjoint action, $`p(h)_m`$ transforms as
$`p(g^1hg)_1`$ $`=`$ $`p(h)_1,`$ (23)
$`p(g^1hg)_\mu `$ $`=`$ $`p(h)_\mu (1+2a^\alpha a_\alpha )2p(h)^\alpha a_\alpha a_\mu 2ϵ_\mu {}_{}{}^{\alpha \beta }p(h)_\alpha a_\beta a_1,`$ (24)
where $`a_m=l_Pp(g)_m`$.
After a short computation, we can verify that the inner product defined in (14) is just the integration over the $`SL(2,R)`$ group manifold with the invariant measure,
$`d\mu (g)={\displaystyle \frac{1}{\mathrm{2\; 3}!}}\mathrm{tr}(g^1dg)^3=2d^4p\delta (1+l_P^2p^mp_m).`$ (25)
Hence we have
$`d\mu (g)=d\mu (gh)=d\mu (hg).`$ (26)
The above formulas will be used in computing one-loop diagrams in section 4.
## 3 Action and Feynman rules
In this section, we construct the actions of interacting scalar field theories in the noncommutative space-time and derive the Feynman rules<sup>3</sup><sup>3</sup>3A similar derivation of action and Feynman rules was carried out for a deformed Minkowski space in ..
A scalar field in the noncommutative space-time is defined by associating momentum space wave functions $`\mathrm{\Phi }(g)`$ to “vertex operators” as follows:
$`\widehat{\mathrm{\Phi }}`$ $`=`$ $`{\displaystyle 𝑑\mu (g)\mathrm{\Phi }(g)e^{ik(g)\widehat{x}}},`$ (27)
$`|\mathrm{\Phi }`$ $`=`$ $`{\displaystyle 𝑑\mu (g)\mathrm{\Phi }(g)e^{ik(g)\widehat{x}}|0}={\displaystyle 𝑑\mu (g)\mathrm{\Phi }(g)|p(g)},`$ (28)
where $`p(g)`$ and $`k(g)`$ are defined by (19). We impose the reality condition on the field $`\widehat{\mathrm{\Phi }}`$, i.e. $`\widehat{\mathrm{\Phi }}^{}=\widehat{\mathrm{\Phi }}`$. We can define the product of the vertex operators by Hausdorff formula, following the line of . Since the Hausdorff formula is nothing but the group multiplication, we obtain
$`e^{ik(g_1)\widehat{x}}e^{ik(g_2)\widehat{x}}=e^{ik(g_1g_2)\widehat{x}}.`$ (29)
Making use of this $``$-product, we can construct interaction terms. For example, the action of noncommutative $`\varphi ^3`$-theory can be defined in the following way:
$$S[\mathrm{\Phi }]=\frac{1}{2}0|\widehat{\mathrm{\Phi }}(\widehat{p}^2+m^2)\widehat{\mathrm{\Phi }}|0+\frac{\lambda }{3}0|\widehat{\mathrm{\Phi }}\widehat{\mathrm{\Phi }}\widehat{\mathrm{\Phi }}|0,$$
(30)
where $`\widehat{p}^2=\widehat{p}_\mu \widehat{p}^\mu `$.
For practical computation, it is convenient to rewrite this action in the momentum space representation with $`SL(2,R)`$ elements, and, for this, we need to evaluate $`0|e^{ik(g)\widehat{x}}|0=0|p(g)`$. From the definition of the inner product, we can show that $`0|p(g)`$ is a $`\delta `$-function with respect to the invariant measure $`d\mu (g)`$, which has a support at the unit element of $`SL(2,R)`$:
$$0|\mathrm{\Phi }𝑑\mu (g)0|p(g)p(g)|\mathrm{\Phi }\mathrm{\Phi }(g=1)=𝑑\mu (g)0|p(g)\mathrm{\Phi }(g).$$
(31)
Henceforth, we denote this $`\delta `$-function as $`\delta (g)`$.
Now we can write down the momentum representation of the action as
$`S[\mathrm{\Phi }]`$ $`=`$ $`{\displaystyle \frac{1}{2}}0|\widehat{\mathrm{\Phi }}(\widehat{p}^2+m^2)\widehat{\mathrm{\Phi }}|0+{\displaystyle \frac{\lambda }{3}}0|\widehat{\mathrm{\Phi }}\widehat{\mathrm{\Phi }}\widehat{\mathrm{\Phi }}|0`$ (32)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle 𝑑\mu (g)\mathrm{\Phi }(g^1)(p^2(g)+m^2)\mathrm{\Phi }(g)}`$ (34)
$`+{\displaystyle \frac{\lambda }{3}}{\displaystyle \underset{i=1}{\overset{3}{}}d\mu (g_i)\delta (g_1g_2g_3)\mathrm{\Phi }(g_1)\mathrm{\Phi }(g_2)\mathrm{\Phi }(g_3)}.`$
Assuming an appropriate path integral measure, we can perform perturbative computation.
$`Z[J]`$ $`=`$ $`{\displaystyle [D\mathrm{\Phi }]\mathrm{exp}\left(iS_{\mathrm{free}}[\mathrm{\Phi }]+iS_{\mathrm{int}}[\mathrm{\Phi }]+i𝑑\mu (g)J(g^1)\mathrm{\Phi }(g)\right)}`$ (35)
$`=`$ $`\mathrm{exp}\left(iS_{\mathrm{int}}\left[{\displaystyle \frac{\delta }{i\delta J}}\right]\right)\mathrm{exp}\left({\displaystyle \frac{1}{2}}{\displaystyle 𝑑\mu (g)iJ(g^1)\frac{i}{p^2(g)+m^2}iJ(g)}\right)`$ (37)
$`\times {\displaystyle }[D\mathrm{\Phi }]\mathrm{exp}\left({\displaystyle \frac{i}{2}}{\displaystyle }d\mu (g)(\mathrm{\Phi }{\displaystyle \frac{1}{p^2+m^2}}J)(p^2+m^2)(\mathrm{\Phi }{\displaystyle \frac{1}{p^2+m^2}}J)\right)`$
We may assume the last Gaussian path integral just gives a constant. Thus, generalising to arbitrary $`\varphi ^n`$ interactions, we obtain the following Feynman rules,
$`\mathrm{propagator}:`$ $`{\displaystyle 𝑑\mu (g)\frac{i}{p^2(g)+m^2}},`$ (38)
$`n\mathrm{vertex}:`$ $`i\lambda _n\delta (g_1\mathrm{}g_n)`$ (39)
Note that, since momentums are $`SL(2,R)`$ elements, we have noncommutativity at the vertex. In the following section, we compute some one-loop diagrams, using these rules.
## 4 One-loop computation
In this section we shall compute the one-loop diagrams of the two-point functions from $`\varphi ^3`$ and $`\varphi ^4`$-interactions defined in the previous section. We will show that the non-planar one-loop diagrams are finite and have infrared singularities from the UV/IR mixing -. We will also find that those diagrams cause a problem concerning the conservation of momentums. In this section, we set $`l_P=1`$ for simplicity.
### 4.1 The planar diagram of the two-point function from $`\varphi ^4`$-interaction
In the first place, we shall compute the simplest graph to show our computational strategy: the one-loop planar diagram of the two-point function from $`\varphi ^4`$-interaction as in fig.(1).
From the Feynman rules in the preceding section, the amplitude of this diagram is given by
$`\mathrm{\Gamma }_P^{(2)}`$ $`=`$ $`{\displaystyle 𝑑\mu (g)i\delta (g^1h_2h_1g)\frac{i}{p^2(g)+m^2}}=\delta (h_2h_1){\displaystyle 𝑑\mu (g)\frac{1}{p^2(g)+m^2}}.`$ (40)
The polar coordinate (10) is convenient for the explicit evaluation of this integral, and we further perform the change of the variable, $`x=\mathrm{cosh}^2\chi `$. With this parameterisation, we obtain
$`p(g)`$ $`=`$ $`(x^{\frac{1}{2}}\mathrm{sin}\theta ,(1+x)^{\frac{1}{2}}\mathrm{cos}\varphi ,(1+x)^{\frac{1}{2}}\mathrm{sin}\varphi ),`$ (41)
$`p^2(g)+m^2`$ $`=`$ $`(p(g)_1)^21+m^2=x\mathrm{cos}^2\theta \mathrm{cos}^2\mu ,`$ (42)
$`d\mu (g)`$ $`=`$ $`{\displaystyle _0^{2\pi }}𝑑\varphi {\displaystyle _0^{2\pi }}𝑑\theta {\displaystyle _1^{\mathrm{}}}{\displaystyle \frac{dx}{2}},`$ (43)
where $`\mathrm{sin}\mu =m(0\mu \pi /2)`$.
Since this integration turns out to be divergent, we introduce a momentum cut-off $`\mathrm{\Lambda }`$ in the $`x`$-integration.
$`\mathrm{\Gamma }_P^{(2)}`$ $`=`$ $`\delta (h_2h_1){\displaystyle _0^{2\pi }}𝑑\varphi {\displaystyle _0^{2\pi }}𝑑\theta {\displaystyle _1^{\mathrm{\Lambda }^2}}{\displaystyle \frac{dx}{2}}{\displaystyle \frac{1}{x\mathrm{cos}^2\theta \mathrm{cos}^2\mu }}`$ (44)
$`=`$ $`2\pi ^2\delta (h_2h_1){\displaystyle _1^{\mathrm{\Lambda }^2}}{\displaystyle \frac{dx}{x}}{\displaystyle _0^{2\pi }}{\displaystyle \frac{id\theta }{2\pi i}}{\displaystyle \frac{1}{\mathrm{cos}^2\theta \mathrm{cos}^2\mu _x}}`$ (45)
$`=`$ $`2\pi ^2\delta (h_2h_1){\displaystyle _1^{\mathrm{\Lambda }^2}}{\displaystyle \frac{dx}{x}}{\displaystyle _0^{2\pi }}{\displaystyle \frac{id\theta }{2\pi i}}{\displaystyle \frac{1}{\mathrm{sin}(\mu _x+\theta )\mathrm{sin}(\mu _x\theta )}},`$ (46)
where we have defined $`\mathrm{cos}\mu _x=\mathrm{cos}\mu /x`$.
By the change of the integration variable $`z=e^{i\theta }`$, the integration over $`\theta `$ is now a contour integration on a unit circle in the $`z`$-plane. There are poles at $`z=\pm e^{i\mu _x},\pm e^{i\mu _x}`$. Since they are on the unit circle, we adopt $`(m^2iϵ)`$-prescription to make the integral well-defined. Then we just pick up the residues at $`z=\pm e^{i\mu _x}`$ as shown in fig.(2). Thus the result of the integration is given by
$`\mathrm{\Gamma }_P^{(2)}`$ $`=`$ $`2\pi ^2i\delta (h_2h_1){\displaystyle _1^{\mathrm{\Lambda }^2}}{\displaystyle \frac{dx}{x}}{\displaystyle \frac{x}{\sqrt{x(1m^2)}\sqrt{1m^2}}}`$ (47)
$`=`$ $`4\pi ^2i\delta (h_2h_1)\left(\sqrt{{\displaystyle \frac{m^2}{1m^2}}}\sqrt{{\displaystyle \frac{\mathrm{\Lambda }^2}{1m^2}}}\right),`$ (48)
where we have neglected the terms that vanish as $`\mathrm{\Lambda }\mathrm{}`$. We need to renormalize the divergence of the second term by a mass counter term.
### 4.2 The non-planar diagram of the two-point function from $`\varphi ^4`$-interaction
Next we consider the non-planar diagram of fig.(3). The amplitude is given by
$`\mathrm{\Gamma }_{NP}^{(2)}={\displaystyle 𝑑\mu (g)i\delta (h_2g^1h_1g)\frac{i}{p^2(g)+m^2}}.`$ (49)
The integration over $`g`$ is performed for the solutions to $`h_2g^1h_1g=1`$. Thus this integral is non-vanishing only if $`h_1`$ and $`h_2`$ belong to the same conjugacy class, so both $`p(h_1)_\mu `$ and $`p(h_2)_\mu `$ must be time-like or space-like, simultaneously. We shall evaluate this integral when both the momentums are time-like. Without loss of generality, we may assume $`p(h_1)_\mu `$ is in the time direction, and $`p(h_2)_\mu `$ is represented by a vector obtained by boosting a vector $`p(h_2^{})_\mu `$ which is in the time direction:
$`p(h_2^{})_\mu `$ $`=`$ $`(p(h_2^{})_0,0,0),`$ (50)
$`p(h_2)_\mu `$ $`=`$ $`p(g_2^1h_2^{}g_2)`$ (51)
$`=`$ $`(p(h_2^{})_0\mathrm{cosh2}\chi _2,p(h_2^{})_0\mathrm{sinh2}\chi _2\mathrm{sin}\varphi _2,p(h_2^{})_0\mathrm{sinh2}\chi _2\mathrm{cos}\varphi _2).`$ (52)
We can rewrite
$`\mathrm{\Gamma }_{NP}^{(2)}`$ $`=`$ $`{\displaystyle 𝑑\mu (g)\delta (h_2^{}g_2g^1h_1gg_2^1)\frac{1}{p^2(g)+m^2}}`$ (53)
$`=`$ $`{\displaystyle 𝑑\mu (g)\delta (h_2^{}g^1h_1g)\frac{1}{p^2(gg_2)+m^2}},`$ (54)
where we have changed the integration variable $`g`$ to $`gg_2`$. An adjoint action is simplified when acting on a vector in the time-direction:
$`p(g^1h_1g)_0`$ $`=`$ $`p(h_1)_0(1+2(p_1^2+p_2^2)),`$ (55)
$`p(g^1h_1g)_i`$ $`=`$ $`2p(h_1)_0(p_0p_i+ϵ_i{}_{}{}^{j}p_{1}^{}p_j),`$ (56)
where we have abbreviated $`p_m=p(g)_m`$. We can show the following formula <sup>4</sup><sup>4</sup>4Abbreviating $`p(g)=p`$ and $`p(h)=q`$, $`{\displaystyle 𝑑\mu (g)\delta (gh)\varphi (g)}`$ $`=`$ $`2{\displaystyle d^4p\delta (1+p^mp_m)|p_1|\theta (p_1q_1)\delta ^{(3)}(p+q)\varphi (p)}`$ (57) $`=`$ $`2{\displaystyle 𝑑p_1\delta (q_1^2p_1^2)|p_1|\theta (p_1q_1)\varphi (p_1,q_\mu )}`$ (58) $`=`$ $`{\displaystyle 𝑑p_1\{\delta (q_1p_1)+\delta (q_1+p_1)\}\theta (p_1q_1)\varphi (p_1,q_\mu )}`$ $`=`$ $`\varphi (q_1,q_\mu )=\varphi (h^1),`$ where the step function $`\theta (p_1q_1)`$ is needed to discard the contribution from $`g=h`$.:
$`\delta (gh)=|p(h)_1|\theta (p(g)_1p(h)_1)\delta ^{(3)}(p(g)+p(h)).`$ (60)
Making use of it, we obtain
$`\mathrm{\Gamma }_{NP}^2`$ $`=`$ $`|p(h_1)_1|2{\displaystyle d^4p\delta (1+p^mp_m)}`$ (62)
$`\times \theta (p(g^1h_1g)_1p(h_2^{})_1)\delta ^{(3)}(p(g^1h_1g)+p(h_2^{})){\displaystyle \frac{1}{p^2(gg_2)+m^2}}`$
$`=`$ $`|p(h_1)_1|2{\displaystyle d^4p\delta (1+p^mp_m)}`$ (65)
$`\times \theta (p(h_1)_1p(h_2^{})_1)\delta (p(h_1)_0(1+2p_1^2+2p_2^2)+p(h_2^{})_0)`$
$`\times \delta (2p(h_1)_0(p_0p_1p_1p_2))\delta (2p(h_1)_0(p_0p_2+p_1p_1)){\displaystyle \frac{1}{p^2(gg_2)+m^2}}.`$
The integrations over $`p_1`$ and $`p_2`$ can be trivially performed by the 2nd and 3rd $`\delta `$-functions, picking up the values $`p_1=p_2=0`$ and giving the Jacobian factor $`1/4(p(h_1)_0)^2(p_1^2+p_0^2)=1/(4p^2(h_1))(p_1^2+p_0^2)`$:
$`\mathrm{\Gamma }_{NP}^{(2)}`$ $`=`$ $`{\displaystyle \frac{|p(h_1)_1|}{4p^2(h_1)}}\theta (p(h_1)_1p(h_2^{})_1)\delta (p(h_1)_0+p(h_2^{})_0)`$ (67)
$`\times 2{\displaystyle 𝑑p_1𝑑p_0\delta (1p_1^2p_0^2)\frac{1}{p^2(gg_2)+m^2}}|_{p_1=p_2=0}.`$
Thus the integration reduces to the contour integral on a unit circle in $`(p_1,p_0)`$-plane. We evaluate this integral by $`(m^2iϵ)`$-prescription:
$`\mathrm{\Gamma }_{NP}^{(2)}`$ $`=`$ $`{\displaystyle \frac{|p(h_1)_1|}{4p^2(h_1)}}\theta (p(h_1)_1p(h_2^{})_1)\delta (p(h_1)_0+p(h_2^{})_0){\displaystyle _0^{2\pi }}𝑑\theta {\displaystyle \frac{1}{\mathrm{cosh}^2\chi _2\mathrm{cos}^2(\theta +\frac{\varphi _2}{2})+\mathrm{cos}^2\mu }}`$ (68)
$`=`$ $`{\displaystyle \frac{|p(h_1)_1|}{4p^2(h_1)}}\theta (p(h_1)_1p(h_2^{})_1)\delta (p(h_1)_0+p(h_2^{})_0){\displaystyle \frac{2\pi i}{\sqrt{(\mathrm{sinh}^2\chi _2+\mathrm{sin}^2\mu )\mathrm{cos}^2\mu }}}`$ (69)
Taking into account the relation $`\mathrm{sinh}^2\chi _2=\frac{1}{2}\left(\frac{p(h_1)p(h_2)}{p(h_1)^2}+1\right)`$ and an identity,
$`\theta (p_1q_1)\delta (p_0+q_0)={\displaystyle \frac{\sqrt{p^2}}{|p_1|}}\theta (p_0q_0)\delta (p_1q_1),`$ (70)
we finally obtain the following Lorentz invariant result:
$$i\mathrm{\Gamma }_{NP}^{(2)}=\frac{\pi }{\sqrt{2(1m^2)}}\frac{\theta (p(h_1)_0p(h_2)_0)\delta (p(h_1)_1p(h_2)_1)}{\sqrt{(12m^2)p^2(h_1)+p(h_1)p(h_2)}}.$$
(72)
This result shows that the non-planar diagram has an infrared singularity coming from the UV/IR mixing - irrespective of massive theory. We can also see that the momentum is not conserved while the momentum square is. One might expect that the momentum would be conserved in the commutative limit $`l_p0`$, since, in this limit, the commutation relations (1) have the translational symmetry $`\widehat{x}^\mu \widehat{x}^\mu +v^\mu `$ with a $`c`$-number vector $`v^\mu `$. However this is not true. In the commutative limit, $`m`$ approaches to zero, and the factor $`1/\sqrt{p^2(h_1)+p(h_1)p(h_2)}`$ in the amplitude $`\mathrm{\Gamma }_{NP}^{(2)}`$ has an infinite peak at $`p_\mu (h_1)+p_\mu (h_2)=0`$. Although the momentums are conserved at this peak, we can easily show that the factor $`1/\sqrt{p^2(h_1)+p(h_1)p(h_2)}`$ is not a $`\delta `$-function under the measure $`d\mu (g)`$. Thus, in this theory, the momentum conservation is violated substantially even in the low-energy limit. The problem of the violation of momentum conservation was observed in another noncommutative space-time in . We shall discuss this problem further in section 5.
### 4.3 The two-point function from $`\varphi ^3`$-interaction
The one-loop contributions to the two-point function from $`\varphi ^3`$-interaction can be computed in the same way as in $`\varphi ^4`$-interaction. There are two diagrams as in figs.(4).
The contribution from the diagram (4)-(a) is given by
$`\mathrm{\Gamma }_P^{(2)}`$ $`=`$ $`{\displaystyle 𝑑\mu (g_1)𝑑\mu (g_2)i\delta (g_2^1h_1g_1)i\delta (g_1^1h_2g_2)\frac{i}{p^2(g_1)+m^2}\frac{i}{p^2(g_2)+m^2}}`$ (73)
$`=`$ $`{\displaystyle 𝑑\mu (g)\delta (h_2h_1)\frac{1}{p^2(g)+m^2}\frac{1}{p^2(h_1g)+m^2}}.`$ (74)
In this case the three-momentum $`p_\mu `$ conserves, so we may assume, without loss of generality, $`p(h_1)_\mu =p(h_2)_\mu =(\mathrm{sin}\theta _1,0,0)`$ for time-like three-momenta. Parameterising the loop momentum by the polar coordinate as in eq.(41), we have $`p^2(h_1g)=\mathrm{cosh}^2\chi \mathrm{cos}^2(\theta +\theta _1)1`$. Thus
$`\mathrm{\Gamma }_P^{(2)}`$ $`=`$ $`\delta (h_2h_1){\displaystyle _0^{2\pi }}𝑑\varphi {\displaystyle _0^{2\pi }}𝑑\theta {\displaystyle _1^{\mathrm{}}}{\displaystyle \frac{dx}{2}}{\displaystyle \frac{1}{x\mathrm{cos}^2\theta \mathrm{cos}^2\mu }}{\displaystyle \frac{1}{x\mathrm{cos}^2(\theta +\theta _1)\mathrm{cos}^2\mu }}`$ (75)
$`=`$ $`{\displaystyle \frac{\pi }{1m^2}}\delta (h_2h_1){\displaystyle \frac{1}{\mathrm{sin}\theta _1}}I(\theta _1),`$ (76)
where
$`I(\theta _1)={\displaystyle _0^{2\pi }}𝑑\theta {\displaystyle \frac{1}{\mathrm{sin}2\theta }}\mathrm{log}\left({\displaystyle \frac{\mathrm{cos}^2(\theta \frac{1}{2}\theta _1)\mathrm{sin}(\theta +\frac{1}{2}\theta _1+\mu )\mathrm{sin}(\theta +\frac{1}{2}\theta _1\mu )}{\mathrm{cos}^2(\theta +\frac{1}{2}\theta _1)\mathrm{sin}(\theta \frac{1}{2}\theta _1+\mu )\mathrm{sin}(\theta \frac{1}{2}\theta _1\mu )}}\right).`$ (77)
Now let us compute the first derivative $`I^{}(\theta _1)`$ with respect to $`\theta _1`$. The $`\theta `$-integral can be rewritten as a contour integral on a unit circle in the $`z=e^{i\theta }`$ plane. There appears several poles on the contour. These poles originate from the poles of the propagators in (75) at $`\mathrm{cos}(\theta \pm \frac{1}{2}\theta _1)=\pm \sqrt{(1m^2)/x}`$, where we have shifted $`\theta `$ by $`\frac{1}{2}\theta _1`$ from that of (75). These poles are treated by the $`(m^2iϵ)`$-prescription as the preceding subsections. Then, since the poles of $`I^{}(\theta _1)`$ at $`\mathrm{cos}(\theta \pm \frac{1}{2}\theta _1)=0`$ come from the $`x=\mathrm{}`$ contributions of the integral (75), we see that these poles are in fact degenerate pairs of poles in the inside and those in the outside of the unit circle. By an explicit computation, it turns out that the contributions from those poles cancel. Out of the poles of $`I^{}(\theta _1)`$ at $`\mathrm{sin}(\theta \pm \frac{1}{2}\theta _1\pm \mu )=0`$, which come from the $`x=1`$ contributions of the integral (75), only the poles at $`\mathrm{sin}(\theta \pm \frac{1}{2}\theta _1+\mu )`$ are inside. Evaluating the residues, we obtain
$`I^{}(\theta _1)`$ $`=`$ $`2\pi i\left({\displaystyle \frac{1}{\mathrm{sin}(\theta _12\mu )}}{\displaystyle \frac{1}{\mathrm{sin}(\theta _1+2\mu )}}\right).`$ (78)
Using $`1/\mathrm{sin}x=(\mathrm{log}\mathrm{tan}(x/2))^{}`$, we obtain
$`I(\theta _1)`$ $`=`$ $`2\pi i\mathrm{log}\left({\displaystyle \frac{\mathrm{tan}(\mu \frac{1}{2}\theta _1)}{\mathrm{tan}(\mu +\frac{1}{2}\theta _1)}}\right),`$ (79)
where the integration constant is determined by $`I(0)=0`$. Thus we reach the result
$`i\mathrm{\Gamma }_P^{(2)}`$ $`=`$ $`{\displaystyle \frac{2\pi ^2}{1m^2}}\delta (h_2h_1){\displaystyle \frac{1}{\mathrm{sin}\theta _1}}\mathrm{log}\left({\displaystyle \frac{\mathrm{tan}(\mu \frac{1}{2}\theta _1)}{\mathrm{tan}(\mu +\frac{1}{2}\theta _1)}}\right),`$ (80)
where $`\mu `$ should be understood as $`\mu iϵ`$. This expression has the branch cuts which represent the continuous spectrum of two-body states as in the fig.(5).
As in $`\varphi ^4`$-interaction, the contribution from fig.(4)-(b) also violates the conservation of the three-momentum. We have computed the amplitude in the case that the external momenta are time-like. The result is
$`i\mathrm{\Gamma }_{NP}^{(2)}`$ $`=`$ $`{\displaystyle \frac{\pi }{\sqrt{1m^2}}}\theta (p_0p_0^{})\delta (p_1p_1^{})`$ (82)
$`\times {\displaystyle \frac{1}{\sqrt{\frac{1}{4}(p+p^{})^2m^2p^2}}}{\displaystyle \frac{2(1m^2)p^2\frac{1}{4}(p+p^{})^2}{(\frac{1}{4}(p+p^{})^2)^24(1m^2)(\frac{1}{4}(p+p^{})^2m^2p^2)}}`$
where $`p_m=p(h_1)_m,p_m^{}=p(h_2)_m`$. This amplitude has the same problems discussed for the non-planar diagram from $`\varphi ^4`$-interaction.
### 4.4 Some comments on higher loops and renormalization
In the preceding subsections, we see that the one-loop contributions to the two-point function conserve $`p^2`$ while they do not conserve $`p_\mu `$. In higher loops, the $`p^2`$ is not conserved either. To see this, let us consider the diagram of fig.(4.4) as an example. The amplitude is proportional to
$$𝑑\mu (g_1)𝑑\mu (g_2)𝑑\mu (g_3)\delta (g_1g_2^1h_2h_1)\delta (g_1^1g_3g_2g_3)𝑑\mu (g_1)𝑑\mu (g_3)\delta (g_1g_3^1g_1^1g_3h_2h_1).$$
(83)
Thus $`h_1`$ and $`h_2`$ do not need to belong to the same conjugacy class, and therefore $`p^2`$ is not conserved.
Figure 6: A $`p^2`$-violating graph
Figure 7: A divergent tadpole graph
Figure 8: A non-planar tadpole graph
As for the renormalization at one-loop order, the only divergent graphs are planar tadpole graphs such as fig.(4.4). These can be calculated simply by replacing the external momenta of the result of subsection 4.1. Then their external momentum dependence becomes $`\delta (h_1\mathrm{}h_n)`$, and these divergences can be absorbed by $`n`$-point counter terms and renormalizable. No non-planar graphs such as fig.(4.4) diverge. This is because those graphs can be evaluated only by replacing the external momenta $`(h_1,h_2)`$ of the result of subsection 4.2 with those of appropriate channel $`(h_1\mathrm{}h_k,h_{k+1}\mathrm{}h_n)`$. Notice that the group manifold $`\delta `$-function at the vertex reduces the dimensions of loop-momentum integration and has an effect similar to the Moyal phase factor which serves as a damping factor. For example, in the non-planar graphs such as figs.(3) and (4.4), the loop momentum runs only over the one-dimensional subgroup of $`SL(2,R)`$ determined from the external momenta.
Beyond the one-loop order, we do not know whether the noncommutative scalar field theory we are discussing is renormalizable or not. In the case of scalar field theory with the Moyal type noncommutativity, renormalizability is shown up to two-loop and also a convergence theorem of Feynman integral has recently been well understood . For the noncommutative YM and QED cases, renormalizability is shown up to one-loop -, . In order to investigate the issue in our field theory from perturbative analysis, we need first a systematic evaluation of Feynman integral as was done in . In the present work, this is out of our scope.
## 5 Translational symmetry
In the previous section, we have shown that the contributions from the non-planar diagram violates the momentum conservation. This violation comes from the fact that the commutation relations of the coordinate operators (1) do not respect the naive translational invariance, $`\widehat{x}^\mu \widehat{x}^\mu +v^\mu `$ with a c-number vector $`v^\mu `$. Since the translational invariance is recovered in the commutation relations in the commutative limit $`l_P0`$, one naively expects that the violation should vanish in this limit. But this naive expectation is not true as we saw for the non-planar one-loop results in the preceding section. Since this violation exists substantially in the low energy limit while the momentum conservation is one of the fundamental components in the present theoretical physics, we need an exact symmetry in the noncommutative space-time, which corresponds to the translational symmetry in the commutative limit.
To introduce such symmetry, let us define a field
$$\widehat{\mathrm{\Psi }}=\underset{n=0}{}𝑑\mu (g)\mathrm{\Phi }(g)_{\mu _1\mathrm{}\mu _n}\widehat{p}^{\mu _1}\mathrm{}\widehat{p}^{\mu _n}\mathrm{exp}(ik_\mu (g)\widehat{x}^\mu ),$$
(84)
where $`\mathrm{\Phi }(g)_{\mu _1\mathrm{}\mu _n}`$ are c-number tensor fields. This can be regarded as a non-local generalisation of the scalar field defined in (27).
Let us define a unitary operator
$$U(v)=\mathrm{exp}(iv_\mu \widehat{p}^\mu ).$$
(85)
The $`p_\mu =0`$ eigen-state $`|0`$ is invariant under this unitary operator
$$U|0=|0.$$
(86)
We define the translational transformation of the field $`\widehat{\mathrm{\Psi }}`$ by
$$\widehat{\mathrm{\Psi }}U^{}\widehat{\mathrm{\Psi }}U.$$
(87)
Then this translational transformation is a symmetry of the following type of action
$$S=0|(\widehat{\mathrm{\Psi }}^{})^mf(\widehat{p})\widehat{\mathrm{\Psi }}^n|0.$$
(88)
In usual commutative cases, the translational transformation just generates a phase shift $`\mathrm{exp}(ik_\mu v^\mu )`$, and the tensor fields $`\widehat{\mathrm{\Phi }}(g)_{\mu _1\mathrm{}\mu _n}`$ do not mix. Hence we may truncate the fields only to the scalar sector. However, in our noncommutative case, since the commutation relations are given by (1), the tensor fields mix. Thus the translational symmetry transforms local fields into non-local ones in general.
This coordinate dependence of locality would not be ridiculous when we recall the initial motivation of considering the noncommutative space-time. One of the authors regarded the noncommutative space-time (1) as a space-time realizing a space-time uncertainty relation derived from a gedanken experiment . We have no means to obtain a coordinate system more accurate than the limit specified by the space-time uncertainty relation. The uncertainty will become larger if a space-time event is farther from the origin of the coordinate system. Thus the space-time spread of an event is not an invariant notion anymore under the translational symmetry.
In places where the general relativity and quantum mechanics play major roles simultaneously, similar kinds of issues have already been observed in the literatures. The subtlety of the locality of an event has appeared already as an aspect of the black hole complementarity of Susskind . For instance, he argued that the location of the baryon violation in a black hole geometry is an observer dependent notion. Moreover, in the holographic description of the world , a space-time event is mapped to a screen, and the spread of the image will depend on the relative locations of the event and the screen .
Lastly we see that the action (88) may be derivable from a pregeometric action<sup>5</sup><sup>5</sup>5A pregeometric action in string theory was proposed in . in the form
$$S_{preg}=0|g(\widehat{\mathrm{\Psi }}^{},\widehat{\mathrm{\Psi }})|0.$$
(89)
An action with a kinetic term would be generated in a certain background of the field $`\widehat{\mathrm{\Psi }}=\widehat{\mathrm{\Psi }}_0`$. For example, starting from a cubic pregeometric action and considering the background of $`\widehat{\mathrm{\Psi }}_0`$ with $`\mathrm{\Phi }_0=m^2`$ and $`\mathrm{\Phi }_0^{\mu \nu }=\eta ^{\mu \nu }`$, we obtain an action with a kinetic term and cubic interactions for the fluctuation field $`\widehat{\mathrm{\Psi }}_1`$ in $`\widehat{\mathrm{\Psi }}=\widehat{\mathrm{\Psi }}_0+\widehat{\mathrm{\Psi }}_1`$. The action (89) would be more interesting than (88), since it has more symmetry.
## 6 Thermodynamics
In this section, we will discuss an analogue of thermodynamics of free scalar field theory in the noncommutative space-time. Reflecting the noncommutativity of the space-time, we find a non-trivial behaviour at high temperature. Following the common trick in usual commutative field theories, we will take the Euclidean metric and calculate the partition function, imposing the periodicity of the inverse temperature in the time-direction.
In the Euclidean metric with signature $`\eta ^{\mu \nu }=(1,1,1)`$ <sup>6</sup><sup>6</sup>6The results in this section are the same for the signature $`\eta ^{\mu \nu }=(1,1,1)`$., the consistency with the Jacobi identity changes the last commutation relation of (1) to
$$[\widehat{p}^\mu ,\widehat{x}^\nu ]=i\eta ^{\mu \nu }\sqrt{1l_P^2\widehat{p}^2}+il_Pϵ^{\mu \nu \rho }\widehat{p}_\rho .$$
(90)
Thus the coordinate operators in the momentum representation are given by
$$\widehat{x}^\mu =i\sqrt{1l_P^2p^2}\frac{}{p_\mu }il_Pϵ^{\mu \nu \rho }p_\nu \frac{}{p^\rho },$$
(91)
where we have ignored the possible ambiguity .
We now use a parameterisation
$`p_0`$ $`=`$ $`{\displaystyle \frac{1}{l_P}}\mathrm{cos}\omega \mathrm{sin}\theta ,`$ (92)
$`p_1`$ $`=`$ $`{\displaystyle \frac{1}{l_P}}\mathrm{sin}\omega \mathrm{cos}\phi ,`$ (93)
$`p_2`$ $`=`$ $`{\displaystyle \frac{1}{l_P}}\mathrm{sin}\omega \mathrm{sin}\phi ,`$ (94)
where the ranges of the parameters are given by
$$0\theta <2\pi ,0\omega \frac{\pi }{2},0\phi <2\pi .$$
(95)
Then the measure is given by
$$(1l_P^2p^2)^{\frac{1}{2}}d^3p=d\theta d\omega d\phi \mathrm{sin}\omega \mathrm{cos}\omega .$$
(96)
Using this parameterisation, the coordinate operators are rewritten as
$`x^0`$ $`=`$ $`il_P\left({\displaystyle \frac{}{\theta }}{\displaystyle \frac{}{\phi }}\right),`$ (97)
$`x^1`$ $`=`$ $`il_P\left(\mathrm{tan}\omega \mathrm{sin}(\theta \phi ){\displaystyle \frac{}{\theta }}+\mathrm{cos}(\theta \phi ){\displaystyle \frac{}{\omega }}\mathrm{cot}\omega \mathrm{sin}(\theta \phi ){\displaystyle \frac{}{\phi }}\right),`$ (98)
$`x^2`$ $`=`$ $`il_P\left(\mathrm{tan}\omega \mathrm{cos}(\theta \phi ){\displaystyle \frac{}{\theta }}+\mathrm{sin}(\phi \theta ){\displaystyle \frac{}{\omega }}+\mathrm{cot}\omega \mathrm{cos}(\theta \phi ){\displaystyle \frac{}{\phi }}\right).`$ (99)
Because of the noncommutativity of the coordinate operators, it is impossible to find an operator which shifts $`x^0`$ by a certain c-number while $`x^{1,2}`$ remain intact. Below we shall just find an operator which shifts $`x^0`$ by $`\beta `$ but does not change the expectation values $`\varphi |x^{1,2}|\varphi `$, where $`|\varphi `$ is an eigen-state of $`x^0`$.
From (99), we see that the eigen-function of the operator $`x^0`$ is given by
$$\varphi =e^{id(\theta \phi )}f(\theta +\phi ,\omega ).$$
(100)
The periodicities of the wave function under $`\theta \theta +2\pi ,\phi \phi +2\pi `$ constrain the values of $`d`$ as
$$d=\frac{n_d}{2},$$
(101)
where $`n_d`$ is an integer, and leads to that $`f`$ must satisfy
$$f(y,\omega )=(1)^{n_d}f(y+2\pi ,\omega ).$$
(102)
A candidate of the shift operator with $`U^{}x^0U=x^0+\beta `$ is given by
$$U(\beta )=\mathrm{exp}\left(\frac{i}{l_P}\beta \theta \right),$$
(103)
where, in order to satisfy the periodicity with respect to $`\theta `$,
$$\beta =l_Pn_\beta $$
(104)
with a positive integer $`n_\beta `$. Since $`U^{}x^1U=x^1+\beta \mathrm{tan}\omega \mathrm{sin}(\theta \phi )`$, we find
$`\varphi |\delta x^1|\varphi `$ $`=`$ $`\varphi |(U^{}x^1Ux^1)|\varphi `$ (105)
$`=`$ $`\beta {\displaystyle 𝑑\theta 𝑑\omega 𝑑\phi \mathrm{sin}\omega \mathrm{cos}\omega |f(\theta +\phi ,\omega )|^2\mathrm{tan}\omega \mathrm{sin}(\theta \phi )}`$ (106)
$`=`$ $`0.`$ (107)
Similarly $`\varphi |\delta x^2|\varphi =0`$. The shift operator $`U`$ could contain $`\phi `$-derivative in order to satisfy only the above features. However, since $`\phi `$-derivative generates a spatial rotation which we do not need, we assume that (103) is the appropriate operator.
The thermodynamic partition function for the inverse temperature $`\beta `$ is obtained by taking the trace with respect to the states satisfying $`U|\varphi =|\varphi `$. Hence we sum up over the states with $`\theta =2\pi l_Pn/\beta `$, where $`n`$ is an arbitrary integer. Thus we obtain
$`\mathrm{ln}Z`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{Tr}_{U=1}\left(\mathrm{ln}(p^2+m^2)\right)`$ (108)
$`=`$ $`{\displaystyle \frac{1}{8\pi }}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle 𝑑\theta d^2q\delta \left(\theta \frac{2\pi l_Pn}{\beta }\right)\mathrm{ln}\left(\frac{1}{l_P^2}\mathrm{sin}^2\theta +q^2\mathrm{cos}^2\theta +m^2\right)}.`$ (109)
The overall numerical constant is determined by comparing with the usual commutative case in the limit $`\beta \mathrm{}`$<sup>7</sup><sup>7</sup>7Since the momentum space is doubly degenerate, we equate (109) in the $`\beta \mathrm{}`$ limit with the usual thermodynamic partition function of two real scalar fields.. The integration over $`\theta `$ can be done by the contour integration with the new variable $`z=e^{i\theta }`$ after applying the formula $`_me^{2\pi mti}=_n\delta (tn)`$ to (109). The result is
$`\mathrm{ln}Z`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{\frac{1}{l_P}}}q𝑑q\mathrm{ln}\left(y_+^{\beta /2l_P}y_+^{\beta /2l_P}\right),`$ (110)
$`y_+`$ $`=`$ $`\sqrt{1+{\displaystyle \frac{q^2+1/l_P^2}{q^2+m^2}}}+\sqrt{{\displaystyle \frac{q^2+1/l_P^2}{q^2+m^2}}}.`$ (111)
The logarithm of the partition function (111) has the usual behaviour $`1/\beta ^2`$ in the low temperature. If we ignore the quantisation of $`\beta `$ (104), it has the behaviour $`\mathrm{log}\beta `$ at the high temperature $`\beta l_P`$. This suggests the reduction of the degrees of freedom in the high energy region. If we can apply the usual first law of thermodynamics to this system, we obtain the entropy density $`s\mathrm{ln}\beta `$ and the energy density $`e1/\beta `$ at the high temperature $`\beta l_P`$. Thus $`s/\sqrt{e}`$ decreases when the temperature increases. On the other hand, in the low temperature region, $`s/\sqrt{e}1/\sqrt{\beta }`$, and it increases with the temperature. Thus there is an upper bound of $`s/\sqrt{e}`$, in agreement with the qualitative argument given in .
However, it is not clear whether we can really regard $`\beta `$ as the inverse of the temperature, and can use the first law of thermodynamics to obtain the energy and entropy. To have a reliable discussion on these issues, we have to construct statistical thermodynamics in the noncommutative space-time. To do so, we need a Hamiltonian and its spectra. Moreover, in usual commutative cases, a thermodynamic system is put in an imaginary box and the states are well regularized by the infrared cut-off given by the size of the box. Since these issues are non-trivial in the noncommutative space-time, the construction of statistical thermodynamics remains unsolved.
## 7 Summary and discussions
In this paper, we have analysed scalar field theories in the noncommutative three-dimensional space-time characterised by $`[x^\mu ,x^\nu ]=2il_Pϵ^{\mu \nu \rho }x_\rho `$. The one-particle Hilbert space in momentum representation is represented by the $`SL(2,R)`$ group manifold. Since momentums are elements of $`SL(2,R)`$, the interaction vertices have definite ordering of legs coming from the noncommutativity of the group elements. We have performed some one-loop computations, which can be evaluated explicitly by contour integrations in complex planes. The non-planar one-loop diagrams of the two-point functions from $`\varphi ^4`$ and $`\varphi ^3`$-interactions were shown to be finite and have the infrared singularity coming from the UV/IR mixing.
The most peculiar feature of the noncommutative space-time is that the commutation relations among the coordinates do not respect the translational symmetry. This violation of translational symmetry is natural from the view point of a fuzzy space-time . If an event is farther from the origin of a reference frame, the location of the event will become more fuzzy by the quantum fluctuation of the space-time. Since the translational symmetry is recovered in the $`l_P0`$ limit in the commutation relations (1), one might expect that the momentum conservation would be recovered in this limit. However we have shown explicitly that this is not true for the non-planar contributions. To remedy this defect, we had to introduce an infinite number of tensor fields. We leave the analysis of the messy theory with an infinite number of tensor fields for future work.
On the other hand, the noncommutativity at the boundary of a membrane in the $`C`$-field background in M-theory is the loop-space noncommutativity $`[x^\mu (\sigma ),x^\nu (\sigma ^{})]iϵ^{\mu \nu \rho }x_\rho ^{}(\sigma )\delta (\sigma \sigma ^{})`$ . This noncommutativity obviously respects the translational symmetry. Since the noncommutative string theory comes from a limit of M-theory, the theory might be controllable. Thus, concerning the interests in noncommutative field theories in more than two-dimensions, the noncommutative string theory might be interesting.
We have finally discussed an analogue of thermodynamics of noncommutative free scalar field theory. What we have computed is the partition function of free scalar field theory in a noncommutative three-dimensional space with Euclidean signature, following the common trick in usual commutative field theories. The result shows the reduction of the degrees of freedom in the ultraviolet, which comes essentially from the compactness of the momentum space in the case of Euclidean signature. However it is not clear to us whether the partition function is really related to the thermodynamics of the noncommutative field theory, since the noncommutativity of the time coordinate makes it hard to define statistical thermodynamics in a convincing way.
There seem to exist several puzzles on the three-dimensional noncommutative field theory such as unitarity and renormalizability. To analyse these issues perturbatively, we need to develop further the computational technic of $`SL(2,R)`$ group theoretic Feynman integral. We hope that the present work may be meaningful as a primitive step toward consistent treatment of quantum field theories in more than two-dimensional noncommutative space-times.
Acknowledgements
The authors would like to thank N. Ikeda for valuable discussions. N.S. is supported in part by Grant-in-Aid for Scientific Research (#12740150), and in part by Priority Area: “Supersymmetry and Unified Theory of Elementary Particles” (#707), from Ministry of Education, Science, Sports and Culture.
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# Gravitational Lensing by Nearby Clusters of Galaxies
## 1 Introduction
Gravitational lensing is a powerful technique to probe distant galaxies as well as for studying the matter distribution in galaxy clusters. Indeed, the analysis of bright arcs and arclets or other distortions (weak lensing), induced by the gravitational lensing of background sources by a galaxy cluster, has allowed the determination of the mass distribution in these structures, independently of other more common techniques, like the application of the virial theorem or the analysis of the X-ray emission.
Most studies of gravitational lensing by clusters have been focused in distant objects, resulting in the discovery of rather high $`z`$ arcs and arclets. Galaxy lensing by clusters, assuming a non-evolving mass profile and a reasonable redshift distribution for the faint galaxy population, has its maximum around $`z0.2`$ (e.g. Natarajan & Kneib, 1997). However, this does not imply that the lensing efficiency of nearby clusters is totally negligible and, in fact, several groups have recently found evidence of strong lensing effects in low-redshift clusters. For instance, Allen, Fabian & Kneib (1996) discovered a $`z=0.43`$ arc in the cluster PKS0745-191 (at $`z=0.103`$), that has been successfully modeled as a gravitational lens image. Shaya, Baumn & Hammergren (1996) found an arc-like structure, still without redshift information, near NGC4881 in the Coma cluster ($`z=0.024`$). Campusano & Hardy (1996) found an arc-like object at $`z=0.073`$ in A3408, a cluster at $`z=0.042`$. Lens models of this structure are discussed by Campusano, Kneib & Hardy (1998). Blakeslee & Metzger (1999) discovered an arc-like object in A2124 ($`z=`$0.066) that is probably the lensed image of a galaxy at $`z`$=0.573. Campusano, Kneib & Hardy (1998) predicted the detection of weak-shear in low ($`z<0.1`$) clusters, which has been recently confirmed by Joffre et al. (1999) and Kneib et al. (2000, in preparation).
It is worth to point out that the study of gravitational lensing by low redshift clusters presents, in principle, an important advantage when compared with those of more distant clusters: due to the large angular diameter that nearby clusters have, gravitational lensing may allow examining in great spatial detail the mass distribution of their central regions. Although the lensing efficiency of a cluster depends strongly on its central mass distribution, the latter is usually not well known. Observations of lensing effects allow probing low redshift clusters with high spatial resolution and, consequently, they can help to add new constraints on the mass distribution in the centers of these structures. For the weak lensing regime (see Mellier, 1999, for a review), the magnitude of the effects in high redshift clusters depends strongly on the imprecisely known redshift distribution of the faint background galaxies, while in low-$`z`$ clusters the weak lensing effects are almost independent of it.
In this paper we present a simple estimation of the expected number of arcs and arclets in low redshift clusters, as well as the results of an analysis of a sample of nearby clusters ($`z0.076`$), where we have looked for arc-like structures that may be produced by gravitational lensing, either by the central cluster potential as a whole or by substructures in the mass distribution related to a galaxy not located at the fiducial center of the cluster. Following Hattori, Kneib & Makino (1999) we call an arc a structure distorted by gravitational lensing with axial ratio (length-to-width) larger than 10, and arclet a structure with axial ratio smaller than 10.
The layout of the paper is as follows. The sample of galaxy clusters analyzed here is discussed in Section 2. The estimation of the number of luminous arcs and arclets expected in a sample of clusters is presented in Section 3 (and Appendix A). In Section 4 we describe the search for arcs and arclets in the sample. The two arclet-like structures found in our search are presented and discussed in Section 5. Follow up observations of these arclets are also presented in this section. Finally, we summarize our conclusions in Section 6. When necessary, we adopt $`H_0=50`$ h$`{}_{50}{}^{}{}_{}{}^{1}`$ km s<sup>-1</sup> Mpc<sup>-1</sup> and, unless where explicitly stated, $`\mathrm{\Omega }_0`$ = 1 and $`\lambda `$ = 0.
## 2 The sample of nearby galaxy clusters
The images analyzed in this project were originally obtained as part of the PhD thesis of Daniel A. Dale on peculiar motions of clusters with $`z<0.1`$ (Dale et al., 1997, 1998, 1999a, 1999b, 1999c). The sample of galaxy clusters with some relevant properties is presented in Table 1.
The observational material that is analyzed here consists of several Kron-Cousins I-band images obtained with the 0.9m CTIO telescope. The details of the observations are discussed in (Dale et al., 1997, 1998) and here we only summarize them. The detector used was the 2k$`\times `$2k Tek2k No.3 CCD, with a scale of 0.396 arcsec per pixel, resulting in a field of 13.5$`{}_{}{}^{}\times 13.5^{}`$ per image. The exposure times were of 600 seconds in all cases. The images reach $``$ 23.7 I mag arcsec<sup>-2</sup> at the 1.0 $`\sigma `$ level over the sky background. In Section 4 we estimate that the isophotal completeness limit for extended sources in these images is $``$19.5 mag (1.5 $`\sigma `$ in 40 connected pixels), corresponding to a total magnitude limit of $``$ 19.0 mag. The average seeing of the images is 1.4<sup>′′</sup>. Although these images are not very deep, this material is the same where Campusano & Hardy (1996) found an arc in A3408. These images in general do not uniformly cover a cluster, since they were taken in regions near spiral galaxies, but the central region of the clusters are always covered. The number of images analyzed per cluster are also presented in Table 1.
This sample is a good representation of the cluster distribution in the nearby universe. Its richness distribution is presented in Table 2. It is consistent with the whole Abell catalog (Abell, Corwin & Olowin, 1989) , with a slight excess of high richness clusters. In Figure 1 we compare the X-ray luminosity distribution of our sample with that of the XBAC catalogue (Ebeling et al., 1996) in the same redshift range (i.e., 0.014 $`z`$ 0.076). The XBAC catalogue is a X-ray flux limited catalog of Abell clusters from the ROSAT all sky survey. Note that the clusters that were not detected by ROSAT have been included in the first bin in Figure 1. This figure indicates that the X-ray luminosity distribution of our sample presents an overall agreement with the cluster distribution of the nearby Universe.
If richness and X-ray luminosity are proportional to cluster mass, we conclude that the mass distribution of our cluster sample is representative of that actually present at low redshifts.
## 3 Probabilities of strong lensing by nearby clusters
In this section we will discuss how often one should expect to observe any signature of strong lensing effects in a sample of nearby clusters, taking into account some observational constraints. The lensing model adopted here assumes that the cluster mass profile may be described by a singular isothermal sphere (SIS) and is presented in detail in Appendix A. The model has 10 parameters. Each cluster is characterized by two parameters: its redshift $`z`$ and the one-dimensional velocity dispersion $`\sigma _v`$. The luminosity function of the field galaxies is described by a Schechter function and has three parameters: $`\varphi ^{}`$, $`M^{}`$, and $`\alpha `$. The analysis of the observations also has three parameters: the adopted limit magnitude $`m_l`$, the seeing of the observations $`\sigma _{seeing}`$, and the minimum flux amplification by lensing, $`A_{min}`$. Note that in the SIS model for strong lensing $`A`$ is also equal to the tangential stretching of the arc or arclet. Once the cosmological model is specified by the density parameters associated to the mass and to the vacuum, $`\mathrm{\Omega }_m`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$, all the parameters for the calculation are fixed and there is no dependence on the value of $`H_0`$.
We adopt here the luminosity function derived in the Stromlo-APM Redshift Survey (Loveday et al., 1992), that is representative of the local field galaxies. This luminosity function is well fitted by a Schechter function with parameters $`M_{b_j}^{}=19.50+5\mathrm{log}h`$, $`\alpha =0.97`$, and $`\varphi ^{}(0)=1.40\times 10^2h^3`$ Mpc<sup>-3</sup>. Assuming a color $`(b_jI)=1.76`$ (Fukugita, Shimasaku & Ichikawa, 1995), that is appropriate for a Sbc galaxy (that is a good average of the local morphological mix of galaxies), we have that $`M_I^{}=22.8`$ for $`H_0=50`$ km s<sup>-1</sup> Mpc<sup>-1</sup>. For simplicity, we neglect any evolution of the parameters $`\alpha `$ and $`M^{}`$ with $`z`$. This luminosity function is consistent with galaxy counts in the I-band (Smail et al., 1995).
We first discuss how the expected number of arcs varies with the magnitude limit of the observations, neglecting seeing effects. Figure 2 shows how the expected number of arcs with A<sub>min</sub> =2 (the minimum amplification produced by a SIS, see Appendix A), for a single cluster at $`z=0.05`$ with $`\sigma _v=1000`$ km s<sup>-1</sup> (a massive, Coma-like cluster), varies with the magnitude limit $`m_l`$ of the arc search. The results are presented for three cosmological models: standard CDM (SCDM: $`\mathrm{\Omega }_M=1`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$); open CDM (OCDM: $`\mathrm{\Omega }_M=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$); and cosmological constant CDM ($`\mathrm{\Lambda }`$CDM: $`\mathrm{\Omega }_M=0.3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$). With $`m_{l,I}=19.0`$ (the limit adopted in the search), we have that $`<N>`$ is equal to 0.30, 0.31 and 0.37, for SCDM, OCDM and $`\mathrm{\Lambda }`$CDM, respectively. A value of $`<N>=1`$ is achieved at $`m_{l,I}=20.2,20.1`$, and 19.9, for these 3 cosmologies. Figure 2 reveals that the expected number of arcs increases strongly with $`m_l`$, and searches going one magnitude deeper than ours may plausibly find 2 to 3 times more evidence of strong lensing than our own search. These results also indicate that different cosmological models lead to similar results, at least for bright values of $`m_l`$. A major source of uncertainty in this kind of calculation is due to the normalization of the luminosity function, $`\varphi ^{}`$, as evidenced by galaxy number counts in different directions. Assuming an uncertainty of a factor 2 in $`\varphi ^{}`$, $`<N>=1`$ would be attained for $`19.3<m_{l,I}<21.0`$, where the brighter and the fainter limits are for $`\mathrm{\Lambda }`$CDM and SCDM, respectively.
Figure 3 shows the dependence with the redshift of the expected number of arcs with A<sub>min</sub> =2, for a cluster with $`\sigma _v=1000`$ km s<sup>-1</sup>, adopting m<sub>l</sub> as 19.0, for the same three cosmologies. Here too we are neglecting seeing effects. This figure indicates that nearby clusters are more efficient than far ones to produce arcs brighter than some magnitude limit. This is due to the fact that low redshift clusters project larger angular cross sections on to the plane of the sky than more distant clusters. This result is not in disagreement with the statistics of arcs as a function of the lens redshift, which has a maximum in the range 0.2 $`<z<`$ 0.4 (e.g. (e.g. Fort & Mellier, 1994). Indeed, Figure 3 shows the expected number of arcs per cluster, and those arc statistics not only consider the lensing efficiency of a cluster at a given redshift, but also the total number of clusters (or the comoving volume if we consider a constant density of clusters) per bin of redshift too.
To calculate $`<N>`$ for our whole cluster sample, we need the velocity dispersion of the clusters. Unfortunately this quantity is not known for 5 clusters in the sample. For A3408 we used the value of Campusano, Kneib & Hardy (1998) “dark halo” scenario. The remaining 4 clusters were removed from this analysis (although they have been included in the arc search described in the next section). Since their velocity dispersions are probably small (due to their low richness and absence of detectable X-ray emission), their impact on the results should be negligible. The values of $`\sigma _v`$ adopted in the calculation were corrected to their rest frame values (i.e., the actual value is the observed value that appears in Table 1 times $`(1+z)^1`$).
Assuming $`A_{min}`$ =2 and the same magnitude limit adopted in the arc search, $`m_{l,I}=19.0`$, the expected number of arcs and arclets in our cluster sample is 2.9 for SCDM, 3.0 for OCDM, and 3.6 for $`\mathrm{\Lambda }`$CDM. Seven clusters with $`\sigma _v900`$ km s<sup>-1</sup> contribute with more than 55% for $`<N>`$. Note that these are upper limits, actually, since this estimate does not include the seeing.
The seeing may dramatically affect this kind of estimate, because it tends to circularize non-circular sources. In order to quantify seeing effects, we have made simulations, using the IRAF package artdata. We have simulated “arcs” as exponential profiles with central surface brightness 19.9 I-band mag arcsec<sup>-2</sup> (typical of a Freeman, 1970, galaxy disc) with several total apparent magnitudes and axial ratios, assuming the same sky level and noise of our images. These simulated images have been convolved with a Moffat PSF with FWHM 1<sup>′′</sup>.4 (the average seeing of our images) and then their magnitudes and axial ratios were measured with the software SExtractor (Bertin & Arnouts, 1996), using the same parameters used in the search for arcs (Section 4).
Let us assume that we can identify a gravitational arc if it has an axial ratio of at least 1.5. For a given value of its intrinsic axial ratio, this arc would be identified only if its magnitude is lower than some limit $`m_{max}`$, because fainter images would appear with amplification inferior to 1.5 due to the seeing. Large and bright arcs are not strongly affected by the seeing, but faint and small arclets are. With 500 simulations for 50 values of the intrinsic axial ratio, we have estimated the mean value of $`m_{max}`$. The results are illustrated in Figure 4. Another result of the simulations is that the difference between total and isophotal magnitudes near the limit of detection is about 0.5. Considering that the observed limit of completness is $``$ 19.5, we conclude that the total magnitude limit is $``$ 19.0 mag.
We present in Figure 5 the expected number of arcs for a cluster with $`\sigma _v=1000`$ km s<sup>-1</sup> at $`z=0.05`$. This figure shows how dramatic are seeing effects for lensing detection. The results displayed in Figure 5 for each cosmological model may be parametrized as $`N_0\mathrm{exp}(\beta \sigma _{seeing})`$, where $`N_0`$ is the expected number of arcs in the absence of image degradation due to the seeing. A value of $`\beta =1.5`$ provides a good fit to the data.
We have recalculated the estimates of the number of arcs in our sample taking the seeing into account. The expected values are 0.32, 0.34, and 0.42, for SCDM, OCDM and $`\mathrm{\Lambda }`$CDM models, respectively. Note that the inclusion of the seeing in the analysis has lead to a substantial decrease in the number of expected arcs; these numbers are $``$9 times smaller than when seeing effects are neglected. Using the OCDM results and Poissonian statistics, the probability of one arc detection in the sample is 24.2%, two detections is 4.1%, three detections is 0.5%, and no detections is 71.2%. Hence, our estimation is more consistent with no detections, but the probability of finding at least one arc is not negligible.
## 4 Search of bright arcs and arclets
In this section we describe the procedure we have adopted in the search for arcs and arclets produced by gravitational lensing in the sample of clusters presented in Section 2. We have looked for evidence of strong lensing not only in the regions corresponding to the central parts of the clusters, where the projected mass density is high and, consequently, the probability of lensing is higher than in other regions, but also around bright galaxies.
Our strategy for the search was the following. Initially, we made a catalog of galaxies using SExtractor. We adopted a detection threshold of 1.5 $`\sigma `$ over the sky level (which corresponds to $`I`$ 23.3 mag arcsec <sup>-2</sup>) and a minimum detection area of 40 pixels (or 6.27 arcsec<sup>2</sup>). The distribution of magnitudes of the galaxy catalog presents a cutoff at $`I20`$, indicating that its completeness limit is at $`I19.5`$.
Afterwards, all galaxy images with semi-major axis larger than 8 pixels (3.17 arcsec) were modeled with the STSDAS/Ellipse package, and the model images were subtracted from the actual galaxy images. The aim here was to reveal any arc-like structure superposed on to a galaxy image. In general this procedure worked well, showing the presence of several objects within the galaxy image. This procedure tends to produce a central residual, due to the seeing and pixelization but, since these residuals are usually restricted to the very central regions of a galaxy image, they do not have a relevant impact in our arc survey. During the process of image subtraction we inspected by eye most of the images in each field, paying special attention to objects with axial ratio larger or equal to 1.5.
After the subtraction of the galaxy images, we created a new SExtractor galaxy catalog and visually inspected all new objects contained in this second catalog, focusing again on the more elongated ones. At first we selected about twenty candidates, elongated objects which could not be morphologically identified, by visual inspection, as edge-on spirals. Then we verified whether these objects were tangentially or radially disposed with respect to the cluster center or some bright, nearby galaxy. After this stage, only two candidates remained. The first is the same arclet discovered by Campusano & Hardy (1996) in A3408 and discussed by Campusano, Kneib & Hardy (1998). The second candidate was found in the cluster A3266. It is not in the central region of the cluster but near one bright elliptical galaxy. It is, hence, a candidate for lensing by a cluster substructure, instead of lensing by the cluster overall potential.
## 5 Discussion
In this section we discuss the main characteristics of the two arclet candidates, taking into account some new follow-up observations.
### 5.1 The arclet in A3266
The cluster A3266 (also known as Sérsic 40/6) is apparently regular, with type I-II in the Bautz-Morgan system. Its center contains a very tight dumbbell pair, at $`\alpha =4^h31^m14.25^s`$ and $`\delta =61^{}27^{}11.3^{\prime \prime }`$ (J2000). The recession velocity of the cluster relative to the CMBR is 17782 km s<sup>-1</sup>. A detailed analysis of this cluster by Quintana, Ramírez & Way (1996) reveals that it has a large velocity dispersion, 1306 $`\pm `$ 73 km s<sup>-1</sup> within $`1h^1`$ Mpc, that has been interpreted as an evidence that A3266 is indeed the result of a recent merger of two structures of comparable masses. This interpretation is also supported by numerical simulations (Flores, Quintana & Way, 1999) as well as by an analysis of the X-ray brightness distribution (Mohr, Fabricant & Geller, 1993).
We present in Figure 6 part of the I-band image of A3266, taken with the 0.9m CTIO telescope. The candidate arclet is indicated with an arrow. Its centroid is at $`\alpha =4^h31^m15.53^s`$ and $`\delta =61^{}30^{}3.7^{\prime \prime }`$ (J2000). It is at 16.6 arcsec (29 h$`{}_{}{}^{1}{}_{50}{}^{}`$ kpc at the cluster distance) from the center of a nearby, bright elliptical, and at 2.89 arcminutes (303 h$`{}_{}{}^{1}{}_{50}{}^{}`$ kpc) from the center of the cluster (at the position of the dumbbell pair). At the isophotal level of 23.3 mag arcsec<sup>-2</sup> (1.5$`\sigma `$ over the sky background) its magnitude, semi-major axis, axial ratio, and position angle are $`I=18.89\pm 0.03`$, $`a=4^{\prime \prime }.1`$, $`a/b=2.11\pm 0.34`$, and $`\theta =3.5^{}\pm 2.9^{}`$, respectively.
The elliptical galaxy near the object is located at $`\alpha =4^h31^m16.6^s`$ and $`\delta =61^{}30^{}8^{\prime \prime }`$ (J2000) and is the second brightest galaxy of the cluster in the I-band (apart from the central dumbbell pair). Its heliocentric radial velocity is 15819 $`\pm `$ 30 km s<sup>-1</sup> and its apparent total B magnitude is 15.40 $`\pm `$ 0.20 (de Vaucouleurs et al., 1991), corresponding to an absolute magnitude of -22.38.
Spectroscopic observations of the arclet by W. Kunkel (private communication) with the 2.5m Dupont telescope have shown that this object has an heliocentric radial velocity of 21900 km s<sup>-1</sup>. With this velocity, this object may be inside the cluster, since most of the cluster members have velocities between 15000 km s<sup>-1</sup> and 21000 km s<sup>-1</sup> (Quintana, Ramírez & Way, 1996) On the other hand it can also be behind the cluster, which strengths the probability that it has been lensed either by the cluster potential or by the nearby galaxy, or both. In particular we want to test the hypothesis that this is a background galaxy being lensed by a mass peak associated with the bright elliptical near the arc.
Let us assume a SIS model for the mass distribution centered at the center of the nearby, bright elliptical galaxy. From the observed elongation of the arc candidate (that is equal to its amplification), the distance between the center of the lens and the arc image is about two critical radius (see Appendix, Equation A2). Hence, the SIS should have a velocity dispersion of 1265 km s<sup>-1</sup>, approximatelly the same of the cluster itself. The mass enclosed within the arc radius is $`3.0\times 10^{13}`$ h$`{}_{}{}^{1}{}_{50}{}^{}`$ M. The bright galaxy near the arc has $`M_I=24.65`$ at the 23.7 mag arcsec<sup>-2</sup> level. Assuming $`VI=1.31`$, appropriate for an elliptical galaxy (Fukugita, Shimasaku & Ichikawa, 1995) the arc would be explained by our lens model if M/L<sub>V</sub>=163 h<sub>50</sub> M/L for this system. Such a high M/L value could be expected if this galaxy were at the center of a massive sub-structure. However, neither the X-ray map of the cluster (Jones & Forman, 1999) nor the dynamical analysis of Quintana, Ramírez & Way (1996) present any evidence of significant sub-structure at this position. Possibly, this arc-like feature is a disk galaxy (or the bar of a disk galaxy) that is a cluster member (or is not far from it) instead of a real arclet.
### 5.2 The arclet in A3408
This structure has been discussed by Campusano, Kneib & Hardy (1998), who successfully modeled it as a galaxy at $`z=0.073`$ lensed by the cluster A3408, at $`z=0.042`$. The adopted cluster mass distribution is a scaled version of mass profiles derived from the study of high redshift clusters. Their preferred model has a component that follows the brightness profile of the central elliptical galaxy and is immersed in a massive dark halo. From the lens model and the equivalent widths of some prominent emission lines (\[O II\] $`\lambda `$3727, \[O III\] $`\lambda `$5007 and H$`\alpha `$) they have suggested that the source galaxy is probably a spiral with intrinsic diameter 14.6 kpc and magnitude $`M_B=18.2`$.
We have imaged the central part of A3408 with the 0.9m telescope of CTIO with interference filters with the aim of detecting other galaxies at the redshift of the source, which could help to improve the lens model. We have used two different filters. One is centered at 7053 Å, with FWHM of 79 Å, which allows the detection of H$`\alpha `$ at the redshift of the source, covering a velocity range of 1700 km/s. The other is centered at 6961 Å, with FWHM of 79 Å, and samples the continuum near the H$`\alpha `$ line. We have used the Tek2k3 CCD to make 5 images of 15 minutes each in each of the two filters.
The images were reduced using standard procedures with IRAF. They were stacked and normalized so that at the end of the reduction we had two images (one for each filter) where the mean flux (counts) of the stars were the same in both images. After that, we produced a new image by subtracting the image taken with the continuum filter from the image taken with the filter centered in H$`\alpha `$. Figure 7 presents both the H$`\alpha `$ and the residual images of the central region of A3408, centered at the position of the star near the arc position. The image containing the residuals indicates that the image subtraction was good, despite the features that remained at the center of the star and galaxy images (produced by seeing and pixelization effects), since most of the extended regions of the object images were removed. Figure 7 also indicates that the H$`\alpha `$ emission of the arc is not uniform and is strongest at the western side of this object.
An interesting feature present in the residual image is a point-like object, between the side of the arc with the strongest H$`\alpha `$ emission and the star below the arc. This object also appears in broad band images, after removing the image of this star, as shown in Figure 8. This feature is probably a companion galaxy of the arc source, at approximately the same redshift. Unfortunately, the presence of the star precludes further analysis of this object and its use to constrain more sophisticated lens models. No other objects are seen in Figure 7 at the same redshift of the arc. Note that this does not mean that a galaxy group (containing the arc source) can not exist, because only their brightest members could be detected in these images. Moreover the CCD covers an area of 1.52 $`\times `$ 1.52 h$`{}_{50}{}^{}{}_{}{}^{1}`$Mpc at z = 0.073. Some loose groups occupy areas larger than it (up to 5 Mpc of side), so we might not be sampling the entire group; however, in this case the over-density caused by this group behind A3408 might be not so significant and the impact of its presence on our probability calculation would be small.
Overall, A3408 seems to be a an interesting structure. It is a poor cluster, with low X-ray emission but high lensing mass, consistent with its velocity dispersion of 900 km s<sup>-1</sup>. On the other side, it may be in interaction with A3407, since they are very close to each other (Galli et al., 1993). Weak lensing observations of this cluster, combined with other methods of mass determination, like dynamical or X-ray analysis, should give us a better understanding of this system.
### 5.3 Comparison of theoretical and observational results
As discussed in Section 3, the probabilities for 0 and 1 detection of a gravitational arc in the cluster sample discussed here and including seeing effects are 71% and 24%, respectively. On the other side, our search for arcs in this sample produced two candidates, one in A3266 and the other in A3408. The arclet in A3408 is indeed good evidence of gravitational lensing by a nearby cluster (Campusano, Kneib & Hardy, 1998). The arc-like object in A3266, however, is more difficult to be interpreted as a result of gravitational lens distortion, because it would require a very massive substructure around the bright elliptical galaxy near the object. Note that one arc detection is consistent with our theoretical expectations but, of course, this does not allow us to strongly constrain any of the model parameters.
It is worth stressing the relative role played by the seeing and the magnitude limit in arc searches. As shown in Section 3, the expected number of arcs increases fast with the limiting magnitude of a survey, but seeing strongly affects fainter galaxies. Better seeing conditions lead to a substantial increase in the efficiency of these surveys. Searches of strong lensing effects in clusters will greatly benefit of the improvement in image quality that is arriving with active and adaptive optics cameras and new generation telescopes.
It is interesting to compare our results with those of Luppino et al. (1999). They have searched for arcs and arclets in a sample of 38 X-ray selected clusters ($`L_X2\times 10^{44}`$ erg s<sup>-1</sup>) from the Einstein Observatory Medium Sensitivity Survey (EMSS) in the redshift range 0.15 $`z`$ 0.823. Their images were obtained with the University of Hawaii 2.2m telescope in the R band for all clusters, and in the B band for most of them. The median seeing in R was 0.8<sup>′′</sup>. They found arcs and arclets in 8 clusters, or 21% of the sample, that can be compared with our success rate of only 3% (1 arc in 33 clusters). This discrepancy may be explained by taking into account that the two samples are very different regarding redshift distribution and image quality (resolution and magnitude limit). The images of Luppino et al. were obtained under better seeing conditions and are substantially deeper than ours. Besides, the EMSS sample, based on X-ray luminosities, is biased towards massive clusters.
## 6 Summary
We have discussed strong lensing effects produced by nearby clusters of galaxies. Using a simple mass model for the clusters we have shown that the expected number of arcs or arclets expected in a sample of nearby clusters is not as small as is usually thought. The results are strongly dependent on the magnitude limit and the seeing quality of the imaging.
Our search for arcs and arclets in a sample of 33 nearby clusters has resulted in two arc candidates, one in A3408 and the other in A3266. The first was already discovered by Campusano & Hardy (1996) and modeled by Campusano, Kneib & Hardy (1998). Observations in H$`\alpha `$ reported here show the presence of another object near the arc at approximately the same redshift, but contamination from a foreground star precludes its further use to constrain the cluster mass model. Our analysis of the arclet candidate in A3266 suggests or a false detection or that the mass concentration necessary to explain this structure by lensing is very strong.
Our result indicates that deep imaging of nearby clusters under good seeing conditions may be extremely useful for high-resolution mapping of the mass distribution of these structures, since the probability of detection of strong lensing features will be enhanced. Additionally, deep imaging with mosaic CCD detectors (e.g. Joffre et al., 1999) also allow the detection of weak lensing, which may provide more constraints on the mass distribution in the central regions of the galaxy clusters.
ESC and LS gratefully acknowledges support by Brazilian agencies FAPESP, CNPq and PRONEX . RG acknowledges support by NSF grant AST-9617069.
## Appendix A Estimation of the expected number of gravitational arcs and arclets in a sample of clusters
We assume here that the matter distribution of a galaxy cluster can be described by a singular isothermal sphere (SIS). The reason for this choice is twofold. First, this is the simplest model, with only one parameter: the one-dimensional velocity dispersion $`\sigma _v`$. Second, in this model one can easily relate amplification of a background source and its tangential stretch to the impact parameter (angular distance between the source and the center of the lens). We also assume that galaxies follow dark matter and have the same velocity dispersion.
More complex models may not be necessary, given the uncertainties in many parameters that enters in this calculation. For instance, if instead of a SIS we have considered a model with a core radius, the lensing probability would decrease (e.g. Wu & Hammer, 1993). Conversely, an ellipticity in the galaxy distribution, as well as the presence of sub-structures in a cluster, increases the lensing probability, due to enhanced tidal effects introduced by the asymmetry in the mass distribution (Bartelmann, 1995). Hence, the effects due to the inclusion of a core radius and cluster ellipticity may more or less cancel each other. Note also that the determination of these two quantities is not easy, being very sensitive to the choice of the cluster center and the presence of sub-clustering. Of course, even more refined mass models of each cluster are possible, where individual galaxy halos are taken into account (e.g. Kneib et al., 1996; Natarajan et al., 1998; Geiger & Schneider, 1999; Bézecourt et al., 1999), but such models are beyond the scope of this simple calculation. Anyway, we do not expect that the order of magnitude of the results presented here will change dramatically with the use of more sophisticated mass models.
Arcs and arclets (see Fort & Mellier, 1994, for a review) are a result of strong lensing by galaxy clusters of extended sources close to cusps or higher order catastrophes in the source plane. In what follows we will assume that an arc or arclet is produced if a source bright enough falls within the critical circle of the cluster in the source plane. It is well known that a spherically symmetric lens like our SIS model does not produce cusps. However, this drawback does not precludes the use of this sort of model to obtain estimates of the the strength of lensing effects. Indeed, such an approach has been extensively applied to compute lensing cross sections (e.g. Bartewlmann, 1993; Cooray, Quashnock & Miller, 1999; Cooray, 1999).
For a singular isothermal sphere be able to act as a strong lens, the source (in the source plane) should be within the critical circle inside which the mean mass surface density of the cluster is greater than a certain critical density that depends only of the relative distances between observer, cluster, and source (and hence of the cosmology). The angular radius of this critical circle (centered in the cluster center) is
$$\theta _c=4\pi (\sigma _v/c)^2\frac{D(z_l,z_s)}{D(0,z_s)}$$
(A1)
where $`D(z_l,z_s)`$ and $`D(0,z_s)`$ are the angular diameter distances between the lens (at redshift $`z_l`$) and the source (at $`z_s`$), and between the observer ($`z=0`$) and the source, respectively. The diameter distances are computed adopting the analytical filled-beam approximation (Fukugita et al., 1992). Following Blandford & Kochanek (1987), the magnification of the brighter of the two images produced by this lens is given by
$$A=1+\frac{\theta _c}{\theta _s}$$
(A2)
where $`\theta _s`$ is the angular distance on the plane of the sky between the center of the lens and the source (the impact parameter).
The expected number of arcs with magnification larger than a certain value $`A_{min}`$ (the minimum amplification produced by a SIS is 2) due to a cluster at redshift $`z_l`$ can be written as
$$<N(z_l)>=_{A_{min}}^{\mathrm{}}𝑑A_{z_l}^{\mathrm{}}𝑑zN(z,A)$$
(A3)
where $`N(z,A)dAdz`$ is the number of observable sources between $`z`$ and $`z+dz`$ that suffer magnifications between $`A`$ and $`A+dA`$ by the gravitational field of the cluster. These are the sources in that redshift interval with luminosities larger than $`L_{min}(z)/A`$, where $`L_{min}(z)`$ is the minimum luminosity that a source at redshift $`z`$ should have to be included in the sample (in the absence of lensing effects), and that are inside a ring with solid angle $`d\mathrm{\Omega }=2\pi \theta d\theta `$. Here $`\theta `$ is the angular distance to the cluster center corresponding to a magnification $`A`$ of the source luminosity, accordingly to Equation (A2). It is easy to verify that
$$d\mathrm{\Omega }=2\pi \theta _c^2\frac{dA}{(A1)^3}$$
(A4)
We can also write
$$N(z,A)dAdz=n(z)dV(z)\frac{d\mathrm{\Omega }}{4\pi }$$
(A5)
where $`n(z)`$ is the mean number density of the sources at redshift $`z`$ that are bright enough to be detected:
$$n(z)=_{L_{min}(z)/A}^{\mathrm{}}\varphi (L,z)𝑑L$$
(A6)
where the term $`L_{min}(z)/A`$ takes into account the magnification bias, by which some sources are magnified by lensing, acquiring an apparent luminosity high enough to be included in a magnitude limited sample.
The differential luminosity function at $`z`$, $`\varphi (L,z)`$, may be described by a Schechter function in which the comoving density of galaxies at redshift $`z`$ with luminosities between $`L`$ and $`L+dL`$ is
$$\varphi (L,z)dL=\varphi ^{}(z)\left(\frac{L}{L^{}(z)}\right)^{\alpha (z)}\mathrm{exp}\left(\frac{L}{L^{}(z)}\right)\frac{dL}{L^{}(z)},$$
(A7)
where $`\varphi ^{}(z)`$, $`L^{}(z)`$ and $`\alpha (z)`$ are the parameters of the luminosity function at redshift $`z`$. Then,
$$n(z)=\varphi ^{}(z)\mathrm{\Gamma }[1+\alpha (z),\frac{L_{min}(z,m_l)}{AL^{}(z)}]$$
(A8)
Assuming that the number of sources in a comoving volume is conserved, we have that
$$\varphi ^{}(z)=(1+z)^3\varphi ^{}(0)$$
(A9)
where $`\varphi ^{}(0)`$ is the normalization of the local ($`z=0`$) luminosity function.
The luminosity $`L_{min}`$ is related to the magnitude limit $`m_l`$ considered in the analysis and the magnification $`A`$ by
$$\frac{L_{min}(z,m_l)}{L^{}(z)}=A\times 10^{0.4(m_l5\mathrm{log}[(1+z)^2D(0,z)]25M^{})}$$
(A10)
where $`z`$ is the redshift of a source and $`M^{}`$ is the local value of the characteristic magnitude of the luminosity function in the photometric band of interest. Note that we are assuming that the source galaxies and $`M^{}`$ have the same $`k`$ and evolutive corrections.
The volume element $`dV(z)`$ that also appears in Equation (A5) is the proper volume between $`z`$ and $`z+dz`$ and may be written as
$$\frac{dV(z)}{dz}=\frac{4\pi cD^2(0,z)}{H_0}\frac{1}{(1+z)E(z)},$$
(A11)
where
$$E(z)=[(1+z)^3\mathrm{\Omega }_M+(1+z)^2(1\mathrm{\Omega }_M\mathrm{\Omega }_\mathrm{\Lambda })+\mathrm{\Omega }_\mathrm{\Lambda }]^{1/2}$$
(A12)
and where $`\mathrm{\Omega }_M`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ are the density parameters for matter and vacuum energy, respectively.
Putting all together, the expected number of arcs produced by a cluster at redshift $`z_l`$ is
$`<N(z_l,\sigma _v)>={\displaystyle \frac{32\pi ^3\varphi ^{}(0)\sigma _v^4}{c^3H_0}}\times `$
$`\times {\displaystyle _{A_{min}}^{\mathrm{}}}{\displaystyle \frac{dA}{(A1)^3}}{\displaystyle _{z_l}^{z_{max}}}dz{\displaystyle \frac{(1+z)^2D(z_l,z)^2\mathrm{\Gamma }[1+\alpha (z),\frac{L_{min}(z,m(A))}{AL^{}(z)}]}{E(z)}}`$ (A13)
where we have adopted $`z_{max}=3`$ (our results are insensitive to this limit of integration, because the number of sources brighter than $`m_l`$ decreases quickly with $`z`$).
A more realistic model must, however, take into account the effect of the seeing, that tends to circularize object images, specially the faint ones. The simulations described in Section 3 show that, for a given intrinsic stretching (or amplification $`A`$), the observed axial ratio will be larger than 1.5 (the minimum value adopted in the arc search discussed in Section 4) only for objects brighter than a certain apparent magnitude $`m_{max}(A)`$. Hence, seeing effects may be included in the analysis by using the minimum of $`m_l`$ and $`m_{max}(A)`$ instead of $`m_l`$ in Equation (A10). The relation between $`m_{max}`$ and $`A`$ adopted here is shown if Figure 4.
For a sample containing $`N_c`$ clusters, the total expected number of arcs is given by
$$<N>=\underset{i=1}{\overset{N_c}{}}<N(z_i,\sigma _{v,i})>$$
(A14)
where $`z_i`$ and $`\sigma _{v,i}`$ are the redshift and velocity dispersion of the $`i`$th cluster.
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# Cellular solid behaviour of liquid crystal colloids 1. Phase separation and morphology
## 1 Introduction
The phase separation and ordering of multi-phase systems, and their resulting physical properties, have been the subject of much active research over the past twenty years. Such a high interest in this field arises from a combination of challenging, fundamental physics (of phase equilibria and dynamics) with a large number of viable applications. Research in liquid crystals has also remained active for decades, for similar fundamental and practical reasons. However, the applications of thermotropic liquid crystals have always focussed on display technology, somewhat overlooking their unique mechanical properties. In the early 1990’s a new field, now called “liquid crystal colloids”, was opened to research by an experimental and theoretical study of suspensions of colloidal particles in a lyotropic liquid crystal by Poulin et al. roux . The key idea of introducing a large amount of mobile interface into the liquid crystal, which generates nontrivial topological constraints and singularities, has been fruitfully explored by further experimental and theoretical work coll95 ; luben ; rama ; interac .
The problem studied is that of a spherical particle with nematic director anchored on its surface, suspended in a nematic liquid crystal matrix. When the director is anchored rigidly, a closed inner surface creates a topological mismatch between the director field $`𝐧(𝐫)`$ on the particle surface and the uniform director at large distances from it. This mismatch leads to topological defects, i.e. regions where the liquid crystal order and the continuity of $`𝐧`$ break down. A connected closed surface represents a point topological charge $`N=1`$. Since the overall sample has to be topologically neutral, the charge produced by the inner surface must be balanced by an associated opposite charge. There are two basic possibilities. The assumption, that a spherical particle in a quadrupolar nematic medium creates deformations that are quadrupolar too, results in the mismatch balanced by a circular loop of disclination with linear strength $`N=1/2`$ and overall point charge $`N=1`$, in the equatorial plane of the particle: a “Saturn ring” ukra , see fig.1(a).
The other possibility is a dipolar configuration with a satellite point defect, a monopole with $`N=1`$ near one of the particle poles luben , see fig.1(b). It appears that, although both situations are possible, the dipolar structure with satellite defect occurs more readily, e.g. in experiments luben ; poulin98 for water droplets dispersed in a thermotropic nematic matrix. When several such droplets come to a near-contact, the effect of chaining occurs, with satellite defects sandwiched between inclusion particles and shared between them, providing a strong and anisotropic interaction force.
On the other hand, when the strength of director anchoring is rather weak, the particle introduces only a small perturbation into the nematic matrix, see fig.1(c). Albeit small, the non-uniform distortion of n is spreading over long distances from the particle – the exact solution for the angle $`\theta `$ of director deviation from the uniform axis $`𝐧_\mathrm{o}`$ can be obtained in terms of the multipole expansion and reads, e.g. ukra (in spherical polar coordinates),
$$\theta (𝐫)=\frac{1}{4}\frac{WR}{K}\left(\frac{R}{r}\right)^3\mathrm{sin}\mathrm{\hspace{0.17em}2}\theta .$$
(1)
In this expression, $`R`$ is the particle radius and the nematic constants $`W`$ and $`K`$ are the anchoring energy on its surface and the Frank elastic constant respectively degen . Their typical values in a thermotropic nematic liquid crystal are $`K10^{11}\text{J/m}`$ and $`W10^6\text{J/m}^2`$, for homeotropic director anchoring as in fig.1.
The parameters in eq. (1) identify the important dimensionless ratio which controls the outcome of the single particle behaviour. $`(WR/K)`$ measures the relative effect of the particle surface (characteristic energy $`WR^2`$) and the bulk director deformations (Frank energy scale $`KR`$), and can be large or small. When $`WR/K1`$ the anchoring should be considered weak and not capable of producing large deformations in the surrounding nematic matrix. In contrast, when $`WR/K1`$ rigid anchoring can be assumed and topological singularities result. For typical parameters, particles of $`R510\mu `$m represent the border between the two topologically distinct regimes in fig.1: submicron colloid particles create small long range director distortions of quadrupolar symmetry, eq.(1), while larger objects most frequently have a dipolar configuration with a satellite defect.
The second physical idea relevant to our work is using the process of phase separation in a liquid crystal phase as a means of creating internal interfaces. The orientational symmetry breaking and the additional curvature elastic energy in the liquid crystal microconfined by such interfaces should have a profound effect on the whole process of phase ordering. Both phase equilibrium and kinetics are altered in non-trivial ways due to the underlying frustrated liquid crystalline order.
Although the progress in the single-particle description of liquid crystal colloids has been noticeable, little experimental work exists in this area of colloid collective behaviour. One of the main questions here is the macroscopic morphology of a colloid and the change in rheological behaviour on increasing concentration of particles, when the suspending matrix would possess a dense network of interaction forces creating high barriers to deformation. One may expect a glass-like freezing of motion at sufficiently high concentrations – an interesting possibility that needs to be thoroughly explored.
The work of Poulin et al. poulin98 on water droplets in a nematic matrix confirmed the ideas about topological defects and their long-range interactions – but they only studied the structure around a single or few very large droplets. In contrast, Tanaka et al. tanaka studied phase ordering in mixture of a nematic 5CB and a surfactant, which forms very small nano-micelles. They reported some provocative results on a phase they named “transparent nematic”. One could speculate that in this case the colloid parameter $`WR/K`$ is so small that the mean-field description of nematic order would fail altogether. Recently, Zapotocky et al. martin examined the colloid rheology of a cholesteric liquid crystal filled with silica particles of size $`1\mu `$m. They have found a variety of rich behaviour due to topological defect-driven enhancement of dynamic elastic properties and the formation of a weak gel with the estimated static modulus of $`G^{}0.01\text{Pa}`$. However, in their work small silica particles have aggregated into much bigger objects (flocs) that produced strong topological forces and ended up in the nodes of a network of disclination lines, which has been carrying the elastic function of the resulting material.
A simpler system has been discovered recently by Meeker et al. wilson . A “classical” liquid crystal colloid has been prepared by mixing a well-characterised thermotropic nematic liquid crystal, 5CB, with sterically-stabilised PMMA particles, imposing radial boundary conditions in the nematic matrix. Mixtures of $`510\%`$ particle volume fraction resulted in a soft solid with significant storage modulus, $`G^{}10^310^5\text{Pa}`$. Such a remarkable mechanical transformation requires further investigation, which is reported in the companion paper no2 . Bright field microscope observations by Meeker et al. suggest that the soft solid comprises a network of particle aggregates, formed by the exclusion of particles from emergent nematic domains as the mixture is cooled below the isotropic-nematic transition $`T_{\mathrm{ni}}`$. However, their optical observations were limited to low particle volume fractions $`(5\%)`$, and depended on reheating the samples back into the isotropic phase to view the particle aggregate structure, since direct optical study in the bulk of a highly non-uniform birefringent system is difficult due to the strong scattering of light: the nematic colloids are opaque below $`T_{\mathrm{ni}}`$. In this paper we use a confocal microscopy technique to study directly the structures in the bulk of the nematic liquid crystal colloid samples, for a wide range of particle concentrations $`(315\%)`$.
In addition to the remarkable mechanical response, the second aspect of interest in this system is the time required for the formation of an aggregated state below $`T_{\mathrm{ni}}`$. The characteristic time required for particle movement under the influence of long range attraction forces has been estimated in interac and gives, for our example of particle size and concentration in the limit of weak anchoring, $`\tau 10^3\eta (K/W^2)10^3\text{s}`$ (here $`\eta 0.1`$ Pa.s is the relevant viscosity coefficient of the nematic liquid crystal). We shall see that a rigid gel is formed in a much shorter time after the colloid mixture is brought below the nematic transition point $`T_{\mathrm{ni}}`$. It appears that the phenomena observed in our liquid crystal colloid system cannot be purely due to the topological defects and their networks: their elastic energy is not enough to explain the modulus and, more importantly, there should be no topological defects at small $`WR/K`$.
Here we study the formation process of the resulting near-solid liquid crystal colloid aggregates in more detail. We show that a process of continuous phase separation of small particles takes place in our system. This separation occurs immediately after the nematic transition and results in the formation of a metastable but very long-lived rigid cellular structure of very thin interfaces, where the particles are densely packed, encapsulating large volumes of nearly pure nematic liquid. We develop a mean-field theoretical model describing the phase stability and the mechanism of cellular solid formation.
## 2 Sample preparation and experimental methods
In order to investigate the robustness of the process discovered by Meeker et al. wilson , we study two materials: 5CB and MBBA, both archetypical thermotropic nematics lcbase . The case of $`4^{}`$-pentyl-$`4`$-cyanobiphenyl, abbreviated as 5CB, obtained from Aldrich Chemicals Co., has already been examined in wilson . This pure material has a stable nematic phase below $`T_{\mathrm{ni}}=35.8`$C, followed by a crystal phase at $`T_\mathrm{x}=22`$C. No crystallisation was observed down to $`5`$C when the colloid particles were added.
In contrast, the MBBA ($`N`$-$`4`$-methoxybenzylidene-$`4^{}`$-butylaniline, from Aldrich), although one of the most common nematic materials in the literature, has a practical complication of being susceptible to hydrolysis. The $`C=N`$ bond in the rigid molecular section dissociates in the presence of water. The reaction is reversible, meaning that at a given temperature and humidity there is an equilibrium balance of proper MBBA compound and its hydrolysis products, which act as impurities and weaken the nematic order. The effect is most noticeable in the depression of the phase transition points: in the “dry” MBBA the nematic transition point is at $`T_{\mathrm{ni}}^{(o)}=47`$C lcbase , followed by a crystal phase at $`T_\mathrm{x}=22`$C, while we have used MBBA in hydrolysis equilibrium at ambient conditions, with $`T_{\mathrm{ni}}=37.3`$C and no crystallisation down to $`5`$C. We employed this as a model for a “dirty”, weakly nematic liquid.
Theoretical arguments suggest that the best phase separation conditions are achieved for small colloid particles (the ones with $`WR/K1`$, but not too small so that the reversal of this relation can occur near $`T_{\mathrm{ni}}`$), see section 4 below. We use monodisperse PMMA spheres (polymethylmetacrylate, bulk glass transition at $`T_\mathrm{g}=105`$C) of radius $`R=150\text{nm}`$, with a polydispersity of about 0.04 (as determined by dynamic light scattering). Particles were covered with chemically grafted poly-$`12`$-hydroxystearic acid (PHSA) chains (prepared by Dr. A. Schofield in Edinburgh). On grafting, the short stabilising chains adopt a conformation radially extending from the grafting surface, making a “hairy” particle. In isotropic suspending liquids, these particles behave like almost ideal hard spheres. In addition, in the nematic suspending matrix, the grafted chains provide a homeotropic (radial) director anchoring with a typical strength $`W10^6\text{J/m}^2`$. With the typical Frank constant $`K10^{11}\text{J/m}`$ this makes the dimensionless colloid parameter to take the value $`WR/K0.02`$.
The preparation procedure described in wilson was followed. The liquid crystal was added to the dried particles at room temperature, i.e. while in the nematic phase (one does not expect good mixing in this situation). We then raise the temperature to well above $`T_{\mathrm{ni}}`$ and subject the sample to the ultrasonic excitation. In this way the particles were gradually dispersed in the isotropic phase of 5CB or MBBA; the samples were stored at $`T45`$C (above $`T_{\mathrm{ni}}`$) in a tumbling device to ensure that the mixtures were homogeneous before any experiment was started. The particle concentration $`\mathrm{\Phi }`$ has been measured by weight. Although the theoretical arguments require a proper (global average) volume fraction $`\mathrm{\Phi }=Nv_R/V`$, we consider the two adequately close because the densities of PMMA and of nematic liquid are not very different. Differential scanning calorimetry (DSC, Perkin-Elmer Pyris 1) was used to identify the phase behaviour of the resulting mixtures.
For this work the major advantage of the confocal optical microscope is that it can produce a three-dimensional image of a relatively thick and optically inhomogeneous sample. Ordinary polarised microscopy, commonly used in studies of liquid crystals, does not produce results here because the aggregated colloid samples strongly scatter light – the samples are opaque. The employment of a pinhole to eliminate the out-of-focus light from planes above and below the focal plane in a confocal microscope produces a sharp picture of optical contrast in this plane. We use a Laser Scanning Confocal Microscope (LSCM 510, by K. Zeiss) in the reflection mode which does not require fluorescent labelling, illuminated by the monochromatic $`540`$ nm laser, without polarisers. In this mode, the images could be collected at a depth of up to $`100\mu `$m under the top sample surface.
## 3 Experimental results
### 3.1 Calorimetry
The starting point of practically any theoretical description of thermotropic nematic colloid phase behaviour, e.g. roux , is the assumption that the nematic transition temperature in the homogeneously mixed colloid is a linearly decreasing function of particle concentration $`\mathrm{\Phi }`$. In fact, below we shall explicitly calculate the slope of this linear dependence: $`T^{}\left(1\alpha \mathrm{\Phi }\right)`$, eq. (3). Of course, such a linear variation means that the particles are regarded as fully independent, non-interacting – which cannot be true for high concentrations. Also, this estimate for $`T_{\mathrm{ni}}(\mathrm{\Phi })`$ is based on a value reversal of the colloid parameter $`WR/K`$ near the transition. This requires a subtle balance of material constants. If $`WR/K`$ remains large deep in the nematic phase, the particles will quickly aggregate under the action of strong attraction forces, see poulin98 ; martin . If $`WR/K`$ remains small even near $`T_{\mathrm{ni}}`$ in spite of its tendency to grow as the nematic order parameter $`Q(T)`$ decreases (cf. section 4), the particles will not have a strong elastic energy around them and the system will behave more like a molecular mixture, a nematic with impurities – a well studied subject with a different phase morphology mixtures .
In order to determine the coupling energy between the nematic field and the individual colloid particle we study the isotropic-nematic phase transition on cooling from a homogeneously mixed state with a small particle volume fraction $`\mathrm{\Phi }`$. The differential calorimetry shows an exothermal latent heat peak at this weakly first-order transition. Figure 2 shows the evolution of this peak on increasing the particle concentration, in 5CB and MBBA colloids. Two main features have to be noted. First, the dependence of $`T_{\mathrm{ni}}`$ (determined by the peak onset) on concentration shows a monotonic trend, fig. 3. The linear fits appear to be quite good, with the slope consistent with theoretical estimates of section 4. Secondly, one finds that the peaks become broader and of lower total area on increasing the particle concentration. This is not unexpected, the first order transition should become more and more diffuse with added impurities. In fact, there is a theoretical view aw ; cardy that such a transition should become continuous in equilibrium. We did not address this question in detail, which would have required significantly lower DSC cooling rates, at the very least.
### 3.2 Cellular structure
Here we present the confocal images of horizontal cross-sections of aggregated nematic colloids. The samples studied here were prepared by simply depositing a portion of isotropic homogeneous mixture (at $`45`$C) on a slide and allowing it to cool to the room temperature, without top cover. The contrast mechanism between the nearly pure nematic regions and the densely aggregated PMMA particles is due to a different average refractive index in the reflection mode of the microscope.
The cellular morphology of aggregated nematic colloids is apparent in fig. 4 for the 5CB system and in fig. 5 for the MBBA system. The relatively large cavities are separated by thin interfaces providing an optical contrast. The proportion of volume between these two fractions indicates that the cavities are filled with the nematic liquid (the majority phase) and the interfaces are made of PMMA particles.
The sizes of cells appear very irregular. One of the reasons is that the thickness of confocal imaging plane is much smaller than the cells and one obtains their different cross-sections. However, it is also clear that cells are very polydisperse both in size and in shape. One may argue that the way of sample preparation has caused the polydispersity: the relatively slow cooling under ambient conditions passes through the binodal regime of nucleation and growth in the phase diagram, see later (fig. 9), and would result in a variety of sizes of growing nematic nuclei. This should be contrasted with a rapid cooling deep into the spinodal decomposition regime, where a characteristic size is selected. We shall see later that samples cooled in the rheometer (at a rate of $`30^\mathrm{o}/\text{min}`$, with an additional small vibration aiding interface packing) have a significantly more regular cellular superstructure.
Another fact, evident in the confocal images, figs. 4 and 5, is that the cell walls are perforated. This is especially evident in the image (b) for the 7% colloid, where most of the walls are continuous in the cross-section, but some are clearly interrupted. This observation brings our system into the class of “open-cell” structures, cf. ashby . Figure 6 shows the microstructures of a free surface of aggregated cellular nematic colloid prepared in the ambient conditions. The difference with the internal cellular morphology is striking. Nevertheless, in spite of a much more regular surface, representing an effective cell wall meeting the outer interface, the perforations of the interface are evident.
Image analysis, which should provide the statistical distribution of cell sizes, is difficult for the confocal scans of different cell cross-sections. Even if the cells were monodisperse, their areas crossing the focal plane would appear broadly distributed. Although there is a clear decrease in cell size with increasing particle concentration in fig. 4, the distribution is too irregular and asymmetric to allow accurate quantitative conclusions.
## 4 Theoretical model
First of all, let us re-iterate that the proposed regime of phase separation as a mechanism of particle aggregation in the nematic phase is taking place due to the smallness of the liquid crystal colloid parameter $`WR/K`$. In the opposite case, the study of which has been attempted several times over the years, the topological defects around particles result in strong interaction forces. If there are only few large particles in the system, they often form “strings” – chains of alternating spheres and topological defects luben ; poulin98 . When the particle concentration is finite, the colloid undergoes fast aggregation into 3-dimensional flocs. Big flocs with stronger relative anchoring (planar on silica surface, as in martin ) fall into the category of $`WR/K1`$. An intermediate stage of evolution is the network of disclinations connecting the flocs and having an effective modulus $`G^{}0.01\text{Pa}`$. The topological defect network would eventually reach the global equilibrium by clearing the nematic volume altogether. (This final state is the eventual fate of all other metastable modes of aggregation, including our cellular solids).
We now return to the case of weak anchoring, $`WR/K1`$, and consider the behaviour of a system of small particles near the nematic transition, where the amplitude of the nematic order parameter $`Q(T)`$ becomes small (the N-I transition is weakly first order, so that the latent heat and the jump of $`Q`$ at the transition are small). Both relevant elastic parameters, the Frank constant $`K`$ and the anchoring energy $`W`$ vanish as $`Q0`$, but with different rates: $`K\kappa Q^2`$, while $`WwQ`$ in the first approximation. Therefore, in the immediate vicinity of $`T_{\mathrm{ni}}`$ the ratio $`WR/K1/Q`$ should become large and the topological regime of nematic director distortions around even very small particles may prevail. This would certainly be the case if a critical second-order transition occurred at $`T_{\mathrm{ni}}`$. As it is, since $`Q`$ never reaches zero, only a certain combination of material parameters would allow $`WR/K`$ to change from small to large.<sup>1</sup><sup>1</sup>1If $`WR/K`$ remains small throughout the nematic phase, the elastic energy of deformations, $`W^2R^3/K`$ ukra , has only a weak dependence on the nematic order parameter $`Q`$ and results in long aggregation times. The elastic energy of these distortions is, near the N-I transition, $`10KR=10\kappa RQ^2`$. This makes an addition to the Landau free energy density describing the nematic transition, for local particle concentration $`\varphi `$ it takes the form $`10KR\varphi /v_\mathrm{R}`$, with $`v_\mathrm{R}=\frac{4}{3}\pi R^3`$ the particle volume, giving
$$F_\mathrm{n}=\frac{1}{2}A_\mathrm{o}\left[TT_c(\varphi )\right]Q^2\frac{1}{3}BQ^3+\frac{1}{4}CQ^4$$
(2)
Here the shift in the critical point is due to the elastic energy around the particles:
$$T_c(\varphi )=T^{}\left(1\frac{10\kappa R\varphi }{A_\mathrm{o}T^{}v_\mathrm{R}}\right)T^{}\left(1\frac{\xi ^2}{R^2}\varphi \right)$$
(3)
where the second, very approximate, equation uses the fact that the nematic correlation length $`\xi `$ is actually given by the ratio of parameters $`\kappa `$ and $`A_\mathrm{o}`$, see Appendix. The experimental data in fig. 3 show that the slope of $`T_c(\mathrm{\Phi })`$ is indeed noticeable. The experimental estimate obtained from the linear fit is not wholly unreasonable, albeit slightly larger than one could expect for our particles and a typical correlation length $`\xi 10\text{nm}`$.
### 4.1 Phase separation below $`T_{\mathrm{ni}}`$
Immediately after the transition, the growing nematic nuclei start expelling the particles in the drive to reduce its thermodynamic free energy (2). The initial average particle volume fraction $`\mathrm{\Phi }`$ splits into locally different valuesin phase separating regions; we now use the notation $`\varphi `$ for the local concentration. The equilibrium nematic order parameter of the mixed phase is
$$Q^{}=\frac{B}{2C}\left(1+\sqrt{1+\frac{4(A_\mathrm{o}T^{})C}{B^2}[1\tau \alpha \varphi ]}\right)$$
(4)
with the reduced temperature $`\tau =T/T^{}`$ and the coefficient $`\alpha 0.1`$, using the data of fig. 3. Accordingly, the upper limit of stability of homogeneously mixed nematic phase is at the concentration
$$\varphi _\mathrm{N}=\frac{1}{\alpha }\left(\frac{B^2}{4(A_\mathrm{o}T^{})C}+1\tau \right).$$
(5)
The equilibrium free energy density $`F_\mathrm{n}(Q^{})`$ then becomes a function of local particle concentration $`\varphi `$, see eq. (4.1) below. The principal feature of this optimised free energy density of a nematic mean field is the rapid increase of $`F_\mathrm{n}(\varphi )`$ with increasing particle concentration and the non-convex form of this dependence. This feature is the driving force for phase separation.
This energy penalty on having the particles uniformly dispersed in the nematic matrix has to be added to the free energy of mixing. The Carnaham-Starling excess free energy density of a hard-sphere suspension fluid is a very good approximation of the equation of state:
$`F_\mathrm{p}`$ $`=`$ $`{\displaystyle \frac{kT}{v_R}}\left(\varphi \mathrm{ln}\varphi +\varphi ^2{\displaystyle \frac{43\varphi }{(1\varphi )^2}}\right),\varphi <\varphi ^{},\mathrm{liquid}`$ (6)
$`F_\mathrm{p}`$ $`=`$ $`{\displaystyle \frac{kT}{v_R}}(1.79\varphi +3\varphi \mathrm{ln}{\displaystyle \frac{\varphi }{1\varphi /\varphi _c}},),\varphi >\varphi ^{},\mathrm{solid}`$
where $`v_R`$ is the particle volume and the two expressions are matched by an adjustable parameter $`u1.7929`$ at a concentration $`\varphi ^{}0.52`$, with $`\varphi _c0.64`$ the random close-packing fraction. The full free energy of nematic colloid, $`F_\mathrm{n}+F_\mathrm{p}`$, then becomes, in dimensionless form,
$`{\displaystyle \frac{v_R}{kT}}F={\displaystyle \frac{a_2}{\tau }}(1\varphi )[(1\tau \alpha \varphi )^2`$
$`+`$ $`{\displaystyle \frac{b_2^2}{6}}(1+{\displaystyle \frac{6}{b_2}}(1\tau \alpha \varphi )+[1+{\displaystyle \frac{4}{b_2}}(1\tau \alpha \varphi )]^{3/2})]`$
$`\varphi \mathrm{ln}\varphi +\varphi ^2{\displaystyle \frac{43\varphi }{(1\varphi )^2}},`$
at small concentrations ($`\varphi <\varphi ^{}`$), where the dimensionless parameters $`a_21`$ and $`b_21`$ are estimated in Appendix, eq. (11). The factor $`(1\varphi )`$ in front of the nematic mean field energy $`F_\mathrm{n}(Q^{})`$ measures the local proportion of the solvent in the colloid mixture.
This is the simplest possible model of nematic colloid phase ordering, essentially the same as used in roux , with minor deviations in the interpretation of nematic order-composition coupling and material constants. This interpretation, however, determines the results and predictions. All parameters entering the free energy $`F(\varphi )`$ will be determined from experiment, thus providing the phase diagram in the temperature-composition plane with no adjustable free parameters.
At low initial average particle volume fraction, the structure of eq. (4.1) is reminiscent of the Flory-Huggins model for demixing. The nematic mean-field energy plays the role of the effective $`\chi `$-parameter potential term in that model, providing non-convex variation in $`F(\varphi )`$. Figure 8 shows a series of plots of eq. (4.1) for the set of Landau parameters corresponding to 5CB nematic liquid crystal, eqs. (10) and (11) in Appendix. There is always a well-defined minimum at very small concentrations,
$$\varphi _0\mathrm{exp}\left[a_2\frac{(1\tau )(1\tau +2\alpha )}{\tau }\right],$$
due to the entropy contribution $`(\varphi \mathrm{ln}\varphi )`$ in eq. (6). The usual common-tangent construction for the $`F(\varphi )`$ connects two branches of the free energy: one non-convex, at very low $`\varphi `$, in an almost pure liquid crystal, and the other fully convex, in the high-$`\varphi `$ phase where the nematic order is unstable $`(Q=0)`$ and the particles are densely compacted.
The schematic phase diagram of the nematic colloid is shown in fig. 9 in coordinates of reduced temperature $`\tau `$ and concentration $`\varphi `$. At high temperature the colloid behaves as a standard athermal hard-sphere system, which crystallises above $`\varphi 0.5`$ and reaches the maximum random close-packing level at $`\varphi _c0.64`$. At an extremely low particle concentration, corresponding to the expanded insert in fig. 8, there is a region of homogeneous nematic mixture, labelled NL in the phase diagram. This region is artificially expanded in fig. 9, the real numbers make it invisible in the plot. A very good analytical approximation for this phase boundary is
$$\varphi _1(\tau )\frac{\tau }{1\tau +\frac{1}{2}\alpha }\left(\frac{1}{a_2\alpha }\right)10^6\tau \mathrm{}$$
(8)
The spinodal on the high-$`\varphi `$ branch is determined by the point of contact between the nematic (concave, locally unstable) and the isotropic (fully convex) parts of the free energy, which is marked as $`\varphi _2(\tau )`$ in fig. 8. Its analytical estimate,
$$\varphi _2=\frac{1}{\alpha }\left(\frac{2}{9}b_2+1\tau \right),$$
is practically indistinguishable from the stability limit $`\varphi _\mathrm{N}`$ in eq. (5). This expresses the known fact about the “weakness” of first order nematic transition: the width of nematic coexistence is very small in ordinary thermotropic liquid crystals. Below the spinodal line $`\varphi _2(\tau )`$ the homogeneous particle mixture is unstable. Above this line, the shaded region in fig. 9 shows the zone of coexistence between the low-$`\varphi `$ (nematic) and the high-$`\varphi `$ (isotropic) phases. The upper boundary of this region, the high-concentration binodal line is calculated numerically from the common-tangent condition for $`F(\varphi )`$ in fig. 8, for the set of parameters given by eq. (11) (thin solid line) and the artificially scaled down nematic energy (bold solid line) to expose the critical point, where this binodal merges with the low-$`\varphi `$ phase boundary $`\varphi _1(\tau )`$.
Note that the high-concentration phase is in the state of colloidal solid. This is a consequence of very large relative strength of nematic mean field, expressed by the large dimensionless parameter in eq. (11). One may say that the effective pressure from the nematic liquid to expel the particles and compact them in high-$`\varphi `$ regions is rather high.
A word of caution has to be sounded at this point, regarding the continuing usage of the linear law $`T_c(\varphi )`$ at relatively high particle concentrations. This approximation is not likely to be valid in the regime of high interparticle interactions and, therefore, the spinodal $`\varphi _2(\tau )`$ would deflect down from the straight line in fig. 9. We, however, are mostly interested in small particle concentrations where the approximation of independent particles would hold.
### 4.2 Cellular structure
We thus envisage the following mechanism of phase evolution of the nematic colloid on its cooling from the isotropic homogeneously mixed phase:
* After quenching the homogeneous colloid suspension below the nematic transition point, the system phase separates into growing regions of pure nematic liquid, from which the particles are expelled into the boundaries.
* On these interfaces of growing nematic nuclei, the concentration of particles is so high that the colloid solidifies and the remaining mesogenic liquid is in the isotropic state.
Clearly, the fine balance of kinetic effects is essential: the phase separation should have a higher rate to develop within growing nematic nuclei. The driving force for this process is the thermodynamic free energy gain in having the high nematic order in a clean system at a given $`T<T_{\mathrm{ni}}`$, in comparison with the uniformly dispersed nematic colloid, which is expressed by non-convexity of plots in fig. 9.
Eventually, when the clean nematic regions grow to come into the near-contact with each other, the particles are only allowed to pack on narrow interfaces. The sample transforms into a cellular superstructure with rather thin densely-packed walls encapsulating nematic volumes of characteristic size $`\lambda `$. This cell size is determined by the initial average colloid concentration $`\mathrm{\Phi }`$ and the thickness of interfaces. Let us define this thickness as $`dnR`$, several times the particle radius (the smallest possible is $`n2`$ in fig.10). Then the total area of such interface in the sample volume $`V`$ is $`𝒜=\mathrm{\Phi }V/(nR)`$. This leads to the order of magnitude estimate for the mesh size
$$\lambda =(V/𝒜)\frac{nR}{\mathrm{\Phi }}.$$
(9)
This gives $`\lambda 6\mu \text{m}`$ for this perhaps unrealistic case of $`n=2`$, and $`R=150\text{nm}`$ and $`\mathrm{\Phi }=0.05`$. The analysis of structures in the next section obtains $`n20`$, see fig. 7, i.e. interfaces of $`10`$ particles thick and $`\lambda 60\mu \text{m}`$ for $`\mathrm{\Phi }=0.05`$.
The cellular solid with thin high-tension interfaces is not the globally equilibrium state. It is obvious that the free energy can be further lowered by aggregating all particles into a 3-dimensional densely packed group, leaving the whole volume for the pure nematic liquid phase. However, the metastable cellular structure may become frozen, trapped by high energy barriers for disrupting the formed interfaces. Several mechanisms contribute to creating such barriers. Possibly the largest contribution arises from the nematic topological argument similar to that in liquid crystal emulsions emuls95 . Consider the nematic alignment in one cell, regarding the interface as a continuous wall imposing a homeotropic director anchoring energy $`W`$. When the mesh size $`\lambda `$ increases so that the familiar dimensionless parameter $`(W\lambda /K)1`$, the nematic liquid crystal within this cell has to possess the topological charge $`N=1`$. When one attempts to break through the cell interface, two topologically charged volumes come into contact and their coalescence must comply with the law of this charge conservation. Therefore, in the moment of formation of a channel between the cells, a new topological defect of the charge $`N=1`$ must be created emuls95 . Later this defect may travel towards one of existing monopoles and annihilate it to minimise the elastic energy of the final joint volume. However, this cannot happen instantly and the formation of each new topological defect costs an elastic energy $`K\lambda `$. This energy barrier can be very high: $`10^{16}\text{J}`$, compared with $`k_\mathrm{B}T4\times 10^{21}\text{J}`$, and the system may remain trapped in the metastable state of random cellular structure. One then expects to find the effects of ageing and time translation invariance breaking, characteristic of weakly non-ergodic glassy dynamics and rheology bouchaud ; sollich .
## 5 Conclusions
We have reported the results of a structural study of nematic colloids based on a classical thermotropic liquid crystals matrix with small monodisperse polymer particles suspended in it wilson . The hydrophobic sterical stabilisation ensures the radial anchoring of the director on particle surface. At relatively small concentrations of particles, we observe good mixing above the clearing point of isotropic-nematic phase transition $`T_{\mathrm{ni}}`$ and a rapid aggregation of the homogeneous mixture into a rigid gel-like solid, completely opaque optically at and below $`T_{\mathrm{ni}}`$. The properties of the phase ordering and the morphology of the phase-separated aggregates were the primary focus of this study.
On cooling from the homogeneous isotropic mixture, we observe a decrease in the transition temperature $`T_{\mathrm{ni}}`$ as a function of average colloid concentration $`\mathrm{\Phi }`$, which follows a reasonably linear law in the region of small concentrations studied here. This is an expected and frequently observed phenomenon, accompanied by a noticeable diffusion of the weak first-order phase transition, which is otherwise sharp in a pure nematic liquid crystal.
Below $`T_{\mathrm{ni}}`$ the colloids undergo phase separation with the resulting structure best described as an open cellular solid, with the particles densely packed in thin walls and the cavities filled with a pure liquid crystal. The characteristic cell size is of the order of $`10100\mu `$m and appears to follow the inverse-proportionality law $`\lambda 1/\mathrm{\Phi }`$, also confirmed by a simple theoretical estimate. Since the nematic director is anchored on cell walls, the resulting randomly quenched birefringent texture strongly scatters light, giving the material its opaque appearance. The remarkable mechanical rigidity of the resulting cellular solids, stemming from the high effective surface tension of interfaces, is the subject of the companion paper no2 .
Colloid suspensions in nematic liquid crystal matrices, thermotropic and lyotropic, have been studied before. Particle aggregation is an unavoidable result in all cases, since the system below $`T_{\mathrm{ni}}`$ tends to minimise the elastic energy of director distortions around individual particles. However, the cellular solid morphology, first observed in wilson , is new. We believe that two main factors contribute to the reason why we obtain such an effect. The particles have to be small and have sufficiently low anchoring energy to ensure that the colloid parameter $`WR/K`$ is small just below $`T_{\mathrm{ni}}`$. On the other hand, our theoretical model suggests that this parameter has to become substantial again, to provide a required thermodynamic force for phase separation. This rather restricts a range of particle sizes and surface treatments that allow for such a compromise to occur. Secondly, we find that the rate of cooling through the first-order nematic transition has an important effect. The cellular structure is much more regular and robust when the homogeneous isotropic colloid is quenched rapidly deep below $`T_{\mathrm{ni}}`$. The theoretical phase diagram suggests that in order to achieve a good selection of size one needs to avoid a binodal region just below $`T_{\mathrm{ni}}`$. In many cases of slow-rate cooling and a colloid material parameters making the binodal gap wider (such as in our example of MBBA system), the mixture would separate in a different fashion – perhaps totally avoiding the cellular solid regime.
Although the details of aggregation mechanism and forces that hold particles together at interfaces are somewhat unclear, the suggested theoretical “toy model” offers a possible explanation for the formation of cellular structure. Still, many questions remain open. Such an unusual phase behaviour and remarkable rheological properties of the liquid crystal colloid suspensions require further detailed study, both theoretical and experimental.
We appreciate valuable discussions with M.E. Cates, S.M. Clarke, P.D. Olmsted and M. Warner. The help of I. Hopkinson with the confocal microscope is gratefully acknowledged. This research has been supported by the UK EPSRC.
## Appendix: <br>Parameters of Landau free energy
It is useful to find the values of the three phenomenological parameters describing the transition, coefficients $`A_\mathrm{o},B`$ and $`C`$ in eq. (2), although different materials will have these parameters slightly different. It is nevertheless instructive to examine the characteristic orders of magnitude. In order to determine three parameters one needs three independent measurements. We are fortunate that there are, in fact, four available:
1) The jump of order parameter at the weak first-order transition is $`\mathrm{\Delta }Q_{\mathrm{ni}}=2B/3C`$. There might be some error in its determination, which should depend (among other factors) on the rate of cooling through the transition. However, because $`Q`$ only varies between $`0`$ and $`1`$ and because a number of molecular theories predict this jump explicitly, one can take qualitatively $`Q_{\mathrm{ni}}0.4`$, giving $`B0.6C`$.
2) The second measurement can be the interval between the transition temperature and the critical point $`T^{}`$ (the latter may be determined by extrapolating the inverse susceptibility, $`\chi ^1|TT^{}|`$). In usual thermotropic nematic liquid crystals this interval is rather small, $`T_{\mathrm{ni}}T^{}=2B^2/9A_\mathrm{o}C1^\mathrm{o}`$. This gives $`B7.5A_\mathrm{o}1^\mathrm{o}`$K. (Note that the dimensionality of $`B`$ and $`C`$ is energy density, while $`A_\mathrm{o}`$ has the dimensions of $`\text{J/m}^3{}_{}{}^{\mathrm{o}}\text{K}`$).
3) The latent heat of the first-order phase transition is $`\mathrm{\Delta }H=T_{\mathrm{ni}}(2A_\mathrm{o}B^2/9C^2)`$, per unit volume. It is typically obtained from calorimetry by integrating the characteristic peak. Again, a large uncertainty may accompany such a measurement because at any non-zero cooling rate the conditions are not exactly equilibrium. However, keeping up the qualitative approach, the value of transition enthalpy per unit mass is of the order $`1\text{J/g}`$ for 5CB at a reasonably slow cooling rate. Taking the density $`\rho 1\text{g/cm}^3`$ and the $`T_{\mathrm{ni}}310^\mathrm{o}`$K, we obtain
$`A_\mathrm{o}`$ $``$ $`6\times 10^3\text{J/m}^3{}_{}{}^{\mathrm{o}}\text{K}`$ (10)
$`\mathrm{Then}B`$ $``$ $`5\times 10^4\mathrm{and}C1.2\times 10^5\text{J/m}^3.`$
4) It is important that the three measurements above give estimates that agree with the fourth way of accessing $`A_\mathrm{o}`$. The nematic correlation length $`\xi `$ may be determined by the ratio of the bare Frank elastic constant $`\kappa =K/Q^2`$ to the thermodynamic energy density: $`\xi ^2=\kappa /A_\mathrm{o}\mathrm{\Delta }T`$. There are many ways of confirming that the characteristic magnitude of $`\xi `$ is $`10\text{nm}`$. For usual values $`\kappa 10^{11}\text{J/m}`$ and $`\mathrm{\Delta }T10^\mathrm{o}`$ we obtain $`A_\mathrm{o}10^4\text{J/m}^3{}_{}{}^{\mathrm{o}}\text{K}`$, in a very reasonable agreement with the previous estimate.
Finally, we may estimate the dimensionless factors entering the final expression for the free energy (4.1). At $`T^{}309^\mathrm{o}`$K ($`kT^{}4.7\times 10^{21}\text{J}`$) and particle volume $`v_R1.4\times 10^{20}\text{m}^3`$ we obtain
$$a_2=\frac{A_\mathrm{o}^2T_{}^{}{}_{}{}^{2}v_R}{4CkT^{}}2\times 10^7\mathrm{and}b_2=\frac{B^2}{A_\mathrm{o}T^{}C}10^2.$$
(11)
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# Connections between Linear Systems and Convolutional Codes
## 1 Introduction
It is common knowledge that there is a close connection between linear systems over finite fields and convolutional codes. In the literature one finds however a multitude of definitions for convolutional codes, which can make it confusing for somebody who wants to enter this research field with a background in systems theory or symbolic dynamics. It is the purpose of this article to provide a survey of the different points of view about convolutional codes.
The article is structured as follow: In Section 2 we will review the way convolutional codes have often been defined in the coding literature .
Section 3 reviews a definition for convolutional codes that can be found in the literature on symbolic dynamics. From the symbolic dynamics point of view , a convolutional code is a linear irreducible shift space.
In Section 4 we will review the class of time-invariant, complete linear behaviors in the sense of Willems . We will show how these behaviors relate to the definitions given in Section 2 and 3.
In Section 5 we will give a definition for convolutional codes in which it is required that the code words have finite support. Such a definition was considered by Fornasini and Valcher and by the author in collaboration with Schumacher, Weiner and York . The study of behaviors with finite support has been done earlier in the context of automata theory and we refer to Eilenberg’s book . We show in Section 5 how this module-theoretic definition relates to complete, linear and time-invariant behaviors by Pontryagin duality.
In Section 6 we will study different first-order representations connected with the different viewpoints. Finally, in Section 7 we compare the different definitions. We also show how cyclic redundancy check codes can naturally be viewed in the context of finite-support convolutional codes.
Throughout the paper we will emphasize the algebraic properties of the different definitions. We will also restrict ourselves to the concrete setting of convolutional codes defined over finite fields. It is however known that many of the concepts in this paper generalize to group codes and multidimensional convolutional codes . All of the definitions which we are going to give are quite similar, but there are some notable differences.
Since the paper draws from results from quite different research areas, one is faced with the problem that there is no uniform notation. In this paper we will adopt the convention used in systems theory in which vectors are regarded as column vectors. For the convenience of the reader, we conclude this section with a summary of some of the notation used in this paper:
| $`𝔽`$ | A fixed finite field; |
| --- | --- |
| $`𝔽[z]`$ | The polynomial ring over $`𝔽`$; |
| $`𝔽[z,z^1]`$ | The Laurent polynomial ring over $`𝔽`$; |
| $`𝔽(z)`$ | The field of rationals; |
| $`𝔽[[z]]`$ | The ring of formal power series of the form $`_{i=0}^{\mathrm{}}a_iz^i`$; |
| $`𝔽((z))`$ | The field of formal Laurent series having the form $`_{i=d}^{\mathrm{}}a_iz^i`$; |
| $`𝔽[[z,z^1]]`$ | The ring of formal power series of the form $`_{i=\mathrm{}}^{\mathrm{}}a_iz^i`$; |
| $``$ | The integers; |
| $`_+`$ | The nonnegative integers; |
| $`_{}`$ | The nonpositive integers. |
Consider the ring of formal power series $`𝔽[[z,z^1]]`$. We will identify the set $`𝔽[[z,z^1]]`$ with the (two-sided) sequence space $`𝔽^{}`$. We have natural embeddings:
$$𝔽𝔽[z]𝔽[z,z^1]𝔽(z)𝔽((z))𝔽[[z,z^1]].$$
With these embeddings we can view e.g. the set of rationals $`𝔽(z)`$ as a subset of the sequence space $`𝔽^{}`$, and we will make use of such identifications throughout the paper.
The set of $`n`$-vectors with polynomial entries will be denoted by $`𝔽^n[z]`$. Similarly we define the sets $`𝔽^n(z),𝔽^n((z))`$ etc. All these sets are subsets of the two sided sequence space $`\left(𝔽^n\right)^{}=𝔽^n[[z,z^1]]`$. The definitions of convolutional codes which we will provide in the next sections will all be $`𝔽`$-linear subspaces of $`\left(𝔽^n\right)^{}`$.
The idea of writing a survey on the different points of view about convolutional codes was suggested to the author by Paul Fuhrmann during a stimulating workshop on “Codes, Systems and Graphical Models” at the Institute for Mathematics and its Applications (IMA) in August 1999. A first draft of this paper was circulated in October 1999 to about a dozen people interested in these research issues. This generated an interesting ‘Internet discussion’ on these issues, in which the different opinions were exchanged by e-mail. Some of these ideas have been incorporated into the final version of the paper and the author would like to thank Dave Forney, Paul Fuhrmann, Heide Gluesing-Luerssen, Jan Willems and Sandro Zampieri for having provided valuable thoughts. The author wishes also to thank the IMA and its superb staff, who made the above mentioned workshop possible.
## 2 The linear algebra point of view
The theory of convolutional codes grew out and extended the theory of linear block codes into a new direction. Because of this reason we start the section with linear block codes and we introduce convolutional codes in a quite intuitive way.
An $`[n,k]`$ linear block code is by definition a linear subspace $`𝒞𝔽^n`$ having dimension $`dim𝒞=k`$. Let $`G`$ be a $`n\times k`$ matrix with entries in $`𝔽`$. The linear map
$$\phi :𝔽^k𝔽^n,mc=Gm$$
is called an encoding map for the code $`𝒞`$ if $`\mathrm{im}(\phi )=𝒞`$. If this is the case then we say $`G`$ is a generator matrix or an encoder for the block code $`𝒞`$.
Assume that a sequence of message blocks $`m_0,\mathrm{},m_t𝔽^k`$ should be encoded into a corresponding sequence of code words $`c_i=Gm_i𝔽^n,i=0,\mathrm{},t`$. By introducing the polynomial vectors $`m(z)=_{i=0}^tm_iz^i𝔽^k[z]`$ and $`c(z)=_{i=0}^tc_iz^i𝔽^n[z]`$ it is possible to describe the encoding procedure through the module homomorphism:<sup>2</sup><sup>2</sup>2Throughout the paper we use the symbol $`\phi `$ to denote an encoding map. The context will make it clear what the domain and the range of this map is in each situation.
$$\phi :𝔽^k[z]𝔽^n[z],m(z)c(z)=Gm(z).$$
(2.1)
The original idea of a convolutional code goes back to the paper of Elias , where it was suggested to use a polynomial matrix $`G(z)`$ in the encoding procedure $`(`$2.1$`)`$.
Polynomial encoders $`G(z)`$ are physically easily implemented through a feedforward linear sequential circuit. Massey and Sain showed that there is a close connection between linear systems and convolutional codes. Massey and Sain viewed the polynomial encoder $`G(z)`$ as a transfer function. More generally it is possible to realize a transfer function $`G(z)`$ with rational entries by (see e.g. ) a linear sequential circuit whose elements include feedback components. If one allows rational entries in the encoding matrix then it seems natural to extend the possible message sequences to the set of rational vectors $`m(z)𝔽^k(z)`$ and to process this sequence by a ‘rational encoder’ resulting again in a rational code vector $`c(z)𝔽^n(z)`$. With this we have a first definition of a convolutional code as it can be found e.g. in the Handbook of Coding Theory \[35, Definition 2.4\]:
###### Definition A
A $`𝔽(z)`$-linear subspace $`𝒞`$ of $`𝔽^n(z)`$ is called a convolutional code.
If $`G(z)`$ is a $`n\times k`$ matrix with entries in $`𝔽(z)`$ whose columns form a basis for $`𝒞`$, then we call $`G(z)`$ a generator matrix or an encoder for the convolutional code $`𝒞`$. $`G(z)`$ describes the encoding map:
$$\phi :𝔽^k(z)𝔽^n(z),m(z)c(z)=G(z)m(z).$$
The field of rationals $`𝔽(z)`$ viewed as a subset of the sequence space $`𝔽^{}=𝔽[[z,z^1]]`$ consists precisely of those sequences whose support is finite on the negative sequence space $`𝔽^{_{}}`$ and whose elements form an ultimately periodic sequence on the positive sequence space $`𝔽^_+`$. It therefore seems that one equally well could restrict the possible message words $`m(z)𝔽^k(z)`$ to sequences whose coordinates consists of Laurent polynomials only, in other words to sequences of the form $`m(z)𝔽^k[z,z^1]`$.
Alternatively one could allow message words $`m(z)`$ whose coordinates are not ultimately periodic and possibly not of finite support on the negative sequence space $`𝔽^{_{}}`$. This would suggest that one should take as possible message words the whole sequence space $`\left(𝔽^k\right)^{}=𝔽^k[[z,z^1]]`$. The problem with this approach is that the multiplication of an element in $`𝔽[[z,z^1]]`$ with an element in $`𝔽(z)`$ is in general not well defined. If one restricts however the message sequences to the field of formal Laurent series then the multiplication is well defined. This leads to the following definition which goes back to the work of Forney . The definition has been adopted in the book by Piret and the book by Johannesson and Zigangirov , and it appears as Definition 2.3 in the Handbook of Coding Theory :
###### Definition A
A $`𝔽((z))`$-linear subspace $`𝒞`$ of $`𝔽^n((z))`$ which has a basis of rational vectors in $`𝔽^n(z)`$ is called a convolutional code.
The requirement that $`𝒞`$ has a basis with rational entries guarantees that $`𝒞`$ has also a basis with only polynomial entries. $`𝒞`$ can therefore be represented by a $`n\times k`$ generator matrix $`G(z)`$ whose entries consist only of rationals or even even polynomials. The encoding map with respect to $`G(z)`$ is given through:
$$\phi :𝔽^k((z))𝔽^n((z)),m(z)c(z)=G(z)m(z).$$
(2.2)
If $`G(z)`$ is a polynomial matrix, then finitely many components of $`m(z)`$ influence only finitely many components of $`c(z)`$, and the encoding procedure may be physically implemented by a simple feedforward linear shift register.
If $`G(z)`$ contains rational entries, then it is in general the case that a finite (polynomial) message vector is encoded into an infinite (rational) code vector of the form $`c(z)=_{i=s}^{\mathrm{}}c_iz^i`$. This might cause some difficulties in the decoder. For the encoding process, $`G(z)`$ can be physically realized by linear shift registers, in general with feedback (see e.g. ).
From a systems theory point of view, it is classical to view the encoding map $`(`$2.2$`)`$ as an input-output linear system. This was the point of view taken by Massey and Sain and thereafter in most of the coding literature. However unlike in systems theory, the important object in coding theory is the code $`𝒞=\mathrm{im}(\phi )`$. As a result one calls encoders $`\phi `$ which generate the same image $`\mathrm{im}(\phi )`$ equivalent; we will say more about this in a moment. In Sections 3 and 4 we will view $`(`$2.2$`)`$ as an image representation of a time-invariant behavior in the sense of Willems , which we believe captures the coding situation in a more natural way.
Assume that $`G(z)`$ and $`\stackrel{~}{G}(z)`$ are two $`n\times k`$ rational encoding matrices defining the same code $`𝒞`$ with respect to either Definition A or A. In this case we say that $`G(z)`$ and $`\stackrel{~}{G}(z)`$ are equivalent encoders. The following lemma is a simple result of linear algebra:
###### Lemma 2.1
Two $`n\times k`$ rational encoders $`G(z)`$ and $`\stackrel{~}{G}(z)`$ are equivalent with respect to either Definition A or A if and only if there is a $`k\times k`$ invertible rational matrix $`R(z)`$ such that $`\stackrel{~}{G}(z)=G(z)R(z)`$.
It follows from this lemma that Definition A and Definition A are completely equivalent with respect to equivalence of encoders.
From an algebraic point of view we can identify a convolutional code in the sense of Definition A or Definition A through an equivalence class of rational matrices. The following theorem singles out a set of very desirable encoders inside each equivalence class.
###### Theorem 2.2
Let $`G(z)`$ be a $`n\times k`$ rational encoding matrix of rank $`k`$ defining a code $`𝒞`$. Then there is a $`k\times k`$ invertible rational matrix $`R(z)`$ such that $`\stackrel{~}{G}(z)=G(z)R(z)`$ has the properties:
$`\stackrel{~}{G}(z)`$ is a polynomial matrix.
$`\stackrel{~}{G}(z)`$ is right prime.
$`\stackrel{~}{G}(z)`$ is column reduced with column degrees $`\{e_1,\mathrm{},e_k\}`$.
Furthermore, every polynomial encoding matrix of $`𝒞`$ which is right prime and column-reduced has (unordered) column degrees $`\{e_1,\mathrm{},e_k\}`$. Thus these indices are invariants of the convolutional code.
The essence of Theorem 2.2 was proved by Forney \[6, Theorem 3\]. In Forney related the indices appearing in (iii) to the controllability and observability indices of a controllable and observable system. Paper had an immense impact in the linear systems theory literature. We will follow here the suggestion of McEliece and call these indices the Forney indices of the convolutional code, despite the fact that Theorem 2.2 can be traced back to the last century, when Kronecker, Hermite and in particular Dedekind and Weber studied matrices over the rationals and more general function fields. In Sections 4 and 5 we will make a distinction between the Forney indices as defined above and the Kronecker indices of a submodule of $`𝔽^n[z]`$.
In the coding literature , an encoder satisfying conditions (i), (ii) and (iii) of Theorem 2.2 is called a minimal basic encoder.
So far we have used encoding matrices to describe a convolutional code. As is customary in linear algebra, one often describes a linear subspace as the kernel of a matrix. This leads to the notion of a parity-check matrix. The following theorem is well known (see e.g. ).
###### Theorem 2.3
Let $`𝒞𝔽^n((z))`$ be a rank-$`k`$ convolutional code in the sense of Definition A. Then there exists an $`r\times n`$ matrix $`H(z)`$ such that the code is equivalently described as the kernel of $`H(z)`$:
$$𝒞=\{c(z)𝔽^n((z))H(z)c(z)=0\}.$$
Moreover, it is possible to choose $`H(z)`$ in such a way that:
$`H(z)`$ is a polynomial matrix.
$`H(z)`$ is left prime.
$`H(z)`$ is row-reduced having row degrees $`\{f_1,\mathrm{},f_r\}`$.
Furthermore, every polynomial parity check matrix of $`𝒞`$ which is left prime and row reduced will have (unordered) row degrees $`\{f_1,\mathrm{},f_r\}`$. Thus these indices are invariants of the convolutional code.
Properties (i)–(iii) essentially follow from the fact that the transpose $`H^t(z)`$ is a generator matrix for the dual (orthogonal) code $`𝒞^{}`$.
The set of indices $`\{e_1,\mathrm{},e_k\}`$ and $`\{f_1,\mathrm{},f_r\}`$ differ in general, their sum is however always the same, and is called the degree of the convolutional code. One says that a rank-$`k`$ code $`𝒞𝔽^n((z))`$ has transmission rate $`k/n`$, controller memory $`m:=\mathrm{max}\{e_1,\mathrm{},e_k\}`$ and observer memory $`n:=\mathrm{max}\{f_1,\mathrm{},f_r\}`$.
Another important code parameter is the free distance. The free distance of a code measures the smallest distance between any two different code words, and is formally defined as:
$$d_{\mathrm{free}}(𝒞):=\underset{\genfrac{}{}{0pt}{}{u,v𝒞}{uv}}{\mathrm{min}}\underset{t}{}d_H(u_t,v_t),$$
(2.3)
where $`d_H(,)`$ denotes the usual Hamming distance on $`𝔽^n`$.
## 3 The symbolic dynamics point of view
In this section we present a definition of convolutional codes as it can be found in the symbolic dynamics literature . Convolutional codes in this framework are exactly the linear, compact, irreducible and shift-invariant subsets of $`𝔽^n[[z,z^1]]`$. In order to make this precise, we will have to develop some basic notions from symbolic dynamics.
In the sequel we will work with the finite alphabet $`𝒜:=𝔽^n`$. A block over the alphabet $`𝒜`$ is a finite sequence $`\beta =x_1x_2\mathrm{}x_k`$ consisting of $`k`$ elements $`x_i𝒜`$. If $`w=w(z)=_iw_iz^i𝔽^n[[z,z^1]]`$ is a sequence, one says that the block $`\beta `$ occurs in $`w`$ if there is some integer $`j`$ such that $`\beta =w_jw_{j+1}\mathrm{}w_{k+j1}`$. If $`X𝔽^n[[z,z^1]]`$ is any subset, we denote by $`(X)`$ the set of blocks which occur in some element of $`X`$.
The fundamental objects in symbolic dynamics are the shift spaces. For this let $``$ be a set of blocks, possibly infinite.
###### Definition 3.1
The subset $`X𝔽^n[[z,z^1]]`$ consisting of all sequences $`w(z)`$ which do not contain any of the (forbidden) blocks of $``$ is called a shift space.
The left-shift operator is the $`𝔽`$-linear map
$$\sigma :𝔽[[z,z^1]]𝔽[[z,z^1]],w(z)z^1w(z).$$
(3.1)
Let $`I_n`$ be the $`n\times n`$ identity matrix. The shift map $`\sigma `$ extends to the shift map
$$\sigma I_n:𝔽^n[[z,z^1]]𝔽^n[[z,z^1]].$$
One says that $`X𝔽^n[[z,z^1]]`$ is a shift-invariant set if $`(\sigma I_n)(X)X`$. Clearly shift spaces are shift-invariant subsets of $`𝔽^n[[z,z^1]]`$.
It is possible to characterize shift spaces in a topological manner. For this we will introduce a metric on $`𝔽^n[[z,z^1]]`$:
###### Definition 3.2
If $`v(z)=_iv_iz^i`$ and $`w(z)=_iw_iz^i`$ are both elements of $`𝔽^n[[z,z^1]]`$ we define their distance through:
$$d(v(z),w(z)):=\underset{i}{}2^{|i|}d_H(v_i,w_i).$$
(3.2)
In this metric two elements $`v(z),w(z)`$ are ‘close’ if they coincide over a ‘large block around zero’. One readily verifies that $`d(,)`$ indeed satisfies all the properties of a metric and therefore induces a topology on $`𝔽^n[[z,z^1]]`$. Using this topology we can characterize shift spaces:
###### Theorem 3.3
A subset of $`𝔽^n[[z,z^1]]`$ is a shift space if and only if it is shift-invariant and compact.
* The metric introduced in Definition 3.2 is equivalent to the metric described in \[29, Example 6.1.10\]. The induced topologies are therefore the same. The result follows therefore from \[29, Theorem 6.1.21\].
The topological space $`𝔽^n[[z,z^1]]`$ is a typical example of a linearly compact vector space, a notion introduced by S. Lefschetz. There is a large theory on linearly compact vector spaces, and several of the results which we are going to derive are valid in this broader context. We refer the interested reader to \[25, §10\] for more details.
A further important concept is irreducibility which will turn out to be equivalent to the concept of controllability in our concrete setting.
###### Definition 3.4
A shift space $`X𝔽^n[[z,z^1]]`$ is called irreducible if for every ordered pair of blocks $`\beta ,\gamma `$ of $`(X)`$ there is a block $`\mu `$ such that the concatenated block $`\beta \mu \gamma `$ is in $`(X)`$.
We are now prepared to give the symbolic dynamics definition for a convolutional code and to work out the basic properties for these codes.
###### Definition B
A linear, compact, irreducible and shift-invariant subset of $`𝔽^n[[z,z^1]]`$ is called a convolutional code.
This is an abstract definition and it is not immediately clear how one should encode messages with such convolutional codes. The following will make this clear.
Let $`G(z)`$ be a $`n\times k`$ matrix with entries in the ring of Laurent polynomials $`𝔽[z,z^1]`$. Consider the encoding map:
$$\phi :𝔽^k[[z,z^1]]𝔽^n[[z,z^1]],m(z)c(z)=G(\sigma )m(z).$$
(3.3)
In terms of polynomials the map $`\phi `$ is simply described through $`m(z)c(z)=G(z^1)m(z)`$.
Recall that a continuous map is called closed if the image of a closed set is closed. Using the fact that $`𝔽^n[[z,z^1]]`$ is compact, one (easily) proves the following result:
###### Lemma 3.5
The encoding map $`(`$3.3$`)`$ is $`𝔽`$-linear, continuous and closed.
Clearly $`\mathrm{im}(\phi )`$ is also shift-invariant, and one shows that the image of an irreducible set under $`\phi `$ is irreducible again.
In summary we have shown that $`\mathrm{im}(\phi )`$ describes a convolutional code in the sense of Definition B. Actually the converse is true as well:
###### Theorem 3.6
$`𝒞𝔽^n[[z,z^1]]`$ is a convolutional code in the sense of Definition B if and only if there exists a Laurent polynomial matrix $`G(z)`$ such that $`𝒞=\mathrm{im}(\phi )`$, where $`\phi `$ is the map in $`(`$3.3$`)`$.
A proof of this theorem will be given in the next section after Theorem 4.8.
The question now arises how Definition B relates to Definition A and Definition A. The following theorem will provide a partial answer to this question.
###### Theorem 3.7
Assume that $`𝒞𝔽^n[[z,z^1]]`$ is a nonzero convolutional code in the sense of Definition A or Definition A. Then $`𝒞`$ is not closed, but the closure of $`\overline{𝒞}`$ of $`𝒞`$ is a convolutional code in the sense of Definition B.
* Let $`G(z)`$ be a minimal basic encoder of $`𝒞`$ and let $`w(z)𝔽^n[z]`$ be the first column of $`G(z)`$. Note that $`w(z)𝒞`$ and that there is at least one entry of $`w(z)`$ which does not contain the factor $`(z1)`$. Let $`\varphi _N(z):=_{i=N}^Nz^i𝔽[z,z^1]`$ and consider the sequence of code words $`w^N(z):=\varphi _N(z)w(z)`$. For each $`N>0`$ one has that $`w^N(z)𝒞`$. However $`lim_N\mathrm{}w^N(z)`$ is in $`𝔽^n[[z,z^1]]𝔽^n((z))𝔽^n[[z,z^1]]𝒞`$. This shows that $`𝒞`$ is not a closed set inside $`𝔽^n[[z,z^1]]`$. The closure $`\overline{𝒞}`$ is obtained by extending the input space $`F^k((z))`$ to all of $`F^k[[z,z^1]]`$. The image of $`F^k[[z,z^1]]`$ under the encoding map $`(`$3.3$`)`$ is closed by Lemma 3.5, hence the closure is a code in the sense of Definition B.
Actually one can show that there is a bijective correspondence between the convolutional codes in the sense of Definition A (respectively Definition A) and the convolutional codes in the sense of Definition B, as we will show in Theorem 7.1 and Theorem 7.2. It is also worthwhile to remark that already in 1983 Staiger published a paper where he studied the closure of convolutional codes generated by a polynomial generator matrix.
In analogy to Lemma 2.1, one has:
###### Lemma 3.8
Two $`n\times k`$ encoding matrices $`G(z)`$ and $`\stackrel{~}{G}(z)`$ defined over the Laurent polynomial ring $`𝔽[z,z^1]`$ are equivalent with respect to Definition B if and only if there is a $`k\times k`$ invertible rational matrix $`R(z)`$ such that $`\stackrel{~}{G}(z)=G(z)R(z)`$.
We leave the proof again as an exercise for the reader. We remark that rational transformations of the form $`R(z)`$ are needed to describe the equivalence, even though it is in general not possible to use a rational encoder $`G(z)`$ in the encoding procedure $`(`$3.3$`)`$. This is simply due to the fact that in general the multiplication of an element of $`𝔽(z)`$ with an element of $`𝔽[[z,z^1]]`$ is not defined. The following example should make this clear. (Compare also with Remark 4.4.)
###### Example 3.9
Consider $`f(z)=\frac{1}{1z}=_{i=0}^{\mathrm{}}z^i𝔽(z)`$ and $`g(z)=_{i=\mathrm{}}^{\mathrm{}}z^i𝔽[[z,z^1]]`$. Trying to multiply the two power series $`f(z),g(z)`$ would result in a power series in which each coefficient would be infinite.
In the same way as at the end of Section 2 we define the transmission rate, the degree, the memory and the free distance of a convolutional code $`𝒞`$ in the sense of Definition B.
## 4 Linear time-invariant behaviors
In this section we will take the point of view that a convolutional code is a linear time-invariant behavior in the sense of Willems . Of course behavioral system theory is quite general, allowing all kinds of time axes and signal spaces. In order to relate the behavioral concepts to the previous points of view, we will restrict our study to linear behaviors in $`\left(𝔽^n\right)^{}=𝔽^n[[z,z^1]]`$ and $`\left(𝔽^n\right)^_+=𝔽^n[[z]]`$.
Let $`\sigma `$ be the shift operator defined in $`(`$3.1$`)`$. One says that a subset $`𝔽^n[[z,z^1]]`$ is time-invariant if $`(\sigma I_n)()`$. The concept therefore coincides with the symbolic dynamics concept of shift-invariance.
In addition to linearity and time-invariance, there is a third important concept usually required of a time-invariant behavior:
###### Definition 4.1
A behavior $`𝔽^n[[z,z^1]]`$ is said to be complete if $`w𝔽^n[[z,z^1]]`$ belongs to $``$ whenever $`w|_J`$ belongs to $`|_J`$ for every finite subinterval $`J`$.
The definition simply says that $``$ is complete if membership can be decided on the basis of finite windows. Completeness is an important well behavedness property for linear time-invariant behaviors, as Willems \[50, p. 567\] emphasized with the remark:
> As such, it can be said that the study of non-complete systems does not fall within the competence of system theorists and could be left to cosmologists or theologians.
In Definition 3.2 we introduced a metric on the vector space $`𝔽^n[[z,z^1]]`$. We remark that with respect to this metric a subset $`𝔽^n[[z,z^1]]`$ is complete if and only if every Cauchy sequence converges inside $``$. In other words, the completeness notion of Definition 4.1 coincides with the usual topological notion of completeness.
The following result is known for linearly compact vector spaces, a proof can be found in :
###### Lemma 4.2
A linear subset $`𝔽^n[[z,z^1]]`$ is complete if and only if it is closed and hence compact.
With these preliminaries we can define a convolutional code as follows:
###### Definition C
A linear, time-invariant and complete subset $`𝔽^n[[z,z^1]]`$ is called a convolutional code.
It is immediate from Lemma 4.2 that the convolutional codes defined in Definition B are complete and that Definition C is more general than Definition B, since no irreducibility is required. It also follows from Theorem 3.7 and Lemma 4.2 that the convolutional codes defined in Definition A and Definition A are in general not complete.
Before we elaborate on these differences we would like also to treat the situation when the time axis is $`_+`$ since traditionally a large part of linear systems theory has been concerned with systems defined on the positive time axis.
We first define the left-shift operator acting on $`\left(𝔽^n\right)^_+=𝔽^n[[z]]`$ through:
$$\sigma :𝔽[[z]]𝔽[[z,z^1]],w(z)z^1(w(z)w(0)).$$
(4.1)
We have used the same symbol as in $`(`$3.1$`)`$ since the context will always make it clear if we work over $``$ or $`_+`$. In analogy to $`(`$3.1$`)`$ $`\sigma `$ extends to the shift map $`\sigma I_n:𝔽^n[[z]]𝔽^n[[z]]`$, and one says a subset $`X𝔽^n[[z]]`$ is time-invariant if $`(\sigma I_n)(X)X`$. Notice however that the map of $`(`$4.1$`)`$, unlike that of $`(`$3.1$`)`$, is not invertible.
With this we have:
###### Definition C
A linear, time-invariant and complete subset $`𝔽^n[[z]]`$ is called a convolutional code.
The following fundamental theorem was proved by Willems \[50, Theorem 5\].
###### Theorem 4.3
A subset $`𝔽^n[[z,z^1]]`$ (respectively a subset $`𝔽^n[[z]]`$) is linear, time-invariant and complete if and only if there is a $`r\times n`$ matrix $`P(z)`$ having entries in $`𝔽[z]`$ such that
$$=\{w(z)P(\sigma )w(z)=0\}.$$
(4.2)
By Lemma 3.5 the linear map $`\psi :𝔽^n[[z,z^1]]𝔽^n[[z,z^1]],w(z)P(\sigma )w(z)`$ is continuous and its kernel is therefore a complete set. It is therefore immediate that the behavior defined in $`(`$4.2$`)`$ is linear, time-invariant and complete. The harder part of Theorem 4.3 is the converse statement.
Equation $`(`$4.2$`)`$ is often referred to as a kernel (or AR) representation of a behavioral system. We will denote a behavior having the form $`(`$4.2$`)`$ by $`\mathrm{ker}P(\sigma )`$. By contrast, the encoding map $`\phi `$ defined in $`(`$3.3$`)`$ describes an image (or MA) representation of the behavior $`\mathrm{im}(\phi )=\mathrm{im}G(\sigma )`$.
The most general representation is an ARMA representation. For this let $`P(z)`$ and $`G(z)`$ be matrices of size $`r\times n`$ and $`r\times k`$ respectively, having entries in the Laurent polynomial ring $`𝔽[z,z^1]`$. Then
$$=\{w(z)𝔽^n[[z,z^1]]m(z)𝔽^k[[z,z^1]]:P(\sigma )w(z)=G(\sigma )m(z)\}$$
(4.3)
is called an ARMA model. One immediately verifies that the set $``$ is linear and time-invariant. It is a direct consequence of Lemma 3.5 that $``$ is also closed and hence complete. Theorem 4.3 therefore states that it is possible to eliminate the so called ‘latent variable’ $`m(z)`$ and describe the behavior $``$ by a simpler kernel representation of the form $`(`$4.2$`)`$. It follows in particular that the code $`\mathrm{im}(\phi )=\mathrm{im}G(\sigma )`$ defined in $`(`$3.3$`)`$ has an equivalent kernel representation of the form $`(`$4.2$`)`$ but that in general the converse is not true.
###### Remark 4.4
As we explained in Section 2 it is quite common to use rational encoders for convolutional codes. In the ARMA model $`(`$4.3$`)`$ we required that the entries of $`P(z)`$ and $`G(z)`$ be from the Laurent polynomial ring. If $`P(z)`$ and $`G(z)`$ were rational matrices, then the behavior $`𝔽^n[[z,z^1]]`$ appearing in $`(`$4.3$`)`$ might not be well defined, as we showed in Example 3.9. On the other hand if one restricts the behavior to the positive time axis $`_+`$, i.e. if one assumes that $`𝔽^n[[z]]`$, then the set $`(`$4.3$`)`$ is defined even if $`P(z)`$ and $`G(z)`$ are rational encoders. This is certainly one reason why much classical system theory focused on shift spaces $`𝔽^n[[z]]`$ or $`𝔽^n((z))`$.
In the sequel we will concentrate on representations of the form $`(`$4.2$`)`$. Again the question arises, when are two kernel representations equivalent?
###### Lemma 4.5
Two $`r\times n`$ matrices $`P(z)`$ and $`\stackrel{~}{P}(z)`$ defined over the Laurent polynomial ring $`𝔽[z,z^1]`$ describe the same behavior $`\mathrm{ker}P(\sigma )=\mathrm{ker}\stackrel{~}{P}(\sigma )𝔽^n[[z,z^1]]`$ if and only if there is a $`r\times r`$ matrix $`U(z)`$, unimodular over $`𝔽[z,z^1]`$, such that $`\stackrel{~}{P}(z)=U(z)P(z)`$.
* \[52, Proposition III.3\].
Similarly, if $`P(z)`$ and $`\stackrel{~}{P}(z)`$ are defined over $`𝔽[z]`$, then these matrices define the same behavior $`\mathrm{ker}P(\sigma )=\mathrm{ker}P(\sigma )𝔽^n[[z]]`$ if and only if there is a matrix $`U(z)`$, unimodular over $`𝔽[z]`$, such that $`\stackrel{~}{P}(z)=U(z)P(z)`$.
The major difference between Definition B and Definition C seems to be that Definition C does not require irreducibility. This last concept corresponds to the term controllability (see ) in systems theory. We first start with some notation taken from :
For a sequence $`w=_{\mathrm{}}^{\mathrm{}}w_iz^i𝔽^n[[z,z^1]]`$, we use the symbol $`w^+`$ to denote the ‘right half’ $`_0^{\mathrm{}}w_iz^i`$ and the symbol $`w^{}`$ to denote the ‘left half’ $`_{\mathrm{}}^0w_iz^i`$.
###### Definition 4.6
A behavior $``$ defined on $``$ is said to be controllable if there is some integer $`\mathrm{}`$ such that for every $`w`$ and $`w^{}`$ in $``$ and every integer $`j`$ there exists a $`w^{\prime \prime }`$ such that $`(z^jw^{\prime \prime })^{}=(z^jw)^{}`$ and $`(z^{j+\mathrm{}}w^{\prime \prime })^+=(z^{j+\mathrm{}}w^{})^+`$.
###### Remark 4.7
Loeliger and Mittelholzer speak of strongly controllable if a behavior satisfies the conditions of Definition 4.6. ‘Weakly controllable’ in contrast requires an integer $`\mathrm{}`$ which may depend on the trajectories $`w`$ and $`w^{}`$. The notions are equivalent in our concrete setting.
We leave it as an exercise for the reader to show that irreducibility as introduced in Definition 3.4 is equivalent to controllability for linear, time-invariant and complete behaviors $`𝔽^n[[z,z^1]]`$. The next theorem gives equivalent conditions for a behavior to be controllable.
###### Theorem 4.8
(cf. \[51, Prop. 4.3\]) Let $`P(z)`$ be a $`r\times n`$ matrix of rank $`r`$ defined over $`𝔽[z,z^1]`$. The following conditions are equivalent:
The behavior $`=\mathrm{ker}P(\sigma )=\{w(z)𝔽^n[[z,z^1]]P(\sigma )w(z)=0\}`$ is controllable.
$`P(z)`$ is left prime over $`𝔽[z,z^1]`$.
The behavior $``$ has an image representation. This means there exists an $`n\times k`$ matrix $`G(z)`$ defined over $`𝔽[z,z^1]`$ such that
$$=\{w(z)𝔽^n[[z,z^1]]m(z)𝔽^k[[z,z^1]]:w(z)=G(\sigma )m(z)\}.$$
Combining the theorem with the facts that completeness corresponds to compactness and irreducibility corresponds to controllability gives a proof of Theorem 3.6.
We conclude the section by defining some parameters of a linear, time-invariant and complete behavior. For simplicity we will do this in an algebraic manner. We will first treat behaviors $`𝔽^n[[z]]`$, i.e. behaviors in the sense of Definition C. In Remark 4.10 we will explain how the definitions have to be adjusted for behaviors defined on the time axis $``$.
Assume that $`P(z)`$ is a $`r\times n`$ polynomial matrix of rank $`r`$ defining the behavior $`=\mathrm{ker}P(\sigma )`$. There exists a matrix $`U(z)`$, unimodular over $`𝔽[z]`$, such that $`\stackrel{~}{P}(z)=U(z)P(z)`$ is row-reduced with ordered row degrees $`\nu _1\mathrm{}\nu _r`$. The indices $`\nu =(\nu _1,\mathrm{},\nu _r)`$ are invariants of the row module of $`P(z)`$ (and hence also invariants of the behavior $``$), and are sometimes referred to as the Kronecker indices or observability indices of $``$. The invariant $`\delta :=_{i=1}^r\nu _i`$ is called the McMillan degree of the behavior $``$. If we think of $``$ as a convolutional code in the sense of Definition C then we say that $``$ has transmission rate $`\frac{nr}{n}`$. Finally, the free distance of the code is defined as in $`(`$2.3$`)`$.
###### Remark 4.9
The Kronecker indices $`\nu `$ are in general different from the minimal row indices (in the sense of Forney ) of the $`𝔽(z)`$-vector space generated by the rows of $`P(z)`$. They coincide with the minimal row indices if and only if $`P(z)`$ is left prime.
###### Remark 4.10
If $`𝔽^n[[z,z^1]]`$ is a linear, time-invariant and complete behavior, then we can define parameters like the Kronecker indices and the McMillan degree in the following way: Assume $`P(z)`$ has the property that $`=\mathrm{ker}P(\sigma )`$. There exists a matrix $`U(z)`$, unimodular over $`𝔽[z,z^1]`$, such that $`\stackrel{~}{P}(z)=U(z)P(z)`$ is row-reduced and $`P(0)`$ has full row rank $`r`$. One shows again that the row degrees of $`\stackrel{~}{P}(z)`$ are invariants of the behavior. The McMillan degree, the transmission rate and the free distance are then defined in the same way as for behaviors $`𝔽^n[[z]]`$.
## 5 The module point of view
Fornasini and Valcher and the present author in joint work with Schumacher, Weiner and York proposed a module-theoretic approach to convolutional codes. The module point of view simplifies the algebraic treatment of convolutional codes to a large degree, and this simplification is probably almost necessary if one wants to study convolutional codes in a multidimensional setting .
From a systems theoretic point of view, the module-theoretic approach studies linear time-invariant systems whose states start at zero and return to zero in finite time. Such dynamical systems have been studied by Hinrichsen and Prätzel-Wolters , who recognized these systems as convenient objects for the study of systems equivalence.
In our development we will again deal with the time axes $``$ and $`_+`$ in a parallel manner.
###### Definition D
A submodule $`𝒞`$ of $`𝔽^n[z,z^1]`$ is called a convolutional code.
We like the module-theoretic language. If one prefers to define everything in terms of trajectories then one could equivalently define $`𝒞`$ as $`𝔽`$-linear, time-invariant subset of $`𝔽^n[[z,z^1]]`$ whose elements have finite support.
The analogous definition for codes supported on the positive time axis $`_+`$ is:
###### Definition D
A submodule $`𝒞`$ of $`𝔽^n[z]`$ is called a convolutional code.
Since both the rings $`𝔽[z,z^1]`$ and $`𝔽[z]`$ are principal ideal domains (PID), a convolutional code $`𝒞`$ has always a well-defined rank $`k`$, and there is a full-rank matrix $`G(z)`$ of rank $`k`$ such that $`𝒞=\mathrm{colsp}_{𝔽[z,z^1]}G(z)`$ (respectively $`𝒞=\mathrm{colsp}_{𝔽[z]}G(z)`$ if $`𝒞`$ is defined as in Definition D). We will call $`G(z)`$ an encoder of $`𝒞`$, and the map
$$\phi :𝔽^k[z,z^1]𝔽^n[z,z^1],m(z)c(z)=G(z)m(z)$$
(5.1)
an encoding map.
###### Remark 5.1
In contrast to the situation of Section 3, it is possible to define a convolutional code in the sense of Definition D (respectively Definition D) using a rational encoder. For this, assume that $`G(z)`$ is an $`n\times k`$ matrix with entries in $`𝔽(z)`$. Then
$$𝒞=\{c(z)𝔽^n[z,z^1]m(z)𝔽^k[z,z^1]:c(z)=G(z)m(z)\}$$
defines a submodule of $`𝔽^n[z,z^1]`$. Note that the map $`(`$5.1$`)`$ involving a rational encoding matrix $`G(z)`$ has to be ‘input-restricted’ in this case.
In analogy to Lemma 3.8 we have:
###### Lemma 5.2
Two $`n\times k`$ matrices $`G(z)`$ and $`\stackrel{~}{G}(z)`$ defined over the Laurent polynomial ring $`𝔽[z,z^1]`$ (respectively over the polynomial ring $`𝔽[z]`$) generate the same code $`𝒞𝔽^n[z,z^1]`$ (respectively $`𝒞𝔽^n[z]`$) if and only if there is a $`k\times k`$ matrix $`U(z)`$, unimodular over $`𝔽[z,z^1]`$ (respectively over $`𝔽[z]`$), such that $`\stackrel{~}{G}(z)=G(z)U(z)`$.
As we already mentioned earlier convolutional codes in the sense of Definitions D and D are linear and time-invariant. The following theorem answers any question about controllability (i.e. irreducibility) and completeness.
###### Theorem 5.3
A nonzero convolutional code with either Definition D or D is controllable and incomplete.
* The proof of the completeness part of the Theorem is analogous to the proof of Theorem 3.7. In order to show controllability, let $`G(z)`$ be an encoding matrix for a code $`𝒞𝔽^n[z]`$ and consider two code words $`w(z)=G(z)(a_0+a_1+\mathrm{}+a_sz^s)`$ and $`w^{}(z)=G(z)(b_0+b_1+\mathrm{}+b_sz^s)`$. The codeword $`w^{\prime \prime }(z)`$ required by Definition 4.6 can be constructed in the form
$$G(z)(a_0+a_1+\mathrm{}+a_jz^j+b_{j+\mathrm{}}z^{j+\mathrm{}}+\mathrm{}+\mathrm{}+b_sz^s).$$
Submodules of $`𝔽^n[z,z^1]`$ (respectively of $`𝔽^n[z]`$) form the Pontryagin dual of linear, time-invariant and complete behaviors in $`𝔽^n[[z,z^1]]`$ (respectively $`𝔽^n[[z]]`$). In the following we follow and explain this in a very explicit way when the time axis is $``$. Of course everything can be done mutatis mutandis when the time axis is $`_+`$.
Consider the bilinear form:
$$\begin{array}{ccc}\hfill (,):𝔽^n[[z,z^1]]\times 𝔽^n[z,z^1]& & 𝔽\hfill \\ \hfill (w,v)& & \underset{i=\mathrm{}}{\overset{\mathrm{}}{}}w_i,v_i,\hfill \end{array}$$
(5.2)
where $`,`$ represents the standard dot product on $`𝔽^n`$. One shows that $`(,)`$ is well defined and nondegenerate, in particular because there are only finitely many nonzero terms in the sum. For any subset $`𝒞`$ of $`𝔽^n[z,z^1]`$ one defines the annihilator
$$𝒞^{}=\{w𝔽^n[[z,z^1]](w,v)=0,v𝒞\}$$
(5.3)
and the annihilator of a subset $``$ of $`𝔽^n[[z,z^1]]`$ is
$$^{}=\{v𝔽^n[z,z^1](w,v)=0,w\}.$$
(5.4)
The relation between these two annihilator operations is given by:
###### Theorem 5.4
If $`𝒞𝔽^n[z,z^1]`$ is a convolutional code with generator matrix $`G(z)`$, then $`𝒞^{}`$ is a linear, left-shift-invariant and complete behavior with kernel representation $`P(z)=G^t(z)`$. Conversely, if $`𝔽^n[[z,z^1]]`$ is a linear, left-shift-invariant and complete behavior with kernel representation $`P(z)`$, then $`^{}`$ is a convolutional code with generator matrix $`G(z)=P^t(z)`$.
###### Remark 5.5
An elementary proof of Theorem 5.4 in the case of the positive time axis $`_+`$ is given in .
###### Remark 5.6
Theorem 5.4 is a special instance of a broad duality theory between solution spaces of difference equations on the one hand and modules on the other, for which probably the most comprehensive reference is Oberst . In this article Oberst \[37, p. 22\] works with a bilinear form which is different from $`(`$5.2$`)`$. This bilinear form induces however the same duality as shown in . Extensions of duality results to group codes were derived by Forney and Trott in .
For finite support convolutional codes in the sense of Definition D or Definition D the crucial issue is observability. In the literature there have been several definitions of observability and it is not entirely clear how these definitions relate to each other.
In the sequel we will follow .
###### Definition 5.7
(cf. \[4, Prop. 2.1\]) A code $`𝒞`$ is observable if there exists an integer $`N`$ such that, whenever the supports of $`v`$ and $`v^{}`$ are separated by a distance of at least $`N`$ and $`v+v^{}𝒞`$, then also $`v𝒞`$ and $`v^{}𝒞`$.
With this we have the ‘Pontryagin dual statement’ of Theorem 4.8:
###### Theorem 5.8
(cf. \[42, Prop. 2.10\]) Let $`G(z)`$ be a $`n\times k`$ matrix of rank $`k`$ defined over $`𝔽[z,z^1]`$. The following conditions are equivalent:
The convolutional code $`𝒞=\mathrm{colsp}_{𝔽[z,z^1]}G(z)`$ is observable.
$`G(z)`$ is right prime over $`𝔽[z,z^1]`$.
The code $`𝒞`$ has a kernel representation. This means there exists an $`r\times n`$ ‘parity-check matrix’ $`H(z)`$ defined over $`𝔽[z,z^1]`$ such that
$$𝒞=\{v(z)𝔽^n[z,z^1]H(z)v(z)=0\}.$$
###### Remark 5.9
The concept of observability is clearly connected to the coding concept of non-catastrophicity. Indeed an encoder is non-catastrophic if and only if the code generated by this encoder is observable. In the context of Definition A (respectively Definition A) every code has a catastrophic as well as a non-catastrophic encoder. In the module setting of Definition D every encoder of an observable code is non-catastrophic and every encoder of an non-observable code is catastrophic. If one defines a convolutional code by Definition D then one could talk of a ‘non-catastrophic convolutional code’. The term observable seems however much more appropriate.
As at the end of Section 4, we now define the code parameters. We do it only for codes given by Definition D and leave it to the reader to adapt the definitions to codes given by Definition D.
Assume that $`G(z)`$ is an $`n\times k`$ polynomial matrix of rank $`k`$ defining the code $`𝒞=\mathrm{colsp}_{𝔽[z]}G(z)`$. There exists a unimodular matrix $`U(z)`$ such that $`\stackrel{~}{G}(z)=G(z)U(z)`$ is column-reduced with ordered column degrees $`\kappa _1\mathrm{}\kappa _k`$. The indices $`\kappa =(\kappa _1,\mathrm{},\kappa _k)`$ are invariants of the code $`𝒞`$, which we call the Kronecker indices or controllability indices of $`𝒞`$. The invariant $`\delta :=_{i=1}^r\kappa _i`$ is called the degree of the code $`𝒞`$. The free distance of the code is defined as in $`(`$2.3$`)`$. Finally we say that $`𝒞`$ has transmission rate $`\frac{k}{n}`$.
## 6 First-order representations
In this section we provide an overview of the different first-order representations (realizations) associated with the convolutional codes and encoding maps which we have defined.
We start with the encoding map $`(`$2.2$`)`$. As is customary in most of the coding literature, we view the map $`(`$2.2$`)`$ as an input-output operator from the message space to the code space. The existence of associated state spaces and realizations can be shown on an abstract level. Kalman first showed how the encoding map $`(`$2.2$`)`$ can be ‘factored’ resulting in a realization of the encoding matrix $`\phi `$. Fuhrmamnn refined the realization procedure in an elegant way. (Compare also .)
In the sequel we will simply assume that a realization algorithm exists. We summarize the main results in the following two theorems:
###### Theorem 6.1
Let $`T(z)`$ be a $`p\times m`$ proper transfer function of McMillan degree $`\delta `$. Then there exist matrices $`(A,B,C,D)`$ of size $`\delta \times \delta `$, $`\delta \times m`$, $`p\times \delta `$ and $`p\times m`$ respectively such that
$$T(z)=C(zIA)^1B+D.$$
(6.1)
The minimality conditions are that $`(A,B)`$ forms a controllable pair and $`(A,C)`$ forms an observable pair. Finally $`(`$6.1$`)`$ is unique in the sense that if $`T(z)=\stackrel{~}{C}(zI\stackrel{~}{A})^1\stackrel{~}{B}+\stackrel{~}{D}`$ with $`(\stackrel{~}{A},\stackrel{~}{B})`$ controllable and $`(\stackrel{~}{A},\stackrel{~}{C})`$ observable, then there is a unique invertible matrix $`S`$ such that
$$(\stackrel{~}{A},\stackrel{~}{B},\stackrel{~}{C},\stackrel{~}{D})=(SAS^1,SB,CS^1,D).$$
(6.2)
Consider the encoding map $`(`$2.2$`)`$ with generator matrix $`G(z)`$. Let $`m(z)=_{i=s}^tm_iz^i𝔽^k((z))`$ and $`c(z)=_{i=s}^tc_iz^i𝔽^n((z))`$ be the sequence of message and code symbols respectively. Then one has:
###### Theorem 6.2
Assume that $`G(z)`$ has the property that $`\mathrm{rank}G(0)=k`$. Then $`G(z^1)`$ is a proper transfer function, and by Theorem 6.1 there exist matrices $`(A,B,C,D)`$ of appropriate sizes such that $`G(z^1)=C(zIA)^1B+D`$. The dynamics of $`(`$2.2$`)`$ are then equivalently described by:
$$\begin{array}{ccc}\hfill x_{t+1}& =& Ax_t+Bm_t,\hfill \\ \hfill c_t& =& Cx_t+Dm_t.\hfill \end{array}$$
(6.3)
The realization $`(`$6.3$`)`$ is useful if one wants to describe the dynamics of the encoder $`G(z)`$. It is however less useful if one is interested in the construction of codes having certain properties. The problem is that every code $`𝒞`$ has many equivalent encoders whose realizations appear to be completely different.
###### Example 6.3
The encoders
$$G(z)=\left(\begin{array}{c}\frac{1z}{z4}\\ \frac{1+z}{z4}\end{array}\right)\text{ and }\stackrel{~}{G}(z)=\left(\begin{array}{c}\frac{1z}{(z2)(z+3)}\\ \frac{1+z}{(z2)(z+3)}\end{array}\right)$$
are equivalent since they define the same code in the sense of Definition A. The transfer functions $`G(z^1)`$ and $`\stackrel{~}{G}(z^1)`$ are however very different from a systems theory point of view. Indeed, they have different McMillan degrees, and over the reals the first is stable whereas the second is not. The state space descriptions are therefor very different for these encoders.
This example should make it clear that for the purpose of constructing good convolutional codes, representation $`(`$6.3$`)`$ is not very useful.
We are now coming to the realization theory of the behaviors and codes of Section 4 and 5. We will continue with our algebraic approach. The results are stated for the positive time axis $`_+`$, but they hold mutatis mutandis for the time axis $``$.
###### Theorem 6.4 (Existence)
Let $`P(z)`$ be an $`r\times n`$ matrix of rank $`r`$ describing a behavior $``$ of the form $`(`$4.2$`)`$ with McMillan degree $`\delta `$. Let $`k=nr`$. Then there exist (constant) matrices $`G,F`$ of size $`\delta \times (\delta +k)`$ and a matrix $`H`$ of size $`n\times (\delta +k)`$ such that $``$ is equivalently described by:
$$=\{w(z)𝔽^n[[z]]\zeta (z)𝔽^{\delta +k}[[z]]:(\sigma GF)\zeta (z)=0,w(z)=H\zeta (z)\}.$$
(6.4)
Moreover the following minimality conditions will be satisfied:
$`G`$ has full row rank;
$`\left[\genfrac{}{}{0pt}{}{G}{H}\right]`$ has full column rank;
$`\left[\genfrac{}{}{0pt}{}{zGF}{H}\right]`$ is right prime.
For a proof, see \[26, Thm. 4.3\] or . Equation $`(`$6.4$`)`$ describes the behavior locally in terms of a time window of length 1. The computation of the matrices $`G,F,H`$ from a kernel description is not difficult. It can even be done ‘by inspection’, i.e., just by rearranging the data . The next result describes the extent to which minimal first-order realizations are unique. A proof is given in \[26, Thm. 4.34\].
###### Theorem 6.5 (Uniqueness)
The matrices $`(G,F,H)`$ are unique in the following way: If $`(\stackrel{~}{G},\stackrel{~}{F},\stackrel{~}{H})`$ is a second triple of matrices describing the behavior $``$ through $`(`$6.4$`)`$ and if the minimality conditions (i), (ii) and (iii) are satisfied, then there exist unique invertible matrices $`S`$ and $`T`$ such that
$$(\stackrel{~}{G},\stackrel{~}{F},\stackrel{~}{G})=(SGT^1,SFT^1,HT^1).$$
(6.5)
The relation to the traditional state-space theory is as follows: Assume that $`P(z)`$ can be partitioned into $`P(z)=(Y(z)U(z))`$ with $`U(z)`$ a square $`r\times r`$ matrix and $`\mathrm{deg}detU(z)=\delta `$, the McMillan degree of the behavior $``$. Assume that $`(G,F,H)`$ provides a realization for $``$ through $`(`$6.4$`)`$. Then one shows that the pencil $`\left[\genfrac{}{}{0pt}{}{zGF}{H}\right]`$ is equivalent to the pencil:
$$\left[\begin{array}{cc}zI_\delta A& B\\ 0& I_k\\ C& D\end{array}\right].$$
(6.6)
The minimality condition (iii) simply translates into the condition that $`(A,C)`$ forms an observable pair, showing that the behavior $``$ is observable. One also verifies that the matrices $`(A,B,C,D)`$ form a realization of the proper transfer function $`U(z)^1Y(z)`$ and that this is a minimal realization if and only if $`(A,B)`$ forms a controllable pair. Finally $`(A,B)`$ is controllable if and only if the behavior $``$ is controllable.
The Pontryagin dual statements of Theorem 6.4 and 6.5 are (see ):
###### Theorem 6.6 (Existence)
Let $`G(z)`$ be an $`n\times k`$ polynomial matrix generating a rate $`\frac{k}{n}`$ convolutional code $`𝒞𝔽^n[z]`$ of degree $`\delta `$. Then there exist $`(\delta +nk)\times \delta `$ matrices $`K,L`$ and a $`(\delta +nk)\times n`$ matrix $`M`$ (all defined over $`𝔽`$) such that the code $`𝒞`$ is described by
$$𝒞=\{v(z)𝔽^n[z]x(z)𝔽^\delta [z]:zKx(z)+Lx(z)+Mv(z)=0\}.$$
(6.7)
Moreover the following minimality conditions will be satisfied:
K has full column rank;
$`[KM]`$ has full row rank;
$`[zK+LM]`$ is left prime.
Equation $`(`$6.7$`)`$ describes the behavior again locally in terms of a time window of length 1.
###### Theorem 6.7 (Uniqueness)
The matrices $`(K,L,M)`$ are unique in the following way: If $`(\stackrel{~}{K},\stackrel{~}{L},\stackrel{~}{M})`$ is a second triple of matrices describing the code $`𝒞`$ through $`(`$6.7$`)`$ and if the minimality conditions (i), (ii) and (iii) are satisfied, then there exist unique invertible matrices $`T`$ and $`S`$ such that
$$(\stackrel{~}{K},\stackrel{~}{L},\stackrel{~}{M})=(TKS^1,TLS^1,TM).$$
(6.8)
If $`G(z)`$ can be partitioned into $`G(z)=\left[\begin{array}{c}Y(z)\\ U(z)\end{array}\right]`$ with $`U(z)`$ a square $`k\times k`$ matrix and $`\mathrm{deg}detU(z)=\delta `$, the degree of the code $`𝒞`$, then the pencil $`[zK+LM]`$ is equivalent to the pencil:
$$\left[\begin{array}{ccc}zI_\delta A& 0_{\delta \times (nk)}& B\\ C& I_{nk}& D\end{array}\right].$$
(6.9)
The minimality condition (iii) then translates into the condition that $`(A,B)`$ forms a controllable pair, showing that the code $`𝒞`$ is controllable. One also verifies that the matrices $`(A,B,C,D)`$ form a realization of the proper transfer function $`Y(z)U(z)^1`$, that this is a minimal realization if and only if $`(A,C)`$ forms an observable pair, and that this is the case if and only if the code $`𝒞`$ is observable. Finally, the Kronecker indices of $`𝒞`$ coincide with the controllability indices of the pair $`(A,B)`$ .
The systems-theoretic meaning of the representation $`(`$6.9$`)`$ is as follows (see ). Partition the code vector $`v(z)`$ into:
$$v(z)=\left[\begin{array}{c}y(z)\\ u(z)\end{array}\right]𝔽^n[z]$$
and consider the equation:
$$\left[\begin{array}{ccc}zI_\delta A& 0_{\delta \times (nk)}& B\\ C& I_{nk}& D\end{array}\right]\left[\begin{array}{c}x(z)\\ y(z)\\ u(z)\end{array}\right]=0.$$
(6.10)
Let
$`x(z)`$ $`=`$ $`x_0z^\gamma +x_1z^{\gamma 1}+\mathrm{}+x_\gamma ;x_t𝔽^\delta ,t=0,\mathrm{},\gamma ,`$
$`u(z)`$ $`=`$ $`u_0z^\gamma +u_1z^{\gamma 1}+\mathrm{}+u_\gamma ;u_t𝔽^k,t=0,\mathrm{},\gamma ,`$
$`y(z)`$ $`=`$ $`y_0z^\gamma +y_1z^{\gamma 1}+\mathrm{}+y_\gamma ;y_t𝔽^{nk},t=0,\mathrm{},\gamma .`$
Then $`(`$6.10$`)`$ is satisfied if and only if
$`x_{t+1}`$ $`=`$ $`Ax_t+Bu_t,`$
$`y_t`$ $`=`$ $`Cx_t+Du_t,`$ (6.11)
$`v_t`$ $`=`$ $`\left({\displaystyle \genfrac{}{}{0pt}{}{y_t}{u_t}}\right),x_0=0,x_{\gamma +1}=0,`$
is satisfied. Note that the state-space representation $`(`$6$`)`$ is different from the representation $`(`$6.3$`)`$. Equation $`(`$6$`)`$ describes the dynamics of the systematic and rational encoder
$$G(z)U^1(z)=\left[\begin{array}{c}Y(z)U(z)^1\\ I_k\end{array}\right].$$
The encoding map $`u(z)y(z)=G(z)U^1(z)u(z)`$ is input-restricted, i.e. $`u(z)`$ must be in the column module of $`U(z)`$ in order to make sure that $`y(z)`$ and $`x(z)`$ have finite support. In terms of systems theory, this simply means that the state should start at zero and return to zero in finite time. Linear systems satisfying these requirements have been studied by Hinrichsen and Prätzel-Wolters .
## 7 Differences and similarities among the definitions
After having reviewed these different definitions for convolutional codes, we would like to make some comparison.
The definitions of Section 2 and Section 3 viewed convolutional codes as linear, time-invariant, controllable and observable behaviors, not necessarily complete. Definition C and Definition C were more general in the sense that non-controllable behaviors were accepted as codes. Definition D and Definition D were more general in the sense that non-observable codes were allowed.
In the following subsection we show that all definitions are equivalent for all practical purposes if one restricts oneself to controllable and observable codes.
### 7.1 Controllable and observable codes
Consider a linear, time-invariant, complete behavior $`𝔽^n[[z,z^1]]`$, i.e. a convolutional code in the sense of Definition C. Let
$$𝒞:=𝔽^n((z)).$$
Then one has
###### Theorem 7.1
$`𝒞`$ is a convolutional code in the sense of Definition A, and its completion $`\overline{𝒞}`$ is the largest controllable sub-behavior of $``$. Moreover, one has a bijective correspondence between controllable behaviors $`𝔽^n[[z,z^1]]`$ and convolutional codes $`𝒞𝔽^n((z))`$ in the sense of Definition A.
* Let $`=\mathrm{ker}P(\sigma )=\{w(z)𝔽^n[[z,z^1]]P(\sigma )w(z)=0\}`$. If $``$ is not controllable, then $`P(z)`$ is not left prime and one has a factorization $`P(z)=V(z)\stackrel{~}{P}(z)`$, where $`\stackrel{~}{P}(z)`$ is left prime and describes the controllable sub-behavior $`\mathrm{ker}\stackrel{~}{P}(\sigma )`$. Since $`\mathrm{ker}V(\sigma )`$ is an autonomous behavior it follows that
$$𝒞=𝔽^n((z))=\mathrm{ker}P(\sigma )𝔽^n((z))=\mathrm{ker}\stackrel{~}{P}(\sigma )𝔽^n((z)).$$
It follows (compare with Theorem 3.7) that the completion $`\overline{𝒞}=\mathrm{ker}\stackrel{~}{P}(\sigma )`$.
Consider now a convolutional code $`𝒞𝔽^n((z))`$ in the sense of Definition A. Define:
$`\stackrel{ˇ}{𝒞}`$ $`:=`$ $`𝒞𝔽^n[z,z^1]`$
$`\stackrel{ˇ}{\stackrel{ˇ}{𝒞}}`$ $`:=`$ $`𝒞𝔽^n[z].`$
Conversely if $`𝒞𝔽^n[z]`$ is a convolutional code in the sense of Definition D, then define:
$`\widehat{𝒞}`$ $`:=`$ $`\mathrm{span}_{𝔽[z,z^1]}\{v(z)v(z)𝒞\}.`$
$`\widehat{\widehat{𝒞}}`$ $`:=`$ $`\mathrm{span}_{𝔽((z))}\{v(z)v(z)𝒞\}.`$
By definition it is clear that $`\widehat{𝒞}\widehat{\widehat{𝒞}}`$ are convolutional codes in the sense of Definition D and Definition A respectively.
###### Theorem 7.2
Assume that $`𝒞𝔽^n((z))`$ is a convolutional code in the sense of Definition A. Then $`\stackrel{ˇ}{\stackrel{ˇ}{𝒞}}𝔽^n[z]`$ is an observable code in the sense of Definition A. Moreover the operations $`\widehat{\widehat{}}`$ and $`\stackrel{ˇ}{\stackrel{ˇ}{}}`$ induce a bijective correspondence between the observable codes $`𝒞𝔽^n[z]`$ and convolutional codes $`𝒞𝔽^n((z))`$ in the sense of Definition A.
Theorem 7.2 is essentially the Pontryagin dual statement of Theorem 7.1; we leave it to the reader to work out the details. Theorem 7.1 and 7.2 together show that there is a bijection between controllable and observable codes in the sense of one definition and another definition. For controllable and observable codes the code parameters like the rate $`k/n`$, the degree $`\delta `$ and the Forney (Kronecker) indices are all the same. Moreover the free distance is in every case the same as well. For all practical purposes one can therefore say that the frameworks are completely equivalent, if one is only interested in controllable and observable codes.
The advantage of Definition D (respectively Definition D) over the other definitions lies in the fact that non-observable codes become naturally part of the theory. It also seems that for construction purposes the relation between quasi-cyclic codes and convolutional codes is best described in a module-theoretic framework.
Definition C (respectively Definition C) allows one to introduce non-controllable codes in a natural way.
A Laurent series setting as in Definition A seems to be most natural if one is interested in the description of the encoder and/or syndrome former. Extensions of the Laurent series framework to multidimensional convolutional codes is however much less natural than the polynomial framework, which is why the theory of multidimensional convolutional codes has mainly been developed in a module-theoretic framework .
### 7.2 Duality
In $`(`$5.2$`)`$ we introduced a bilinear form which induced a bijection between behaviors $`𝔽^n[[z,z^1]]`$ and modules $`𝒞𝔽^n[z,z^1]`$. This duality is a special instance of Pontryagin duality, and generalizes to group codes and multidimensional systems .
In this subsection we show that the bilinear form $`(`$5.2$`)`$ can also be used to obtain a duality between modules and modules (both in $`𝔽^n[z,z^1]`$) or between behaviors and behaviors (both in $`𝔽^n[[z,z^1]]`$).
For this let $`𝒞𝔽^n[z,z^1]`$ be a submodule. Define:
$$𝒞^{}:=𝒞^{}𝔽^n[z,z^1].$$
(7.1)
One immediately verifies that $`𝒞^{}`$ is a submodule of $`𝔽^n[z,z^1]`$, which necessarily is observable. One always has $`𝒞(𝒞^{})^{}`$.
One can do something similar for behaviors. For this let $`𝔽^n[[z,z^1]]`$ be a behavior. Define:
$$^{}:=\left(𝔽^n[z,z^1]\right)^{}=\overline{^{}}.$$
(7.2)
Then it is immediate that $`^{}`$ is a controllable behavior, $`(^{})^{}`$ and $`(^{})^{}`$ describes the controllable sub-behavior of $``$.
It is also possible to adapt $`(`$5.2$`)`$ for a duality of subspaces $`𝒞𝔽^n((z))`$. For such a subspace we define:
$$𝒞^{}:=\left(𝒞𝔽^n[z,z^1]\right)^{}𝔽^n((z)).$$
(7.3)
The duality $`(`$7.1$`)`$ does not in general correspond to the linear algebra dual of the $`R=𝔽[z,z^1]`$ module $`𝒞R^n`$ since there is some ‘time reversal’ involved. The same is true for the duality $`(`$7.3$`)`$, which does not correspond to the linear algebra dual of the $`𝔽((z))`$ vector space $`𝒞`$ without time reversal.
If one works however with the ‘time-reversed’ bilinear form:
$$\begin{array}{ccc}\hfill [,]:𝔽^n[[z,z^1]]\times 𝔽^n[z,z^1]& & 𝔽\hfill \\ \hfill (w(z),v(z))& & \underset{i=\mathrm{}}{\overset{\mathrm{}}{}}w_i,v_i\hfill \end{array}$$
(7.4)
then the definitions $`(`$7.1$`)`$ and $`(`$7.3$`)`$ do correspond to the module dual (and the linear algebra dual respectively), used widely in the coding literature . In this case one has: If $`G(z)`$ is a generator matrix of $`𝒞^{}`$ then $`H(z):=G^t(z)`$ is a parity check matrix of $`(𝒞^{})^{}`$.
In the Laurent-series context it is also possible to induce the duality $`(`$7.3$`)`$ directly through the time-reversed bilinear form defined on $`𝔽^n((z))\times 𝔽^n((z))`$:
$$\begin{array}{ccc}\hfill [,]:𝔽^n((z))\times 𝔽^n((z))& & 𝔽\hfill \\ \hfill (w(z),v(z))& & \underset{i=\mathrm{}}{\overset{\mathrm{}}{}}w_i,v_i.\hfill \end{array}$$
(7.5)
Note that the sum appearing in $`(`$7.5$`)`$ is always well defined. This bilinear form has been widely used in functional analysis and in systems theory .
### 7.3 Convolutional codes as subsets of $`𝔽[[z,z^1]]`$, a case study.
In this subsection we illustrate the differences of the definitions in the peculiar case $`n=1`$.
If one works with Definition A or Definition B then there exist only the two trivial codes having the $`1\times 1`$ generator matrix $`(1)`$ and $`(0)`$ as subsets of $`𝔽[[z,z^1]]`$.
The situation of Definition C is already more interesting. For each polynomial $`p(z)`$ one has the associated ‘autonomous behavior’:
$$=\{w(z)p(\sigma )w(z)=0\}.$$
(7.6)
Autonomous behaviors are the extreme case of uncontrollable behaviors. If $`\mathrm{deg}p(z)=\delta `$, then $``$ is a finite-dimensional $`𝔽`$-vector space of dimension $`\delta `$. For coding purposes $``$ is not useful at all. Indeed, the code allows only $`\delta `$ symbols to be chosen freely, say the symbols $`w_0,w_1,\mathrm{},w_{\delta 1}`$. With this the codeword $`w(z)=_{i=\mathrm{}}^{\mathrm{}}w_iz^i`$ is determined, and the transmission of $`w(z)`$ requires infinite symbols in the past and infinite symbols in the future. In other words, the code has transmission rate $`0`$. The distance of the code is however very good, namely $`d_{\mathrm{free}}()=\mathrm{}`$. If $``$ is defined on the positive time axis, i.e. $`𝔽[[z]]`$ then the situation is only slightly better. Indeed in this situation, one sends first $`\delta `$ message words and then an infinite set of ‘check symbols’. As these remarks make clear, a code of the form $`(`$7.6$`)`$ is not very useful.
The most interesting situation happens in the setup of Definition D and Definition D. In this situation the codes are exactly the ideals $`<g(z)>𝔽[z,z^1]`$ (respectively $`<g(z)>𝔽[z]`$). We now show that ideals of the form $`<g(z)>`$ are of interest in the coding context.
###### Example 7.3
Let $`𝔽=𝔽_2=\{0,1\}`$. Consider the ideal generated by $`g(z)=(z+1)`$. $`<g(z)>𝔽[z,z^1]`$ consists in this case of the even-weight sequences, namely the set of all sequences with a finite and even number of ones. This code is controllable but not observable.
Ideals of the form $`<g(z)>`$ are the extreme case of non-observable behaviors. In principle this makes it impossible for the receiver to decode a message. However with some additional ‘side-information’ decoding can still be performed, as we now explain.
One of the most often used codes in practice is probably the cyclic redundancy check code (CRC code). These codes are the main tool to ensure error-free transmissions over the Internet. They can be defined in the following way: Let $`g(z)𝔽[z]`$ be a polynomial. Then the encoding map is simply defined as:
$$\phi :𝔽[z]𝔽[z],m(z)c(z)=g(z)m(z).$$
(7.7)
The code is then the ideal $`<g(z)>=\mathrm{im}(\phi )`$. The distance of this code is $`2`$, since there exists an integer $`N`$ such that $`(z^N1)<g(z)>`$. As we already mentioned the code is not observable. Assume now that the sender gives some additional side information indicating the start and the end of a message. This can be either done by saying: “I will send in a moment 1 Mb”, or it can be done by adding some ‘stop signal’ at the end of the transmission. Once the receiver knows that the transmission is over, he applies long division to compute
$$c(z)=\stackrel{~}{m}(z)g(z)+r(z),\mathrm{deg}r(z)<\delta .$$
If $`r(z)=0`$ the receiver accepts the message $`\stackrel{~}{m}(z)`$ as the transmitted message $`m(z)`$. Otherwise he will ask for retransmission.
The code performs best over a channel (like the Internet) which has the property that the whole message is transmitted correctly with probability $`p`$ and with probability $`1p`$ whole blocks of the message are corrupted during transmission. One immediately sees that the probability that a corrupted message $`\stackrel{~}{m}(z)`$ is accepted is $`q^\delta `$, where $`q=|𝔽|`$ is the field size.
One might argue that the code $`<g(z)>=\mathrm{im}(\phi )`$ is simply a cyclic block code, but this is not quite the case. Note that the protocol does not specify any length of the code word and in each transmission a different message length can be chosen. In particular the code can be even used if the message length is longer than $`N`$, where $`N`$ is the smallest integer such that $`(z^N1)<g(z)>`$.
###### Example 7.4
Let $`𝔽=𝔽_2=\{0,1\}`$ and let $`g(z)=z^{20}+1`$. Assume transmission is done on a channel with very low error probability where once in a while a burst error might happen destroying a whole sequence of bits. Assume that the sender uses a stop signal where he repeats the 4 bits $`0011`$ for 100 times. Under these assumptions the receiver can be reasonably sure once a transmission has been complete. The probability of failure to detect a burst error is in this case $`2^{20}`$ which is less than $`10^6`$. Note that $`g(z)`$ is a very poor generator for a cyclic code of any block length.
###### Remark 7.5
CRC codes are in practice often implemented in a slightly different way than we described it above (see e.g. ). The sender typically performs long division on $`z^\delta m(z)`$ and computes
$$z^\delta m(z)=f(z)g(z)+r(z),\mathrm{deg}r(z)<\delta .$$
He then transmits the code word $`c(z):=z^\delta m(z)r(z)<g(z)>`$. Clearly the schemes are equivalent. The advantage of the latter is that the message sequence $`m(z)`$ is transmitted in ‘plain text’, allowing processing of the data immediately.
### 7.4 Some geometric remarks
One motivation for the author to take a module-theoretic approach to convolutional coding theory has come from algebraic-geometric considerations. As is explained in , a submodule of rank $`k`$ and degree $`\delta `$ in $`𝔽^n[z]`$ describes a quotient sheaf of rank $`k`$ and degree $`\delta `$ over the projective line $`^1`$. The set of all such quotient sheaves having rank $`k`$ and degree at most $`\delta `$ has the structure of a smooth projective variety denoted by $`X_{k,n}^\delta `$. This variety has been of central interest in the recent algebraic geometry literature. In the context of coding theory, it has actually been used to predict the existence of maximum-distance-separable (MDS) convolutional codes .
The set of convolutional codes in the sense of Definition A or A or B having rate $`\frac{k}{n}`$ and degree at most $`\delta `$ form all proper Zariski open subsets of $`X_{k,n}^\delta `$. The points in the closure of these Zariski open sets are exactly the non-observable codes if the rate is $`\frac{1}{n}`$. These geometric considerations suggest that non-observable convolutional codes should be incorporated into a complete theory of convolutional codes. The following example will help to clarify these issues:
###### Example 7.6
Let $`\delta =2`$, $`k=1`$ and $`n=2`$, i.e., consider $`X_{1,2}^2`$. Any code of degree at most 2 then has an encoder of the form:
$$G(z)=\left(\begin{array}{c}g_1(z)\\ g_2(z)\end{array}\right)=\left(\begin{array}{c}a_0+a_1z+a_2z^2\\ b_0+b_1z+b_2z^2\end{array}\right)$$
We can identify the encoder through the point $`(a_0,a_1,a_2,b_0,b_1,b_2)^5`$. The variety $`X_{1,2}^2`$ is in this example exactly the projective space $`^5`$. For codes in the sense of Definition A or A or B, $`G(z)`$ must be taken as a basic minimal encoder in order to have a unique parameterization. This requires that $`g_1(z)`$ and $`g_2(z)`$ are coprime polynomials. The set of coprime polynomials $`g_1(z),g_2(z)`$ viewed as a subset of $`^5`$ forms a Zariski open subset $`U^5`$ described by the resultant condition
$$det\left(\begin{array}{cccc}a_0& 0& b_0& 0\\ a_1& a_0& b_1& b_0\\ a_2& a_1& b_2& b_1\\ 0& a_2& 0& b_2\end{array}\right)0.$$
For codes in the sense of Definition D, we require that $`a_0`$ and $`b_0`$ are not simultaneously zero in order to have a unique parameterization. Definition D leads to a larger Zariski open set $`V`$, i.e. $`UV^5`$. Only with Definition D does one obtain the whole variety $`X_{1,2}^2=^5`$.
In the general situation $`X_{k,n}^\delta `$ naturally contains the non-observable codes as well. If $`k=1`$, then $`X_{1,n}^\delta =^{n(\delta +1)1}`$, and the codes in the sense of Definition D having rate $`\frac{1}{n}`$ and degree at most $`\delta `$ are exactly parameterized by $`X_{1,n}^\delta `$.
## 8 Conclusion
The paper surveys a number of different definitions of convolutional codes. All definitions have in common that a convolutional code is a subset $`𝒞𝔽^n[[z,z^1]]`$ which is both linear and time-invariant. The definitions differ in requirements such as controllability, observability, completeness and restriction to finite support.
If one requires that a code be both controllable and observable, then the restriction to any finite time window will result in equivalent definitions. Actually Loeliger and Mittelholzer define a convolutional code locally in terms of one trellis section and they require in their definition that a code is controllable and observable. Algebraically such a trellis section is simply described through the generalized first order description $`(`$6.4$`)`$ or $`(`$6.7$`)`$.
If one wants to have a theory which allows one to work with rational encoders, then it will be necessary that the code has finite support on the negative time axis $`_{}`$ (or alternatively on the positive time axis $`_+`$). This is one reason why a large part of the coding literature works with the field of formal Laurent series.
If one wants in addition to have a theory which can accommodate non-observable codes (and such a theory seems to have some value) then it is best to work in a module-theoretic setting.
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# Period map for non-compact holomorphically symplectic manifolds
## 1 Introduction
Deformations and moduli of compact Kähler manifolds are a well studied subject, dating back to Kodaira-Spencer \[KS\]. The moduli of non-compact manifolds are rarely mentioned, mostly because they are much harder to define and study.
The work on the moduli of compact holomorphically symplectic manifolds and Calabi-Yau manifolds is still far from the conclusion; the local case is due to F.Bogomolov, A.Beauville, G.Tian, A.Todorov, P.Deligne and Z.Ran (\[Bo\], \[Bea\], \[T\], \[To\], \[R\]). Extreme importance of this subject is highlighted by thousands of papers on Mirror Symmetry, which appeared since then.
In the non-compact case, some work in this direction was done by M.Kontsevich and S.Barannikov (\[BK\]) and others, but, for the most part, this territory is still uncharted. However, there are many examples that suggest that at least for some non-compact holomorphically symplectic manifolds a good local deformation theory does exist. In particular, in the well-studied case of smooth crepant resolutions of symplectic quotient singularities in dim $`2`$ (the so-called Du Val points), a likely candidate for the universal local deformation is provided by the simultaneous resolution of Brieskorn \[Br\].
In this paper we extend to the non-compact case the algebraic construction of the local deformation space of a Calabi-Yau manifold $`M`$, due to Z.Ran. Unfortunately, our results are valid only when the manifold $`M`$ is holomorphically symplectic.
Our approach is essentially the same as the original approach of Bogomolov. It is based on the so-called period map. Instead of deformations of a holomorphically symplectic manifold $`M`$, one considers deformations of the pair $`M,\mathrm{\Omega }`$, where $`\mathrm{\Omega }`$ denotes the holomorphic symplectic form, that is, a nowhere degenerate closed $`(2,0)`$-differential form. Given a local deformation $`\pi :\stackrel{~}{M}S`$ of $`M,\mathrm{\Omega }`$ with a simply connected base, the cohomology of the individual fibers of $`\pi `$ are identified by the Gauss-Manin connection. Taking the cohomology class of the holomorphic symplectic form of each fiber, one obtains a map $`\mathrm{𝖯𝖾𝗋}:SH^2(M)`$. Bogomolov and Beauville have shown that for $`M`$ compact, the map $`\mathrm{𝖯𝖾𝗋}`$ induces a holomorphic immersion of the coarse marked moduli space $``$ of $`M,\mathrm{\Omega }`$ into $`H^2(M)`$. The image of $`\mathrm{𝖯𝖾𝗋}`$ belongs to a certain quadric $`𝒞H^2(M)`$, cut by the so-called Bogomolov-Beauville form. Moreover, the period map $`\mathrm{𝖯𝖾𝗋}:𝒞`$ is locally an isomorphism.
Bogomolov extended these results to Calabi-Yau manifolds in an unpublished I.H.E.S preprint (1982). In 1987, Tian and Todorov published a different proof of Bogomolov’s theorem. Their proof of Bogomolov-Tian-Todorov theorem was based on Hodge theory. An algebraic version of their arguments was proposed by Z.Ran (\[R\]). Ran’s argument uses the degeneration of the $`E_2`$-term of the Dolbeault spectral sequence (proven in algebraic case by Deligne and Illusie, \[DI\]). However, this spectral sequence is not degenerate in non-compact case, hence this argument does not work for open Calabi-Yau manifolds.
We found that a version of Z.Ran’s argument is valid for holomorphic symplectic manifolds (under some additional, quite weak, assumptions). One can explain this heuristically as follows. For a complex manifold $`M`$, deformations are classified by $`H^1(𝒯(M))`$, where $`𝒯(M)`$ is the tangent sheaf. When $`M`$ is Calabi-Yau, $`𝒯(M)`$ is isomorphic to $`\mathrm{\Omega }^{n1}(M)`$, where $`n=\mathrm{𝖽𝗂𝗆}_{}M`$. To show that the deformations of $`M`$ have no obstructions, we would need to prove that the $`E_2`$-term of the Dolbeault spectral sequence degenerates in $`H^1(\mathrm{\Omega }^{n1}(M))`$. This is very far from truth in the non-compact case. However, if $`M`$ is holomorphic symplectic, we have $`𝒯(M)\mathrm{\Omega }^1(M)`$, and instead of $`H^1(M,\mathrm{\Omega }^{n1}(M))`$ we have to consider $`H^1(M,\mathrm{\Omega }^1(M))`$. The only differential in the spectral sequence that maps into $`H^1(M,\mathrm{\Omega }^1(M))`$ starts at $`H^1(M,𝒪_M)`$. If we assume for simplicity that $`H^i(M,𝒪_M)=0`$ for $`i1`$, then this differential vanishes tautologically. The differentials that start at $`H^1(M,\mathrm{\Omega }^1(M))`$ can still be non-trivial, but it turns out that they become irrelevant if instead of deformations of $`M`$ one considers deformations of the pair $`M,\mathrm{\Omega }`$.
Thus in the case $`H^i(M,𝒪_M)=0`$, $`i1`$ we obtain the following result.
###### Theorem 1.1
Let $`M`$ be a holomorphic symplectic manifold such that for every $`i1`$ we have $`H^i(𝒪_M)=0`$. Then there exists a coarse moduli space of formal deformations
$$\pi :\stackrel{~}{M}Spl(M,\mathrm{\Omega })$$
of pairs $`M,\mathrm{\Omega }`$. Moreoved, the period map $`\mathrm{𝖯𝖾𝗋}:Spl(M,\mathrm{\Omega })H^2(M)`$ gives an isomorphism of $`Spl(M,\mathrm{\Omega })`$ and the formal completion of $`H^2(M)`$ in $`[\mathrm{\Omega }]H^2(M)`$. $`\mathrm{}`$
There is a more general version of this result (Theorem 3.6), which works for a larger class of non-compact holomorphically symplectic manifolds. Here we also obtain a coarse moduli space $`Spl(M,\mathrm{\Omega })`$, which is smooth and finite-dimensional, but it is no longer isomorphic to $`H^2(M)`$. However, the period map remains an immersion.
We will now give a semi-rigorous sketch of the proof of this result. First, consider an easy but important example of affine holomorphic symplectic manifold $`M`$. We have $`H^1(𝒯(M))=0`$, hence any first-order complex deformation of $`M`$ is trifial. Using induction, it is easy to show that any formal complex deformation of $`M`$ is also trivial, that is, for any formal deformation $`\pi :\stackrel{~}{M}S`$ of $`M`$, we have $`\stackrel{~}{M}S\times M`$. However, a symplectic deformation needs not to be trivial, because we may have non-trivial variations of holomorphic symplectic structure. A formal deformation of the pair $`M,\mathrm{\Omega }`$ is determined by the deformation of a closed $`(2,0)`$-form $`\mathrm{\Omega }`$.
By Grothendieck’s theorem, the topological cohomology of $`M`$ is isomorphic to the hypercohomology of the algebraic de Rham complex
(1.1)
$$\begin{array}{ccccccccc}0& & 𝒪_M& \stackrel{}{}& \mathrm{\Omega }^1M& \stackrel{}{}& \mathrm{\Omega }^2M& \stackrel{}{}& \mathrm{}\end{array}$$
Since $`M`$ is affine, $`H^i(\mathrm{\Omega }^j(M))=0`$ for $`i>0`$. Therefore, $`H^2(M)`$ is isomorphic to the second cohomology of the complex 1.1. A first order infinitesimal automorphism of $`M`$ is given by the section of $`TM`$, which is isomorphic to $`\mathrm{\Omega }^1(M)`$. A (co-)vector field $`\gamma \mathrm{\Omega }^1(M)𝒯(M)`$ acts on $`\mathrm{\Omega }^iM`$ by the Lie derivative. Therefore, $`\gamma `$ acts on a closed form $`\mathrm{\Omega }\mathrm{\Omega }^2(M)`$ by adding $`d\gamma `$, and the first-order deformations of $`M,\mathrm{\Omega }`$ are classified by closed 2-forms up to exact 2-forms. Using induction, it is easy to check that this is true in any order. Therefore, the period map gives an isomorphism of the coarse moduli of $`M,\mathrm{\Omega }`$ and the formal neighbourhood of the class $`[\mathrm{\Omega }]`$ in $`H^2(M)`$.
The general proof is deduced from the affine version as follows.
Let $`S`$ be a spectrum of an Artin ring over $``$. We consider the deformations of $`M,\mathrm{\Omega }`$ over $`S`$ as a stack of groupoids over $`M`$. We introduce another stack, called Kodaira-Spencer stack, which – whenever it is defined – classifies the maps from $`M`$ to $`H^2(M)`$. Using the affine version of the main theorem, we show that these stacks are equivalent over affine open subsets of $`M`$. This immediately implies that these stacks are equivalent over the whole $`M`$, which concludes the study of the formal deformations of $`M,\mathrm{\Omega }`$ and shows that these deformations are classified by the maps form the base of deformation to $`H^2(M)`$.
The above account is greatly simplified. To make it at least approximately workable, we have to do big technical adjustments. The problems are twofold: first of all, the Kodaira-Spencer groupoid is not defined in a general situation – we have to go step-by-step through what we call elementary base extensions to define it properly. Secondly, the Kodaira-Spencer groupoid classifies not the maps from the base $`S`$ to $`H^2(M)`$, but rather a certain maps of complexes of sheaves, which are reduced to the maps from $`S`$ to $`H^2(M)`$ when $`H^1(𝒪_M)=H^2(𝒪_M)=0`$.
Now we give a more precise version of the definition of the Kodaira-Spencer stack. Let $`S`$ be an Artin scheme over $``$, and $`S_0S`$ a closed subscheme defined by an ideal $`I𝒪_S`$, such that $`I^2=0`$. Assume that the ideal $`I`$ is sufficiently small (see Definition 4.3 for the precise condition on $`I`$). Fix a deformation $`\pi _0:\stackrel{~}{M}_{S_0}S_0`$ of $`M,\mathrm{\Omega }`$. Consider the set $`\mathrm{𝖣𝖾𝖿}(\stackrel{~}{M}_{S_0},S)`$ of all deformations $`\stackrel{~}{M}_S`$ of $`M,\mathrm{\Omega }`$ over $`S`$, equipped with an isomorphism
$$\stackrel{~}{M}_S\times _SS_0\stackrel{~}{M}_{S_0}.$$
Clearly, $`\mathrm{𝖣𝖾𝖿}(\stackrel{~}{M}_{S_0},S)`$ is a stack of groupoids over $`M`$. Consider a truncated relative de Rham complex $`F^1\mathrm{\Omega }^{}(\stackrel{~}{M}_{S_0}/S_0)`$,
$$\begin{array}{ccccc}\mathrm{\Omega }^1(\stackrel{~}{M}_{S_0}/S_0)& \stackrel{}{}& \mathrm{\Omega }^2(\stackrel{~}{M}_{S_0}/S_0)& \stackrel{}{}& \mathrm{}\end{array}$$
A holomorphic symplectic structure on the deformation $`\stackrel{~}{M}_{S_0}`$ defines a morphism of complexes $`\pi _0^{}𝒪_{S_0}[2]F^1\mathrm{\Omega }^{}(\stackrel{~}{M}_{S_0}/S_0)`$. Taking its derivative along the Gauss-Manin connection on $`S_0`$, we obtain the so-called Kodaira-Spencer map of the deformation $`\stackrel{~}{M}_{S_0}`$
$$\pi _0^{}TS_0[2]\stackrel{\theta _0}{}F^1\mathrm{\Omega }^{}(\stackrel{~}{M}_{S_0}/S_0).$$
The Kodaira-Spencer stack $`KS(\stackrel{~}{M}_{S_0},S)`$ is defined as follows. The objects of $`KS(\stackrel{~}{M}_{S_0},S)`$ are all morphisms of complexes
$$\pi _0^{}TS𝒪_S𝒪_{S_0}[2]\stackrel{\theta }{}F^1\mathrm{\Omega }^{}(\stackrel{~}{M}_{S_0}/S_0)$$
such that the restriction of $`\theta `$ to $`TS_0`$ is equal to $`\theta _0`$. The morphisms between any two such $`\theta _1`$, $`\theta _2`$ are chain homotopies between them – that is, maps
$$\pi _0^{}TS𝒪_S𝒪_{S_0}[1]\stackrel{\gamma }{}F^1\mathrm{\Omega }^{}(\stackrel{~}{M}_{S_0}/S_0)$$
satisfying $`d\gamma =\theta _1\theta _2`$. This obviously defines a groupoid. The definition of the period map from $`\mathrm{𝖣𝖾𝖿}(\stackrel{~}{M}_{S_0},S)`$ to $`KS(\stackrel{~}{M}_{S_0},S)`$ is straightforward; using the deformation theory for affine $`M`$, we show that this is an isomorphism. In other words, there are exactly as many ways to extend the deformation from $`S_0`$ to $`S`$ as there are ways to extend the corresponding morphism to $`H^2(M)`$. Using induction, this leads immediately to the classification result for the deformations stated above.
The paper is organized as follows. We start with a semi-heuristic analytic proof of the main theorem obtained by the use of the Cartan-Maurer equation, as suggested by Kontsevich and Barannikov. This proof cannot be made precise in the non-compact situation, but we have decided to include it anyway, since it exhibits nicely the general idea behind the rigorous proof. Then we turn to purely algebaric methods. Section 3 contains the necessary definitions and the precise statement of our result in the strongest form. In Section 4, we introduce the approriate symplectic version of the Kodaira-Spencer class and define elementary extensions. In Section 5, we rework our construction in the language of stacks and prove the main theorem. Finally, Section 6 is a short postface which explains how our results are related to the known facts (in particular, to the results of Z. Ran).
## 2 Deformations in mixed formal complex-analytic category
One of the way to study deformations us through the $`DG`$-algebra approach suggested by Kontsevich and Barannikov (\[Ba\], \[BK\]). The main advantage of this method is that it foregoes the tedious step-by-step constructions of Grothendieck’s local deformations theory, and gives the results in the form of an explicit power series. Unfortunately, it is not quite as useful in non-compact case, when we have no means to check that these series converge.
The deformations one obtains from the Dolbeault complex lie a weird mixed formal-complex-analytic category. To obtain something definite, one needs some kind of integrability-type conditions. Kontsevich and Barannikov work with compact manifolds, so that in their situation this is not a problem – integrability can be obtained directly from functional analysis. However, since we study open manifolds, the functional analysis does not help, and we stay in the mostly useless mixed category.
Nevertheless, Kontsevich’s approach (which goes back to Kodaira) is beautiful and quite useful as a heuristic tool. Therefore we decided to express some of our results in this language before going to the rigourous step-by-step proof.
Since the $`DG`$-algebra approach is used only for heuristics, this section will be quite sketchy; a cursory knowledge of \[BK\] is required.
We start with a review of the deformation theory as it is given in \[Ba\] and \[BK\]. Let $`M`$ be a complex manifold, $`M_{}`$ the underlying real analytic manifold and $`M_{}[[t]]:=M_{}\times \mathrm{Spec}([[t]])`$ the “mixed formal-real analytic” manifold obtained as a product of $`M_{}`$ and the formal disk $`\mathrm{\Delta }:=\mathrm{Spec}([[t]])`$. Using $`DG`$-algebras, one may classify the complex deformations of $`M`$ over $`\mathrm{\Delta }`$, that is, the complex structures $`J`$ on $`M_{}[[t]]`$ such that the zero fiber of $`(M_{}[[t]],J)`$ is isomorphic to $`M`$. It is well known that such deformations are classified by the solutions of the Maurer-Cartan equation
(2.1)
$$\overline{}\gamma (t)=\frac{1}{2}[\gamma (t),\gamma (t)],$$
where $`\gamma (t)\mathrm{\Lambda }^{0,1}(TM)[[t]]`$ is a $`C[[t]]`$-valued $`(0,1)`$-form with coefficients in holomorphic vector fields, and
$$[,]:\mathrm{\Lambda }^{0,1}(TM)[[t]]\times \mathrm{\Lambda }^{0,1}(TM)[[t]]\mathrm{\Lambda }^{0,2}(TM)[[t]]$$
is the Schouten bracket.
We explain in a few words how the complex structures are classified by the solutions of Maurer-Cartan.
Consider the sheaf $`A^\text{},\text{}=\mathrm{\Lambda }^\text{},\text{}(M)`$ as an abstract sheaf of algebras. A complex structure on $`M`$ is defined by an identification between $`\mathrm{\Lambda }^1(M)`$ and $`A^{0,1}`$. Since $`\mathrm{\Lambda }^1(M)`$ classifies the derivations of the sheaf of smooth functions, to give an identification $`\mathrm{\Lambda }^1(M)A^{0,1}`$ is the same as to give a derivation $`\overline{}:C^{\mathrm{}}(M)A^{0,1}`$. The difference between two such operators is given by $`\gamma \mathrm{\Lambda }^{0,1}(TM)`$. An integrability condition is written as $`(\overline{}+\gamma )^2=0`$, which is rewritten as the Maurer-Cartan equation. Deformations are equivalent if the corresponding operators are exchanged by an automorphism of $`A^\text{},\text{}`$.
We are going to write a similar interpretation for holomorphic symplectic deformations.
Fix a holomorphic symplectic form $`\mathrm{\Omega }A^{2,0}`$. A holomorphic symplectic deformation is defined by an operator
$$\overline{}+\gamma :C^{\mathrm{}}(M)A^{0,1},\gamma (\mathrm{\Omega })=0$$
such that $`(\overline{}+\gamma )^2`$, or, equivalently, $`\gamma `$ satisfies the Maurer-Cartan equation 2.1. Deformations are equivalent if the corresponding operators are exchanged by an automorphism of $`A^\text{},\text{}`$ preserving $`\mathrm{\Omega }`$.
In this spirit, Barannikov \[Ba\] describes the deformations of Calabi-Yau manifolds, with $`\mathrm{\Omega }`$ a nowhere degenerate section of the sheaf of top-degree $`(p,0)`$-forms.
For the holomorphic symplectic manifolds, the deformations are described in terms the sheaf $`\mathrm{𝑎𝑚}`$ of holomorphic Hamiltonian vector fields. First of all, automorphisms of $`A^\text{},\text{}`$ preserving $`\mathrm{\Omega }`$ correspond to the Hamiltonian vector fields $`\beta \mathrm{𝑎𝑚}(M)`$. Secondly, the condition $`\gamma \mathrm{\Omega }=0`$ means that $`\gamma \mathrm{\Lambda }^{0,1}(\mathrm{𝑎𝑚}(M))`$.
We obtain that the first-order deformations of $`M`$ are described by the cohomology of the sheaf of Hamiltonian vector fields. To prove that the deformations of $`M`$ are unobstructed (that is, every one-parametric first-order deformation is extended to a deformation over a formal disk), we need to show the following. Let $`\gamma _1\mathrm{\Lambda }^{0,1}(\mathrm{𝑎𝑚}(M))`$ be a $`\overline{}`$-closed $`(0,1)`$-form with coefficients in $`\mathrm{𝑎𝑚}`$. Then there exists a series $`\gamma _2,\gamma _3,\mathrm{}\mathrm{\Lambda }^{0,1}(\mathrm{𝑎𝑚}(M))`$ such that
$$(\overline{}+t\gamma _1+t^2\gamma _2+\mathrm{})^2=0,$$
that is, the formal series $`\gamma _1+t\gamma _2+t^2\gamma _3+\mathrm{}`$ satisfies the Maurer-Cartan equation 2.1. In this setting, Maurer-Cartan can be written as follows:
(2.2)
$$\overline{}\gamma _{n+1}=\frac{1}{2}\underset{i+j=n}{}[\gamma _i,\gamma _j],n=2,3,\mathrm{}$$
Consider the following commutative diagram
(2.3)
$$\begin{array}{ccccccc}& \overline{}& \mathrm{\Lambda }^{2,1}& \overline{}& \mathrm{\Lambda }^{2,2}& \overline{}& \\ & & & & & & \\ & \overline{}& \mathrm{\Lambda }^{1,1}& \overline{}& \mathrm{\Lambda }^{1,2}& \overline{}& \\ & & & & & & \\ & \overline{}& \mathrm{\Lambda }^{0,1}& \overline{}& \mathrm{\Lambda }^{0,2}& \overline{}& \end{array}$$
Identifying $`TM`$ and $`\mathrm{\Lambda }^{1,0}`$, we can realize the polyvector fields as differential $`(p,q)`$-forms. In particular, the solutions $`\gamma _i`$ of 2.2 are considered now as sections of $`\mathrm{\Lambda }^{1,1}`$. The Hamiltonian condition is written as $`\gamma _i=0`$, where $`:\mathrm{\Lambda }^{1,1}\mathrm{\Lambda }^{2,1}`$ is an operator in 2.3. Consider the “symplectic Hodge operator”
$$\mathrm{\Lambda }:\mathrm{\Lambda }^{p,q}\mathrm{\Lambda }^{p2,q}$$
which is adjoint to the multiplication by the holomorphic symplectic form $`\mathrm{\Omega }`$ via the non-degenerate product defined by $`\mathrm{\Omega }`$. The Schouten bracket is written in terms of $`\mathrm{\Lambda }`$ and $``$ as follows.
###### Lemma 2.1
(Tian-Todorov lemma for holomorphic symplectic manifolds.) Let $`\gamma ,\gamma ^{}\mathrm{\Lambda }^{1,1}`$, and let $`[\gamma ,\gamma ^{}]\mathrm{\Lambda }^{2,1}`$ denote the Schouten bracket (we identify $`TM`$ and $`\mathrm{\Lambda }^{1,0}`$ using the holomorphic symplectic form). Then
$$[\gamma ,\gamma ^{}]=\mathrm{\Lambda }(\gamma \gamma ^{})\mathrm{\Lambda }(\gamma \gamma ^{})\mathrm{\Lambda }(\gamma \gamma ^{}).$$
###### Proof.
This is a symplectic version of the standard Tian-Todorov lemma, and it is proven in exactly the same fashion as the usual one (\[To\]). $`\mathrm{}`$
The main result of this Section is the following theorem, which states that the deformations of holomorphic symplectic manifold are unobstructed, given that $`H^2(𝒪_M)=0`$. This assumptions is not essential, but it makes the statements much simpler.
###### Theorem 2.2
Let $`M`$ be a complex holomorphic symplectic manifold, and let $`\gamma _1\mathrm{\Lambda }^{0,1}(\mathrm{𝑎𝑚})`$ be a $`\overline{}`$-closed 1-form with coefficients in Hamiltonian vector fields. Assume that $`H^2(𝒪_M)=0`$, that is, the holomorphic cohomology of the structure sheaf vanish. Then 2.2 has a solution.
###### Proof.
Let us write 2.2 in terms of the diagram 2.3, using the isomorphism $`TM\mathrm{\Lambda }^{1,0}`$. Using induction, we may assume that
$$\overline{}\underset{i+j=n}{}[\gamma _i,\gamma _j]=\underset{i+j+k=n}{}[\gamma _i,[\gamma _j,\gamma _k]]+[[\gamma _i,\gamma _j,]\gamma _k],$$
which is equal zero by Jacoby identity. Therefore,
(2.4)
$$\underset{i+j=n}{}[\gamma _i,\gamma _j]$$
is $`\overline{}`$-closed; to solve 2.2 we need to show that it is $`\overline{}`$-exact, that is, to prove that it represents a zero class in the cohomology $`H^2(\mathrm{𝑎𝑚}(M))`$. The Hamiltonian vector fields are identified with $``$-closed $`(1,0)`$-forms. Poincare lemma gives an exact sequence of sheaves
(2.5)
$$0𝒪_M\stackrel{}{}\mathrm{𝑎𝑚}(M)\mathrm{\hspace{0.25em}0},$$
Consider the piece
$$H^2(𝒪_M)H^2(\mathrm{𝑎𝑚})\stackrel{h}{}H^3(M)$$
of the corresponding long exact sequence. By Lemma 2.1, the expression 2.4 is $``$-exact; therefore, it represents zero in $`H^3(M)`$. By our assumptions, $`H^2(𝒪_M)=0`$, and therefore, the map $`h:H^2(\mathrm{𝑎𝑚})H^3(M)`$ is an embedding. This proves that the sum 2.4 represents zero in $`H^2(\mathrm{𝑎𝑚}(M))`$, and there exists $`\gamma _n\mathrm{\Lambda }^{0,1}(\mathrm{𝑎𝑚}(M))`$ such that $`\overline{}\gamma _n=_{i+j=n}[\gamma _i,\gamma _j]`$. We proved Theorem 2.2. $`\mathrm{}`$
In Theorem 2.2 we identified the deformation space of $`M,\mathrm{\Omega }`$ with the cohomology of $`\mathrm{𝑎𝑚}(M)`$. These cohomology spaces are very easy to write down explicitly. Writing the long exact sequence corresponding to 2.5, we obtain
$$H^1(M)H^1(𝒪_M)H^1(\mathrm{𝑎𝑚}(M))H^2(M)H^2(𝒪_M).$$
This gives an identification of the deformation space with $`H^2(M)`$ when $`H^1(𝒪_M)=H^2(𝒪_M)=0`$, and, more generally, allows one to express the holomorphic symplectic deformations through $`H^i(M)`$, $`H^i(𝒪_M)`$ $`(i=1,2)`$.
In the remaining part of the paper, we combine the intuition of the $`DG`$-algebra approach with the hard science of Grothendieck’s local deformations theory, and obtain essentially the same results in a much more rigorous setting.
## 3 Statement of the results.
### 3.1 Admissible manifolds.
To begin with, we will describe the restrictions that we need to impose on the given holomorphically symplectic manifold $`X`$.
Let $`X`$ be a smooth algebraic manifold over $``$. By a deformation of $`X`$ over a pointed scheme $`S,oS`$ we will understand a scheme $`\stackrel{~}{X}/S`$ smooth over $`S`$ and equipped with an isomorphism $`o\times _S\stackrel{~}{X}X`$. We will only consider deformations over spectra $`S=\mathrm{Spec}A`$ of local Artin $``$-algebra $`A`$, so that the fixed point is given by the maximal ideal $`𝔪A`$. Unless otherwise mentioned, all deformations will be assumed to be of this type.
Recall that the de Rham complex $`\mathrm{\Omega }^{\text{}}(X)`$ is equipped with the Hodge, a.k.a. stupid filtration $`F^{\text{}}\mathrm{\Omega }^{\text{}}(X)`$ given by
$$F^i\mathrm{\Omega }^j(X)=\{\begin{array}{cc}\mathrm{\Omega }^j(X),\hfill & ji,\hfill \\ 0,\hfill & \text{otherwise}.\hfill \end{array}$$
This filtration is also defined, and by the same formula, for the relative de Rham complex $`\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}/S)`$ of an arbitrary deformation $`\stackrel{~}{X}/S`$.
For symplectic deformations, it is the first term $`F^1\mathrm{\Omega }^{\text{}}(X)`$ of the Hodge filtration that plays the crucial role.
The complex $`F^1\mathrm{\Omega }^{\text{}}(X)`$ can be included in an obvious exact triangle
(3.1)
$$\begin{array}{ccccccc}F^1\mathrm{\Omega }^{\text{}}(X)& & \mathrm{\Omega }^{\text{}}(X)& & 𝒪_X& & \end{array}$$
where $`𝒪_X`$ is the structure sheaf of the manifold $`X`$. This triangle induces an exact triangle on cohomology.
###### Definition 3.1
A smooth algebraic manifold $`X`$ over $``$ is called admissible if for any deformation $`\pi :\stackrel{~}{X}S`$ the relative cohomology sheaf
$$R^2\pi _{}F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}/S)$$
is a flat sheaf on $`S`$ and the canonical map
$$R^2\pi _{}F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}/S)R^2\pi _{}\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}/S)$$
is injective.
This definition is pretty technical, because it is given in the most general form. For all practical applications that we see at the moment, it suffices to assume the stronger condition on $`X`$ provided by the following easy lemma.
###### Lemma 3.2
Let $`X`$ be a smooth complex algebraic manifold. If for all $`p1`$ the canonical map
$$H^p(X,)H^p(X,𝒪_X)$$
is surjective, then the manifold $`X`$ is admissible.
###### Proof.
Let $`\pi :\stackrel{~}{X}S`$ be an arbitrary deformation. The existence of the Gauss-Manin connection implies that
$$R^p\pi _{}\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}/S)H^p(X,)𝒪_S$$
for every $`p0`$. We will prove that
1. for every $`p2`$, the canonical map $`R^p\pi _{}\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}/S)R^p\pi _{}𝒪(\stackrel{~}{X})`$ is surjective and the sheaf $`R^p\pi _{}𝒪(\stackrel{~}{X})`$ is flat on $`S`$.
Use downward induction on $`p`$, starting with an arbitrary $`p>2\mathrm{𝖽𝗂𝗆}X`$. Assume (A) proved for all $`p>k`$. Denote by $`i:oS`$ the embedding of the base point. Consider the restriction of the cohomology exact triangle induced by 3.1 to the base point $`oS`$. Then base change and the Nakayma Lemma immediately imply that the map
$$R^k\pi _{}\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}/S)R^k\pi _{}𝒪(\stackrel{~}{X})$$
is surjective. This together with the inductive assumption implies that
1. the canonical map $`R^p\pi _{}F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}/S)R^p\pi _{}\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}/S)`$ is injective and the sheaf $`R^p\pi _{}F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}/S)`$ if flat over $`S`$ for every $`pk+1`$.
Therefore if the sheaf $`R^k\pi _{}𝒪(\stackrel{~}{X})`$ is not flat, then the coboundary map
$$L^1i^{}R^k\pi _{}𝒪(\stackrel{~}{X}))i^{}R^k\pi _{}F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}/S)$$
is not zero, which contradicts the assumption. Carrying the induction half-step further, we derive that the map
$$R^1\pi _{}\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}/S)R^1\pi _{}𝒪(\stackrel{~}{X})$$
is surjective. Therefore (B) also holds for $`p=2`$, which proves the lemma. $`\mathrm{}`$
This lemma shows that a manifold $`X`$ is admissible in the following cases:
1. $`X`$ is compact (Hodge theory).
2. $`X`$ is affine ($`H^i(X,𝒪_X)=0`$ for $`i1`$).
3. More generally, $`X`$ admits a proper generically one-to-one map $`\pi :XY`$ into an affine variety $`Y`$ with rational singularities (again $`H^i(X,𝒪_X)=0`$ for $`i1`$).
4. $`X`$ admits a proper generically one-to-one map $`\pi :XY`$ into an affine variety $`Y`$ and has trivial canonical bundle $`K_X`$ ($`H^i(X,𝒪_X)=H^i(X,K_X)=0`$ for $`i>0`$. <sup>1</sup><sup>1</sup>1This follows immediately from the Grauert-Riemenschneider Vanishing Theorem, \[GPR\].
If $`X`$ is an admissible manifold, one can choose a splitting $`H^2(X,)H^2(X,F^1\mathrm{\Omega }^{\text{}}(X))`$ of the canonical surjection $`H^2(X,F^1\mathrm{\Omega }^{\text{}}(X))H^2(X,)`$. Together with the Gauss-Manin connection, this splitting defines an isomorphism
$$R^2\pi _{}F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}/S)H^2(X,F^1\mathrm{\Omega }^{\text{}}(X))𝒪_S$$
for every deformation $`\pi :\stackrel{~}{X}S`$. We will always assume given such a splitting, keeping in mind that this introduces into our constructions an element of choice. Note that this choice does not appear in the case of the affine $`X`$, since in this case for every deformation we have $`R^2\pi _{}𝒪(\stackrel{~}{X})=0`$ and $`R^2\pi _{}F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}/S)R^2\pi _{}\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}/S)`$.
### 3.2 Symplectic deformations and the period map.
Let $`X`$ be an admissible manifold. Assume from now on that the manifold $`X`$ is equipped with a non-degenerate closed $`2`$-form $`\mathrm{\Omega }\mathrm{\Omega }^2(X)`$.
###### Definition 3.3
A symplectic deformation $`\stackrel{~}{X}/S`$ of the symplectic manifold $`X`$ over a base $`S`$ is a usual deformation $`\pi :\stackrel{~}{X}S`$ equipped with a closed relative $`2`$-form $`\mathrm{\Omega }\mathrm{\Omega }^2(\stackrel{~}{X}/S)`$ which becomes the given $`2`$-form under the isomorphism $`o\times _S\stackrel{~}{X}X`$.
For every local Artin scheme $`S=\mathrm{Spec}A`$, we will denote by $`\mathrm{𝖣𝖾𝖿}(X,S)`$ or simply by $`\mathrm{𝖣𝖾𝖿}(S)`$ the set of isomorphism classes of symplectic deformations $`\stackrel{~}{X}/S`$ of $`X`$ over $`S`$.
Choose once and for all a splitting $`H^2(X,)H^2(X,F^1\mathrm{\Omega }^{\text{}}(X))`$, so that for every deformation $`\pi :\stackrel{~}{X}S`$ we have an isomorphism
$$R^2F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}/S)H^2(X,F^1\mathrm{\Omega }^{\text{}}(X))𝒪_S.$$
If the deformation $`\stackrel{~}{X}/S`$ is symplectic, then the relative $`2`$-form $`\mathrm{\Omega }\mathrm{\Omega }^2(\stackrel{~}{X}/S)`$ defines a canonical cohomology class
$$[\mathrm{\Omega }]H^2(X,F^1\mathrm{\Omega }^{\text{}}(X))𝒪_S.$$
This class gives a scheme map
$$\mathrm{𝖯𝖾𝗋}(\stackrel{~}{X}):STot(H^2(X,F^1\mathrm{\Omega }^{\text{}}(X)),$$
where $`Tot(H^2(X,F^1\mathrm{\Omega }^{\text{}}(X))`$ denotes the total space of $`H^2(X,F^1\mathrm{\Omega }^{\text{}}(X))`$ considered as a scheme. Further on, we shall simplify notation by omitting “$`Tot`$”.
###### Definition 3.4
The period domain of the admissible symplectic manifold $`X`$ is the completion of the vector space $`H^2(X,F^1\mathrm{\Omega }^{\text{}}(X))`$ near the point $`[\mathrm{\Omega }]H^2(X,F^1\mathrm{\Omega }^{\text{}}(X))`$ corresponding to the symplectic form $`\mathrm{\Omega }\mathrm{\Omega }^2(X)`$.
The map $`\mathrm{𝖯𝖾𝗋}(\stackrel{~}{X})`$ is called the period map of the deformation $`\stackrel{~}{X}/S`$.
Note that the period map by definition maps $`S`$ into the period domain $`\mathrm{𝖯𝖾𝗋}H^2(X,F^1\mathrm{\Omega }^{\text{}}(X))`$.
###### Remark 3.5
This definition of the period domain is essentially cheating: in Bogomolov’s theory the period domain is not a formal scheme, but a globally (and non-trivially) defined quadric in the projectivization of the vector space $`H^2(X,)`$. However, since we work only with infinitesemal deformation, Definition 3.4 is sufficient for our purposes.
For any local Artin scheme $`S`$, taking the period map defines a map
$$\mathrm{𝖯𝖾𝗋}:\mathrm{𝖣𝖾𝖿}(S)\mathrm{𝖯𝖾𝗋}(S)$$
from the set of deformation classes over $`S`$ to the set of $`S`$-points of the formal scheme $`\mathrm{𝖯𝖾𝗋}`$. This map (which we will, by abuse of the language, also call the period map) is functorial in $`S`$ (where $`\mathrm{𝖣𝖾𝖿}(S)`$ is considered as a functor by taking pullbacks). We can now state our main result.
###### Theorem 3.6
Let $`X`$ be an admissible manifold equipped with a symplectic $`2`$-form $`\mathrm{\Omega }\mathrm{\Omega }^2(X)`$. Then for any local Artin scheme $`S`$, the period map
$$\mathrm{𝖯𝖾𝗋}:\mathrm{𝖣𝖾𝖿}(S)\mathrm{𝖯𝖾𝗋}(S)$$
induces a set bijection from the set of isomorphism classes of symplectic deformations $`\stackrel{~}{X}/S`$ to the set of $`S`$-points of the period domain $`\mathrm{𝖯𝖾𝗋}`$.
In particular, there exists a (formal) symplectic deformation $`𝔛/\mathrm{𝖯𝖾𝗋}`$ such that any deformation $`\stackrel{~}{X}/S`$ is isomorphic to the pullback of $`𝔛`$ by means of the period map $`\mathrm{𝖯𝖾𝗋}(\stackrel{~}{X}):S\mathrm{𝖯𝖾𝗋}`$.
###### Remark 3.7
Let $`M`$ satisfy $`H^i(M,𝒪_N)=0`$ for $`i1`$. Then Theorem 3.6 holds. Moreover, the period map $`\mathrm{𝖯𝖾𝗋}:\mathrm{𝖣𝖾𝖿}(S)H^2(M)`$ is locally an isomorphism.
###### Remark 3.8
For $`M`$ affine, we obviously have $`H^i(M,𝒪_M)=0`$ for $`i1`$. $`H^1(𝒪_N)=H^2(𝒪_M)=0`$ obviously holds. In this case, Theorem 3.6 can be proved by a standard inductive argument (see Introduction).
## 4 Elementary extensions and the Kodaira-Spencer class.
### 4.1 Elementary extensions.
To prove Theorem 3.6, we will use induction on the length of the local Artin algebra $`A=𝒪(S)`$. To set up the induction, we introduce the following.
###### Definition 4.1
Let $`S_0S`$ be a closed embedding of local Artin schemes, and let $`\stackrel{~}{X}_0/S_0`$ be a symplectic deformation of a holomorphically symplectic manifold $`X`$ over $`S_0`$.
Then by $`\mathrm{𝖣𝖾𝖿}(\stackrel{~}{X}_0,S)`$ we will denote the set of isomorphism classes of symplectic deformations $`\stackrel{~}{X}/S`$ equipped with a symplectic isomorphism $`\stackrel{~}{X}_SS_0\stackrel{~}{X}_0`$.
This is consistent with our earlier notation. Indeed, by Definition 3.3, every symplectic deformation is canonically trivialized over the base point $`oS`$. Therefore $`\mathrm{𝖣𝖾𝖿}(X,S)`$ in the sense of Definition 4.1 is still the set of isomorphism classes of all symplectic deformations of the manifold $`X`$.
The corresponding notion on the “period” side is the following.
###### Definition 4.2
Let $`p_0:S_0\mathrm{𝖯𝖾𝗋}`$ be a map from the closed subscheme $`S_0S`$ to the period domain $`\mathrm{𝖯𝖾𝗋}`$. Then by $`\mathrm{𝖯𝖾𝗋}(p_0,S)`$ we will denote the set of all maps $`f:S\mathrm{𝖯𝖾𝗋}`$ such that $`F|_{S_0}=p_0`$.
It is customary in deformation theory to prove theorems step-by-step, starting with the case of square-zero extensions. However, we will need a slightly smaller class of extensions $`S_0S`$. Namely, let $`A`$ be a local Artin algebra, and let $`IA`$ be a square-zero ideal, so that $`I^2=0`$. Then we have the usual exact sequence of the modules of Kähler differentials over $``$
(4.1)
$$\begin{array}{ccccccc}I& & \mathrm{\Omega }^1(A)/I& & \mathrm{\Omega }^1(A/I)& & 0.\end{array}$$
###### Definition 4.3
The extension $`\mathrm{Spec}A/I\mathrm{Spec}A`$ will be called elementary if the sequence 4.1 is also exact on the left.
The following lemma immediately implies that every local Artin scheme $`S`$ admits a filtration $`oS_0\mathrm{}S_k=S`$ such that all extensions $`S_iS_{i+1}`$ are elementary.
###### Lemma 4.4
Let $`A`$ be a local Artin algebra with the maximal ideal $`𝔪`$. Assume that $`𝔪^{p+1}=0`$, while $`𝔪^p0`$. Then the extension $`\mathrm{Spec}A/𝔪^p\mathrm{Spec}A`$ is elementary.
###### Proof.
We have to prove that the canonical map $`𝔪^p\mathrm{\Omega }^1(A)/𝔪^p`$ is injective. It suffices to prove this for $`A`$ replaced with its associated graded quotient with respect to the $`𝔪`$-adic filtration. Thus we can assume that $`A`$ and $`\mathrm{\Omega }^1(A)`$ are graded. The map $`𝔪^p\mathrm{\Omega }^1(A)/𝔪^p`$ is the composition of the de Rham differential $`d:𝔪^p\mathrm{\Omega }^1(A)`$ and the projection $`\mathrm{\Omega }^1(A)\mathrm{\Omega }^1(A)/𝔪^p`$. Since $`d`$ is obviously injective, it suffices to prove that $`d(𝔪^p)𝔪^p\mathrm{\Omega }^1(A)=0`$. But this is trivial: $`d(𝔪^p)`$ has degree $`p`$ with respect to the grading on $`\mathrm{\Omega }^1(A)`$, while $`𝔪^p\mathrm{\Omega }^1(A)`$ is of degree $`(p+1)`$. $`\mathrm{}`$
This lemma reduces Theorem 3.6 to the following claim.
###### Proposition 4.5
Let $`X`$ be an admissible symplectic manifold. Let $`S_0S`$ be an elementary extension of local Artin schemes, and let $`\stackrel{~}{X}_0/S_0`$ be an arbitrary symplectic deformation of the manifold $`X`$. Denote by $`p_0:S_0\mathrm{𝖯𝖾𝗋}`$ the period map of the deformation $`\stackrel{~}{X}_0/S_0`$.
Then the period map
$$\mathrm{𝖯𝖾𝗋}:\mathrm{𝖣𝖾𝖿}(\stackrel{~}{X}_0,S)\mathrm{𝖯𝖾𝗋}(p_0,S)$$
is as isomorphism.
### 4.2 The Kodaira-Spencer class.
In order to start the proof of Proposition 4.5, we will need a convenient description of the set $`\mathrm{𝖯𝖾𝗋}(p_0,S)`$. To give such a description, we will consider not the period map itself, but its differential.
###### Definition 4.6
Let $`\stackrel{~}{X}/S`$ be a symplectic deformation of a symplectic manifold $`X`$. Then the class
$$\theta =\mathrm{\Omega }H^2(\stackrel{~}{X},F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}/S))_{𝒪(S)}\mathrm{\Omega }^1(S)$$
obtained by application of the Gauss-Manin connection $``$ to the relative $`2`$-form $`\mathrm{\Omega }`$ is called the Kodaira-Spencer class of the deformation $`\stackrel{~}{X}/S`$.
Note that the symplectic form $`\mathrm{\Omega }`$ is a cohomology class of the complex $`F^2\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}/S)`$. Since the Gauss-Manin connection decreases the Hodge filtration at most by $`1`$, the Kodaira-Spencer class $`\theta `$ is well-defined for an arbitrary symplectic manifold $`X`$.
When the symplectic manifold $`X`$ is admissible, the Kodaira-Spencer class essentially coincides with the codifferential of the period map $`\mathrm{𝖯𝖾𝗋}(\stackrel{~}{X}):S\mathrm{𝖯𝖾𝗋}`$. More precisely, the codifferential
(4.2)
$$\delta \mathrm{𝖯𝖾𝗋}(\stackrel{~}{X}):\mathrm{𝖯𝖾𝗋}^{}\mathrm{\Omega }^1(\mathrm{𝖯𝖾𝗋})\mathrm{\Omega }^1(S)$$
of the period map is given by
$$\delta \mathrm{𝖯𝖾𝗋}(\stackrel{~}{X})(\alpha )=\alpha ,\theta ,$$
where the bundle $`\mathrm{𝖯𝖾𝗋}^{}\mathrm{\Omega }^1(\mathrm{𝖯𝖾𝗋})`$ is identified with the trivial bundle
$$\left(H^2(X,F^1\mathrm{\Omega }^{\text{}}(X))\right)^{}𝒪_S,$$
$`\alpha `$ is an arbitrary section of this trivial bundle, and $`,`$ means the pairing on the first factor in
$$H^2(\stackrel{~}{X},F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}/S))_{𝒪(S)}\mathrm{\Omega }^1(S).$$
Let now $`i:S_0S`$ be an elementary extension, let $`\stackrel{~}{X}_0/S_0`$ be a symplectic deformation of a symplectic manifold $`X`$, and let
$$\theta _0H^2(\stackrel{~}{X}_0,F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}_0/S_0))_{𝒪(S_0)}\mathrm{\Omega }^1(S_0)$$
be its Kodaira-Spencer class. Denote by $`\eta :i^{}\mathrm{\Omega }^1(S)\mathrm{\Omega }^1(S_0)`$ the canonical surjection of the modules of differentials.
###### Definition 4.7
By $`KS(\theta _0,S)`$ we will denote that set of all cohomology classes
$$\theta H^2(\stackrel{~}{X}_0,F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}_0/S_0))_{𝒪(S_0)}i^{}\mathrm{\Omega }^1(S)$$
such that
$$\eta (\theta )=\theta _0H^2(\stackrel{~}{X}_0,F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}_0/S_0))_{𝒪(S_0)}\mathrm{\Omega }^1(S_0).$$
Assume that $`X`$ is admissible, so that for every symplectic deformation we have the period map, and denote by $`p_0:S_0\mathrm{𝖯𝖾𝗋}`$ the period map of the deformation $`\stackrel{~}{X}_0/S_0`$. It is the set $`KS(\theta _0,S)`$ which we will use as a model for the set $`\mathrm{𝖯𝖾𝗋}(p_0,S)`$. To do this, notice that every element $`p\mathrm{𝖯𝖾𝗋}(p_0,S)`$ defines an element $`\theta (p)KS(\theta _0,S)`$ by the formula 4.2. This correspondence is in fact one-to-one.
###### Lemma 4.8
The correspondence $`p\theta (p)`$ is a bijection between the set $`\mathrm{𝖯𝖾𝗋}(p_0,S)`$ and the set $`KS(\theta _0,S)`$.
###### Proof.
Indeed, both sets are torsors over the group
$$H^2(\stackrel{~}{X}_0,F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}_0/S_0))_{𝒪(S_0)}I,$$
where $`I𝒪(S)`$ is the kernel of the map $`𝒪(S)𝒪(S_0)`$. For the group $`\mathrm{𝖯𝖾𝗋}(p_0,S)`$, this is obvious from the exact sequence
$$\begin{array}{ccccccccc}0& & I& & 𝒪(S)& & 𝒪(S_0)& & 0,\end{array}$$
while for the set $`KS(\theta _0,S)`$ this follows from the exact sequence 4.1 of the differentials (exact on the left since the extension $`S_0S`$ is elementary). $`\mathrm{}`$
The advantages of the set $`KS(\theta _0,S)`$ over the set $`\mathrm{𝖯𝖾𝗋}(p_0,S)`$ are twofold. Firstly, Definition 4.7 can be refined so that it takes account of automorphisms of deformation of $`X`$. This will be explained in the next section. Secondly, and this we will state now, Definition 4.7 is essentially local on $`X`$. Namely, we have the following obvious fact.
###### Lemma 4.9
The set $`KS(\theta _0,S)`$ can be equivalently defined as the set of all maps
$$\theta \mathrm{Hom}(𝒪(S_0)[2],F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}_0/S_0)_{𝒪(S_0)}i^{}\mathrm{\Omega }^1(S))$$
in the derived category of sheaves of $`𝒪(S_0)`$-modules on $`X`$, such that $`\eta (\theta )=\theta _0`$. Here $`𝒪(S_0)`$ is considered as the constant sheaf. $`\mathrm{}`$
Because of this, we can use Definition 4.7 to restate Proposition 4.5 without the admissibility condition on the manifold $`X`$.
###### Proposition 4.10
Let $`S_0S`$ be an elementary extension of local Artin algebras. Let $`X`$ be a symplectic manifold, let $`\stackrel{~}{X}_0/S_0`$ be a symplectic deformation, and let
$$\theta _0\mathrm{Hom}(𝒪(S_0)[2],F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}_0/S_0)_{𝒪(S_0)}\mathrm{\Omega }^1(S_0))$$
be its Kodaira-Spencer class.
Then associating to a symplectic deformation $`\stackrel{~}{X}/S`$ its Kodaira-Spencer class defines a bijection
$$\mathrm{𝖯𝖾𝗋}(\stackrel{~}{X}_0,S):\mathrm{𝖣𝖾𝖿}(\stackrel{~}{X}_0,S)KS(\theta _0,S).$$
## 5 Localization, stacks and the proof of the main theorem.
### 5.1 Groupoids.
We can now begin the proof of Proposition 4.10, hence also of Theorem 3.6.
Assume given a symplectic manifold $`X`$, an elementary extension $`i:S_0S`$ and a symplectic deformation $`\stackrel{~}{X}_0/S_0`$. Denote by
$$\theta _0\mathrm{Hom}(𝒪(S_0)[2],F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}_0/S_0)_{𝒪(S_0)}\mathrm{\Omega }^1(S_0))$$
the Kodaira-Spencer class of the deformation $`\stackrel{~}{X}_0/S_0`$.
To begin with, we will refine the definitions of the sets $`\mathrm{𝖣𝖾𝖿}(\stackrel{~}{X}_0,S)`$ and $`KS(\theta _0,S)`$ and of the map
$$\mathrm{𝖯𝖾𝗋}(\stackrel{~}{X}_0,S):\mathrm{𝖣𝖾𝖿}(\stackrel{~}{X}_0,S)KS(\theta _0,S)$$
so as to take into account possible automorphisms of the deformations $`\stackrel{~}{X}\mathrm{𝖣𝖾𝖿}(\stackrel{~}{X}_0,S)`$. Such deformations naturally form a category. Moreover, this category is obviously a groupoid (i.e. all morphisms are invertible).
###### Definition 5.1
By $`𝒟\mathrm{𝑒𝑓}(\stackrel{~}{X}_0,S)`$ we will denote the groupoid of all symplectic deformations $`\stackrel{~}{X}/S`$ equipped with an isomorphism $`\stackrel{~}{X}\times _SS_0\stackrel{~}{X}_0`$.
The set $`\mathrm{𝖣𝖾𝖿}(\stackrel{~}{X}_0,S)`$ is by definition the set of objects in the groupoid $`𝒟\mathrm{𝑒𝑓}(\stackrel{~}{X}_0,S)`$.
To construct a groupoid version of the set $`KS(\theta _0,S)`$, note that the elements
$$\theta \mathrm{Hom}(𝒪(S_0)[2],F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}_0/S_0)_{𝒪(S_0)}\mathrm{\Omega }^1(S_0))$$
naturally classify exact sequences
$$\begin{array}{ccccccccc}0& & F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}_0/S_0)_{𝒪(S_0)}i^{}\mathrm{\Omega }^1(S)& & & & 𝒪(S_0)& & 0\end{array}$$
in the (abelian) category of complexes of sheaves of $`𝒪(S_0)`$-modules on $`X`$. Such an element satisfies $`\eta (\theta )=\theta _0`$ if and only if there exists a commutative diagram
(5.1)
$$\begin{array}{ccccccccc}0& & F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}_0/S_0)_{𝒪(S_0)}i^{}\mathrm{\Omega }^1(S)& & & & 𝒪(S_0)& & 0\\ & & \mathrm{𝗂𝖽}\eta & & & & & & & & \\ 0& & F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}_0/S_0)_{𝒪(S_0)}\mathrm{\Omega }^1(S_0)& & & & 𝒪(S_0)& & 0\end{array},$$
whose the bottom row is the exact sequence corresponding to the given class
$$\theta _0\mathrm{Hom}(𝒪(S_0)[2],F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}_0/S_0)_{𝒪(S_0)}\mathrm{\Omega }^1(S_0)),$$
and $``$ is the associated extension.
This motivates the following.
###### Definition 5.2
By $`𝒦𝒮(\theta _0,S)`$ we will denote the groupoid whose objects are commutative diagrams of the type 5.1, and whose morphisms are maps between these commutative diagrams identical everywhere except for $``$.
Again, by definition the set $`KS(\theta _0,S)`$ is the set of isomorphism classes of objects in the groupoid $`𝒦𝒮(\theta _0,S)`$. Moreover, the period map
$$\mathrm{𝖯𝖾𝗋}(\stackrel{~}{X}_0,S):\mathrm{𝖣𝖾𝖿}(\stackrel{~}{X}_0,S)KS(\theta _0,S)$$
lifts to a functor
(5.2)
$$𝒫\mathrm{𝑒𝑟}(\stackrel{~}{X}_0,S):𝒟\mathrm{𝑒𝑓}(\stackrel{~}{X}_0,S)𝒦𝒮(\theta _0,S).$$
This is not immediately obvious, since an element in a first $`\mathrm{Ext}`$-group defines a short exact sequence only up to a non-canonical isomorphism. However, in our case there is a canonical choice for a short exact sequence 5.1. Namely, for every deformation $`\pi :XS`$ we have a two-step filtration $`\pi ^{}\mathrm{\Omega }^1(S)\mathrm{\Omega }^1(X)`$ on the sheaf of Kähler differentials $`\mathrm{\Omega }^1(X)`$, with the quotient $`\mathrm{\Omega }^1(X/S)`$. This filtration induces a filtration on the total de Rham complex $`\mathrm{\Omega }^{\text{}}(X)`$. Taking only the top two quotients of this filtration, we obtain a canonical short exact sequence
(5.3)
$$\begin{array}{ccccccccc}0& & \pi ^{}\mathrm{\Omega }^1(S)\mathrm{\Omega }^{\text{}}(X/S)& & & & \mathrm{\Omega }^{\text{}+1}(X/S)& & 0\end{array}$$
of complexes on $`X`$ and the corresponding extension class
$$\eta \mathrm{Ext}^1(\mathrm{\Omega }^{\text{}+1}(X/S),\pi ^{}\mathrm{\Omega }^1(S)\mathrm{\Omega }^{\text{}}(X/S)).$$
The class $`\eta `$ essentially induces the Gauss-Manin connection: for every relative cohomology class
$$\alpha H^{\text{}}(X,\mathrm{\Omega }^{\text{}}(X/S))=\mathrm{Ext}^{\text{}}(𝒪(X),\mathrm{\Omega }^{\text{}}(X/S)),$$
the class
$$(\alpha )H^{\text{}}(X,\mathrm{\Omega }^{\text{}}(X/S))\mathrm{\Omega }^1(S)$$
is equal to $`\rho \alpha `$. Moreover, the sequence 5.3 is compatible with the Hodge filtration. In particular, we have a canonical short exact sequence
(5.4)
$$\begin{array}{ccccccccc}0& & \pi ^{}\mathrm{\Omega }^1(S)F^1\mathrm{\Omega }^{\text{}}(X/S)& & & & F^2\mathrm{\Omega }^{\text{}}(X/S)& & 0.\end{array}$$
Now, by definition of the Kodaira-Spencer class $`\theta `$ we have
$$\theta =(\mathrm{\Omega })=\eta \mathrm{\Omega }.$$
Thus the short exact sequence 5.1 corresponding to the deformation $`X/S`$ is obtained by composing the canonical short exact sequence 5.4 with the map $`𝒪(X)[2]F^2\mathrm{\Omega }^{\text{}}(X/S)`$ given by $`\mathrm{\Omega }`$.
The reader can easily see that this construction is completely functorial, so that we indeed have a functor 5.2. We will call it the period functor.
Our third and final reformulation of Theorem 3.6 is the following.
###### Proposition 5.3
Under the assumption of Proposition 4.10, the period functor
$$𝒫\mathrm{𝑒𝑟}(\stackrel{~}{X}_0,S):𝒟\mathrm{𝑒𝑓}(\stackrel{~}{X}_0,S)𝒦𝒮(\theta _0,S)$$
is an equivalence of categories.
This proposition implies Proposition 4.10, hence also Proposition 4.5 and Theorem 3.6.
### 5.2 Reduction to the affine case.
The following lemma explains why we introduced the groupoids.
###### Lemma 5.4
Assume that Proposition 5.3 holds for affine symplectic manifolds $`X`$. Then it holds for an arbitrary symplectic manifold $`X`$.
###### Proof.
For every open subset $`UX`$, the deformation $`\stackrel{~}{X}_0/S_0`$ induces a symplectic deformation $`\stackrel{~}{U}_0/S_0`$ of the manifold $`U`$. The collection of groupoids $`𝒟\mathrm{𝑒𝑓}(\stackrel{~}{U}_0,S)`$ is a stack on $`X`$ in Zariski topology.
Moreover, for every open $`UX`$ any diagram of the type 5.1 induces by restriction a diagram of the same type for the manifold $`U`$, and the collection of groupoids $`𝒦𝒮(\theta _0|_U,S)`$ is also a stack on $`X`$ in Zariski topology (it is to insure this that we have chosen to work with exact sequences of complexes on $`X`$ instead of using the derived category).
The period functors $`𝒫\mathrm{𝑒𝑟}(\stackrel{~}{U}_0,S):𝒟\mathrm{𝑒𝑓}(\stackrel{~}{U}_0,S)𝒦𝒮(\theta _0|_U,S)`$ is a functor between these stack. Since by assumption it is an equivalence for affine $`UX`$, it is an equivalence for an arbitrary $`UX`$ – in particular, for $`X`$ itself. $`\mathrm{}`$
### 5.3 The affine case.
It remains to prove Proposition 5.3 in the case of an affine symplectic manifold $`X`$. This is essentially done in the Introduction. Here we state the same argument in a more refined language.
Assume given an affine symplectic manifold $`X`$, an elementary extension $`S_0S`$, and a symplectic deformation $`\stackrel{~}{X}_0/S_0`$ with Kodaira-Spencer class $`\theta _0`$.
Since $`X`$ is affine and smooth, every deformation $`\stackrel{~}{X}/S`$ is trivial as an algebraic manifold. Thus we can choose an isomorphism $`\stackrel{~}{X}_0X\times S_0`$, and every object $`\stackrel{~}{X}\mathrm{𝖮𝖻}𝒟\mathrm{𝑒𝑓}(\stackrel{~}{X}_0,S)`$ is isomorphic as an algebraic manifold to $`X\times S`$. We introduce the following notation.
###### Definition 5.5
If a group $`G`$ acts on a set $`N`$, then the quotient groupoid $`N/G`$ is the groupoid whose objects are elements $`nN`$ and whose morphisms are given by
$$\mathrm{Hom}(n_1,n_2)=\{gGgn_1=n_2\},n_1,n_2N.$$
Then the groupoid $`𝒟\mathrm{𝑒𝑓}(\stackrel{~}{X}_0,S)`$ is equivalent to the quotient groupoid
$$\mathrm{𝖲𝗒𝗆𝗉𝗅}/\mathrm{𝖠𝗎𝗍},$$
where $`\mathrm{𝖲𝗒𝗆𝗉𝗅}`$ is the set of all relative closed $`2`$-forms $`\mathrm{\Omega }\mathrm{\Omega }^2(X\times S/S)`$ whose restriction to $`\stackrel{~}{X}_0=X\times S_0X\times S`$ coincides with the given symplectic form $`\mathrm{\Omega }_0\mathrm{\Omega }^2(\stackrel{~}{X}_0/S_0)`$, and $`\mathrm{𝖠𝗎𝗍}`$ is the group of automorphisms of the algebraic manifold $`X\times S`$ which commute with the projection $`X\times SS`$ and which are identical on $`\stackrel{~}{X}_0X\times S`$.
Every form $`\mathrm{\Omega }\mathrm{𝖲𝗒𝗆𝗉𝗅}`$ can be represented as
$$\mathrm{\Omega }=\mathrm{\Omega }_0+\beta ,$$
where $`\beta \mathrm{\Omega }^2(X)I`$ is a closed $`2`$-form on $`X`$ with values in the ideal $`I=\mathrm{Ker}(𝒪(S)𝒪(S_0)`$. Therefore we have a canonical identification $`\mathrm{𝖲𝗒𝗆𝗉𝗅}=\mathrm{\Omega }_{cl}^2(X)I`$ of $`\mathrm{𝖲𝗒𝗆𝗉𝗅}`$ with the space $`\mathrm{\Omega }_{cl}^2X)I`$ of closed $`I`$-valued $`2`$-forms on $`X`$. Moreover, every automorphism $`g\mathrm{𝖠𝗎𝗍}`$ is an automorphism of the function ring $`𝒪(X\times S)`$ of the form
$$g=\mathrm{𝗂𝖽}+\xi ,$$
where $`\xi 𝒯(X)I`$ is an $`I`$-valued vector field on $`X`$. Since $`I𝒪(S)`$ is a square-zero ideal, the group $`\mathrm{𝖠𝗎𝗍}`$ is commutative and isomorphic to the abelian group $`𝒯(X)I`$ – that is, we have
$$(\mathrm{𝗂𝖽}+\xi _1)(\mathrm{𝗂𝖽}+\xi _2)=\mathrm{𝗂𝖽}+\xi _1+\xi _2$$
for every $`\xi _1,\xi _2𝒯(X)I`$. Finally, by Cartan homotopy formula, the action of $`\mathrm{𝖠𝗎𝗍}`$ on $`\mathrm{𝖲𝗒𝗆𝗉𝗅}`$ is given by
$$(\mathrm{𝗂𝖽}+\xi )(\mathrm{\Omega }_0+\beta )=\mathrm{\Omega }_0+\beta +d(\mathrm{\Omega }_0\mathrm{}\xi ).$$
where $`d`$ is the de Rham differential (all the other terms vanish since $`I^2=0`$). To sum up, we have
$$𝒟\mathrm{𝑒𝑓}(\stackrel{~}{X}_0,S)\mathrm{\Omega }_{cl}^2(X)I/𝒯(X)I,$$
and the action is given by
$$\xi \beta =\beta +\mathrm{\Omega }_0\mathrm{}\xi .$$
We will now describe the right-hand side of the hypothetical equivalence 5.2 – that is, the groupoid $`𝒦𝒮(\theta _0,S)`$ – in a similar way. To do this, notice that since $`X`$ is affine, we can replace sheaves on $`X`$ with the modules of their global sections. Moreover, the trivialization $`\stackrel{~}{X}_0X\times S_0`$ provides identifications
$`F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}_0/S_0)_{𝒪(S_0)}i^{}\mathrm{\Omega }^1(S)`$ $`F^1\mathrm{\Omega }^{\text{}}(X)_{}\mathrm{\Omega }^1(S)/I,`$
$`F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}_0/S_0)_{𝒪(S_0)}\mathrm{\Omega }^1(S_0)`$ $`F^1\mathrm{\Omega }^{\text{}}(X)_{}\mathrm{\Omega }^1(S_0),`$
of the complexes in the left column of a commutative diagram of type 5.1.
In every commutative diagram of complexes of this type, both rows split as exact sequences of graded vector spaces. The only possibly non-trivial extension data are contained in the differential of the complex $``$. Denote the set of possible differentials by $`\mathrm{𝖣𝗂𝖿𝖿}`$. Analogously, every map between two diagrams of the type 5.1 must be an automorphism of $``$ which is upper-triangular with respect to the splitting
$$=(F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}_0/S_0)_{𝒪(S_0)}\mathrm{\Omega }^1(S)/I)𝒪(S_0).$$
If we denote the group of all such automorphisms by $`\mathrm{𝖳𝗋}`$, then the groupoid $`𝒦𝒮(\theta _0,S)`$ is equivalent to the quotient groupoid $`\mathrm{𝖣𝗂𝖿𝖿}/\mathrm{𝖳𝗋}`$.
To identify the set $`\mathrm{𝖣𝗂𝖿𝖿}`$, note that every possible differential $`\mathrm{𝖣𝗂𝖿𝖿}`$, being a map of $`𝒪(S_0)`$-modules, is completely determined by its value
$$(1)\mathrm{\Omega }^2(X)\mathrm{\Omega }^1(S)/I$$
on the unity $`1𝒪(S_0)`$. Since the differential in the complex $``$ from the bottom row of 5.1 is fixed, every two such differentials $`_1,_2\mathrm{𝖣𝗂𝖿𝖿}`$ must differ by a $`2`$-form
$$\beta =_1(1)_2(1)\mathrm{\Omega }^2(X)I\mathrm{\Omega }^2(X)\mathrm{\Omega }^1(S)/I.$$
Moreover, since every differential $`\mathrm{𝖣𝗂𝖿𝖿}`$ must satisfy $`^2=0`$, the difference $`\beta =_1(1)_2(1)`$ must be a closed $`I`$-valued $`2`$-form. Thus the set $`\mathrm{𝖣𝗂𝖿𝖿}`$ is a torsor over the abelian group $`\mathrm{\Omega }_{cl}^2(X)I`$.
Analogously, every triangular map $`g\mathrm{𝖳𝗋}`$ is of the form
$$g=\left(\begin{array}{cc}\mathrm{𝗂𝖽}& a\\ 0& \mathrm{𝗂𝖽}\end{array}\right)$$
for some
$$a\mathrm{Hom}(𝒪(S_0),\mathrm{\Omega }^1(X)I)\mathrm{Hom}(𝒪(S_0),\mathrm{\Omega }^1(X)\mathrm{\Omega }^1(S)/I),$$
and composition of maps $`g_1`$, $`g_2`$ simply adds the associated elements $`a_1`$, $`a_2`$. Since $`a`$ must be a map of $`𝒪(S_0)`$-modules, it is completely determined by
$$\alpha =a(1)\mathrm{\Omega }^1(X)I.$$
Therefore $`\mathrm{𝖳𝗋}`$ is the abelian group $`\mathrm{\Omega }^1(X)I`$. Under these identifications, the action of $`\mathrm{𝖳𝗋}`$ on $`\mathrm{𝖣𝗂𝖿𝖿}`$ is given by
$$\alpha \beta =\beta +d\alpha ,$$
where $`d`$ is the de Rham differential.
Having done these identifications, we notice that the period functor $`𝒫\mathrm{𝑒𝑟}:𝒟\mathrm{𝑒𝑓}(\stackrel{~}{X}_0,S)𝒦𝒮(\theta _0,S)`$ is defined on the level of quotient groupoids by maps
$`\mathrm{𝖲𝗒𝗆𝗉𝗅}`$ $`\mathrm{𝖣𝗂𝖿𝖿}`$
$`\mathrm{𝖠𝗎𝗍}`$ $`\mathrm{𝖳𝗋}`$
The first map is tautologically an isomorphism $`\mathrm{\Omega }_{cl}^2(X)I\mathrm{\Omega }_{cl}^2(X)I`$. The second is the map
$$𝒯(X)I\mathrm{\Omega }^1(X)I$$
given by $`\xi \mathrm{\Omega }_0\mathrm{}\xi `$. It is an isomorphism because the symplectic $`2`$-form $`\mathrm{\Omega }_0\mathrm{\Omega }^2(\stackrel{~}{X}_0/S_0)`$ is by assumption non-degenerate.
This finishes the proof of Proposition 5.3, hence, by a long chains of reductons, also proves Theorem 3.6. $`\mathrm{}`$
## 6 Postface
To finish the paper, we would to give a few comments (perhaps a bit vague) as to how the theory of symplectic deformations is related to the usual deformation theory, and what is the relation of this paper to the known results.
The period map is a purely symplectic phenomenon – it has no analogs in the usual deformation theory (although there might be a useful version for the deformation theory of Calabi-Yau manifolds). On the other hand, its differential, which we called the Kodaira-Spencer class, is a fairly general thing. Namely, recall that for every smooth family $`\pi :XS`$, say over an affine base $`S`$, there exist a canonical class
$$\theta \mathrm{Ext}^1(\mathrm{\Omega }^1(X/S),\pi ^{}\mathrm{\Omega }^1(S)),$$
– the extension class given by the exact sequence
$$\begin{array}{ccccccccc}0& & \pi ^{}\mathrm{\Omega }^1(S)& & \mathrm{\Omega }^1(X)& & \mathrm{\Omega }^1(X/S)& & 0\end{array}$$
of differentials for the map $`\pi :XS`$. Since $`X/S`$ is smooth, the sheaf $`\mathrm{\Omega }^1(X/S)`$ is flat, and this class can be reinterpreted as a class in the first cohomology group
$$H^1(X,𝒯(X/S)\pi ^{}\mathrm{\Omega }^1(S))H^1(X,𝒯(X/S))\mathrm{\Omega }^1(S)$$
of the relative tangent bundle $`𝒯(X/S)`$. It is this class that is usually called the Kodaira-Spencer class of the deformation $`X/S`$.
If the deformation $`X/S`$ is symplectic, then the symplectic form $`\mathrm{\Omega }\mathrm{\Omega }^2(X/S)`$ identifies the relative tangent bundle $`𝒯(X/S)`$ with the relative cotangent bundle $`\mathrm{\Omega }^1(X/S)`$ by means of the correspondence $`\xi \mathrm{\Omega }\mathrm{}\xi `$. This makes it possible to compare the usual Kodaira-Spencer class
$$\theta H^1(X,𝒯(X/S))\mathrm{\Omega }^1(S)H^1(X,\mathrm{\Omega }^1(X/S)\mathrm{\Omega }^1(S))$$
with the “symplectic” Kodaira-Spencer class $`\stackrel{~}{\theta }`$ introduced in Definition 4.6. These classes essentially coincide:
###### Lemma 6.1
For every symplectic deformation $`X/S`$, the usual Kodaira-Spencer class
$$\theta =\xi \mathrm{}\mathrm{\Omega }H^1(X,\mathrm{\Omega }^1(X/S)\mathrm{\Omega }^1(S))$$
is obtained from the symplectic Kodaira-Spencer class
$$\stackrel{~}{\theta }H^2(X,F^1\mathrm{\Omega }^{\text{}}(X/S))\mathrm{\Omega }^1(S))=H^1(X,F^1\mathrm{\Omega }^{\text{}}(X/S)[1])\mathrm{\Omega }^1(S))$$
by the tautological projection
$$F^1\mathrm{\Omega }^{\text{}}(X/S)[1]\mathrm{\Omega }^1(X/S).$$
###### Proof.
Indeed, by definition of the Gauss-Manin connection $``$ and the usual Kodaira-Spencer class $`\xi H^1(X,𝒯(X/S))\mathrm{\Omega }^1(S)`$, for every smooth family $`X/S`$ and every global relative $`k`$-form $`\alpha \mathrm{\Omega }^k(X/S)`$ we have
$$(\alpha )=\alpha \mathrm{}\xi .$$
Applying this to our symplectic deformation $`X/S`$ and to the symplectic form $`\mathrm{\Omega }`$ gives the result. $`\mathrm{}`$
In fact, Definition 4.7 of the groupoid $`𝒦𝒮(\stackrel{~}{X}_0,S)`$ is also general and works in the usual deformation theory. Moreover, the corresponding versions of Proposition 4.10 and Proposition 5.3 are also true (and proofs are more or less the same). However, in the usual case we do not have the period map, and the groupoid $`𝒦𝒮(\stackrel{~}{X}_0,S)`$ is not easy to describe. In particular, it might be empty – that is, there might be a homological obstruction to the existence of a commutative diagram of type 5.1. This is a very well-known phenomenon. In the symplectic case, the existence of the period map ensures that (for admissible manifolds) there are no obstructions for deformation at any step.
When one tries to describe usual deformations by means of the associated Kodaira-Spencer class (instead of a period map which does not exist in general), one enters a closed loop: the class $`\theta `$, which theoretically should uniquely define a deformation $`X/S`$, lies in the group $`H^1(X,𝒯(X/S))`$ which itself depends on the deformation. The main technical idea in our proof of Theorem 3.6 is to avoid it by going step-by-step through elementary extensions and using the exact sequence of differentials
$$\begin{array}{ccccccccc}0& & I& & \mathrm{\Omega }^1(A)/I& & \mathrm{\Omega }^1(A/I)& & 0\end{array}$$
for an elementary extension $`\mathrm{Spec}A/I\mathrm{Spec}A`$. This allows one to describe extension of a deformation $`X_0/\mathrm{Spec}(A/I)`$ to a deformation $`X/\mathrm{Spec}A`$ in terms of the lifting of the Kodaira-Spencer class from $`\mathrm{\Omega }^1(A/I)`$ to $`\mathrm{\Omega }^1(A)/I`$ – and this works, because the module $`\mathrm{\Omega }^1(A)/I`$ is already defined over $`A/I`$. This idea (at least for one-parameter elementary extensions $`\mathrm{Spec}[t]/t^k\mathrm{Spec}[t]/t^{k+1}`$) is due entirely to Z. Ran \[R\]. We believe that it is this technique that he called the $`T_1`$-lifting property.
We would also like to notice that all the obstructions to symplectic deformations vanish essentially because for an admissible manifold $`X`$, the sheaf $`R^2\pi _{}(\stackrel{~}{X},F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}/S))`$ is flat on $`S`$ for every deformation $`\pi :XS`$. The same thing happens in the Ran’s proof of the Tian-Todorov Lemma – that is, the lack of obstructions for deformations of a compact Calabi-Yau manifold is due to the flatness of some canonically defined sheaves. However, Ran works with the usual deformations, and he (in the notation as above) needs the flatness of two sheaves: $`\pi _{}\mathrm{\Omega }^n(\stackrel{~}{X}/S)`$ and $`R^1\pi _{}\mathrm{\Omega }^{n1}(\stackrel{~}{X}/S)`$ (here $`n=\mathrm{𝖽𝗂𝗆}X`$). This flatness is provided by Hodge theory. In the symplectic version of the proof, one would use instead the sheaves $`\pi _{}\mathrm{\Omega }^2(\stackrel{~}{X}/S)`$ and $`R^1\pi _{}\mathrm{\Omega }^1(\stackrel{~}{X}/S)`$.
From this point of view, the only new thing in our paper is the following observation: if one agrees to consider deformations that are a priori symplectic, then one can combine $`\pi _{}\mathrm{\Omega }^2(\stackrel{~}{X}/S)`$ and $`R^1\pi _{}\mathrm{\Omega }^1(\stackrel{~}{X}/S)`$ into $`R^2\pi _{}F^1\mathrm{\Omega }^{\text{}}(\stackrel{~}{X}/S)`$ – and the latter sheaf is flat for a much wider class of manifolds $`X`$.
Finally, there’s another, perhaps more conceptual explanation for the exceptional role played by the complex $`F^1\mathrm{\Omega }^{\text{}}(X)`$ in the symplectic deformation theory. This explanation comes from the general deformation theory of algebras over an operad, sketched for example in \[G1\], \[G2\]. Symplectic manifolds can not be described by operad. However, a more general class of Poisson manifolds admits such a description. The general deformation theory for algebras over an operad works in a way completely parallel to the usual one, but the tangent bundle (more generally, the tangent complex) is replaced by its operadic version. For the Poisson operad and a symplectic manifold $`X`$, the Poisson tangent complex is precisely $`F^1\mathrm{\Omega }^{\text{}}(X)`$ (independently of the symplectic form). For a more general Poisson structure, this is the complex $`\mathrm{\Lambda }^1𝒯(X)`$ of polyvector fields of degree $`1`$ on $`X`$, and the differential is given by the commutator with the Poisson bivector field $`\mathrm{\Theta }\mathrm{\Lambda }^2𝒯(X)`$.
Acknowledgements: This work is greatly influenced by T.Pantev, who explained to us the algebraic proof of Bogomolov-Tian-Todorov theorem. The second author is grateful to D. Kazhdan, M. Kontsevich, Yu.I.Manin, S. Merkulov and A. Todorov for interesting talks on the deformation theory.
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# 1 Introduction
## 1 Introduction
A key ingredient in the study of CP violation in the Standard Model is the theoretical prediction of the $`K^0`$$`\overline{K}^0`$ mixing amplitude. This involves the computation of the $`\mathrm{\Delta }S=2`$ matrix element
$$\overline{K}^0|O^{\mathrm{\Delta }S=2}|K^0\frac{8}{3}f_K^2m_K^2B_K$$
(1)
of the operator
$$O^{\mathrm{\Delta }S=2}O_1=(\overline{s}^A\gamma _\mu (1\gamma _5)d^A)(\overline{s}^B\gamma _\mu (1\gamma _5)d^B),$$
(2)
where $`s`$ and $`d`$ stand for strange and down quarks and $`A,B`$ are colour indices.
In general, important information on the physics beyond the Standard Model, such as various SUSY extensions (MSSM, NMSSM,…), can be obtained by studying $`\mathrm{\Delta }F=2`$ transitions (see and references therein for a discussion). For neutral kaons, such processes require, besides $`O_1`$, also the knowledge of the matrix elements of the operators (we adopt here the notation of ref. )
$`O_2=(\overline{s}^A(1\gamma _5)d^A)(\overline{s}^B(1\gamma _5)d^B)`$
$`O_3=(\overline{s}^A(1\gamma _5)d^B)(\overline{s}^B(1\gamma _5)d^A)`$
$`O_4=(\overline{s}^A(1\gamma _5)d^A)(\overline{s}^B(1+\gamma _5)d^B)`$ (3)
$`O_5=(\overline{s}^A(1\gamma _5)d^B)(\overline{s}^B(1+\gamma _5)d^A)`$
On the lattice, the matrix elements of the four-fermion operators above are typically extracted from the large-time asymptotic behaviour of three-point correlation functions of the form $`K_P^0(x_1)O^{\mathrm{\Delta }S=2}(0)K_P^0(x_2)`$, where $`K_P^0`$ are pseudoscalar sources with suitable quark flavour, $`K_P^0(x)=\overline{d}^A(x)\gamma _5s^A(x)`$. Expressed in terms of traces of quark propagators, these correlation functions correspond to the so-called “eight”- shaped quark diagrams given in fig. 1.
In order to obtain the physical amplitudes, it is necessary to compute the matrix elements of the renormalized operators corresponding to those defined in eqs. (2) and (3). With Wilson-like fermions, the renormalization procedure is complicated by the presence of explicit chiral-symmetry breaking ($`\chi SB`$) in the lattice fermion action: because of the Wilson term, dimension-six operators belonging to different chiral representations can mix with each other.
The testing ground for the restoration of chiral symmetry has been the chiral behaviour of $`\overline{K}^0|O^{\mathrm{\Delta }S=2}|K^0`$, which, if properly renormalized, vanishes when the $`K`$-meson becomes massless . Although several attempts with Wilson fermions have given reasonable measurements of $`B_K`$, it remains true that the control of the renormalization of the relevant operator is rather problematic <sup>1</sup><sup>1</sup>1$`B_K`$ has also been obtained with staggered fermions mainly in the quenched approximation (see ref. for a review). The (surviving) chiral symmetry in the staggered fermion formalism ensures the vanishing of the relevant matrix element in the chiral limit..
The root of the problem is the operator subtraction outlined above. $`O^{\mathrm{\Delta }S=2}`$ mixes with other operators $`O_i`$ of the same dimension but with “wrong naïve chirality”. Thus, the ($`\mu `$-dependent) $`K^0`$$`\overline{K}^0`$ matrix element of the renormalized operator $`\widehat{O}^{\mathrm{\Delta }S=2}`$ is given in terms of the ($`a`$-dependent) bare matrix elements by:
$$\overline{K}^0|\widehat{O}^{\mathrm{\Delta }S=2}(\mu )|K^0=\underset{a0}{lim}\overline{K}^0|Z_0^{\mathrm{\Delta }S=2}(a\mu ,g_0^2)\left[O^{\mathrm{\Delta }S=2}(a)+\underset{i}{}\mathrm{\Delta }_i(g_0^2)O_i(a)\right]|K^0$$
(4)
The overall renormalization constant $`Z_0^{\mathrm{\Delta }S=2}(a\mu ,g_0^2)`$ is logarithmically divergent and its determination does not affect the chiral behaviour of the matrix element. The constants $`\mathrm{\Delta }_i(g_0^2)`$’s are finite mixing coefficients which depend on the lattice bare coupling (also expressed as $`\beta =6/g_0^2`$) only . They have been either computed in perturbation theory (PT) , or fixed by using the non-perturbative (NP) method of refs. , or determined using the relevant Ward identities (WI) on quark states (QWI) <sup>2</sup><sup>2</sup>2 Another method, suggested in refs. , is the use of gauge-invariant WI on hadronic states. This method has never been implemented though..
When PT is used for the calculation of the mixing coefficients $`\mathrm{\Delta }_i`$, the mass dependence of the renormalized-operator matrix element, $`\overline{K}^0|\widehat{O}^{\mathrm{\Delta }S=2}|K^0`$, shows large deviations from the expected chiral behaviour. Past experience suggests, instead, that the chiral behaviour of $`\overline{K}^0|\widehat{O}^{\mathrm{\Delta }S=2}|K^0`$ is satisfactory if a non-perturbative method, either the NP renormalization or the QWI , is used.
In spite of this progress, the determination of the mixing coefficients is a long and painful procedure. The non-perturbative renormalization techniques proposed in the past have their limitations and, for this reason, an accurate and systematic study of the uncertainties in the determination of the mixing coefficients would be necessary before keeping the final error on the matrix element fully under control.
In this paper, we propose a new method which allows the calculation of $`\mathrm{\Delta }S=2`$ amplitudes with Wilson fermions without determining the mixing coefficients. The method is based on the lattice Ward identities and it is non-perturbative. It can be applied to any Wilson-like formulation of the action (Wilson action, tree-level improved, NPI and alike). Its accuracy, in the absence of improvement (of the action and of the operators), is of $`𝒪(a)`$, which is the best possible accuracy, attainable only with a perfect determination of the mixing coefficients $`\mathrm{\Delta }_i`$ in previous approaches.
The same method can be used to compute $`\mathrm{\Delta }I=3/2`$ $`K`$-$`\pi `$ matrix elements, e.g. those relative to the electro-penguin operators, and $`\mathrm{\Delta }I=3/2`$ parity-conserving amplitudes in hyperon decays. It fails, unfortunately, for $`\mathrm{\Delta }I=1/2`$ transitions because of the mixing with lower-dimension operators.
An alternative approach, based on the twisted mass QCD formalism , also allows the determination of the $`\mathrm{\Delta }S=2`$ mixing amplitude without lattice subtractions .
## 2 Description of the Method
In this section, we show how it is possible to compute the physical $`K^0`$$`\overline{K}^0`$ mixing amplitude without knowing, or determining, the mixing coefficients induced by the explicit $`\chi SB`$ of the lattice action. We first recall some basic notions about operator renormalization and mixing with Wilson fermions and then explain our proposal.
### 2.1 Operator Renormalization with Wilson Fermions
Following (see also and ), we classify the complete basis of dimension-six, four-fermion operators which mix under renormalization, relying on general symmetry arguments based on the vector-flavour symmetry, which survives on the lattice. To this purpose it is convenient to consider separately the parity-even and parity-odd parts of the operators in eqs. (2) and (3). Thus, for example, the parity-even part of the operator $`O_1=Q_1𝒬_1`$ is given by
$$Q_1=O_{[VV+AA]}=(\overline{s}^A\gamma _\mu d^A)(\overline{s}^B\gamma _\mu d^B)+(\overline{s}^A\gamma _\mu \gamma _5d^A)(\overline{s}^B\gamma _\mu \gamma _5d^B),$$
(5)
whereas the parity odd is defined as
$$𝒬_1=O_{[VA+AV]}=2(\overline{s}^A\gamma _\mu d^A)(\overline{s}^B\gamma _\mu \gamma _5d^B).$$
(6)
On the basis of $`CPS`$ symmetries, it can been shown that the renormalization of the $`𝒬_i`$ operators is not affected by the explicit $`\chi SB`$ of the lattice action and proceeds exactly as in the “continuum” theory. The corresponding renormalization matrix $`𝒵_{ij}`$ (which obviously depends on the renormalization-scheme) is a block diagonal matrix :
$$\left(\begin{array}{c}\widehat{𝒬}_1\\ \widehat{𝒬}_2\\ \widehat{𝒬}_3\\ \widehat{𝒬}_4\\ \widehat{𝒬}_5\end{array}\right)=\left(\begin{array}{ccccc}\hfill 𝒵_{11}& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 𝒵_{22}& \hfill 𝒵_{23}& \hfill 0& \hfill 0\\ \hfill 0& \hfill 𝒵_{32}& \hfill 𝒵_{33}& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 𝒵_{44}& \hfill 𝒵_{45}\\ \hfill 0& \hfill 0& \hfill 0& \hfill 𝒵_{54}& \hfill 𝒵_{55}\end{array}\right)\left(\begin{array}{c}𝒬_1\\ 𝒬_2\\ 𝒬_3\\ 𝒬_4\\ 𝒬_5\end{array}\right).$$
(7)
Thus the lattice does not induce extra subtractions ($`\mathrm{\Delta }_i`$) for the parity-odd sector ($`𝒬_k`$; $`k=1,\mathrm{},5`$), since the mixing in eq. (7) is the same as it would appear in the absence of $`\chi SB`$. In particular, $`𝒬_1`$ renormalizes multiplicatively.
The parity-even sector, on the contrary, is not protected by $`CPS`$ symmetries, and all the five relevant operators get mixed because of the $`\chi SB`$ of the lattice action. In this case, it is convenient to separate the operator mixing into two classes: i) the first which consists in correcting the operator mixing induced by the breaking of chiral symmetry; ii) the second is the renormalization which survives in the continuum limit. In the absence of explicit $`\chi SB`$, the mixing structure is the same as the one considered above for the parity-odd counterparts. The corresponding parity-even operators ($`\stackrel{~}{Q}_k`$; $`k=1,\mathrm{},5`$) would renormalize according to:
$$\left(\begin{array}{c}\widehat{Q}_1\\ \widehat{Q}_2\\ \widehat{Q}_3\\ \widehat{Q}_4\\ \widehat{Q}_5\end{array}\right)=\left(\begin{array}{ccccc}\hfill Z_{11}& \hfill 0& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill Z_{22}& \hfill Z_{23}& \hfill 0& \hfill 0\\ \hfill 0& \hfill Z_{32}& \hfill Z_{33}& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill Z_{44}& \hfill Z_{45}\\ \hfill 0& \hfill 0& \hfill 0& \hfill Z_{54}& \hfill Z_{55}\end{array}\right)\left(\begin{array}{c}\stackrel{~}{Q}_1\\ \stackrel{~}{Q}_2\\ \stackrel{~}{Q}_3\\ \stackrel{~}{Q}_4\\ \stackrel{~}{Q}_5\end{array}\right),$$
(8)
where the $`\stackrel{~}{Q}_i`$ represent the bare operators, which transform as elements of irreducible representations of the chiral group (obviously up to terms of $`𝒪(a)`$). In the presence of the Wilson term, the $`\stackrel{~}{Q}_i`$ are defined as
$$\left(\begin{array}{c}\stackrel{~}{Q}_1\\ \stackrel{~}{Q}_2\\ \stackrel{~}{Q}_3\\ \stackrel{~}{Q}_4\\ \stackrel{~}{Q}_5\end{array}\right)=\left(\begin{array}{c}Q_1\\ Q_2\\ Q_3\\ Q_4\\ Q_5\end{array}\right)+\left(\begin{array}{ccccc}\hfill 0& \hfill \mathrm{\Delta }_{12}& \hfill \mathrm{\Delta }_{13}& \hfill \mathrm{\Delta }_{14}& \hfill \mathrm{\Delta }_{15}\\ \hfill \mathrm{\Delta }_{21}& \hfill 0& \hfill 0& \hfill \mathrm{\Delta }_{24}& \hfill \mathrm{\Delta }_{25}\\ \hfill \mathrm{\Delta }_{31}& \hfill 0& \hfill 0& \hfill \mathrm{\Delta }_{34}& \hfill \mathrm{\Delta }_{35}\\ \hfill \mathrm{\Delta }_{41}& \hfill \mathrm{\Delta }_{42}& \hfill \mathrm{\Delta }_{43}& \hfill 0& \hfill 0\\ \hfill \mathrm{\Delta }_{51}& \hfill \mathrm{\Delta }_{52}& \hfill \mathrm{\Delta }_{53}& \hfill 0& \hfill 0\end{array}\right)\left(\begin{array}{c}Q_1\\ Q_2\\ Q_3\\ Q_4\\ Q_5\end{array}\right).$$
(9)
In other words, first the lattice subtraction is performed, followed by the renormalization of the remaining logarithmic divergencies. The above mixing pattern is abbreviated, in matrix form, as
$`\widehat{Q}`$ $`=`$ $`Z\stackrel{~}{Q}`$
$`\stackrel{~}{Q}`$ $`=`$ $`[I+\mathrm{\Delta }]Q`$ (10)
where $`I`$ is the $`5\times 5`$ unit matrix.
### 2.2 $`K^0`$$`\overline{K}^0`$ Mixing without Subtractions
For the sake of illustration, we discuss here the determination of the $`K^0`$$`\overline{K}^0`$ matrix element of the operator $`O^{\mathrm{\Delta }S=2}=O_1`$ only. The extension to the other operators ($`O_i`$, $`i=2,\mathrm{}5`$) is straightforward.
Le us consider the Ward identities which can be derived from the $`\tau _3`$ axial rotation
$`\delta u=\gamma _5u,\delta \overline{u}=\overline{u}\gamma _5,`$
$`\delta d=\gamma _5d,\delta \overline{d}=\overline{d}\gamma _5,`$ (11)
where $`u`$ and $`d`$ are the up and down quarks taken with degenerate masses, $`m_u=m_d=m`$. For further use, we also introduce the following bilinear operators <sup>3</sup><sup>3</sup>3 As kaon source, we may as well use the fourth component of the axial current $`\overline{d}\gamma _0\gamma _5s`$ instead of the pseudoscalar density.
$`\mathrm{\Pi }^0(x)`$ $`=`$ $`\overline{d}(x)\gamma _5d(x)\overline{u}(x)\gamma _5u(x),K_P^0(t)={\displaystyle \underset{\stackrel{}{x}}{}}\overline{d}(\stackrel{}{x},t)\gamma _5s(\stackrel{}{x},t),`$
$`K_S^0(t)`$ $`=`$ $`{\displaystyle \underset{\stackrel{}{x}}{}}\overline{d}(\stackrel{}{x},t)s(\stackrel{}{x},t),`$ (12)
and the corresponding renormalized quantities $`\widehat{K}_P^0(t)=Z_PK_P^0(t)`$ and $`\widehat{K}_S^0(t)=Z_SK_S^0(t)`$. The useful Ward identity in our case is then
$`\delta \left[\widehat{K}_P^0(t_1)\widehat{O}_1(0)\widehat{K}_P^0(t_2)\right]\delta \left[S\right]\widehat{K}_P^0(t_1)\widehat{O}_1(0)\widehat{K}_P^0(t_2)=0,`$ (13)
where $`\delta \left[\mathrm{}\right]`$ denotes the rotation of the argument of $`\delta `$ and $`\delta \left[S\right]`$ is the rotation of the action under the transformation defined in eq. (11<sup>4</sup><sup>4</sup>4 In general, different choices of the flavour content of the operators, sources and rotation are possible or may be necessary. For example for the operators $`O_4`$ and $`O_5`$, the suitable choice is $`O_{4,5}=(\overline{s}\mathrm{\Gamma }u)(\overline{s}\mathrm{\Gamma }d)`$ ($`\mathrm{\Gamma }`$ matrices are not specified) with $`K_1=\overline{d}\gamma _5s`$ and $`K_2=\overline{u}\gamma _5s)`$ and the same rotation as in eq. (11)..
In terms of the fields defined in eq. (12), and of the parity-even and parity-odd operators $`Q_1`$ and $`𝒬_1`$, the Ward identity reads
$`2\widehat{K}_P^0(t_1)\widehat{Q}_1(0)\widehat{K}_P^0(t_2)=2m{\displaystyle \underset{x}{}}\mathrm{\Pi }^0(x)\widehat{K}_P^0(t_1)\widehat{𝒬}_1(0)\widehat{K}_P^0(t_2)`$ (14)
$`\widehat{K}_S^0(t_1)\widehat{𝒬}_1(0)\widehat{K}_P^0(t_2)\widehat{K}_P^0(t_1)\widehat{𝒬}_1(0)\widehat{K}_S^0(t_2)+𝒪(a),`$
where, see eqs. (7)–(9),
$$\widehat{Q}_1=Z_{11}\left(Q_1+\underset{i=2,5}{}\mathrm{\Delta }_{1i}Q_i\right),\widehat{𝒬}_1=𝒵_{11}𝒬_1.$$
(15)
The term on the l.h.s. of eq. (14), corresponding to the rotation of the operator $`𝒬_1`$, is the quantity from which we may extract the physical $`K^0`$$`\overline{K}^0`$ mixing amplitude. The first term on the r.h.s. corresponds to the insertion of the rotation of the action. In terms of Feynman diagrams, it can be seen as the “decay” of a neutral (zero four-momentum) pion, $`\pi ^0`$, into two $`\overline{K}^0`$s under the action of $`𝒬_1`$. In the $`SU(2)`$ isospin symmetric case ($`m_u=m_d`$), only the emission diagrams shown in fig. 2 must be considered. The last two terms in eq. (14) correspond to the rotation of the pseudoscalar kaon sources. These terms are necessary to saturate the Ward identity.
One could envisage the following method to extract the coefficients $`\mathrm{\Delta }_{1i}`$ and the ratio $`\mathrm{\Delta }_{11}=Z_{11}/𝒵_{11}`$, thus determining the subtracted operator $`\stackrel{~}{Q}_1`$. Since eq. (14) must be satisfied for any value of $`t_1`$ and $`t_2`$, the Ward identity corresponds to a system of linear equations in the unknown quantities $`\mathrm{\Delta }_i`$ ($`i=11,12,\mathrm{},15`$. At least in principle, one has an independent equation for any assigned value of $`t_1`$ and $`t_2`$. In practice, the equations become dependent when the Ward identity, for large values of $`t_1`$ and $`t_2`$, is saturated by the contribution of the lowest lying states, namely the two $`K`$-mesons.
This method to extract the $`\mathrm{\Delta }_i`$, which in practice may be very difficult to implement, is however unnecessary and a much easier procedure gives us directly the wanted quantity, namely the correlation function of $`\widehat{Q}_1`$. Let us consider the Ward identity (14) in the limit $`t_1\mathrm{}`$ and $`t_2\mathrm{}`$ (in practice for large values of the time distances). In this limit, we may safely neglect the last two terms of eq. (14), because they correspond to the propagation of scalar states, which are exponentially suppressed with respect to the kaon contribution <sup>5</sup><sup>5</sup>5 This point can be explicitly checked by computing $`\widehat{K}_S^0(t_1)\widehat{𝒬}_1(0)\widehat{K}_P^0(t_2)`$ in the same numerical simulation as the other correlation functions appearing in eq. (14).. Then, up to exponentially suppressed terms, we have ($`\mathrm{\Delta }_{11}=Z_{11}/𝒵_{11}`$ and we have divided all the terms of the Ward identity by a factor of two)
$`\underset{t_1\mathrm{},t_2\mathrm{}}{lim}\mathrm{\Delta }_{11}K_P^0(t_1)\stackrel{~}{Q}_1(0)K_P^0(t_2)`$ (16)
$`{\displaystyle \frac{|0|K_P^0|K^0|^2}{4m_K^2}}e^{m_K(t_1+|t_2|)}\times 𝒵_{11}^1\overline{K}^0|\widehat{Q}_1|K^0`$
$`m{\displaystyle \underset{x}{}}\mathrm{\Pi }^0(x)K_P^0(t_1)𝒬_1(0)K_P^0(t_2)+𝒪(a).`$
Note that the last correlation function, $`\mathrm{\Pi }^0(x)K_P^0(t_1)𝒬_1(0)K_P^0(t_2)`$, is expressed in terms of bare quantities only. Indeed a single constant is sufficient to obtain the physical amplitude, namely $`𝒵_{11}`$ which relate the bare lattice parity-odd operator to the continuum one, renormalized in a specified renormalization scheme. Obviously this constant cannot be determined from the Ward identity, since its values depend on the renormalization condition imposed to the renormalized operator .
In terms of the Feynman diagrams defined in figs. 1 and 2, eq. (16) can be written as
$`2\left(C8(t_1,t_2)+D8(t_1,t_2)\right)=`$
$`2m\left(CE(t_1,t_2)+CE(t_2,t_1)+DE(t_1,t_2)+DE(t_2,t_1)\right)+𝒪(a).`$
In summary, the strategy to obtain the physical matrix element $`\overline{K}^0|\widehat{Q}_1|K^0`$ is extremely simple:
* one computes the correlation function
$`G_3(t_1,t_2)=m{\displaystyle \underset{x}{}}\mathrm{\Pi }^0(x)K_P^0(t_1)𝒬_1(0)K_P^0(t_2)`$
of the bare parity-odd operator at large time distances $`t_{1,2}`$; as “mass” $`m`$, it is more convenient to use the quark mass defined as $`m=Z_A\rho `$, where $`\rho `$ is defined using the axial Ward identity
$`2\rho ={\displaystyle \frac{_\mu A_\mu }{P}}.`$
In the above equation $`A_\mu `$ and $`P`$ are the bare (eventually improved) axial current and pseudoscalar densities and the matrix elements are usually taken between the vacuum and a pion (in our case the $`\pi ^0`$) at rest. $`Z_A`$ is the axial-current renormalization factor.
* then one divides $`G_3(t_1,t_2)`$ by the factor
$`(t_1,t_2)={\displaystyle \frac{Z_5}{4m_K^2}}e^{m_K(t_1+|t_2|)},`$
where $`Z_5=|0|K_P^0|K^0|^2`$, obtaining the quantity $`=𝒵_{11}^1\overline{K}^0|\widehat{Q}_1|K^0`$. $``$ can be readily computed from the study of the kaon two-point correlator $`K_P^0(t)K_P^0(0)`$. Alternatively, one may construct the ratio
$`R(t_1,t_2)=Z_5{\displaystyle \frac{G_3(t_1,t_2)}{K_P^0(t_1)K_P^0(0)K_P^0(t_2)K_P^0(0)}},`$
at large time distances.
* finally, the physical amplitude is given by
$`\overline{K}^0|\widehat{Q}_1|K^0=𝒵_{11}\times .`$
The constant $`𝒵_{11}`$ cannot be determined from the Ward identity and has to be fixed either using perturbation theory, or non-perturbatively on quark states or with the Schrödinger functional method .
The application of the method discussed above to the operators of eq. (3), appearing in extensions of the Standard Model, is so easy that it does not require further discussion. The same approach can be used to compute the $`\mathrm{\Delta }I=3/2`$ $`K`$-$`\pi `$ matrix elements, $`\pi |O_i^{\mathrm{\Delta }S=1}|K`$. A further interesting application is related to hyperon decays. In hyperon decays both the parity-odd and parity-even terms contribute. For the parity-odd case, as explained in 2.1, there are no subtractions induced by the chiral symmetry of the lattice action; for the parity-even case, as done before, we add to the relevant correlation function a soft-pion (zero four momentum) field, i.e. a pseudoscalar density summed over $`x`$, multiplied by the factor $`2m`$ and replace the parity-even operator with the corresponding parity-odd one. Unfortunately, this method works only for $`\mathrm{\Delta }I=3/2`$ transitions since, in the $`\mathrm{\Delta }I=1/2`$ case, the presence of power divergences, due to mixing with lower dimensional operators, makes the strategy proposed in this paper almost impossible to implement in practice.
## 3 Conclusion
With Wilson-like fermions, the calculation of the $`K^0`$$`\overline{K}^0`$ amplitude is complicated by the necessity of accurately determining many mixing coefficients (either perturbatively or non-perturbatively), which arise from the explicit chiral symmetry breaking in the lattice action. In this paper, we have shown that, by using suitable Ward identities, it is possible to compute the physical amplitudes without any subtraction. The error on the final answer is of $`𝒪(a)`$, which corresponds to the best accuracy attainable with a perfect determination of the mixing coefficients in previous approaches. The extension of this method to the improved case, i.e. to obtain matrix elements with an accuracy of $`𝒪(a^2)`$, does not seem possible at present, but it is worth being investigated. Our method can also be applied to $`\mathrm{\Delta }I=3/2`$ hyperon decays. A pioneering numerical calculation of the $`K^0`$$`\overline{K}^0`$ matrix elements using the strategy proposed in this paper is already underway.
## Acknowledgements
G.M. warmly thank R. Frezzotti, S. Sint and A. Vladikas for illuminating discussions at the 2000 Ringberg workshop. Their proposal to compute the mixing amplitude from the parity-odd operator using twisted-mass fermions combined with a chiral rotation, has triggered the present investigation. We thank C.T. Sachrajda, S. Sharpe and M. Testa for stimulating discussions and useful comments. V. G. has been supported by CICYT under the Grant AEN-96-1718, by DGESIC under the Grant PB97-1261 and by the Generalitat Valenciana under the Grant GV98-01-80. We acknowledge the M.U.R.S.T. and the INFN for partial support.
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# 1 Introduction
## 1 Introduction
During the end stages in the evolution of certain supermassive stars general relativity indicates that all the material of the star will collapse into a singularity. This is one of the difficulties with classical general relativity, and it is often suggested that quantum gravity effects will somehow prevent the formation of true singularities. Rhoades and Ruffini <sup>?</sup> have shown that even if the material of the star “stiffens” to the point where the speed of sound in the material becomes equal to the speed of light, the formation of a singularity can not be avoided if the star’s final mass is $`3.2M_{{\scriptscriptstyle }}6.4\times 10^{30}`$ kg. The Hawking-Penrose theorems <sup>?</sup> show that the formation of such singularities is a generic feature of classical general relativity. A rough argument for why a singularity inevitably forms for certain collapsing stars can be given as follows : For a star in which the thermonuclear fire has gone out, the gravitational attraction can be counterbalanced by the quantum mechanical Pauli exclusion pressure. To get an estimate of how the quantum mechanical pressure balances gravity one can use energy considerations <sup>?</sup> with the total energy of the star taken as the sum of the energies of all the particles and the gravitational binding energy. In the relativisitic regime the average energy of each particle is on the order of $`cp_F`$ where $`p_F=(3\pi ^2\mathrm{}^3N/V)^{1/3}`$ is the Fermi momentum, and $`N`$ is the total number of particles contained in the volume $`V`$. The total energy coming form this source is $`E_F=Ncp_F`$ or
$$E_F=N\left(\frac{3\pi ^2c^3\mathrm{}^3N}{V}\right)^{1/3}$$
(1)
The gravitational binding energy is of the order, $`GM^2/R`$, where $`M`$ is the total mass and $`R`$ is the radius of the star. Combining the gravitational binding energy and the energy of Eq. (1) gives an estimate for the total energy of the system (ignoring the rest mass)
$$E_{total}=\left(\frac{9\pi c^3N^4\mathrm{}^3}{4}\right)^{1/3}\frac{1}{R}\frac{GM^2}{R}$$
(2)
In the nonrelativistic case the first term goes as $`1/R^2`$ <sup>?</sup> and there is a radius where stable equilibrium is achieved. For the relativisitic case given in Eq. (2), the quantum mechanical exclusion pressure becomes too “soft”so that no stable equilibrium exists and the star collapses. More rigorous work bears out the conclusion of this rough estimate. In addition to Ref. <sup>?</sup>, Buchdahl <sup>?</sup> <sup>?</sup> has shown that a star of mass $`M`$ with a radius $`R=9M/4`$ or smaller, can not reach static equilibrium. These works indicate that the formation of a singularity can not be prevented by the mechanical forces that the material of the star could exert.
Another example of how general relativity results in particles being inevitably forced to the central singularity of a gravitating point mass can be seen by considering a test particle moving in the Schwarzschild field of some point mass $`M`$. The effective potential per unit mass is <sup>?</sup>
$$V_{eff}=\frac{c^2}{2}\frac{GM}{r}+\frac{L^2}{2r^2}\frac{GML^2}{c^2r^3}$$
(3)
The second term is the standard Newtonian gravitational potential per unit mass, and the third term is the usual centripetal barrier. The last term is a general relativistic addition. It has the effect that if $`r`$ becomes too small there is no stable orbit, and the particle ends up at $`r=0`$. This is to be contrasted with Newtonian gravity where as long as $`L^20`$ the test particle will not be pulled into the singularity
One option for avoiding these singualrities is that gravity must somehow be modified, and it is usually hypothesized that quantum gravity will somehow accomplish this. In particular one would like the strength of the gravitational interaction to decrease at small distance, or large energies. In the following sections we will present various arguments that point to the possibility that the gravitational interaction does decrease with decreasing distance.
## 2 Scaling of $`G`$ by Analogy with Particle Physics
Gauge theories play an important role in modern physics. In the Standard Model <sup>?</sup> of particle physics, matter particles interact via gauge interactions of the group $`SU(3)\times SU(2)\times U(1)`$. General relativity can also be cast in the form of a gauge theory <sup>?</sup>. One key difference between the gauge theories of particle physics and general relativity is that the former have been successfully quantized, but not the latter.
When the gauge theories of particle physics are quantized certain phenomenon occur. In particular, the coupling strength of the gauge theory becomes energy or scale dependent. For the electromagnetic part of the Standard Model the coupling strength increases with increasing energy for the energies so far probed. This is seen experimentally <sup>?</sup> where at low energies $`e^2/4\pi =\alpha _{em}1/137`$ while at energies around 100 GeV $`\alpha _{em}1/128`$. For the SU(3), strong interaction part of the Standard Model one finds that the coupling strength decreases with increasing energy. This decrease of the coupling strength with increasing energy is called asymptotic freedom <sup>?</sup>, and its discovery was one of the first successes of QCD, since it gave an explanation for why, at high energies, the quarks inside the nucleon behaved as if they were essentially free (Bjorken scaling) <sup>?</sup>. Thus, quantized gauge theories have coupling strengths which are scale dependent, and non-Abelian gauge theories can exhibit a coupling strength which becomes weaker at short distances or high energies.
If general relativity is viewed as a gauge theory <sup>?</sup> one can speculate that, as in the case of other quantized gauge theories, the coupling strength of a quantum theory of gravity may become scale dependent. Since general relativity shares similarities with non-Abelian gauge theories, it could be conjectured that general relativity may also be asymptotically free. One of the first theoretical questions which any full theory of quantum gravity, such as string theory or loop quantum gravity <sup>?</sup>, should address is whether the gravitational interaction strength scales with energy, and the nature of the scaling.
## 3 Scaling in effective theories of general relativity
Effective field theory techniques allow one to discuss the quantum corrections to field theories even if the field theories are conventionally non-renormalizable. Recently Donoghue <sup>?</sup> applied effective field theory methods to general relativity to calculate the lowest order quantum corrections to the Newtonian potential. For two point mass, $`M_1`$ and $`M_2`$ separated by a distance $`r`$, quantum corrections modify the Newtonian potential as
$$V(r)=\frac{GM_1M_2}{r}\left[1\frac{127}{30\pi ^2}\frac{G\mathrm{}}{c^3r^2}\right]$$
(4)
In Ref. <sup>?</sup> there is a further correction which goes as $`G(M_1+M_2)/rc^2`$. However this is just a post-Newtonian correction from classical general relativity rather than a quantum correction, since it does not contain $`\mathrm{}`$. When a correct theory of quantum gravity is found it should yield the same kind of quantum corrections in the regime where the effective field theory calculation is valid (i.e. for $`r`$ significantly larger than the Planck length). From Eq. (4) it can be seen that the quantum effects tend to decrease the strength of the gravitational interaction as $`r`$ gets smaller. At ordinary distances this effective decrease of the gravitational coupling is currently unmeasurable since the second term in Eq. (4) is extremely small. It is possible to write Eq. (4) in the usual form, $`V(r)=G(r)M_1M_2/r`$, with a $`r`$ dependent $`G`$
$$G(r)=G_{\mathrm{}}\left(1\frac{127}{30\pi ^2}\frac{G_{\mathrm{}}\mathrm{}}{c^3r^2}\right)$$
(5)
where $`G_{\mathrm{}}6.67\times 10^{11}Nm^2/kg^2`$ is the gravitational coupling constant determined as $`r\mathrm{}`$. For distances not too close to the Planck scale, Eq. (5) implies that Newton’s constant decreases with decreasing distance. Usually the running of the coupling constant in gauge theories is given in terms of a scaling with energy rather than with distance. In the appropriate units one can replace distances $`r`$ for energies $`k`$ via $`r1/k`$, in terms of which the running $`G`$ from Eq. (5) would become $`G(k)=G_0(AG_0^2)k^2`$ <sup>?</sup> where $`A>0`$ is a constant, and $`G_0=G_{\mathrm{}}`$ is the gravitational coupling determined as $`k0`$.
In terms of the effective potential of Eq. (3) one can replace $`G`$ by $`G(r)`$ of Eq. (5) so that the effective potential becomes
$$\stackrel{~}{V}_{eff}=\frac{c^2}{2}\frac{G(r)M}{r}+\frac{L^2}{2r^2}\frac{G(r)ML^2}{c^2r^3}$$
(6)
Now, whereas $`V_{eff}`$ of Eq. (3) had no stable minimum if $`r`$ became too small, $`\stackrel{~}{V}_{eff}`$ of Eq. (6) always has a stable minimum at some small $`r`$. This is most easily seen in the $`L=0`$ case where, by using Eq. (6) to calculate $`d\stackrel{~}{V}_{eff}/dr=0`$, one finds that the effective potential with the distance dependent $`G`$ has a minimum at
$$r_{min}=\sqrt{\frac{127G_{\mathrm{}}\mathrm{}}{10\pi ^2c^3}}1.8\times 10^{35}meters$$
(7)
The numerical value for $`r_{min}`$ shows the weak point of this hypothesized asymptotic freedom of general relativity : this distance is at the Planck distance scale where the effective theory used to calculate $`G(r)`$ of Eq. (5) is suspect. At this scale one really needs a full theory of quantum gravity in order to calculate the scale dependence of $`G`$ with confindence. One can still speculate that this asymptotic freedom, indicated by the effective field theory at low energies, continues to all energy scales for a complete theory of quantum gravity. This is the reverse of speculations in QCD, where the running of $`\alpha _{QCD}`$ in the high energy regime is often said to imply the increase of the coupling strength at low energies, and therefore confinement. Also it can be pointed out that superfically the effective field theory result is not completely unreasonable. Plugging $`r_{min}`$ back into Eq. (5) gives a value for the second term of $`0.3`$ compared to the value of $`1`$ for the first term, so that the first quantum correction is still smaller than the zeroth order classical term.
To make a connection to the Fermi energy argument we need to relate the distance $`r`$ between two point particles with the radius $`R`$ of the star. For a star of radius $`R`$ composed of $`N`$ particles, the average distance between any two of the particles will be roughly $`r=R/(N)^{1/3}`$. With this, Eq. (5) can be written as
$$G(R)=G_{\mathrm{}}\left(1\frac{127}{30\pi ^2}\frac{G_{\mathrm{}}\mathrm{}N^{2/3}}{c^3R^2}\right)$$
(8)
Replacing $`G`$ of Eq. (2) with the scale dependent $`G(R)`$ of Eq. (8) and calculating $`dE_{total}/dR`$ now gives
$$\frac{dE_{total}}{dR}=\frac{1}{R^2}\left(G_{\mathrm{}}M^2\left[\frac{9\pi c^3\mathrm{}^3N^4}{4}\right]^{1/3}\right)\frac{127G_{\mathrm{}}^2\mathrm{}N^{2/3}M^2}{10\pi ^2c^3R^4}$$
(9)
The last term, which arises from the quantum corrections of the effective gravitational field theory, ensures that it is always possible to find some $`R`$ so that $`dE/dR=0`$. This implies that a stable balance between gravity and the quantum mechanical pressure can be achieved due to the weakening of the gravitational interaction. Aside from the heuristic nature of this argument (a more serious calculation would use the Oppenheimer-Volkoff equation with the scale dependent $`G(R)`$ of Eq. (8)) it is found that the radius $`R`$, for which Eq. (9) gives an equilibrium, is again outside the regime where the effective field theory calculation can be trusted. For a star of mass $`M=5M_{{\scriptscriptstyle }}1.0\times 10^{31}kg`$ with $`N=M/m_n=5.97\times 10^{57}`$, it is found that Eq. (9) gives an equilibrium radius of $`R=2.70\times 10^{15}`$ m. This implies an average spacing between the particles of $`r=R/N^{1/3}=1.49\times 10^{34}`$ m, which is only one order of magnitude above the Planck scale of $`10^{35}`$ m. Again, one can hypothesize that the weakening of the gravitational coupling, $`G`$, implied by the low energy effective theory will continue at higher energy scales.
This hypothesized asymptotic freedom for general relativity would not prevent the formation of a black hole, since in the example given above the horizon forms at a distance around 15 km. The scaling of $`G`$ only replaces the singularity at the center of the black hole with an extemely dense, but non-singular mass.
## 4 Scaling in Kaluza-Klein Theories
If the gravitational interaction is eventually unified with the Standard Model interactions, then the scaling of the various coupling strengths may be interrelated. This is similar to grand unified theories such as SU(5) <sup>?</sup>, where the scaling of the various non-gravitational couplings are related. Kaluza-Klein theories offer a simple and direct example of how the scaling of the gravitational and non-gravitational coupling strengths may be related. In the original Kaluza-Klein theory <sup>?</sup> a relationship exists between the electric coupling $`e`$ and Newton’s constant
$$G=\frac{r_5^2e^2}{16\pi }=\frac{r_5^2\alpha _{em}}{4}$$
(10)
where $`r_5`$ was the radius of the curled up fifth dimension. To get non-Abelian gauge fields it is necessary to have more than one compactified dimension. In Ref. <sup>?</sup> a relationship similar to Eq. (10) is given except with the electromagnetic coupling $`e`$ is replaced by the non-Abelian coupling $`g`$, and $`r_5`$ replaced by some rms circumference of the curled up dimensions. The key point about Eq. (10) or its non-Abelian version, is that Newton’s constant is proportional to the square of some non-gravitational coupling constant. Thus $`G`$ should scale with distance or energy in the same manner as $`g^2`$. For non-Abelian theories the coupling strength, $`\alpha =g^2/4\pi `$, usually decreases in strength with decreasing distance scale in a logarithmic way (i.e. $`\alpha (r)=\alpha _o[1+c\alpha _o\mathrm{ln}(r/r_o)]^1`$ where $`c`$ is some positive constant which depends on the non-Abelian gauge group, and $`r_o`$ is a reference distance at which the coupling is measured) so that $`G(r)`$ should also decrease logarithmically. The running of $`G`$ here is different than in the previous section. First, as noted in Ref. <sup>?</sup> the running of $`G`$ implied by the effective field theory treatment goes as a power of energy or inverse power of distance, whereas in the present example, the running is logrithmic. Secondly, in the effective field theory approach the direct quantum corrections of gravity were discussed. Here, the direct quantum effects of four dimensional gravity are ignored, but one still finds that $`G`$ runs if $`g`$ runs. At energies far from the Planck scale the description of the compactified dimensions in terms of a non-Abelian gauge field theory is reasonable, especially if these Kaluza-Klein fields of the compactified dimensions are to describe real non-Abelian fields. If the effective, non-Abelian coupling, $`g`$, exhibits asymptotic freedom (as it should if it is to model the behaviour of non-Abelian fields of the Standard Model) then so will $`G`$. As in the previous section, when Planck scale energies and distances are approached, this treatment of the curled up dimensions by an effective non-Abelian gauge theory breaks down, and a complete, non-perturbative method of quantizing this higher dimensional gravitation theory is required. This running of the gravitational coupling in Kaluza-Klein models again opens up the possibility that the formation of singularities in gravitational collapse may be avoided.
One worry about this Kaluza-Klein argument is that if the running of the non-Abelian coupling, $`g`$, is experimentally observed then the running of $`G`$ should also be seen. For example, let $`g`$ be the QCD coupling. The perturbative running of the QCD coupling is experimentally observed at energy scales greater than about 2 GeV (see Ref. <sup>?</sup> p. 82). If the running of $`G`$ were tied to $`g`$ then one might think that some experimental signature of this running of $`G`$ should have be seen. In accelerator experiments, however, one deals with such small quantities of matter, gravitationally speaking, that any kind of running of $`G`$ would be undetectable. Even inside an apparently high energy environment like the interior of the Sun, where there is a gravitationally significant amount of matter, one has a temperature $`1.6\times 10^7`$ K, which corresponds to an energy scale of about 1.4 keV. This is not in the energy range where the perturbative running of $`g`$ (and therefore $`G`$) would apply.
Situations where a gravitationally significant amount of matter at a high enough energy could exist, occur in situations of gravitational collaspe. For example, taking a stellar mass of $`M=5M_{{\scriptscriptstyle }}=1.0\times 10^{31}`$ kg, so that $`N=M/m_n=5.97\times 10^{57}`$, and taking $`R=1000`$ m gives, an average energy per particle of $`E_F/N6.9`$ GeV (see Eq. (1)), which is an energy range where the perturbative scaling of $`g`$ should apply.
## 5 Conclusions
We have argued that the singularities which occur in general relativity in certain situations, such as gravitational collapse, could possibly be avoided if quantum gravity exhibits a scaling of Newton’s constant. In the case of stellar collapse it has been shown <sup>?</sup> that even if the material of the dead star exerts the maximum possible resistive force, it can not counterbalance the inward push of gravity. Thus the only obvious way to avoid these singularities would be to somehow modify the gravitational interaction, which is essentially the idea behind the common statements that a full quantum theory of gravity would somehow prevent the formation of these singularities. In this article we have presented motivations that a quantum theory of gravity should have a coupling strength which weakens at short distance, or large energy scales, thus allowing an equilibrium to be reached between the quantum mechanical exclusion pressure and the weakened gravitational interaction. First, by analogy with other non-Abelian gauge theories, which when quantized exhibit asymptotic freedom, we argued that a quantum theory of gravity may also exhibit asymptotic freedom. Second, from recent effective field calculations it is found that the gravitational interaction does grow weaker with decreasing distance scale, at least for scales which are not too close to the Planck scale. Finally, if gravity is unified with the other interactions, as in Kaluza-Klein theories, then the running of the different couplings should be related; if the non-Abelian coupling $`g`$ exhibits asymptotic freedom then so should $`G`$. All of these arguments are only good for energies and distances far from the Planck scale. However, the idea that quantum gravity may exhibit asymptotic freedom provides a concrete mechanism of how the singularities of classical general relativity may be avoided. There are currently theories, such as string theory or loop quantum gravity, which hold out the hope of giving a complete quantum theory of gravity. One of the first questions that could be asked of such a complete theory of quantum gravity would be the nature of the non-perturbative scaling, if any, that it gives for $`G`$.
## 6 Acknowledgements
I would like to thank Vladimir Dzhunushaliev for discussions and comments during the writing of this article. This work was supported in part by a Collaboration in Basic Science and Engineering grant from the National Research Council.
References
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# Li2VO(Si,Ge)O4, a prototype of a two-dimensional frustrated quantum Heisenberg antiferromagnet
\[
## Abstract
NMR and magnetization measurements in Li<sub>2</sub>VOSiO<sub>4</sub> and Li<sub>2</sub>VOGeO<sub>4</sub> are reported. The analysis of the susceptibility shows that both compounds are two-dimensional $`S=1/2`$ Heisenberg antiferromagnets on a square lattice with a sizeable frustration induced by the competition between the superexchange couplings $`J_1`$ along the sides of the square and $`J_2`$ along the diagonal. Li<sub>2</sub>VOSiO<sub>4</sub> undergoes a low-temperature phase transition to a collinear order, as theoretically predicted for $`J_2/J_1>0.5`$. Just above the magnetic transition the degeneracy between the two collinear ground states is lifted by the onset of a structural distortion.
\]
In recent years one has witnessed an extensive investigation of quantum phase transition in low-dimensional $`S=1/2`$ Heisenberg antiferromagnets (QHAF) as a function of doping, magnetic field and disorder. For example, two-dimensional QHAF (2DQHAF) have been widely studied to show the occurrence of a phase transition from the renormalized classical to the quantum disordered regime upon charge doping . Another possibility to drive quantum phase transitions in a 2DQHAF is to induce a sizeable frustration. In particular, for a square lattice with an exchange coupling along the diagonal $`J_2`$ about half of the one along the sides of the square $`J_1`$ (see Fig. 1a), a crossover to a quantum disordered phase with a finite gap between the singlet ground state and the first excited state is expected . For $`J_2/J_10.5`$ a Néel order is envisaged, while for $`J_2/J_10.5`$ a collinear order should develop. The collinear order (see Fig. 1a), which can be considered as formed by two interpenetrating Néel sublattices with staggered magnetization $`𝐧_\mathrm{𝟏}`$ and $`𝐧_\mathrm{𝟐}`$, is characterized by an Ising order parameter $`\sigma =𝐧_\mathrm{𝟏}.𝐧_\mathrm{𝟐}=\pm 1`$ . The two values of $`\sigma `$ correspond to the two collinear configurations, one with spins ferromagnetically aligned along the $`x`$ axis, with a magnetic wave-vector $`𝐐=(0,\pi /a)`$, the other with spins ferromagnetically aligned along the $`y`$ axis $`(𝐐=(\pi /a,0))`$. At a certain temperature an Ising phase transition occurs and the system choses among the $`x`$ or $`y`$ collinear configurations. The precise boundaries of the $`J_2/J_1`$ phase diagram for a frustrated 2DQHAF are unknown and could be modified by the presence of a finite third neighbour coupling . These theoretical predictions have not found an experimental support so far, mainly due to the absence of a system which can be regarded as a prototype of a frustrated 2DQHAF.
In this letter we present NMR and magnetization measurements that prove that the isostructural compounds Li<sub>2</sub>VOSiO<sub>4</sub> (LSVO for short) and Li<sub>2</sub>VOGeO<sub>4</sub> (LGVO) , formed by layers of V<sup>4+</sup> ($`S=1/2`$) ions on a square lattice (see Fig. 1b), are prototypes of frustrated 2DQHAF with significant coupling between both first ($`J_1`$) and second ($`J_2`$) neighours. Moreover we show that LSVO undergoes a phase transition to a low temperature collinear order, as expected for $`J_2/J_1>0.5`$. The phase transition is triggered by a lattice distortion which lifts the degeneracy between the two possible collinear ground states and could belong to the Ising universality class.
<sup>29</sup>Si NMR and magnetization measurements have been performed on powder samples while <sup>7</sup>Li NMR measurements, thanks to <sup>7</sup>Li sensitivity, have been carried out also on a $`1\times 1\times 0.2`$ mm<sup>3</sup> LSVO single crystal. NMR spectra and nuclear spin-lattice relaxation rate $`1/T_1`$ have been measured by using standard pulse sequences. The field-cooled magnetization $`M`$ was measured with a commercial Quantum Design MPMS-XL7 SQUID magnetometer.
The structure of $`V^{4+}`$ layers suggests that both the couplings between first and second neighbours can be significant. It is however difficult a priori to decide which one should dominate: first neighbours are connected by two superexchange channels, but they are located in pyramids looking in opposite directions and are not exactly in the same plane, whereas second neighbours are connected by only one channel, but they are located in pyramids looking in the same directions and are in the same plane. It would thus be highly desirable to extract information on the relative value of these exchange integrals from the susceptibility ($`\chi =M/H`$) (see Fig. 2). Although the temperature dependence of the susceptibility of the $`J_1J_2`$ model is not known accurately as a function of $`J_1`$ and $`J_2`$, it turned out to be possible to obtain useful information from the following considerations. If the system was not frustrated, i.e. if $`J_1J_2`$ or $`J_2J_1`$, the susceptibility would be that of a regular Heisenberg AF on the square lattice with coupling $`J`$. In that case, Quantum Monte Carlo simulations can be used to determine the temperature dependence of the susceptibility, and the maximum occurs at $`T_{\mathrm{max}}0.935J`$. Since in that case the Curie-Weiss temperature $`\mathrm{\Theta }=J`$ one has a ratio $`T_{\mathrm{max}}/\mathrm{\Theta }0.935`$. Now, while $`T_{\mathrm{max}}`$ is known very accurately from our measurements, the precise determination of $`\mathrm{\Theta }`$ is more problematic. A simple fit with $`\chi (T)=\chi _{VV}+C/(T+\mathrm{\Theta })`$, with $`\chi _{VV}`$ Van-Vleck susceptibility, is not good enough because there is an important dependence of the results on the lowest temperature used in the fit. To overcome this problem, we have performed a fit of the high temperature part of the susceptibility up to third order. The coefficients are then consistent with a $`J_1J_2`$ model only in a small window for the lowest temperature. Within this window, $`\mathrm{\Theta }`$ depends very weakly on the lowest temperature and a precise estimate of the Curie-Weiss temperature can be achieved. The results are $`\mathrm{\Theta }7.4K`$ for LSVO and $`\mathrm{\Theta }5.2K`$ for LGVO . Accordingly, the ratios $`T_{\mathrm{max}}/\mathrm{\Theta }`$ are equal to 0.72 and 0.67, respectively. In both cases, this ratio is significantly lower than the value 0.935, supporting the presence of a sizeable frustration. Besides, with a smaller ratio, LGVO is expected to be closer to the fully frustrated point $`J_2/J_1=1/2`$ than LSVO, in qualitative agreement with the fact that no phase transition was found in that system down to 1.9 K (see below). What we cannot say however on the basis of this analysis is which of the couplings $`J_1`$ and $`J_2`$ is larger.
To be more quantitative, we need to know the ratio $`T_{\mathrm{max}}/\mathrm{\Theta }`$ for the $`J_1J_2`$ model as a function of $`J_2/J_1`$. It turned out to be impossible to get accurate estimates in the strongly frustrated region $`J_2/J_11/2`$, but exact diagonalizations with 3 sizes available (4, 8 and 16 sites) give a reliable estimate for $`J_2/J_1<.4`$, while Quantum Monte Carlo simulations, which suffer from the minus sign problem, provide useful information down to $`J_2/J_12`$ on the other side, where only two sizes can be used for exact diagonalizations due to the type of order (8 and 16 sites). The experimental ratio $`T_{\mathrm{max}}/\mathrm{\Theta }=.72`$ for LSVO then implies that $`J_2/J_1`$ is approximately equal to .1 or to 3.5, while for LGVO, $`T_{\mathrm{max}}/\mathrm{\Theta }=.67`$ implies that $`J_2/J_1`$ is either close to .25 or to 2.5.
A significant difference between the two compounds is discernable if one reports the derivative $`d\chi /dT`$ vs. T (see the inset to Fig. 2). One observes that around $`T_c2.83`$ K a peak is present for LSVO, while no anomaly in $`d\chi /dT`$ is detected for LGVO, down to 1.9 K. The peak occurs at the same temperature where a peak in <sup>7</sup>Li NMR $`1/T_1`$ is observed (see Fig. 3), signaling a phase transition to a magnetically ordered state. Remarkably, $`T_c`$ was found independent on the magnetic field intensity, within $`\pm 0.15`$ K (i.e. $`\pm 5`$% ), up to $`H=7`$ Tesla.
In LSVO, for $`H=1.8`$ Tesla, one observes that <sup>7</sup>Li $`1/T_1`$ is constant between $`3.5`$ and $`293`$ K. In the high temperature limit ($`T\mathrm{\Theta }`$), by resorting to the usual Gaussian form for the spin correlation function one has
$$(1/T_1)_{\mathrm{}}=\frac{\gamma ^2}{2}\frac{S(S+1)}{3}\frac{\sqrt{2\pi }}{\omega _E}\times \underset{k,i,j}{}(A_{ij}^k)^2$$
(1)
with $`A_{ij}`$ ($`i=x,y,z`$, $`j=x,y`$) the components of the hyperfine tensor due to the $`k^{th}`$ V<sup>4+</sup>, $`\gamma `$ the gyromagnetic ratio and $`\omega _E=Jk_B\sqrt{2zS(S+1)/3}/\mathrm{}`$. $`z=8`$ is the number of nearest neighbour spins of a V<sup>4+</sup> coupled via an effective superexchange coupling $`J`$ related to $`J_1`$ and $`J_2`$. <sup>7</sup>Li hyperfine coupling constants in LSVO have been estimated by reporting the temperature dependence of the paramagnetic shift, both for the single crystal and for the powders, as a function of the susceptibility. It turned out that <sup>7</sup>Li nuclei are coupled to V<sup>4+</sup> ions both via a dipolar and a transferred hyperfine coupling $`A_T=1700`$ Gauss , which is attributed to the two V<sup>4+</sup> nearest neighbours. The dipolar term is close to the one estimated on the basis of lattice sums. From Eq. 1, by using the high temperature value of $`1/T_1`$, one finds $`J=6.4`$ K, a value consistent with the Curie-Weiss temperature derived from the analysis of the susceptibility.
For $`T<\mathrm{\Theta }`$, V<sup>4+</sup> spins are strongly correlated and the behaviour of $`1/T_1`$ depends on which regime LSVO is: quantum critical, renormalized classical or quantum disordered . However, one has to notice that the temperature dependence of $`1/T_1`$ can be strongly influenced by the $`q`$-dependent hyperfine form factor . Hence, in order to make significant statements on the correlated spin dynamics of this frustrated 2DQHAF we have first calculated the hyperfine form factor, which was found only weakly $`q`$dependent in view of the sizeable transferred coupling. Therefore, the temperature dependence of $`1/T_1`$ is fully determined by the one of the correlation length and the fact that $`1/T_1`$ is constant down to $`3`$ K suggests that LSVO, for $`H=1.8`$ Tesla, is in the quantum critical regime .
Below $`T_c`$ one observes an activated decrease of $`1/T_1`$ which is typical of a magnetically ordered system with a gap in the spin wave spectrum . An estimate of the gap amplitude can be done using a simple Arrhenius fit, which yields $`\mathrm{\Delta }=18`$ K (see the inset to Fig. 3). Although the estimate of $`\mathrm{\Delta }`$ might not be very accurate in view of a certain proximity to the phase transition it should be noticed that the estimated value is considerably larger than what one would expect for a magnetic system with a coupling constant of a few degrees kelvin, as deduced from susceptibility and $`1/T_1`$ measurements for $`T\mathrm{\Theta }`$. This fact could suggest a modification of the superexchange coupling at low temperatures, possibly involving a lattice distortion. The occurrence of a lattice distortion is, in fact, corroborated by the modifications in <sup>29</sup>Si NMR powder spectra below $`3.4`$ K (see Fig. 4a). One observes, just above $`T_c`$ the appearence of a shifted narrow peak in <sup>29</sup>Si NMR powder spectrum. On decreasing temperature the low-frequency peak progressively disappears while the intensity of the high frequency one increases. Two aspects should be remarked: $`1)`$ <sup>29</sup>Si NMR line does not broaden below $`T_c`$ indicating that the local field at <sup>29</sup>Si nuclei is zero; $`2)`$ the modification in the shift has to be associated with a modification in the chemical shift or hyperfine coupling, suggesting the occurrence of a structural distortion. The absence of a magnetic field at <sup>29</sup>Si site can be accounted for either by an AF order with the spins parallel to the $`c`$ axis or by a collinear order with V<sup>4+</sup> spins along the sides of the square lattice. In order to exclude the presence of an AF order we studied the angular dependence of the magnetic field at <sup>7</sup>Li nuclei below $`T_c`$ . By considering the same hyperfine coupling tensor determined above $`T_c`$ we found that for the AF order the magnetic field intensity should always increase on turning the magnetic field from parallel to the $`c`$-axis to parallel to the $`ab`$ plane, at variance with the experimental findings. Moreover, the fact that the spins are parallel to the $`ab`$ plane below $`T_c`$ is supported by a recent EPR analysis of the $`g`$ tensor , showing that there is a larger in-plane magnetic anisotropy. Therefore, we conclude that the magnetic order is collinear with the spins in the $`ab`$ plane. This is the first evidence of a collinear order in a frustrated 2DQHAF, whose existence has been theoretically put forward long ago by Chandra and coworkers . It should be observed that only one of the two possble collinear orders, with $`\sigma =\pm 1`$, is compatible with zero magnetic field at <sup>29</sup>Si, the one with the spins parallel to the staggered modulation, i.e. for spins along the $`x`$-axis the one with magnetic vector $`𝐐=(\pi /a,0)`$. The fact that in LSVO always one type of collinear order develops indicates that the four-fold symmetry of the square lattice was broken, possibly by the lattice distortion occurring just above $`T_c`$.
Further information on the collinear phase in LSVO can be derived from the temperature dependence of <sup>7</sup>Li NMR spectra below $`T_c`$ (see Fig. 4c). The temperature dependence of the order parameter, proportional to V<sup>4+</sup> average magnetic moment, is obtained from the splitting of the satellites of <sup>7</sup>Li NMR spectrum. The central peak and the two satellites do not correspond to the $`1/21/2`$ and $`\pm 3/2\pm 1/2`$ transitions, respectively, but correspond to Li sites where the local field is either zero or non-zero (parallel or antiparallel to the external one). We have ruled out the possibility of a quadrupolar splitting by checking that both the length of the RF pulse yielding a $`\pi /2`$ rotation and the recovery curve of nuclear magnetization were the same for all lines. Moreover the observed splitting is nearly an order of magnitude larger than the quadrupolar one calculated on the basis of a point charge approximation. One notices (see Fig. 4b) a rather sharp, but continuous, decrease of the order parameter close to $`T_c`$. An accurate determination of the critical exponent $`\beta `$ would require a temperature stability better than $`5\times 10^3`$ K which could not be achieved with our cryogenic apparatus. Still an upper limit for $`\beta `$ can be estimated from our measurements which are carried out in steps of $`10^2`$ K for $`TT_c`$. We find that $`\beta 0.25`$, a value compatible with a 2D Ising phase transition, where $`\beta =1/8`$. It should be observed that the relative amplitude of the central and satellite lines varies with decreasing temperature. This could be due to a modification of the interplanar correlation, to which <sup>7</sup>Li spectrum is sensitive. To elucidate this aspect further investigation of the collinear order is demanded.
In conclusion, we have presented for the first time susceptibility measurements showing that Li<sub>2</sub>VOSiO<sub>4</sub> and Li<sub>2</sub>VOGeO<sub>4</sub> can be considered as prototypes of frustrated 2DQHAF and NMR spectra demonstrating that in the former a collinear phase is established at low-temperature, as predicted for $`J_2/J_1>1/2`$ . Finally <sup>29</sup>Si NMR spectra suggest the occurrence of a structural distortion, just above the magnetic transition, which lifts the degeneracy between the two collinear ground states.
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# Untitled Document
hep-th/0005007 TIFR-TH/00-21
A Note on Supergravity Duals of Noncommutative Yang-Mills Theory
Sumit R. Das <sup>1</sup> das@theory.tifr.res.in, and Bahniman Ghosh <sup>2</sup> bghosh@theory.tifr.res.in
Tata Institute of Fundamental Research
Homi Bhabha Road, Mumbai 400 005, INDIA
A class of supergravity backgrounds have been proposed as dual descriptions of strong coupling large-N noncommutative Yang-Mills (NCYM) theories in $`3+1`$ dimensions. However calculations of correlation functions in supergravity from an evaluation of relevant classical actions appear ambiguous. We propose a resolution of this ambiguity. Assuming that some holographic description exists - regardless of whether it is the NCYM theory - we argue that there should be operators in the holographic boundary theory which create normalized states of definite energy and momenta. An operator version of the dual correspondence then provides a calculation of correlators of these operators in terms of bulk Green’s functions. We show that in the low energy limit the correlators reproduce expected answers of the ordinary Yang-Mills theory.
April, 2000
1. Introduction
Noncommutative gauge theories appear as limits of D-brane open string theories in the presence of nonvanishing NSNS B fields ,. In a precise decoupling limit was defined in which the string tension is scaled to infinity and the closed string metric is scaled to zero keeping the dimensionful NSNS B field fixed, and it was shown that in this limit the open string theory on D-branes reduces to precisely noncommutative Yang-Mills (NCYM) theory. Furthermore, NCYM appears naturally in the IKKT matrix theory , a connection which has led to further insight. Noncommutative field theories have novel perturbative behaviour including an IR/UV mixing .
It is natural to ask whether the large-N limit of NCYM is dual to closed string theory in some background, analogous to the duality of usual Yang-Mills theory to string theories ,. In fact, in and supergravity duals of strong ’t Hooft coupling limits of NCYM were proposed. Further evidence for such duality came from the study of D-instantons in such supergravity backgrounds and their relationship to instantons of the NCYM . Aspects of such duality have been explored in
However, the nature of this duality has been rather confusing. In usual AdS/CFT duality, there is a well defined connection between boundary correlators in supergravity and correlators of local operators of the boundary conformal field theory. Consider $`AdS_{d+1}`$ with a metric
$$ds^2=\frac{1}{z^2}(dt^2+dz^2+d\stackrel{}{x}d\stackrel{}{x})$$
where $`\stackrel{}{x}=x^i,i=1\mathrm{}(d1)`$. The boundary of this space is at $`z=z_0`$ with $`z_0<<1`$. In the strong ’t Hooft coupling limit the correspondence may be written as
$$e^{\mathrm{\Gamma }_{eff}[\varphi _0^i(x)]}=𝒟A_\mu e^{S_{CFT;z_0}_i{\scriptscriptstyle 𝑑x\varphi _0^i(x)𝒪_i(x)}}$$
where $`S_{CFT;z_0}`$ is the action of a $`d`$ dimensional theory with a cutoff $`z_0`$. $`\mathrm{\Gamma }_{eff}[\varphi _0^i(x)]`$ is the effective action of supergravity in the background (1.1) as a functional of the values of the various fields $`\varphi ^i(z,x)`$ on the boundary,
$$\varphi ^i(z_0,x)=z_0^{\mathrm{\Delta }_id}\varphi _0^i(x)$$
$`𝒪_i(x)`$ is the operator in the boundary theory which is dual to the field $`\varphi ^i`$. This particular behaviour of the fields near the boundary given in (1.1) is guaranteed by the conformal isometries of the $`AdS`$ space-time and $`\mathrm{\Delta }_i`$ turns out to be the conformal dimension of the operator $`𝒪_i(x)`$. Thus in this case there is a holographic relationship between the bulk theory and a local field theory on the boundary.
In a similar approach was tried for the spacetimes proposed to be duals of NCYM. It was found that to make sense of the results, the relationship between the boundary values of the supergravity fields and the sources of some operators in the proposed dual boundary theory is necessarily nonlocal in position space. In other words the relationship between momentum space quantities involve nontrivial functions of momentum. Similar momentum dependent factors were encountered in the study of one point functions in instanton backgrounds in . It is not a priori clear what these factors should be and the entire procedure appears ambiguous since we can get any answer we want by considering different renormalization factors.
In this paper we assume that a holographic description of the bulk theory in such backgrounds exists in terms of a theory living on the boundary - which may or may not be the NCYM theory. This boundary theory will be generically nonlocal and there would not be local operators creating physical states. However, because of translation invariance, there are well defined operators which create normalized states of definite energy and momenta. A dual correspondence in this context means that such states are to be identified with states of supergravity modes propagating in this background. Consequently, bulk correlation functions in supergravity, with points taken to lie on the boundary, define a set of momentum space boundary correlators unambiguously. This gives a supergravity prediction for such correlators which should reduce to usual Yang-Mills correlators in the low energy limit. We show that this is indeed true.
One of the motivations to study the dual correspondence in the decoupling limit of is that some of the dual space-times are asymptotically flat in terms of the Einstein metric. The hope is that some understanding of the holographic correspondence may teach us something about holographic relationships in other asymptotically flat spaces. In fact there is a close relationship between the type of problems encountered here and those encountered in the holographic correspondence for NS five branes . Here again the dual spacetimes are asymptotically flat and the boundary theory is nonlocal. Our prescription of reading off correlators from bulk Green’s function is in fact closely related to the procedure adopted in where the correlators are identified with relevant S-matrix elements.
In section 2 we present the supergravity solutions we are dealing with and the equations for decoupled modes propagating in this background. In section 3 we calculate the classical action for this on-shell mode as in and state the ambiguity one encounters in trying to extract boundary correlators from this action <sup>3</sup> We correct some errors in the treatment of .. In section 4 we give the operator version of the standard AdS/CFT correspondence for $`B=0`$ and show how boundary correlators of operators creating normalized states in momentum space can be read off from bulk Green’s function. In section 5 we repeat this for our backgrounds with $`B0`$ and show how boundary correlators of momentum space operators can be obtained unambiguously. Section 6 contains concluding remarks.
2. The supergravity solution and modes
We will consider two kinds of IIB supergravity backgrounds with nonzero $`B`$ fields obtained in . The D3 brane has a worldvolume along $`(x^0\mathrm{}x^3)`$.
The first kind of background has a nonzero NSNS $`B_{23}`$ with all other components set to zero. In the decoupling limit which corresponds to the low energy limit of the string metric is
$$ds^2=\alpha ^{}R^2[u^2(dx_0^2+dx_1^2)+\frac{u^2}{1+a^4u^4}(dx_2^2+dx_3^2)+\frac{du^2}{u^2}+d\mathrm{\Omega }_5^2]$$
where the dilaton $`\varphi `$, NS field $`B`$ and the RR 2-form field $`\stackrel{~}{B}`$ and the five form field strength $`F`$ are given by
$$\begin{array}{cc}& e^{2\varphi }=\frac{g^2}{1+a^4u^4}B_{23}=\frac{\alpha ^{}R^2}{a^2}\frac{a^4u^4}{1+a^4u^4}\hfill \\ & \stackrel{~}{B}_{01}=\frac{\alpha ^{}a^2R^2}{g}u^4F_{0123u}=\frac{\alpha ^2}{g(1+a^4u^4)}_u(u^4R^4)\hfill \end{array}$$
where $`R^4=4\pi gN`$ and $`g`$ is the open string coupling. In the infrared, $`u0`$ the space time is $`AdS_5\times S^5`$.
In this background, the graviton fluctuation $`h_{01}`$ with zero momenta along $`x^0,x^1`$ and zero angular momenta along the $`S^5`$ satisfy a simple decoupled equation. With $`\varphi =g^{00}h_{01}`$ this is
$$_\mu (\sqrt{g}e^{2\varphi }g^{\mu \nu }_\nu \varphi )=0$$
which becomes in terms of modes
$$\varphi (u,x_2,x_3)=[\frac{d^2k}{(2\pi )^4}]\varphi (\stackrel{}{k},u)e^{ik_2x^2ik_3x^3}$$
$$_u(u^5_u\varphi (\stackrel{}{k},u))k^2u(1+a^4u^4)\varphi (\stackrel{}{k},u)=0$$
where in (2.1) $`k^2=k_2^2+k_3^2`$. Such zero energy perturbations do not make sense in Lorentzian signature. We will therefore work in the euclidean signature.
The second kind of background has self dual $`B`$ fields and has in addition a nontrivial axion $`\chi `$. The decoupling limit solution in euclidean signature is given by the Einstein metric $`ds_E^2`$
$$ds_E^2=\frac{\alpha ^{}R^2}{\sqrt{\widehat{g}}}[(f(u))^{1/2}(d\stackrel{~}{x}_0^2+\mathrm{}+d\stackrel{~}{x}_3^2)+(f(u))^{1/2}(du^2+u^2d\mathrm{\Omega }_5^2)]$$
while the other fields are
$$\begin{array}{cc}& e^{\varphi _0}=\frac{1}{\widehat{g}}i\chi _0=\frac{1}{\widehat{g}}u^4f(u)\hfill \\ & \stackrel{~}{B}_{01}=\stackrel{~}{B}_{23}=\frac{i}{\widehat{g}}B_{01}=\frac{i}{\widehat{g}}B_{23}=\frac{i\alpha ^{}a^2R^2}{\widehat{g}}(f(u))^1\hfill \\ & F_{0123u}=\frac{4i(\alpha ^{})^2R^4}{\widehat{g}u^5}(f(u))^2\hfill \end{array}$$
where
$$f(u)=\frac{1}{u^4}+a^4,R^4=4\pi \widehat{g}N$$
An interesting feature of this solution is that the space-time is asymptotically flat (in Einstein metric) even in the decoupling limit. In fact (2.1) is exactly the full D3 brane metric. Furthermore near the boundary $`u=\mathrm{}`$ the string coupling vanishes.
In the background (2.1) - (2.1), it was shown in that there are special fluctuations which satisfy decoupled equations. These are fluctuations of dilaton $`\delta \varphi `$ and axion $`\delta \chi `$ which obey the condition
$$\delta \chi +ie^\varphi \delta \varphi =0.$$
The corresponding equation for the fluctuation $`\delta \varphi `$ is
$$_E^2(e^{\varphi _0}\delta \varphi )=0$$
where $`_E^2`$ is the laplacian in the Einstein metric $`ds_E^2`$ given in (2.1). These fluctuations are now allowed to carry momenta along all the brane worldvolume directions, but for simplicity we will consider zero $`S^5`$ angular momenta. The modes are
$$e^{\varphi _0}\delta \varphi (u,x)=[\frac{d^4k}{(2\pi )^4}]\mathrm{\Phi }(\stackrel{}{k},u)e^{i\stackrel{}{k}\stackrel{}{x}}$$
where
$$\stackrel{}{k}\stackrel{}{x}=\underset{i=0}{\overset{3}{}}k_ix^i$$
and the equation (2.1) once again becomes
$$_u(u^5_u\mathrm{\Phi }(\stackrel{}{k},u))k^2u(1+a^4u^4)\mathrm{\Phi }(\stackrel{}{k},u)=0$$
which is identical to (2.1). In (2.1),
$$k^2=\underset{i=0}{\overset{3}{}}(k_i)^2$$
3. Solutions and the boundary action
The solutions of the equation (2.1) or (2.1) may be written in terms of Mathieu functions . Introducing the coordinates
$$u=\frac{1}{a}e^w$$
the two independent solutions may be chosen to be
$$\frac{1}{u^2}H^{(1)}(\nu ,w+\frac{i\pi }{2}),\frac{1}{u^2}H^{(2)}(\nu ,w\frac{i\pi }{2})$$
Here the parameter $`\nu `$ is determined in terms of the combination $`(ka)`$. It has a power series expansion given by
$$\nu =2\frac{i\sqrt{5}}{3}(\frac{ka}{2})^4+\frac{7i}{108\sqrt{5}}(\frac{ka}{2})^8+\mathrm{}$$
The Mathieu functions $`H^{(i)}`$ have the asymptotic property
$$H^{(i)}(\nu ,z)H_\nu ^{(i)}(kae^z)z\mathrm{}$$
where $`H_\nu ^{(i)}(z)`$ denotes Hankel functions. Furthermore, in this region, only $`(ka)<<1`$ contribute significantly, so that $`\nu 2`$. Thus near $`u=0`$ where the geometry is $`AdS_5\times S^5`$ the solutions become the standard solutions of the massless wave equation in $`AdS`$, viz $`(1/u^2)K_2(k/u)`$ and $`(1/u^2)I_2(k/u)`$
¿From the asymptotic behavior of the Hankel functions it is clear that the solution $`H^{(1)}(\nu ,w+\frac{i\pi }{2})`$ is well behaved and goes to zero in the interior of the spacetime at $`u=0`$, while the solution $`H^{(2)}(\nu ,w\frac{i\pi }{2})`$ is well behaved at $`u=\mathrm{}`$.
$$\begin{array}{cc}& H^{(1)}(\nu ,w+\frac{i\pi }{2})e^{i\frac{\pi }{2}(\nu +1)}\sqrt{\frac{2}{\pi kae^w}}e^{kae^w}w\mathrm{}(u0)\hfill \\ & H^{(2)}(\nu ,w\frac{i\pi }{2})e^{i\frac{\pi }{2}(\nu +1)}\sqrt{\frac{2}{\pi kae^w}}e^{kae^w}w\mathrm{}(u\mathrm{})\hfill \end{array}$$
3.1. Supergravity actions and correlators in the usual AdS/CFT correspondence
Let us recall the standard way of obtaining correlators in the AdS/CFT correspondence using (1.1). For a minimally coupled scalar field of mass $`m`$ in $`AdS_{d+1}`$ with a metric given by (1.1) one considers a solution (in euclidean signature) which is smooth in the interior
$$\varphi (x,z)=[\frac{d^dk}{(2\pi )^d}]k^\nu z^{d/2}K_\nu (kz)e^{ikx}\varphi _0(k)$$
Here $`x`$ denotes all the four directions $`x=(\stackrel{}{x},t)`$ and $`k^2=k_0^2+\stackrel{}{k}^2`$. One then computes the supergravity action after putting in a boundary at $`z=z_0<<1`$. Note that the fourier modes $`\varphi _0(k)`$ have been defined so that as one approaches the boundary
$$\mathrm{Lim}_{z0}\varphi (z,k)=z^{\frac{d}{2}\nu }\varphi _0(k)$$
$`\varphi _0(k)`$ are then taken to be sources for operators $`O(k)`$ conjugate to the supergravity field in the Yang-Mills theory living on the boundary. The action is purely a boundary term with the leading nonanalytic piece
$$S[\frac{d^dk}{(2\pi )^d}]\varphi _0(k)\varphi _0(k)k^{2\nu }\mathrm{log}(kz_0)$$
The correspondence (1.1) then leads to a two point function of $`𝒪(k)`$ with a leading nonanalytic piece
$$<𝒪(k)𝒪(k)>k^{2\nu }\mathrm{log}(kz_0)$$
which shows that the dimension of this operator is
$$\mathrm{\Delta }=d/2+\nu $$
This specific relation between the sources and the boundary values of the field is simple - the power of the infrared cutoff which appears in (3.1) is $`(d\mathrm{\Delta })`$. By the IR/UV correspondence this is precisely the power of ultraviolet cutoff necessary to add a perturbation $`𝒪(k)`$ to the boundary action.
3.2. Supergravity actions for $`B0`$
The proposal of is to consider $`u=\mathrm{}`$ as the boundary also for the classical solutions in the presence of a $`B`$ field. The solution to be used in calculating the classical action has to be regular in the interior which means we have to take
$$\varphi _k(u)=\frac{1}{u^2}H^{(1)}(\nu ,w+\frac{i\pi }{2})e^{i\frac{\pi }{2}(\nu +1)}\varphi _0(k)$$
Once again the classical action is a boundary term. To evaluate this term we need to find the behavior of solution (3.1) for large values of $`u`$. This can be done using the relation
$$H^{(1)}(\nu ,w)=\frac{1}{C(ka)}[(\chi (ka)\frac{1}{\chi (ka)})H^{(1)}(\nu ,w)+(\chi (ka)\frac{1}{\eta (ka)^2\chi (ka)})H^{(2)}(\nu ,w)]$$
where we have defined
$$\eta (ka)=e^{i\pi \nu }C(ka)=\eta (ka)\frac{1}{\eta (ka)}$$
and the function $`\chi (ka)`$ has been defined in in terms of relations between various Mathieu functions. Defining further the functions
$$A(ka)=\chi (ka)\frac{1}{\chi (ka)}B(ka)=\eta (ka)\chi (ka)\frac{1}{\eta (ka)\chi (ka)}$$
the asymptotic form of the solution $`\varphi _k(u)`$ becomes
$$\varphi _k(u)\frac{1}{u^2C(ka)}\sqrt{\frac{2}{\pi ka^2u}}[iA(ka)e^{ka^2u}\widehat{B}(ka)e^{ka^2u}]\varphi _0(k)(u\mathrm{})$$
where $`\widehat{B}(ka)`$ denotes the real part of $`B(ka)`$ <sup>4</sup> The real part has to be taken in this asymptotic expansion since the wave functions are real. This is similar to what happens in the asymptotic expansions of modified Bessel’s functions.. We will see shortly that $`A(ka)`$ is purely imaginary so that the expression in (3.1) is real. As expected the solution diverges exponentially at the boundary $`u=\mathrm{\Lambda }`$ with $`\mathrm{\Lambda }>>1`$. The contribution to the classical action from this solution, which becomes the boundary term
$$[u^5\varphi _k(u)_u\varphi _k(u)]_{u=\mathrm{\Lambda }}$$
is clearly divergent. Subtracting this infinite piece we get a term where the exponentials cancel leaving a contribution
$$S_B=\frac{5}{2\mathrm{\Lambda }}\frac{2}{\pi ka^2}\frac{iA(ka)\widehat{B}(ka)}{C^2(ka)}\varphi _0(k)\varphi _0(k)$$
Note that the boundary action given in has an error and differs from the above by essentially a factor of $`1/\mathrm{\Lambda }`$.
Mimicking the standard procedure in the AdS/CFT correspondence we might want to relate the functions $`\varphi _0(k)`$ to source terms in the dual theory living on the boundary, so that derivatives of the classical action with respect to these would give correlation functions of the dual operators. So far, however, we have no clue about this precise relationship. The only guide we have at this stage is that in the low energy limit $`ka<<1`$ the correlators should reproduce the known answers in the $`AdS_5\times S^5`$ case - for these minimally coupled scalars the two point function should go as $`k^4\mathrm{log}(k)`$.
3.3. Low energy limits
We need to find the low $`ka`$ expansion of the various functions. This may be done using the results of as follows. First note that the function $`\eta (ka)`$ is purely real, as follows from the expression for $`\nu `$ in (3.1). The functions $`A(ka),B(ka)`$ and $`C(ka)`$ satisfy a unitarity relation
$$|B(ka)|^2=|A(ka)|^2+|C(ka)|^2$$
Using the reality of $`\eta (ka)`$ and hence $`C`$ it may be easily shown from (3.1) that $`\chi (ka)`$ must be a pure phase, so that $`A(ka)`$ is purely imaginary. Denoting
$$\eta (ka)=e^{\beta (ka)}\chi (ka)=e^{i\gamma (ka)}$$
and using (3.1) the various functions may be expressed in terms of the function
$$P(ka)=\frac{|C(ka)|^2}{|B(ka)|^2}$$
which is the absorption probability in the full D3 brane background computed in . We give below some expressions which we will need
$$\begin{array}{cc}& \frac{\widehat{B}(ka)}{iA(ka)}=[\frac{P(ka)\mathrm{cosh}^2\beta (ka)\mathrm{sinh}^2\beta (ka)}{1P(ka)}]^{1/2}\hfill \\ & \frac{iA(ka)}{C(ka)}=[\frac{1}{P(ka)}1]^{1/2}\hfill \end{array}$$
The expansion of $`P(ka)`$ given in is
$$P(ka)=A_0(ka)^8[1+A_1(ka)^4+A_2(ka)^4\mathrm{log}(ka)+\mathrm{}]$$
where $`A_i`$ are numerical constants. The expansion of $`\beta (ka)`$ is of the form
$$\beta (ka)=\beta _0(ka)^4[1+\beta _1(ka)^4+\beta _2(ka)^8+\mathrm{}]$$
where $`\beta _i`$ are numerical coefficients. This leads to the following expansions
$$\begin{array}{cc}& \frac{\widehat{B}(ka)}{iA(ka)}(ka)^4[1+\alpha _1(ka)^4+\alpha _2(ka)^4\mathrm{log}(ka)+\mathrm{}]\hfill \\ & \frac{iA}{C(ka)}\frac{1}{(ka)^4}[1+\gamma _1(ka)^4+\gamma _2(ka)^4\mathrm{log}(ka)+\mathrm{}]\hfill \end{array}$$
where the coefficients $`\alpha _i`$ and $`\gamma _i`$ may be obtained from the expansions in (3.1) and (3.1). Using these, the action $`S_B`$ is seen to behave as
$$S_B\frac{1}{a(ka)^5}[1+(ka)^4+(ka)^4\mathrm{log}(ka)]\frac{1}{\mathrm{\Lambda }}\varphi _0(k)\varphi _0(k)$$
for small momenta. Therefore if we define a renormalized boundary value of the field
$$\mathrm{\Phi }_0(k)=F(ka)\varphi _0(k)$$
such that
$$F(ka)\frac{1}{(ka)^{5/2}}\frac{1}{\mathrm{\Lambda }^{1/2}}$$
for small momenta, and declare that $`\mathrm{\Phi }_0(k)`$ are the sources which couple to the boundary theory operator dual to the bulk field $`\varphi `$, we would certainly get the correct low momentum behavior for the nonanalytic piece of the two point function <sup>5</sup> In it is claimed that the two point function of the operators defined there (which differs from ours) has the correct low energy behavior. We havent been able to see how this follows.
Clearly, this is an arbitrary procedure. Once we use momentum dependent factors to renormalize fields, we can get any answer we want ! At this stage there is no obvious principle which determines this factor. Any statement about holography would be an empty statement.
4. Bulk and Boundary Green’s functions in $`AdS/CFT`$
The usual $`AdS_{d+1}/CFT_d`$ correspondence may be understood in terms of the modes of bulk field operators. Our treatment follows with some differences. Consider quantization of a massive scalar field $`\varphi (z,\stackrel{}{x},t)`$ of mass $`m`$ which is minimally coupled to the AdS metric (1.1). The field has the following mode expansion
$$\varphi (z,\stackrel{}{x},t)=\frac{z^{d/2}}{2R^{(d1)/2}}_0^{\mathrm{}}d\alpha \frac{d^{d1}k}{(2\pi )^d}(\frac{\alpha }{\omega })^{1/2}J_\nu (\alpha z)[a(\stackrel{}{k},\alpha )e^{i(\omega t\stackrel{}{k}\stackrel{}{x})}+(h.c.)]$$
where
$$\omega ^2=\stackrel{}{k}^2+\alpha ^2\nu =\frac{1}{2}(d^2+4m^2)^{\frac{1}{2}}$$
The modes are normalized so that the annihilation/creation operators satisfy the standard commutators
$$[a(\stackrel{}{k},\alpha ),a^{}(\stackrel{}{k}^{},\alpha ^{})]=\delta ^{(d1)}(\stackrel{}{k}\stackrel{}{k}^{})\delta (\alpha \alpha ^{})$$
We can now make a change of integration variables from $`(\alpha ,\stackrel{}{k})`$ to $`(\omega ,\stackrel{}{k})`$ and define new operators <sup>6</sup> There is a subtely here. In Lorentzian signature we must have $`\omega ^2>k^2`$ so that the range of integration over the four momenta is strictly restricted. However, this fact has no consequence for two point functions. We thank E. Martinec for discussions about this point.
$$b(\stackrel{}{k},\omega )=(\frac{\omega }{\alpha })^{1/2}a(\stackrel{}{k},\alpha )$$
which satisfy the commutation relations
$$[b(\stackrel{}{k},\omega ),b^{}(\stackrel{}{k}^{},\omega ^{})]=\delta ^{(d1)}(\stackrel{}{k}\stackrel{}{k}^{})\delta (\omega \omega ^{})$$
and rewrite the expansion (4.1) as
$$\varphi (z,\stackrel{}{x},t)=\frac{z^{d/2}}{2R^{(d1)/2}}d\omega \frac{d^{d1}k}{(2\pi )^d}J_\nu (\alpha z)[b(\stackrel{}{k},\omega )e^{i(\omega t\stackrel{}{k}\stackrel{}{x})}+(h.c.)]$$
In (4.1) $`\alpha `$ is determined in terms of $`\omega `$ and $`\stackrel{}{k}`$ by the relation (4.1). The states created by $`b^{}(\omega ,\stackrel{}{k})`$, denoted as $`|\omega ,\stackrel{}{k}>=b^{}(\omega ,\stackrel{}{k})|0>`$ are normalized according to $`d`$-dimensional delta functions, as follows from (4.1),
$$<\omega ,\stackrel{}{k}|\omega ^{},\stackrel{}{k}^{}>=\delta ^{(d1)}(\stackrel{}{k}\stackrel{}{k}^{})\delta (\omega \omega ^{})$$
The holographic correspondence then implies that these states are also states in the $`d`$ dimensional boundary theory and there are composite operators which create these states. We can define these operators in momentum space as
$$𝒪(\omega ,\stackrel{}{k})=\frac{2\pi }{R^{(d1)/2}}(\omega ^2\stackrel{}{k}^2)^{\nu /2}[\theta (\omega )b(\omega ,k)+\theta (\omega )b^{}(\omega ,\stackrel{}{k})]$$
The overall power of $`\alpha =(\omega ^2\stackrel{}{k}^2)^{1/2}`$ follows from the fact that as we approach the boundary $`z0`$ the radial wave function $`J_\nu (\alpha z)(\alpha z)^\nu `$. This allows us to define in an unambiguous way a momentum space correlation function of the boundary theory in terms of the boundary values of the fourier transform of the Feynman Green’s function of the bulk field
$$<𝒪(\omega ,\stackrel{}{k})𝒪(\omega ,\stackrel{}{k})>=\mathrm{Lim}_{z,z^{}0}(\frac{\alpha ^{2\nu }}{\psi _{\omega ,\stackrel{}{k}}(z)\psi _{\omega ,\stackrel{}{k}}(z^{})})G_F(z,z^{};\omega ,k)$$
where we have used translation invariance along the $`\stackrel{}{x}`$ directions and $`\psi _{\omega ,\stackrel{}{k}}(z)`$ denote the bulk radial wavefunctions
$$\psi _{\omega ,\stackrel{}{k}}(z)=z^{d/2}J_\nu (\alpha z)$$
The corresponding wavefunctions $`\psi _{\omega ,\stackrel{}{k}}(t,\stackrel{}{x},z)=z^{d/2}J_\nu (\alpha z)e^{i(\omega t\stackrel{}{k}\stackrel{}{x})}`$ are normalized in terms of the standard Klein-Gordon norm
$$(\psi _{\omega ,\stackrel{}{k}}(\stackrel{}{x},t,z),\psi _{\omega ^{},\stackrel{}{k}^{}}(\stackrel{}{x},t,z))=\delta (\omega \omega ^{})\delta (\stackrel{}{k}\stackrel{}{k}^{})$$
We can now consider a Wick rotation to obtain a relation between Euclidean Green’s functions
$$<𝒪(k)𝒪(k)>_E=\mathrm{Lim}_{z,z^{}0}(\frac{k^{2\nu }}{\psi _k(z)\psi _k(z^{})})G_E(z,z^{};\omega ,k)$$
where the wave functions are also rotated to euclidean space and $`k`$ without a vector sign denotes the $`d`$ dimensional euclidean momenta.
It may be easily checked that the euclidean bulk Green’s function leads to the euclidean correlators on the boundary obtained according to the procedure of . The euclidean bulk Green’s function is
$$G_E(z,k;z^{},k)=(zz^{})^{d/2}K_\nu (kz)I_\nu (kz^{})z^{}<z$$
Using the asymptotic form of the modified Bessel functions for $`z,z^{}0`$, we easily see that the leading nonanalytic piece of the Green’s function is
$$G_E(z,z^{},k)(zz^{})^{\nu +d/2}k^{2\nu }\mathrm{log}(kz)$$
Using (4.1) we get the boundary correlator
$$<𝒪(k)𝒪(k)>k^{2\nu }\mathrm{log}(kz)$$
in agreement with (3.1).
We can in fact define local operators on the boundary by taking the fourier transform of the momentum space field $`𝒪(\omega ,\stackrel{}{k})`$. In fact the power of $`\alpha `$ in (4.1) has been chosen such that this boundary field is in fact the boundary value of the bulk field upto the value of the bulk wavefunction at the boundary
$$\mathrm{Lim}_{z0}\varphi (z,\stackrel{}{x},t)=(z)^{\nu +d/2}𝒪(\stackrel{}{x},t)$$
thus defining $`𝒪(\stackrel{}{x},t)`$. The reason why this can be done in an unambiguous fashion is that the wavefunctions decay as the same power of $`z`$ regardless of the value of the momenta. This in turn is a consequence of the conformal isometries of $`AdS`$ space. For IIB supergravity on $`AdS_5\times S^5`$ it is known that the boundary operators thus defined are in fact local gauge invariant operators of $`N=4`$ SYM theory.
5. Bulk and Boundary Green’s functions for $`B0`$
For our supergravity backgrounds with $`B0`$ the conjectured dual theory - NCYM - is not a local quantum field theory in the conventional sense. The backgrounds we are dealing with are not asymptotically AdS - in fact the second class of euclidean backgrounds are asymptotically flat. However we do have translation invariance along the brane directions - so that physical states of the NCYM can be still labelled by the energy and momenta. If this duality is indeed correct, we should be able to represent such states by on shell states of supergravity with the same values of energy and momentum - essentially repeating the treatment of section 2. In this section we argue that such a correspondence between operators in momentum space leads to an unambiguous supergravity prediction of two point functions for the corresponding operators in the NCYM theory. We will then verify that this has the correct low energy behavior.
For the kind of fields which satisfy the minimally coupled massless Klein Gordon equation as we have been studying the mode expansion in our background may be written down in analogy to (2.6) with the Bessel function being replaced by the appropriate Mathieu function. With Lorentzian signature there are two independent wavefunctions which are (delta function) normalizable, corresponding to incoming and outgoing waves in the full D3 brane geometry. In the following we will consider only the wave function whose euclidean continuation does not blow up at infinity - this is the wave function given by
$$\mathrm{\Psi }_k(u)=N(ka)e^{i\frac{\pi }{2}(\nu +1)}\frac{1}{u^2}H^{(2)}(\nu ,w)$$
where $`N(ka)`$ is a normalization factor. The mode expansion may be then written as
$$\varphi (\stackrel{}{x},u,t)=\frac{d^3k}{(2\pi )^3}\frac{d\omega }{2\pi }N(ka)\frac{1}{u^2}H^{(2)}(\nu ,w)e^{\frac{i\pi \nu }{2}}[e^{i(\omega t\stackrel{}{k}\stackrel{}{x})}\widehat{a}(\omega ,k)+c.c.]$$
$`N(ka)`$ is determined by requiring that the operators $`\widehat{a}(\omega ,k)`$ satisfy the standard commutation relation
$$[\widehat{a}(\omega ,k),\widehat{a}^{}(\omega ^{},k^{})]=\delta ^{(3)}(\stackrel{}{k}\stackrel{}{k}^{})\delta (\omega \omega ^{})$$
Of course there is no ambiguity in this expansion. Following the logic of section 2 we then conclude that there are operators which create well defined states and the two point function of these operators are related to the bulk Green’s functions by a relation similar to (4.1)
$$<𝒪(k)𝒪(k)>_E=\mathrm{Lim}_{u,u^{}\mathrm{}}(\frac{k^4}{\mathrm{\Psi }_k(u)\mathrm{\Psi }_k(u^{})})𝒢_E(u,u^{};k)$$
where $`𝒢_E`$ is the euclidean Green’s function in this background. (In this expression $`k`$ stands for the four vector)
The Green’s function $`𝒢_E`$ can be easily computed since we know the two independent classical solutions - in fact this has already been computed in to obtain the D-instanton solution in these backgrounds. This is given by
$$\begin{array}{cc}& 𝒢_E(u,u^{};k)=\frac{\pi C(ka)}{4i(uu^{})^2A(ka)}H^{(1)}(\nu ,w+\frac{i\pi }{2})H^{(2)}(\nu ,w^{}\frac{i\pi }{2})u^{}>u\hfill \\ & 𝒢_E(u,u^{};k)=\frac{\pi C(ka)}{4i(uu^{})^2A(ka)}H^{(1)}(\nu ,w^{}+\frac{i\pi }{2})H^{(2)}(\nu ,w\frac{i\pi }{2})u>u^{}\hfill \end{array}$$
Using the asymptotic expressions for the Mathieu functions (equation (3.1)) one gets for $`u^{}>u>>1`$
$$𝒢_E(u,u^{};k)=𝒢_0(u,u^{};k)+\frac{1}{2ka^2}\frac{1}{(uu^{})^{5/2}}\frac{i\widehat{B}(ka)}{A(ka)}e^{ka^2(u+u^{})}$$
where $`𝒢_0`$ denotes the Green’s function in flat space
$$𝒢_0(u,u^{};k)=\frac{1}{2ka^2}\frac{1}{(uu^{})^{5/2}}e^{ka^2(u^{}u)}$$
In (5.1) we will have to substitute this asymptotic $`𝒢_E`$. In view of the fact that the Green’s function became a sum of the free piece and a “connected” piece suggests that it is natural to subtract out the free piece in (5.1). We will adopt this prescription.
Note that the expression (5.1) is unambiguous - there is no room for momentum dependent wave function renormalizations here. The Green’s function $`𝒢_E`$ is completely determined once the wave equation is known, including all normalizations and the wave functions are determined upto phases which in any case cancel in this expression.
If (5.1) is to make any sense, the $`u`$ dependence should cancel in the right hand side. From the asymptotic form of $`H^{(2)}`$ given in (3.1) and the form of the connected bulk Green’s function in (5.1) we see that this indeed happens. The final answer is
$$<𝒪(k)𝒪(k)>=k^4\frac{i\widehat{B}(ka)}{A(ka)}(\frac{1}{N(ka)})^2$$
We do not know how to obtain an explicit expression for $`N(ka)`$ for all $`ka`$ \- though we emphasize that it can be obtained in principle. However it is straightforward to find the behavior of this normalization factor for small values of $`ka`$. This is because in this regime the various Mathieu functions become Bessel functions. Using the relation (notations are those of )
$$H^{(2)}(\nu ,z)=\frac{\eta (ka)}{\chi (ka)}H^{(1)}(\nu ,z)\frac{2A(ka)}{C(ka)}J(\nu ,z)$$
we see that in this regime
$$H^{(2)}(\nu ,w\frac{i\pi }{2})\frac{2iA(ka)}{C(ka)}e^{i\frac{\pi }{2}(\nu +1)}I_\nu (k/u)$$
Since $`I_\nu `$ is the euclidean continuation of the normalizable solution in the $`AdS`$ limit we now know that for $`ka<<1`$ we must have
$$N(ka)\frac{C(ka)}{iA(ka)}$$
Plugging in the small $`ka`$ expansions for the various functions, given in (3.1) we easily verify that the leading nonanalytic term for small $`ka`$ is given by
$$<𝒪(k)𝒪(k)>k^4\mathrm{log}(ka)$$
which is the answer in the absence of a $`B`$ field.
The fact that we have obtained the correct low $`(ka)`$ behavior is nontrivial. This is because we have taken the boundary at $`u=\mathrm{}`$ before taking any low energy limit. We consider this result to be an evidence in favor of the holographic correspondence proposed in .
The procedure adopted above to obtain the correlation functions is in fact similar to the procedure adopted in for the case of NS five branes where again one has an asymptotically flat geometry in the decoupling limit. In this work the two point function is related to the S-matrix. This can be of course obtained from the euclidean Green’s function using a reduction procedure. Our procedure is of course valid for the standard AdS case where the space-time is not asymptotically flat.
6. Conclusions
The boundary theory proposed to be dual to the supergravity backgrounds we have considered in this paper - viz. NCYM theory - is nonlocal. It is likely that there are no local gauge invariant operators (in terms of the usual noncommutative gauge fields ) in this theory and correlation functions in position space do not make any obvious sense. In this paper we have argued that whether or not the boundary theory is indeed NCYM theory, states are still specified by values of the energy and (three dimensional) momenta. Thus there must be some operators $`𝒪(k)`$ in momentum space which create such states normalized in the standard fashion. The holographic correspondence identifies these states with normalized states in supergravity. We have shown that this implies an unambiguous prediction of the correlators of $`𝒪(k)`$. We verified that for small momenta they reproduce the expected result for usual YM theory.
If the boundary theory is indeed the NCYM theory, it is crucial to find the momentum space dual operators. In the standard $`AdS/CFT`$ correspondence one useful way to read off the dual operators is to consider the coupling of a three brane to background supergravity fields, using, e.g. a nonabelian version of the DBI-WZ action . One possibility is to explore these couplings in the presence of a $`B`$ field and in the low energy limit of . Rewriting DBI-WZ actions will not be sufficient in this case since the nontrivial modifications come when the supergravity backgrounds carry momenta in the brane direction so that derivatives of gauge fields are important. Some couplings of closed string modes to open string modes in the presence of $`B`$ field have been studied in . However it remains to be seen whether one could extract the couplings in terms of the standard fields of the NCYM theory. In particular it appears that the effect of $`B`$ fields may not be all encoded in the star product once closed string couplings are included.
7. Acknowledgements
We would like to thank R. Gopakumar, Y. Kiem, J. Maldacena, S. Minwalla and S. Trivedi for discussions. S.R.D. would like to thank the String Theory Group at Harvard University for hospitality during the final stages of this work.
8. Note Added:
While the paper was being typed, a related paper appeared on the net which has some overlap with our work.
References
relax A. Connes, M. Douglas and A. Schwarz, JHEP 9802 (1998) 003, hep-th/9711162; M. Douglas and C. Hull, JHEP 02 (1998) 008, hep-th/9711165. relax Y.K.E. Cheung and M. Krogh, Nucl. Phys. B528 (1998) 185, hep-th/9803031; F. Ardalan, H. Arfaei and M.M. Sheikh-Jabbari, JHEP 02:016 (1999), hep-th/9810072; C.S. Chu and P.M. Ho, Nucl. Phys. B550 (1999) 151, hep-th/9812219; M.M. Sheikh-Jabbari, Phys. Lett. B 455 (1999) 129, hep-th/9901080 and JHEP 06 (1999) 015, hep-th/9903205; V. Schomerus, JHEP 9906:030 (1999), hep-th/9903205; D. Bigatti and L. Susskind, hep-th/9908056. relax N. Seiberg and E. Witten, JHEP 9909:032 (1999), hep-th/9908142. relax H. Aoki, N. Ishibashi, S. Iso, H. Kawai, Y. Kitazawa and T. Tada, hep-th/9908141; N. Ishibashi, S. Iso, H. Kawai and Y. Kitazawa, hep-th/9910004; I. Bars and D. Minic, hep-th/9910091; J. Ambjorn, Y. Makeenko, J. Nishimura and R. Szabo, hep-th/9911041; S. Iso, H. Kawai and Y. Kitazawa, hep-th/0001027. relax S. Minwalla, M. Raamsdonk and N. Seiberg, hep-th/9912072; I. Ya. Aref’eva, D.M. Belov and A.S. Koshelev, hep-th/9912075 and hep-th/0001215; M. Raamsdonk and N. Seiberg, hep-th/0002186; M. Hayakawa, hep-th/9912094,hep-th/9912167; W. Fischler, E. Gorbatov, A. Kashani-Poor, S. Paban, P. Pouliot and J. Goomis, hep-th/0002067; A. Matusis, L. Susskind and N. Toumbas, hep-th/0002075, O. Andreev and H. Dorn, hep-th/0003113; Y. Kiem and S. Lee, hep-th/0003145; I. Ya. Aref’eva, D.M. Belov A.S. Koshelev and O.A. Rytchkov, hep-th/0003176; A. Bilal, C. Chu and R. Russo, hep-th/0003180; J. Gomis, M. Kleban, T. Mehen, M. Rangamani and S. Shenker, hep-th/0003215; A. Rajaraman and M. Rozali, hep-th/0003227; H. Liu and J. Michelson, hep-th/0004013; G. Arcioni, J. Barbon, J. Goomis and M. Vazquez-Mozo, hep-th/0004080; F. Zanora, hep-th/0004085; J. Gomis, K. Landsteiner and E. Lopez, hep-th/0004115. relax J. Maldacena, Adv.Theor. Math. Phys. 2 (1998) 231, hep-th/9711200 relax S. Gubser, I. Klebanov and A. Polyakov, Phys. Lett. B428 (1998) 105,hep-th/9802109; E. Witten, Adv. Theor. Math. Phys. 2 (1998) 253, hep-th/9802150; D. Freedman, S. Mathur, A. Matusis and L. Rastelli, Nucl. Phys. B546 (1999) 96, hep-th/9804058. relax A. Hashimoto and N. Itzhaki, hep-th/9907166. relax J. Maldacena and J. Russo, hep-th/9908134. relax S.R. Das, S. Kalyana Rama and S. Trivedi, JHEP 0003 : 004, 2000; hep-th/9911137. relax M. Li, hep-th/9909085; A. Dhar, G. Mandal, S. Wadia and K.P. Yogendran, hep-th/9910194; T. Harmark and N. Obers, JHEP 0003:024 (2000), hep-th/9911169; J. Liu and S. Roy, hep-th/9912165; J. Barbon and A. Pasquinucci, hep-th/0002187; N. Ishibashi, S. Iso, H. Kawai and Y. Kitazawa, hep-th/0004038. relax S. Minwalla and N. Seiberg, JHEP 9906 (1999) 007, hep-th/9904142 relax S. Gubser and A. Hashimoto, Comm. Math. Phys. 203 (1999) 325, hep-th/9803023. relax V. Balasubramanian, P. Kraus and A. Lawrence, Phys.Rev. D59 (1999) 04600, hep-th/9805171; T. Banks, M. Douglas, G. Horowitz and E. Martinec, hep-th/9808016; V. Balasubramanian, P. Kraus, A. Lawrence and S. Trivedi ,Phys.Rev. D59 (1999) 104021, hep-th/9808017 relax S.R. Das and S. Trivedi, Phys. Lett. B445 (1998) 142, hep-th/9804149; S. Ferrara, M. Lledo and A. Zaffaroni, Phys. Rev. D58 (1998) 105029, hep-th/9805082. relax S. Hyun, Y. Kiem, S. Lee and C.Y. Lee, hep-th/9909059 relax U. Danielsson, A. Guijosa, M. Kruczenski and B. Sundborg, hep-th/0004187.
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# Dense gas in the dust lane of Centaurus ABased on observations collected at the European Southern Observatory, La Silla, Chile
## 1 Introduction
Centaurus A (NGC 5128) is the closest radio galaxy (distance 3.5 Mpc, 1$`{}_{}{}^{\prime \prime }=`$ 17 pc, Soria et al. 1996, Hui et al. 1993, Israel 1998, Ebneter & Balick 1983, de Vaucouleurs 1979) and exhibits a very prominent dust lane. Centaurus A is a strong radio galaxy with a milliarcsecond nuclear continuum source (Kellermann et al. 1997; Shaffer & Schilizzi 1975; Kellermann 1974) and two giant radio lobes. Absorption against the nuclear source has been found in HI (van der Hulst et al. 1983) and many molecular species and transitions (Gardner & Whiteoak 1976; Whiteoak & Gardner 1971; Bell & Seaquist 1988; Seaquist & Bell 1986, 1990; Phillips et al. 1987; Eckart et al. 1990a; Israel et al. 1990, 1991; Wild et al. 1997; Wiklind & Combes 1997).
Several studies of Centaurus A in the millimeter wavelength range have been carried out. Although other elliptical galaxies with dust lanes have recently been detected in CO (Sage & Galleta 1993), Centaurus is the best object for a detailed study due to its proximity and corresponding large angular size. Phillips et al. (1987) and Quillen et al. (1992) observed several positions in the CO(2-1) line along the dust lane at a resolution of 30<sup>′′</sup>. In previous papers we presented a fully sampled map of the <sup>12</sup>CO(1-0) emission together with IRAS observations of the FIR continuum (Eckart et al. 1990b), measurements of the millimeter absorptions lines towards the nucleus (Eckart et al. 1990a; see also Wiklind & Combes 1997), and a <sup>12</sup>CO(2-1) map along the dust lane at a resolution of 22<sup>′′</sup> (Rydbeck et al. 1993).
In Wild et al. (1997) we presented the first fully sampled <sup>13</sup>CO(1-0) map along the dust lane of Centaurus A, as well as single spectra of the <sup>13</sup>CO(2-1) emission in the disk and C<sup>18</sup>O(1-0) emission at the central position. These new data allowed us, in combination with the <sup>12</sup>CO(1-0) and <sup>12</sup>CO(2-1) maps obtained earlier, to study the excitation conditions of the molecular gas in detail throughout the dust lane. Using different CO line ratios and their variation across the disk of Centaurus A, we inferred the physical parameters of the molecular ISM and their spatial variations.
Here we extend our investigation to the very dense (10<sup>4</sup> to 10<sup>5</sup>cm<sup>-3</sup>) phase of the molecular gas by observing line emission of the density tracers HCN and CS towards the nucleus and the off-nuclear dust lane in Centaurus A.
## 2 Observations
The observations were carried out in two observing runs in January 1996 and January 1998 with the 15m Swedish ESO Submillimeter Telescope (SEST) on La Silla, Chile. The adopted central position for Centaurus A was $`\alpha (1950)`$ $`=13^h22^m31.8^s`$ and $`\delta (1950)=42^{}45^{}30^{\prime \prime }`$. We observed the <sup>12</sup>CO(1-0), HCN(1-0), CS(2-1), and CS(3-2) molecular line emission and absorption towards the central non-thermal continuum source and selected positions in the dust lane. We used the 3 mm and 2 mm SIS receivers with system temperatures around 140 K and 180 K, respectively, and beam widths ranging from 34<sup>′′</sup> (FWHM) at 147 GHz to 56<sup>′′</sup> at 89 GHz. We observed in the dual beam switching mode, i.e. a chopper wheel switched between two positions on the sky displaced by about 12 arcminutes in azimuth. First the source was placed in one beam and then in the other beam. The backend was the low resolution acousto-optical spectrometer with a bandwidth of about 1 GHz in 1440 channels. The separation between channels was 0.7 MHz, and the actual spectral resolution 1.4 MHz (4.7 km/s at 89 GHz). We adopted the LSR velocity scale. The system was flux calibrated with the chopper wheel method (Kutner & Ulich 1981). We used main beam efficiencies $`\eta _{\mathrm{MB},89\mathrm{G}\mathrm{H}\mathrm{z}}`$ = 0.75, $`\eta _{\mathrm{MB},97\mathrm{G}\mathrm{H}\mathrm{z}}`$ = 0.73, $`\eta _{\mathrm{MB},115\mathrm{G}\mathrm{H}\mathrm{z}}`$ = 0.70, and $`\eta _{\mathrm{MB},147\mathrm{G}\mathrm{H}\mathrm{z}}`$ = 0.66 for the conversion from antenna temperature to Rayleigh-Jeans main beam brightness temperatures. The pointing was accurate to within 3<sup>′′</sup>. It was checked frequently on the nuclear continuum source of Centaurus A, and the nearby (distance $`17^{}`$) SiO maser W Hya. Fig. 1 shows the positions in the dust lane of Centaurus A at which the dense molecular gas has been investigated. Fig. 2 and Fig. 3 show the HCN J=1–0 and CS spectra along the dust lane. Table 1 gives parameters of the observed spectra.
## 3 Analysis
In order to measure the mass of dense molecular gas and therefore the star formation potential in the dust lane of Centaurus A we observed HCN and CS line emission. Both molecules trace higher density gas than CO, since they have large dipole moments and require n(H<sub>2</sub>)$``$10<sup>4</sup>cm<sup>-3</sup> for significant excitation. The analysis is done by calculating the line luminosities and estimating from them the amount of dense molecular gas. These properties are then discussed and compared to data from the Milky Way and external galaxies.
Line luminosities can be obtained from the integrated line intensities $`I_{\mathrm{line}}=T_{\mathrm{MB}}\times \delta v`$ via
$$L_{\mathrm{line}}=I_{\mathrm{line}}D^2\mathrm{\Omega },$$
(1)
where $`T_{\mathrm{MB}}`$ is the main beam brightness temperature, D is the distance to the source and $`\mathrm{\Omega }`$ is the solid angle of the beam convolved with the source.
Using radiative transfer solutions (e.g. Kwan & Scoville 1975, Linke & Goldsmith 1980), assuming that the HCN and CS line emission traces gravitationally bound or virialized clouds with a density range of a few 10<sup>4</sup> cm<sup>-3</sup> to a few 10<sup>5</sup> cm<sup>-3</sup> and kinetic temperatures between 10 to 60 K. Solomon et al. (1990, 1992) present relations between the mass of molecular gas at these densities and the corresponding observed line luminosity as defined above. For HCN(1-0) they find in their 1992 paper
$$M_{\mathrm{HCN}}20_{10}^{+30}\times L_{\mathrm{HCN}}\mathrm{M}_{}(\mathrm{K}\mathrm{km}\mathrm{s}^1)^1$$
(2)
and for CS(2-1) they derive in the 1990 paper
$$M_{\mathrm{CS}}(35150)\times L_{\mathrm{CS}}\mathrm{M}_{}(\mathrm{K}\mathrm{km}\mathrm{s}^1)^1.$$
(3)
The assumptions made above are likely to be valid for the molecular ISM in Centaurus A as well. This can be based on the fact that, according to radiative transport calculations (Eckart et al. 1990a), the lower density molecular gas - as traced by CO isotopic emission - has similar properties compared to those in Galactic GMCs and galaxies with indications for enhanced star formation activity.
Although the mass estimates derived using the relations given above include considerable uncertainties they can be used as a general guide line to investigate the properties of the dense interstellar medium. The assumed kinetic temperature range appears to be appropriate for the dense star forming ISM in Centaurus A. Eckart et al. (1990b), Joy et al. (1988), as well as Marston & Dickens (1988) give dust temperatures in the dust lane in the range of 40 K. From the HCN(1-0)/HNC(1-0) ratio of about 2, Israel (1992) introduces kinetic temperatures of the order of $``$25 K.
Table 2 lists the integrated line intensities $`T_{\mathrm{MB}}\times \delta v`$ for the HCN(1-0) and the CO(1-0) line emission, as well as the ratios I(HCN)/I(CO) corrected for main beam efficiencies. In Table 3 we list the CO(1-0) and HCN(1-0) line luminosities as well as the estimated FIR luminosities $`L_{\mathrm{FIR}}`$ as a function of position and as estimated quantities integrated over the dust lane. Table 4 lists line intensities $`T_{\mathrm{MB}}\times \delta v`$ and luminosities (corrected for main beam efficiency) for CS(2-1) and CS(3-2) lines we measured at a few positions. Table 5 lists the mass of dense molecular gas as derived from the line luminosities and Eqs. (2) and (3).
## 4 Results and Discussion
### 4.1 Dense gas traced by HCN
HCN traces molecular gas at much higher density of about $``$10<sup>4</sup>cm<sup>-3</sup> than CO ($``$500cm<sup>-3</sup>). In Cen A the HCN emission within the dust lane is clearly detected over a similar velocity range as CO ($`v_{\mathrm{LSR}}`$=250 - 850 km/s). The HCN emission, however, is stronger peaked on the nucleus compared to the CO emission. In Fig. 4 we compare the spatial variation of the intensity ratio I(HCN)/ I(CO) as it is observed in Cen A and M 51 (Kohno et al. 1996). At the central position, the integrated intensity ratio I(HCN)/ I(CO) peaks at 0.06, and decreases to $``$0.02 to 0.04 in the dust lane. Ratios of the order of 0.1 or higher are only observed in active nuclear regions of Seyfert galaxies and ULIRGs (Kohno et al. 1996; Solomon et al. 1992; Nguyen-Q-Rieu et al. 1992; Helfer & Blitz 1993, 1995; Jackson et al. 1993; Tacconi et al. 1994; Sternberg et al. 1994).
The ratio of HCN to CO luminosity is 1/6 for ULIRGs, but only 1/80 in normal spirals (Solomon et al. 1992). For Cen A the HCN to CO luminosity is $`L_{\mathrm{CO}}/L_{\mathrm{HCN}}`$=1/13 and therefore closer to the value for ULIRGs than normal spirals.
Fig. 5 shows that Centaurus A has the same ratio of FIR to HCN luminosity as ULIRGs and normal spirals, including the Milky Way. This is also true for the center and the off-positions in the dust lane. This suggests that in Cen A the star formation rate per mass of dense gas is in good approximation independent on the position and infrared luminosity within the dust lane.
Detailed observations (Lee et al. 1990) of Galactic GMC’s show that the average CS(2-1)/CO(1-0) intensity ratio is $``$1/300. Since HCN(1-0) is usually 1.5 to 2 times stronger than CS(2-1) and has larger source sizes Solomon et al. (1992) adopt an HCN(1-0)/CO(1-0) ratio of $``$1/100 for the molecular ring in the Galactic disk. This results in a HCN luminosity for the Milky Way of about 4$`\times `$10<sup>6</sup> K km s$`^1`$pc<sup>2</sup> which is quite comparable to what we find for Centaurus A with 5.5$`\times `$10<sup>6</sup> K km s$`^1`$pc<sup>2</sup>.
In Fig. 6 we show the position of Centaurus A in a plot of $`log(L_{\mathrm{HCN}}/L_{\mathrm{CO}})`$ which measures the fraction of dense gas and $`log(L_{\mathrm{FIR}}/L_{\mathrm{CO}})`$ which measures the efficiency with which molecular gas is transformed into OB stars (Solomon et al. 1992). This plot includes ULIRGs as well as normal and infrared luminous galaxies. We find that as a whole Cen A fits very well into this correlation and is located between both groups. This is also true for the disk at angular separations of about 60<sup>′′</sup> and 90<sup>′′</sup> from the center. However, towards the nucleus the fraction of dense molecular gas measured via the line luminosity ratio L(HCN)/L(CO) as well as the star formation efficiency L<sub>FIR</sub>/L<sub>CO</sub> is more comparable to ULIRGs rather than normal and infrared luminous galaxies. This suggests that most of the FIR luminosity of Centaurus A originates in regions of very dense molecular gas and high star formation efficiency.
The HCN line luminosity can now be used to estimate the amount of dense (10<sup>4</sup> cm<sup>-3</sup>) gas via Eq. (2). From the HCN line luminosity we find a total mass of molecular gas that must be at densities $`>`$10<sup>4</sup>cm<sup>-3</sup> of $`8.6_4^{+13}\times 10^7`$M. This has to be compared to the total molecular gas mass derived from the <sup>12</sup>CO(1-0) line of 2$`\times `$10<sup>8</sup>M (Eckart et al. 1990a). The comparison shows that a large fraction - approximately one third - but at least one sixth - of the molecular line emission in Cen A must originate from sites with abundant dense molecular gas. The ratio of molecular gas mass at densities of $`>`$10<sup>4</sup>cm<sup>-3</sup> to gas at $``$300cm<sup>-3</sup> is about 1/20 in the Galaxy and almost unity for ULIRGs (Solomon et al. 1992). Based on this quantity the physical conditions of the dense molecular ISM in Cen A are apparently closer to those of ULIRGs than to our Milky Way. The presence of a large amount of dense molecular gas is also supported by the fact that if one multiplies $`M_{\mathrm{HCN}}(H_2)`$ by the mean Galactic $`L_{\mathrm{FIR}}/M_{\mathrm{HCN}}(H_2)`$ ratio of 71 $`M_{}(\mathrm{K}\mathrm{km}\mathrm{s}^1)^1\mathrm{pc}^2`$ (Solomon et al. 1992) one obtains about 6$`\times `$10<sup>9</sup>L which equals the observed FIR luminosity. The implicite assumption here is that the mean Galactic conversion factor is applicable. Accepting this, the result implies that active star formation in the dust lane of Cen A is the actual source of its FIR luminosity and that the AGN which is responsible for strong radio and X-ray radiation is not contributing substantially.
### 4.2 Dense gas traced by CS
The CS(2-1) and CS(3-2) lines trace even denser molecular gas at 10<sup>5</sup>cm<sup>-3</sup>. The problem is that CS line emission is quite weak. Detections of CS in extra-galactic sources are sparse. In the sample of normal and interacting galaxies and ULIRGs Solomon et al.(1992) detected CS line emission only in Arp 220. For Cen A we find a total CS(3-2) line luminosity of $`L_{\mathrm{CS}(32)}5.4\times 10^5\mathrm{K}\mathrm{km}\mathrm{s}^1\mathrm{pc}^2`$. The ratio of the CS to CO line luminosity is then $`L_{\mathrm{CS}(32)}/L_{\mathrm{CO}}1/170`$, very similar to the value of 1/250 given by Solomon et al. (1990) for the Milky Way. For the Galactic Center the same authors derive a corresponding ratio of about 1/20 to 1/30. This indicates that the Galactic Center is approximately 8 to 10 times more luminous in the CS(3-2) line that the molecular gas in the disk. At the center of Cen A the $`L_{\mathrm{CS}(32)}`$ is 2 times higher compared to the off-center dust lane measurement.
The CS(3-2) line luminosity results in a total mass estimate of gas at densities of 10<sup>5</sup>cm<sup>-3</sup> of about $`M_{\mathrm{CS}(32)}(H_2)`$ = (2-8)$`\times `$10<sup>7</sup>M. This number is in very good agreement with the corresponding value derived from the HCN(1-0) line emission.
## 5 Conclusion
While the line luminosities of the HCN(1-0) and CS(3-2) as well as the FIR luminosity of the molecular gas in the dust lane of Cen A are quite comparable to each other there are definite differences in the overall fraction of dense molecular gas and the efficiency with which stars are formed from it. This star formation activity is also the source of the FIR luminosity of Cen A. About 40% or even more of the total molecular line luminosity in Cen A originates in dense gas. This suggests that star formation as well as the bulk of the dense molecular gas is mostly concentrated in GMC complexes rather than in a more diffuse molecular gas component. This is already indicated by the ring-like distribution of HII region found by Graham (1979) as well as the MEM deconvolved <sup>12</sup>CO(2-1) line emission mass by Rydbeck et al. (1993).
We also note that with respect to other positions in the dust lane the I(HCN)/I(CO) ratio is larger at separations of about 100<sup>′′</sup> from the nucleus rather than at separations of about 60<sup>′′</sup>. The larger distance is close to the inner edge of the ring of HII regions and corresponds well with the position of the folds in the warped molecular gas disk of Centaurus A (Quillen et al. 1992, 1993; Sparke 1996) and an increased intensity in the 15$`\mu `$m continuum dust emission (Block & Sauvage 2000; Mirabel et al. 1999). Therefore a higher I(HCN)/I(CO) ratio may be due to a combination of enhanced star formation efficiency at these positions and an increase in column density due to the folds. Due to the low intrinsic velocity dispersion of the thin molecular disk (Quillen et al. 1992) and due to the fact that the <sup>12</sup>CO line emission is originating in optically thick molecular gas (Wild et al. 1997) this may lead to shadowing of molecular clouds along the line of sight toward the folds. This effect will be stronger for <sup>12</sup>CO than for the small, dense cloud cores seen in the less abundant HCN line emission. This effect may therefore lead to an intensity decrease in the <sup>12</sup>CO line and an increase in the HCN(1-0) line, resulting in the observed variation of the I(HCN)/I(CO) line ratio.
Future interferometric measurements will allow us to study the distribution of molecular gas in the dust lane of Centaurus A in much greater detail. Line ratios, luminosities and star formation can then be investigated for individual GMC complexes.
###### Acknowledgements.
We are grateful to the SEST team and the ESO staff on La Silla and in Garching for their support and hospitality. We thank Lars-Åke Nyman for taking an additional central CO spectrum.
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# Dirac, Anderson, and Goldstone on the Kagomé
## I Introduction
The spin 1/2 Heisenberg antiferromagnet on the kagomé lattice is a good candidate for a two-dimensional quantum system with a spin disordered ground state. While it appears that on square and triangular lattices an antiferrmomagnet will acquire Néel order, on the kagomé lattice strong numerical evidence has accumulated that the system is spin disordered, as seen by the existence of a gap to triplet excitations and through consideration of the spectra of finite size samples. Numerically, one finds a continuum of low energy states below the triplet gap. The continuum of low energy excitations provides a great puzzle to theory in the absence of an obvious broken symmetry.
There are good experimental realizations of kagomé systems, despite the presence of additional couplings, including the jarosites and $`SrCrGaO`$. While in iron jarosites, these additional couplings produce long-range order, in deuteronium jarosite and $`SrCrGaO`$ no long range order is seen. Additionally, in $`SrCrGaO`$ a quadratic specific heat and very weak field dependence of the specific heat are in agreement with the picture of a continuum of low-energy singlet excitations seen in numerics, suggesting that the latter two compounds provide good realizations of the kagomé antiferromagnet.
Given the lack of spin order, RVB ideas seem natural for this system, and indeed have stimulated much theoretical work on the system. Large $`N`$ calculations based on $`SU(N)`$ have been used to suggest a spin-Peierls state. Calculations based on $`Sp(2N)`$ have suggested a phase with deconfined, gapped, bosonic spinons. Chiral states have also been proposed, but do not account for the excitation spectrum and also are in disagreement with the rapid decay of chirality-chirality correlation functions seen in numerics. States with BCS pairing have been suggested but again do not account for the excitation spectrum; due to the non-bipartite nature of the kagomé lattice, these states are not equivalent to flux states. In addition to the long-range states, short-range RVB states based on a reduced Hilbert space of dimers have also been considered and provide some explanation for the gapless continuum.
An RVB state on the kagomé lattice would be particularly attractive, given the intensive work on RVB states on the square lattice, especially in connection with high-$`T_c`$ materials. In the absence of doping, the square magnet eventually acquires Néel order and the spinon excitations disappear from the system. Since the kagomé lattice does not acquire Néel order, it could be a very important model for a spin liquid or spin solid state.
The idea behind the present approach is a to start with a long-range RVB treatment of the kagomé lattice and consider various ways of gapping the spinon excitation spectrum. We will first construct a “parent state” which will be the best RVB state that does not break time-reversal symmetry or any lattice symmetry.
We will then demonstrate an interesting massless Dirac structure for this state. Various other known RVB states can be obtained by perturbing the parent state, lowering the symmetry and giving mass to the Dirac particles, so that the parent state unifies a wide class of states. Physically, we expect that the system will attempt to give mass to the Dirac particles and open a gap, picking out one of these other lower symmetry states. We will discuss the symmetry breaking through a renormalization group treatment. We will obtain some kind of spin solid state, and some low energy Goldstone and gauge excitations, which we will argue provide the low energy degrees of freedom seen experimentally.
The RVB states can be thought of by a decoupling procedure, in which we decompose spin-1 operators into pairs of spin-1/2 operators. Take a Hamiltonian
$$H=\underset{<i,j>}{}J\stackrel{}{S}_i\stackrel{}{S}_j$$
(1)
where the sum extends over neighboring sites $`i,j`$.
Introduce the spinon fields $`\psi _a^{}(i),\psi _a(i)`$, where $`a=u,d`$ labels up and down spinon fields. We can then introduce a Hubbard-Stratonovich field $`t_{ij}=t_{ji}^{}`$ such that
$$H=\underset{<i,j>}{}(\psi _a^{}(i)t_{ij}\psi _a(j)+h.c.)+\frac{2}{J}\underset{<i,j>}{}|t_{ij}|^2$$
(2)
By taking a mean-field in $`t`$, minimizing the total energy of the fermions and the Hubbard-Stratonovich field, we obtain an RVB state. One must at some point project results onto the physical space in which each site is singly occupied.
Later, we will find this projection to be extremely important. In the absence of projection, the ideal mean-field state is almost always found by taking a dimer covering of the lattice, with $`t_{ij}`$ nonvanishing only on the given dimers. Projection can stabilize RVB states, so although our first calculations will ignore the effects of projection, in a naive mean-field, we will later discuss a projected mean-field that includes some of the essential effects of projection.
We will then have to proceed beyond mean-field solutions. We will consider a functional integral with fields $`\psi (i)`$ and $`t_{ij}`$, fluctuating about a saddle-point of the action. There are a large number of possible fluctuations in $`t_{ij}`$, including a set of pure gauge fluctuations, as well as a set of gauge fields. Most other fluctuations can be ignored because they do not contribute to the low-energy dynamics. However, there will be a particular set of fluctuations in $`t_{ij}`$ that produce a mass for the fermion field. Although these fluctuations are not gapless, we will retain these fluctuations due to their impact on the low energy dynamics of the fermion field. We will see using a renormalization group that the effective action of these fields can differ greatly from that suggested by the mean-field.
To outline the paper, we will first describe the parent state, and then discuss how to perturb the parent state to obtain other proposed RVB states. Then we will discuss naive and projected mean-field theory treatments of these states. We will the proceed to a field theoretic treatment of fluctuations about the mean-field and a renormalization group that will suggest one particular symmetry breaking pattern. We will discuss the pseudo-Goldstone and gauge modes that arise from this symmetry breaking, and the mechanism that ultimately gives them a very small energy gap.
Next, we proceed to a discussion of finite system size effects as a first step in comparison with numerics. These effects lead to an additional flux for odd system sizes which leads to a nonvanishing Chern number for odd system sizes. We will then compare the low energy bosonic modes from the field theory to the low energy singlet modes found in numerical calculations as well as checking dimer-dimer correlation functions and many-body density of states.
## II The Parent RVB State
Although the short range RVB calculations provide one starting point for the kagomé lattice antiferromagnet, we will be interested in looking at long range states instead. Certainly, the short range RVB calculations themselves suggest that long range antiferromagnetic correlations are important; the variational energy of these states improve when second neighbor dimers are included. Further, while the kagomé lattice has a gap to triplet excitations, this gap is about an order of magnitude smaller than $`J`$; from a short range calculation one might expect a gap of order $`J`$ as that is the energy to break a dimer. However, Mila has suggested that within a short-range state it is possible to have a triplet gap much smaller than $`J`$, so the small triplet gap does not necessarily rule against a short-range state.
The best RVB state on the kagomé lattice antiferromagnet is a chiral spin liquid. A similar chiral state was obtained using a hardcore boson representation of the spins and transmuting the statistics from bosonic to fermionic using a Chern-Simons field. However, numerical calculations do not support a large chirality-chirality correlation function or expectation value of the chirality operator, which would seem to rule these states out. So, we will look for the best RVB state that is not chiral.
Assuming that we are looking for a long-range RVB state in which all $`t_{ij}`$ have the same magnitude, the only choice we have is how much flux to put into the system. The state we choose involves putting $`\pi `$ flux through the hexagons, and no flux through the triangles. This state offers a better mean field energy than any other nonchiral RVB state, including the state with no flux through the system at all.
The unit cell of the kagomé lattice consists of three sites on a triangle. Once we add flux to the system, the unit cell doubles, and requires six sites on two triangles. We will find it convenient to double the unit cell again, to twelve sites, including six sites on a hexagon and the six sites which neighbor the hexagon. This cell is shown in figure 1. The kagomé lattice is made up of a triangular lattice of these 12-site unit cells. We have numbered the points in the cell for later reference.
For now, let us assume that we pick $`J`$ such that $`|t|=1`$ within our RVB state. Then, the band structure for our RVB state is shown in figure 2 scanning along the given line of momenta in the Brillouin zone. There is a degeneracy of states: the bottom line in the figure actually consists of four bands, while the other four lines in the figure consist of two bands each, providing a total of twelve bands. At $`(0,0)`$, four bands meet at energy less than zero and another four meet at positive energy. Near this point the spectrum becomes relativistic.
The particles occupy the lowest six bands of the system, meaning that where the bands meet the spectrum becomes gapless. The system can gain energy by perturbing about our given RVB state; we expect that the greatest gain in energy comes from opening a gap. For this reason, we will study the Dirac point and look at possible perturbations to the Dirac equation.
At the Dirac point, the Schrödinger equation for the fermions becomes
$$E\psi =v_f(\alpha _xk_x+\alpha _yk_y)\psi $$
(3)
where $`\psi `$ is a four component spinor, and the matrices $`\alpha _x,\alpha _y`$ are anti-commuting $`\alpha `$-matrices. The particular basis chosen for $`\psi `$ and for the $`\alpha `$-matrices is unimportant. We find by explicit computation that
$$v_f=(0.408248\mathrm{})|t|$$
(4)
Given the Dirac equation (3), we would like to consider the effect of perturbations $`\delta t_{ij}`$ on the low energy structure. This analysis will enable us to focus on those fluctuations in $`t_{ij}`$ which have the greatest impact on the low-energy dynamics and which must be kept when we proceed to a field-theory treatment of fluctuations.
From equation (2), the system pays an energy cost equal to $`\frac{1}{J}|\delta t_{ij}|^2`$, but it can gain energy by opening a gap for the Dirac particles. As a result, we look for the perturbations which open the greatest gap for the Dirac particles for given $`|\delta t_{ij}|^2`$. Let us first proceed algebraically, considering possible perturbations to the Dirac equation which will open a gap, and only then ask how to obtain these perturbations from $`\delta t_{ij}`$.
A given perturbation $`\delta t_{ij}`$ will perturb the Dirac equation to
$$E\psi =\left(v_f(\alpha _xk_x+\alpha _yk_y)+M\right)\psi $$
(5)
where $`M`$ is some matrix, the projection of $`\delta t_{ij}`$ onto the space of the four states at the Dirac point. It may be shown that the perturbation $`M`$ will be most efficient at opening a gap when it anti-commutes with $`\alpha _x,\alpha _y`$. By efficient, we mean that we wish to maximize the gap for given $`\mathrm{Tr}(M^2)`$, as a first step to maximizing the gap for given $`|\delta t_{ij}|^2`$. Since there is only a sixteen dimensional space of matrices $`M`$, we can easily characterize all matrices that have the needed anti-commutation property; it is a four dimensional vector space.
We will write three of the perturbations as matrices $`M_i`$, for i=1,2,3. In terms of $`\alpha `$-matrices, they will be $`M_1=\alpha _z,M_2=\beta ,M_3=\beta \alpha _x\alpha _y\alpha _z`$. These perturbations anti-commute with each other; in fact, one can make a change in basis in the Dirac equation which leaves $`\alpha _x,\alpha _y`$ unchanged, but produces continuous $`O(3)`$ rotations in the space of $`M_1,M_2,M_3`$. This continuous symmetry is only valid at low energy; it will be broken to a discrete symmetry by lattice effects as discussed below. By taking $`M=_im_iM_i`$, for some numbers $`m_1,m_2,m_3`$, we open a gap equal to $`\sqrt{_i(m_i^2)}`$. We will refer to these as nonchiral mass terms.
The fourth perturbation is of a different sort. It is $`M=m_cM_c`$ with $`M_c=i\alpha _x\alpha _y`$ (here, c stands for chiral and we will refer to this as a chiral mass term). This perturbation breaks parity and time-reversal symmetry. $`M_c`$ commutes with $`M_1,M_2,M_3`$. As mentioned above, we are interested in the most efficient way for the system to open a gap; since $`M_c`$ does not anti-commute with $`M_i`$, it is most efficient for the system to take either purely chiral mass or purely nonchiral mass, so that $`m_i=0`$ or $`m_c=0`$.
Next, we would like to ask what perturbations in the $`t_{ij}`$ will produce the desired mass matrix $`M`$. We will find that to produce $`M_i`$ requires dimerizing the system by making the magnitudes of the $`|t_{ij}|`$ non-uniform; to produce $`M_c`$ requires adding additional fluxes to the system. Clearly, there is a large degeneracy here, as there is a only 16 dimensional space of matrices $`M`$ while in a given unit cell there is a 48 dimensional space of perturbations to $`t_{ij}`$. While some of the degeneracy is due to the large number of possible gauge transforms on $`t_{ij}`$, this does not completely alleviate the problem. Again, the question of efficiency becomes important: for each $`M`$, there is a class of $`t_{ij}`$ which produce the desired perturbation, but only one element in the class minimizes $`_{i,j}|t_{ij}|^2`$. We will find one unique perturbation in $`t_{ij}`$ (up to arbitrariness in gauge) which produces the desired mass matrix.
The perturbation that we will pick for $`M_1`$ is shown in figure 3. One can see that all the horizontal bonds have been either decreased or increased in strength, such that along a horizontal line the bonds alternate in strength while horizontal bonds which are in a vertical column all have the same strength.
This perturbation is a dimerization of the $`t_{ij}`$. The spins form singlets most strongly across the largest $`t_{ij}`$, so that dimerization of $`t_{ij}`$ tends to produce a spin solid and changes the long-range RVB to a short range set of singlets.
This mass term breaks rotational and translational symmetry of the lattice. We will pick $`M_2`$ and $`M_3`$ to be lattice rotations of $`M_1`$. The continuous symmetry of the Dirac equation will then be broken at the lattice level to a discrete symmetry of permutations of $`m_1,m_2,m_3`$ under lattice rotation, while lattice translations change the sign of any 2 of the 3 mass terms $`m_1,m_2,m_3`$. We will find later that while we obtain symmetry breaking and produce a mass, the discrete nature of the lattice group will leave us with only pseudo-Goldstone modes.
There are two mass terms which are symmetric under rotations. They are
$$M_{12}=\frac{M_1+M_2+M_3}{\sqrt{3}}$$
(6)
$$M_6=M_{12}=\frac{M_1M_2M_3}{\sqrt{3}}$$
(7)
where 12 denotes the fact that the $`t_{ij}`$ are strongest on the 12-site loop surrounding the unit cell, while 6 denotes the fact that the bonds are strongest on the hexagons and triangles. We show the perturbation to $`t_{ij}`$ to produce $`M_{12}`$ in figure 4.
To estimate dimerization later, it will be useful to know connect the change in $`t`$ to the eigenvalues of the mass matrix that arises. One finds that if the bonds on the 12-site loop are increased by $`\delta t`$, while those on the hexagons and triangles are decreased by the same amount, then one produces a term $`m_{12}M_{12}`$ in the Dirac equation with $`m_{12}=1.57735\mathrm{}.\delta t`$.
For calculations later, it will be convenient to transform to a basis of gamma matrices. Defining $`\gamma _t=\beta `$, $`\gamma _i=\beta \alpha _i`$, we find that $`M_1=\gamma _3`$, $`M_2=1`$, $`M_3=i\gamma _5`$, while $`M_c=i\gamma _t\gamma _x\gamma _y`$.
It is interesting to compare to above characterization of possible perturbations in terms of gamma matrices to the situation in the $`\pi `$-flux phase on the square lattice, where there is again a Dirac spectrum, and again various mass terms can be introduced. We find that we are able to introduce one nonchiral mass term by dimerizing the horizontal bonds of the square lattice, so that the horizontal bonds alternate in strength as one moves horizontally along the lattice; another mass term can be introduced by dimerizing the vertical bonds. In the limit of extreme dimerization, these states correspond to short-range RVB states in which the dimers are stacked on top of each other, and all lie either horizontally or vertically. By taking sums of these two mass terms, we can produce a short range state in which dimers resonate around a square. The final nonchiral mass term can be obtained by placing an on-site potential on one sublattice of the square lattice; this corresponds to introducing Néel order into the system. Due to the highly frustrated nature of the kagomé lattice, in this paper we will not have any such terms involving introducing on-site potentials.
For the square lattice, the chiral mass term can also be introduced. It requires adding additional couplings to the system, which connect diagonally across a given plaquette, and then inserting $`\pi /2`$ flux through the triangle that is formed when a particle traverses two sides of a plaquette and then crosses the plaquette on the diagonal.
## III Connection to Other Valence Bond States
From the parent state, it is possible to continuously connect to other possible valence bond states, using the mass terms we have found above. Let us first consider the chiral spin liquid and then the spin-Peierls solid of Marston and Zeng.
Let us consider a state such that the flux through each triangle is equal to $`\theta `$, and the flux through each hexagon is equal to $`\pi 2\theta `$. Then, as $`\theta `$ varies from $`0`$ to $`\frac{\pi }{2}`$ we continuously deform from the parent state to the chiral spin liquid. By looking at how $`t_{ij}`$ changes along this deformation, and then projecting this change onto the Dirac point, we find that for small $`\theta `$, the perturbation exactly produces the chiral mass term $`M_c`$! Let us note that at $`\theta =\frac{\pi }{4}`$ there is a highly interesting band structure, discussed in the Appendix, with multiple flat bands.
The chiral spin liquid state improves on our parent state at the mean-field level. Since we will argue below that fluctuations stabilize our state against the chiral mass term, let us here analyze why the chiral state works at the mean-field level, and provide a qualitative explanation of why fluctuations destroy the chiral spin liquid.
The idea behind the chiral state results from the “Rokshar rules”, which argue that one should put flux $`\frac{\pi }{2}`$ through every triangle, no flux through the hexagons, and have a total flux of $`\pi `$ through the loop of length twelve that surrounds a hexagon and its six attached triangles. These rules are derived from considering individual hexagons and triangles in isolation, and minimizing the mean-field energy.
While our parent state appears to violate every one of the rules, except the rule for the length-12 loop, the chiral spin liquid is in perfect agreement with these rules! Let us focus on one isolated triangle, with $`|t|=1`$. If there is no flux through the triangle, there are two negative energy states with energy $`1`$. We can put two particles in one state, and one in another, for a total energy of $`3`$. By adding $`\frac{\pi }{2}`$ flux, we have one state at energy $`\sqrt{3}`$ and another state at energy zero. By putting two particles in the lowest energy state, we improve the energy of the system, and introduce a chirality.
Now, consider the triangle coupled to the rest of the system. If the rest of the system strongly scatters the particles in the triangle, it may no longer by appropriate to think of two particles in one state and one in another. One must instead think of each of the two negative energy states of the triangle as each having average occupation of one-and-a-half particles. In that case, it is most advantageous to put no flux through the triangle.
So, if the system to which the triangle is coupled is chiral, so that all triangles in the system have the same flux $`\frac{\pi }{2}`$, then the chiral spin liquid may work. But if the triangle is coupled to states which are not chiral, then the chiral spin liquid is destroyed. One sees this even at the mean field level; a state in which triangles have alternating flux $`\pm \frac{\pi }{2}`$ is significantly worse in energy than our parent state. Similarly, if one introduces sufficiently strong dimerization $`m_{12}`$, one finds that the system is stable against weak $`m_c`$. Within the RG below, we will consider the fluctuations in the masses $`m_i`$ and show that they help stabilize the parent state against the chiral perturbation.
The spin solid of Marston and Zeng can also be obtained from the parent state. In this spin solid state, the idea is to look for dimer coverings which maximize the number of “perfect hexagons”, hexagons on which three of the bonds are covered by dimers. Attached to these hexagons are “defect triangles”, triangles on which no bonds are covered by triangles. Clearly, we wish to increase $`|t_{ij}|`$ on the perfect hexagons, while decreasing it on the defect triangles. This will project onto the mass term $`M_6`$ on the 12-site cell that includes the given perfect hexagon.
The perfect hexagons are then supposed to form a lattice. We can obtain this lattice by taking the triangular lattice of 12-site cells, and placing perfect hexagons inside the 12-site cells on two out of the three sublattices of the triangular lattice. In this case the cells containing perfect hexagons form a honeycomb lattice. This gives rise to a staggered mass state. We produce a mass term $`M_6`$ on two-thirds of the system, so that the Dirac particles feel a net $`M_6`$ at zero momentum, as well as a fluctuating $`M_6`$ at finite momentum.
## IV Naive Mean-Field and Projected Mean-Field
To compare the energies of possible RVB states, including the various mass perturbations of our parent state, we turn to the RVB mean field. When we look for a mean-field solution of $`t_{ij}`$ in equation (2), it is known that the best mean-field is a dimer covering. Still, let us start by looking at results of the naive mean-field calculation, and then later provide a projected mean field calculation.
Let us introduce the Green’s function between sites, $`G_{ij}`$, defined to be the sum over all occupied fermionic states, $`\psi `$, of
$$\psi ^{}(i)\psi (j)$$
(8)
The fermionic energy is then equal to
$$2\underset{<i,j>}{}G_{ij}t_{ji}$$
(9)
where the factor of 2 arises from the presence of up and down species of fermion.
For our parent state, explicit calculation shows that, for nearest neighbors $`i,j`$,
$$|G_{ij}|=0.221383\mathrm{}$$
(10)
From the mean-field condition for equation (2), we find for the parent state that $`t=0.221383J`$, so by equation (4)
$$v_f=(0.0904\mathrm{})J$$
(11)
We find that within the projected mean-field that the parent stable is unstable to all of the massive fluctuations, including the chiral mass fluctuation which will drive the system to a chiral spin liquid.
For infinitesimal perturbations, the different nonchiral masses all provide an equivalent improvement in mean-field energy. To some extent this is due to the approximate low energy symmetry of the Dirac equation to rotating continuously among the different mass terms. However, it is interesting that lattice effects do not break this symmetry. The reason is the discrete lattice symmetry. The change in energy for taking $`M=m_1M_1+m_2M_2+m_3M_3`$ is, for small $`m`$, a quadratic form in $`m_i`$. Let this form be
$$\underset{i,j}{}c_{ij}m_im_j$$
(12)
The coefficients $`c_{11},c_{22},c_{33}`$ in this form must all be the same due to lattice rotation symmetry. Lattice translation symmetry permits one to change the sign of any two of the $`m_i`$, and prevents a nonvanishing $`c_{ij}`$ for $`ij`$. Therefore, for small perturbations the energy gain for introducing a mass must be dependent only on the magnitude of the mass, not the particular mass term used. For larger perturbations, the energy gains may depend on the particular mass term used, and of all the nonchiral mass terms, the system gains the most energy by a mass $`M_{12}`$.
Now let us turn to the projected mean-field. Instead of doing a full Gutzwiller projection, we will use an approximation introduced by Hsu. Within a variational Gutzwiller projection, one minimizes the energy of the Hamiltonian (1). Hsu’s idea at the lowest level of approximation is to note that the Hamiltonian is a sum of terms $`J\stackrel{}{S}_i\stackrel{}{S}_j`$ over different neighbors $`i,j`$, and, when evaluating the expectation value of each of these terms, to perform the projection only on the given sites $`i,j`$. At this level, the variational principle corresponds to minimizing
$$\underset{<i,j>}{}(\stackrel{}{S}_i\stackrel{}{S}_j)\underset{<i,j>}{}6\frac{|G_{ij}|^2}{1+16|G_{ij}|^4}$$
(13)
over all possible $`t_{ij}`$, where $`G_{ij}`$ is determined by $`t_{ij}`$.
For our parent state, we find that $`6\frac{|G_{ij}|^2}{1+16|G_{ij}|^4}=0.2832\mathrm{}`$ By going to the chiral spin liquid, the system improves the ground state energy by roughly $`2.9`$ percent within the projected mean-field approximation. By going to a state with staggered $`\pm \frac{\pi }{2}`$ flux through each triangle, the systems worsens the ground state energy by roughly $`1.4`$ percent. Within this approximation the system is stable against the nonchiral mass perturbations as all the nonchiral mass perturbations worsen the ground state energy at this level of approximation. Again due to discrete lattice symmetry, the energy cost is independent of the particular mass term for small mass, while for larger perturbations, the energy costs differ, and the $`M_{12}`$ perturbation costs the least energy.
## V Field Theory of Fluctuations
Having discussed the naive and projected mean-fields for the problem, we must include fluctuations about the mean-field. To do this, we will use an $`SU(N)`$ generalization of the original problem, such that the projection procedure of Hsu becomes exact. After discussing how to do this in the abstract, we will present the field theory for our specific problem: it will have a number of interacting fields, including fermions, gauge fields, the nonchiral mass terms discussed above, which we will refer to as “pion” fields, and the chiral mass term.
The approximation of Hsu ammounts to minimizing equation (13). By introducing an auxiliary field $`\lambda _{ij}`$ we can “decouple” this sum of functions of $`G_{ij}`$ and instead extremize the function
$$\underset{i,j}{}G_{ij}\lambda _{ji}+f(|\lambda _{ij}|^2)$$
(14)
where $`f(|\lambda |^2)`$ is a Legendre transform of $`\frac{|G|^2}{1+16|G|^4}`$.
Then, we can interchange the order of extremizations, and extremize this quantity over $`t_{ij}`$ before extremizing over $`\lambda _{ij}`$. We find that this is extremized at $`t_{ij}=\lambda _{ij}`$. Then we proceed to extremizing over the one remaining set of variables $`\lambda _{ij}`$. But since $`t_{ij}=\lambda _{ij}`$, we are are equivalently trying to extremize the function
$$\underset{i,j}{}G_{ij}t_{ji}+f(|t_{ij}|^2)$$
(15)
over all $`t_{ij}`$. Note, now, that $`G_{ij}t_{ji}`$ is exactly the kinetic energy of the fermions. So, finally, we are trying to extremize
$$H=\underset{<i,j>}{}(\psi _a^{}(i)t_{ij}\psi _a(j)+h.c.)+\underset{<i,j>}{}f(|t_{ij}|^2)$$
(16)
Returning to the language of functional integrals, we can introduce an $`SU(N)`$ field theory for which the Hsu projection procedure becomes exact. We take a large $`N`$ limit in the number of fermion fields in equation (16), so that $`a=1\mathrm{}N`$. Then, we integrate over all possible $`t_{ij}`$ in that equation, undoing the decoupling procedure above, and rewrite the result in terms of spin operators. We find
$$H=\underset{<i,j>}{}\frac{\stackrel{}{S}_i\stackrel{}{S}_j}{1+16\stackrel{}{S}_i\stackrel{}{S}_j}$$
(17)
We should note a few facts about this procedure. When we demonstrate the equivalence of the large $`N`$ mean-field with the Hsu mean-field, it is the large $`N`$ limit that permits us to ignore fluctuations in $`t_{ij},\lambda _{ij}`$, so that the decouplings amounts exactly to taking a Legendre transforms; at finite $`N`$, the decoupling of an interaction is not exactly a Legendre transform. Further, we used the word “extremize” above with care: in some cases we maximize while in other cases we minimize, as in some places the function $`\frac{|G|^2}{1+16|G|^4}`$ has positive curvature while in other cases it has negative curvature. This does not provide any formal problems when performing the decoupling at the level of functional integrals, so long as we correctly choose the integration contour of $`\lambda _{ij}`$.
The fractional operator, $`_{<i,j>}\frac{\stackrel{}{S}_i\stackrel{}{S}_j}{1+16\stackrel{}{S}_i\stackrel{}{S}_j}`$ in equation (17) may be interpreted as a formal power series, so that it includes operators of the form $`(\stackrel{}{S}_i\stackrel{}{S}_j)^k`$ for all $`k`$. At $`N=2`$, this operator is equivalent to $`\stackrel{}{S}_i\stackrel{}{S}_j`$, up to a constant factor.
Finally, the above procedure is similar to the technique of introducing biquadratic interactions, $`(\stackrel{}{S}_i\stackrel{}{S}_j)^2`$, into the Hamiltonian to stabilize RVB states against dimerization. We simply prefer the above Hamiltonian as it reproduces exactly our desired mean-field theory.
Having defined a large $`N`$ theory with no fluctuations in the decoupling fields, we next add in fluctuations. Formally, this can be handled by a $`1/N`$ expansion. We will directly write the field theory at $`N=1`$ without including explicit factors of $`N`$.
The theory includes several modes. There is the Dirac fermion field, $`\psi (x)`$. This is coupled to a fluctuating $`U(1)`$ gauge field, $`A^\mu (x)`$, $`\mu =t,x,y`$. By considering other fluctuations in $`t_{ij}`$ we will also obtain a fluctuating chiral mass field that we will refer to as $`\sigma (x)`$, and a triplet of fluctuating nonchiral mass terms that we will group into one “pion” field $`\pi ^a(x)`$, $`a=1,2,3`$. If the pion field acquires an expectation value, then the fermions will acquire a mass $`m_a=\pi ^a`$.
For the field theory, we will suppress the velocity $`v_f`$. At the level of the bare action the pi, sigma, and gauge fields can have different velocities from the fermi fields. However, the greatest contribution to the effective action of the bosonic fields arises from integrating over the relativistic fermions, so that at low energies the velocity of the bosonic fields must be roughly equal to that of the fermionic fields.
The Lagrangian $`L`$ we will take is
$$L=d^3xL_f+L_A+L_M$$
(18)
where
$$L_f=\overline{\psi }_{u,d}(x)\left(\gamma ^\mu (A_\mu +i_\mu )+\gamma _0M_a\pi ^a+\gamma _0M_c\sigma \right)\psi _{u,d}(x)$$
(19)
$$L_A=\frac{1}{4\mathrm{\Lambda }g_a^2}F_{\mu \nu }F^{\mu \nu }$$
(20)
$$L_M=\frac{1}{2\mathrm{\Lambda }g_\sigma ^2}\sigma (x)\left(_\mu ^2+m_\sigma ^2\right)\sigma (x)+\frac{1}{2\mathrm{\Lambda }g_\pi ^2}\pi ^a(x)\left(_\mu ^2+(m_\pi ^2)^{ab}\right)\pi ^b(x)$$
(21)
We have written the mass for the pion field as a matrix. However, following the arguments for equation (12), the masses of the different pion modes must be the same. It is only after condensation of the pion field, breaking the lattice symmetry, that the masses can differ.
We have inserted factors of $`\mathrm{\Lambda }`$, representing a lattice cutoff scale, into the action to make the coupling constants dimensionless. We have chosen to scale the bosonic fields so that all coupling constant dependence appears in the action $`L_A`$ and $`L_M`$, not $`L_f`$.
In addition to the terms we have written, there must be a quartic interaction term for the $`\pi `$ and $`\sigma `$ fields. This term is necessary to stabilize the action if the system spontaneously breaks a symmetry and has either $`m_\pi <0`$ or $`m_\sigma <0`$. Since we will be initially starting the renormalization group of the next section with both such masses positive, we can temporarily ignore the quartic term at high energy under the assumption that this term is small. If the system acquires an expectation value for the $`\pi `$ fields, the quartic term will break the continuous symmetry down to the lattice symmetry, and give a small mass $`m_\pi ^{}`$ for the approximate Goldstone modes. There will also be cubic terms that give a mass to these modes.
Another interesting term we have left out is $`\overline{\psi }\gamma _\mu \gamma _\nu F^{\mu \nu }\psi `$, which can be added to change the $`g`$-factor of the Dirac fermions. In the absence of external fields, there are two degenerate states of the Dirac equation at each energy. For physical electrons, this reflects a spin degeneracy. For the spinons we consider, which already have a definite spin, this degeneracy instead reflects a chirality degeneracy, and we will refer to it as such. If the system has an odd number of sites, and hence an unpaired spinon, not only does the system have a net spinon, but it also has a net chirality, which can be taken to be positive or negative. Generically the $`g`$-factor will be non-zero.
## VI Renormalization Group
We will consider the one-loop RG from the field theory. We will see that it is indeed possible for fluctuations to lead to a condensation of the pion field.
For the gauge, pion, and sigma fields we will use a simple mode elimination RG, with a cutoff $`\mathrm{\Lambda }`$. For the fermion fields, we will introduce a set of massive regulator fields with masses of order $`\mathrm{\Lambda }`$ and reduce the regulator mass. The choice of the particular mass terms for the regulator fields represents a lattice breaking of the Goldstone symmetry. It is possible to preserve the needed lattice symmetry of equation (12) by introducing seven regulator fields. Four are ghost fields with masses proportional to $`M_1+M_2+M_3,M_1M_2M_3,M_1+M_2M_3,M_1M_2+M_3`$ and the other three are not ghosts and have masses proportional to $`M_1,M_2,M_3`$.
Initially the theory will have a cutoff $`\mathrm{\Lambda }_0`$, defining the lattice scale. As we renormalize, we lower the cutoff $`\mathrm{\Lambda }`$, and rescale all distances and fields to keep $`\mathrm{\Lambda }`$ fixed.
We must take into account self-energy corrections to the fermions from interactions with the bosonic fields, vertex corrections, and self-energy corrections to the bosonic fields from vacuum polarization bubbles. If we were to take into account only the self-energy corrections to the bosonic fields, and not the vertex and fermionic self-energy terms, we would find that we are considering just the mean-field theory in the bosonic fields.
We find the following RG equations:
$$\frac{d\mathrm{ln}g_A}{d\mathrm{ln}(\mathrm{\Lambda }_0/\mathrm{\Lambda })}=12\frac{g_A^2}{3}$$
(22)
$$\frac{d\mathrm{ln}g_\pi }{d\mathrm{ln}(\mathrm{\Lambda }_0/\mathrm{\Lambda })}=1+\frac{3g_A^2+g_\pi ^2g_\sigma ^2}{2\pi ^2}+\mathrm{vacuum}\mathrm{polarization}$$
(23)
$$\frac{d\mathrm{ln}g_\sigma }{d\mathrm{ln}(\mathrm{\Lambda }_0/\mathrm{\Lambda })}=1+\frac{3g_A^23g_\pi ^2g_\sigma ^2}{2\pi ^2}+\mathrm{vacuum}\mathrm{polarization}$$
(24)
$$\frac{dm_\pi ^{ab}}{d\mathrm{ln}(\mathrm{\Lambda }_0/\mathrm{\Lambda })}=2+\mathrm{vacuum}\mathrm{polarization}$$
(25)
$$\frac{dm_\sigma }{d\mathrm{ln}(\mathrm{\Lambda }_0/\mathrm{\Lambda })}=2+\mathrm{vacuum}\mathrm{polarization}$$
(26)
We have avoided explicitly writing the vacuum polarization contributions to the sigma and pion fields. The vacuum polarization contribution to the mass is regularization dependent, while the vacuum polarization contribution to the coupling constant is ultraviolet convergent and is dominated by the infrared contribution.
Fluctuations in the gauge field increase the coupling constants for the pion and sigma fields. This reflects the binding force due to the gauge field between charged spinons, and the resulting tendency to break chiral symmetry. Further, we see that the coupling constant for the pion field increases more rapidly than that for the sigma field, reflecting the destabilization of the chiral state by fluctuations in the pion field.
Thus, we see from the renormalization group that there is a range of bare parameters such that the theory will condense the pion field, producing a nonchiral mass term for the fermions, even though at the mean-field level the theory would rather produce a nonchiral mass term for the fermions.
In the next section we will consider the low energy action after condensation. We will first discuss the mass term for the fermions that appears.
Unfortunately, it is beyond our ability to calculate the bare parameters in the field theory with any precision, and so the mass of the fermion field is not something we can compute. Let us instead take the mass of the fermion as an input from numerics, and use that to check for consistency of our theory. Extrapolating finite size results from systems of up to 36 sites, one finds that the system has a gap to triplet excitations which is of order $`J/20`$ or less. While the gap is decreasing even at the largest sizes, it appears to be bounded below by roughly $`J/40`$. Assuming that the triplet excitations are made of pairs of spinons, we can estimate the spinon gap as being half the triplet gap. Further evidence for the spinon gap being roughly half the triplet gap comes from odd-even studies of the energy dependence on N. The fermion mass is half the spinon gap, or one quarter the triplet gap.
Using this estimated spinon gap, and the calculated velocity of our Dirac particles from RVB theory, we can obtain the correlation length of the Dirac particles. Taking the estimate of $`J/20`$ for the triplet gap, we find that the correlation length is roughly $`8`$ of our twelve-site unit cells, large enough to include the $`N=36`$ numerical studies.
We can also estimate the strength of dimerization at the mean field level, by asking how large a change in $`t_{ij}`$ is needed to produce the desired mass term. Assuming that the dimerization is provided by a perturbation of the form $`M_{12}`$, one finds that the $`t_{ij}`$ 12-site loops are increased by approximately $`3.5`$ percent, while the other bonds are decreased by $`3.5`$ percent. This is a relatively small amount of dimerization, and we expect that only after a significant increase in system size will numerical studies be able to detect this directly from a dimer-dimer correlation function.
## VII Low Energy Modes
The remaining low energy modes after the pion field condenses are the Goldstone excitations of the pion field and the gauge excitations, which we will argue provide the low energy singlet modes seen in numerics. While numerical calculations have only probed systems up to $`N=36`$ sites, which is relatively small considering that we take a unit cell of 12 sites, experiment also reveals a quadratic low temperature specific heat. This specific heat suggests that a bosonic mode with a linear density of states survives to much larger scales, while the insensitivity of the specific heat to weak magnetic fields suggests that the mode is still a singlet. In this section we will first discuss the nature of the low energy modes and then the ultimate fate of our pion and gauge excitations at large distances, including a gap for the pion from lattice effects as well as a confining phase for the gauge fields.
Once the pion field condenses, the system is left with two pseudo-Goldstone pion modes as well as gauge modes. The gauge action is
$$L_A=\frac{1}{4\mathrm{\Lambda }g_A^2}F_{\mu \nu }F^{\mu \nu }$$
(27)
with
$$g_A^2m$$
(28)
where $`m=|\pi |`$ is the fermion mass. The pion action will be a sigma model. If we change the normalization on the pion field so that $`|\pi |=1`$, we get the model
$$L_M=\frac{1}{2g^2}(_\mu \pi ^a)^2+\frac{(m_\pi ^{})^2}{2g^2}\pi ^1\pi ^2\pi ^3$$
(29)
where the coupling constant $`g^2`$ is proportional to $`m^1`$ and the mass $`m_\pi ^{}`$ represents the breaking of continuous symmetry by lattice effects.
We have chosen the mass term for the pion to cause the pion to prefer to condense in a way that gives rise to a mass $`M_{12}`$. In the projected mean-field calculation above, we considered states invariant under the rotational symmetry, so that $`|m_1|=|m_2|=|m_3|`$. This provided two inequivalent perturbations. In one, we increased $`|t_{ij}|`$ on hexagons and triangles; in the other we increased $`|t_{ij}|`$ on a loop of length twelve. While at the mean-field level symmetry breaking does not occur, from the projected mean-field calculation we can still argue that the preferred symmetry breaking pattern would be given by a mass matrix $`M_{12}`$. Other patterns are of course possible, and comparison with numerics provides some evidence that a staggered mass is also a possibility, at least for small systems; in the conclusion we discuss possibilities of numerically testing the preferred mass pattern.
In the continuum field theory, the pion mass $`(m_\pi ^2)^{ab}`$ is ultraviolet divergent. However, lattice symmetry forces the masses of the pion modes to be the same before condensation. In order to use the continuum theory to estimate $`m_\pi ^{}`$ after condensation, we need to turn to the cubic interaction terms in $`\pi `$. These are
$$d^3xg_3\pi _1\pi _2\pi _3$$
(30)
with a cubic coupling constant $`g_3`$ that is generically of order unity. Inserting an expectation value of $`\pi `$ of order $`m`$, we obtain a quadratic term in $`\pi `$. Including this quadratic term in $`(m_\pi ^2)^{ab}`$, this will cause the masses of the different pion modes to differ by order $`m`$ so that $`m_\pi ^{}`$ will be of order $`m`$.
However, we can obtain a better estimate of the mass difference numerically from the the projected mean-field calculation of energy; while we did not obtain pion condensation at the mean-field level, for a given expectation value of the pion field, we can use the difference in mean-field energies for various continuous rotations of the pion field to obtain an estimate of the pion gap. Using the numerical estimate for the triplet gap, and hence the estimate for the fermionic mass $`m`$, we have calculated the projected mean-field energies for taking $`M=mM_{12}`$ and $`M=mM_6`$. The difference in energies is $`0.000396J`$ per 12-site cell, so that $`m_\pi ^{}\sqrt{0.000396mJ}`$. This is small enough that we can ignore this mass for most purposes. Evidently, the cubic coupling constant $`g_3`$ is very small.
While the pion is already gapped by lattice effects, instantons will gap the gauge field, leading to confinement of the spinons. The gauge field describes compact QED in 2+1 dimensions, which is confining for all $`g_A`$. The gauge coupling is proportional to $`m`$, so that the action for an instanton will be of order
$$S\frac{\mathrm{\Lambda }_0}{m}$$
(31)
The instanton density is proportional to $`e^S`$. In the weak coupling limit, the instantons lead to a gap for the gauge field of order $`\mathrm{\Lambda }_0e^S`$. As this is exponentially small, we can ignore the gap in the gauge field.
## VIII Finite Size Systems and Chern Numbers
In this section we will consider some effects of finite size systems to begin comparison with numerics. First we will consider some complications in defining the parent state on systems with an odd number of sites, which force the system to have some net flux. Then we will show how this leads to a degeneracy in the spectrum and nonvanishing Chern numbers for the states under twist in the spin boundary conditions. Finally, we will discuss some effects of finite size for even size systems.
One of the most interesting results found in numerical studies of odd size systems is a non-vanishing Chern number for the ground state of $`\pm 1`$. This is a quantity that provides an analogue for a spin system of the quantum Hall effect. Since the Hamiltonian of equation (1) does not explicitly break time-reversal symmetry, a non-vanishing Chern number requires a spontaneous breaking of time-reversal symmetry. However, the spontaneous breaking of time-reversal symmetry is not enough, as other spin systems that break this symmetry have vanishing Chern number; the kagomé antiferromagnet may be the first Hamiltonian with time-reversal and parity symmetry where a non-vanishing Chern number has been observed. In order to understand the appearance of the Chern number, we must first understand how to form the parent state on an odd size system.
The systems studied numerically have periodic boundary conditions, so that they live on a torus. On system defined on a torus the net flux penetrating the surface must be an integer multiple of $`2\pi `$. One can also add solenoid fluxes, $`\theta _1,\theta _2`$, defining the phase that the spinon acquires when traversing a topologically nontrivial closed loop around the torus. For simplicity, we will use coordinates on the torus which range from 0 to $`2\pi `$ in both directions, although in actuality for the kagomé lattice the systems considered are not square.
The parent state has $`\pi `$ flux through each hexagon. On a system with $`N`$ sites, there are $`N/3`$ hexagons, and so on a system with an odd $`N`$, one would like to have a net flux through the system that is an odd multiple of $`\pi `$. This is not possible, and so the system must have some additional flux so that the total is a multiple of $`2\pi `$. For example, on a system with $`N=27`$, there are 9 hexagons, so the system can put $`\pi \pm \frac{\pi }{9}`$ flux through each so that the total flux is either $`8\pi `$ or $`10\pi `$. The system then must become chiral and break time reversal symmetry since it cannot construct the parent state.
A qualitative way of describing this effect is to say that for a system with an odd number of sites, there is an unpaired spinon, which has a chirality. The spinon then couples to the gauge field and produces a flux.
Given this net flux through the system, let us consider the effective Dirac equation for the spinons. The results we get for the Chern number do not rely on the Dirac description, and can be derived directly from the lattice model; we feel that the Dirac method is more elegant and gives more physical insight.
The Dirac particles feel the extra flux that has been added, and so the spinons move in a magnetic field, such that the net flux the spinons feel is exactly $`\pm \pi `$. Again, there seems to be a contradiction, since it is not possible for the system to have a net flux of $`\pm \pi `$ flux through the torus. The answer to the contradiction is that the Dirac particles have an extra chirality index. So, in addition to including solenoid fluxes for the Dirac equation, the Dirac particles can change chirality when completing a loop around the system. Let us then generalize the solenoid flux to a pair of 4-by-4 matrices $`U_1`$, $`U_2`$, describing the change in the wavefunction when the particle completes a loop.
Then, when the particle traverses a loop around the torus from $`(0,0)`$ to $`(0,2\pi )`$ to $`(2\pi ,2\pi )`$ to $`(2\pi ,0)`$ to $`(0,0)`$, the wavefunction gets multiplied by
$$U_1U_2U_1^{}U_2^{}$$
(32)
where the minus sign is from the magnetic flux through the torus. Since equation (32) must be equal to 1, we find that $`U_1`$,$`U_2`$ necessarily anti-commute.
One may regard the matrices $`U_1,U_2`$ as arising from a non-Abelian gauge field connecting opposite chiralities of the spinons. The commutator of the matrices represents an additional flux of $`\pi `$ from the non-Abelian field, giving a total flux of $`2\pi `$ on the torus. The extra $`\pi `$ flux from the non-Abelian field is the flux that arises from having an odd number of hexagons on the lattice, so that when the particle completes the given loop around the lattice it has enclosed an odd number of $`\pi `$ fluxes. One sees that the non-Abelian flux is localized at a point, although one must be careful that this localization at a point does not imply a breaking of translational symmetry.
The addition of matrices $`U_1,U_2`$ is natural from the lattice point of view. The unit cell which includes both chiralities of Dirac particles is 12 sites, while the smallest unit cell possible for the parent state is 6 sites. Since there is no way to cover an odd size lattice with 6 or 12 site unit cells, something must scatter between chiralities, as when the particle completes a loop it has changed between chiralities.
To give a very simple example of this, consider a one-dimensional ring with an odd number of sites. The natural unit cell for the a one-dimensional chain is two-sites, to include both Dirac points. If the particle moves around an odd-length ring, two sites at a time, it must return to the starting point displaced by half a unit cell.
To give a slightly more complicated example, consider the $`\pi `$-flux phase of the square lattice for a system of 9 sites, shown in figure 5. Solid lines represent bonds within the cell of 9 sites, while dotted lines represent bonds to provide toroidal boundary conditions. There are 9 squares in the system, and so there will be $`\pi \pm \frac{\pi }{9}`$ flux through each square. 4 of the squares lie within the cell and are labeled 1-4, another 4 lie to the sides and are labeled 5-8, while the 9th square lies in the corner. The natural unit cell for the Dirac particles is 4 sites, so when a particle completes a loop around the torus it is displaced by half a unit cell. The wavefunction is multiplied by a matrix $`U_1`$ for a loop in the $`\widehat{x}`$-direction and a matrix $`U_2`$ for a loop in the $`\widehat{y}`$-direction. Precisely due to the odd number of squares, one finds that the matrices $`U_1,U_2`$ anti-commute. It is natural to think of the non-Abelian flux as arising from the $`\pi `$-flux through the 9th square, on the corners.
Returning to the kagomé lattice, let us now look at the wavefunctions of the Dirac equation with this magnetic flux. It is convenient to find the wavefunctions by enlarging the torus by a factor of two in each direction, as shown in figure 5. The $`+`$ and $``$ symbols in the figure denote the chirality of the particle in each quadrant. When the particle completes a circuit on the original torus, it changes chirality, and hence moves into a different quadrant of the enlarged torus, while picking up a phase.
The net flux on the enlarged torus is equal to 4 times the flux on the original torus, or $`4\pi `$. One might imagine that there will be extra sources of $`\pi `$-flux on the enlarged torus at the points where the quadrants meet. However, since the $`\pi `$ non-Abelian flux is purely a result of an odd number of hexagons on the original torus, we can drop the extra sources of flux on the enlarged torus, and we are left with an explicitly translationally invariant problem of a Dirac particle moving in a constant magnetic field.
On the enlarged torus, the Dirac equation becomes a two-component equation
$$\left(iE\sigma _z+\sigma _x(A_x+i_x)+\sigma _y(A_y+i_y)\right)\psi =0$$
(33)
Taking the square we find
$$E^2\psi =\left((i_xA_x)^2+(i_yA_y)^2+\sigma _zB\right)\psi $$
(34)
This is the well known equation for a Dirac particle in a magnetic field, and has Landau levels.
Since the net flux through the enlarged torus is equal to $`4\pi `$, there are two-states in each Landau level for each $`\sigma _z`$, hence four states for each Landau level in total. As we are dealing with a two component equation, only one sign of $`E`$ is allowed in equation (34) for given $`\sigma _z`$.
Therefore, the energy levels on the enlarged torus are doubly degenerate. However, the enlarged torus has an unphysical degree of freedom: opposite quadrants describe the state of the particle on the original torus. So, the energy levels on the original torus are only singly degenerate and the spectrum is discrete with one level at zero energy. This is the relativistic generalization of Landau levels. For an odd size system, all Landau levels below the zero energy are occupied, and hence filled, for both spin up and spin down particles, while the zero energy level is occupied only by one unpaired spinon.
Numerically, the ground state of the many-body system has been seen to have an extra degeneracy factor of two, beyond the trivial spin degeneracy. This is a consequence of the spontaneous generation of the magnetic field, so that the system can pick either sign for the field.
An interesting way of viewing the Landau levels is that we have $`\pi `$ Abelian flux, implying that there are $`1/2`$ states per Landau level. Multiplying the $`1/2`$ by a factor of two for chirality degeneracy, we get one state per Landau level.
We can now introduce the Chern number of the system, which characterizes a transverse response of spin currents. Let us adjust the boundary conditions of the system so that
$$S^\pm (x,y)=e^{\pm i\varphi _1}S^\pm (x+2\pi ,y))$$
(35)
$$S^\pm (x,y)=e^{\pm i\varphi _2}S^\pm (x,y+2\pi )$$
(36)
where $`\varphi _1,\varphi _2`$ are angles.
If the ground state is a wavefunction $`\mathrm{\Psi }`$, then the Chern number is defined as the integral
$$\frac{1}{2\pi }\frac{\mathrm{\Psi }}{\varphi _1}|\frac{\mathrm{\Psi }}{\varphi _2}𝑑\varphi _1𝑑\varphi _2$$
(37)
This number is quantized, and non-vanishing only for complex states. Since the Hamiltonian does not break time-reversal symmetry, complex conjugate states are degenerate with opposite Chern numbers.
In the presence of these boundary conditions, the spinon boundary conditions, with additional self-generated fluxes $`\theta _1,\theta _2`$, become
$$\psi _u^{}(x,y)=e^{i\frac{\rho _1^u}{2}}\psi _u^{}(x+2\pi ,y)$$
(38)
$$\psi _d^{}(x,y)=e^{i\frac{\rho _1^d}{2}}\psi _d^{}(x+2\pi ,y)$$
(39)
where
$$\rho ^{u,d}=(2\theta _1\pm \varphi _1)$$
(40)
and similarly for the other direction. The $`2\pi `$ periodicity in $`\varphi `$ is not obvious from equation (38,39), but the ability of the system to adjust $`\theta `$ produces the desired periodicity.
To give a simple example of how a system can adjust $`\theta `$, consider a system of 4 sites on a ring. In the absence of a twist in boundary conditions $`\varphi `$, the system places $`\theta =\pi `$ flux through the ring. As $`\varphi `$ increases, $`\theta `$ remains equal to $`\pi `$, and the mean-field energy of the system gradually increases until $`\varphi =\pi `$. At this point, $`\theta `$ jumps to $`0`$, and the mean-field energy begins to decrease for increasing $`\varphi `$. So, as $`\varphi `$ varies from $`0`$ to $`2\pi `$ and $`\theta `$ jumps as described above, we find that $`\frac{\rho }{2}=\frac{(2\theta +\varphi )}{2}`$ varies from $`\pi `$ to $`\frac{3\pi }{2}`$ to $`\frac{\pi }{2}`$ to $`\pi `$. In order for $`\theta `$ to jump like this, the spinon states with $`\frac{\rho }{2}=\pi \pm \frac{\pi }{2}`$ must be degenerate.
Now, we can look at the Chern number of the system, assuming non-interacting spinons. It then amounts to a Chern number calculation of the fermionic wavefunctions. Assuming non-interacting spinons, we can get the change in $`\mathrm{\Psi }`$ in equation (37) from the change in the spinon wavefunctions.
While in general the wavefunction gets multiplied by a matrix on moving around the original torus, only the $`U(1)`$ part of this matrix adjusts in response to changes in $`\varphi `$. The $`U(1)`$ part of the matrix is just the angle $`\theta `$. Carrying out the calculation on the enlarged torus, we find that the boundary conditions become
$$\psi _u^{}(x,y)=e^{i\rho _1^u}\psi _u^{}(x+4\pi ,y)$$
(41)
and similarly for down spinons.
So, we wish to compute
$$\frac{1}{2\pi }\frac{\psi }{\rho _1^{u,d})}|\frac{\psi }{\rho _2^{u,d})}\frac{\rho _1^{u,d}}{\varphi _1}\frac{\rho _2^{u,d}}{\varphi _2}𝑑\varphi _1𝑑\varphi _2$$
(42)
summed over all spinon wavefunctions $`\psi `$.
In equation (41) the periodicity in $`\varphi `$ seems obvious even without $`\theta `$, as on the enlarged torus the periodicity of the spinon wavefunctions in response to a twist in boundary conditions is halved in both directions. However, we have introduced the degeneracy of two on the enlarged torus representing the fact that on the original torus the wavefunctions are periodic in $`(\rho _1,\rho _2)`$ with periods $`(0,2\pi )`$ and $`(\pi ,\pi )`$ but not with period $`(0,\pi )`$, and as a result only half the possible wavefunctions on the enlarged torus are physical.
In order to keep the wavefunction in the physical sector, $`\theta `$ must jump discontinuously by $`\pi `$ as $`\varphi `$ changes, and as a result for a given spinon state we only integrate equation (42) over half the torus of possible phases $`(\rho _1,\rho _2)`$. The fact of integrating over half the torus, or equivalently the fact that the Landau levels contain one physical and one unphysical wavefunction, does not prevent a defining of the Chern number for the spinon wavefunctions. When $`\theta `$ jumps it connects two degenerate states, and so the contribution of equation (42) to equation (37) must still be quantized as an integer which we can still refer to as a Chern number for the spinon.
Within the lattice formulation there are no further conceptual problems and we must simply compute the integrals, but within the continuum Dirac equation we must account for the negative energy sea. The correct understanding of this was found by Haldane.
One must add a massive regulator field, and compare the difference in Chern number between the massless and massive fields. The massive field has the same Landau level spectrum, but no zero mode. So, the difference in Chern numbers is due to the zero mode, which sits in the lowest Landau level. There are two states in the lowest Landau level, one physical and one unphysical. It is the physical state that carries the Chern number of $`\pm 1`$, giving the ground state of the spin system a net Chern number of $`\pm 1`$, as seen numerically.
We expect that the low-lying states will continue to have an odd Chern number, in agreement with numerical results. If a particle-hole pair is excited within the Dirac band near the Dirac point, the Chern number will not change. If the particle is excited from the band edge, the Chern number can change by $`\pm 2`$. Only if a particle is excited from the flat band to the Dirac band can the Chern number change by $`\pm 1`$, giving rise to an even Chern number. However, these states will be much higher in energy. One can also consider excited states with net Abelian flux be equal to $`3\pi ,5\pi ,\mathrm{}`$
It is very interesting to think about these possible excited states which may have more than $`\pi `$ flux for the Dirac particles. The two-component particles carry a $`\sigma _z`$ index, which will couple to the magnetic field. If a large field is induced, a number of spinons of the same $`\sigma _z`$ will be produced in the zero mode, so that the total number of spinons in the zero mode is odd. For one spinon we had one filled Landau level, with one particle. With several spinons one might be able to construct fractional Hall states of spinons.
We have argued that in the thermodynamic limit the system will acquire a mass. On an odd size lattice, the mass term must change sign somewhere, as the lattice cannot be tiled with 12-site unit cells. At the domain wall where the mass changes sign, one expects to trap a midgap state, so there still should be a zero mode, even with mass. This may permit the nonvanishing Chern number to survive.
Returning to even system size, let us consider the size dependence of the triplet gap. The energy of the spinon is $`E=\sqrt{(v_fk)^2+m^2}`$. In the absence of a solenoid flux, the smallest $`k`$ would be equal to zero, but by creating a solenoid flux the energy can be improved and the smallest $`k`$ will be of order the inverse linear dimension of the system, or $`N^{1/2}`$. As a result, the triplet gap is decreasing with system size, in agreement with numerics. By twisting the spin boundary conditions one may be able to reduce the gap to $`S_z=\pm 1`$ excitations. It would be interesting to look for this effect.
Further, in the presence of these solenoid fluxes, other fermionic states with $`k0`$ will become approximately degenerate with the $`k=0`$ fermionic state. This means that a spatially varying mass term which scatters between $`k`$ states can open a gap just as well as the spatially constant mass term can. This will be important when we consider the low energy Goldstone excitations on finite size systems, below.
## IX Goldstone Modes, Tower States, and Numerics
After breaking a symmetry, and giving mass to the fermions, the system is left with low energy pion and gauge modes. Above, we argued that the gap for these modes is too small to be seen in numerical calculations. In this section, we will treat these modes as gapless and discuss the energy spectrum that results for finite size systems to compare with numerical calculations.
It is known that breaking a continuous symmetry gives rise to two kinds of low energy modes. First, there are the Goldstone modes with non-zero wavevector $`k`$. In the case of our pion and gauge modes, the energy is then proportional to $`k`$. For a $`2+1`$ dimensional system with $`N`$ sites, the lowest $`k`$ is proportional to $`N^{1/2}`$ and so the lowest Goldstone excitation has energy proportional to $`N^{1/2}`$.
Second, there is the “tower” of $`k=0`$ modes. These correspond to global rotations of the entire system, and have an energy proportional to $`N^1`$.
Numerical diagonalization of the triangular lattice Heisenberg anti-ferromagnet, which has Néel order, shows very clearly the distinction between the tower of states, and the $`k0`$ states (spin waves). However, no such distinction is found in the kagomé lattice, no separation between low energy modes of energy $`N^1`$ and $`N^{1/2}`$. Within our model, this is to be expected for $`N=36`$. Since the effective action of the pion and gauge fields arises from integrating out the fermions, this action must be approximately relativistic, with the same velocity as the fermions. So, even without explicit calculation, we can obtain the energy of the lowest $`k0`$ mode directly from the velocity appearing in Dirac equation. For the largest numerical diagonalizations, systems with $`N=36`$ total sites or 3 of our 12-site unit cells, this energy turns out to be of order the triplet gap, and so this Goldstone mode is too high in energy to appear in the continuum of low energy singlets.
The only states that will be observed in numerics are states in the tower. We can obtain the energy of these states from equation (29), assuming that $`\pi (x)`$ is constant. Then we get, assuming small $`m_\pi ^{}`$,
$$L=\frac{N\mathrm{\Lambda }_0^2}{g^2}(_t\pi ^a)^2$$
(43)
where $`N\mathrm{\Lambda }_0^2`$ is the area of the system. The states of this will be spherical harmonics, perturbed by mass term. With $`g^2m^1`$, these states will have an energy proportional to
$$\frac{\mathrm{\Lambda }_0^2}{Nm}$$
(44)
and so for sufficiently large $`N`$ will be below the triplet gap. A more precise knowledge of the prefactors will be needed to determine whether this is low enough to correspond to the low energy modes seen numerically.
There will also be “tower” states for the gauge field, which correspond to different solenoid fluxes through the system. The energy for these can be obtained from equation (27) assuming that $`A_\mu `$ is constant over the sample. To get the energy for these we have to remember that the gauge group is compact, and realize that $`F^{\mu \nu }`$ is derived from a set of $`U(1)`$ matrices with a lattice length $`\mathrm{\Lambda }_0`$. Then, the energy of the gauge states is proportional to
$$\frac{m}{N}$$
(45)
which is definitely below the triplet gap and certainly small enough to be the origin of some of the low energy modes in numerics.
There is one puzzle involved in the tower states. It was observed numerically that, on 36 site samples, the energy of the lowest energy state of the system at given total momentum did not vary appreciably across the Brillouin zone. This may seem to be in contradiction to the hypothesis that the low energy states come from the tower. The resolution of this may lie in realizing that there are only 4 inequivalent points in the Brillouin zone of the system. One of these is at $`k=0`$, while including a constant nonvanishing mass term $`M_{12}`$ also breaks translational symmetry on the kagomé lattice, and can give one more point in the Brillouin zone. To obtain the last points in the Brillouin zone, one must remember that for small systems there can be other symmetry breaking patterns with staggered mass, breaking translational symmetry in different ways, as discussed in the section on finite size effects.
This would imply that for sufficiently large system sizes one would find a more significant variation in the energies across the Brillouin zone, as only some of the states could be obtained from the tower. One would also find for 36 site systems that the spinon solenoid fluxes would change under a twist in spin boundary conditions, and so the fermionic states at different $`k`$ would lose their degeneracy, leading to a change in energy of some of the $`k0`$ states when varying boundary conditions.
Another possibility is that, even for infinite systems, short-distance effects lead to the production of a staggered-mass pattern, in the style of the “perfect hexagon” state discussed above. The staggered mass pattern gives rises to an enlarged unit cell, and makes it possible to get low energy states at the other points in the Brillouin zone.
## X Further Comparison With Numerics
In addition to the existence of the low energy states, we make further comparison with numerics for the many-body density of states and the dimer-dimer correlation functions.
Assuming the existence of a low-energy bosonic mode with linear dispersion relation, so that the single-particle density of states scales linearly with energy, one would expect the many-body density of states at energy $`E`$ to scale for large systems as an exponential of $`E^{2/3}`$. The quadratic behavior of experimentally measured specific heat is in agreement with this.
Since we have argued that the low-energy states in numerical calculations are largely “tower” states, it is impossible to extract the density of states in a large system from the finite size density of states. The true exponential growth can only be seen when the $`k0`$ modes become important.
In this regard, the numerically measured quadratic many-body density of states at very low energies does not say anything about the dispersion of the Goldstone modes. Instead, it is a reflection of the fact that if a finite number of “tower” modes are excited then the many-body density of states is a power law.
For the dimer-dimer calculations we can make a more direct comparison with a numerical calculation of these correlations on a 36 site system. In the thermodynamic limit, the dimer-dimer correlation function should be long-ranged, reflecting the existence of a non-zero $`m_{12}`$.
However, there is also a short-range fermionic contribution to the dimer-dimer correlation functions, and for 36 site systems, so that the system size is smaller than the correlation area of the fermions, we can ignore the effect of a non-zero $`m_{12}`$ on the dimer-dimer correlation function and directly study the correlation functions of massless fermions.
If we are interested in a dimer-dimer correlation function
$$C_{(i,j)(k,l)}=(\stackrel{}{S}_i\stackrel{}{S}_j)(\stackrel{}{S}_k\stackrel{}{S}_l)(\stackrel{}{S}_i\stackrel{}{S}_j)(\stackrel{}{S}_k\stackrel{}{S}_l)$$
(46)
we can express this, under the assumption of weak spinon interaction, directly in terms of the spinon Green’s functions. Writing each spin operator in terms of spinons and considering various contractions we obtain
$$C_{(i,j)(k,l)}=4\mathrm{R}\mathrm{e}(G_{ij}G_{jk}G_{kl}G_{li})4\mathrm{R}\mathrm{e}(G_{ij}G_{jl}G_{lk}G_{ki})\mathrm{Re}(G_{ik}G_{kj}G_{jl}G_{li})+\frac{1}{2}|G_{ik}|^2|G_{jl}|^2+\frac{1}{2}|G_{ij}|^2|G_{jl}|^2$$
(47)
To obtain quantitatively accurate answers, we must include the effects of projection within an approximation like that used above, projecting on site $`i,j,k,l`$, requiring that there be one fermion on each of these four sites. We have done this at 3 different levels of approximation.
At the lowest level, we have noted that the wavefunction before projection will have one fermion on each of these sites is roughly the product of the probability that it will have one fermion on each of $`i,j`$ by the probability that it will have one fermion on each of $`k,l`$. Given that $`|G_{ij}|=0.221383\mathrm{}`$ for neighboring $`i,j`$, we should replace equation (47) by
$$C_{(i,j)(k,l)}=\kappa \left(4\mathrm{R}\mathrm{e}(G_{ij}G_{jk}G_{kl}G_{li})4\mathrm{R}\mathrm{e}(G_{ij}G_{jl}G_{lk}G_{ki})\mathrm{Re}(G_{ik}G_{kj}G_{jl}G_{li})+\frac{1}{2}|G_{ik}|^2|G_{jl}|^2+\frac{1}{2}|G_{ij}|^2|G_{jl}|^2\right)$$
(48)
where
$$\kappa =(1.9264)^2$$
(49)
is the desired factor of $`(\frac{4}{1+16|G|^4})^2`$.
At a more refined level, we have carried out the computation projecting on all four sites exactly. At the third level of approximation we have started to project out onto additional sites as well. We have calculated the correlation functions in this approximation, in the limit of an infinite system size, for pairs of bonds that can both be written in the same 12-site unit cell. We compare the result to results from numerics on a 36-site system. We do not consider pairs of bonds that cannot be written in the same unit cell, as at this separation, finite-size effects will become important in the numerics and the comparison will become impossible. We could in principle improve on our comparison with numerics by computing the spinon correlation functions in a 36-site system also, in which case it should be possible to compare all bonds, but we have not done this.
The results are shown in table I, where in the column Theory I we show the simplest approximation, and in the column Theory II we show the second level of approximation. Qualitatively, the theory works quite well on the signs even at this level, getting 9 out of 12 correct. The only signs that the theory gets wrong at this level occur when the theory predicts a very small value ($`<.01`$) for the correlation function.
To improve this result, we included the third level of approximation in which we also project out onto a fifth site $`(m)`$ for those dimer correlation functions such that there is one and only one site $`(m)`$ which neighbors both dimers. In the column Theory III we show this level of approximation, as well as the particular site $`(m)`$ that we picked. Once this is done all the signs work out for 11 out of 12 correlations, and the qualitative agreement is almost perfect.
The magnitudes work out less well, as most of the dimer-dimer correlations are far too small within the spinon calculation. However, RVB calculations are quite poor at getting long range correlations without including some gauge fluctuations. For example, in the one-dimensional Heisenberg antiferromagnet, the spins on the same sublattice are uncorrelated. By including gauge fluctuations, this result can be substantially improved.
For the largest positive correlation function, $`C_{(6,7)(1,8)}`$, the magnitude does work out well even at the simplest approximation. However, for the largest negative function, $`C_{(6,7)(11,5)}`$, the magnitude is off by roughly a factor of 5, until we go to the third level of approximation, at which point the magnitude becomes roughly correct. We can hope that a better inclusion of fluctuations will improve these results, just as it has done for the one-dimensional chain. Perhaps projecting on an entire 12-site cell would give better results, as suggested by the improvement in the results in column III. Detailed calculations of projection for some trial wavefunctions on the kagomé lattice have been performed by Hsu and Schofield. A similarly detailed calculation for our parent state would be of interest.
## XI Conclusion
In conclusion, we have constructed an RVB state on the kagomé lattice, the “parent state”, which has a Dirac structure. Consideration of the various mass perturbations to the Dirac equation unifies several other previously suggested long-range states of the kagomé lattice Heisenberg antiferromagnet, with the exception of the BCS state and the bosonic $`Sp(2N)`$ state.
While at the projected mean-field level the chiral spin liquid appears to be the best RVB state, we have argued by a renormalization group treatment that fluctuations provide a mechanism for stabilizing a state with a nonchiral mass term. The numerical evidence also argues against a chiral state. We have then proceeded to explicit comparison with numerics, taking as input only one quantity, the triplet mass gap.
The physical idea behind our construction is that, given massless fermions, the system must ultimately try to break some symmetry to give mass to the fermions. There are many ways of doing this, but they all correspond to either introducing chirality, or to introducing some kind of spin-solid state in which the system dimerizes the $`t_{ij}`$. If we ignore the chiral state, we must have some kind of spin solid: we have proposed one possible spin solid and its attendant pseudo-Goldstone excitations. Other spin solids are possible, and may in fact be realized in favor of our proposal, but the general principle that the massless fermions must break a symmetry and produce pseudo-Goldstone excitations should be robust. Two other possibilities for spin solids that must be considered are the constant $`M_6`$ solid, and the perfect hexagon system with spatially varying $`M_6`$.
Experimentally, it may be possible in principle to detect the spin-Peierls order. Although this as appears as long-range order only in 4 spin correlation functions, it should give rise to a short-range oscillatory piece in the spin-spin correlation function. Since the spin-spin correlation function decays exponentially and the dimerization is weak, this would be very difficult to detect, but in principle it is possible via neutron scattering.
Theoretically, more work is needed on the fluctuations about the parent state. Doing a Gutzwiller projection of the wavefunctions on a 36-site lattice should enable much more direct comparison with numerics, especially for the dimer-dimer correlation functions.
Numerically, it may be possible to confirm the identification of the low-energy states with tower states from symmetry breaking. If one computes a dimer-dimer correlation function with the same pair of sites taken at two different times it may be possible to see the oscillations in the mass field. This may be difficult, though, given the relatively weak amount of dimerization present and the problem of extracting the contribution of the mass term to dimer-dimer correlation functions from the background of the fermionic contribution. Similarly, one can try to compute a susceptibility to spin-Peierls ordering by explicitly dimerizing coupling constants $`J`$ in equation (1). If indeed the system wishes to spontaneously order in the thermodynamic limit, then the susceptibility to dimerization should be large. This may make it possible to unambiguously determine whether the system prefers the $`M_{12}`$ or other ordering pattern.
It might also be interesting numerically to look at model systems in which the Hamiltonian has additional terms coupling to the chirality operator on each triangle. This might make it possible to probe the stability of the system to the chiral mass term, as well as providing some interesting states in which the Dirac fermions are moving in a large net magnetic field.
## XII Appendix: An Interesting Flat Band Case
When the system has flux $`\pi /4`$ through each triangle and flux $`\pi /2`$ through each hexagon the band structure becomes very peculiar. Using a 12-site unit cell, we find that the lowest band is doubly degenerate, and almost exactly flat. The next band above that is quadruply degenerate and exactly flat. The higher bands, which are all empty, are not flat.
It is very unlikely that any such state could be stabilized. We have argued above the the kagomé lattice is not a chiral spin liquid. However, it may be possible to add chirality operators to the Hamiltonian to tune the flux through the triangles.
In this case, the physics of the flat band state would be very amusing. We are used to the fact that in, for example, the one-dimensional Heisenberg antiferromagnet, one can deduce that the spin-1 excitations are composite objects of two spinons by looking at the excitation spectrum. Since the energy of the spin-1 object is a sum of two different energies, there is a continuum of possible energies. When, however, one of the two spinons is a hole excitation from a flat band, the energy is constant over the band, and there is no sign of the composite nature of the spin waves when looking at their energy spectrum. One would have a situation with one spinon hopping freely over the lattice, while the other spinon sits unmoving on a given unit cell!
## XIII Acknowledgements
I would like to thank S. Sondhi for suggesting the problem of the kagomé antiferromagnet, and for many useful discussions and insights. I would also like to thank R. Moessner for useful discussions on theory and experiment in frustrated antiferromagnets, V. N. Muthukumar for discussions on RVB ideas, A. Vishwanath for discussions on Chern numbers, and S. Sachdev for clarifying the estimate of the pion gap.
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# Rotational structure of 𝑇=0 and 𝑇=1 bands in the 𝑁=𝑍 nucleus 62Ga
## 1 Introduction
One of the currently most interesting questions in nuclear structure research concerns the interplay between isovector ($`T`$=1) and isoscalar ($`T`$=0) neutron-proton pairing, see e.g. and references therein. Effects from this interplay can be seen mainly in nuclei with $`NZ`$, since the $`pp`$ and $`nn`$ $`T`$=1 pairing forces dominate in nuclei with proton (or neutron) excess. Odd-odd nuclei with an equal number of neutrons and protons are of particular interest since the isoscalar and isovector $`np`$ forces compete, and the position of $`T`$=0 states relative to $`T`$=1 states is found to change with the mass number: Most odd-odd $`N`$=$`Z`$ nuclei with $`A40`$ have $`T`$=0 ($`I>0`$) ground states while most of the heavier nuclei have $`T`$=1 ($`I`$=0) ground states.
Recently, several new states of positive parity were found in the odd-odd $`N`$=$`Z`$ nucleus $`{}_{31}{}^{}{}_{}{}^{62}`$Ga<sub>31</sub> . The states form a $`T`$=0 rotational band from $`I`$=1 up to $`I`$=17 that backbends at $`I9`$, see fig.1. The $`I`$=1 $`T`$=0 state decays to the $`I`$=0 $`T`$=1 ground state, and no excited $`T`$=1 states were identified. The odd-odd nature of <sup>62</sup>Ga implies blocking of proton-proton and neutron-neutron isovector ($`T`$=1) pairing properties, and both $`T`$=1 and $`T`$=0 neutron-proton pairing may be important . From shell model calculations it seems that the relative importance of these two pair fields change as a function of angular momentum, since $`T`$=0 states are more favored than $`T`$=1 states for $`I\stackrel{>}{}\mathrm{\hspace{0.33em}2}`$. We shall address these questions by comparing experimental data to a shell model (SM) calculation and an unpaired cranked Nilsson-Strutinsky (CNS) calculation.
In the ground state of $`{}_{31}{}^{}{}_{}{}^{62}`$Ga<sub>31</sub> the $`f_{7/2}`$ shell is filled, and the three protons and three neutrons outside the <sup>56</sup>Ni core occupy the lowest-lying orbitals emerging from $`p_{3/2}`$ and $`f_{5/2}`$. In the deformed shell model this implies a small prolate or triaxial deformation ($`\epsilon 0.17`$) and thus collective rotation. The relatively small number of active particles allows extended shell model calculations utilizing the valence shells $`pf_{5/2}g_{9/2}`$. Collective rotation and backbending can then be studied and compared in the laboratory frame (spherical shell model) and in the intrinsic frame (deformed mean field) description. Similar studies have recently been performed for the even-even nuclei <sup>48</sup>Cr and <sup>36</sup>Ar .
Results from a shell model calculation for <sup>62</sup>Ga in the valence space $`pf_{5/2}g_{9/2}`$, utilizing a slightly modified version of the effective interaction , are presented in section 2. Low-lying $`T`$=0 and $`T`$=1 bands are calculated, and we discuss the role of isovector and isoscalar pairing for the different bands as a function of angular momentum. In section 3, low-lying $`\alpha `$=1 (odd spins) and $`\alpha `$=0 (even spins) bands in <sup>62</sup>Ga calculated in CNS neglecting pairing are presented, and it is discussed how the bands can be approximately identified with $`T`$=0 and $`T`$=1 states. In section 4 the CNS and SM results are compared to experimental energies. Also shell occupation numbers are calculated and compared in both models, and the nature of the observed backbending is subsequently discussed. In section 5 we report on calculations of $`B(E2)`$-values and spectroscopic quadrupole moments in SM as well as in CNS. Finally, a short summary of the results is given in section 6.
## 2 Shell model calculation
At present large scale shell model codes can easily handle 6 particles in the model space $`pf_{5/2}g_{9/2}`$. For <sup>62</sup>Ga this implies an assumption of a closed $`f_{7/2}`$ shell that is a good approximation for most of the states discussed here. However, later we will see that this assumption is rather severe for the description of the higher spin states.
The shell model calculation was performed using the effective interaction derived from a realistic G matrix whose monopole part has been phenomenologically adjusted . The interaction was previously used to describe e.g. <sup>76</sup>Ge and <sup>82</sup>Se . The single-particle energies were taken from the experimental spectrum of <sup>57</sup>Ni: $`\epsilon _{p_{3/2}}`$=0.0, $`\epsilon _{f_{5/2}}`$=0.77, $`\epsilon _{p_{1/2}}`$=1.113, $`\epsilon _{g_{9/2}}`$=3.0 MeV (estimated). Recently, the single-particle energy of the $`g_{9/2}`$ shell in <sup>57</sup>Ni was experimentally established . The measured value, $`\epsilon _{g_{9/2}}`$=3.7 MeV, is larger than previously used value, 3.0 MeV. However, it was found that the agreement of the calculated spectrum with the experimental energies in <sup>62</sup>Ga is better if the single-particle energy of the $`g_{9/2}`$-shell is lowered, thus we used $`\epsilon _{g_{9/2}}`$=2.5 MeV, and this was the only modification of the interaction .
### 2.1 $`T`$=0 and $`T`$=1 bands
The $`T`$=0 band was observed up to $`I`$=17 while only one $`T`$=1 state, the $`I`$=0 ground state, was observed. The $`T`$=0 states form a rotational band that backbends at $`I`$=9, see fig.1. The 246 keV, 376 keV and 1241 keV $`\gamma `$-rays were assigned experimentally as corresponding to stretched $`E2`$ transitions, and the $`T`$=0 levels $`1^+`$, $`3^+`$, $`5^+`$ and $`7^+`$ could thus be assigned , while the observed levels above the backbending could be assigned mainly from energy systematics and comparisons to SM calculations. No $`B(E2)`$ values have so far been measured, except for the $`3^+`$ state that was found to be isomeric with a life-time of 4.6(1.6) ns corresponding to $`B(E2)=197(69)`$ e<sup>2</sup>fm<sup>4</sup> .
Experimental data is compared to calculated energies of low-lying $`T`$=0 and $`T`$=1 states in fig.2. Measured $`T`$=1 energies above the $`I`$=0 state are taken from the analogue states in <sup>62</sup>Zn . In agreement with data the ground state is found to have $`I`$=0 and isospin $`T`$=1 in the SM calculation. Also the excitation energy from the $`T`$=1 ground state to the first excited $`T`$=0 $`I`$=1 state is well reproduced. With increasing angular momentum the $`T`$=0 states become more favored, and the $`T`$=0 band crosses the $`T`$=1 band already at $`I2`$. At low spins the agreement with data is very good for the $`T`$=0 states while the higher spin states ($`I>11`$) are less well reproduced. This is presumably due to the restricted model space used in the SM calculation, as will be discussed in sect.4. Fig.2a shows two calculated states at spin $`I`$=9. The higher-lying state is connected to the yrast states with lower spin by a strong $`B(E2)`$ value, and is constructed only from $`fp`$ shell orbits. The configuration of the yrast $`I`$=9 state involves the excitation of about two particles into the $`g_{9/2}`$ shell and is connected by a strong $`E2`$ matrix element to the yrast $`I`$=11 state. There is a strong $`B(E2)`$ transition from the yrast 9<sup>+</sup> to the yrare 7<sup>+</sup> state that also has two particles excited to the $`g_{9/2}`$ shell. This state, however, is calculated to have larger excitation energy than $`9_1^+`$ state. Thus in the sequence of odd spin states with $`T`$=0 we see a change of configuration at $`I9`$ that later will be discussed in a greater detail. A similar configuration change appears at $`I10`$ in the $`T`$=1 band, see fig.2c. In the $`T`$=1 band constructed from $`fp`$-shell orbits only, the energies of states above spin $`I`$=4 are not very well reproduced in the SM calculation, while the calculated energies of the $`I`$=10, 12, 14 and 16 states, which involve the excitation of two particles to the $`g_{9/2}`$ shell, come out quite close to the experimental values.
The increased favoring of $`T`$=0 states relatively to $`T`$=1 states implies a moment of inertia that is larger for the $`T`$=0 band than for the $`T`$=1 band. The origin of this difference may be found in different pairing properties of the two bands. We shall therefore study $`T`$=0 and $`T`$=1 pairing energy contributions to different states in the $`T`$=1 and $`T`$=0 bands.
### 2.2 Pairing in the SM
The pairing energy in the shell model is calculated by taking energy difference of states calculated using the full Hamiltonian and the Hamiltonian from which a schematic $`L`$=0 pairing interaction is subtracted (cf. ). The normalized form of the pairing interaction was used with the strength parameters $`\overline{G}_{01}`$=2.98 and $`\overline{G}_{10}`$=4.75 for the $`T`$=1 and $`T`$=0 pairing, respectively. Since the excited $`T`$=0 band has a very different configuration as compared to the lower $`T`$=0 band, it was possible to calculate the pairing energy for each configuration separately even in the band-crossing region.
In fig.3 $`T`$=0 and $`T`$=1 pairing energy contributions in the orbital angular momentum $`L`$=0 channel is shown for $`T`$=0 odd spin states (fig.3a), for $`T`$=1 even spin states (fig.3b), and for $`T`$=0 even spin states (fig.3c). The $`T`$=1 $`I`$=0 ground state is seen to have approximately equal amounts of pairing energy contributions from the $`T`$=0 and $`T`$=1 channels. For the lowest $`T`$=0 state ($`I`$=1), on the other hand, the isoscalar neutron-proton ($`T`$=0) pairing plays a larger role than the isovector ($`T`$=1) pairing, although the latter has contributions from neutron-neutron, proton-proton and neutron-proton pairing. The same is valid for the lowest $`T`$=0 state with even spin, $`I`$=2.
Due to isospin symmetry all parts of the isovector pairing, $`nn`$, $`pp`$ and $`np`$, are identical in the $`T`$=0 states. The odd-odd nature of <sup>62</sup>Ga implies that at least one isovector pair is broken in the $`T`$=0 states, i.e. the isovector pairing is reduced. From fig.3a it is thus seen that each component of the isovector pairing mode contributes to the $`T`$=0 ($`I`$=1) state only about 0.25 MeV.
The isoscalar pairing energy is calculated as about 1.4 MeV in lowest $`T`$=0 as well as in lowest $`T`$=1 states, while the total isovector pairing energy contributions are about 0.7 MeV and 1.5 MeV, respectively. In total, the $`T`$=1 ($`I`$=0) ground state thus gets a pairing energy contribution of about 2.9 MeV (fig.3b), while the corresponding energy is about 2.1 MeV for the lowest $`T`$=0 state ($`I`$=1), see fig.3a. This favoring by about 0.8 MeV pairing energy of the lowest $`T`$=1 state compared to lowest $`T`$=0 state plays an important role in making the $`T`$=1 state to become the ground state.
With increasing spin the pairing energy, both isoscalar and isovector, is approximately constant for the lowest odd-$`I`$, $`T`$=0 states, while for the $`T`$=1 states pairing energy decreases rather drastically, especially the isovector part. That is, the moment of inertia for the odd-$`I`$ $`T`$=0 states is more or less unaffected by pairing correlations (for $`I<9`$), while it is reduced by about 50% for the $`T`$=1 states. This causes the $`T`$=1 band to cross the $`T`$=0 band already at $`I2`$. The pairing energy decreases with increasing spin also for even-$`I`$ $`T`$=0 states, see fig.3c. However, since $`T`$=1 pairing is partly blocked the effect is less drastical for these states as compared to $`T`$=1 states.
In SM calculations of the ground band in $`{}_{24}{}^{}{}_{}{}^{48}`$Cr<sub>24</sub> it was found that isoscalar pairing correlations are less dependent on increasing spin than isovector pairing correlations , i.e. the moment of inertia is less sensitive to isoscalar pairing. The present results for <sup>62</sup>Ga do, however, suggest similar energy correlations with increasing spin from isoscalar and isovector pairing, or even a larger energy change due to isoscalar pairing, see fig.3. The different pairing behavior found in the two nuclei is explained by the size of the pairing energy in the ground state. In <sup>48</sup>Cr isovector pairing gives a much larger energy contribution to the $`I`$=0 ground state, $`E_{01}`$= 3.6 MeV, as compared to isoscalar pairing, $`E_{10}`$=2.0 MeV . In <sup>62</sup>Ga, on the other hand, the two contributions are similar for the $`T`$=1 ground state ($`E_{01}E_{10}1.5`$ MeV), while isoscalar pairing energy is largest for $`T`$=0 states. Since isoscalar and isovector pairing energy both decrease with increasing spin, and are approximately zero at the same spin value (when seniority is large enough), the energy change with spin is mainly determined by the energy contribution to the starting spin state. However, the pairing energy may decrease in different ways for different bands, see fig.3.
The configuration of the low-spin states can be followed smoothly to the non-yrast $`I`$=9 and $`I`$=10 states (for $`T`$=0 and $`T`$=1). These states constitute approximately maximum seniority states ($`v=6`$), and all studied pairing energy contributions vanish there. The yrast states with spin values higher or equal to $`I`$=9 (or $`I`$=10) correspond to states where approximately two particles are excited to the $`g_{9/2}`$ shell, i.e. they have the configuration $`(pf_{5/2})^4g_{9/2}^2`$ (or, if also the particles in the filled $`f_{7/2}`$ shell are counted, $`(fp)^{20}g_{9/2}^2`$). For these states pairing correlations are again active, see fig.3. In a similar way as for the $`(fp)^{22}`$ configuration, isoscalar and isovector pairing decrease with increasing spin and are zero for the $`I`$=17 and $`I`$=16 states, which are the maximum spin states in the $`(pf_{5/2})^4g_{9/2}^2`$ configuration. The main difference between the two configurations, $`(pf_{5/2})^6`$ and $`(pf_{5/2})^4g_{9/2}^2`$, is that the isovector pairing is of a similar size as the isoscalar pairing in the $`(fp)^{20}g_{9/2}^2`$ configuration. This may be explained in the following way. Since the two odd nucleons excited to the $`g_{9/2}`$ shell can couple to isospin zero (preferentially with $`I`$=9), the two pairs in the $`fp`$ shell are fully active for $`T`$=0 pairing as well as for $`T`$=1 pairing.
### 2.3 High-spin behavior
The somewhat irregular energy behavior of the $`(fp)^{20}g_{9/2}^2`$ band between $`I`$=9 and 17, as seen in experimental data (fig.2), can be understood from a schematic seniority model discussed in . Two neutrons and two protons occupy $`(p_{3/2},f_{5/2},p_{1/2})`$ and one neutron and one proton occupy $`g_{9/2}`$. The lowest energy of a given spin is obtained by first minimizing the total seniority, $`v=v_p+v_n`$, and then maximizing the reduced isospin, $`t=\frac{1}{2}|v_pv_n|`$, where $`v_p`$ and $`v_n`$ are the proton and neutron seniority, respectively. With one neutron and one proton into the $`g_{9/2}`$ shell, $`9^+`$ is the highest spin that can be constructed for $`v`$=2. By breaking the neutron (or proton) pair in $`(p_{3/2},f_{5/2},p_{1/2})`$ at most 4 more units of $`\mathrm{}`$ can be obtained, and the $`11^+`$ and $`13^+`$ states correspond to $`v`$=4 and $`t`$=1 in this simple model. In the same way the remaining proton (or neutron) pair is broken giving $`15^+`$ and $`17^+`$ with $`v`$=6 and $`t`$=0. The arc-like structure seen in the energy sequence $`I=9^+13^+`$, and much more evident for $`I=13^+17^+`$, thus reminds of a seniority coupling scheme in $`(p_{3/2},f_{5/2},p_{1/2})`$. The smaller collectivity (smaller configuration mixing) for the higher spin states implies a coupling scheme more similar to a pure seniority scheme. The band termination at $`17^+`$ is thus smooth ($`\mathrm{\Delta }v`$=0 and $`\mathrm{\Delta }t`$=0) and favored , and is very similar to the band terminations seen e.g. at $`8^+`$ in $`{}_{30}{}^{}{}_{}{}^{60}`$Zn<sub>30</sub> ($`(p_{3/2}f_{5/2})^4`$), at $`12^+`$ in $`{}_{22}{}^{}{}_{}{}^{44}`$Ti<sub>22</sub> ($`f_{7/2}^4`$) and at $`15^+`$ in $`{}_{23}{}^{}{}_{}{}^{46}`$V<sub>23</sub> ($`f_{7/2}^6`$). However, as will be shown in next section, the CNS calculation suggests that also the $`f_{7/2}`$ shell is important for the description of the higher spin states in <sup>62</sup>Ga, and the band termination scenario may be more complicated than discussed here.
## 3 Cranked Nilsson-Strutinsky calculation
The cranked Nilsson-Strutinsky calculation is performed with standard parameters , and with the possibility to fix the configuration in the minimization procedure . We allow for a free minimization of the two quadrupole degrees of freedom, $`\epsilon `$ and $`\gamma `$, as well as one hexadecapole degree of freedom, $`\epsilon _4`$. All kinds of pairing interactions are neglected. It is the experience that also unpaired CNS calculations give a good insight into the nuclear structure, in particular at higher spins. Below we present the results of a CNS calculation of the rotational behavior of <sup>62</sup>Ga. Since SM states are additionally classified with the isospin quantum number, we first need an understanding of isospin in the CNS model (subsection 3.1). In subsection 3.2 we discuss the behavior of the lowest rotational bands with even as well as odd spins, i.e., with signature quantum numbers $`\alpha `$=0 and $`\alpha `$=1, respectively.
### 3.1 Isospin in CNS
In the cranked Nilsson-Strutinsky model isospin is not a good quantum number. The Coulomb force (included in the present model through different sets of Nilsson parameters for protons and neutrons) breaks isospin, and total wave-functions with good isospin are not constructed. If pairing is neglected one may still construct states of approximately good isospin for nuclei with $`NZ`$ (cf. ). For odd-odd $`N`$=$`Z`$ nuclei we may approximately identify a configuration with isospin $`T`$=0 if it cannot be realized (due to the Pauli principle) in the even-even neighbor. That is, if all protons are placed in “exactly” the same single-particle orbits as the neutrons (the orbits are nearly the same since the Coulomb force plays a minor role), the configuration is identified with isospin $`T`$=0. This means that the two odd nucleons in <sup>62</sup>Ga have the same signature and parity, i.e. they form a pair with the total signature $`\alpha `$=1 and parity $`\pi `$=+. The same applies for each neutron-proton pair, and due to an odd number of pairs the total signature and parity is thus $`\alpha `$=1 and $`\pi `$=+, corresponding to a positive parity rotational band with odd spins. The shell filling of the three least bound protons and neutrons in <sup>62</sup>Ga for the $`\alpha `$=1 and $`\alpha `$=0 bands at appropriate equilibrium deformation, is shown in fig.4. On the right-hand side (positive $`\gamma `$) we show filling of proton-neutron pairs in (approximately) identical orbitals, leading to the total signature $`\alpha `$=1 which may be identified with an isospin $`T`$=0 configuration.
If the two odd nucleons are placed in orbits with different signatures, the proton (neutron) can be replaced by a neutron (proton). Consequently, the state can be realized also in the even-even neighbors. This situation is shown in the left-hand part of the single-particle diagram (fig.4). The last filled neutron and proton now occupy orbitals with different signatures, and a state with total signature $`\alpha `$=0 is formed. Since the orbits of the odd neutron and proton can be interchanged, two almost degenerate bands are formed in the CNS. Both bands have $`\alpha `$=0, i.e. even spins, and form the basis of one $`T`$=0 band (signature partner of the $`\alpha `$=1 $`T`$=0 band) and one $`T`$=1 band. Without additional correlations between particles the two bands should thus come very close in energy, and we may compare the calculated $`\alpha `$=0 band(s) with the lowest-lying $`T`$=1 band as well as with the lowest $`T`$=0 band with even spins.
With this simple classification, neglecting pair correlations, the $`T`$=0, $`\alpha `$=1 band will in general be lowest in energy, and the excitation energy to the $`T`$=1 band approximately correspond to the calculated signature splitting of the last filled orbital. However, the two signature partners may correspond to nuclear shapes with different deformations, implying a more complicated relation between the energies of the two bands also in unpaired calculations.
### 3.2 Total signature $`\alpha `$=0 and $`\alpha `$=1 bands
The calculated yrast $`\alpha `$=0 (“$`T`$=1/$`T`$=0”) and $`\alpha `$=1 (“$`T`$=0”) band energies are shown in fig.2, and equilibrium deformations are shown in fig.5. In figs.2b and 2d a rotational reference has been subtracted to facilitate the reading of the figure. At low spin values the configuration with all 22 particles outside <sup>40</sup>Ca placed in the $`fp`$ shell, $`(fp)^{22}`$, comes out as lowest in energy for even as well as for odd spin values in the CNS calculation. At $`I\stackrel{>}{}\mathrm{\hspace{0.33em}9}`$ it is more favorable to excite two particles to the $`g_{9/2}`$ shell and the configuration $`(fp)^{20}g_{9/2}^2`$ becomes yrast for states in the $`\alpha `$=0 as well as $`\alpha `$=1 bands.
All states in the yrast $`(fp)^{22}`$ configuration, the $`\alpha `$=1 band as well the $`\alpha `$=0 band, are calculated to be triaxial with rather large values of the triaxiality parameter $`\gamma `$, see fig.5. The $`\alpha `$=1 states correspond to rotation around the smallest axis (positive $`\gamma `$), while the $`\alpha `$=0 states correspond to rotation around the (classically forbidden) intermediate axis (negative $`\gamma `$). The different types of triaxiality for the two bands can be understood from the behavior of single-particle energies as functions of triaxiality parameter, $`\gamma `$, see fig.4. For the $`\alpha `$=1 band the two odd nucleons occupy an orbit with the signature $`\alpha `$=+1/2 which prefers positive $`\gamma `$, and the energy for the total configuration obtains a minimum at $`\gamma >0`$. In the $`\alpha `$=0 states one of the two odd nucleons occupies an orbital with $`\alpha `$=$``$1/2 which prefers negative $`\gamma `$, and the total configuration ends up with a deformation $`\gamma <0`$.
With increasing spin the $`\gamma `$ value increases for both the $`\alpha `$=1 and $`\alpha `$=0 bands; for the $`\alpha `$=1 band from $`\gamma 25^{}`$ for the $`I`$=3 state to $`\gamma =60^{}`$ (i.e. band termination in a non-collective state, cf. ) at $`I`$=9, and for the $`\alpha `$=0 band from $`\gamma 25^{}`$ for the $`I`$=2 state to $`\gamma =5^{}`$ for $`I`$=8, see fig.5. The termination of the $`\alpha `$=1 band at $`I^\pi `$=9<sup>+</sup> corresponds to an oblate shape where the angular momentum components $`\mathrm{\Omega }`$=3/2 and 1/2 of $`p_{3/2}`$ and the $`\mathrm{\Omega }`$=5/2 component of $`f_{5/2}`$ are occupied by the three protons and three neutrons outside the $`f_{7/2}`$ shell. Electromagnetic properties of these two bands are discussed in section 5 below.
States with the configuration $`(fp)^{20}g_{9/2}^2`$, $`\alpha `$=1 (“$`T`$=0”) are also found to be triaxial with rotation around the smallest axis (positive $`\gamma `$), see fig.5. The degree of triaxiality increases smoothly along the band from $`\gamma 10^{}`$ at $`I`$=7 to the oblate symmetric state, $`\gamma =60^{}`$, at $`I`$=17 which terminates the rotational band within this configuration. The $`\alpha `$=0 (“$`T`$=0/$`T`$=1”) $`(fp)^{20}g_{9/2}^2`$ band behaves in a similar way as the $`\alpha `$=1 $`(fp)^{20}g_{9/2}^2`$ configuration, and terminates at $`I`$=16. The CNS calculated energies of $`\alpha `$=1 isobaric analog states in <sup>62</sup>Zn were reported in .
## 4 Comparison between unpaired CNS and SM
In this section we compare results obtained in the two models. In subsection 4.1 energies are compared, and occupation numbers in spherical $`j`$-shells are compared in subsection 4.2. The backbending observed in the $`T`$=0 band (fig.1) is discussed in subsection 4.3.
### 4.1 Energies
In fig.2 energies calculated in the CNS model are compared with both SM results and experimental values: the $`\alpha `$=1 states in figs.2a and 2b (experimental data from <sup>62</sup>Ga), and the $`\alpha `$=0 states in figs.2c and 2d (experimental data from <sup>62</sup>Zn). In particular the low-spin $`\alpha `$=1 states agree quite well with data for $`T`$=0 states (as well as with SM results), while the even spin states show large deviations with the $`T`$=1 states. In fact, the moment of inertia of the $`T`$=0 band appears very similar in unpaired CNS and SM calculations (and in experimental data). This is since pairing correlations are relatively stable for most $`T`$=0 states, while it is strongly changing with increasing spin for $`T`$=1 states, as discussed in sect.2.
Since the CNS calculation is performed without pairing it seems more relevant to compare CNS energies with unpaired SM energies. This is done by subtracting isovector and isoscalar ($`L`$=0) pairing contributions, taken from fig.3, from the SM energies for each individual state, and in fig.6 we compare energies of the lowest states calculated in the CNS model with SM results. The energies are indeed very similar in the unpaired CNS and the “unpaired” SM. Note, however, that it is not straightforward that the removed pairing energies in SM correspond to the (not included) pairing in CNS.
From fig.6a it is seen that the yrast sequence of odd spin states ($`\alpha `$=1), calculated in unpaired CNS, come quite close to the lowest odd spin states in SM calculation ($`T`$=0 states) when pairing has been removed. The fact that low-spin CNS energies agree with full SM results (fig.2) as well as with SM results without pairing (fig.6a) is explained by the rather constant pairing energy contribution for the low spin, odd-$`I`$, $`T`$=0 states, see fig.3a. However, the non-yrast 9<sup>+</sup> state, that terminates the $`(fp)^{22}`$ configuration has a pairing energy contribution which is approximately zero. The SM results without pairing thus deviate rather strongly from the CNS results if all states in the $`(fp)^{22}`$ band are studied. The origin of this deviation is not clear.
Since $`\alpha `$=0 states (even spins) calculated in CNS can be compared either to $`T`$=1 even spin states or to $`T`$=0 even spin states, both types of SM states are shown in fig.6b. In the full SM calculation the rotational bands with $`T`$=0 and $`T`$=1 even spin states show quite different energy behavior. However, when pairing has been subtracted they are similar, and also similar to unpaired CNS results. The unpaired $`T`$=1 states come approximately 1 MeV higher in energy for most spin states compared to the $`T`$=0 states, and the $`T`$=0 states are close to unpaired CNS energies. The approximately constant energy difference between $`T`$=0 and $`T`$=1 states that remains in the SM calculation after all $`L`$=0 pairing has been subtracted, originates from other parts of the interaction.
### 4.2 Sub-shell occupancy
A more detailed understanding of the structure of the $`T`$=0 and $`T`$=1 bands can be achieved by studying wave functions. In SM the wave functions are extremely complex and composed by very many components. A simple view can, however, be obtained by studying occupation numbers in different $`j`$ shells. Such studies are also feasible in CNS, and some aspects of wave functions can then be compared in the two models.
The added neutron and proton occupation numbers in spherical $`j`$-shells for some selected states are compared in Table 1. The method described in ref. was used to calculate the occupancies in CNS. In SM, the total sum of occupation numbers in the model space ($`p_{3/2}`$, $`f_{5/2}`$, $`p_{1/2}`$ and $`g_{9/2}`$) is 6, while CNS does not have such limitation for the configuration space. The $`f_{7/2}`$ shell belongs to the core in SM and is consequently always “occupied” by 16 particles.
The total occupancy in the $`fp`$-shell orbits $`f_{7/2}`$, $`p_{3/2}`$, $`f_{5/2}`$, $`p_{1/2}`$, and in $`g_{9/2}`$ covers 21.45-21.95 particles in CNS, and the remaining part, 0.55-0.05 particles, is distributed over other shells due to $`\mathrm{\Delta }N`$=2 oscillator shell mixing. It is seen that the two bands described as $`T`$=0 (exemplified by $`I^\pi `$=3<sup>+</sup> and 9<sup>+</sup>) and $`T`$=1 ($`I^\pi `$=2<sup>+</sup> and 8<sup>+</sup>) with the same configuration, $`(fp)^{22}`$, have rather similar occupation numbers. This is also the case in the SM calculation. In general, the occupation numbers are quite similar in CNS and SM calculations. For example, the occupation of $`f_{5/2}`$ increases with increasing spin in both models. The major difference between the two calculations is the relative distribution of particles in $`p_{3/2}`$ and $`p_{1/2}`$. All calculated states in SM have almost 1 particle less occupying $`p_{3/2}`$ and 0.5 occupying $`p_{1/2}`$ more, as compared to the corresponding states in CNS. This may be related to the difference in the single-particle spin-orbit energy splitting between $`p_{3/2}`$ and $`p_{1/2}`$ in SM (1.1 MeV) and in CNS (2.8 MeV).
Largest deviations between SM and CNS occupancies appear above the band-crossing, in particular at the 17<sup>+</sup> state which terminates the rotational band. In the corresponding configuration the three “valence” protons and neutrons occupy the $`\mathrm{\Omega }`$=9/2 orbital from $`g_{9/2}`$, $`\mathrm{\Omega }`$=5/2 orbital from $`f_{5/2}`$ and $`\mathrm{\Omega }`$=3/2 orbital from $`p_{3/2}`$, giving $`I=2(9/2+5/2+3/2)=17`$. Due to a rather large oblate deformation of the 17<sup>+</sup> state ($`\epsilon 0.27`$, $`\gamma =60^{}`$) predicted in CNS, the orbitals emerging from $`p_{3/2},`$ $`f_{5/2}`$ and $`p_{1/2}`$ shells contain large components of the (spherical) $`f_{7/2}`$ shell. Due to this mixing the $`17^+`$ state contains approximately two holes in the spherical $`f_{7/2}`$ shell, see Table 1. Note that this band is described as in the notation of ref. , i.e., with no holes in orbitals emerging from the (deformed) $`f_{7/2}`$ shell in the CNS calculation.
Thus from the results presented in Table 1, one may see that, according to the CNS calculation, it might be important to include $`f_{7/2}`$ in the SM calculation for an accurate description of states in the $`(fp)^{20}g_{9/2}^2`$ band, particularly close to band termination.
### 4.3 Backbending
Both unpaired CNS calculations and full SM calculations reproduce the observed backbending very well, see fig.1. In fact, also “unpaired” SM calculation describes the backbending quite well (fig.6a). From the CNS calculation we find that for $`I9`$ it is energetically favorable to excite the two odd nucleons from the $`fp`$ shell into $`g_{9/2}`$ orbits. Also the SM calculation suggests that the lowest $`T`$=0 states with $`I`$$``$9 are formed by exciting two nucleons to $`g_{9/2}`$, see Table 1. The backbending is thus suggested to be caused by an unpaired band-crossing involving a two-particle-two-hole excitation. The origin of the backbending in <sup>62</sup>Ga is thus very different from backbending in heavier nuclei, where it is understood as rotation-induced alignment of angular momentum vectors of a pair of neutrons (or protons) in a high-$`j`$ orbit . This type of backbending is an effect from pairing, and does not appear in unpaired solutions. Unpaired band-crossings have been observed in the $`A160`$ mass region at high spin (where pairing plays a minor role) .
In unpaired CNS calculations the backbending in <sup>62</sup>Ga originates from a sharp single-particle level crossing (for neutrons as well as for protons) between the $`fp`$-shell orbit and $`g_{9/2}`$ orbit. At low spin, when the deformation is close to prolate, the energetically most favored $`g_{9/2}`$ orbit can be approximately assigned as the \[440 1/2\] Nilsson orbit with signature $`\alpha `$=1/2. The $`fp`$-shell orbit has mixed $`p_{3/2}f_{5/2}`$ components (\[321 1/2\], $`\alpha `$=1/2 in the low-spin limit for prolate deformations). The excitation (band-crossing) thus leaves the filled (deformed) $`f_{7/2}`$ shell intact<sup>1</sup><sup>1</sup>1 The band is thus different from a recently observed excited band in $`{}_{29}{}^{}{}_{}{}^{58}`$Cu which has also one neutron and one proton promoted to the $`g_{9/2}`$ shell, but where the excitation takes place from deformed orbits emerging from $`f_{7/2}`$ , giving rise to a larger deformation ($`\epsilon 0.4`$). Compared to the discussed band in <sup>62</sup>Ga, the band in <sup>58</sup>Cu is obtained by removing two protons and two neutrons from the $`f_{7/2}`$ shell (in the notation of ref. , the configuration is ).. Also this configuration has approximate isospin $`T`$=0. In the single-particle diagram (fig.4) the excitation can be seen in the $`\alpha `$=1 configuration as a lifting of one proton as well as one neutron from the highest-lying filled negative-parity orbital to the lowest un-occupied positive-parity orbital. At the rotational frequency for which fig.4 is drawn (corresponding to $`I`$$``$6$`\mathrm{}`$), the excitation is energetically unfavored, but at higher rotational frequency the positive-parity orbital comes down in energy below the filled negative-parity orbital, and the excitation is energetically favorable for $`I9`$. The excitation implies an increased quadrupole deformation, from $`\epsilon 0.17`$ to $`\epsilon 0.23`$, see fig.5.
At $`I`$=9 the two configurations come very close in energy, see fig.2. In CNS the two configurations can be seen as separate energy minima in the potential-energy surface shown in fig.7. The minimum of the $`(fp)^{20}g_{9/2}^2`$ configuration at ($`\epsilon `$,$`\gamma `$)= (0.23,$`17^{}`$) comes 330 keV lower in energy than the $`(fp)^{22}`$ state at ($`\epsilon `$,$`\gamma `$)=(0.13,$`60^{}`$), and the two minima are separated by a barrier having height of about 500 keV. The $`(fp)^{20}g_{9/2}^2`$ configuration is formed from the $`(fp)^{22}`$ configuration by exciting one proton and one neutron from the $`fp`$-shell to $`g_{9/2}`$, i.e. by a 2p-2h excitation, cf fig.4. Since the single-particle levels involved in the excitation have different parity they exhibit a sharp level crossing when plotted versus rotational frequency. Thus, in the mean field approximation the interaction between the two configurations is identically zero. The interaction may, however, be described beyond the mean field approximation as a dynamical tunneling phenomenon. Tunneling between the two minima gives rise to mixing of the two wave functions, and two new eigenstates appear. The SM in principle contains all kinds of dynamical correlations, and the two 9<sup>+</sup> states shown in fig.2 are already mixed.
Backbending is calculated also in the $`T`$=1 and $`T`$=0 bands with even spins ($`\alpha `$=0 band in CNS), see fig.6. Also in these cases the backbending corresponds to a band-crossing between bands with configurations $`(fp)^{22}`$ and $`(fp)^{20}g_{9/2}^2`$. In <sup>62</sup>Zn these states have been observed and CNS calculation presented in .
## 5 Electromagnetic properties
In previous sections we have shown that the structure of the $`T`$=0 and $`T`$=1 rotational bands changes with increasing spin. The changes occur both within a given configuration and, at the backbending, from one configuration to another. These changes have been seen in both SM and CNS calculations, see e.g. Table 1. In the CNS calculation large differences were also seen between the $`\alpha `$=1 ($`T`$=0) and $`\alpha `$=0 ($`T`$=0/1) bands in terms of different kinds of triaxiality (fig.5). These differences are not seen in sub-shell occupancies (Table 1) but may be significant for other properties.
Electromagnetic properties, such as stretched $`B(E2)`$ values and spectroscopic quadrupole moments, are expected to be sensitive to triaxiality and the axis around which the rotation occurs. In addition, they are also measurable, at least in principle. We shall therefore study electromagnetic properties of the considered states in this section, both in CNS and SM. Electromagnetic properties were calculated in the SM using effective charges $`e_p`$=1.5 and $`e_n`$=0.5. In CNS it is not straightforward to calculate electromagnetic properties, and in next subsection we describe the approximative method we used.
### 5.1 EM properties in CNS
Angular momentum is not a good quantum number in the cranking model, and there is no direct way to calculate electromagnetic properties. For axially symmetric shapes with collective rotation, the rotor model can be assumed to be valid, and electromagnetic properties approximately calculated from the electric quadrupole deformation of the mean field. For triaxial shapes the $`K`$ quantum numbers are mixed and the simple rotor model cannot be applied. However, in the high-spin limit it is possible to derive expressions for $`B(E2)`$ and spectroscopic quadrupole moment, $`Q_{\mathrm{spec}}`$ at any value of the triaxiality parameter $`\gamma `$ . The combined expressions of the axially symmetric rotor formulae and the high-spin formulae from were given in :
$$B(E2;I+2,KI,K)=\frac{5}{6\pi }I+2K20|IK^2Q_2(\widehat{x})^2,$$
(1)
$$Q_{\mathrm{spec}}(I,K)=2II20|IIIK20|IKQ_0(\widehat{x}),$$
(2)
where $`Q_0(\widehat{x})`$ and $`Q_2(\widehat{x})`$ are electric quadrupole moments around the rotation axis ($`x`$ axis). They are calculated in CNS from the proton wave-functions at the appropriate equilibrium deformation. These expressions are thus valid either at axial-symmetric shapes for any $`I`$ and $`K`$, or at triaxial shapes for high-spin values, when the Clebsch-Gordan coefficients have taken their asymptotic values (and are independent of $`K`$). However, as was seen in ref., where $`B(E2)`$ and $`Q_{\mathrm{spec}}`$ values for different states in <sup>48</sup>Cr were calculated from eqs.(1) and compared to SM results, the expressions seem to work quite well also outside the region of validity. With these restrictions in mind we shall use eqs.(1) to estimate $`B(E2)`$ and $`Q_{\mathrm{spec}}`$ values for different states calculated in <sup>62</sup>Ga. For $`B(E2)`$-values the change of deformation between mother state and daughter state is neglected; properties of the mother states are simply used.
### 5.2 Comparison between CNS and SM
In figs.8 and 9 we compare SM with CNS results for $`B(E2)`$ values and $`Q_{\mathrm{spec}}`$ values, respectively. The $`K`$ value appearing in eqs.(1) has been set to $`K`$=1 for the bands with odd spins, and $`K`$=0 for the band with even spins. The use of eqs.(1) for the low-spin parts of the bands is outside their expected region of validity, but we still believe the trends suggested by the calculations contain important information. In the CNS calculation the different kinds of triaxiality of the $`(fp)^{22}`$ bands with $`T`$=0 (positive $`\gamma `$) and $`T`$=1 (negative $`\gamma `$) implies about 50 percent larger $`B(E2)`$ values for the $`T`$=1 band compared to the $`T`$=0 band. In addition, $`|Q_{\mathrm{spec}}|`$ is calculated as much larger for states in the $`T`$=0 band than in the $`T`$=1 band. The increasing $`\gamma `$ deformation, finally leading to band termination of the $`(fp)^{22}`$ band at $`I`$=9 (fig.5), implies a gradual loss of collectivity and consequently decreasing $`B(E2)`$ values, while $`|Q_{\mathrm{spec}}|`$ reaches its maximum value. The different behavior of $`B(E2)`$ and $`Q_{\mathrm{spec}}`$ in the two bands can be understood from the $`\gamma `$-dependence of the two quadrupole moments which appear in eqs.(1): $`Q_2(\widehat{x})\mathrm{cos}(\gamma +30^{})`$ and $`Q_0(\widehat{x})\mathrm{sin}(\gamma +30^{})`$.
The $`B(E2)`$ values in the $`(fp)^{22}`$ bands are calculated as somewhat larger in the SM than in the CNS model. This might be explained by a lack of dynamical correlations in CNS. Quantum fluctuations around the equilibrium deformation, particularly in the $`\epsilon `$-direction, are expected to increase the $`B(E2)`$ values, see . The general trends of $`B(E2)`$ and $`Q_{\mathrm{spec}}`$ are, however, similar in the two models. The difference in $`B(E2)`$ values and $`Q_{\mathrm{spec}}`$ values between the $`T`$=0 and $`T`$=1 bands in the $`(fp)^{22}`$ configuration, as suggested from the different types of triaxiality obtained in the CNS calculation, is also seen in the SM calculation. In both models $`B(E2)`$ values are larger for the $`T`$=1 band than for the $`T`$=0 band (fig.8), while the opposite is true for $`|Q_{\mathrm{spec}}|`$ (fig.9). Although the effect is smaller in SM than in CNS, the CNS view of different types of triaxiality is supported by the SM results. Indeed, if $`\epsilon `$ and $`\gamma `$ deformations are deduced from SM calculated $`B(E2)`$ and $`Q_{\mathrm{spec}}`$ values through eqs.(1), both $`T`$=0 and $`T`$=1 bands are found to correspond to triaxial shapes ($`\gamma 0`$). Furthermore, the $`T`$=1 band corresponds to ($`\gamma \stackrel{<}{}\mathrm{\hspace{0.33em}0}`$), while states in the $`T`$=0 band has $`\gamma >0`$.
The decrease of $`B(E2)`$ with increasing spin, particularly seen in the $`T`$=0 band, is much less accentuated in the SM. $`E2`$ transitions along the band with the $`(fp)^{20}g_{9/2}^2`$ configuration, that is yrast above the band-crossing, are stronger than for the $`(fp)^{22}`$ band in both CNS and SM calculations. But, as discussed in section 3, also states within this configuration lose collectivity as angular momentum is increased, and in the CNS calculation the rotational band terminates at $`I^\pi `$=17<sup>+</sup> in an oblate shape with the symmetry axis coinciding with the rotation axis. The alignment sets in smoothly and the $`\gamma `$-deformation changes in a gradual way from $`8^{}`$ for the $`7^+`$ state to $`60^{}`$ for the $`17^+`$ state, see fig.5. This implies a gradual decrease of $`B(E2)`$ with increasing spin, as seen in fig.8. The decreasing behavior of $`B(E2)`$ in the band above the band-crossing also appears in the SM results, however, much less drastic than in CNS.
The restricted model space in the SM calculation, particularly the lack of holes in $`f_{7/2}`$, is more severe for the description of states in the $`(fp)^{20}g_{9/2}^2`$ configuration than in the $`(fp)^{22}`$ configuration, especially at higher spins, see Table 1. Effects from the restricted configuration space are seen in energies (see fig.2), but become much larger for the electromagnetic properties. Indeed, $`|Q_{\mathrm{spec}}|`$ is in general considerably smaller in SM than in CNS. States within the $`(fp)^{20}g_{9/2}^2`$ configuration are calculated in CNS to have $`\epsilon 0.24`$ (and $`\gamma `$ between $`10^{}`$ and $`60^{}`$) while corresponding $`\epsilon `$ values in the SM calculation, as deduced from eqs.(1), are considerably smaller, $`\epsilon 0.16`$. This is explained by the deformation driving effect of holes in $`f_{7/2}`$ that is then missing in SM wave functions.
## 6 Summary
The rotational behavior of low-lying $`T`$=0 and $`T`$=1 bands in the odd-odd $`N`$=$`Z`$ nucleus <sup>62</sup>Ga have been studied and compared in the cranked Nilsson-Strutinsky model and the spherical shell model. Both models have certain advantages and disadvantages that are analyzed in this comparative study. Their results were also compared with the observed $`T`$=0 band in <sup>62</sup>Ga as well $`T`$=1 band in <sup>62</sup>Zn.
CNS is an approximative mean field theory without dynamical correlations. In addition, pairing is neglected in the present study. On the other hand, CNS provides a good and intuitive picture of different kinds of collective phenomena, and works well in more or less all mass regions of nuclei, at all deformations, and in an extended region of angular momenta . The SM works very well in restricted areas where the valence space is not too large. The wave functions are extremely complex and the intuitive picture of collective phenomena is missing. By comparing SM and CNS results we have obtained an understanding of different collective phenomena in <sup>62</sup>Ga, such as rotation, backbending/band-crossing and band termination. Furthermore, we have studied in SM how isoscalar and isovector pairing affect $`T`$=0 and $`T`$=1 states differently with increasing spin.
Our main conclusions about <sup>62</sup>Ga properties are:
* The backbending seen in the $`T`$=0 band of <sup>62</sup>Ga is caused by an unpaired band-crossing between two bands with different configurations. Below the band-crossing all 22 particles outside the <sup>40</sup>Ca core occupy the $`fp`$ shell, and above the band-crossing one proton and one neutron are excited to $`g_{9/2}`$. This picture is supported both by CNS and SM calculations, and both models give energy states for the $`T`$=0 rotational band in good agreement with data, see fig.1.
* Backbending is predicted at $`I`$$``$10 in both $`T`$=0 and $`T`$=1 even spin bands.
* Pairing energy correlations in the (odd spin) $`T`$=0 band are less sensitive to increasing spin than in the (even spin) $`T`$=1 band, see fig.3. The effect is quite similar for isoscalar and isovector pairing, although isoscalar pairing is slightly more stable. Consequently, the moment of inertia is approximately constant for the $`T`$=0 band while it increases with increasing spin for the $`T`$=1 band. This results in a crossing between the $`T`$=0 band and $`T`$=1 band at $`I2`$ seen in fig.2.
* The blocking of the isovector pairing for $`T`$=0 states implies a smaller (negative) energy contribution from isovector pairing to the $`I`$=1 state ($`T`$=0), than to the $`I`$=0 state ($`T`$=1) (fig.3). This favors the $`I`$=0 ($`T`$=1) state that becomes the ground state.
* The excitation energy of the $`I`$=1 ($`T`$=0) state is well reproduced in the SM calculation.
* Energies of $`T`$=0 band with odd spins come out very similar in unpaired CNS calculation (lowest signature $`\alpha `$=1 band), and in SM calculation, when isoscalar and isovector pairing in the $`L`$=0 channel has been removed, see fig.6. In a similar way energies in the $`T`$=1 and $`T`$=0 bands with even spins come out similar to CNS band with $`\alpha `$=0.
* In CNS both bands involved in the band-crossing show band termination, at $`I`$=9 and $`I`$=17, respectively. At the band-crossing at $`I`$=9 the ground band terminates its rotational structure, and only one band continues after the band-crossing. The termination of the $`(fp)^{20}g_{9/2}^2`$ band at $`I=17`$ is smooth and favored.
* The size of the $`E2`$ transition matrix element decreases with increasing spin in both models, see fig.8, although the effect is weaker in SM.
* From CNS it is suggested that the $`T`$=0 and $`T`$=1 bands at low spins correspond to triaxial shapes with rotation around the intermediate and the smallest axes, respectively. This implies quite different electromagnetic properties for the two bands, see figs.8 and 9, and is partly supported by SM calculations.
* At higher spins the used SM configuration space becomes insufficient, and the need for involving holes in $`f_{7/2}`$ was stressed.
Acknowledgments
We thank Ingemar Ragnarsson and Dirk Rudolph for useful discussions and comments on the manuscript. We would like to thank Etienne Caurier for access to the shell model computer code . A. J. thanks the Swedish Institute (“The Visby programme”) for financial support, and S. Å. thanks the Swedish Research Council (NFR).
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# Four-Fermion Production in Electron-Positron Collisions
## 1 Introduction
During the year 1999 an informal workshop on Monte Carlo (MC) generators and programs took place at CERN, concentrating on processes in $`e^+e^{}`$ interactions at LEP 2 centre-of-mass energies (161 GeV to 210 GeV). One of the goals was to summarize and review critically the progress made in theoretical calculations and their implementation in computer programs since the 1995 workshop on Physics at LEP2. One of the reasons for this report was the need of having an official statement on various physics processes and the accuracy of their predictions, before deciding on LEP 2 activities in the year 2000.
This part of the workshop report summarizes the findings in the area of Four-Fermion final states. At the beginning of the workshop the following goals were identified for the Four-Fermion sub-group:
* Describe the new calculations and improvements in the theoretical understanding and in the upgraded MC implementations.
* Indicate where new contributions have changed previous predictions in the MC adopted by the collaborations, and specify why, how and by how much.
* In those cases where a substantial discrepancy has been registered and the physical origin has been understood, recommendations should be made on what to use.
* In those cases where we have found incompleteness of the existing MC, but no complete improvement is available, we should be able to indicate a sound estimate of the theoretical uncertainty, and possibly way and time scale for the solution.
Our strategy is determined by the physics issues arising in the experimental analyses performed at LEP 2. Therefore, the four LEP Collaborations have been asked to provide a list of relevant processes together with the level of theoretical accuracy needed.
Clearly, the LEP experiments investigate many different processes. For theoretical predictions we thus have to manage with lots of different sets of cuts. At the beginning of our activities the four experiments have presented us with lists that reflect rather diverse styles and different approaches: The complexity of the observables varied greatly, ranging from those defined by simple phase-space cuts on four-fermion (+ photon) level to complete event-selection procedures requiring parton shower and hadronization of quark systems.
An effort was made to settle as much as possible on a set of quasi-realistic but simple cuts for each process. We have collected processes and/or phase space regions where improved theoretical predictions are desirable. A weight has been assigned to each process according to its relevance and urgency.
The focus of activity has been on improving the theoretical predictions for the relevant processes and/or phase space regions. Also, all contributors have been asked to give an estimate for the remaining theoretical uncertainty. As a consequence, the output of the whole operation should not be a mere collection of comparison tables but a coherent attempt in assessing the theoretical uncertainty to be associated to any specific process.
The realm of theoretical uncertainty is ill defined and in order to reach a general consensus one cannot be satisfied with just some statement on the overall agreement among different programs. Whenever differences are found, one has to make sure that they are due to physics, and not to some different input. So our project had to foresee a preliminary phase with more of a technical benchmark. Once trivial discrepancies are understood and sorted out, one can start digging into inevitable differences arising from different implementations of common theoretical wisdom.
In a vast majority of cases the main theoretical problem is represented by the inclusion of QED radiation. Therefore, one of the main questions was: can we improve upon our treatment of QED radiation and/or give some safe estimate of the theoretical uncertainty associated with it?
Below we will present our reference table of four-fermion processes. It is an idealised common ground where, in principle, all theoretical predictions should be compared. More advanced setups would be accessible only to a more limited number of generators, built for that specific purpose.
It is useful to recall that the ultimate, perfect program does not exist and, most likely, will never exist. Roughly speaking, programs belong to two quite distinct classes. On one side there are event generators, usually interfaced with parton shower and hadronization packages. They may miss some fine points of the theoretical knowledge but represent an essential ingredient in the experimental analyses concerning the evaluation of signal efficiencies and backgrounds. Thus they create the necessary bridge between the raw data recorded by the detectors and the background-subtracted efficiency-corrected results published. At the other end of this cosmos we have semi-analytical programs that are not meant to generate events. Rather, they show their power in dealing with the signal, furnishing the implementation of (almost) everything available in the literature concerning the calculation of specific processes. In either case, we want to know about the theoretical uncertainty, process by process, to make clear which program is able to achieve that level of accuracy under which configuration. For $`W`$-pair production, however, the scenario is slightly changed: We have now MC event generators that, at the same time, represent a state-of-the-art calculation. Nevertheless, we do not have yet the ultimate MC: the one with radiative corrections, virtual/soft/hard photons, DPA, complete phase-space including single-$`W`$, single-$`Z`$, $`Z\gamma ^{}`$ and able to produce weight-1 events in finite time.
The results presented in this report are based on several different approaches and on comparisons of their numerical predictions. They are calculated with the following computer codes: BBC, CompHEP, GENTLE, grc4f, KORALW/YFSWW/YFSZZ, NEXTCALIBUR, PHEGAS/HELAC, RacoonWW, SWAP/WRAP, WPHACT and WTO/ZZTO.
This article is organised as follows. In Sect. 2 we present the four-fermion processes looked at in detail. Then we review the most recent theoretical developments in four-fermion physics in $`e^+e^{}`$ interactions. In Sect. 4 we discuss the CC03 $`\sigma _{WW}`$ cross-section and predictions based on the DPA. Here different approaches are compared. In Sect. 5 we discuss the radiative process with $`4\mathrm{f}+\gamma `$ final states. In Sect. 6 the single-$`W`$ production is critically discussed. Finally the NC02 cross-section, $`\sigma _{ZZ}`$ is analysed in Sect. 7 Conclusions and outlook are presented in Sect. 8
## 2 Four-fermion processes
Here we present our basic reference table and specify the calculational setup. One should read it as summarizing our original manifest. After reading the following sections, it will be instructive to come back here with a critical eye: not all the items and questions listed below have found a satisfactory answer. This was, somehow, foreseeable. If one thinks carefully one will easily discover some important message also for those issues that remain unsolved: they cannot be solved in any reasonable time scale and the associated effect is a real source of uncertainty.
### 2.1 List of processes
The following list provides the observables together with precision tags in $`\%`$, as requested by the experimental Collaborations. The accuracy of MC simulations should be better than the requested precision tag, i.e. the physics uncertainty should be smaller and at worst the one indicated. How much better is left to the contributors. For benchmarking it is certainly advisable to use the maximum available precision.
In general, radiative corrections and radiative photons in the final state should be considered for all processes, including the discussion of photon energy and polar-angle spectra. Typical minimal requirements on real photons are: energy $`E_\gamma >1`$GeV; polar angle $`|\mathrm{cos}\theta _\gamma |<0.985,0.997,0.9995`$ depending on channel; and minimal angle between photon and any charged final-state fermion $`\xi >5^{}`$.
* $`WW`$ and $`ZZ`$ type signal:
1. $`e^+e^{}WW`$ all (CC03). The full phase space is needed and the inclusive cross-section accuracy is $`0.2\%`$, which is $`1/3`$ of experimental accuracy combining all LEP 2 energies,
The spectrum for the photon energy and the polar angle is needed ($`|\mathrm{cos}\theta _\gamma |<0.985(0.997)`$).
2. $`e^+e^{}ZZ`$ all (NC02). The full phase space is needed and the inclusive cross-section accuracy is $`1\%`$. The spectrum for the photon energy and the polar angle is needed ($`|\mathrm{cos}\theta _\gamma |<0.985(0.997)`$).
3. $`e^+e^{}l\nu l\nu (\gamma )`$ where all $`\{e/\mu /\tau \}\{e/\mu /\tau \}`$ combinations are requested with the following conditions: ($`|\mathrm{cos}\theta _{l_1/l_2}|<0.985`$, $`E_{l_1/l_2}>5`$GeV, $`M(l^+l^{})>10(45)`$GeV (full and high-mass region). The inclusive cross-section accuracy is $`4\%`$ for individual combination; the inclusive cross-section accuracy is $`1\%`$ for the summed one; photon energy and polar angle spectrum ($`|\mathrm{cos}\theta _\gamma |<0.985(0.997)`$).
4. $`e^+e^{}\overline{q}qe\nu (\gamma )`$ (CC20), $`q`$-flavour blind, $`|\mathrm{cos}\theta _e|<0.985`$, $`E_e>5`$GeV, $`M(q\overline{q})>10(45)`$GeV (full and high-mass region); inclusive cross-section accuracy is $`1\%`$ (5% for low-mass region); photon energy and polar angle spectrum ($`|\mathrm{cos}\theta _\gamma |<0.985(0.997)`$).
5. $`e^+e^{}\overline{q}q\mu \nu (\gamma )`$ and $`e^+e^{}\overline{q}q\tau \nu (\gamma )`$ (incl. tau polarization in tau decay) (CC10), $`|\mathrm{cos}\theta _{\mu /\tau }|<0.985,E_{\mu /\tau }>5`$GeV, $`M(q\overline{q})>10(45)`$GeV (full and high-mass region), inclusive cross-section accuracy $`1\%`$. Photon energy and polar angle ($`|\mathrm{cos}\theta _\gamma |<0.985(0.997)`$) spectrum.
6. $`e^+e^{}q\overline{q}q\overline{q}(\gamma )`$, flavour blind, $`bbq\overline{q}`$, $`bbbb`$. At least two pairs with $`M(q_i,q_j)>10(45)`$GeV (full and high-mass region), inclusive cross-section accuracy $`1\%`$. photon energy and polar angle ($`|\mathrm{cos}\theta _\gamma |<0.985(0.997)`$) spectrum.
7. $`e^+e^{}q\overline{q}l^+l^{}(\gamma )`$, $`q`$-flavour blind, heavy $`q`$-flavors, $`l=e/\mu /\tau `$, $`|\mathrm{cos}\theta _{l_1}|<0.985`$, no cut on 2nd lepton (only one lepton tagged), $`M(q\overline{q})>10(45)`$GeV (full and high-mass region), inclusive cross-section accuracy $`2\%`$. Photon energy and polar angle ($`|\mathrm{cos}\theta _\gamma |<0.985(0.997)`$) spectrum.
8. $`e^+e^{}q\overline{q}l^+l^{}(\gamma )`$, $`q`$-flavour blind, heavy $`q`$-flavors, $`|\mathrm{cos}\theta _{l_1}|,|\mathrm{cos}\theta _{l_2}|<0.985`$ (both leptons tagged), full and high-mass regions: $`M(l^+l^{})>10(45)`$GeV, $`M(q\overline{q})>10(45)`$GeV, inclusive cross-section accuracy $`2\%`$. Photon energy and polar angle ($`|\mathrm{cos}\theta _\gamma |<0.985(0.997)`$) spectrum.
9. $`e^+e^{}q\overline{q}e^+e^{}(\gamma )`$, $`q`$-flavour blind, heavy $`q`$-flavors, with one electron in the beam pipe, $`|\mathrm{cos}\theta _e|>0.997`$, and one electron tagged, $`|\mathrm{cos}\theta _e|<0.985`$, $`M(q\overline{q})>10(45)`$GeV (full and high-mass region) . Photon energy and polar angle ($`|\mathrm{cos}\theta _\gamma |<0.985(0.997)`$) spectrum.
10. $`e^+e^{}q\overline{q}\nu \overline{\nu }(\gamma )`$, $`q`$-flavour blind, heavy $`q`$-flavors, $`M(q\overline{q})>10(45)`$GeV, inclusive cross-section accuracy $`4\%`$ (10% for low-mass region). Photon energy and polar angle ($`|\mathrm{cos}\theta _\gamma |<0.985(0.997)`$) spectrum.
11. $`e^+e^{}l^+l^{}L^+L^{}(\gamma )`$ and $`e^+e^{}l^+l^{}l^+l^{}(\gamma )`$ (all possible charged lepton flavour combinations): 3 or 4 leptons within acceptance $`|\mathrm{cos}\theta |<0.985`$, $`M(l^+l^{})`$ and $`M(L^+L^{})>10(45)`$GeV (full and high-mass region). Photon energy and polar angle ($`|\mathrm{cos}\theta _\gamma |<0.985(0.997)`$) spectrum.
* Single-$`W`$ type signal:
1. $`e^+e^{}q\overline{q}e\nu (\gamma )`$, $`|\mathrm{cos}\theta _e|>0.997`$, either $`M(q\overline{q})>45`$GeV or $`E_{q_1},E_{q_2}>15`$GeV, inclusive cross-section accuracy $`3\%`$, photon energy and polar angle ($`|\mathrm{cos}\theta _\gamma |<0.997(0.9995)`$) spectrum.
2. $`e^+e^{}e\nu e\nu (\gamma )`$, $`|\mathrm{cos}\theta _e|>0.997`$, $`E_e>15`$GeV, $`|\mathrm{cos}\theta _e|<0.7(0.95)`$, inclusive cross-section accuracy $`5\%`$, photon energy and polar angle ($`|\mathrm{cos}\theta _\gamma |<0.997(0.9995)`$) spectrum.
3. $`e^+e^{}e\nu \mu \nu (\gamma )`$ and $`e^+e^{}e\nu \tau \nu (\gamma )`$, $`|\mathrm{cos}\theta _e|>0.997`$, $`E_{\mu /\tau }>15`$GeV, $`|\mathrm{cos}\theta _{\mu /\tau }|<0.95`$, inclusive cross-section accuracy $`5\%`$, photon energy and polar angle ($`|\mathrm{cos}\theta _\gamma |<0.997(0.9995)`$) spectrum.
This list deserves already few words of comment.
For hadronic systems (CC or NC), there is usually a requirement of at least $`45`$ GeV invariant mass ($`W`$ and $`Z`$ signal) or at least $`10`$ GeV (background for other processes). Even lower invariant masses, say down to $`1`$GeV, should be handled by the dedicated $`\gamma \gamma `$ subgroup. For leptons, there should be no problem to go down to lower invariant masses or energies than listed above.
We consider as radiative events those events with real photons where at least one photon passes the photon requirements listed above, and as non-radiative events those with no photon or only photons below the minimal photon requirements. In case of non-radiative and radiative events, the cross section and its accuracy is needed. In case of non-radiative events, this amounts to adding up virtual and soft radiative corrections. In case of radiative events, some distributions are needed in addition, in particular photon energy and polar angle, and photon angle with respect to the nearest charged final-state fermion.
### 2.2 Questions to theory
We now elaborate in more detail on specific questions associated to specific processes.
* $`𝒪(\alpha )`$ electroweak corrections to $`e^+e^{}WW4\mathrm{f}`$.
Until 1999, the LEP experiments were using a $`2\%`$ theoretical uncertainty on the calculation of the CC03 $`W`$-pair cross section, not changed since the 1995 LEP 2 workshop. Although no complete one-loop $`𝒪\left(\alpha \right)`$ EW calculation exist yet for off-shell $`e^+e^{}WW4`$f production, we wish the theoretical uncertainty to be below $`1\%`$ ($`0.5\%`$ if possible) with justification. Also the uncertainties in CC03 vs. $`4\mathrm{f}`$ corrections when measuring the $`WW`$ cross section should be understood.
* Photon radiation (ISR) with $`p_t`$ in $`WW`$ and $`ZZ`$-dominated channels.
The principle effects will be on the selection efficiency and on the differential distributions used for W mass and triple gauge boson coupling (TGC) studies. The interest in photons is twofold: photons explicitly identified as such - usually at larger polar angles - and photons which simply create noticeable activity in the detector. The latter is, for example, also important in single-$`W`$ type analysis, therefore the photon angular range is extended to very low polar angles.
* Single $`W`$ channels.
For the single-$`W`$ process there are several issues to be addressed. In the region of high invariant masses of the $`W`$ boson (above $`45`$GeV) this process is important for both searches and TGC measurements. One topic of investigation should be ISR: this process is dominated by $`t`$-channel diagrams, whereas the current MC program implement ISR assuming s-channel reactions. A second issue is the treatment of the $`\alpha _{\mathrm{QED}}`$ scale, not only for single-$`W`$ but also for single-$`Z`$ and for $`Z\gamma ^{}`$. Is it better to re-weight on a event by event basis or on a diagram basis?
One of the outcomes of the workshop should be a recommendation on the mass cut which distinguishes the high mass region (more reliable) from the low mass region, i.e. the lower value to which the $`5\%`$ (or better) precision tag applies.
The importance of ISR in this channel is threefold: (a) change in total cross-section due to normal radiative corrections, (b) change in event distributions used to make cuts which changes the fraction of the total that fall inside our cuts, (c) fraction of events with identified photons - this forms a background to some of the chargino searches where a detected gamma is required.
Since the single-$`W`$ topology is defined as the one where a high mass object is found in the detector and the electron is not observed, we would like to know how the presence of $`p_t`$ ISR changes the fraction of events where the electron gets kicked out of the beam-pipe, how the differential distributions are distorted for TGC studies and what the explicit hard photon rate is.
The low mass region (below $`45`$ GeV) is mostly important for searches and studied within the $`\gamma \gamma `$ sub-group. One would like to trust the MC predictions down to $`510`$GeV invariant mass for the hadronic system. The required precision should also be around $`5`$ to $`10\%`$.
### 2.3 Input parameter set
A set of parameters must be specified for the calculation of $`𝒪\left(\alpha \right)`$ predictions (CC03 and to some extent also NC02). Once radiative corrections are included, the question of Renormalization Scheme (RS) and of Input Parameter Set (IPS) becomes relevant. For calculation, the following input parameters are used:
$`M__Z`$ $`=`$ $`91.1867\mathrm{GeV},1/\alpha (0)=137.0359895,`$
$`G_F`$ $`=`$ $`1.16637\times 10^5\mathrm{GeV}^2.`$ (1)
As far as masses are concerned one should use:
PDG values, i.e.
$`m_e`$ $`=`$ $`0.51099907\mathrm{MeV},m_\mu =105.658389\mathrm{MeV},`$
$`m_\tau `$ $`=`$ $`1.77705\mathrm{GeV}.`$ (2)
for light quarks one should make a distinction; for phase space:
$$m_u=5\mathrm{MeV},m_d=10\mathrm{MeV},\text{only relevant for single-}W,$$
(3)
while, in principle, these masses should not be used in deriving $`\alpha _{\mathrm{QED}}(s)`$ from $`\alpha _{\mathrm{QED}}(0)`$.
Here the recommendation follows the agreement in our community on using the following strategy for the evaluation of $`\alpha _{\mathrm{QED}}`$ at the mass of the $`Z`$. Define:
$$\alpha (M__Z)=\frac{\alpha (0)}{1\mathrm{\Delta }\alpha ^{(5)}(M__Z)\mathrm{\Delta }_{\mathrm{top}}(M__Z)\mathrm{\Delta }_{\mathrm{top}}^{\alpha \alpha __S}(M__Z)},$$
(4)
where one has $`\mathrm{\Delta }\alpha ^{(5)}(M__Z)=\mathrm{\Delta }\alpha _{\mathrm{lept}}+\mathrm{\Delta }\alpha _{\mathrm{had}}^{(5)}`$.
The input parameter should be $`\mathrm{\Delta }\alpha _{\mathrm{had}}^{(5)}`$, as it is the contribution with the largest uncertainty, while the calculation of the top contributions to $`\mathrm{\Delta }\alpha `$ is left for the code. This should become common to all codes. Codes should include, for $`\mathrm{\Delta }\alpha _{\mathrm{lept}}`$, the recently computed $`𝒪\left(\alpha ^3\right)`$ terms of and use as default $`\mathrm{\Delta }\alpha _{\mathrm{had}}^{(5)}=0.0280398`$, taken from . Using the default one obtains $`1/\alpha ^{(5)}(M__Z)=128.877`$, to which one must add the $`t\overline{t}`$ contribution and the $`𝒪\left(\alpha \alpha __S\right)`$ correction induced by the $`t\overline{t}`$ loop with gluon exchange, . Therefore, light quark masses should not appear in the evaluation of $`\alpha _{\mathrm{QED}}(M__Z)`$ and one should end up with:
$`1/\alpha (M__Z)`$ $`=`$ $`128.887,`$
for $`M__Z=91.1867GeV,m_t=175\mathrm{GeV},`$ (5)
$`M__H=150\mathrm{GeV},\alpha __S(M__Z)=0.119.`$
Furthermore, one should use:
$$\alpha __S(M__Z)=0.119,M__H=150\mathrm{GeV},M__W=80.350\mathrm{GeV}.$$
(6)
The quantities $`\mathrm{\Gamma }__Z,\mathrm{\Gamma }__W`$ should be understood as computed in the minimal standard model, e.g. $`\mathrm{\Gamma }__Z=2.49471\mathrm{GeV}`$ and $`\mathrm{\Gamma }__W=2.08699\mathrm{GeV}`$ for our IPS.
Now we come to the most important point, what to do with IPS in the presence of radiative corrections. In principle, all RS and all IPS are equally good and accepted, and differences are true estimates of some component of the theoretical uncertainty. However, we want to make sure that differences are not due to technical precision. The IPS that we want to specify is over-complete, let us repeat,
$`G_F`$ $`=`$ $`1.16637\times 10^5\mathrm{GeV}^2,1/\alpha (M__Z)=128.887,`$
$`M__Z`$ $`=`$ $`91.1867\mathrm{GeV},M__W=80.350\mathrm{GeV},`$
$`\alpha __S(M__Z)`$ $`=`$ $`0.119,M__H=150\mathrm{GeV}.`$ (7)
Clearly, once radiative correction are on, $`s_\theta =s_\theta ^{\mathrm{xxx}\mathrm{scheme}}`$ and we don’t care anymore since enough radiative corrections should be included to make all schemes equivalent to $`𝒪\left(\alpha \right)`$. Thus, for $`𝒪\left(\alpha \right)`$ numbers $`s_\theta `$ drops out. Perhaps we should give the highest marks to schemes where $`M__W`$ is in the IPS; after all, experiments measure $`M__W`$ at LEP 2 and any scheme where $`M__W`$ is not a primary quantity in the IPS is as bad as a scheme for LEP 1 where $`M__Z`$ is a derived quantity.
Nevertheless, we can use the over-completeness of the present IPS to set some internal consistency: it is a good idea to have an over-complete set of IPS, nevertheless consistent, so that everybody can make his favourite choice of the RS. Since we include values for $`\alpha (M__Z)`$ and for $`G_F`$ we can, as well, fine-tune the numbers so that the internal relations hold, to the best of our knowledge. The recommendation, in this case, is as follows:
* write down your favorite equation
$$f(M__Z,M__W,m_t,M__H,\alpha __S(M__Z),\alpha (0),G_F)=0,$$
(8)
* keep everything fixed but $`m_t`$ which, in turn, is derived as a solution of the consistency equation (for OMS this involves typically $`\mathrm{\Delta }r`$).
Even this solution is RS-dependent but variation should be minimal, sort of irrelevant. For instance, one could use the following result (derived from TOPAZ0 ):
$$m_t=174.17\mathrm{GeV}\text{Default for CC03}𝒪\left(\alpha \right).$$
(9)
With $`M__W=80.350GeV`$ and $`M__H=150`$GeV we are in a lucky situation, $`m_t`$ doesn’t change too much. For more solutions, we refer to Tab.(1).
### 2.4 Comparisons for $`4\mathrm{f}`$ results
There is an old tradition in LEP physics, new theoretical ideas and improvements should always be cross-checked before being adopted in the analysis of the experimental data. In this Report we present accurate and detailed comparisons between different generators. In most cases the authors have agreed to coordinate their action in understanding the features of the generators, their intrinsic differences and the goodness of their agreement or disagreement for the predictions. In so doing, and for the attuned comparisons, they can exclude that eventual disagreement may originate from trivial sources, like different input parameters.
Before entering into a detailed study of the numerical results it is important to underline how an estimate of the theoretical uncertainty emerges from the many sets of numbers obtained with the available generators. First of all one may distinguish between intrinsic and parametric uncertainties. The latter are normally associated with a variation of the input parameters according to the precision with which they are known. These uncertainties will eventually shrink when more accurate measurements will become available.
In this Report we are mainly devoted to a discussion of the intrinsic uncertainties associated with the choice of one scheme versus another. With one generator alone one cannot simulate the shift of a given quantity due to a change in the renormalization scheme. Thus the corresponding theoretical band in that quantity should be obtained from the differences in the prediction of the generators. On top of that we should also take into account the possibility of having different implementations of radiative corrections within one code. Many implementations of radiative corrections and of DPA are equally plausible and differ by non-leading higher order contributions, which however may become relevant in view of the achieved or projected experimental precision. This sort of intrinsic theoretical uncertainty can very well be estimated by staying within each single generator. However, since there are no reasons to expect that these will be the same in different generators, only the full collection of different sources will, in the end, give a reliable information on how accurate an observable may be considered from a theoretical point of view.
## 3 Phenomenology of unstable particles
In order to extract the $`WW`$ signal from the full set of $`\mathrm{e}^+\mathrm{e}^{}4f`$ processes, the CC03 cross-section was introduced and discussed in . In lowest order, this cross-section is simply based on the three $`WW`$ signal diagrams with the full four-particle kinematics with off-shell $`W`$ bosons. Compared to the full set of diagrams, the CC03 subset depends only trivially on the final state and allows to combine all channels easily. However, since the CC03 cross-section is based on a subset of diagrams, it is gauge-dependent and usually defined in the ’t Hooft–Feynman gauge. While the CC03 cross-section is not an observable, it is nevertheless a useful quantity at LEP 2 energies where it can be classified as a pseudo-observable. It contains the interesting physics, such as the non-abelian couplings and the sensitivity of the total cross section to $`M__W`$ near the $`W`$-pair threshold. The goal of this common definition is to be able to combine the different final state measurements from different experiments so that the new theoretical calculations can be checked with data at a level better than $`1\%`$. Note, however, that the CC03 cross-section will become very problematic at linear-collider energies, where the background diagrams and the gauge dependences are much larger.
It is worth summarizing the status of the $`WW`$ cross-section prior to the 2000 Winter Conferences. Nominally, any calculation for $`\mathrm{e}^+\mathrm{e}^{}\mathrm{WW}4f`$ was a tree level calculation and one could try the standard procedure of including, in a reasonable way, as much as possible of the universal corrections by constructing an improved Born approximation (hereafter IBA). This is the way the data have been analyzed so far, mostly with the help of GENTLE. Different programs have been compared for CC03, see Ref. : when one puts the same input parameters, renormalization scheme, etc, a technical agreement at the $`0.1\%`$ level is found. The universal corrections are not enough, since we wish the theoretical uncertainty to be below $`1\%`$ ($`0.5\%`$ seems possible) with justification.
Indeed, we have clear indications that non-universal electroweak corrections for $`WW`$(CC03) cross-section are not small and even larger than the experimental LEP accuracy. GENTLE will produce a CC03 cross-section, typically in the $`G_F`$-scheme, with universal ISR QED and non-universal ISR/FSR QED corrections, implemented with the so-called current-splitting technique. The corresponding curve has been used for the definition of the Standard Model prediction with a $`\pm 2\%`$ systematic error assigned to it. This error estimate is based on the knowledge of both leading and full $`𝒪\left(\alpha \right)`$ corrections to on-shell $`W`$-pair production. Note that, in GENTLE, the non-universal ISR correction with current-splitting technique reads as $`+0.4\%`$ effect at LEP 2 energies.
Recently, a new electroweak $`𝒪\left(\alpha \right)`$ CC03 cross-section has become available, in the framework of DPA, showing a result that is $`2.5÷3\%`$ smaller than the CC03 cross-section from GENTLE. This is a big effect since the combined experimental accuracy of LEP experiments is even smaller. It is, therefore, of the upmost importance to understand the structure of a DPA-corrected CC03 cross-section.
The double-pole approximation (DPA) of the lowest-order cross-section emerges from the CC03 diagrams upon projecting the $`W`$-boson momenta in the matrix element to their on-shell values. This means that the DPA is based on the residue of the double resonance, which is a gauge-invariant quantity, because it is directly related to the sub-processes of on-shell $`W`$-pair production and on-shell $`W`$ decay. In contrast to the CC03 cross-section, the DPA is theoretically well-defined. The price to be paid for this is the exclusion of the threshold region, where the DPA is not valid. On the other hand, the DPA provides a convenient framework for the inclusion of radiative corrections.
### 3.1 Dealing with unstable particles
Most of our technical problems originate from the complications naturally pertaining to the gauge structure of the theory and to the presence of unstable particles. As an interlude, we would like to summarize the nominal essence of the theoretical basis of all generators. In this respect one should remember that several, new, theoretical ideas were fully developed also as a consequence of the previous workshop on $`WW`$-physics (Physics at LEP2, Yellow report CERN/96-01, February 1996) and, in turn, many generators have profited from the most recent theoretical development. Furthermore, this Section will be a natural place where to add some consideration about the fine points in the DPA-procedure.
Four-fermion production processes, with or without radiative corrections, all involve fermions in the initial and final state and unstable gauge bosons as intermediate particles. Sometimes a photon is also present in the final state. If complete sets of graphs contributing to a given process are taken into account, the associated matrix elements are in principle gauge-invariant, i.e. they are independent of gauge fixing and respect Ward identities. This is, however, not guaranteed for incomplete sets of graphs like the ones corresponding to the off-shell $`W`$-pair production process (CC03). Indeed this process has been found to violate the $`SU(2)`$ Ward identities .
In addition, the unstable gauge bosons that appear as intermediate particles can give rise to poles $`1/(p^2M^2)`$ in physical observables if they are treated as stable particles. In view of the high precision of the LEP 2 experiments, the proper treatment of these unstable particles has become a demanding exercise, since on-shell approximations are simply not good enough anymore. A proper treatment of unstable particles requires the re-summation of the corresponding self-energies to all orders. In this way the singularities originating from the poles in the on-shell propagators are regularized by the imaginary parts contained in the self-energies, which are closely related to the decay widths ($`\mathrm{\Gamma }`$) of the unstable particles. The perturbative re-summation itself involves a simple geometric series and is therefore easy to perform. However, this simple procedure harbours the serious risk of breaking gauge invariance. Gauge invariance is guaranteed order by order in perturbation theory. Unfortunately one takes into account only part of the higher-order terms by re-summing the self-energies. This results in a mixing of different orders of perturbation theory and thereby jeopardizes gauge invariance, even if the self-energies themselves are extracted in a gauge-invariant way. Apart from being theoretically unacceptable, gauge-breaking effects can also lead to large errors in the MC predictions. At LEP 2 energies this problem occurs for instance in the reactions $`e^+e^{}e^{}\overline{\nu }_eu\overline{d},e^+\nu _e\overline{u}d`$ for forward-scattered beam particles .
Based on this observation, it is clear that a gauge-invariant scheme is required for the treatment of unstable particles. It should be stressed, however, that any such scheme is arbitrary to a greater or lesser extent: since the Dyson summation must necessarily be taken to all orders of perturbation theory, and we are not able to compute the complete set of all Feynman diagrams to all orders, the various schemes differ even if they lead to formally gauge-invariant results. Bearing this in mind, we need besides gauge invariance some physical motivation for choosing a particular scheme. In this context two options can be mentioned. Either one can try to subtract gauge-violating terms or one can try to add gauge-restoring terms to the calculation.
The first option is the so-called pole scheme . In this scheme one decomposes the complete amplitude by expanding around the poles. As the physically observable residues of the poles are gauge-invariant, gauge invariance is not broken if the finite width is taken into account in the pole terms $`1/(p^2M^2)`$. In reactions with multiple unstable-particle resonances it is rather awkward to perform the complete pole-scheme expansion with all its subtleties in the treatment of the mapping of the off-shell phase space on the on-shell phase space.
Therefore one usually approximates the expansion by retaining only the terms with the highest degree of resonance. This approximation is called the leading-pole approximation and is closely related to on-shell production and decay of the unstable particles. The accuracy of the approximation is typically $`𝒪\left(\mathrm{\Gamma }/M\right)`$, making it a suitable tool for calculating radiative corrections, since in that case the errors are further suppressed by powers of the coupling constant. Since diagrams with a lower degree of resonance do not feature in the leading-pole approximation, it is not an adequate approach for describing lowest-order reactions. So, for lowest-order reactions one needs an alternative approach.
The second option is based on the fundamentally different philosophy of trying to determine and include the minimal set of Feynman diagrams that is necessary for compensating the gauge violation caused by the self-energy graphs. This is obviously a theoretically very satisfying solution, but it may cause an increase in the complexity of the matrix elements and consequently a slowing down of the numerical calculations. Two methods have been developed along these lines.
First of all, for the gauge bosons we are guided by the observation that the lowest-order decay widths are exclusively given by the imaginary parts of the fermion loops in the one-loop self-energies. It is therefore natural to perform a Dyson summation of these fermionic one-loop self-energies and to include the other possible one-particle-irreducible fermionic one-loop corrections (fermion-loop scheme. For the lowest-order LEP 2 process $`e^+e^{}4\mathrm{f}`$ this amounts to adding the fermionic corrections to the triple gauge-boson vertex. The complete set of fermionic contributions forms a manifestly gauge-invariant subset, since it involves the closed subset of all $`𝒪\left([N_c^f\alpha /\pi ]^n\right)`$ contributions (with $`N_c^f`$ denoting the colour degeneracy of fermion $`f`$). Moreover, it obeys all Ward identities exactly, even with re-summed propagators, as shown in Ref. for two- and four-fermion production. For any particle reaction this can be deduced from the fact that the Ward identities of the underlying gauge symmetry, which are obeyed by the fermion loops, survive such a consistent Dyson summation, in contrast to the Slavnov–Taylor identities of the BRS symmetry, as shown in Ref. in the framework of the background-field formalism . The limitation of the fermion-loop scheme is due to the fact that it does not apply to particles with bosonic decay modes and that on resonance one perturbative order is lost. This in turn disqualifies it as a candidate for handling radiative corrections. Moreover, the inclusion of a full-fledged set of one-loop corrections in a lowest-order amplitude tends to over-complicate things for reactions like $`e^+e^{}4\mathrm{f}\gamma `$.
Recently a novel non-diagrammatic technique has been proposed for arbitrary tree-level reactions, involving all possible unstable particles and an unspecified amount of stable external particles . By using gauge-invariant non-local effective Lagrangians, it is possible to generate the self-energy effects in the propagators as well as the required gauge-restoring terms in the multi-particle (3-point, 4-point, etc.) interactions. In this way the full set of Ward identities can be solved, while keeping the gauge-restoring terms to a minimum.
A simplified version of this non-diagrammatic technique is the complex-mass scheme, which was introduced in Ref. for the reactions $`e^+e^{}4\mathrm{f}`$ and $`e^+e^{}4\mathrm{f}\gamma `$. In this scheme, the modifications of the vertices that are necessary to compensate the width effects of the propagators are obtained by analytically continuing the corresponding mass parameters in all Feynman rules consistently, leading to complex couplings. The complex-mass scheme preserves all Ward identities and works for arbitrary lowest-order predictions. As a small drawback we note, that for space-like gauge-boson momenta the propagators are complex in the complex-mass scheme, whereas perturbation theory in fact predicts the absence of any imaginary contribution to the propagator. This leads to complex couplings through gauge restoration and it will change, potentially, the CP structure of the theoretical predictions, whenever imaginary parts are redistributed between vertex functions.
We must admit that the effect on the CP structure has not been investigated in any scheme. However, for the Fermion-Loop scheme one does not see any problem with CP and for the non-local approach the modifications of the vertices have the feature that no imaginary parts are generated for space-like particles. One can also use the non-local approach starting from proper imaginary parts for time-like and unproper ones for space-like propagators and then look for a solution. One finds the complex mass scheme. As such it is confirmed by the non-local method, but only when one starts with an ad-hoc ansatz.
### 3.2 The leading-pole approximation
As mentioned above, the pole scheme consists in decomposing the complete amplitude by expanding around the poles of the unstable particles. The residues in this expansion are physically observable and therefore gauge-invariant. The pole-scheme expansion can be viewed as a gauge-invariant prescription for performing an expansion in powers of $`\mathrm{\Gamma }/M`$. It should be noted that there is no unique definition of the residues. Their calculation involves a mapping of off-shell matrix elements with off-shell kinematics on on-resonance matrix elements with restricted kinematics. This mapping, however, is not unambiguously fixed. After all, it involves more than just the invariant masses of the unstable particles and one thus has to specify the variables that have to be kept fixed in the mapping. The resulting implementation dependence manifests itself in differences of sub-leading nature, e.g. $`𝒪\left(\mathrm{\Gamma }/M\right)`$ suppressed deviations in the leading pole-scheme residue. In special regions of phase space, where the matrix elements vary rapidly, the implementation dependence can take noticeable proportions. This happens in particular near phase-space boundaries, like thresholds.
In order to make these statements a bit more transparent, we sketch the pole-scheme method for a single unstable particle. In this case the Dyson re-summed lowest-order matrix element is given by
$`^{\mathrm{}}`$ $`=`$ $`{\displaystyle \frac{W(p^2,\omega )}{p^2\stackrel{~}{M}^2}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{\stackrel{~}{\mathrm{\Sigma }}(p^2)}{p^2\stackrel{~}{M}^2}}\right)^n={\displaystyle \frac{W(p^2,\omega )}{p^2\stackrel{~}{M}^2+\stackrel{~}{\mathrm{\Sigma }}(p^2)}}`$ (10)
$`=`$ $`{\displaystyle \frac{W(M^2,\omega )}{p^2M^2}}{\displaystyle \frac{1}{Z(M^2)}}+\left[{\displaystyle \frac{W(p^2,\omega )}{p^2\stackrel{~}{M}^2+\stackrel{~}{\mathrm{\Sigma }}(p^2)}}{\displaystyle \frac{W(M^2,\omega )}{p^2M^2}}{\displaystyle \frac{1}{Z(M^2)}}\right],`$
where $`\stackrel{~}{\mathrm{\Sigma }}(p^2)`$ is the unrenormalized self-energy of the unstable particle with momentum $`p`$ and unrenormalized mass $`\stackrel{~}{M}`$. The renormalized quantity $`M^2`$ is the pole in the complex $`p^2`$ plane, whereas $`Z(M^2)`$ denotes the wave-function factor:
$$M^2\stackrel{~}{M}^2+\stackrel{~}{\mathrm{\Sigma }}(M^2)=0,Z(M^2)=1+\stackrel{~}{\mathrm{\Sigma }}^{}(M^2).$$
(11)
The first term in the last expression of Eq. (10) represents the single-pole residue, which is closely related to on-shell production and decay of the unstable particle. The second term between the square brackets has no pole and can be expanded in powers of $`p^2M^2`$. The argument $`\omega `$ denotes the dependence on the other variables, i.e. the implementation dependence. After all, the unstable particle is always accompanied by other particles in the production and decay stages.
For instance, consider the LEP1 reaction $`e^+e^{}\overline{f}f`$. In the mapping $`p_Z^2M^2`$ one can either keep $`t=(p_e^{}p_f)^2=p_Z^2(1\mathrm{cos}\theta )/2`$ fixed or $`\mathrm{cos}\theta `$. In the former mapping $`\mathrm{cos}\theta _{\text{pole}}`$ is obtained from the on-shell relation $`\mathrm{cos}\theta _{\text{pole}}=1+2t/M^2`$, whereas in the latter mapping $`t_{\text{pole}}=M^2(1\mathrm{cos}\theta )/2`$. It may be that a particular mapping leads to an unphysical point in the on-shell phase space. In the present example $`t_{\text{pole}}`$ will always be physical when $`\mathrm{cos}\theta `$ is kept fixed in the mapping. However, since $`|\mathrm{cos}\theta _{\text{pole}}|>1`$ for $`t<\mathrm{Re}M^2`$, it is clear that mappings with fixed Mandelstam variables harbour the potential risk of producing such unphysical phase-space points.<sup>1</sup><sup>1</sup>1In the resonance region, $`|p_Z^2M^2||M^2|`$, the unphysical on-shell phase-space points occur near the edge of the off-shell phase space, since $`t<\mathrm{Re}M^2`$ requires $`\mathrm{cos}\theta 1`$.
This can have repercussions on the convergence of the pole-scheme expansion. Therefore it is recommended to use implementations that are free of unphysical on-shell phase-space points.
The issue of taking angles instead of Mandelstam variables was raised in Ref. (see text after Eq.(58) there) and in the second reference of (see paragraph after Eq.(2)). For the DPA presented in Ref. , in discussing the treatment of the mapping of the off-shell phase space on the on-shell phase space, angles and completely decoupled off-shell invariant masses for the $`W`$ bosons were used. Finally, in Ref. the numerical effects coming from different phase-space treatments was considered also numerically. Specifically, the non-factorizable corrections were considered for different choices of Mandelstam variables used in the DPA.
The at present only workable approach for evaluating the radiative corrections to resonance-pair-production processes, like $`W`$-pair production, involves the so-called leading-pole approximation (LPA). This approximation restricts the complete pole-scheme expansion to the term with the highest degree of resonance. In the case of $`W`$-pair production only the double-pole residues are hence considered. This is usually referred to as the DPA. The intrinsic error associated with this procedure is $`\alpha \mathrm{\Gamma }_W/(\pi M__W)\times \mathrm{ln}(\mathrm{})\stackrel{<}{}0.5\%`$, except far off resonance, where the pole-scheme expansion cannot be viewed as an effective expansion in powers of $`\mathrm{\Gamma }/M`$, and close to phase-space boundaries, where the DPA cannot be trusted to produce the dominant contributions. In the above error estimate, the $`\mathrm{ln}(\mathrm{})`$ represents leading logarithms or other possible enhancement factors in the corrections. In the latter situations also the implementation dependence of the double-pole residues can lead to enhanced errors. Close to the nominal (on-shell) $`W`$-pair threshold, for instance, the intrinsic error is effectively enhanced by a factor $`M__W/(\sqrt{s}2M__W)M__W/\mathrm{\Delta }E`$. In view of this it is wise to apply the DPA only if the energy is several $`\mathrm{\Gamma }_W`$ above the threshold.
In the DPA one can identify two types of contributions. One type comprises all diagrams that are strictly reducible at both unstable $`W`$-boson lines (see Fig. 1). These corrections are therefore called factorizable and can be attributed unambiguously either to the production of the $`W`$-boson pair or to one of the subsequent decays. The second type consists of all diagrams in which the production and/or decay sub-processes are not independent and which therefore do not seem to have two overall $`W`$ propagators as factors (see Fig. 2). We refer to these effects as non-factorizable corrections.<sup>2</sup><sup>2</sup>2It should be noted that the exact split-up between factorizable and non-factorizable radiative corrections requires a precise (gauge-invariant) definition. We will come back to this point.
In the DPA the non-factorizable corrections arise exclusively from the exchange or emission of photons with $`E_\gamma \stackrel{<}{}\mathrm{\Gamma }_W`$ . Hard photons as well as massive-particle exchanges do not lead to double-resonant contributions. The physical picture behind all of this is that in the DPA the $`W`$-pair process can be viewed as consisting of several sub-processes: the production of the $`W`$-boson pair, the propagation of the $`W`$ bosons, and the subsequent decay of the unstable $`W`$ bosons. The production and decay are hard sub-processes, which occur on a relatively short time interval, $`𝒪\left(1/M__W\right)`$. They are in general distinguishable as they are well separated by a relatively big propagation interval, $`𝒪\left(1/\mathrm{\Gamma }_W\right)`$. Consequently, the corresponding amplitudes have certain factorization properties. The same holds for the radiative corrections to the sub-processes. The only way the various stages can be interconnected is via the radiation of soft photons with energy of $`𝒪\left(\mathrm{\Gamma }_W\right)`$.
As is clear from the above-given discussion of the DPA, a specific prescription has to be given for the calculation of the DPA residues. Or, in other words, we have to fix the implementation of the mapping of the full off-shell phase space on the kinematically restricted (on-resonance) one. Two strategies have been adopted in the literature . One can opt to always extract pure double-pole residues . This means in particular that after the integration over decay kinematics and invariant masses has been performed the on-shell cross-section should be recovered. Alternatively, one can decide to exclude the off-shell phase space from the mapping and apply the residue only to the matrix elements . We will come back to the conceptual and numerical differences between these two implementation strategies in the detailed discussion of the DPA programs. At this point we merely note that the numerical differences can serve as an estimate of the theoretical uncertainty of the DPA procedure. Ref. also used the approach in which the full off-shell phase space is maintained and the residue is only applied to the matrix elements.
In the rest of this section we will explain those aspects of the DPA procedure that are common to both implementation methods. To this end we focus on the lowest-order reaction
$$e^+(q_1)e^{}(q_2)W^+(p_1)W^{}(p_2)\overline{f}_1(k_1)f_1^{}(k_1^{})f_2(k_2)\overline{f}_2^{}(k_2^{}),$$
(12)
involving only those diagrams that contain as factors the Breit–Wigner propagators for the $`W^+`$ and $`W^{}`$ bosons. Here $`\overline{f}_1`$ and $`f_1^{}`$ are the decay products of the $`W^+`$ boson, and $`f_2`$ and $`\overline{f}_2^{}`$ those of the $`W^{}`$ boson. It should be noted that a large part of the radiative corrections in DPA to this reaction can be treated in a way similar to the lowest-order case, which is therefore a good starting point. The amplitude for process (12) takes the form
$$=\underset{\lambda _1,\lambda _2}{}\mathrm{\Pi }_{\lambda _1\lambda _2}(M_1,M_2)\frac{\mathrm{\Delta }_{\lambda _1}^{(+)}(M_1)}{D_1}\frac{\mathrm{\Delta }_{\lambda _2}^{()}(M_2)}{D_2},$$
(13)
where any dependence on the helicities of the initial- and final-state fermions has been suppressed, and
$$D_i=M_i^2M__W^2+iM__W\mathrm{\Gamma }_W,M_i^2=(k_i+k_i^{})^2.$$
(14)
The quantities $`\mathrm{\Delta }_{\lambda _1}^{(+)}(M_1)`$ and $`\mathrm{\Delta }_{\lambda _2}^{()}(M_2)`$ are the off-shell $`W`$-decay amplitudes for specific spin-polarization states $`\lambda _1`$ (for the $`W^+`$) and $`\lambda _2`$ (for the $`W^{}`$), with $`\lambda _i=(1,0,+1)`$. The off-shell $`W`$-pair production amplitude $`\mathrm{\Pi }_{\lambda _1\lambda _2}(M_1,M_2)`$ depends on the invariant fermion-pair masses $`M_i`$ and on the polarizations $`\lambda _i`$ of the virtual $`W`$ bosons. In the limit $`M_iM__W`$ the amplitudes $`\mathrm{\Pi }`$ and $`\mathrm{\Delta }^{(\pm )}`$ go over into the on-shell production and decay amplitudes.
In the cross-section the above factorization leads to
$$\underset{\text{fermion helicities}}{}||^2=\underset{\lambda _1,\lambda _2,\lambda _1^{},\lambda _2^{}}{}𝒫_{[\lambda _1\lambda _2][\lambda _1^{}\lambda _2^{}]}(M_1,M_2)\frac{𝒟_{\lambda _1\lambda _1^{}}(M_1)}{|D_1|^2}\frac{𝒟_{\lambda _2\lambda _2^{}}(M_2)}{|D_2|^2}.$$
(15)
In Eq. (15) the production part is given by a $`9\times 9`$ density matrix
$$𝒫_{[\lambda _1\lambda _2][\lambda _1^{}\lambda _2^{}]}(M_1,M_2)=\underset{e^\pm \text{ helicities}}{}\mathrm{\Pi }_{\lambda _1\lambda _2}(M_1,M_2)\mathrm{\Pi }_{\lambda _1^{}\lambda _2^{}}^{}(M_1,M_2).$$
(16)
Similarly the decay part is governed by $`3\times 3`$ density matrices
$$𝒟_{\lambda _i\lambda _i^{}}(M_i)=\underset{\text{fermion helicities}}{}\mathrm{\Delta }_{\lambda _i}(M_i)\mathrm{\Delta }_{\lambda _i^{}}^{}(M_i),$$
(17)
where the summation is performed over the helicities of the final-state fermions.
It is clear that Eq. (16) is closely related to the absolute square of the matrix element for stable unpolarized $`W`$-pair production. In that case the cross-section contains the trace of the above density matrix
$$\text{Tr}𝒫(M__W,M__W)=\underset{\lambda _1,\lambda _2}{}𝒫_{[\lambda _1\lambda _2][\lambda _1\lambda _2]}(M__W,M__W)=\underset{\text{all polarizations}}{}|\mathrm{\Pi }_{\lambda _1\lambda _2}(M__W,M__W)|^2.$$
(18)
The decay of an unpolarized on-shell $`W`$ boson is determined by
$$\text{Tr}𝒟(M__W)=\underset{\lambda _i}{}𝒟_{\lambda _i\lambda _i}(M__W)=\underset{\text{all polarizations}}{}|\mathrm{\Delta }_{\lambda _i}(M__W)|^2.$$
(19)
Note, however, that also the off-diagonal elements of $`𝒫(M__W,M__W)`$ and $`𝒟(M__W)`$ are required for determining Eq. (15) in the limit $`M_iM__W`$.
As a next step it is useful to describe the kinematics of process (12) in a factorized way, i.e. using the invariant masses $`M_1`$ and $`M_2`$ of the fermion pairs. The differential cross-section takes the form
$$d\sigma =\frac{1}{2s}||^2d\mathrm{\Gamma }_{4\mathrm{f}}=\frac{1}{2s}||^2d\mathrm{\Gamma }_{\text{pr}}d\mathrm{\Gamma }_{\text{dec}}^+d\mathrm{\Gamma }_{\text{dec}}^{}\frac{dM_1^2}{2\pi }\frac{dM_2^2}{2\pi },$$
(20)
where $`d\mathrm{\Gamma }_{4\mathrm{f}}`$ indicates the complete four-fermion phase-space factor and $`s=(q_1+q_2)^2`$ the centre-of-mass energy squared. The phase-space factors for the production and decay sub-processes, $`d\mathrm{\Gamma }_{\text{pr}}`$ and $`d\mathrm{\Gamma }_{\text{dec}}^\pm `$, read
$`d\mathrm{\Gamma }_{\text{pr}}`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^2}}\delta (q_1+q_2p_1p_2){\displaystyle \frac{d\stackrel{}{p}_1}{2p_{10}}}{\displaystyle \frac{d\stackrel{}{p}_2}{2p_{20}}},`$
$`d\mathrm{\Gamma }_{\text{dec}}^+`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^2}}\delta (p_1k_1k_1^{}){\displaystyle \frac{d\stackrel{}{k}_1}{2k_{10}}}{\displaystyle \frac{d\stackrel{}{k}_1^{}}{2k_{10}^{}}},`$
$`d\mathrm{\Gamma }_{\text{dec}}^{}`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^2}}\delta (p_2k_2k_2^{}){\displaystyle \frac{d\stackrel{}{k}_2}{2k_{20}}}{\displaystyle \frac{d\stackrel{}{k}_2^{}}{2k_{20}^{}}}.`$ (21)
When the factorized form for $`||^2`$ is inserted one obtains
$`d\sigma `$ $`=`$ $`{\displaystyle \frac{1}{2s}}{\displaystyle \underset{\lambda _1,\lambda _2,\lambda _1^{},\lambda _2^{}}{}}𝒫_{[\lambda _1\lambda _2][\lambda _1^{}\lambda _2^{}]}(M_1,M_2)d\mathrm{\Gamma }_{\text{pr}}\times 𝒟_{\lambda _1\lambda _1^{}}(M_1)d\mathrm{\Gamma }_{\text{dec}}^+\times 𝒟_{\lambda _2\lambda _2^{}}(M_2)d\mathrm{\Gamma }_{\text{dec}}^{}\times `$ (22)
$`\times {\displaystyle \frac{1}{2\pi }}{\displaystyle \frac{dM_1^2}{|D_1|^2}}\times {\displaystyle \frac{1}{2\pi }}{\displaystyle \frac{dM_2^2}{|D_2|^2}},`$
which is the common starting point for any of the DPA implementations.
### 3.3 Radiative corrections in double-pole approximation
A full calculation of the complete electroweak $`𝒪\left(\alpha \right)`$ corrections to $`\mathrm{e}^+\mathrm{e}^{}4f(+\gamma )`$ for all four-fermion final states is beyond present possibilities. While the real bremsstrahlung corrections induced by $`\mathrm{e}^+\mathrm{e}^{}4f\gamma `$ are known exactly , there are severe technical and conceptual problems with the virtual corrections to four-fermion production. Fortunately, the full account of the $`𝒪\left(\alpha \right)`$ corrections is not needed at the level of accuracy demanded by LEP 2. For $`W`$-pair-mediated processes, $`\mathrm{e}^+\mathrm{e}^{}\mathrm{WW}4f`$, the required accuracy of predictions is of the order of $`0.5\%`$ for integrated quantities. At this level, the corrections to $`W`$-pair production can be treated in the DPA. In regions of phase space where two resonant $`W`$ bosons do not dominate the cross-sections, such as in the $`WW`$-threshold region or in the single-$`W`$ domain, the DPA is, of course, not valid and one should resort to other approximations as the Weizsäcker-Williams for single-$`W`$ .
Since only diagrams with two nearly resonant $`W`$ bosons are relevant for the DPA, the number of graphs is reduced considerably, and a generic treatment of all four-fermion final states is possible. Obviously all diagrams that appear for the pair production and the decay of on-shell $`W`$ bosons are also relevant for the pole expansion in the DPA. Since such contributions involve a product of two independent Breit–Wigner factors for the $`W`$ resonances, they are called factorizable corrections. However, there exist also doubly-resonant corrections in which the production and decay sub-processes do not proceed independently. Power counting reveals that such corrections are only doubly-resonant if the particle that is exchanged by the sub-processes is a low-energetic photon. Owing to the complicated off-shell behaviour of these corrections, they are called non-factorizable.
While the definition of the DPA is straightforward for the virtual corrections, it is problematic for the real corrections. The problem is due to the momentum carried away by photon radiation. The invariant masses of the $`W`$ bosons in contributions in which the photon is emitted in the $`W`$-pair production subprocess differ from those where the photon is emitted in the $`W`$-decay sub-processes. The corresponding Breit–Wigner resonances overlap if the energy of the emitted photon is of the order of $`\mathrm{\Gamma }_\mathrm{W}`$. It is not obvious how to define the DPA for such photons. Therefore, the results based on a DPA for the real corrections have to be treated with some caution.
According to the above classification, there are four categories of contributions to $`𝒪\left(\alpha \right)`$ corrections in DPA: factorizable and non-factorizable ones both for virtual and real corrections. In the following the salient features of those four parts are described.
#### 3.3.1 Virtual corrections
As a first step we discuss how to separate the virtual corrections into a sum of factorizable and non-factorizable virtual corrections. The diagrammatic split-up according to reducible and irreducible $`W`$-boson lines is an illustrative way of understanding the different nature of the two classes of corrections, but since the double-resonant diagrams are not gauge-invariant by themselves the precise split-up needs to be defined properly.
We can make use of the fact that there are effectively two scales in the problem: $`M__W`$ and $`\mathrm{\Gamma }_W`$. Let us now consider virtual corrections coming from photons with different energies:
* soft photons, $`E_\gamma \mathrm{\Gamma }_W`$,
* semi-soft photons, $`E_\gamma =𝒪\left(\mathrm{\Gamma }_W\right)`$,
* hard photons, $`\mathrm{\Gamma }_WE_\gamma =𝒪\left(M__W\right)`$.
Only soft and semi-soft photons contribute to both factorizable and non-factorizable corrections. The latter being defined to describe interactions between different stages of the off-shell process. The reason for this is that only these photons can induce relatively long-range interactions and thereby allow the various sub-processes, which are separated by a propagation interval of $`𝒪\left(1/\mathrm{\Gamma }_W\right)`$, to communicate with each other. Virtual corrections involving the exchange of hard photons or of massive particles contribute exclusively to the factorizable corrections. In view of the short range of the interactions induced by these particles, their contribution to the non-factorizable corrections are suppressed by at least $`𝒪\left(\mathrm{\Gamma }_W/M__W\right)`$.
As hard photons contribute to the factorizable corrections only, we merely need to define a split-up for soft and semi-soft photons. It is impossible to do this in a consistent gauge-invariant way on the basis of diagrams. In Refs. it was shown that only part of particular diagrams should be attributed to the non-factorizable corrections, the rest being of factorizable nature. The complete set of non-factorizable corrections was obtained by collecting all terms that contain the ratios $`D_i/[D_i\pm 2kp_i]`$, where $`k`$ denotes the momentum of the (semi-)soft photon. The so-defined non-factorizable corrections read
$$_{\text{nf}}^{\text{virt}}=i_0^{\text{DPA}}\frac{d^4k}{(2\pi )^4[k^2+io]}\left[\left(𝒥_0^\mu +𝒥_{}^\mu \right)𝒥_{+,\mu }+\left(𝒥_0^\mu +𝒥_{}^\mu \right)𝒥_{,\mu }+𝒥_+^\mu 𝒥_{,\mu }\right],$$
(23)
which contains the gauge-invariant currents
$$𝒥_0^\mu =e\left[\frac{p_1^\mu }{kp_1+io}+\frac{p_2^\mu }{kp_2+io}\right],$$
$$𝒥_{}^\mu =e\left[\frac{q_1^\mu }{kq_1+io}\frac{q_2^\mu }{kq_2+io}\right],𝒥_{}^\mu =+e\left[\frac{q_1^\mu }{kq_1+io}\frac{q_2^\mu }{kq_2+io}\right]$$
(24)
for photon emission from the production stage of the process, and
$`𝒥_+^\mu `$ $`=`$ $`e\left[{\displaystyle \frac{p_1^\mu }{kp_1+io}}+Q_{f_1}{\displaystyle \frac{k_1^\mu }{kk_1+io}}Q_{f_1^{}}{\displaystyle \frac{k_{1}^{}{}_{}{}^{\mu }}{kk_1^{}+io}}\right]{\displaystyle \frac{D_1}{D_1+2kp_1}},`$
$`𝒥_{}^\mu `$ $`=`$ $`e\left[{\displaystyle \frac{p_2^\mu }{kp_2+io}}+Q_{f_2}{\displaystyle \frac{k_2^\mu }{kk_2+io}}Q_{f_2^{}}{\displaystyle \frac{k_{2}^{}{}_{}{}^{\mu }}{kk_2^{}+io}}\right]{\displaystyle \frac{D_2}{D_22kp_2}}`$ (25)
for photon emission from the decay stages of the process. Here $`_0^{\text{DPA}}`$ is the lowest-order matrix element in DPA and $`Q_f`$ stands for the charge of fermion $`f`$ in units of $`e`$. Since Eq. (23) contains (at least) all contributions from diagrams with irreducible $`W`$-boson lines, it can be viewed as a gauge-invariant extension of the set of $`W`$-irreducible diagrams. In general one has to calculate all of the integrals appearing in the above expressions. The complete set of integrals has been given in Ref. and explicit expressions for the full set of virtual factorizable corrections can be found in . However, if one is interested in the sum of virtual corrections and real-photon radiation, then some simplifications occur depending on the treatment of the photon<sup>3</sup><sup>3</sup>3Note that Eq. (23) is UV-finite and contains $`4`$\- and $`5`$-point integrals. In fact it was observed that certain combinations of these $`4`$\- and $`5`$-point integrals are equal to a simple (Coulomb-like) $`3`$-point integral plus a constant. This simple $`3`$-point integral has an artificial UV divergence, which cancels against the constant and can be regulated by either a cut-off (BBC) or by keeping the DPA-subleading $`k^2`$ contributions in the denominators (RACOONWW). The final answer of course does not depend on this..
If the radiated (real) photon is treated inclusively, then many of the terms in Eq. (23) cancel . In this context the difference in the signs of the $`io`$ parts appearing in the currents $`𝒥_{}`$ and $`𝒥_{}`$ are crucial. These signs actually determine which interference terms give rise to a non-vanishing non-factorizable contribution after virtual and real-photon corrections have been added. As a result of such considerations only a very limited subset of ‘final-state’ interferences survives for inclusive photons: the virtual corrections corresponding to Figs. 2 and 3 as well as the associated real-photon corrections.
The sum of virtual and real non-factorizable corrections has been calculated, Refs. <sup>4</sup><sup>4</sup>4The original result of the older calculation does not agree with the two more recent results , which are in mutual agreement. As known from the authors of Ref. , their corrected results also agree with the ones of Refs. .. It has been shown in Ref. that this sum vanishes if the invariant masses of both $`W`$ bosons are integrated over, i.e. in particular that the full non-factorizable correction to the total cross-section is zero in DPA.
In Refs. the full non-factorizable corrections have also been discussed numerically. They vanish on top of the double resonance and are of the order of $`1\%`$ in its vicinity. The shift in the $`W`$ invariant-mass distributions is only of the order of a few MeV. These results can be reproduced by a simple approximation based on the so-called screened Coulomb ansatz. However, it is important to note that all these numerical results on non-factorizable corrections are based on the DPA for real corrections and have been obtained in idealized treatment of phase space, namely the assumption that the $`W`$-boson momenta can be reconstructed from the fermion momenta alone, i.e. without photon recombination. It is not clear how these results change in physical situations with photon recombination.
The virtual factorizable corrections consist of all hard contributions and the left-over part of the semi-soft ones. The so-defined factorizable corrections have the nice feature that they can be expressed in terms of corrections to on-shell sub-processes, i.e. the production of two on-shell W bosons and their subsequent on-shell decays. The corresponding matrix element can be expressed in the same way as described at lowest-order:
$`_{\text{fact}}^{\text{virt}}={\displaystyle \underset{\lambda _1,\lambda _2}{}}\mathrm{\Pi }_{\lambda _1\lambda _2}(M_1,M_2){\displaystyle \frac{\mathrm{\Delta }_{\lambda _1}^{(+)}(M_1)}{D_1}}{\displaystyle \frac{\mathrm{\Delta }_{\lambda _2}^{()}(M_2)}{D_2}}.`$ (26)
Here two of the amplitudes are taken at lowest order, whereas the remaining one contains all possible one-loop contributions, including the $`W`$ wave-function factors that appear in Eq. (10). In this way the well-known on-shell radiative corrections to the production and decay of pairs of $`W`$ bosons appear as basic building blocks of the factorizable corrections.<sup>5</sup><sup>5</sup>5Note that the complete density matrix is required in this case, in contrast to the pure on-shell calculation which involves the diagonal elements of the density matrix only. In the semi-soft limit the photonic virtual factorizable corrections to the production stage, contained in $`\mathrm{\Pi }`$, cancel against the corresponding real-photon corrections. Non-vanishing contributions from $`\mathrm{\Pi }`$ occur as soon as the $`k^2`$ terms in the propagators cannot be neglected anymore. An example of this is the factorizable correction from the Coulomb graph Fig. 3. For the on-shell (factorizable) part of the Coulomb effect photons with momenta $`k_0=𝒪\left(\mathrm{\Delta }E\right)`$ and $`|\stackrel{}{k}|=𝒪\left(\sqrt{M__W\mathrm{\Delta }E}\right)`$ are important , i.e. $`k^2`$ cannot be neglected in the propagators of the unstable particles. Since we stay well away from the $`W`$-pair threshold ($`\mathrm{\Delta }E\sqrt{s}2M__W\mathrm{\Gamma }_W`$), this situation occurs outside the realm of the semi-soft photons. This fits nicely into the picture of the production stage being a hard subprocess, governed by relatively short time scales as compared with the much longer time scales required for the non-factorizable corrections, which interconnect the different sub-processes.
#### 3.3.2 Real-photon radiation
In this subsection we discuss the aspects of real-photon radiation in the DPA as used in . To this end we consider the process
$$e^+(q_1)e^{}(q_2)W^+(p_1)W^{}(p_2)\left[\gamma (k)\right]\overline{f}_1(k_1)f_1^{}(k_1^{})f_2(k_2)\overline{f}_2^{}(k_2^{})\gamma (k),$$
(27)
where in the intermediate state there may or may not be a photon. We will show how to extract the gauge-invariant double-pole residues in different situations. The exact cross-section for process (27) can be written in the following form
$$d\sigma =\frac{1}{2s}|_\gamma |^2d\mathrm{\Gamma }_{4\mathrm{f}\gamma }=\frac{1}{2s}\left[2\mathrm{R}\mathrm{e}\left(_0_+^{}+_0_{}^{}+_+_{}^{}\right)+|_0|^2+|_+|^2+|_{}|^2\right]d\mathrm{\Gamma }_{4\mathrm{f}\gamma },$$
(28)
where $`d\mathrm{\Gamma }_{4\mathrm{f}\gamma }`$ indicates the complete five-particle phase-space factor, and the matrix elements $`_0`$ and $`_\pm `$ correspond to the diagrams where the photon is attached to the production or decay stage of the three $`W`$-pair diagrams, respectively. This split-up can be achieved with the help of the partial-fraction decomposition
$$\frac{1}{D_i(D_i+2pk)}=\frac{1}{2pk}\left(\frac{1}{D_i}\frac{1}{D_i+2pk}\right).$$
(29)
Each contribution to the cross-section can be written in terms of polarization density matrices, which originate from the amplitudes
$$_0=\mathrm{\Pi }_\gamma (M_1,M_2)\frac{\mathrm{\Delta }^{(+)}(M_1)}{D_1}\frac{\mathrm{\Delta }^{()}(M_2)}{D_2},$$
(30)
$$_+=\mathrm{\Pi }(M_{1\gamma },M_2)\frac{\mathrm{\Delta }_\gamma ^{(+)}(M_{1\gamma })}{D_{1\gamma }}\frac{\mathrm{\Delta }^{()}(M_2)}{D_2},$$
(31)
$$_{}=\mathrm{\Pi }(M_1,M_{2\gamma })\frac{\mathrm{\Delta }^{(+)}(M_1)}{D_1}\frac{\mathrm{\Delta }_\gamma ^{()}(M_{2\gamma })}{D_{2\gamma }},$$
(32)
where all polarization indices for the W bosons and the photon have been suppressed, and
$$D_{i\gamma }=D_i+2kk_i+2kk_i^{},M_{i\gamma }^2=M_i^2+2kk_i+2kk_i^{},M_i^2=(k_i+k_i^{})^2.$$
(33)
The matrix elements $`\mathrm{\Pi }_\gamma `$ and $`\mathrm{\Delta }_\gamma ^{(\pm )}`$ describe the production and decay of the $`W`$ bosons accompanied by the radiation of a photon. The matrix elements without subscript $`\gamma `$ have been introduced in Eq. (13).
In the calculation of the Born matrix element and virtual corrections only two poles could be identified in the amplitudes, originating from the Breit–Wigner propagators $`1/D_i`$. The pole-scheme expansion was performed around these two poles. In contrast, the bremsstrahlung matrix element has four in general different poles, originating from the four Breit–Wigner propagators $`1/D_i`$ and $`1/D_{i\gamma }`$. As mentioned above, the matrix element can be rewritten as a sum of three matrix elements ($`_0,_+,_{}`$), each of which only contain two Breit–Wigner propagators. For these three individual matrix elements the pole-scheme expansion is fixed, as before, to an expansion around the corresponding two poles. However, when calculating cross-sections \[see Eq. (28)\] the mapping of the five-particle phase space introduces a new type of ambiguity. The interference terms in Eq. (28) involve two different double-pole expansions simultaneously. One might think this will pose a problem, since there is no natural choice for the phase-space mapping in those cases. As we will see later, however, only photons with $`E_\gamma \stackrel{<}{}\mathrm{\Gamma }_WM_W`$ give noticeable contributions to these interference terms. This means that one can apply a soft-photon-like (semi-soft) approximation (see below).
In Ref. it was argued that the resulting ambiguity in the phase-space mapping will not have significant repercussions on the quality of the DPA calculation, in the same way as stable-particle calculations are not significantly affected by the photon momentum in the soft-photon regime. We note, however, that there is still some controversy on this issue.
Let us return now to the three earlier-defined regimes for the photon energy:
* for hard photons \[$`E_\gamma \mathrm{\Gamma }_W`$\] the Breit–Wigner poles of the W-boson resonances before and after photon radiation are well separated in phase space (see $`M_{i\gamma }^2`$ and $`M_i^2`$ defined above). As a result, the interference terms in Eq. (28) can be neglected. This leads to three distinct regions of on-shell contributions, where the photon can be assigned unambiguously to the W-pair-production subprocess or to one of the two decays. This assignment is determined by the pair of invariant masses (out of $`M_i^2`$ and $`M_{i\gamma }^2`$) that is in the $`M__W^2`$ region. Therefore, the double-pole residue can be expressed as the sum of the three on-shell contributions without increasing the intrinsic error of the DPA. Note that in the same way it is also possible to assign the photon to one of the sub-processes, since misassignment errors are suppressed, assuming for convenience that all final-state momenta can ideally be measured.
* for semi-soft photons \[$`E_\gamma =𝒪(\mathrm{\Gamma }_W)`$\] the Breit–Wigner poles are relatively close together in phase space, resulting in a substantial overlap of the line shapes. The assignment of the photon is now subject to larger errors. Moreover, since the interference terms in Eq. (28) cannot be neglected, a proper prescription for calculating the DPA residues (i.e. the phase-space mapping) is required .
* for soft photons \[$`E_\gamma \mathrm{\Gamma }_W`$\] the Breit–Wigner poles are on top of each other, resulting in a pole-scheme expansion that is identical to the one without the photon.
Let us first consider the hard-photon regime in more detail. Due to the fact that the poles are well separated in the hard-photon regime, it is clear that the interference terms are suppressed and can be neglected:
$$d\sigma =\frac{1}{2s}\left[|_0|^2+|_+|^2+|_{}|^2\right]d\mathrm{\Gamma }_{4\mathrm{f}\gamma }.$$
(34)
Note that each of the three terms has two poles, originating from two resonant propagators. However, these poles are different for different terms. The phase-space factor can be rewritten in three equivalent ways. The first is
$$d\mathrm{\Gamma }_{4\mathrm{f}\gamma }=d\mathrm{\Gamma }_\text{0}^\gamma =d\mathrm{\Gamma }_{\text{pr}}^\gamma d\mathrm{\Gamma }_{\text{dec}}^+d\mathrm{\Gamma }_{\text{dec}}^{}\frac{dM_1^2}{2\pi }\frac{dM_2^2}{2\pi },$$
(35)
with
$$d\mathrm{\Gamma }_{\text{pr}}^\gamma =\frac{1}{(2\pi )^2}\delta (q_1+q_2p_1p_2k)\frac{d\stackrel{}{p}_1}{2p_{10}}\frac{d\stackrel{}{p}_2}{2p_{20}}\frac{d\stackrel{}{k}}{(2\pi )^32k_0}.$$
(36)
The two others are
$$d\mathrm{\Gamma }_{4\mathrm{f}\gamma }=d\mathrm{\Gamma }_+^\gamma =d\mathrm{\Gamma }_{\text{pr}}d\mathrm{\Gamma }_{\text{dec}}^{+\gamma }d\mathrm{\Gamma }_{\text{dec}}^{}\frac{dM_{1\gamma }^2}{2\pi }\frac{dM_2^2}{2\pi },$$
(37)
with
$$d\mathrm{\Gamma }_{\text{dec}}^{+\gamma }=\frac{1}{(2\pi )^2}\delta (p_1k_1k_1^{}k)\frac{d\stackrel{}{k}_1}{2k_{10}}\frac{d\stackrel{}{k}_1^{}}{2k_{10}^{}}\frac{d\stackrel{}{k}}{(2\pi )^32k_0},$$
(38)
and a similar expression for $`d\mathrm{\Gamma }_{}^\gamma `$. The phase-space factors $`d\mathrm{\Gamma }_{\text{pr}}`$ and $`d\mathrm{\Gamma }_{\text{dec}}^\pm `$ are just the lowest-order ones. The cross-section can then be written in the following equivalent form
$$d\sigma =\frac{1}{2s}\left[|_0|^2d\mathrm{\Gamma }_0^\gamma +|_+|^2d\mathrm{\Gamma }_+^\gamma +|_{}|^2d\mathrm{\Gamma }_{}^\gamma \right].$$
(39)
In order to extract gauge-invariant quantities, the DPA limit should be taken. This amounts to taking the limit $`p_{1,2}^2M__W^2`$, using a particular prescription for mapping the full off-shell phase space on the kinematically restricted on-resonance one. Note however that $`p_{1,2}`$ can be different according to the $`\delta `$-functions in the decay parts of the different phase-space factors. To be specific, the production term $`|_0|^2`$ has poles at $`p_i^2=M_i^2=M__W^2`$, $`|_+|^2`$ has poles at $`p_1^2=M_{1\gamma }^2=M__W^2`$ and $`p_2^2=M_2^2=M__W^2`$, and $`|_{}|^2`$ has poles at $`p_1^2=M_1^2=M__W^2`$ and $`p_2^2=M_{2\gamma }^2=M__W^2`$.
We complete our survey of the different photon-energy regimes by considering semi-soft and soft photons. The split-up of factorizable and non-factorizable real-photon corrections proceeds in the same way as described in the previous subsection for virtual corrections. The result reads in semi-soft approximation
$$d\sigma =\frac{1}{2s}|_\gamma |^2d\mathrm{\Gamma }_{\text{4f}\gamma }d\sigma _{\text{DPA}}^0\frac{d\stackrel{}{k}}{(2\pi )^32k_0}\left[2\mathrm{R}\mathrm{e}\left(_0^\mu _{+,\mu }^{}+_0^\mu _{,\mu }^{}+_+^\mu _{,\mu }^{}\right)+|_0^2|+|_+^2|+|_{}^2|\right].$$
(40)
The gauge-invariant currents $`_0`$ and $`_\pm `$ are given by
$`_0^\mu `$ $`=`$ $`e\left[{\displaystyle \frac{p_1^\mu }{kp_1}}{\displaystyle \frac{p_2^\mu }{kp_2}}{\displaystyle \frac{q_1^\mu }{kq_1}}+{\displaystyle \frac{q_2^\mu }{kq_2}}\right],`$
$`_+^\mu `$ $`=`$ $`e\left[{\displaystyle \frac{p_1^\mu }{kp_1}}+Q_{f_1}{\displaystyle \frac{k_1^\mu }{kk_1}}Q_{f_1^{}}{\displaystyle \frac{k_{1}^{}{}_{}{}^{\mu }}{kk_1^{}}}\right]{\displaystyle \frac{D_1}{D_1+2kp_1}},`$
$`_{}^\mu `$ $`=`$ $`+e\left[{\displaystyle \frac{p_2^\mu }{kp_2}}+Q_{f_2}{\displaystyle \frac{k_2^\mu }{kk_2}}Q_{f_2^{}}{\displaystyle \frac{k_{2}^{}{}_{}{}^{\mu }}{kk_2^{}}}\right]{\displaystyle \frac{D_2}{D_2+2kp_2}}.`$ (41)
The first three interference terms in Eq. (40) correspond to the real non-factorizable corrections. The last three squared terms in Eq. (40) belong to the factorizable real-photon corrections. They constitute the semi-soft limit of Eq. (39).
### 3.4 A hybrid scheme – virtual corrections in DPA and real corrections from full matrix elements
The reliability of the error estimate of $`(\alpha /\pi )\times (\mathrm{\Gamma }_\mathrm{W}/M__W)\times \mathrm{ln}(\mathrm{})\stackrel{<}{}0.5\%`$ for the accuracy of the DPA can, of course, only be controlled by a comparison to calculations that are based on the full matrix elements. While for the virtual corrections such results do not exist yet, the situation for the real corrections is much better, since full matrix-element calculations for the processes $`\mathrm{e}^+\mathrm{e}^{}4f\gamma `$ are available . The latter results seem to be of particular importance, because the above error estimate for real corrections in DPA is subject of some controversy.
Although it deserves some care, it is possible to combine the virtual $`𝒪\left(\alpha \right)`$ corrections in DPA with real corrections from the full $`\mathrm{e}^+\mathrm{e}^{}4f\gamma `$ lowest-order matrix elements. The non-trivial point in this combination lies in the relations of IR and mass singularities in virtual and real corrections. The singularities have the form of a universal radiator function multiplied or convoluted with the respective lowest-order matrix element $`_0`$ of the non-radiative process. Since $`_0`$ appears in DPA for the virtual correction ($`_0^{\mathrm{DPA}}`$), but as full matrix element for the real ones, a simple summation of virtual and real corrections would lead to a mismatch in the singularity structure and eventually to totally wrong results. A solution of this problem is to extract those singular parts from the real photon contribution that exactly match the singular parts of the virtual photon contribution, then to replace the full $`|_0|^2`$ by $`|_0^{\mathrm{DPA}}|^2`$ in these terms and finally to add this modified part to the virtual corrections. This modification is allowed in the range of validity of the DPA and leads to a proper matching of all IR and mass singularities. The described approach for such a hybrid DPA scheme is followed in the RacoonWW program . More details of this approach can also be found in Sect. 4.1.
A particular advantage of this method is due to the fact that the leading ISR logarithms, which are part of the extracted singularities of the real corrections, can be easily kept with the full matrix element $`_0`$ (see for details). In this way, the logarithmic enhancement factor $`\mathrm{ln}(\mathrm{})`$ does not involve large contributions from the electron mass, i.e. corrections like $`\mathrm{ln}(m_\mathrm{e}^2/s)`$. In the hybrid scheme, also the non-factorizable corrections have to be treated carefully. If the full matrix elements for photon radiation is employed, one cannot exploit any cancellations between real and virtual non-factorizable corrections, as it is done in the calculations of . Instead, one needs the full set of non-factorizable virtual corrections, which includes also photons coupling to the initial state. Such results can be derived from Eq. (23) and Ref. , and are explicitly given in Ref. .
### 3.5 Intrinsic ambiguities and reliability of the double-pole approximation
The theoretical accuracy of theoretical predictions is indeed at the core of the workshop. For this reason it has already been discussed extensively in a purely theoretical context. Although only the numerical comparisons can tell us where the present theoretical uncertainty really stands, it is not superfluous that the relevant facts are summarized in one place.
An improved assessment of the theoretical uncertainty can be obtained by varying predictions within the intrinsic freedom of the followed approach for the DPA. For instance, any kind of DPA makes use of an on-shell projection of the off-shell four-fermion phase space to the phase space with on-shell $`W`$ bosons. The difference between different on-shell projections is part of the theoretical uncertainty of the DPA approach and should be considered in predictions (see Sect. 4.2 for a numerical discussion).
It is a fact of life that questions of principle are sometimes of scarce practical relevance. CC03 contains gauge-invariance-breaking terms but what is their numerical impact at LEP 2 energies? It is quite a known fact that, when computed in the ’t Hooft-Feynman gauge, they are unimportant. At least they are for the $`WW`$ total cross-section – the signal – and we can verify this statement by comparing the gauge-dependent CC03 with the full gauge-invariant cross-section (CC11 for instance) including background diagrams. There is a general agreement, dating from the ’95 workshop that the difference is less than $`0.2\%`$ at LEP 2 energies.
It is bizarre that one can render the Born CC03 diagrams gauge-invariant at the prize of large numerical variations; it is enough to project the kinematics in the matrix elements onto the on-shell phase space, while keeping the off-shellness in the Breit-Wigner propagators. However, this changes the cross-section by several per cent! Therefore, the use of DPA at Born level (CC03) is numerically not recommendable. Once more, for lowest-order reactions one needs an alternative approach and for predictions that have a DPA Born and a DPA $`𝒪\left(\alpha \right)`$ and nothing else the expected accuracy is no more than $`\mathrm{\Gamma }_W/M__W2.5\%`$. The difference between Born CC03 and Born DPA should not enter in the discussion of the theoretical uncertainty.
At the Born level one can accept a non-gauge-invariant CC03 cross-section (at least in the ’t Hooft-Feynman gauge) as a reasonable quantity at LEP 2 energies. For higher energies one should be more careful.
The same phenomenon will occur when we include radiative corrections and we would like to add some comment on the DPA procedure, in particular on the choice of projecting the kinematics.
For high enough energies, any process $`e^+e^{}VV`$ will be a dominant source of four-fermion final states due to the double resonant enhancement and hence CC03(NC02) will be a good approximation to the total cross-section for four-fermion production in a situation where we exclude certain regions of the phase space, e.g., a small scattering angle of the outgoing electron in single-$`W`$ production.
Thus, for example, to calculate the cross-section $`e^+e^{}VV`$ one proceeds as described above; one calculates the matrix element for $`e^+e^{}VV4\mathrm{f}`$ and extracts the part resonant in the invariant masses of the pairs, $`k_+^2,k_{}^2`$. The general matrix element takes the form
$$(\mathrm{},k_+,k_{},\mathrm{})=\underset{i}{}M_i(\mathrm{},k_+,k_{},\mathrm{})A_i(\mathrm{},k_+^2,k_{}^2),$$
(42)
where the $`M_i`$ contain the spinor and Lorentz tensor structure of the matrix element, e.g. they have the external fermionic wave-functions attached. The $`A_i`$ are Lorentz scalars that depend on the invariants of the problem and become non-trivial and difficult to compute when higher order corrections are included. One way of looking at the DPA-procedure is to say that the resonant part is extracted from the $`A_i`$, by Laurent expansion. The external particle wave functions, and hence the $`M_i`$, should not be affected by the process hence the kinematics of the problem should be left unchanged because the final state integrations involve only the fermions, stable on-shell particles. The gauge nature of the theory is intimately connected with the $`A_i`$ not with kinematics.
Whenever we have processes with external, unstable, vector-bosons, like in $`WWWW`$ or $`ZZZZ`$, the Higgs resonance will appear in the $`s`$-channel and by shifting e.g. a factor $`s`$ from the $`M_i`$ to the $`A_i`$ one gets factors $`s/M^2`$ which violate unitarity at high energies . This can be avoided by making the splitting between the $`A_i`$ and $`M_i`$ with some care. Here, for $`e^+e^{}WW,ZZ`$, the corresponding factors do not directly violate unitarity. Nevertheless, one could expect that Ward identities are violated by the splitting by terms of the order $`k_\pm ^2/M^21`$, i.e. non double-resonant terms negligible in the DPA approach. If, on the other hand, one includes the $`M_i`$ in the DPA, as commonly done, one has on-shell matrix elements and the WI are fulfilled, at the price of expanding kinematics.
We do not necessarily expect an improvement of the accuracy when taking the $`M_i`$ exactly, but comparing results with DPA applied to $`M_i`$ or not could give an additional estimate on the theoretical uncertainty, of the order of $`\alpha /\pi \left(\mathrm{CC03Born}/\mathrm{CC03DPA}1\right)`$. We expect that, well above threshold, this will not exceed the quoted $`0.5\%`$ DPA precision, which involves logarithmic enhancement factors.
Another questionable point in DPA is connected to the fact that a particular mapping may lead to an unphysical point in the on-shell phase-space (c.f. Sect. 3.2). Even if we do not expand the kinematics in the $`M_i`$ there are Landau singularities in the $`A_i`$ at the edge of the off-shell phase space. If one performs a DPA projection in the $`A_i`$, these Landau singularities move into the on-shell phase space, although only at a distance $`(k^2M^2)/M^2`$ from the boundary . This might happen when the $`A_i`$ are parametrized in terms of invariants. If on the other hand, one parametrizes the $`A_i`$ in terms of angles and energies, this can be more easily avoided.
Note that the formulation of a DPA where the on-shell projection is not applied to the $`M_i`$ has been implemented the formulation of the LPA of Ref. (eqs.(1) and (2)).
### 3.6 Remarks on DPA corrections to distributions inclusive w.r.t. photons
The DPA corrections to distributions that are inclusive w.r.t. photons depend in a very sensitive way on how the four-particle phase space is parametrized, or, in other words, on the way the distributions are defined after the photon has been integrated out. This statement sounds obvious, but nevertheless deserves some special attention.
In particular the invariant-mass distributions ($`W`$ line shapes) are affected. In reactions with two resonances the invariant masses have to be defined from the decay products. Depending on the precise definition of the invariant masses different sources of large Breit–Wigner distortions can be identified , in contrast to the situation at LEP1 where only initial-state radiation (ISR) can cause such distortions.
In Ref. it has been shown that also final-state radiation (FSR) can induce distortions. This is a general property of resonance-pair reactions, irrespective of the adopted scheme for implementing the finite-width effects. The only decisive factor for the distortion to take place is whether the virtuality of the unstable particle is defined with ($`s_V^{}`$) or without ($`s_V`$) the radiated photon (see Fig. 4).
Upon integration over the photon momentum, the former definition (cf $`M_{i\gamma }^2`$ defined in Sec. 3.3.2) is free of large FSR effects from the $`V`$-decay system. It can only receive large corrections from the other (production or decay) stages of the process. The latter definition (cf $`M_i^2`$ defined in Sec. 3.3.2), however, does give rise to large FSR effects from the V-decay system. In contrast to the LEP1 case, where the ISR-corrected line shape receives contributions from effectively lower $`Z`$-boson virtualities, the $`s_V`$ line shape receives contributions from effectively higher virtualities $`s_V^{}`$ of the unstable particle. As was argued above, only sufficiently hard photons ($`E_\gamma \mathrm{\Gamma }_V`$) can be properly assigned to one of the on-shell production or decay stages of the process in the DPA. For semi-soft photons \[$`E_\gamma =𝒪(\mathrm{\Gamma }_V)`$\], however, the assignment is not so clear-cut and will be determined by the experimental event-selection procedure.
Event selection procedures that involve an invariant-mass definition in terms of the decay products without the photon give rise to large FSR-induced distortion effects . These are caused by semi-soft photons, since hard FSR photons move the virtuality $`s_V^{}`$ of the unstable particle far off resonance for near-resonance $`s_V`$ values, resulting in a suppressed contribution to the $`s_V`$ line shape. This picture fits in nicely with the negligible overlap of the three on-shell double-pole contributions for hard photons, discussed above. The reason why the FSR distortions can be rather large lies in the fact that the final-state collinear singularities \[$`\frac{\alpha }{\pi }Q_f^2\mathrm{ln}(m_f^2/M_V^2)\mathrm{ln}(\mathrm{\Gamma }_V/M_V)`$\] do not vanish, even not for fully inclusive photons. After all, a fixed value of $`s_V`$ makes it impossible to sum over all degenerate final states by a mere integration over the photon momentum. So the KLN theorem does not apply in this case. These FSR distortion effects result in shifts in the measurement of the $`W`$-boson mass of the order of 40 MeV, as has been qualitatively confirmed in Ref. .
This situation changes for event-selection procedures in which not all photons can be separated from the charged fermions. If photon recombination has to be taken into account, i.e. if photons within a finite cone around the charged fermions have to be combined with the corresponding fermion into a single particle, the mentioned mass singularities connected to final-state fermions disappear. The KLN theorem applies and the large fermion-mass logarithms are effectively replaced by logarithms depending on the cone size . In Ref. this expectation has been confirmed numerically, showing that the large negative shifts in the peak position of the $`W`$ invariant-mass distribution obtained without photon recombination are reduced. In Ref. it has been shown that the effect of photon recombination can even overcompensate the momentum loss from FSR if the recombination is very inclusive. This is due to the recombination of photons that are radiated off the initial state or off particles belonging to the other decaying $`W`$ boson. The resulting positive peak shifts can amount to several $`10\mathrm{MeV}`$. Explicit numerical results on $`W`$ invariant-mass distributions can also be found in Sect. 4.
Finally we mention a special property of the non-factorizable corrections. When considering pure angular distributions with an inclusive treatment of the photons, one should integrate over the photon phase space and the invariant masses $`M_i^2`$. After integrating out both invariant masses the non-factorizable corrections will vanish, which is a typical feature of the non-factorizable interference effects .
### 3.7 Double-pole approximations in practice
For LEP 2 energies three different groups<sup>6</sup><sup>6</sup>6Another DPA has been discussed in Ref. for linear-collider energies. have formulated versions of a DPA for $`\mathrm{e}^+\mathrm{e}^{}\mathrm{WW}4f(+\gamma )`$. While Beenakker, Berends and Chapovsky , called BBC in the following, formulated a semi-analytic DPA, the other two groups implemented variants of the DPA in the event generators YFSWW and RacoonWW . The basic features of these different implementations are summarized in the following.
#### 3.7.1 The YFSWW approach
YFSWW: $`𝒪\left(\alpha \right)`$ correction to $`\mathrm{e}^+\mathrm{e}^{}\mathrm{W}^+\mathrm{W}^{}`$ in LPA, using the results of Ref. , leading-log corrections to leptonic $`W`$ decays via PHOTOS (up to two radiative photons with finite $`p_t`$ according to the exact $`𝒪\left(\alpha \right)`$ soft limit), $`W`$ decays normalized to branching ratios, quark hadronization with JETSET and $`\tau `$ decays with TAUOLA (including radiative corrections), YFS exponentiation for ISR and photon emission from $`W`$-bosons, off-shell Coulomb singularity, no full non-factorizable corrections – only an approximation in terms of the screened Coulomb ansatz of Ref. , approximate $`W`$ spin correlations (incomplete correlation beyond Born) – they are missing only in a non-IR non-LL part of EW virtual corrections.
#### 3.7.2 The BBC approach
BBC: semi-analytical calculation of complete $`𝒪\left(\alpha \right)`$ corrections in DPA (with both factorizable and non-factorizable corrections and $`W`$ spin correlations), no background. Since the DPA is only valid well above threshold, the on-shell part of the Coulomb singularity is automatically included as part of the factorizable corrections and the off-shell part is contained in the non-factorizable corrections, as discussed in Ref. .
#### 3.7.3 The RacoonWW approach
RacoonWW treats the virtual $`𝒪\left(\alpha \right)`$ corrections to $`\mathrm{e}^+\mathrm{e}^{}\mathrm{WW}4f`$ in DPA. No further approximations beyond the pole expansion of the matrix element are made, i.e. non-factorizable corrections are included, and $`W`$-spin correlations are respected. The Coulomb singularity is part of the virtual corrections, and the corresponding part that goes beyond DPA has been added as discussed in Ref. . The real $`𝒪\left(\alpha \right)`$ corrections are based on the full $`4\mathrm{f}+\gamma `$ matrix element (of the CC11 class), so that the full kinematics is supported also for photon radiation. All matrix elements are based on massless fermions, and fermion masses are introduced only for collinear photon emission that is inclusive within a (small) finite cone for each fermion. Thus, a photon collinear to an outgoing fermion has to be recombined with the corresponding fermion, and a photon close to the beams has to be considered as invisible. Initial-state radiation beyond $`𝒪\left(\alpha \right)`$ is treated in the structure-function approach, including soft-photon exponentiation and leading-log contributions up to $`𝒪\left(\alpha ^3\right)`$.
### 3.8 The fermion-loop and non-local approaches
As was mentioned above, the alternative to subtracting sub-leading gauge-violating terms is to add gauge-restoring terms to the calculation. In order to do this, one has to add to the amplitude those terms that are needed for satisfying the Ward identities. This is not easy to do in general. The following observation helps. The very fact that the perturbative amplitudes require re-summation of the self-energies indicates that either the perturbative expansion parameter (coupling constant) is not the proper one, or alternatively that the quantity that is expanded (i.e. the lowest-order Lagrangian of the Standard Model) is not the best choice. This observation leads one to consider first the one-loop corrected effective potential of the Standard Model before doing Born calculations, in order to avoid Dyson re-summation of the self-energies.
For the discussion of the fermion-loop and non-local approaches it is therefore worthwhile to first have a closer look at the origin of the gauge-invariance problem associated with the re-summation of self-energies. To this end we consider the simple example of an unbroken non-abelian $`SU(N)`$ gauge theory with fermions and subsequently integrate out these fermions .
First we fix the notations and introduce some conventions. The $`SU(N)`$ generators in the fundamental representation are denoted by $`\text{T}^a`$ with $`a=1,\mathrm{},N^21`$. They are normalized according to $`\text{Tr}(\text{T}^a\text{T}^b)=\delta ^{ab}/2`$ and obey the commutation relation $`[\text{T}^a,\text{T}^b]=if^{abc}\text{T}^c`$. In the adjoint representation the generators $`\text{F}^a`$ are given by $`(\text{F}^a)^{bc}=if^{abc}`$. The Lagrangian of the unbroken $`SU(N)`$ gauge theory with fermions can be written as
$$(x)=\frac{1}{2}\text{Tr}\left[𝑭_{\mu \nu }(x)𝑭^{\mu \nu }(x)\right]+\overline{\psi }(x)(iD/m)\psi (x),$$
(43)
with
$$𝑭_{\mu \nu }\text{T}^aF_{\mu \nu }^a=\frac{i}{g}[D_\mu ,D_\nu ],D_\mu =_\mu ig\text{T}^aA_\mu ^a_\mu ig𝑨_\mu .$$
(44)
Here $`\psi `$ is a fermionic $`N`$-plet in the fundamental representation of $`SU(N)`$ and $`A_\mu ^a`$ are the ($`N^21`$) non-abelian $`SU(N)`$ gauge fields, which form a multiplet in the adjoint representation. The Lagrangian (43) is invariant under the $`SU(N)`$ gauge transformations
$$\psi (x)\psi ^{}(x)=\text{G}(x)\psi (x),$$
$$𝑨_\mu (x)𝑨_\mu ^{}(x)=\text{G}(x)𝑨_\mu (x)\text{G}^1(x)+\frac{i}{g}\text{G}(x)\left[_\mu \text{G}^1(x)\right],$$
(45)
with the $`SU(N)`$ group element defined as $`\text{G}(x)=\mathrm{exp}[ig\text{T}^a\theta ^a(x)]`$. The covariant derivative $`D_\mu `$ and field strength $`𝑭_{\mu \nu }`$ both transform in the adjoint representation
$$D_\mu \text{G}(x)D_\mu \text{G}^1(x),𝑭_{\mu \nu }(x)\text{G}(x)𝑭_{\mu \nu }(x)\text{G}^1(x).$$
(46)
Since the Lagrangian is quadratic in the fermion fields, one can integrate them out exactly in the functional integral. The resulting effective action is then given by
$$iS_{\text{eff}}[J]=id^4x\{\frac{1}{2}\text{Tr}\left[𝑭_{\mu \nu }(x)𝑭^{\mu \nu }(x)\right]+J_\mu ^a(x)A^{a,\mu }(x)\}+\text{Tr}\left[\mathrm{ln}(D/im)\right],$$
(47)
with $`J_\mu ^a(x)`$ denoting the gauge-field sources. The trace on the right-hand side has to be taken in group, spinor, and coordinate space. As a next step one can expand the effective action in terms of the coupling constant
$`\text{Tr}\left[\mathrm{ln}(D/im)\right]`$ $`=`$ $`\text{Tr}\left[\mathrm{ln}(/im)\right]+\text{Tr}\left[\mathrm{ln}(1+{\displaystyle \frac{g}{i/m}}𝑨/)\right]`$ (48)
$`=`$ $`\text{Tr}\left[\mathrm{ln}(/im)\right]+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^{n1}}{n}}\text{Tr}\left[({\displaystyle \frac{g}{i/m}}𝑨/)^n\right].`$
Note that the left-hand side of Eq. (48) is gauge-invariant as a result of the trace-log operation. In contrast, the separate terms of the expansion on the right-hand side are not gauge-invariant. This is due to the fact that, unlike in the abelian case, the non-abelian gauge transformation (45) mixes different powers of the gauge field $`A_\mu `$ in Eq. (48). Thus, if one truncates the series on the right-hand side of Eq. (48) one will in general break gauge invariance. From Eq. (48) it is also clear that the fermionic part of the effective action induces higher-order interactions between the gauge bosons.
What are these higher-order interactions? Let us consider the quadratic gauge-field contribution
$$\frac{1}{2}\text{Tr}\left[(\frac{g}{i/m}𝑨/)^2\right]=\frac{1}{2}d^4xd^4y\text{Tr}\left[O(x,y)O(y,x)\right],$$
(49)
where
$$O(x,y)=gS_\text{F}^{(0)}(xy)𝑨/(y)$$
(50)
and $`iS_\text{F}^{(0)}(xy)=<0|T(\psi (x)\overline{\psi }(y))|\mathrm{\hspace{0.17em}0}>_{\text{free}}`$ is the free fermion propagator. The trace on the right-hand side of Eq. (49) has to be taken in group and spinor space. A quick glance at this quadratic gauge-field contribution reveals that it is just the one-loop self-energy of the gauge boson induced by a fermion loop. In the same way, the higher-order terms $`g^nA^n`$ in Eq. (48) are just the fermion-loop contributions to the $`n`$-point gauge-boson vertices.
One can truncate the expansion in Eq. (48) at $`n=2`$, thus taking into account only the gauge-boson self-energy term and neglecting the fermion-loop contributions to the higher-point gauge-boson vertices. This is evidently the simplest procedure for performing the Dyson re-summation of the fermion-loop self-energies. However, as was pointed out above, truncation of Eq. (48) at any finite order in $`g`$ in general breaks gauge invariance. This leads to the important observation that, although the re-summed fermion-loop self-energies are gauge-independent by themselves, the re-summation is nevertheless responsible for gauge-breaking effects in the higher-point gauge-boson interactions through its inherent mixed-order nature. Another way of understanding this is provided by the gauge-boson Ward identities. Since the once-contracted $`n`$-point gauge-boson vertex can be expressed in terms of $`(n1)`$-point vertices, it is clear that gauge invariance is violated if the self-energies are re-summed without adding the necessary compensating terms to the higher-point vertices.
An alternative is to keep all the terms in Eq. (48). Then the matrix elements derived from the effective action will be gauge-invariant. Keeping all the terms means that we will have to take into account not only the fermion-loop self-energy in the propagator, but also all the possible fermion-loop contributions to the higher-point gauge-boson vertices. This is exactly the prescription of the fermion-loop scheme (FLS) . Although the FLS guarantees gauge invariance of the matrix elements, it has disadvantages as well. Its general applicability is limited to those situations where non-fermionic particles can effectively be discarded in the self-energies, as is for instance the case for $`\mathrm{\Gamma }_W`$ and $`\mathrm{\Gamma }_Z`$ at lowest order. Another disadvantage is that in the FLS one is forced to do the loop calculations, even when calculating lowest-order quantities. For example, the calculation of the tree-level matrix element for the process $`e^+e^{}4\mathrm{f}\gamma `$ involves a four-point gauge-boson interaction, which has to be corrected by fermion loops in the FLS. This over-complicates an otherwise lowest-order calculation.
It is clear that the FLS provides more than we actually need. It does not only provide gauge invariance for the Dyson re-summed matrix elements at a given order in the coupling constant, but it also takes into account all fermion-loop corrections at that given order. In the vicinity of unstable particle resonances the imaginary parts of the fermion-loop self-energies are effectively enhanced by $`𝒪(1/g^2)`$ with respect to the other fermion-loop corrections. Therefore, what is really needed is only a minimal subset of the non-enhanced contributions such that gauge invariance is restored. In a sense one is looking for a minimal solution of a system of Ward identities. The FLS provides a solution, but this solution is far from minimal and is only practical for particles that decay exclusively into fermions. Since the decay of unstable particles is a physical phenomenon, it seems likely that there exists a simpler and more natural method for constructing a solution to a system of Ward identities, without an explicit reference to fermions. This is precisely the philosophy behind the non-local approach . This approach consists in using gauge-invariant non-local effective Lagrangians for generating both the self-energy effects in the propagators and the required gauge-restoring terms in the higher-point interactions. In this way the full set of Ward identities can be solved, while keeping the gauge-restoring terms to a minimum.
#### 3.8.1 The fermion-loop scheme
The Fermion-Loop scheme developed in and refined in makes the approximation of neglecting all masses for the incoming and outgoing fermions in the processes $`e^+e^{}n`$fermions. It is possible, however, to go beyond this approximation and give the construction of an exact Fermion-Loop scheme (EFL) , i.e. , a scheme for incorporating the finite-width effects in the theoretical predictions for tree-level, LEP 2 and beyond, processes.
One can work in the ’t Hooft-Feynman gauge and create all relevant building blocks, namely the vector-vector , vector-scalar and scalar-scalar transitions of the theory, all of them one-loop re-summed. The loops, entering the scheme, contain fermions and, as done before in , one allows for a non-zero top quark mass inside loops. There is a very simple relation between re-summed transitions and running parameters, since Dyson re-summation is most easily expressed in terms of running couplings and running mixing angles.
In the EFL generalization, it is particularly convenient to introduce additional running quantities. They are the running masses of the vector bosons, $`M__0^2(p^2)=M^2(p^2)/c^2(p^2)`$, formally connected to the location of the $`W`$ and $`Z`$ complex poles. After introducing these running masses, it is straightforward to prove that all $`𝒮`$-matrix elements of the theory assume a very simple structure. Coupling constants, mixing angles and masses are promoted to running quantities and the $`𝒮`$-matrix elements retain their Born-like structure, with running parameters instead of bare ones, and vector-scalar or scalar-scalar transitions disappear if we employ unitary-gauge–like vector boson propagators where the masses appearing in the denominator of propagators are the running ones. If the $`WW`$ and $`\varphi \varphi `$ transitions are denoted by $`S_W^{\mu \nu }`$ and by $`S_\varphi `$ with,
$`S_W^{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{g^2}{16\pi ^2}}\mathrm{\Sigma }_W^{\mu \nu },\mathrm{\Sigma }_W^{\mu \nu }=\mathrm{\Sigma }_W^0\delta ^{\mu \nu }+\mathrm{\Sigma }_W^1p^\mu p^\nu ,`$
$`S_\varphi `$ $`=`$ $`{\displaystyle \frac{g^2}{16\pi ^2}}\mathrm{\Sigma }_\varphi ,`$ (51)
then, the $`W`$-boson running mass is defined by the following equation (note the metric):
$$\frac{1}{M^2(p^2)}=\frac{1}{M^2}\frac{p^2S_W^0+\frac{M^2}{p^2}S_\varphi }{p^2S_\varphi },$$
(52)
The whole amplitude can be written in terms of a $`W`$-boson exchange diagram, if we make use of the following effective propagator:
$$\mathrm{\Delta }_{\mathrm{eff}}^{\mu \nu }=\frac{1}{p^2+M^2S_W^0}\left[\delta ^{\mu \nu }+\frac{p^\mu p^\nu }{M^2(p^2)}\right].$$
(53)
For the vertices we need that all vector-boson lines be off mass-shell and non-conserved and, moreover, a Ward identity has to be computed and not only the corresponding amplitude. Therefore, the number of terms increases considerably with respect to the standard formulation of the FL-scheme and we refer to Ref. for all details.
The renormalization of ultraviolet divergences can be easily extended to the EFL-scheme by showing that all ultraviolet divergent parts of the one-loop vertices, $`\gamma WW,\gamma W\varphi ,\gamma \varphi W`$ and $`\gamma \varphi \varphi `$ for instance, are proportional to the lowest order part. Therefore, the only combinations that appear are of the form $`1/g^2+VVV`$vertex or $`M^2/g^2+VV\varphi `$ vertex etc. All of them are, by construction, ultraviolet finite.
Equipped with this generalization of the Fermion-Loop scheme, one can prove the fully-massive $`U(1)`$ Ward identity which is required for a correct treatment of the single-$`W`$ processes. As a by-product of the method, the cross-section for single-$`W`$ production automatically evaluates all channels at the right scale, without having to use ad hoc re-scalings and avoiding the approximation of a unique scale for all terms contributing to the cross-section.
The generalization of the Fermion-Loop scheme goes beyond its, most obvious, application to single-$`W`$ processes and allows for a gauge invariant treatment of all $`e^+e^{}n`$fermion processes with a correct evaluation of the relevant scales. Therefore, the EFL-scheme can be applied to several other processes like $`e^+e^{}Z\gamma ^{}`$ and, in general to $`e^+e^{}6`$fermion processes, with the inclusion of a stable, external, top quark, but it does not apply to reactions involving the physical Higgs boson. Furthermore, the scheme misses those corrections to the total decay width in the propagator denominators that are induced by two-loop contributions.
#### 3.8.2 The non-local approach
The main idea of the non-local approach is to rearrange the series on the right-hand side of Eq. (48) in such a way that each term becomes gauge-invariant by itself. Subsequent truncation of the series at a given term is then allowed. It is possible to approximate Eq. (48) by means of an effective Lagrangian in such a way that the resulting effective action has the following properties:
* it generates the Dyson re-summed transverse gauge-boson self-energy in the propagator. This means that it contributes to the gauge-boson two-point function. Hence, the effective lagrangian should depend at least on two gauge-boson fields.
* the Dyson re-summed self-energy is in general not a constant, but rather a function dependent on the interaction between the gauge-bosons and the fermions. This means that the effective Lagrangian should in general be non-local (bi-local) in the gauge fields. Thus the gauge fields should be taken at two different space–time points.
* it is gauge-invariant. As such the effective Lagrangian should have the form of an infinite tower of gauge fields.
For the gauging procedure of the non-local Lagrangians we will need a special ingredient, the path-ordered exponential, which is defined as
$$U(x,y)=U^{}(y,x)=\text{P}\mathrm{exp}\left[ig\underset{x}{\overset{y}{}}𝑨_\mu (\omega )𝑑\omega ^\mu \right]$$
(54)
Here $`d\omega ^\mu `$ is the element of integration along some path $`\mathrm{\Omega }(x,y)`$ that connects the points $`x`$ and $`y`$.<sup>7</sup><sup>7</sup>7In principle we are free to choose this particular path. This freedom is just one out of the many freedoms that characterize the treatment of unstable particles (as mentioned earlier). It just reflects the fact that in a perturbative expansion one is free to pick up additional higher-order contributions, since the answer at any given (truncated) order will not be changed by such additional terms. The so-defined path-ordered exponential transforms as
$$U(x,y)\text{G}(x)U(x,y)\text{G}^1(y)$$
(55)
under the $`SU(N)`$ gauge transformations. It hence carries the gauge transformation from one space-time point to the other.
For a $`SU(N)`$ Yang–Mills theory the non-local action with the above-described properties takes the form
$$𝒮_{\text{NL}}=\frac{1}{2}d^4xd^4y\mathrm{\Sigma }_{\text{NL}}(xy)\text{Tr}\left[U(y,x)𝑭_{\mu \nu }(x)U(x,y)𝑭^{\mu \nu }(y)\right]d^4xd^4y_{\text{NL}}(x,y),$$
(56)
with $`_{\text{NL}}(x,y)`$ the non-local effective Lagrangian. As required, the action contains bilinear gauge-boson interactions. The induced infinite tower of higher-point gauge-boson interactions, which are also of progressively higher order in the coupling constant $`g`$, is needed for restoring gauge invariance.
It is important to stress at this point that this term in the effective action should not be understood as a new fundamental interaction. It is generated by radiative corrections. From the point of view of general properties of non-local Lagrangians, the non-local coefficient $`\mathrm{\Sigma }_{\text{NL}}(xy)`$ is arbitrary. In practice, however, it is fixed by the explicit interaction between the gauge-bosons and the fermions in the underlying fundamental theory. In our simple example this connection is given by Eq. (48).
Let us now derive the two-point function as an example of the Feynman rules generated by Eq. (56):
$$\text{}:i\mathrm{\Sigma }^{a_1a_2,\mu _1\mu _2}(x_1,x_2)=\frac{i\delta ^2(𝒮_\text{L}+𝒮_{\text{NL}})}{\delta A_{\mu _1}^{a_1}(x_1)\delta A_{\mu _2}^{a_2}(x_2)}.|_{A=0},$$
(57)
where the local action $`𝒮_\text{L}`$ follows from the gauge-boson term in Eq. (43). The Fourier transform of this two-point function can be calculated in a straightforward way, since the path-ordered exponentials are effectively unity. The result reads
$$i\stackrel{~}{\mathrm{\Sigma }}^{a_1a_2,\mu _1\mu _2}(q_1,q_2)=i\delta ^{a_1a_2}\left(q_1^\mu q_1^\nu q_1^2g^{\mu \nu }\right)\left[1+\stackrel{~}{\mathrm{\Sigma }}_{\text{NL}}(q_1^2)\right](2\pi )^4\delta ^{(4)}(q_1+q_2).$$
(58)
Note that this two-point interaction is transverse, as it should be for an unbroken theory. The non-local coefficient acts as a (dimensionless) correction to the transverse free gauge-boson propagator, exactly what is needed for the Dyson re-summation of the gauge-boson self-energies. The infinite tower of gauge-restoring higher-point gauge-boson interactions are provided by the gauge-boson fields present in both $`𝑭_{\mu \nu }`$ and the path-ordered exponential occurring in Eq. (56). For explicit Feynman rules we refer to Ref. .
Although the above-described non-local procedure provides a gauge-invariant framework for performing the Dyson re-summation of the gauge-boson self-energies, we want to stress that it is not unique. We have seen above that the FLS provides a different solution of the system of gauge-boson Ward identities. In the context of non-local effective Lagrangians it is always possible to add additional towers of gauge-boson interactions that start with three-point interactions and therefore do not influence the Dyson re-summation of the gauge-boson self-energies.
In the light of the discussion presented in Sect. 3.8, we rearrange the series on the right-hand side of Eq. (48) according to gauge-invariant towers of gauge-boson interactions labelled by the minimum number of gauge bosons that are involved in the non-local interaction. Effectively this constitutes an expansion in powers of the coupling constant $`g`$, since a higher minimum number of particles in the interaction is equivalent to a higher minimum order in $`g`$. In order to achieve minimality we have truncated this series at the lowest effective order. This should not be viewed as some ad hoc recipe, but rather as a systematic expansion of the effective potential.
Up to now we have seen how the non-local effective Lagrangian method works for unstable gauge bosons in a simple $`SU(N)`$ gauge theory with fermions. In the Standard Model there are different types of unstable particles: the top-quark, the massive gauge-bosons, the Higgs boson. In Ref. it was shown how to extend the above-described method in such a way that it allows the description of all the unstable particles in terms of bi-local effective Lagrangians.
## 4 The CC03 cross-section, $`\sigma _{WW}`$
As mentioned before, a new electroweak $`𝒪\left(\alpha \right)`$ CC03 cross-section is available, showing a result that is between $`2.5\%`$ and $`3\%`$ smaller than the old 1995 CC03 cross-section predicted with GENTLE. This is a big effect since the combined experimental accuracy of LEP experiments is even smaller.
In the ’95 workshop predictions for CC03 were produced with variations in the IPS which agreed at the level of $`1\%`$, and then a $`2\%`$ theoretical error was quoted, to be conservative. How does this estimate compare with the present shift of $`2.5÷3\%`$ downwards? This is a $`1.25`$ to $`1.5`$ sigma difference, totally acceptable within the area of statistics. Certainly, this is more of a systematic theoretical uncertainty which is hard to quantify, but still: it is compatible and in agreement. However, a comment is needed here. In ’95 several groups produced tuned comparisons for CC03 agreeing at the level of one part in $`10^4`$. Then they moved to the Best-You-Can approach, defined by switching on all flags to get the best physics description according to the flag description of individual codes. The program GENTLE, in its BYC-mode, was selected to represent the Standard Model. However, if we take other codes, noticeably WPHACT and WTO, we easily discover CC03, Born-like, predictions that have a maximal $`+1.6\%`$ shift with respect to RacoonWW ($`+1.3\%`$ with respect to YFSWW) at the highest energy. Therefore, the old estimate of $`2\%`$ in theoretical accuracy was not underestimated.
It is important to discuss the numerical predictions for the DPA-corrected CC03 cross-section. Therefore, in this Section, we present numerical results and also an accurate description of the comparisons between different approaches, YFSWW, BBC and RacoonWW. In principle, one would like to understand the effect of DPA and, therefore, is interested in the ratio (with DPA)/(without DPA), both with ISR, (naive) QCD etc. for each of the programs. For this Report, however, this was not done and we have to take the old results (e.g. GENTLE) for a comparison newold. By comparing different calculations one can numerically check the quality of the DPA for CC03.
### 4.1 Description of the programs and their results
#### CC03 with RacoonWW
### Authors
#### General description
The program RacoonWW evaluates cross-sections and differential distributions for the reactions $`\mathrm{e}^+\mathrm{e}^{}4\mathrm{f}`$ and $`\mathrm{e}^+\mathrm{e}^{}4\mathrm{f}+\gamma `$ for all four-fermion final states. For the W-pair mediated channels $`\mathrm{e}^+\mathrm{e}^{}\mathrm{WW}4\mathrm{f}(+\gamma )`$ the full virtual $`𝒪\left(\alpha \right)`$ corrections are taken into account in DPA, while for the corresponding real corrections the full $`4\mathrm{f}+\gamma `$ matrix elements are used.
#### Features of the program
* Lowest order: the full matrix elements for all $`4\mathrm{f}`$ final states are included, and the contribution of the CC03 matrix elements or of other subsets of diagrams is provided as an option. All external fermions are assumed to be massless.
* Virtual $`𝒪\mathbf{\left(}\alpha \mathbf{\right)}`$ corrections: the full one-loop corrections are included in DPA, i.e. all factorizable corrections and non-factorizable corrections . In this way, full $`W`$-spin correlations are taken into account.
* Real corrections — $`\mathrm{𝟒}𝐟\mathbf{+}\gamma `$ production: the cross-sections are based on the full matrix-element calculation for all $`4\mathrm{f}+\gamma `$ final states with massless fermions. If the process $`\mathrm{e}^+\mathrm{e}^{}4f+\gamma `$ is investigated with a separable photon, i.e. if the photon is neither soft nor collinear to a charged fermion, all $`4\mathrm{f}+\gamma `$ final states are possible, and subsets of diagrams can be chosen as options (e.g. boson-pair production diagrams, QCD background diagrams). If the real corrections to $`\mathrm{e}^+\mathrm{e}^{}\mathrm{WW}4f`$ are calculated, the full $`4\mathrm{f}+\gamma `$ matrix elements for the CC11 class<sup>8</sup><sup>8</sup>8The CC11 class is the smallest gauge-invariant subset of diagrams for $`\mathrm{e}^+\mathrm{e}^{}4f`$ that contains all graphs with two resonant $`W`$ bosons; in this class only those background diagrams are missing that are peculiar to $`\mathrm{e}^\pm `$, $`\nu _\mathrm{e}`$, $`\overline{\nu }_\mathrm{e}`$, or $`f\overline{f}`$ pairs in the final state. are taken, i.e. photon radiation from background diagrams is partially included.
Depending on the choice of the user, the cancellation of collinear and infrared singularities is performed within the phase-space slicing method or within the subtraction formalism of Ref. . In both cases, care is taken in avoiding mismatch between the singularities of the virtual and the real corrections, which is non-trivial owing to the application of the DPA to the virtual corrections only. The treatment of fermion-mass singularities is described below in more detail.
* ISR: higher-order ISR is implemented via structure functions for the incoming $`e^+`$ and $`e^{}`$. The structure functions used are those of Ref. with the ‘BETA’ choice, i.e. the collinear-soft leading logarithms are exponentiated. If the $`𝒪\left(\alpha \right)`$ corrections to $`\mathrm{e}^+\mathrm{e}^{}\mathrm{WW}4f`$ are included, the $`𝒪\left(\alpha \right)`$ contributions already contained in the structure functions are subtracted, in order to avoid double counting, and the full CC11 Born matrix elements are used in the convolution.
* Treatment of collinear photons: the program is only applicable to observables that involve no mass-singular contributions from the final state. These mass singularities cancel if all photons collinear to a charged final-state fermion are combined with this fermion<sup>9</sup><sup>9</sup>9Note that without photon recombination, only the total cross section (without any cuts) fulfills this requirement, whereas distributions or cuts that make use of fermion momenta in general involve mass-singular corrections.. The recombination procedure is controlled by recombination cuts, i.e. photon emission angles and photon energies, or invariant masses of photon–fermion pairs. Specifically, first the charged fermion that is closest to the photon according to these criteria (emission angle or invariant mass) is selected, and secondly the photon is recombined with this fermion if it is within the recombination cuts for a final-state fermion and discarded for an initial-state fermion. The mass singularities that remain from collinear photon emission off initial-state electron or positron \[i.e. the $`(\alpha \mathrm{ln}m_\mathrm{e})^n`$ terms\] are included in the structure functions.
* Coulomb singularity: within DPA it is fully included in the $`𝒪\left(\alpha \right)`$ corrections. The full off-shell behaviour of the singularity as described in Ref. can be switched on as an option.
* Finite gauge-boson widths: in the tree-level processes $`\mathrm{e}^+\mathrm{e}^{}4\mathrm{f},4\mathrm{f}+\gamma `$ several options are included, such as fixed-width, running-width, and complex-mass scheme . If $`𝒪\left(\alpha \right)`$ corrections are taken into account the fixed width is automatically used.
* Cuts: since each event is completely specified, in principle any conceivable phase-space cut can be implemented. However, since all fermions are taken to be massless, singularities can occur in photon-exchange channels, rendering cuts unavoidable. In particular, if a charged fermion–anti-fermion pair is produced, a lower cut on its invariant mass has to be specified, or if a final-state electron or positron is present, cuts on its minimal angle to the beam and its minimal energy are required. For calculations based on restricted sets of diagrams, not all cuts are necessary; in particular, no cut at all is needed for the CC03 diagrams.
* QCD contributions: gluon-exchange contributions can be switched on in the tree-level processes $`\mathrm{e}^+\mathrm{e}^{}4\mathrm{f},4\mathrm{f}+\gamma `$. Gluon-emission processes $`\mathrm{e}^+\mathrm{e}^{}4\mathrm{f}+g`$ can be calculated for the CC11 class of $`4\mathrm{f}`$ final states.
For the QCD corrections to $`\mathrm{e}^+\mathrm{e}^{}\mathrm{WW}4f`$, one can choose between the naive QCD factors of $`(1+\alpha _\mathrm{s}/\pi )`$ per hadronically decaying W boson and the full $`𝒪\left(\alpha _\mathrm{s}\right)`$ corrections in DPA. The full calculation is performed in the same way as the photonic parts of the $`𝒪\left(\alpha \right)`$ corrections.
* IBA: the program includes, as an option, an improved Born approximation (IBA) , which involves the leading ISR logarithms, the running of the electromagnetic coupling, corrections associated with the $`\rho `$ parameter, and the Coulomb singularity.
* Subsets of diagrams: For lowest-order predictions of $`\mathrm{e}^+\mathrm{e}^{}4f,4\mathrm{f}+\gamma `$ there is the possibility to select subsets of diagrams, such as those including the pair production of $`\mathrm{W}`$, $`\mathrm{Z}`$, $`\mathrm{Z}/\gamma ^{}`$, or $`\mathrm{W}/\mathrm{Z}/\gamma ^{}`$ bosons. Furthermore, all diagrams corresponding to the CC11 process class can be selected.
* Intrinsic ambiguities: the accuracy of the DPA can be studied by changing the DPA within its intrinsic ambiguities. This is described in Sect. 4.2.
#### Program layout
RacoonWW consists of two nearly independent Monte Carlo programs: one uses phase-space slicing and the other the subtraction method of Ref. . Only the main control program, the routines for photon recombination and phase-space cuts, and the calculation of the matrix elements are commonly used. The numerical integration is performed with the multi-channel Monte Carlo technique and adaptive weight optimization . The generator produces weighted events.
#### Input parameters/schemes
RacoonWW needs the following input parameters:
$`\alpha (0),\alpha (M__Z),G_F,\alpha _\mathrm{s},M__W,M__Z,M__H,\mathrm{\Gamma }_W,\mathrm{\Gamma }_Z,`$
$`m_f,f=\mathrm{e},\mu ,\tau ,\mathrm{u},\mathrm{c},\mathrm{t},\mathrm{d},\mathrm{s},\mathrm{b}.`$ (59)
The weak mixing angle is fixed by $`c_\mathrm{w}^2=1s_\mathrm{w}^2=M__W^2/M__Z^2`$, and the quark-mixing matrix is set to unity. The masses of external fermions are consistently set to zero where possible. While the masses of the final-state fermions appear only as regulators, the mass singular logarithms of ISR depend on $`m_\mathrm{e}`$. The user can choose between the externally fixed $`W`$ width $`\mathrm{\Gamma }_\mathrm{W}`$ and an internally calculated value including electroweak and/or QCD one-loop radiative corrections.
The parameter set (59) is over-complete. The program supports three different input schemes, fixing the independent parameters. We recommend to use the $`G_F`$ scheme where the tree level is fixed by $`G_F`$, $`M_\mathrm{W}`$, and $`M_\mathrm{Z}`$ and the relative $`𝒪(\alpha )`$ corrections are calculated with $`\alpha (0)`$.
The code is available from the authors upon request.
#### Numerical results
In Tab.(2) we list the predictions of RacoonWW for the total CC03 cross-section including radiative corrections (best-with-CC03-Born as defined below). We give the results for one leptonic channel, for one semi-leptonic channel, for one hadronic channel, and for the sum of all channels separately. Note that for CC03 and negligible fermion masses the results are independent of the final state within these channels. No cuts are applied. While in all other RacoonWW results in this report LL $`𝒪\left(\alpha ^3\right)`$ corrections according to Ref. are included, in this table only the LL $`𝒪\left(\alpha ^2\right)`$ terms are taken into account. The LL $`𝒪\left(\alpha ^3\right)`$ contributions reduce the cross-sections by only about $`0.02\%`$ The given errors are purely statistical. The error for the total cross section were obtained by adding the (statistically correlated) errors of the various channels linearly.
In the following we show the predictions from RacoonWW for the $`M(W^{})`$ invariant mass distributions in four different configurations:
| 4f-Born: | full $`\mathrm{e}^+\mathrm{e}^{}4f`$ Born without radiative corrections; |
| --- | --- |
| best-with-4f-Born: | full $`\mathrm{e}^+\mathrm{e}^{}4f`$ Born plus radiative corrections |
| | including ISR beyond $`𝒪\left(\alpha \right)`$, |
| | soft photon exponentiation, |
| | LL $`𝒪\left(\alpha ^3\right)`$, and naive QCD ; |
| CC03-Born: | CC03 Born without radiative corrections; |
| best-with-CC03-Born: | CC03 Born plus radiative corrections |
| | including ISR beyond $`𝒪\left(\alpha \right)`$, |
| | soft photon exponentiation, |
| | LL $`𝒪\left(\alpha ^3\right)`$, and naive QCD, |
for the three final states, $`\mu ^+\nu _\mu \tau ^{}\overline{\nu }_\tau ,u\overline{d}\mu ^{}\overline{\nu }_\mu `$ and $`u\overline{d}s\overline{c}`$ at $`\sqrt{s}=200\mathrm{GeV}`$.
As explained in the text, DPA sits only in the virtual correction in the RacoonWW approach. Everything else is (or can be) calculated from full $`4f(\gamma )`$ matrix elements. This means that best-with-4f-Born and best-with-CC03-Born contain the same DPA part (the virtual correction).
All distributions have been obtained with the following cut and photon recombination procedure:
1. All photons within a cone of $`5^{}`$ around the beams are treated as invisible, i.e. their momenta are disregarded when calculating angles, energies, and invariant masses.
2. Next, the invariant masses of the photon with each of the charged final-state fermions are calculated. If the smallest one is smaller than $`M_{\mathrm{rec}}`$ or if the photon energy is smaller than $`1\mathrm{GeV}`$, the photon is combined with the corresponding fermion, i.e. the momenta of the photon and the fermion are added and associated with the momentum of the fermion, and the photon is discarded.
3. Finally, all events are discarded in which one of the final-state charged fermions is within a cone of $`10^{}`$ around the beams. No other cuts are applied.
We consider the cases of a tight recombination cut $`M_{\mathrm{rec}}=5\mathrm{GeV}`$ (bare) and of a loose recombination cut $`M_{\mathrm{rec}}=25\mathrm{GeV}`$ (calo). Born predictions are independent of the recombination cut. The $`W^{}`$ invariant-mass is always defined via the four-momenta (after eventual recombination with the photon) of the $`W^{}`$ decay fermions.
In Fig. 5 (left) we show the CC03-Born predictions for the $`M(W^{})`$ distributions for all three final states. The r.h.s. of Fig. 5 shows best-with-CC03-Born with the bare recombination cut , i.e. the corrections are included in DPA. In Fig. 6 we show the effect of the radiative corrections by computing the ratio of the invariant-mass distributions including radiative corrections and the Born distributions both for bare and calo recombination. In the peak region, i.e. $`|M(W^{})M__W|<\mathrm{\Gamma }_W/2`$, the effects of radiative corrections lower the line-shape by approximately $`3\%`$ $`(5\%)`$ ($`u\overline{d}s\overline{c}`$), $`7\%`$ $`(7\%)`$ ($`u\overline{d}\mu ^{}\overline{\nu }_\mu `$), and $`11\%`$ $`(12\%)`$ ($`\mu ^+\nu _\mu \tau ^{}\overline{\nu }_\tau `$) for bare (calo) distributions. The differences between the final states originate mainly from the (naive) QCD corrections.
The shape of the relative corrections to the invariant-mass distributions can be understood as follows. For small recombination cuts (bare), in most of the events the $`\mathrm{W}^{}`$ bosons are defined from the decay fermions only. If a photon is emitted from the decay fermions and not recombined, the invariant mass of the fermions is smaller than the one of the decaying $`\mathrm{W}^{}`$ boson. This leads to an enhancement of the distribution for invariant masses below the $`\mathrm{W}`$ resonance. This effect becomes smaller with increasing recombination cut $`M_{\mathrm{rec}}`$. On the other hand, if the recombination cut gets large, the probability increases that the recombined fermion momenta receive contributions from photons that are radiated during the $`\mathrm{W}`$-pair production subprocess or from the decay fermions of the $`\mathrm{W}^+`$ boson. This leads to positive corrections above the considered $`\mathrm{W}^{}`$ resonance. The effect is larger for the hadronic invariant mass since in this case, two decay fermions (the two quarks) can be combined with the photon. The effect of the squared charges of the final-state fermions is marginal in this case because the contribution of initial-state fermions dominates.
In Fig. 7 (left) we show the 4f-Born predictions for the $`M(W^{})`$ distributions, without radiative corrections, i.e. the invariant mass distributions are constructed from all diagrams without the restriction to the CC03 diagrams. The ratio 4f-Born/CC03-Born, shown in Fig. 7 (right) for the $`udsc`$ final state, confirms the goodness of the CC03 approximation for final states involving no electrons in describing the $`WW`$ cross-section at LEP 2 energies, especially in the peak region. The ratio best-with-4F-Born/best-with-CC03-Born is nearly the same as the one shown on the r.h.s. of Fig. 7, since the corrections contained in the numerator and the denominator are the same. In Fig. 8 we show the ratio best-with-4f-Born/CC03-Born for both the bare (right) and the calo (left) $`W^{}`$ invariant-mass distributions, exhibiting the combined effect of including radiative corrections and background diagrams.
Further numerical results from RacoonWW can be found in Ref. and, for the same set-up as here, in Sect. 4.
#### CC03 with KORALW/YFSWW
#### Authors
#### General Description
The program KORALW1.42 has been fully documented and published in Ref. . Here one can find the differences between YFSWW3 and KORALW in terms of radiative corrections.
Thus, here we describe YFSWW3 first. This latter program evaluates the the double resonant process $`e^+e^{}W^+W^{}4f`$ in the presence of multiple photon radiation using Monte Carlo event generator techniques. The theoretical formulation is based, in the leading pole approximation (LPA), on the exact $`𝒪\left(\alpha \right)_{\mathrm{prod}}`$ YFS exponentiation, with $`𝒪\left(\alpha \right)`$ corrections (both weak and QED) to the production process taken from Ref. , combined with $`𝒪\left(\alpha ^3\right)`$ LL ISR corrections in the YFS scheme and with FSR implemented in the $`𝒪\left(\alpha ^2\right)`$ LL approximation using PHOTOS . Anomalous $`WWV`$ couplings are supported. The Monte Carlo algorithm used to realize the YFS exponentiation is based on the YFS3 algorithm presented in Ref. and in Ref. . This algorithm is now described in detail in Ref. . In this way, one achieves an event-by-event realization of our calculation in which arbitrary detector cuts are possible and in which infrared singularities are cancelled to all orders in $`\alpha `$. A detailed description of this work can be found in Refs. . The program KoralW 1.42 evaluates all four-fermion processes in $`e^+e^{}`$ annihilation by means of the Monte Carlo techniques. It generates all four-fermion final states with multi-branch dedicated Monte Carlo pre-samplers and complete, massive, Born matrix elements. The pre-samplers cover the entire phase space. Multi-photon bremsstrahlung is implemented in the ISR approximation within the YFS formulation with the $`𝒪\left(\alpha ^3\right)`$ leading-log matrix element. The anomalous $`WWV`$ couplings are implemented in CC03 approximation. The standard decay libraries (JETSET, PHOTOS, TAUOLA) are interfaced. The semi-analytical CC03-type code KorWan for differential and total cross-sections is included. It operates both in weighted (integrator) and unweighted (event generator) modes. The detailed description of this work can be found in Refs. and the long write-up of the program in Ref. .
#### Features of the Program
As the program KORALW1.42 is already published in Ref. , we again start with the features of the YFSWW3 program. The latter code is a complete Monte Carlo event generator and gives for each event the final particle four-momenta for the entire $`4f+n\gamma `$ final state over the entire phase space for each final state particle. The events may be weighted or unweighted, as it is more or less convenient for the user accordingly. The code features two realizations of the LPA, which are described in Refs. wherein we also discuss their respective relative merits.
The operation of the code is entirely analogous to that of the MC’s YFS3 and YFS2 in Refs. . A crude distribution based on the primitive Born level distribution and the most dominant part of the YFS form factors that can be treated analytically is used to generate a background population of events. The weight for these events is then computed by standard rejection techniques involving the ratio of the complete distribution and the crude distribution. As the user wishes, these weights may be either used directly with the events, which have the four-momenta of all final state particles available, or they may be accepted/rejected against a maximal weight WTMAX to produce unweighted events via again standard MC methods. Standard final statistics of the run are provided, such as statistical error analysis, total cross-sections, etc. The total phase space for the process is always active in the code.
The program prints certain control outputs. The most important output of the program is the series of Monte Carlo events. The total cross-section in $`pb`$ is available for arbitrary cuts in the same standard way as it is for YFS3 and YFS2, i.e. the user may impose arbitrary detector cuts by the usual rejection methods. The program is available from the authors via e-mail. The program is currently posted on WWW at http://enigma.phys.utk.edu as well as on anonymous ftp at enigma.phys.utk.edu in the form of a tar.gz file in the /pub/YFSWW/ directory together with all relevant papers and documentation in postscript.
As far as the $`W`$-pair physics is concerned the KoralW is optimized to operate together with the YFSWW program: KoralW provides the complete background (beyond CC03) simulation by including all the Born level Feynman diagrams of a given process, whereas the signal process (CC03) is simulated by YFSWW including first order corrections to $`W`$ production. The final prediction is then obtained by adding and subtracting appropriate results.
In order to facilitate this add and subtract procedure both programs have been re-organized in the following way: (1) The CC03 anomalous Born matrix element and corresponding phase-space generator, covering the entire phase-space, are the same in both codes. (2) The ISR, based on YFS principle, with $`𝒪\left(\alpha ^3\right)`$ leading-log matrix element and finite transverse photon momenta is also the same in both codes (in the case of YFSWW it requires switching off the bremsstrahlung off $`W`$-pair). (3) The FSR is realized in both codes in the same way with the help of PHOTOS library. (4) The input data cards are in the same format for both codes and can be stored in one data file with common data base of parameters along with keys specific for both programs.
The features (1) – (3) guarantee that the common for both programs Born+ISR+FSR CC03 part can be defined and conveniently subtracted. This is a non-trivial feature, as for instance there are a number of different implementations of photonic cascades available amongst four-fermion Monte Carlo codes. The feature (4) is a matter of convenience as it allows for coherent and safe handling of the input parameters. For CC03, we note for clarity that YFSWW3 and KoralW 1.42 differ in that YFSWW3 has the YFS exponentiated exact NL $`𝒪\left(\alpha \right)`$ correction to the production process whereas KoralW 1.42 does not.
#### Numerical results
We start with predictions for the total cross-section, shown in Tabs.(35), where the Born approximation and the best results are shown. These results in Tabs.(35) already show the size of the NL $`𝒪\left(\alpha \right)`$ correction, $`1.52.0\%`$, when compared to the analogous results from programs such as GENTLE, see for example Ref. . In the sub-section below on the comparison between RacoonWW and YFSWW3, results such as those in Tabs.(35) are used to arrive at the current precision on the total $`WW`$ signal cross-section at LEP 2 energies.
Turning now to KORALW, we note that it has multiple-options in the presence multi-photonic events. It can define distributions for
1. visible $`\gamma `$ (radiative/hardest);
2. all photons, i.e. no cuts, in which case one can take only a) the most energetic photon to determine energy and angles (all/hardest), b) the sum (all/sum).
A sample of results is shown in Figs. 911 where we present various differential distributions for $`e^+e^{}\overline{u}d\overline{l}\nu _l`$ including all background graphs and emission of multiple photons with finite transverse momenta from initial and final states generated by KoralW. The following general cuts have been used for all plots: $`M_{ud}10`$ GeV, $`E_l5`$ GeV and $`|\mathrm{cos}\theta _l|0.985`$.
In the first plot of Fig. 9 the photon energy distributions are shown for: the hardest of all photons, the hardest of visible (radiative) photons and the sum of all photons. A visible photon is defined as having energy of at least $`1`$GeV, separated by at least $`5^{}`$ from all charged fermions and having $`|\mathrm{cos}\theta _\gamma |0.985`$. Apart from the natural big difference between visible and invisible photons one can also see a substantial effect due to emission of more than one photon (hardest vs. sum). A similar pattern for the electron final state is shown in Fig. 10. In the second plot of Fig. 9 the angular distributions of the hardest and hardest visible photon are shown.
In Fig. 11 the invariant mass distributions are shown. mass<sub>1</sub> denotes the $`ud`$-system invariant mass and mass<sub>2</sub> the $`\mu \nu _\mu `$ mass. Calo mass includes all photons that have either energy smaller than $`1`$GeV or their angle to any final state charged particle less than $`10^{}`$ for leptons or $`25^{}`$ for quarks. In the case of leptons one can see the familiar pattern of reduction of the cross-section below the peak (and weak change above) due to FSR when going from the Bare to Calo mass definition (cf. eg. Ref. ). In the case of hadrons the FSR is not generated.
#### CC03 with GENTLE
#### Authors
We describe shortly the GENTLE development after v.2.00 (1996). GENTLE v.2.10 (March 2000) , with authors D. Bardin, J. Biebel, D. Lehner, A. Leike, A. Olchevski and T. Riemann can be obtained from: http://www.ifh.de/$``$riemann/doc/Gentle/gentle.html,
/afs/cern.ch/user/b/bardindy/public/Gentle2\_10
Program developments since GENTLE v.2.00 (used in the 1996 LEP 2 workshop):
GENTLE v.2.01 (14 March 1998) compared to v.2.00:
Angular distribution (with anomalous couplings) extended from CC03 class to CC11 class .
GENTLE v.2.02 (11 Sept 1998) compared to v.2.01:
For CC cross-sections, also a constant $`W`$ width may be chosen; minor bugs eliminated.
ZAC v.0.9.4 (12.02.1999) : new package, includes anomalous couplings and calculates the angular distribution for polarized $`Z`$ pair production in the NC08 class.
GENTLE v.2.10 (March 2000) differs from v.2.02 by the following features:
– for the CC cross-sections, above threshold the Coulomb correction was modified.
– the NC cross-sections in package 4fan include now besides the NC32 class also the NC02 process; also some new options introduced, see flag descriptions below.
As is well-known, recent comparisons for the total CC03 cross-section showed that GENTLE v.2.00 overestimated it by about $`2\%`$. The reason was understood in a study made by the RacoonWW collaboration . It was found that the Coulomb correction as computed in references overestimates the FSR QED correction above the $`2W`$ threshold. Such a behaviour was not excluded, of course, because old calculations control only the leading term at threshold $`𝒪\left(1/\beta _W\right)`$, where $`\beta _W^2=14M__W^2/s`$. Only more complete calculations, using e.g. the DPA, may check how precisely the $`1/\beta _W`$ approximation works.
An introduction of a simple suppression factor
$`\text{max}(1{\displaystyle \frac{\beta _W}{\beta _W|_{\sqrt{s}=200GeV}}},\mathrm{\hspace{0.33em}0})`$ (60)
switching off the Coulomb correction smoothly between $`\sqrt{s}=2M__W`$ and $`200`$GeV improves the numerical agreement with RacoonWW considerably. In this sense, the introduction of such a fudge factor is justified by a more complete calculation based on DPA.
Compared to GENTLE v.2.00, new or extended flag regimes in GENTLE v.2.10 allow for:
(a) IFUDGF=0,1: switching the Coulomb suppression factor off/on (for IPROC=1, ie, CC);
(b) IIQCD=0,1: without/with inclusive (naive) treatment of QCD corrections (for IPROC=2, ie, NC);
(c) IIFSR=0,1,2: choice of final state QED corrections \[none, at scale $`s`$=$`M__Z^2`$, or at scale $`s_i`$\] (for IPROC=2);
(d) ICHNNL=0,1: switching between NC02 and NC32 classes (for IPROC=2);
(e) IGAMWS=0,1: switching between constant and $`s`$-dependent $`W`$ width (for IPROC=1);
(f) IINPT=2: use of the $`G_F`$ input scheme (for IPROC=2) See section 2.13, eq. (8) of : $`s_\theta ^2(1M__W^2/M__Z^2)`$, $`g^24\sqrt{2}G_FM__W^2`$.
Further, by calling subroutine WUFLAG, one may redefine the numerical value of $`\alpha _{em}(M__Z^2)`$=ALPHFS (for IIFSR=1).
Remaining electroweak corrections, genuine weak corrections in particular, are not included in GENTLE. Although, we have several choices of input parameters. We may recall here that GENTLE v.2.00 had two options: $`\alpha `$-scheme and $`G_F`$-scheme, as defined by Eqs.(71). In GENTLE v.2.10 this is extended to the NC32 family. A sample of the numerical results is shown in Tab.(6). The CC table is produced with the following GENTLE flag settings:
IPROC,IINPT,IONSHL,IBORNF,IBCKGR,ICHNNL = 1 1 1 1 0 0
IGAMZS,IGAMWS,IGAMW,IDCS,IANO,IBIN = 0 0 1 0 0 0
ICONVL,IZERO,IQEDHS,ITNONU,IZETTA = x x x 0 1
ICOLMB,IFUDGF,IIFSR,IIQCD = 2 1 0 1
IMAP,IRMAX,IRSTP,IMMIN,IMMAX = 1 0 1 1 1
As seen from the Table, there is a very good agreement between GENTLE v.2.10 and RacoonWW. It is important to emphasize, that the introduction of a suppression factor, Eq.(60), is the only modification as compared to v.2.00 which overestimated the total cross-section by about $`2\%`$. In this respect one could say that, following GENTLE’s example, all programs that do not include DPA may, nevertheless, give an effective description of CC03 that emulates the results of DPA, e.g. RacoonWW. Nevertheless, only programs including DPA represent a state-of-the-art calculation. Indeed, the Coulomb correction is just part of the full $`O(\alpha )`$ correction and cannot be split from the rest unambiguously at energies well above threshold. However, an improved Born approximation (IBA) comes significantly closer to the $`𝒪\left(\alpha \right)`$-corrected result if the Coulomb singularity is switched off above threshold with some weight function $`f(\beta _W)`$. This was already done in the IBA of Ref. , where $`f(\beta _W)`$ reduced the Coulomb part from $`2\%`$ to about $`1\%`$ at $`\sqrt{s}=200\mathrm{GeV}`$. The more radical $`f(\beta _W)`$ of (60) reduces the $`2\%`$ to zero at $`\sqrt{s}=200\mathrm{GeV}`$.
Concerning the theoretical uncertainties given in Table 6, one should understand that they are exclusively due to ISR as it is implemented within the GENTLE approach. As seen, they are of the order of $`0.75\%`$. Again, a complete approach, like the DPA, is better suited to provide a safe estimate of theoretical uncertainties.
#### Comparison between RacoonWW and BBC results
#### Authors
| RacoonWW | A. Denner, S. Dittmaier, M. Roth and D. Wackeroth |
| --- | --- |
| BBC | F. Berends, W. Beenakker and A. Chapovsky |
In this section we compare the Monte Carlo generator RacoonWW with the semi-analytical benchmark program of Berends, Beenakker and Chapovsky, called BBC in the following. The numerical comparison has been done for the leptonic channel $`\mathrm{e}^+\mathrm{e}^{}\nu _\mu \mu ^+\tau ^{}\overline{\nu }_\tau `$ and the input parameters of Ref. . As explained in more detail below, in this section the RacoonWW results are not calculated with the preferred options, but rather in a setup as close as possible to the BBC approach.
The two programs include the complete electroweak $`𝒪(\alpha )`$ corrections to $`\mathrm{e}^+\mathrm{e}^{}\mathrm{WW}4f(+\gamma )`$, both including the non-factorizable corrections and $`W`$-spin correlations, which at present is only possible within the DPA formalism. Although both programs use the DPA, nevertheless there are differences between these two calculations. One is technical, the usual difference between a flexible Monte Carlo calculation, which is also meant for experimental use, and a more rigid semi-analytical one, which was constructed as a benchmark for future calculations. The other difference is in the implementation of the DPA. The BBC calculation adheres strictly to DPA definitions, so also the phase space and photon emission are taken in DPA. In RacoonWW the matrix elements for virtual corrections are calculated in the DPA, but the exact off-shell phase space is used. For real photon radiation the DPA is not used. Instead all Born diagrams for $`\mathrm{e}^+\mathrm{e}^{}4f\gamma `$ (including the background) are taken into account and the finite width is introduced in the fixed-width scheme. Formally this procedure is not gauge-invariant, but it has been checked numerically with a gauge-invariant calculation (complex-mass scheme). The matching between the virtual and real corrections, which is necessary in order to cancel the IR and mass singularities, is done in such a way that the leading-logarithmic corrections arising from ISR are taken into account exactly, i.e. not in DPA. By comparing the two calculations one can numerically check the quality of the DPA for real-photon radiation. The expected differences in the relative corrections between both approaches are formally of $`𝒪\left(\alpha /\pi \times \mathrm{\Gamma }_\mathrm{W}/\mathrm{\Delta }E\right)`$, with $`\mathrm{\Delta }E=\sqrt{s}2M_\mathrm{W}`$ near the $`W`$-pair production threshold.
The differences in the approaches have important consequences. With RacoonWW predictions can be obtained for general cuts and physically relevant situations. The fact that the masses of the final-state fermions are neglected restricts the applicability of the program to those observables that are free of mass singularities connected to the final state. This means, in particular, that collinear photons have to be combined with the corresponding fermions. This combination depends on the experimental situation, which in turn depends on the type of final state. The semi-analytical approach is of course less flexible for implementing the experimental cuts. In the benchmark BBC calculation some of the integrations were performed analytically in order to speed up the numerical evaluations. For instance, the invariant-mass distributions were treated differently from observables where the invariant masses have been integrated over. This is not a requirement in general in the DPA if one is prepared to do more of the integrations numerically. On the other hand, a treatment of mass-singular observables, i.e. ones without photon recombination, can be easily performed in the semi-analytical approach.
For the total cross-section, the differences between the two approaches should be of the naively expected DPA accuracy of $`𝒪\left(\mathrm{\Gamma }_W/\mathrm{\Delta }E\right)`$ relative to the $`𝒪\left(\alpha \right)`$ correction. In Figure 12 we show the prediction of BBC as points with error-bars and the prediction of RacoonWW as a curve together with error-bars for some points. All error-bars are purely statistical.
As shown in Figure 12, both calculations agree very well above $`185\mathrm{GeV}`$. Below this energy the differences in the implementation of the DPA become visible, in agreement with the expected relative error of $`𝒪\left(\mathrm{\Gamma }_W/\mathrm{\Delta }E\right)`$. The main effect originates probably from the different treatment of the $`𝒪(\alpha )`$ ISR and the phase space. While BBC treat the complete $`𝒪(\alpha )`$ correction (including ISR) in DPA and use the on-shell phase space consistently, in RacoonWW<sup>10</sup><sup>10</sup>10The exponentiation of ISR has been switched off in RacoonWW for this comparison, the on-shell Coulomb singularity has been used and no naive QCD corrections are included. Moreover, the lowest-order cross-section used for normalization is calculated in DPA with on-shell phase space. This allows to compare directly the relative corrections of both approaches. the universal leading-log part of the $`𝒪(\alpha )`$ ISR correction is applied to the full CC11 cross-section, and the off-shell phase space is used throughout. Below about $`170\mathrm{GeV}`$ the DPA cannot be trusted any more for both virtual corrections and real-photon radiation, since the kinetic energy of the $`W`$ bosons becomes of the order of the $`W`$ width. The large deviations of up to $`2\%`$ in the energy range between $`170`$ and $`180\mathrm{GeV}`$ can be partially attributed to the fact that BBC treats also the leading logarithmic ISR corrections in DPA which is not done in RacoonWW. Therefore this difference cannot be viewed automatically as a theoretical uncertainty of the Monte Carlo programs.
For angular and energy distributions unavoidable differences arise from the definition of the phase-space variables in the presence of photon recombination. When defining the momenta of the $`W`$ bosons for angular distributions, BBC chooses to assign the photon to one of the production/decay sub-processes. If the detected photon is hard, $`E_\gamma \mathrm{\Gamma }_W`$, then this is theoretically possible. The error in the assignment is suppressed by $`𝒪(\mathrm{\Gamma }_W/\mathrm{\Delta }E)`$. If the detected photon is semi-soft, $`E_\gamma \mathrm{\Gamma }_W`$, then it is impossible to assign it to any of the sub-processes, but as the photon momentum is much smaller than the $`W`$-boson momentum, the error associated with this procedure is suppressed by the same relative $`𝒪(\mathrm{\Gamma }_W/\mathrm{\Delta }E)`$. The angles are then determined from the resulting $`W`$-momenta and the original fermion momenta. In RacoonWW, all angles are defined from the fermion momenta after eventual photon recombination. To this end, the invariant masses of the photon with each of the charged initial- or final-state fermions are calculated. If the smallest of these invariant masses is smaller than $`M_{\mathrm{rec}}`$ and the fermion corresponding to this invariant mass is a final-state particle, the photon is recombined with this fermion. The two different angle definitions lead to a redistribution of events in the angular distributions, which arises, in particular, from hard photon emission.
The relative corrections to the distributions in the cosines of the polar production angle, $`\theta _\mathrm{W}=\mathrm{}(\mathrm{e}^+,\mathrm{W}^+)`$, and the decay angle, $`\theta _{\mu \mathrm{W}}=\mathrm{}(\mu ^+,\mathrm{W}^+)`$, are compared for $`\sqrt{s}=184\mathrm{GeV}`$ in Figure 13.
The results of BBC are again shown as points with error-bars. The results of RacoonWW are plotted as histograms for two different photon recombination cuts $`M_{\mathrm{rec}}=5\mathrm{GeV}`$ or $`25\mathrm{GeV}`$. The relative corrections in the two recombination schemes differ at the level of $`0.5÷1\%`$, with the largest differences for large angles where the cross-section is small. The deviations between BBC and RacoonWW are somewhat larger than this and also larger than in the case of the total cross-section, but of the same order of magnitude. A repetition of the analysis at $`\sqrt{s}=250\mathrm{GeV}`$ has shown that the deviations at large angles grow with increasing centre-of-mass energy, since also the hard-photon redistribution effects grow with energy.
Invariant-mass distributions depend crucially on the treatment of the real photons. Since this is fundamentally different in RacoonWW and BBC, it does not make sense to compare these distributions between the two programs. Specifically, BBC define the $`W`$ invariant masses from the fermion momenta only (bare or muon-like) which make them sensitive to the collinear mass singularities. In RacoonWW, the photons are always recombined with the fermions (calorimetric or electron-like). The actual mass shifts crucially depend on the experimental setup. They are of the order of several $`10\mathrm{MeV}`$ and negative for the bare procedure. In the calorimetric treatment these mass shifts are reduced and can even become positive depending on the recombination procedure.
As was already mentioned earlier, the most important difference between the two approaches is the treatment of real-photon radiation. Therefore, it is important to compare distributions that are exclusive in the photon variables. As an example of such a distribution we present in Figure 14 a comparison of the photon spectrum, $`E_\gamma d\sigma /dE_\gamma `$, as a function of photon energy at the CM energy $`184\mathrm{GeV}`$.
The spectrum is shown for two different sets of angular cuts, which restrict the angles between the photon and the beam momenta, $`\mathrm{}(\mathrm{e}^\pm ,\gamma )`$, the photon and final-state lepton momenta, $`\mathrm{}(\mathrm{}^\pm ,\gamma )`$, and the beam and final-state lepton momenta, $`\mathrm{}(\mathrm{}^\pm ,\mathrm{e}^\pm )`$:
1. $`\mathrm{}(\mathrm{e}^\pm ,\gamma )>1\mathrm{deg}`$, $`\mathrm{}(\mathrm{}^\pm ,\gamma )>5\mathrm{deg}`$ and $`\mathrm{}(\mathrm{}^\pm ,\mathrm{e}^\pm )>10\mathrm{deg}`$,
2. $`\mathrm{}(\mathrm{e}^\pm ,\gamma )>50\mathrm{deg}`$, $`\mathrm{}(\mathrm{}^\pm ,\gamma )>50\mathrm{deg}`$ and $`\mathrm{}(\mathrm{}^\pm ,\mathrm{e}^\pm )>10\mathrm{deg}`$.
The first set of cuts is closer to experiment, but the second suppresses the dominant contribution of ISR in the real-photonic factorizable corrections. Since the second set of cuts removes a large part of the phase space, statistics in the first case is about ten times bigger than in the second case. However, the second set of cuts renders non-factorizable and factorizable radiation of a comparable order, thus checking the former. Figure 14 reveals an agreement between the two approaches within $`10\%`$ for both sets of cuts, which is of the order of the naive expectation for the DPA error of $`𝒪\left(\mathrm{\Gamma }_W/\mathrm{\Delta }E\right)`$. Note a peculiar decrease of the photon energy spectrum at lower photon energies for the second set of cuts. It was numerically checked in the BBC approach that this decrease is due to non-factorizable contributions (interference between various stages of the process). More precisely, the non-factorizable part amounts to roughly 20% of the complete contribution and is negative for $`E_\gamma \mathrm{\Gamma }_\mathrm{W}`$; it tends to zero above $`E_\gamma \mathrm{\Gamma }_\mathrm{W}`$.
#### Comparison of RacoonWW and YFSWW3 results
#### Authors
| RacoonWW | A.Denner, S.Dittmaier, M.Roth and D.Wackeroth |
| --- | --- |
| YFSWW3 | S. Jadach, W. Placzek, M. Skrzypek, B. Ward and Z. Was |
In this section we compare results obtained with the Monte Carlo generators RacoonWW and YFSWW3 . The numerical comparison has been done for the LEP 2 input parameter set. This comparison is restricted to the CC03 contributions for $`\mathrm{e}^+\mathrm{e}^{}\mathrm{WW}4f`$, i.e. background diagrams have been omitted<sup>11</sup><sup>11</sup>11Note that the real corrections in RacoonWW include the background diagrams of the CC11 class, and the ISR is convoluted with this class of diagrams. For LEP 2 energies, however, the difference induced by these background diagrams with respect to the Born should be at the per mille level..
First we recall that RacoonWW contains the complete electroweak $`𝒪(\alpha )`$ corrections to $`\mathrm{e}^+\mathrm{e}^{}\mathrm{WW}4f(+\gamma )`$ within the DPA, including the non-factorizable corrections and $`W`$-spin correlations. Real-photon emission is based on the full $`\mathrm{e}^+\mathrm{e}^{}4f\gamma `$ matrix element (of the CC11 class), and ISR beyond $`𝒪(\alpha )`$ is treated in the structure-function approach with soft-photon exponentiation and leading-logarithmic contributions in $`𝒪(\alpha ^3)`$. To be more precise, for $`4f`$ and $`4f\gamma `$ (with a hard non-collinear $`\gamma `$) at tree level all final states are supported, i.e. also Mix43, i.e. $`\overline{u}d\overline{d}u`$. If, however, soft and collinear photons are allowed, the virtual correction to $`e^+e^{}WW4\mathrm{f}`$ is required. In this case, RacoonWW takes photon radiation from the CC11 class into account<sup>12</sup><sup>12</sup>12To do this, in any program, for Mix43 would require virtual corrections to $`Z`$-pair production, which are not implemented.. The singular Coulomb correction is included with its full off-shell behaviour.QCD corrections are taken into account by the naive QCD factors $`(1+\alpha _\mathrm{s}/\pi )`$ for hadronically decaying $`W`$ bosons.
In YFSWW3 the exact $`𝒪(\alpha )`$ electroweak corrections to $`\mathrm{e}^+\mathrm{e}^{}\mathrm{W}^+\mathrm{W}^{}`$ are implemented together with YFS exponentiation of the corresponding soft-photon effects for the production process as defined in the DPA, which is equivalent to the LPA as defined in Ref. for this process. ISR beyond $`𝒪(\alpha )`$ is taken into account up to $`𝒪(\alpha ^3)`$ in leading-logarithmic approximation. The full off-shell behaviour of the singular Coulomb correction is included. The corrections to the $`W`$ decay, including naive QCD corrections, are implemented by using the corrected branching ratios. In this way, the total cross-section receives the full $`𝒪(\alpha )`$ corrections in DPA. Taking this cross section as normalization, final-state radiation with up to two photons is generated by PHOTOS, which is based on a leading-logarithmic (LL) approximation in which finite $`p_\mathrm{T}`$ effects are taken into account in such a way that the soft limit of the respective exact $`𝒪\left(\alpha \right)`$ $`p_\mathrm{T}`$ spectrum is reproduced.
For observables where the decay of the $`W`$ bosons and their off-shellness are integrated out, the expected differences between the two calculations are of the order of the accuracy of the DPA, i.e. of the relative order $`𝒪\left(\alpha \mathrm{\Gamma }_\mathrm{W}/\mathrm{\Delta }E\right)`$, modulo possible enhancement factors. Here $`\mathrm{\Delta }E`$ is a typical energy scale for the considered observable, i.e. $`\mathrm{\Delta }E\sqrt{s}2M_\mathrm{W}`$ for the total cross-section near the $`W`$-pair production threshold. For observables that depend on the momenta of the decay products larger differences can be expected. This holds, in particular, for observables involving a real photon. While such observables are based on the full lowest-order matrix element for $`\mathrm{e}^+\mathrm{e}^{}4f\gamma `$ in RacoonWW, in YFSWW3 the multi-photon radiation in the $`WW`$ production (within the YFS scheme) is combined with $`𝒪(\alpha ^2)`$ LL radiation in $`W`$-decays (done by PHOTOS), i.e. the real photon radiation is treated in DPA and some finite $`𝒪(\alpha )`$ terms from FSR are neglected, but the treatment of the leading logarithms goes beyond strict $`𝒪(\alpha )`$.
For the total cross-section, the differences between the two approaches should be of the naively expected DPA accuracy, i.e. below 0.5% for $`\sqrt{s}>180\mathrm{GeV}`$. In Table 7 we compare the results from both generators for the total cross-section without any cuts. The best numbers correspond to the inclusion of all corrections implemented in the programs. Independently of the channel both programs differ by $`0.2÷0.3\%`$, which is of the order of the intrinsic ambiguity of any DPA implementation, i.e. the numbers are consistent with each other.
The results of YFSWW3 presented here differ from the ones presented at the winter conferences, where still a difference of $`0.7\%`$ between the programs was reported. The main point is that the results in Table 7 are obtained with version 1.14 whereas those presented at the winter conferences were obtained with version 1.13. Version 1.14, which has benefitted from the detailed comparison between the RacoonWW and YFSWW3 virtual corrections, represents, according to renormalization group improved YFS theory , an improved re-summation of the higher order corrections as compared to version 1.13. We stress that we (the RacoonWW and YFSWW3 groups) have also checked that, when we use the same couplings, our $`𝒪\left(\alpha \right)`$ virtual plus soft corrections in the W-pair production building block agree differentially at the sub-per mille level and agree for the total cross section at $`<0.01\%`$. This is an important cross check on both programs. However, as a by-product of this detailed comparison, we have realized that the $`G_F`$ scheme of Refs. has only the IR divergent part of the virtual photonic corrections with coupling $`\alpha (0)`$ whereas the renormalization group equation implies that any photon of 4-momentum $`q`$ should couple completely with $`\alpha (0)`$ when $`q^20`$, where $`\alpha (q^2)`$ is the running renormalized QED coupling. In version 1.14 of YFSWW3, we have made this improvement as implied by the renormalization group equation . The generic size of the resulting shift in the YFSWW3 prediction can be understood by isolating the well-known soft plus virtual LL ISR correction to the process at hand, which has in $`𝒪\left(\alpha \right)`$ the expression
$$\delta _{ISR,LL}^{v+s}=\beta \mathrm{ln}k_0+\frac{\alpha }{\pi }\left(\frac{3}{2}L+\frac{\pi ^2}{3}2\right),$$
(61)
where $`\beta \frac{2\alpha }{\pi }(L1)`$$`L=\mathrm{ln}(s/m_e^2)`$, and $`k_0`$ is a dummy soft cut-off which cancels out of the cross section as usual. In the $`G_F`$ scheme of Refs. which is used in YFSWW3-1.13, only the part $`\beta \mathrm{ln}k_0+(\alpha /\pi )(\pi ^2/3)`$ of $`\delta _{ISR,LL}^{v+s}`$ has the coupling $`\alpha (0)`$ and the remaining part of $`\delta _{ISR,LL}^{\mathrm{v}+\mathrm{s}}`$ has the coupling $`\alpha _{G_F}\alpha (0)/(10.0371)`$. The renormalization group improved YFS theory implies, however, that $`\alpha (0)`$ should be used for all the terms in $`\delta _{ISR,LL}^{v+s}`$. This is done in YFSWW3-1.14 and results in the normalization shift $`\left((\alpha (0)\alpha _{G_F})/\pi \right)(1.5L2)`$, which at 200GeV is $`0.33\%`$. This explains most of the change in the normalization of YFSWW3-1.14 vs that of YFSWW3-1.13. Moreover, it does not contradict the expected total precision tag of either version of YFSWW3 at their respective stages of testing. We stress that, according to the renormalization group equation, version 1.14 is an improvement over version 1.13 – it better represents the true effect of the respective higher order corrections. More details of the actual scheme of renormalization and re-summation used in YFSWW3-1.14 will appear elsewhere .
In RacoonWW, the coupling $`\alpha (0)`$ is used everywhere in the relative $`𝒪\left(\alpha \right)`$ corrections, even in the $`G_F`$ scheme, in order to include the appropriate coupling for the (dominant) photonic corrections. A switch in YFSWW3-1.14 to this scheme shifts the maximal differences between the programs to $`0.34\%`$, somewhat larger than the $`0.27\%`$ shown. This confirms the expectation that the effects of unknown higher-order corrections are at the level of $`0.1\%`$.
It should be noted that the results in Table 7 lie by $`2÷3\%`$ below the LL-type predictions given by GENTLE (see also Section 4). As stated above, however, this consideration only applies to GENTLE in some special setup. The disagreement with all other codes active in the ‘95 workshop is within $`1.5\%`$. The fact that two independent Monte Carlo calculations with physical precision at the level of $`𝒪(\frac{\alpha }{\pi }\frac{\mathrm{\Gamma }_W}{M__W})`$ now agree to $`0.2÷0.3\%`$ at $`200`$GeV for this total cross section is truly an important improvement over the situation in the ’95 workshop .
In the following we consider observables obtained with the cut and photon recombination procedure as given in the description of numerical results of RacoonWW in Sect. 4.1. We again consider the cases of a tight recombination cut $`M_{\mathrm{rec}}=5\mathrm{GeV}`$ (bare) and of a loose recombination cut $`M_{\mathrm{rec}}=25\mathrm{GeV}`$ (calo).
Table 8 shows the analogous cross-sections to Table 7 but now with the described bare cuts applied. The difference of $`0.2÷0.3\%`$ between the two compared programs does not change by the applied cuts. When turning from bare to calo cuts the results for the cross-sections do not change significantly; of course, the lowest-order results do not change at all.
In the following relative corrections to various distributions for the semi-leptonic channel $`\mathrm{e}^+\mathrm{e}^{}\mathrm{u}\overline{\mathrm{d}}\mu ^{}\overline{\nu }_\mu `$ are compared at $`\sqrt{s}=200\mathrm{GeV}`$. All these distributions have been calculated using the above set of separation and recombination cuts.
The corrections to the cosine of the production angle for the $`\mathrm{W}^+`$ and $`\mathrm{W}^{}`$ bosons are shown in Figures 15 and 16, respectively, for the bare (left) and the calo (right) recombination schemes. The distributions are compatible with each other to better than $`1\%`$. The largest differences are of the order of $`1\%`$ and appear in general for large scattering angles.
The corrections to the invariant mass distributions for the $`\mathrm{W}^+`$ and $`\mathrm{W}^{}`$ bosons are shown in Figures 17 and 18 for the bare (left) and the calo (right) recombination scheme. The distributions are statistically compatible with each other everywhere and agree within 1%. It should be noted that the distortion of the distributions is mainly due to radiation off the final state and the $`W`$ bosons. It may seem remarkable that the LL approach of PHOTOS properly accounts for these distortion effects. But one should remember that PHOTOS was fine-tuned to describe the exact $`𝒪\left(\alpha ^1\right)`$ FSR for the radiative $`Z`$ and $`\tau `$ decays, like $`Z\mu ^{}\mu ^+(\gamma )`$ and $`\tau \mu \nu \overline{\nu }(\gamma )`$. PHOTOS was also cross-checked against the exact matrix element for the $`W\mu \nu \gamma `$ process.
Figs. 1920, and Fig. 21 show the distributions in the photon energy $`E_\gamma `$, in the cosine of the polar angle of the photon (w.r.t. the $`\mathrm{e}^+`$ axis), and in the angle between the photon and the nearest final-state charged fermion from the two programs and in the two recombination schemes.
The differences are of the order of $`15÷20\%`$. Differences of this order may be expected, since photonic observables are no corrections anymore, but belong to the class of $`\mathrm{e}^+\mathrm{e}^{}4f\gamma `$ processes, since $`\mathrm{e}^+\mathrm{e}^{}4f`$ does not contribute here. Whether or not the observed differences are consistent with the differences in the treatments of the real photon emission in the two programs is under investigation.
### 4.2 Internal estimate of theoretical uncertainty for CC03
Here we give a quantitative statement on the theoretical precision for DPA-approximation.
#### Estimating the theoretical uncertainty of the DPA with RacoonWW
### Authors
All existing calculations of electroweak corrections to $`\mathrm{e}^+\mathrm{e}^{}\mathrm{WW}4f`$ are based on DPA. A naive estimate of the accuracy of this approach yields $`(\alpha /\pi )\times \mathrm{ln}(\mathrm{})\times \mathrm{\Gamma }_\mathrm{W}/M_\mathrm{W}`$, where $`\mathrm{\Gamma }_\mathrm{W}/M_\mathrm{W}`$ is the generic accuracy of the DPA, $`\alpha /\pi `$ results from considering one-loop corrections, and $`\mathrm{ln}(\mathrm{})`$ represents leading logarithms or other possible enhancement factors in the corrections. This naive estimate suggests that the DPA has an uncertainty of some $`0.1\%`$. Note, however, that this estimate can fail whenever small scales become relevant. In particular near the $`W`$-pair threshold, the estimate should be replaced by $`(\alpha /\pi )\times \mathrm{ln}(\mathrm{})\times \mathrm{\Gamma }_\mathrm{W}/(\sqrt{s}2M_\mathrm{W})`$.
In order to investigate the accuracy of the DPA quantitatively, a number of tests have been performed with RacoonWW. The implementation of the DPA has been modified within the formal level of $`\alpha \mathrm{\Gamma }_\mathrm{W}/M_\mathrm{W}`$, and the obtained results have been compared. Note that in RacoonWW only the virtual corrections are treated in DPA, while real photon emission is based on the full $`\mathrm{e}^+\mathrm{e}^{}4f\gamma `$ matrix element with the exact five-particle phase space. Thus, only the $`24`$ part is effected by the following modifications. Specifically, three types of uncertainties have been considered (see Ref. for details):
* Different on-shell projections:
In order to define a DPA one has to specify a projection of the physical momenta to a set of momenta for on-shell $`W`$-pair production and decay<sup>13</sup><sup>13</sup>13This option only illustrates the effect of different on-shell projections in the four-particle phase space; if real photonic corrections are treated in DPA the impact of different projections may be larger.. This can be done in an obvious way by fixing the direction of one of the $`\mathrm{W}`$ bosons and of one of the final-state fermions originating from either $`\mathrm{W}`$ boson in the CM frame of the incoming $`\mathrm{e}^+\mathrm{e}^{}`$ pair. The default in RacoonWW is to fix the directions of the momenta of the fermions (not of the anti-fermions) resulting from the $`\mathrm{W}^+`$ and $`\mathrm{W}^{}`$ decays (def). A different projection is obtained by fixing the direction of the anti-fermion from the $`\mathrm{W}^+`$ decay (proj) instead of the fermion direction.
* Treatment of soft photons:
In RacoonWW, the virtual photon contribution is treated in DPA, while real photon radiation is fully taken into account. These two contributions have to be matched in such a way that IR and mass singularities cancel. This requirement only fixes the universal, singular parts, but leaves some freedom to treat non-universal, non-singular contributions in DPA or not. For instance, in the branch of RacoonWW that employs the subtraction method of Ref. , the endpoint contributions of the subtraction functions are calculated in DPA and added to the virtual photon contribution as default. As an option, RacoonWW allows to treat also the universal (IR-sensitive) part of the virtual photon contribution off-shell by extracting an U(1)-invariant factor à la YFS from the virtual photon contribution and adding it to the real photon contribution, i.e. this soft+virtual part of the photonic correction is treated off shell (eik). The two described treatments only differ by terms of the form $`(\alpha /\pi )\times \pi ^2\times 𝒪\left(1\right)`$ which are either multiplied with the DPA (def) or with the full off-shell Born cross-sections (eik).
* On-shell versus off-shell Coulomb singularity:
The Coulomb singularity is (up to higher orders) fully contained in the virtual $`𝒪\left(\alpha \right)`$ correction in DPA. Performing the on-shell projection to the full virtual correction leads to the on-shell Coulomb singularity. However, since the Coulomb singularity is an important correction in the LEP 2 energy range and is also known beyond DPA, RacoonWW includes this extra off-shell Coulomb correction as default. Switching the extra off-shell parts of the Coulomb correction off (Coul), yields an effect of the order of the accuracy of the DPA.
In the following table and figures the total cross-section and various distributions have been compared for the different versions of the DPA defined above. The results have been obtained using the LEP 2 input parameter set and the set of separation and recombination cuts as given in the description of numerical results of RacoonWW in Sect. 4.1. The recombination cut is chosen to be $`M_{\mathrm{rec}}=25\mathrm{GeV}`$. As default, we take the RacoonWW results (best-with-4f-Born) of Sect. 4 for the process $`\mathrm{e}^+\mathrm{e}^{}\mathrm{u}\overline{\mathrm{d}}\mu ^{}\overline{\nu }_\mu (\gamma )`$ at $`\sqrt{s}=200\mathrm{GeV}`$, which are based on the above input. The only differences are that the naive QCD factors and ISR corrections beyond $`𝒪\left(\alpha \right)`$ are not included in the results of this section. The results for the total cross-section are shown in Table 9.
Note that these cross-sections are calculated with the above cuts. We find relative differences at the level of $`0.1\%`$. As expected, the prediction that is based on the on-shell Coulomb correction is somewhat higher than the exact off-shell treatment, since off-shell effects screen the positive Coulomb singularity. The results in Table 9 have been obtained using phase-space slicing for the treatment of the IR and collinear singularities. If the subtraction method is used instead, the resulting cross-section is about 0.01% smaller.
In Figures 22 and 23 we show the differences of the proj, eik, and Coul modifications to the default version of the DPA for some distributions. For the distribution in the cosine of the W-production angle $`\theta _{\mathrm{W}^+}`$ and in the W-decay angle $`\theta _{\mathrm{W}^{}\mu ^{}}`$ (see Figure 22) the relative differences are of the order of $`0.1÷0.2\%`$ for all angles, which is of the expected order for the intrinsic DPA uncertainty.
For the $`\mu `$-energy distribution, shown in the l.h.s. of Figure 23, the differences are typically of the same order, as long as $`E_\mu `$ is in the range for $`W`$-pair production, which is $`20.2\mathrm{GeV}<E_\mu <79.8\mathrm{GeV}`$ at $`\sqrt{s}=200\mathrm{GeV}`$. Outside this region, the four-fermion process is not dominated by the $`W`$-pair diagrams, and the DPA is not reliable anymore, which is also indicated by large intrinsic ambiguities. The r.h.s. of Figure 23 shows the DPA uncertainties for the $`\mathrm{u}\overline{\mathrm{d}}`$ invariant-mass distribution. Within a window of $`2\mathrm{\Gamma }_\mathrm{W}`$ around the W resonance the relative differences between the considered modifications are also at the level of $`0.1÷0.2\%`$. The differences grow with the distance from the resonance point.
The discussed results illustrate that the intrinsic ambiguities of the DPA, as applied in RacoonWW, are at the level of a few per mil, whenever resonant $`W`$-pair production dominates the considered observable.
#### Estimating the theoretical uncertainty of the DPA with YFSWW3-KoralW
The accuracy of the combined result from YFSWW3 1.13 and our all 4-fermion process MC KoralW 1.42 as presented in Ref. is expected to be below $`0.5\%`$ for the total cross-section when all tests are finished. These tests are currently in progress.
### 4.3 Summary and conclusions
In this Section we have compared different theoretical predictions for the CC03 cross-section that have been used to analyze the data in terms of all $`W`$-pair final states, $`4q(qqqq)`$ and non-$`4q(qq\mathrm{l}\nu ,\mathrm{l}\nu \mathrm{l}\nu )`$. The major achievement in this area is represented by inclusion of radiative corrections in DPA for the $`WW`$ cross-section.
Data are collected from $`161`$GeV up to $`210`$GeV. One should remember that below some threshold ($`170\mathrm{GeV}`$) the DPA cannot be trusted any more for both virtual corrections and real-photon radiation, since the kinetic energy of the $`W`$ bosons becomes of the order of the $`W`$ width. RacoonWW has shown that the intrinsic ambiguities of its implementation of the DPA are at the level of a few per mille.
For the total CC03 cross-section, the differences between RacoonWW and YFSWW should be of the naively expected DPA accuracy, i.e. below 0.5% for $`\sqrt{s}>180\mathrm{GeV}`$. And, indeed, independently of the channel, the two MC differ by $`0.2÷0.3\%`$ in the results presented herein and this increases to 0.4% if uncertainties from unknown higher-order corrections are taken into account. Note that, with bare cuts applied, the difference of $`0.2÷0.3\%`$ shown here between the two compared programs does not change.
The corrections to the distribution in the cosine of the production angle for the $`\mathrm{W}^+`$ and $`\mathrm{W}^{}`$ bosons have also been analyzed for the bare and the calo recombination algorithms. They are compatible with each other at a level below 1%. Although compatible with the statistical accuracy, the deviations seem to become somewhat larger for large scattering angles. The corrections to the invariant mass distributions for the $`\mathrm{W}^+`$ and $`\mathrm{W}^{}`$ bosons, again with bare and calo recombinations are statistically compatible between the two Monte Carlo programs everywhere and agree within 1%.
Another comparison, shown in Figure 12, indicates that RacoonWW and BBC calculations agree very well for the total $`W`$-pair production cross-section above $`185\mathrm{GeV}`$. Below this energy the differences in the implementation of the DPA become visible, in agreement with the expected relative error of $`𝒪\left(\alpha /\pi \times \mathrm{\Gamma }_W/\mathrm{\Delta }E\right)`$. However, for angular and energy distributions unavoidable differences at the level of $`1÷2\%`$ arise between the two predictions, as a consequence of the definition of the phase-space variables in the presence of photon recombination. Although the BBC-calculation has not been implemented in a MonteCarlo it can be used for obtaining a relative $`𝒪\left(\alpha \right)`$ correction factor where one has an estimated internal accuracy ranging from $`1.5\%`$ at lower energies to $`0.3\%`$ at $`210`$GeV.
In conclusion, from the direct comparisons of RacoonWW and YFSWW3, supported by BBC, we can estimate an overall theoretical uncertainty of the current predictions for the total $`WW`$ cross-section at $`0.4\%`$ at $`200\mathrm{GeV}`$. The $`0.4\%`$ precision tag is an important conclusion of this Workshop.
For other energies no complete investigations of the theoretical uncertainty have been performed. However, based on the error estimate of $`0.4\%`$ for $`200\mathrm{GeV}`$, the intrinsic uncertainty of the DPA of $`0.2\%`$ at $`200\mathrm{GeV}`$ and the generic energy dependence of this uncertainty given by $`\mathrm{\Gamma }__W/(E_{\mathrm{CMS}}2M_\mathrm{W})`$ we estimate an uncertainty of the predictions of RacoonWW and YFSWW3 of $`0.5\%`$ for $`180\mathrm{GeV}`$ and $`0.7\%`$ for $`170\mathrm{GeV}`$. This could be somewhat further reduced, if the sources of the differences between the different programs are found.
Results for the $`WW`$ cross-section at $`𝒪\left(\alpha \right)`$ are also available from GRACE but a comparison with the other codes is not yet at the level of those already presented where a considerable amount of time was invested to try to understand differences towards a safe estimate of theoretical uncertainty.
## 5 Four fermions plus a visible photon
The class of processes that are investigated at LEP 2 are $`e^+e^{}W^+W^{}4`$f, single-$`W`$ production, $`Z`$-boson-pair production, single-$`Z`$ production. LEP 2 and also future linear colliders will allow us to study a new class of processes, $`e^+e^{}4\mathrm{f}+\gamma `$.
The physical interest of the latter is twofold. They can be used to obtain informations on the quartic gauge-boson couplings and include the production processes of three gauge-bosons, $`W^+W^{}\gamma `$, $`ZZ\gamma `$ and $`Z\gamma \gamma `$. In this case the photon is visible by definition and we term the corresponding process radiative, i.e. we consider as radiative events those events with photons where at least one photon passes the experimental photon requirements, for instance $`E_\gamma >1`$GeV, $`\mathrm{cos}\theta _\gamma <0.985(0.997)`$ and $`\theta _{f\gamma }>5^{}`$.
Note that for all final states, the invariant mass needs a more precise definition in case radiative photons are present in the event. From a calculational point of view, there is always a minimal invariant mass (energy and separation angle) below which photons are not resolved. Thus we need to specify fermion-photon invariant mass or fermion-photon energies and separation angles, below which the the photon are combined with the fermion and above which the photons are not included in the mass calculation. A bare mass would set these cuts rather tight, excluding photons from the $`f\overline{f}`$ mass, a calo mass would set the separation cuts looser. Theorists like cuts on $`M(\gamma \mathrm{nearestf})`$. Experimentalists like cuts on energies and angles. In the following we list both TH(eory)-cuts and EXP(erimental)-cuts.
:
* $`M(\overline{f}f+(\gamma ))`$ including photons if $`M(f+\gamma )<5`$GeV,
* $`M(\overline{f}f+(\gamma ))`$ including photons if $`M(f+\gamma )<25`$GeV.
:
* : $`M(\overline{f}_1f_2+(\gamma ))`$, photons less than 1 GeV or less than $`1^{}`$ away from $`f_1`$ or $`f_2`$ are included;
* : $`M(\overline{l}_1l_2+(\gamma ))`$, photons less than 1 GeV or less than $`10^{}`$ away from charged leptons are included,
$`M(\overline{q}_1q_2+(\gamma ))`$, photons less than 1 GeV or less than $`25^{}`$ away from either quark $`q_1,q_2`$ are included, which takes at least the major difference between fermions - quarks versus leptons - into account.
These definitions serve for benchmarking distributions, not so much to mimic an actual experimental strategy, which is of course fermion dependent. In other words this is an approximation to the experimental side: if the fermion is a muon, even $`0^{}`$ opening angles can be separated experimentally. In addition, for identified photons one still may or may not choose to recombine the photon with the fermion.
Furthermore, $`e^+e^{}4\mathrm{f}+\gamma `$ is an important building block for the radiative corrections to the Born process $`e^+e^{}4`$f, hence non-radiative events are those with no photon or only photons below the minimal photon requirements. In case of non-radiative events, this amounts to adding up virtual and soft radiative corrections. The effect of $`𝒪\left(\alpha \right)`$ QED corrections very often amounts to several percent, mostly originating from collinear photon radiation off highly energetic particles and from virtual photon exchange. For initial state radiation, for instance, we have three types of corrections, a) $`𝒪\left(\alpha /\pi \mathrm{ln}(m_e/Q)\right)`$ with $`Qm_e`$ being the typical scale at which the process occur, b) $`𝒪\left(\alpha /\pi \right)`$ from hard photons that must, nevertheless, be included for a $`1\%`$ precision tag, c) leading $`𝒪\left(\alpha ^2\right)`$, or higher corrections that becomes relevant for a precision tag below the $`1\%`$ thresholds.
Owing to the fact that a theoretical prediction with a typical accuracy of some fraction of a percent must include all QED corrections, we face the complexity of it. Handling the singularities of the squared matrix element represents a formidable task; in any bremsstrahlung process the integrand blows up for arbitrary small photon energies and similar problems arise from collinear emission off the charged particles.
A general comment about this section is that some of the programs, but not all, implement $`4\mathrm{f}+\gamma `$ at the level of (exact) matrix elements. Few programs have only an effective treatment of photons via structure functions, with or without $`p_t`$. Furthermore we also have to distinguish between massless vs. massive calculations.
### 5.1 Description of the programs and their results
#### $`4\mathrm{f}+\gamma `$ with RacoonWW
#### Authors
#### General description
The program RacoonWW evaluates cross-sections and differential distributions for the reactions $`\mathrm{e}^+\mathrm{e}^{}4\mathrm{f}`$ and $`\mathrm{e}^+\mathrm{e}^{}4\mathrm{f}+\gamma `$ for all four-fermion final states. The long write-up has already been presented in Sect. 4.1, so that we only stress the features that are peculiar to $`4\mathrm{f}+\gamma `$ production with a separated hard photon. The calculation is based on full $`4\mathrm{f}+\gamma `$ matrix elements for all final possible states. Since fermion masses are neglected, lower cuts on the invariant mass of $`\mathrm{f}\overline{\mathrm{f}}`$ pairs and on $`\mathrm{e}^\pm `$ emission angles have to be imposed, in addition to the angular and energy cuts for the hard photon. RacoonWW supports different ways to treat finite gauge-boson widths (fixed and running widths, complex-mass scheme) and allows to select subsets of graphs ($`VV\gamma `$ signal diagrams, QCD background). Detailed numerical results on $`4\mathrm{f}+\gamma `$ production with RacoonWW can be found in Ref. and in Sect. 5.2.
#### $`4\mathrm{f}+\gamma `$ with PHEGAS/HELAC
#### Author
This section refers to a novel Monte Carlo program that is capable to deal with any tree-order process involving any particle and interaction described by the Standard Model, including QCD.
The program consists of two modules:
1. HELAC which is a matrix element computation-tool based on Dyson-Schwinger equations, and
2. PHEGAS an automatic phase-space generator capable to simulate all peaking structures of the amplitude.
The over all code is using a Monte Carlo integration based on multichannel optimization .
#### HELAC
The matrix element is evaluated using a recursive approach based on Dyson-Schwinger equations. The computational cost exhibits an exponential growth ($`3^n`$) as a function of the number of external particles ($`n`$) which for multi-particle processes results to a very important increase in the efficiency as compared with the traditional Feynman-graph approach whose computational cost grows factorially ($`n!`$). In order to optimize code’s efficiency the computational strategy consists of two phases. In the first phase a solution to the recursive equations is established in terms of an integer array containing all relevant information for the process under consideration. This is the initialization phase and is performed once at the beginning of the execution of the program. In the second phase, using the already generated information, the actual computation is performed resulting to the numerical evaluation of the amplitude for each specific phase-space point provided.
In order to consistently describe unstable particles the fixed width as well as the complex width schemes have been included. ISR and running couplings are also an option and work is in progress to implement higher order corrections within the approach of reference .
In order to deal with numerical stability problems, besides the double precision, a quadruple as well as a multi-precision version is available. This makes HELAC able to deal with processes exhibiting strong collinear singularities, like $`e^{}e^+e^{}e^+\mu ^{}\mu ^+`$ at zero scattering angles. Moreover all particle masses and vertices of the Standard Model, including QCD, in both the Feynman and unitary gauges are incorporated.
#### PHEGAS
Although several matrix element computational tools were available in the past that can deal with arbitrary processes, to the one or to the other extent , phase-space generators were always developed according to a specific process or a class of processes . PHEGAS is a phase space generator that incorporates in an automatic way all possible kinematical mappings for any given process, using the relevant information provided by HELAC. To this end each Feynman graph contributing to the process under consideration gives rise to a kinematical mapping. The integration is performed via a Monte Carlo multichannel approach and during the computation, weight optimization selects automatically those kinematical mappings that are relevant for the process under consideration.
As a first highly non-trivial test PHEGAS/HELAC has been used to produce results for four-fermion plus a visible photon within the current study. Nevertheless, it is worthwhile to emphasize that PHEGAS/HELAC is able to deal with any process involving any Standard Model particle and is by no means restricted to $`e^+e^{}\text{4 f}+\gamma `$ reactions. A detailed presentation of the code, the implemented algorithms as well as the incorporated physics effects will be available in the near future .
#### $`4\mathrm{f}+\gamma `$ with WRAP
#### Authors
#### Description of the Method.
Contributions of the Pavia/ALPHA group to the subject of four fermions plus gamma final states are summarized.
#### Hard-scattering matrix element
The exact tree-level matrix elements for the processes with four fermions plus a visible photon in the final state are computed by means of the ALPHA algorithm . At present, the processes which can be mediated by two $`W`$-bosons ( CC processes) or by two $`Z`$-bosons ( NC processes) are accounted for. The effect of finite fermion masses is taken into account exactly both in the kinematics and dynamics. The contribution of anomalous trilinear gauge couplings can be also simulated, after having implemented in ALPHA and cross-checked the parameterization in terms of $`\mathrm{\Delta }k_\gamma `$, $`\lambda _\gamma `$, $`\delta _Z`$, $`\mathrm{\Delta }k_Z`$ and $`\lambda _Z`$ of refs. . The genuinely anomalous quartic gauge boson couplings, involving at least one photon and relevant for this process at tree-level, are also included, according to the parameterization of Ref. . Final cross-checks on anomalous quartic couplings are in progress. The fixed-width scheme is adopted as gauge-restoring approach, as motivated in comparison with other gauge-invariance-preserving schemes in Ref. .
#### Radiative corrections
The phenomenologically relevant Leading Log (LL) QED radiative corrections, due to initial-state radiation (ISR), are implemented via the Structure Function (SF) formalism , according to the two following options:
* collinear SF $`D(x,s)`$;
* $`p_t`$-dependent SF $`\stackrel{~}{D}(x,\mathrm{cos}\theta _\gamma ;s)`$, i.e. a combination of the collinear SF $`D(x,s)`$ with an angular factor for photon radiation inspired by the leading behaviour $`1/(pk)`$ .
In fact, as discussed in detail in refs. , due to the presence of an observed photon in the final state, the treatment of ISR in terms of collinear SF turns out to be inadequate because affected by double counting between the pre-emission photons (described by the SF) and the observed one (described by the hard-scattering matrix element).<sup>14</sup><sup>14</sup>14In the tuned comparison with RacoonWW the effect of ISR SF was switched off. By keeping under control also the transverse degrees of freedom of ISR, as allowed by $`p_t`$-dependent SF, it is possible to remove the double-counting effects, following the procedure for the calculation of the QED corrected cross-section discussed in Ref. , i.e.
$$\sigma _{QED}^{4\mathrm{f}+1\gamma }=𝑑x_1𝑑x_2𝑑c_\gamma ^{(1)}𝑑c_\gamma ^{(2)}\stackrel{~}{D}(x_1,c_\gamma ^{(1)};s)\stackrel{~}{D}(x_2,c_\gamma ^{(2)};s)\mathrm{\Theta }(\mathrm{cuts})𝑑\sigma ^{4\mathrm{f}+1\gamma },$$
(62)
where $`c_\gamma ^{(i)}\mathrm{cos}\theta _\gamma ^{(i)}`$, $`i=1,2`$. According to eq. (62), an equivalent photon is generated for each colliding lepton and accepted as a higher-order ISR contribution if:
* the energy of the equivalent photon is below the threshold for the observed photon $`E_\gamma ^{\mathrm{min}}`$, for arbitrary angles; or
* the angle of the equivalent photon is outside the angular acceptance for the observed photons, for arbitrary energies.
Within the angular acceptance of the detected photon, the cross-section is evaluated by means of the exact matrix element for the processes $`e^+e^{}4\mathrm{f}+\gamma `$. Therefore, eq. (62) applies to the signature of four fermions plus exactly one photon in the final state, corrected by the effects of undetected soft and/or collinear ISR. The $`Q^2`$-scale entering the QED SF is fixed to be $`Q^2=s`$.
#### Computational tool and obtained results
The theoretical features sketched above have been implemented into a massive MonteCarlo (MC) program, named WRAP (W Radiative process with Alpha & Pavia). The multi-channel importance sampling technique is employed to perform the phase-space integration, paying particular attention to the infrared and collinear peaking structures due to photon emission. The code supports realistic event selections and can be employed either as a cross-section calculator or as a true event generator. Results obtained in the present study can be summarized as follows: We have performed a critical analysis of the effect of ISR (see Figs. 2426) and a study of the impact of finite fermion masses (see Tab.(10)), Finally, we have tuned comparisons with the predictions of other codes, especially with RacoonWW (see Sec. 5.2).
The impact of ISR via collinear SF on the $`4\mathrm{f}+\gamma `$ integrated cross-section of the CC10 final state $`\mu ^{}\overline{\nu }_\mu u\overline{d}\gamma `$ is shown in Figs. 2425, as a function of the LEP 2 c.m.s. energy (Fig. 24) and of the photon energy threshold at $`\sqrt{s}=192`$ GeV (Fig. 25). Fig. 24 shows that ISR in the collinear approximation reduces the Born cross-section between $`1612\%`$ in the c.m.s. range $`180190`$GeV and at the $`10\%`$ level close to $`200`$GeV, for the considered photon separation cuts. In particular, at $`\sqrt{s}=`$ 192 GeV the reduction factor as due to ISR is $`1213\%`$, almost independent of the photon detection threshold, as shown in Fig. 25.
Note that collinear SF contradicts photon detection criteria, as discussed before. However, in order to get a first estimate of the correction due to ISR, collinear SF can be used, since the error introduced by this treatment (double-counting effects) is estimated in Fig. 26, by comparing collinear and $`p_t`$ structure functions.
As far as fermion masses are concerned we show in Tab.(10) a comparison between the cross-section for the final state $`\mu ^{}\overline{\nu }_\mu c\overline{s}\gamma `$ in the massless approximation is compared with the same cross-section in the presence of finite masses for the final state fermions. The mass values and cuts used are: $`m_\mu =0.105`$ GeV, $`m_s=0.3`$ GeV, $`m_c=1.55`$ GeV, with $`M_{cs}3`$ GeV. In the considered channel with a muon in the final state, the minimum separation angle between the quarks and the photon is maintained fixed at $`5^{}`$, while the separation angle between the muon and the photon is varied from $`1^{}`$ down to zero. It can be seen that the mass effects on the the integrated cross section are of the order of $`1\%`$ for not too small separation angles, but it may reach, not surprisingly, the $`10\%`$ level in more stringent conditions, where only a massive $`4\mathrm{f}+\gamma `$ calculation can provide a reliable prediction in the presence of muons in the final state.
#### $`4\mathrm{f}+\gamma `$ with CompHEP
#### Authors
#### General description
The program CompHEP calculates cross-sections and distributions for all channels $`e^+e^{}4f`$ and $`e^+e^{}4f+\gamma `$. The calculation is based on a tree-level matrix element for the complete set of diagrams. Finite fermion masses are taken into account both in the matrix element and in the four or five particle phase space parameterization. The fixed-width prescription is used for the gauge boson propagators. In so far as CompHEP uses the squared diagrams technique, the calculation for the five particle states with radiative gamma is CPU time consuming and in the following only the results for the channel $`e^+e^{}\gamma \mu \overline{\nu }\overline{\nu }_\mu u\overline{d}`$ (2556 squared diagrams) are presented (Fig. 27, Fig. 28, where the factor $`\alpha (0)/\alpha _{G_F}`$ is not accounted for). We used the standard set of cuts including EXP-cuts for the distributions in the bare and calo mass.
#### On-shell $`W`$ boson approximation for $`e^+e^{}\gamma \mu \overline{\nu }_\mu u\overline{d}`$
In the $`24`$ approximation of the on-shell $`W`$ boson $`e^+e^{}\gamma \mu \overline{\nu }_\mu W^+`$ for the $`25`$ process $`e^+e^{}\gamma \mu \overline{\nu }_\mu u\overline{d}`$ the number of diagrams is much smaller ($`31`$ for the $`4`$-body and $`71`$ for the $`5`$-body final state). It is interesting to find out if a simpler on-shell $`W`$ approximation reproduces with enough likelihood the total rate and distributions given by the exact $`25`$ tree level amplitude. The possibility to describe quantitatively the $`5`$-body distributions of radiative events by some trivial change of the normalization in the $`4`$-body results could be attractive.
We calculated the cross section of the process $`e^+e^{}\gamma \mu \overline{\nu }_\mu W^+`$ multiplied by a factor given by the following on-shell $`W`$ isotropic decay to $`u\overline{d}`$. Vectors of the $`u`$, $`\overline{d}`$ quarks momenta generated randomly in the rest frame of the $`W`$ were boosted to the $`e^+e^{}`$ c.m.s., where the standard kinematical cuts were introduced: $`E_\gamma `$ 1 GeV, $`E_\mu 5`$GeV, $`|\mathrm{cos}\theta (\gamma e)|0.985`$. Furthermore, $`|\mathrm{cos}\theta (\mu e)|0.985`$, $`\theta (\gamma ,\mu )`$, $`\theta (\gamma ,u)`$, and $`\theta (\gamma ,\overline{d})`$ $``$ 5. Such a scheme of calculation is based on the well-known approximation of infinitely small $`W`$ width $`M__W\mathrm{\Gamma }_{\mathrm{tot}}/[(M_{u\overline{d}}^2M__W^2)^2+M__W^2\mathrm{\Gamma }_{tot}^2]\pi \delta (M_{u\overline{d}}^2M__W^2)`$ and have been widely used for the simulation of the $`3`$\- and $`4`$-body final states in many generators. The simulation by PYTHIA generator follows slightly better scheme, when the $`W`$ decay products invariant mass is distributed according to the Breit-Wigner and gamma radiation from quarks can be switched on in the approximation of final state shower.
The total rate of the $`\sigma (e^+e^{}\gamma \mu \overline{\nu }_\mu W^+)`$ Br$`(W^+u\overline{d})`$ is equal to $`49.4(2)`$fb to be compared with the exact $`25`$ result $`69.1(9)`$ fb. Missing contribution of the omitted diagrams, especially from the phase space regions near the collinear and infrared poles of the photons radiated from the initial state and the $`u`$, $`d`$ quarks leads to substantial underestimate of the rate. Peaks of the forward and back-scattered photons (Fig. 27), radiated from the initial $`e^+`$, $`e^{}`$, are much stronger underestimated than the photon distribution in the central rapidity region. Distributions in the quark energy and transverse momentum (upper plots in Fig. 28) are rather different in the exact and approximate calculation. For the exact calculation the quark energy spectrum more rapidly decreases than for the approximation where the photon radiation from quarks is not accounted for. In the exact $`5`$-body consideration the $`W`$ boson is created in a rather well defined polarization state, so the approximation of an isotropic on-shell $`W`$ decay could be unsatisfactory for angular variables. Large difference of the distributions in the photon-fermion (muon or quark) angle (lower plot in Fig. 27) is caused by a simple combinatorial reason. Calo jet-jet mass (lower plot in Fig. 28) contains the unresolved photon radiated from the initial state or from the muon, so only $`M_{u\overline{d}\gamma }M__W`$ is possible.
It follows that in the case of four fermion events with radiative photon the approximation of the on-shell $`W`$ isotropic decay does not, generally speaking, satisfactorily describe both the total rate and the full set of final particle distributions.
#### $`4\mathrm{f}+\gamma `$ via Structure Functions with NEXTCALIBUR
#### Authors
In this Section we show illustrative results for the processes $`e^+e^{}\mu ^{}\mu ^+u\overline{u}(\gamma )`$ ($`ZZ`$ signal) and $`e^+e^{}\mu ^{}\overline{\nu }_\mu u\overline{d}(\gamma )`$ ($`WW`$ signal). Analogous results for the single-$`W`$ case can be found in section 6.
NEXTCALIBUR does not contain the exact matrix element for $`e^+e^{}4f+\gamma `$, therefore we generate photons always through $`p_t`$-dependent $`ISR`$ Structure Functions. We used the set of cuts specified in the proposal at $`\sqrt{s}=200`$ GeV, all diagrams and fermion masses included. In tables 11 and 12 four values of cross section (in pb) are shown.
The first value, labelled by tot, is the sum of radiative and non radiative events (within the specified separation cuts for the generated photons). The second one nrad corresponds to non-radiative events and the third one srad to single-radiative events, namely events with only one radiated photon outside the separation cuts. We also include a fourth entry that represents the small fraction of radiative events with $`2`$ photons (drad).
To check the sensitivity of the distributions to the chosen form of Structure Function, we run again the above processes with a slightly different implementation of the sub-leading terms, without observing any significant deviation with respect to the previous results.
#### $`4\mathrm{f}+\gamma `$ with GRACE
#### Authors
In this Section we present results from GRACE for the $`4\mathrm{f}+\gamma `$ processes with $`W`$-pair and single-$`W`$ cuts. Parameters and cuts used are the same as those of the WRAP and RacoonWW collaborations, except that we used $`\alpha _{G_F}`$ for all vertices. Unfortunately, GRACE results cannot be compared directly with those of RacoonWW and WRAP; indeed, when GRACE numbers are be compared with the others one should multiply by a factor $`\alpha (0)/\alpha _{G_F}`$. To check the calculations, the following tests have been performed for the processes $`e^+e^{}\mu \overline{\nu }_\mu u\overline{d}\gamma `$ at $`\sqrt{s}=200`$GeV:
* Gauge parameter independence check; the amplitude generated by GRACE keeps gauge parameters in covariant gauge. It has been checked numerically that the amplitude is independent of gauge parameters at several phase-space points.
* Ward Identity check; when the polarization vectors of the external photons are replaced by their four-momentum, the amplitude must be zero due to Ward-Identity. We have checked it numerically at several phase-space points.
* Soft photon check; the cross-sections with soft-photon emission can be easily calculated by non-radiation cross-section and the soft-photon emission function. We have calculated the soft-photon emission cross-section by two methods;
1. Using $`4\mathrm{f}+\gamma `$ matrix elements with cuts, $`10^4\mathrm{GeV}<E_\gamma <10^2\mathrm{GeV}`$, no angular cut on the photon, $`|\mathrm{cos}\theta _\mu |<0.985,E_\mu >5\mathrm{GeV},M(ud)>10\mathrm{GeV}`$, giving $`\sigma =0.5105\pm 0.0002`$pb;
2. Using $`4\mathrm{f}+\gamma `$ matrix elements with soft-photon function with cuts, $`10^4\mathrm{GeV}<E_\gamma <10^2\mathrm{GeV}`$, no angular cut on the photon, $`|\mathrm{cos}\theta _\mu |<0.985,E_\mu >5\mathrm{GeV},M(ud)>10\mathrm{GeV}`$ giving $`\sigma =0.5109\pm 0.0005`$pb.
The two methods, therefore, give consistent results. We have used exact matrix elements for the calculations of $`4\mathrm{f}+\gamma `$. For $`W`$-pair processes, we simply used fixed width for the gauge-boson propagator in the unitary gauge. For single-$`W`$ processes we used a special gauge for the $`t`$-channel photon, which shows very small effects from the gauge violation due to the gauge-boson width.
Distributions from GRACE are shown in Fig. 29-34.
### 5.2 Comparisons for $`4\mathrm{f}+\gamma `$
A first set of comparisons between the predictions of several independent codes, namely WRAP, CompHEP, GRACE and RacoonWW was performed at the beginning of the workshop.
This comparison covers integrated cross-sections and various differential distributions, essentially for a CC10 final state. Discrepancies observed at that stage are mainly to be ascribed to non-tuned comparisons. In fact, a detailed tuned comparison between WRAP, RacoonWW and PHEGAS/HELAC presently shows a beautiful agreement for several distributions and final states.
Input parameters and cuts used to carry out this tuned comparison correspond to those of the $`4\mathrm{f}`$ proposal (in the approximation of massless fermions). In particular, the photon cuts are: $`E_\gamma ^{\mathrm{min}}=1`$ GeV, $`|\mathrm{cos}\theta _\gamma |<`$ 0.985, at $`\sqrt{s}=200`$ GeV. In the PHEGASRacoonWWWRAP comparison, the following final states have been considered:
* $`\mu \overline{\nu }_\mu u\overline{d}\gamma `$
* $`e^{}\overline{\nu }_eu\overline{d}\gamma `$
* $`\mu \overline{\nu }_\mu \tau ^+\nu _\tau \gamma `$
* $`e^{}\overline{\nu }_e\tau ^+\nu _\tau \gamma `$
* $`s\overline{c}u\overline{d}\gamma `$
The observables studied in the tuned comparison are:
* integrated cross-sections;
* $`E_\gamma `$ distribution, $`d\sigma /dE_\gamma `$ \[fb/GeV\];
* distribution in the cosine of the photon angle $`\theta _\gamma `$, $`d\sigma /d\mathrm{cos}\theta _\gamma `$ \[fb\];
* distribution in the opening angle $`\theta _{f\gamma }`$ between the photon and the nearest charged final-state fermion, $`d\sigma /d\theta _{\gamma f}`$ \[fb\];
* distributions in the bare invariant masses of the $`W^+`$ and $`W^{}`$ bosons, $`M_+=M_{u\overline{d}},M_{\tau ^+\nu _\tau }`$, $`d\sigma /dM_+`$ \[fb/GeV\]; $`M_{}=M_{s\overline{c}},M_{\mu ^{}\overline{\nu }_\mu }`$, $`M_{e^{}\overline{\nu }_e}`$, $`d\sigma /dM_{}`$ \[fb/GeV\].
All the observables are calculated for $`\sqrt{s}=200`$GeV in the fixed width scheme. The squared matrix element is calculated in the $`G_F`$ scheme and subsequently multiplied by $`\alpha (0)/\alpha _{G_F}`$, to take exactly into account of the scale of the real photon.
The applied cuts are:
* common to all processes: $`E_\gamma >1`$GeV, $`|\mathrm{cos}(\theta (\gamma ,\mathrm{beam})|<0.985`$, $`\theta (\gamma ,\mathrm{f})>5^{}`$, f = charged fermion.
* for $`ud\mu \nu _\mu \gamma `$ and $`ude\nu _e\gamma `$: $`M(ud)>10`$GeV, $`|\mathrm{cos}\theta (\mathrm{l},\mathrm{beam})|<0.985`$ $`E_\mathrm{l}>5`$GeV, where l is a charged lepton,
* for $`\tau \nu _\tau \mu \nu _\mu \gamma `$ and $`\tau \nu _\tau e\nu _e\gamma `$: $`|\mathrm{cos}\theta (\mathrm{l},\mathrm{beam})|<0.985`$, $`E_\mathrm{l}>5`$GeV, $`M(l^+l^{})>10`$GeV,
* for $`udcs\gamma `$: at least two pairs with $`M(q_iq_j)>10`$GeV.
The generators have produced a huge collection of results and only a small sample will be shown here.
The total cross-sections are reported in Table 13 where the differences between the predictions of WRAP, RacoonWW and PHEGAS/HELAC are around $`0.1\%`$, signalling perfect technical agreement.
In the following we will show few example of predictions. By comparing the three different codes with a tuned comparison we get a rough estimate of the associated technical uncertainty also for distributions. Besides the distributions compared in plots we also be present ratio-plots, as the distributions themselves are too close to show a difference between programs in the actual scale.
First we consider the angular distribution, i.e. the $`\mathrm{cos}\theta _\gamma `$ distribution in the range $`[1,1]`$ for various final states, as shown in Figures 35 to 39, where we also plot the ratios bewteen the predictions.
Similarly, the $`E_\gamma `$ distributions and ratios in the range GeV are shown for various processes in Figures 40 to 43. Note that virtual corrections are not included, therefore, the photon spectrum starts at some lower boundary of $`1`$GeV. Deviations are of the order of $`1\%`$ for soft photons and tend to deteriorate for harder ones. Statistically the deviations are compatible with zero. Note that for very hard photons the cross section and therefore the accuracy of the numerical integration of the programs becomes poorer.
In Fig. 44 we show the fermion-photon opening angle $`\theta (\gamma ,\mathrm{f})`$ (where f is a charged fermion) distribution. In the same figure we show the percentage deviation between the three predictions. The most interesting region occurs for small angles, i.e. towards the collinear region, where a reasonable agreement is registered, of the order of a percent. For the used statistics the deviations are not yet significant. The agreement deteriorates for a larger separation between the photon and the charged particles. However, in this region the cross-section is an order of magnitude lower. Note that the peculiar behavior of the distribution towards $`0^{}`$ is only due to the fact that the third bin is between $`3.6^{}`$ and $`5.4^{}`$ with a cut at $`5^{}`$.
Finally, we compare the distributions in the bare invariant masses. First the $`W^{}`$ one as predicted by WRAP and RacoonWW. Results for all considered channels and for the $`W^{}`$-distribution are shown in Fig. 45(left). Note that the curves for the two purely leptonic channels and the two semi-leptonic final states are almost identical. In Fig. 45 (right) we also present the percentage deviations for the process $`u\overline{d}s\overline{c}\gamma `$. In Fig. 46 we give the corresponding $`W^+`$ invariant-mass distribution including results from the $`3`$ programs. In Fig. 47 we show the ratio between the $`W^+`$ and the $`W^{}`$ invariant-mass distributions from WRAP and RacoonWW respectively.
### 5.3 Estimate of theoretical uncertainty
No global statement can be given, at the moment, on this issue. The following programs have agreed to make individual statements:
#### RacoonWW
Since the program has only tree-level precision for $`e^+e^{}4\mathrm{f}+\gamma `$, a reliable estimate for the theoretical uncertainty cannot be given with the present version. This could be done if leading corrections such as ISR were included, which is planned in future extensions of the program.
#### WRAP
WRAP has tried to estimate the theoretical uncertainty in $`4\mathrm{f}+\gamma `$ processes coming from variations in the renormalization scheme. The selected process is $`e^+e^{}u\overline{d}\mu ^{}\overline{\nu }_\mu \gamma `$ with the cuts used in the tuned comparisons. The following two schemes have been adopted:
$`I)s__W^2`$ $`=`$ $`1{\displaystyle \frac{M__W^2}{M__Z^2}},\alpha ={\displaystyle \frac{4\sqrt{2}G_FM__W^2s__W^2}{4\pi }},g^2=4\pi {\displaystyle \frac{\alpha }{s__W^2}},`$
$`II)s__W^2`$ $`=`$ $`{\displaystyle \frac{\pi \alpha (2M__W)}{\sqrt{2}G_FM__W^2}},g^2=4\sqrt{2}G_FM__W^2,\text{with}\alpha (2M__W)=128.07.`$ (63)
The cross section is always rescaled by the factor $`\alpha (0)/\alpha `$ in order to take into account of the scale $`\alpha (0)`$ for the emitted real photon. Here, $`\alpha `$ is the value computed in the corresponding renormalization scheme. The results are shown in Tab.(14).
Note that the overall theoretical uncertainty for $`4\mathrm{f}+\gamma `$ production cannot be below the level of $`1÷2\%`$. In this respect the numbers given in Tab.(14) are only a partial indication of possible sources of uncertainty. As shown by the previous analysis, ISR needs to be taken into account in programs for a realistic analysis of $`4\mathrm{f}+\gamma `$ final states. Furthermore, in order to avoid double-counting between pre-emission and matrix-element radiation, the implementation of QED corrections in computational tools for $`4\mathrm{f}+\gamma `$ processes should rely upon methods, such parton shower, YFS or $`p_t`$-dependent structure functions, able to keep under control photon $`p_t`$ effects. Effects due to finite fermion masses can become important at some percent level for small photon-charged fermion separation cuts.
In order to better understand the uncertainty associated to the implementation of collinear ISR in $`4\mathrm{f}+\gamma `$ processes, a comparison between the effects of ISR via collinear SF and $`p_t`$-dependent SF, respectively, is shown in Fig. 26 for the cross section of the the CC10 final state $`\mu ^{}\overline{\nu }_\mu u\overline{d}\gamma `$, as a function of the minimum energy of the observed photon, at $`\sqrt{s}=192`$ GeV. As can be seen, the two prescriptions for ISR can differ at $`5\%`$ level for $`E_\gamma ^{\mathrm{min}}`$ close to $`12\mathrm{GeV}`$, while the difference becomes smaller and smaller as $`E_\gamma ^{\mathrm{min}}`$ increases. In general, the difference between collinear and $`p_t`$-dependent SF is stronger near the soft and collinear regions, as a priori expected, and it gives an estimate of the size of the double-counting effect at the level of ISR.
#### NEXTCALIBUR
To check the sensitivity of various distributions to the chosen form of the Structure Functions, the processes $`e^+e^{}\mu ^{}\mu ^+u\overline{u}(\gamma )`$ and $`e^+e^{}e^{}\overline{\nu }_eu\overline{d}(\gamma )`$ have been considered with a slightly different implementation of the sub-leading terms, without observing any significant deviation, at the per mille level, with respect to the previous results.
### 5.4 Summary and conclusions
While the technical precision in $`e^+e^{}4\mathrm{f}+\gamma `$ does not represent a problem anymore for all those programs that implement an (exact) matrix element, very little effort has been devoted in analyzing the overall theoretical uncertainty. Some of the programs also include the large effect of initial state radiation at the leading logarithmic level. When this is done, the bulk of large radiative corrections is included. Since however in general non-logarithmic $`𝒪\left(\alpha \right)`$ corrections are not known, the theoretical accuracy is at the level of $`2.5\%`$ on integrated cross-sections and on inclusive distributions.
## 6 Single-$`W`$
Another interesting process at LEP 2 is the so-called single-$`W`$ production, $`e^+e^{}We\nu `$ which can be seen as a part of the CC20 process, $`e^+e^{}\overline{q}q(\mu \nu _\mu ,\tau \nu _\tau )e\nu _e`$, or as a part of the Mix56 process, $`e^+e^{}e^+e^{}\nu _e\overline{\nu }_e`$. For a more detailed theoretical review we refer to and to . All processes in the CC20/Mix56 families are usually considered in two regimes, $`|\mathrm{cos}\theta (e^{})|c`$ or SA and $`|\mathrm{cos}\theta (e^{})|c`$ or LA. In the list of observables, the single $`W`$ production is defined by those events that satisfy $`|\mathrm{cos}\theta (e^{})|0.997`$ and therefore is a SA.
The LA cross-section has been computed by many authors and references can be found in . It represents a contribution to the $`e^+e^{}W^+W^{}`$ total cross-section. From a theoretical point of view the evaluation of a LA cross-section is free of ambiguity, even in the approximation of massless fermions, as long as a gauge-preserving scheme is applied and $`\theta (e^{})`$ is not too small.
For SA instead, one cannot employ the massless approximation anymore. In other words, in addition to double-resonant $`W`$-pair production with one $`W`$ decaying into $`e\nu _e`$, there are $`t`$-channel diagrams that give a sizeable contribution for small values of the polar scattering angle of the $`t`$-channel electron. Single-$`W`$ processes are sensitive to the breaking of $`U(1)`$ gauge invariance in the collinear limit, as described in Ref. (see also ). The correct way of handling them is based on the so-called Fermion-Loop (FL)scheme , the gauge-invariant treatment of the finite-width effects of $`W`$ and $`Z`$ bosons in LEP 2 processes. However, till very recently, the Fermion-Loop scheme was available only for the LA-regime. For $`e^+e^{}e^{}\overline{\nu }_ef_1\overline{f}_2`$, the $`U(1)`$ gauge invariance becomes essential in the region of phase space where the angle between the incoming and outgoing electrons is small, see the work of and also an alternative formulations in . In this limit the superficial $`1/Q^4`$ divergence of the propagator structure is reduced to $`1/Q^2`$ by $`U(1)`$ gauge invariance. In the presence of light fermion masses this gives raise to the familiar $`\mathrm{ln}(m_e^2/s)`$ large logarithms. Furthermore, keeping a finite electron mass through the calculation is not enough. One of the main results of was to show that there are remaining subtleties in CC20, associated with the zero mass limit for the remaining fermions.
In a generalization of the Fermion-Loop scheme (hereafter EFL) is introduced to account for external, non-conserved, currents. Another extension has been given in for the imaginary parts of Fermion-Loop contributions, which represents the minimal set for preserving gauge invariance.
The most recent numerical results produced for single-$`W`$ production are from the following codes : CompHEP, GRC4F, NEXTCALIBUR, SWAP, WPHACT and WTO.
In view of a requested, inclusive cross-section, accuracy of $`2\%`$ we must include radiative corrections to the best of our knowledge, at least the bulk of any large effect. As we know, the correct scale of the couplings and their differentiation between $`s`$\- and $`t`$-channel is connected to the real part of the corrections, so that the imaginary FL is not enough, we need a complete FL for single-$`W`$, or EFL. Having all the parts, the tree-level couplings are replaced by running couplings at the appropriate momenta and the massive gauge-boson propagators are modified accordingly. The vertex coefficients, entering through the Yang–Mills vertex, contain the lowest order couplings as well as the one-loop fermionic vertex corrections.
Each calculation aimed to provide some estimate for single-$`W`$ production is, at least nominally, a tree level calculation. Among other things it will require the choice of some Input Parameter Set (IPS) and of certain relations among the parameters. Thus, different choices of the basic relations among the input parameters can lead to different results with deviations which, in some case, can be sizeable and should be included in the theoretical uncertainty. Here, more work is needed.
For instance, a possible choice is to fix the coupling constant $`g`$ as
$$g^2=\frac{4\pi \alpha }{s__W^2},s__W^2=\frac{\pi \alpha }{\sqrt{2}G_FM__W^2},$$
(64)
where $`G_F`$ is the Fermi coupling constant. Another possibility would be to use
$$g^2=4\sqrt{2}G_FM__W^2,$$
(65)
but, in both cases, we miss the correct running of the coupling. Ad hoc solutions should be avoided, and the running of the parameters must always follow from a fully consistent scheme.
Another important issue in dealing with single-$`W`$ production is connected with the inclusion of QED radiation. It is well known that universal, $`s`$-channel structure functions are not adequate enough to include the radiation since they generate an excess of ISR bremsstrahlung. In $`t`$-channel dominated processes the interference between incoming fermions becomes very small while the destructive interference between initial and final states becomes strong.
It is quite a known fact that, among the electroweak corrections, QED radiation gives the largest contribution and the needed precision requires a re-summation of the large logarithms. For annihilation processes, $`e^+e^{}\overline{f}f`$, initial state radiation is a definable, gauge-invariant concept and we have general tools to deal with it; the structure function approach and also the parton-shower method. However, when we try to apply the algorithm to four-fermion processes that include non-annihilation channels we face a problem: it is still possible to include the large universal logarithms by making use of the standard tools but an appropriate choice of scale is mandatory. Such is the case in single-$`W`$. The problem of the correct scale to be used in QED corrections has been tackled by two groups, GRACE and SWAP and additional results will be shown in Sub-Sects. 6.2.16.2 and in Sub-Sect. 6.2.
### 6.1 Signal definition in single-$`W`$
The experimental requirements on single-$`W`$ are:
* CC20 – Mix56 calculations with some detector acceptance that are used for a) triple gauge coupling determination, b) standard model background to searches;
* the LEP EWWG cross-section definition that is used to combine the cross-section measurements from the four LEP experiments.
During the last WW99 Crete Workshop a proposal has been made to reach a common signal definition for the LEP EWWG cross-section . The persons who participated in the WW99 workshop agreed on some setup to define the single-$`W`$ production and now this has been formalized in one of the LEP EWWG meetings; there, it was decided to have a combination of the single-$`W`$ cross-section using the signal definitions of Tab.(15) for $`e^+e^{}e^{}\overline{\nu }_ef^{}\overline{f}`$: The set of $`t`$-channel diagrams, all for CC20, are shown in Fig. 48. The signal definition uses $`10`$ diagrams for CC20, $`9`$ for CC18 and $`37`$ for Mix56.
Note that charge-conjugate state should be taken into account and that an asymmetric cut has been introduced for $`ee\nu \nu `$; the latter is due to the fact that the process itself is CP-even when no cut is applied, but an ambiguity remains if one starts to discuss single-$`W`$ with $`e^{}`$ in the forward direction. Then we should multiply this process by a factor $`2`$ as well. The goal of this common definition is to be able to combine the different $`e\nu \overline{q}q,e\nu \mu \nu ,e\nu \tau \nu ,e\nu e\nu `$ measurements from different experiments so that the new theoretical calculations can be checked with data at a level better than $`10\%`$.
Signal definition has a longstanding tradition in LEP physics, the most celebrated being the $`t`$-channel subtraction in Bhabha and the most recent being the CC03 cross-section. Here we have a different situation. First of all, nobody has radiative corrections for single-$`W`$ production, hence the usual argument of the availability of a sophisticated semi-analytical calculation for the signal does not apply. We could avoid a definition of the signal in terms of diagrams and have recourse to a definition in terms of cuts since, in a very narrow cone around the beam axis, the single-$`W`$ family is fully dominated by the $`t`$-channel photons.
#### 6.1.1 A study of single-$`W`$ signal definition with CompHEP
#### Authors
#### Single-$`W`$ signal definition in the reaction $`e^+e^{}e^+e^{}\nu _e\overline{\nu }_e`$
It is well-known for a long time how the single $`W`$ signal can be separated with the help of kinematical cuts . The typical set of cuts used by ALEPH, DELPHI and L3 collaborations for the leptonic four fermion states $`e^{}\overline{\nu }_el^+\nu _l`$ separates the configurations with very forward $`e^{}`$ and a rather energetic $`l^+`$ produced at a sufficiently large angle with the beam. For instance, the L3 cuts to be used in the following calculations are $`|\mathrm{cos}\theta _e^{}|0.997`$, $`E_l15\mathrm{GeV}`$ and $`|\mathrm{cos}\theta _{l^+}|0.997`$. In the case of the semi-leptonic states $`e^{}\overline{\nu }_eq\overline{q}^{^{}}`$ an additional cut $`M(q\overline{q}^{^{}})`$ 45 GeV have been applied by OPAL. In so far as different collaborations are using not exactly the same cuts (defined by the optimal detector acceptance), the definition of the $`W`$ signal in terms of angular cuts is not universal and some standardization procedure is needed. In the recent proposal by LEP experiments the OPAL collaboration considered the possibility to introduce the definition of the $`W`$ signal in terms of diagrams. Angular cuts on the forward electron and the corresponding anti-lepton are not imposed, so the single $`W`$ cross-section depends only on the $`E_l`$ energy cut and is defined by the gauge invariant subset of the $`t`$-channel single resonant diagrams. The universality of such definition is satisfactory if the interferences between the gauge invariant subsets of diagrams in the channels $`e^{}\overline{\nu }_el^+\nu _l`$ and $`e^{}\overline{\nu }_eq\overline{q}^{^{}}`$ are always negligible. Then indeed the single $`W`$ cross-section in terms of diagrams could be meaningful.
We performed a detailed calculation of the contributions from various diagram sets of the Mix56 channel $`e^+e^{}e^+e^{}\nu _e\overline{\nu }_e`$ (see Appendix for Fig. 63-Fig. 70 referred to in the following). Using the general approach to the amplitude decomposition into gauge invariant classes , we found ten gauge invariant subsets of diagrams (see Fig. 63-Fig. 64). In Tab.(16) $`18W`$ denotes two gauge invariant subsets of $`9`$ diagrams with single $`W`$ (see Fig. 63), $`8Z`$ denotes two gauge invariant subsets of $`4`$ diagrams with single $`Z`$ (see Fig. 64), $`9W^+W^{}`$ stands for the double-resonant subset (Fig. 65) and so on. Main contribution to the final configurations with forward electron come from the single $`W`$ and the single $`Z`$ production, while various $`\gamma ,Ze^+e^{}`$ conversion corrections (Fig. 68-Fig. 70) to the $`e^+e^{}e^+e^{},\nu _e\overline{\nu }_e`$ are negligible. For the case of angular cuts on the forward electron the interference between the single $`W`$ and single $`Z`$ subsets $`18W`$ and $`8Z`$ is negative and equal to several fb. However, if the angular cuts are removed, the destructive interference modulo increases rather considerably (Table 16). This is not an unexpected fact since both single-$`W`$ and single-$`Z`$ (NC processes with one lost electron) subsets have a similar $`t`$-channel pole structure. Other interferences are also not negligible. So in the case of $`e^+e^{}\nu _e\overline{\nu }_e`$ channel the diagram-based definition of single $`W`$ signal is not completely satisfactory.
### 6.2 Description of the programs, results and comparisons
#### WTO and EFL
#### Author
#### The Fermion-Loop scheme (EFL)
The EFL scheme for non-conserved currents is described in Ref. and briefly discussed in Sect. 3.8.1. It consists of the re-summation of the fermionic one-loop corrections to the vector-vector, vector-scalar and scalar-scalar propagators and of the inclusion of all remaining fermionic one-loop corrections, in particular those to the Yang–Mills vertices.
In the original formulation, the Fermion-Loop scheme requires that vector bosons couple to conserved currents, i.e. , that the masses of all external fermions can be neglected. There are several examples where fermion masses must be kept to obtain a reliable prediction. As already stated, this is the case for the single-$`W`$ production mechanism, where the outgoing electron is collinear, within a small cone, with the incoming electron. Therefore, $`m_e`$ cannot be neglected. Furthermore, among the $`20`$ Feynman diagrams that contribute (for $`e\overline{\nu }_eu\overline{d}`$ final states, up to $`56`$ for $`e^+e^{}\nu _e\overline{\nu }_e`$) there are multi-peripheral ones that require a non-vanishing mass also for the other outgoing fermions.
As well known in the literature, the Fixed-Width scheme behaves properly in the collinear and high-energy regions of phase space, to the contrary of the Running-Width scheme, but it completely misses the running of the couplings, an effect that is expected to be above the requested precision tag of $`2\%`$. To be specific the name of Fixed-Width scheme is reserved for the following: the cross-section is computed using the tree-level amplitude. The massive gauge-boson propagators are given by $`1/(p^2m^2+i\mathrm{\Gamma }m)`$. This gives an unphysical width in $`t`$-channel, but retains $`U(1)`$ gauge invariance in the CC20 process.
The correct way of handling this problem is to apply the EFL-scheme and, by considering the impact of the EFL-scheme on the relevant observables, one is able to judge on the goodness of naive rescaling procedures or of any incomplete FL-scheme. One of the problem with the latter is that vertices, although chosen to respect gauge-invariance, are not uniquely defined. Furthermore, couplings other than $`\alpha _{\mathrm{QED}}`$ usually do not evolve with the scale and complex poles, the truly gauge-invariant quantities, are never introduced or explicitly computed. Finally, programs than cannot split diagrams and apply an overall rescaling, both in $`s`$\- and $`t`$-channel, mistreat single-$`W`$ and/or violates $`SU(2)`$ invariance.
#### Numerical results and recommendations.
Numerical results for EFL have been shown in Ref. . Here, we limit the presentation to some useful recommendations:
* the bulk of the effect is in the running of the e.m. coupling constant;
* one can compute the single-$`W`$ cross-sections for a fixed mass of the top quark, such as $`m_t=173.8`$GeV, without finding any significative difference with the case where $`m_t`$ is fixed by a consistency relation. We are using complex-mass renormalization but we only include fermionic corrections. Therefore, we can start with the Fermi coupling constant but also with $`M__W`$ as an input parameter. Equating the corresponding renormalization conditions yields a relation between $`M__Z`$, $`G_F`$, $`\mathrm{Re}\{\alpha (M__Z^2)^1\}`$, $`M__W`$, and $`m_t`$, see . This relation can be solved iteratively for $`m_t`$. For the following input parameter set, $`M__W=80.350\mathrm{GeV}`$, $`M__Z=91.1867\mathrm{GeV}`$ and $`G_F=1.16639\times \mathrm{\hspace{0.17em}10}^5\mathrm{GeV}^2`$, we obtain the following solution:
$$\mu _W=\sqrt{\mathrm{Re}\left(p__W\right)}=80.324\mathrm{GeV},\gamma _W=\frac{\mathrm{Im}\left(p__W\right)}{\mu _W}=2.0581\mathrm{GeV},m_t=148.62\mathrm{GeV},$$
(69)
with $`26\mathrm{MeV}`$ difference between $`M__W`$ and $`\mu _W`$. See Sect. 6.3.1 for the inclusion of QCD effects. This type of effect should be included in any incomplete FL-scheme;
* the main accent in the EFL-scheme is on putting the correct scale in the running of $`\alpha _{\mathrm{QED}}`$. The latter is particularly important for the $`t`$-channel diagrams, dominated by a scale $`q^20`$ and not $`q^2M__W^2`$. However, a consistent implementation of radiative corrections does more than evolving $`\alpha _{\mathrm{QED}}`$ to the correct scale, other couplings are also running, propagators are modified and vertices are included;
* the effective FW-scheme describes considerably well the hadronic final state with a cut of $`M(u\overline{d})>45`$GeV. However, the diminution induced by $`\alpha _{\mathrm{QED}}(q^2)`$ is too large for the leptonic final state. The latter is a clear sign that other effects are relevant and a naive rescaling does not suffice in reproducing a realistic approximation in all situations, at least not within the $`2\%`$ level of requested theoretical accuracy;
* Modifications induced by the fermionic loops are sensitive to the relative weight of single-resonant terms and of multi-peripheral peaks. Furthermore, the effect of radiative corrections inside the $`W`$-propagators ($`\rho `$-factors of Ref. ) is far from being negligible and tends to compensate the change due to the running of $`\alpha _{\mathrm{QED}}`$.
These recommendations are better illustrated by few examples. At $`\sqrt{s}=183`$GeV we consider the angular distribution, $`d\sigma /d\theta _e`$ for the $`u\overline{d}e^{}\overline{\nu }_e`$ final states. The results are shown in Fig. 49.
From Fig. 49 we see that the EFL prediction is lower than the FW one, from $`7.46\%`$ in the bin $`0^{}0.1^{}`$ to $`5.56\%`$ in the bin $`0.3^{}0.4^{}`$. Correspondingly, the first bin is $`6.78`$ higher than the second one, $`11.60(16.37)`$ than the third(fourth) one. This is not a surprise, since the first bin represents $`50\%`$ of the total single-$`W`$ cross-section.
Always in the same figure, we have reported the behavior of $`\left[\alpha (q^2)/\alpha _{G_F}1\right]^2`$ as a function of $`\theta _e`$ for three values of $`y`$, using the appropriate relation: $`q^2=q^2(\theta _e,y)`$, $`y`$ being the fraction of the electron energy carried by the photon. The behavior of EFL/FW-1, when we vary $`\theta _e`$, is very similar to the one given by the ratio of coupling constants, indicating that the bulk of the effect is in the running of the e.m. coupling constant.
For completeness we have reported the numerical results for the three energies in Tab.(17), where the first entry is Fixed-Width distribution and the second entry is EFL one. Only the first four bins are shown, owing to the fact that they are the most significant in the distribution. The third entry in Tab.(17) gives EFL/FW-1 in percent.
Next we consider $`e^+e^{}e\nu \mu \nu `$, with $`|\mathrm{cos}\theta _e|>0.997`$, $`E_\mu >15`$GeV, and $`|\mathrm{cos}\theta _\mu |<0.95`$. In Tab.(18) we report the comparison between the EFL distribution and the FW one for $`\sqrt{s}=183`$GeV. As before, only the most significant bins are shown ($`0.0^{}÷0.4^{}`$).
As for the hadronic case, the EFL prediction is considerably lower than the FW one, although the percentage difference between the two is approximately $`2.2\%÷2.4\%`$ smaller than in the previous case. Useful comparisons will be presented in the WPHACT description of this Section.
A final comment will be devoted to QED ISR. Very often one can find the statement that the choice of the appropriate scale in the structure functions is mandatory. This is a jargon for ‘implementing the correct exponentiation factor in multi-photon emission’. Note that the usual infrared exponent $`\alpha B`$ is represented by
$`\alpha B`$ $`=`$ $`{\displaystyle \frac{2\alpha }{\pi }}\left[{\displaystyle \frac{1+r^2}{1r^2}}\mathrm{ln}\left({\displaystyle \frac{1}{r}}\right)1\right]{\displaystyle \frac{2\alpha }{\pi }}\left(\mathrm{ln}{\displaystyle \frac{Q^2}{m^2}}1\right),\text{for}Q^2m^2,`$
$`{\displaystyle \frac{m^2}{Q^2}}`$ $`=`$ $`{\displaystyle \frac{r}{(1r)^2}}r,`$ (70)
where $`Q^2`$ is the Mandelstam invariant associated with the emitting pair. For $`|t|m_e^2`$ the photon radiation is governed by $`\mathrm{ln}(|t|/m_e^2)`$ rather than by $`\mathrm{ln}(s/m_e^2)`$. The difference is again a large log and explain the excess of radiation generated by $`s`$-channel SF. However, the whole expression for $`B`$ is known and not only its asymptotic behavior (the scale). Therefore, for vanishing scattering angles, the correct behavior should be read from Eq.(70). In this respect one should remember that $`|t_\gamma |_{\mathrm{min}}`$ in single-$`W`$ can be much lower than $`m_e^2`$, being $`m_e^2y^2/(1y)`$ where $`y=M^2(\nu _ef_1\overline{f}_2)/s`$.
#### Single-$`W`$ with WPHACT
#### Authors
A new version of WPHACT is now available. It includes all massive matrix elements in addition to the previous ones which accounted for $`b`$-quark masses only. As before, the matrix elements are computed with the method of Ref. , which has proved to be fast and reliable in particular for massive calculations. New mappings of the phase space have been added, in order to account in an efficient way for the peaking structure of contributions like single-$`W`$, single-$`Z`$ and $`\gamma \gamma `$ contributions. With the new version one has, therefore, the choice of using fully massive or massless calculations. The former are needed in various processes which diverge for massless fermions , while the latter are faster and give an excellent approximation for most cases. We start with the introduction of the IFL-scheme showing comparisons with alternative solutions, designed to deal with gauge-invariance issues. However, the most important part is contained in the second Subsection where the effective scaling induced by $`\alpha _{\mathrm{QED}}`$ is presented.
#### IFL-scheme
The Imaginary Part Fermion-Loop scheme has been generalized to the fully massive case of non-conserved weak currents in Ref. . The results obtained have been compared with other gauge restoring schemes used in single-$`W`$ processes computations. The following schemes have been considered in the analysis:
* Imaginary-part FL scheme(IFL): The imaginary part of the fermion-loop corrections, as computed in Ref. , are used. Fermion masses are neglected only in loops but not in the rest of the diagrams.
* Fixed width(FW): The $`\mathrm{W}`$-boson propagators show an unphysical width for $`p^2<0`$, but retains $`U(1)`$ gauge invariance in the CC20 process .
* Complex Mass(CM): All weak boson masses squared $`M__B^2,B=W,Z`$ are changed to $`M_B^2iM_B\mathrm{\Gamma }_B`$ ($`\mathrm{\Gamma }_B`$ is the on-shell $`B`$ width), including when they appear in the definition of the weak mixing angle. This scheme, which again gives an unphysical width in some cases, has however the advantage of preserving both $`U(1)`$ and $`SU(2)`$ Ward identities.
* Overall scheme(OA): The diagrams for $`e^{}e^+e^{}\overline{\nu }_eu\overline{d}`$ can be split into two sets that are separately gauge invariant under $`U(1)`$. In the actual implementation of the OA-scheme, $`t`$ channel diagrams are computed without any width and are then multiplied by $`(q^2M^2)/(q^2M^2+iM\mathrm{\Gamma })`$ where $`q`$, $`M`$ and $`\mathrm{\Gamma }`$ are the momentum, the mass and the width of the possibly-resonant $`\mathrm{W}`$-boson. This scheme retains $`U(1)`$ gauge invariance at the expenses of a mistreatment of the non-resonant terms.
In order to asses the relevance of current non-conservation, the imaginary part of the fermion-loop corrections have also been implemented with the assumption that all currents that couple to the fermion-loop are conserved. In this case the expressions of Ref. reduce to those computed in . Note that the masses of external fermions are nonetheless taken into account in the calculation of the matrix elements. This scheme violates $`U(1)`$ gauge-invariance by terms which are proportional to the fermion masses squared, as already noted in Ref. . However they are enhanced at high energy by large factors and can be numerically quite relevant. This scheme will be referred to as the imaginary-part FL scheme with conserved currents (hereafter IFLCC). All schemes described above have been implemented in the new version of WPHACT with the fully massive option.
In Tab.(19) the cross-sections for CC20 are given for the different gauge restoring schemes at LEP 2 and LC energies. From it, one can immediately deduce that the IFL, FW, CM and the OA schemes agree within $`2\sigma `$ in almost all cases. The IFLCC scheme agrees with the other ones at LEP 2 energies but already at 800 $`\mathrm{GeV}`$ it overestimates the total cross-section by about $`6\%`$. At $`1.5\mathrm{TeV}`$ the error is almost a factor of two. On the contrary, even in the presence of non–conserved currents, i.e. of massive external fermions, the FW CM and OA calculations give predictions which are in agreement, within a few per mil, with the IFL scheme. The agreement with the results of a self-consistent approach justifies, from a practical point of view, the ongoing use of the FW, CM and OA schemes.
The possible dependence of this agreement on the particular single-$`W`$ process considered has been examined and we compare in Tab.(20) the cross-sections obtained in the IFL and FW scheme at $`\sqrt{s}=200`$ GeV. In this case, as in the following ones, the standard cuts have been applied: the electron angle is limited in all processes by $`|\mathrm{cos}\theta _e^{}|>0.997`$, the other charged lepton by $`|\mathrm{cos}\theta _l|<0.95`$, its energy has to be $`E_l>15`$ GeV. These results confirm that, at LEP 2, there is no dependence of the cross-sections on the scheme. Distribution of several observables have also been studied with WPHACT in the IFL and FW schemes. In most variables like the electron angle and energy no difference has been found. However, the mass spectrum of the $`u\overline{d}`$ pair shows some scheme dependence, as reported in Fig. 50. The physical motivation for this difference can be traced to the fact that the IFL scheme uses, correctly, a running $`W`$-width. In fact, comparing IFL mass distribution with a FW calculation in which $`W`$ mass and width are properly shifted , the difference is reduced to a small overall factor, as expected, and should not be viewed as a theoretical uncertainty.
In any case, in view of possible discrepancies, the use of IFL has to be preferred among the schemes analyzed in this section.
#### Running of $`\alpha `$, comparisons with EFL.
The EFL scheme implemented in for the massive (non-conserved currents) case solves the gauge-invariance problems exactly, as IFL does, but in addition it computes the real part of Fermion-Loop radiative corrections. These terms are known to determine the running of the couplings involved in single-$`W`$ processes. One may argue, therefore, that considering the running of $`\alpha _{\mathrm{QED}}`$ at an appropriate physical scale might account for the most relevant part of EFL corrections. To test the correctness of this argument, a proper $`\alpha _{\mathrm{QED}}`$ evolution has been introduced as an option in WPHACT. For every set of final momenta, $`\alpha _{\mathrm{QED}}`$ is evaluated at the scale $`t`$, the virtuality of the photon emitted by the electron line, and used for two vertices in the $`t`$-channel contributions only.
The separate gauge invariance of $`s`$\- and $`t`$-channel diagrams makes it possible to use a different $`\alpha `$ for them: $`\alpha (t)`$ for t-channel and $`\alpha _{G_F}`$ for $`s`$-channel. Such a separation, which can be implemented in codes computing Feynman diagrams as WPHACT, should automatically account for the relative weight of $`s`$ and $`t`$ contributions for any set of cuts.
Computations performed with this choice will be referred to as IFL<sub>α</sub>. Several comparisons have been performed between the IFL and IFL<sub>α</sub> schemes and with the FW/EFL predictions by WTO.
The good agreement of the two codes as far as FW and IFL schemes are concerned is documented in Tabs.(2125) for the cross-sections, the electron angular distribution and the quark invariant mass distribution. However, this has to be considered as a technical agreement more than a physical one. Whether IFL<sub>α</sub> can satisfactory reproduce the EFL complete calculations seems to depend on the process considered. Note, in Tab.(21), the agreement between IFL<sub>α</sub> and EFL for the total cross-section of the process $`e^{}e^+e^{}\overline{\nu }_eu\overline{d}`$. Only at $`200`$GeV there is a disagreement of less than $`0.5\%`$ . Moreover, the angular distribution studied in Tab.(22), for the most relevant bins, never shows a higher discrepancy.
The variation of the cross-section of the process at hand with the invariant mass $`M(u\overline{d})`$ cut is reported in Fig. 51 from which one deduces that the IFL and the IFL<sub>α</sub> schemes practically coincide when the cut reaches the mass of the $`W`$-boson. In Tab.(23) one sees that, even varying the cuts, the difference between FL and IFL<sub>α</sub> is at most of the order of $`1\%`$.
The conclusion is, therefore, that at LEP 2 and for $`e^{}e^+e^{}\overline{\nu }_eu\overline{d}`$ the IFL<sub>α</sub> scheme is reliable at the percent level. The same does not apply to $`e^{}e^+e^{}\overline{\nu }_e\mu ^+\nu _\mu `$, as can be verified with the help of Tab.(24) and Tab.(25). From these one sees that the discrepancy is of the order of $`2\%`$ or worse. This confirms that varying the scale of $`\alpha _{\mathrm{QED}}`$, on an event by event basis, is not completely satisfactory. These numerical results point towards an estimate of about $`3\%`$ theoretical error for single-$`W`$ predictions via the IFL<sub>α</sub>-scheme. One can try to apply the running of $`\alpha _{\mathrm{QED}}`$ to only one vertex of the $`t`$-channel diagrams; the agreement obtained with this approximation (hereafter IFL<sub>α1</sub>) is much better for $`e^{}e^+e^{}\overline{\nu }_e\mu ^+\nu _\mu `$. Of course, it becomes worse for $`e^{}e^+e^{}\overline{\nu }_eu\overline{d}`$. At $`183,189`$ and $`200`$GeV the cross-sections for $`e^{}e^+e^{}\overline{\nu }_e\mu ^+\nu _\mu `$ are respectively $`25.65(1),28.80(2),34.86(2)`$ fb, to be compared with the EFL results of Tab.(24). The first bins of the angular distribution are also very close to EFL. No physical meaning has to be attributed to this fact: there is no theoretical reason for using running $`\alpha _{\mathrm{QED}}`$ just at one vertex. The agreement may be accidental and it is probably due to the fact that with the cuts used for $`e^{}e^+e^{}\overline{\nu }_e\mu ^+\nu _\mu `$ the contribution of multi-peripheral diagrams is suppressed.
Since the IFL<sub>α</sub> and IFL<sub>α1</sub> schemes are, in turn, in good agreement with complete EFL for different processes and cuts, the difference between their results will be used as an estimate of the theoretical error for $`e^{}e^+e^{}e^+\nu _e\overline{\nu }_e`$ and $`e^{}e^+e^{}e^+\nu _\mu \overline{\nu }_\mu `$, where EFL predictions are not available. The cross-sections for such processes are presented in Tab.(26) and Tab.(27). The angular distributions for the four processes that we have discussed so far are reported in bins of $`0.01`$ degrees in Fig. 52. Note that the relevant part of the cross-section is concentrated in the first three or four bins, also for $`e^{}e^+e^{}e^+\nu _e\overline{\nu }_e`$(Mix) and $`e^{}e^+e^{}e^+\nu _\mu \overline{\nu }_\mu `$(NC), as well as for the two CC processes. Finally the comparison between WPHACT and WTO has been extended to cover the LEP 2 signal definition for the hadronic decays of the $`W`$-boson, but the results will not be presented here.
#### Single-$`W`$ and SWAP.
#### Authors
#### Description of the Method
Contributions of the Pavia/ALPHA group to the subject of single-$`W`$ production are summarized. The exact matrix elements for single-$`W`$ production are computed by means of the ALPHA algorithm for the automatic evaluation of the Born scattering amplitudes. Fermion masses are exactly accounted for in the kinematics and dynamics. The contribution of anomalous trilinear gauge couplings is also taken into account. The anomalous gauge boson couplings $`\mathrm{\Delta }k_\gamma `$, $`\lambda _\gamma `$, $`\delta _Z`$, $`\mathrm{\Delta }k_Z`$ and $`\lambda _Z`$ are implemented according to the parameterization of refs. . The fixed-width scheme is adopted as gauge-restoring approach, as motivated in comparison with other gauge-invariance-preserving schemes in Ref. .
#### Radiative corrections
Leading-log (LL) QED radiative corrections are implemented via the Structure Function (SF) formalism in the collinear approximation . The $`Q^2`$-scale entering the SF $`D(x,Q^2)`$ is fixed by comparing the $`𝒪\left(\alpha \right)`$ expansion of the SF method with the analytic results obtained for the $`𝒪\left(\alpha \right)`$ double-log photonic corrections as given by soft-photon bremsstrahlung from the external legs, its virtual counterpart and hard-photon radiation collinear to the final-state particles. Notice that, since the goal is to determine the scale entering the SF, only the contribution of real photons is explicitly calculated, because the virtual corrections, in order to preserve the cancellations of infrared singularities, must share the same leading collinear structure of the real part itself. More details about the derivation in the present approach of the soft/collinear limit of the $`𝒪\left(\alpha \right)`$ correction can be found in .
For example, for the process $`e^+e^{}e^{}\overline{\nu }u\overline{d}`$, this comparison translates in the following two $`Q^2`$-scales: (two initial-state (IS) SF are assumed: $`Q_{}^2`$ refers to the SF attached to the incoming electron, while $`Q_+^2`$ to the SF attached to the incoming positron)
$$Q_{}^2=4E^2\frac{(1c_{})^2}{\delta ^2},Q_+^2=2^{\frac{14}{9}}E^2\frac{\left((1c_{\overline{d}})(1c_u)^2\right)^{\frac{2}{3}}}{\left((1c_{u\overline{d}})^2\delta ^5\right)^{\frac{2}{9}}}$$
(71)
where $`E`$ is the beam energy, $`c_{}`$ the cosine of the electron scattering angle, $`c_u`$ and $`c_{\overline{d}}`$ the cosine of the quark scattering angles with respect to the initial positron, $`c_{u\overline{d}}`$ the cosine of the relative angle between the quarks, $`\delta `$ the half-opening angle of the electromagnetic jet (calorimetric angular resolution). It is worth noticing that in the numerical implementation, whenever one of the two scales is less than a small cut-off ($`\mathrm{\Lambda }_{\mathrm{cut}\mathrm{off}}^2=4m_e^2`$, where $`m_e`$ is the electron mass), the radiation from the corresponding leg is switched off, in accordance with the expected power law behaviour.<sup>15</sup><sup>15</sup>15Although this behavior is exactly known and could be implemented, SWAP has evidence for a corresponding small effect. It was carefully tested that variations of the cut-off do not alter the numerical results.
Also a naive ansatz for the two scales, as motivated by an analysis of the single-$`W`$ process in terms of the Weizsäcker-Williams approximation, can be given as follows:
$$Q_{,\mathrm{naive}}^2=|q_\gamma ^{}^2|,Q_{+,\mathrm{naive}}^2=M__W^2$$
(72)
where $`q_\gamma ^{}^2`$ is the squared momentum transfer in the $`ee\gamma ^{}`$ vertex and $`M__W`$ is the mass of the $`W`$ boson.
The effect of vacuum polarization is also taken into account in the calculation, by including the contribution of leptons, heavy quarks and light quarks, the latter according to the standard parameterization of Ref. .
#### Computational tool and obtained results
The theoretical features sketched above have been implemented into a massive MonteCarlo (MC) program, named SWAP (Single W process with Alpha & Pavia). The multi-channel importance sampling technique is employed to perform the phase-space integration. The code supports realistic event selections and can be employed either as a cross-section calculator or as a true event generator. The main results obtained in the present study can be summarized as follows: we have performed a critical analysis of the energy scale for QED radiation (see Fig. 53); Next, we have evaluated the effect of a running of $`\alpha _{\mathrm{QED}}`$ (see Fig. 54); Finally we have performed a tuned comparisons with other codes.
Input parameters and cuts used to obtain the numerical results shown in the following are those of the $`4\mathrm{f}`$ proposal for the process $`e^+e^{}e^{}\overline{\nu }u\overline{d}`$ ($`|\mathrm{cos}\vartheta _e|>0.997,M_{u\overline{d}}>45`$ GeV). For Fig. 53 the value of $`\delta `$ parameter entering eq. (71) is $`\delta =5^{}`$, but it has been checked that the numerical results are very marginally affected by its actual value.
Scales & QED radiation. In Fig. 53 the numerical impact of different choices of the $`Q^2`$-scale on the cross-section of the single-$`W`$ process $`e^+e^{}e^{}\overline{\nu }u\overline{d}`$ is shown. The marker $``$ represents the Born cross-section, $``$ represents the correction given by $`Q_\pm ^2=s`$ scale for both IS SF(s), $`\mathrm{}`$ represents the correction given by $`Q_\pm ^2=|q_\gamma ^{}^2|`$ scale for both IS SF(s), $`\mathrm{}`$ the correction given by the scales of eq. (71), the correction given by the naive scales of eq. (72). It can be seen that neither the $`s`$ scale, as implemented in computational tools used for the analysis of the single-$`W`$ process, nor the $`|q_\gamma ^{}^2|`$ scale, as recently proposed in Ref. , are able to reproduce the effects due to the scales of eq. (71) and eq. (72). These two scales are in good agreement and both predict a lowering of the Born cross-section of about $`8\%`$, almost independent of the c.m.s. LEP 2 energy. Note the $`4\%`$ difference between ISR with $`s`$-scale and the new scale.
Running of $`\alpha _{\mathrm{QED}}`$. Because $`G_F`$, $`M__W`$ and $`M__Z`$ are the agreed input parameters in the $`4f`$ proposal, the value of the e.m. coupling constant $`\alpha `$ is fixed at tree-level to a high energy value as specified by the $`G_F`$-scheme. On the other hand, the single-$`W`$ process is a $`q_\gamma ^{}^20`$ dominated process and therefore the above high-energy evaluation of $`\alpha `$, $`\alpha _{G_F}`$, needs to be rescaled to its correct value at small momentum transfer. In order to take into account the effect of the running of $`\alpha _{\mathrm{QED}}`$ in a gauge invariant way, a re-weighting procedure can be adopted, by simply rescaling the differential cross section $`d\sigma /dt`$ ($`tq_\gamma ^{}^2`$) in the following way
$$\frac{d\sigma }{dt}\frac{\alpha ^2(0)}{\alpha _{G_F}^2}\frac{d\sigma }{dt},\frac{d\sigma }{dt}\frac{\alpha ^2(t)}{\alpha _{G_F}^2}\frac{d\sigma }{dt},$$
(73)
where $`\alpha (0),\alpha (t)`$ is the QED running coupling computed at virtuality $`q_\gamma ^{}^2`$ equal to $`0`$ and $`t`$, respectively.
Fig. 54 shows the effects of the above re-weighting procedure. The $`\mathrm{}`$ represent the relative difference between the integrated cross-section computed in terms of $`\alpha _{G_F}`$ and the cross-section computed in terms of $`\alpha (0)`$, while $`\mathrm{}`$ is the relative difference between the integrated cross-section computed in terms of $`\alpha _{G_F}`$ and the cross-section computed in terms of $`\alpha (t)`$. As can be seen, the rescaling from $`\alpha _{G_F}`$ to $`\alpha (t)`$ introduces a negative correction of about $`56\%`$ in the LEP 2 energy range. The difference between $`\mathrm{}`$ and $`\mathrm{}`$, which is about $`23\%`$, is a measure of the running of $`\alpha _{\mathrm{QED}}`$ from $`q_\gamma ^{}^2=0`$ to $`q_\gamma ^{}^2=t`$. A detailed numerical analysis of the effect of the running couplings in single-$`W`$ production has been very recently performed in Ref. , based on the theoretical results of the massive fermion-loop scheme of Ref. . The results for the running of $`\alpha _{\mathrm{QED}}`$, as shown in Fig. 54, are in agreement with those of Ref. , as far as the effect of $`\alpha _{\mathrm{QED}}`$ is concerned, which is the bulk of the EFL contribution, leaving residual differences at the level of $`12\%`$, depending on the considered channel and event selection, see also the discussion in the WPHACT part of this Section.
#### NEXTCALIBUR
#### Authors
This section describes the features of a new Monte Carlo program NEXTCALIBUR , which aims at keeping the advantages of EXCALIBUR , but tries to improve on its shortcomings. The advantages, which should be kept are the high speed of the program and the applicability to all possible 4-fermion final states. The shortcomings of EXCALIBUR, which are partly related to its assets, are the massless nature of its fermions, the inclusive treatment of ISR QED corrections (no $`p_t`$ from a photon in an event) and the neglect of any running of coupling constants.
#### The strategy of the code
To start with, it should be noted that unless stated otherwise complex gauge boson masses and a complex weak mixing angle are used to ensure gauge invariant matrix elements . This procedure has been shown to work well . The various wanted improvements will now be successively discussed.
#### Inclusion of fermion masses
Inside the program a massive matrix element is needed, for the calculation of which a recursive method is used.
This massive matrix element now exists in the whole phase space, since the singularities of the massless case are regularized. Nevertheless serious numerical cancellations take place in very specific situations. The most dramatic case is caused by the photonic multi-peripheral diagrams which blow up for forward scattering. When at the same time both electron and positron move in the forward direction, it becomes necessary to perform the calculation in quadruple precision. When only one is moving in the forward direction the usual double precision is sufficient. A version of the program using double precision in all possible situations is currently under study.
The phase space generation is an extension of the treatment in EXCALIBUR, i.e. a self-adjusting multi-channel approach, now including the multi-peripheral situation in an improved form.
With the above mentioned ingredients one indeed has an event generator for any massive four-fermion final state. In particular, for the potentially dangerous kinematical situations events can now be generated, like forward single W-production or $`\gamma \gamma `$ processes. Also all channels, where Higgs exchange can take place now indeed contain Higgs exchanges.
To demonstrate the ability of the program to cover all phase-space regions, without loosing efficiency, we show, in Tabs.(2829), the total cross-sections for the processes $`e^+e^{}e^+e^{}\mu ^+\mu ^{}`$ and $`e^+e^{}e^+e^{}e^+e^{}`$. Where available, we compare our predictions with the QED numbers published in Ref. .
NEXTCALIBUR contains all electroweak diagrams, and can therefore be used to compute the electroweak background to the above $`\gamma \gamma `$ processes. By looking at the last entry of the tables, the latter is found to be less than $`1\%`$ at LEP 2 energies, at least for totally inclusive quantities.
All numbers have been produced at the Born level, but ISR and running $`\alpha _{\mathrm{QED}}`$ can be included as described in the next sections.
#### Taking into account the correct scale
As mentioned above, the matrix element calculation can easily be modified. One option would be to take into account Fermion-Loop corrections, which becomes relevant when there are different scales in the matrix element, e.g. due to small $`t`$-channel scales. A possible solution is the Fermion-Loop approach of Refs. , where all fermion corrections are consistently included by introducing running couplings $`g(s)`$ and $`e(s)`$, together with the re-summed bosonic propagators.
In presence of the $`WW\gamma `$ vertex, the above ingredients are not sufficient to ensure gauge invariance, because loop mediated vertices have to be consistently included. On the contrary, when no $`WW\gamma `$ vertex is present, the neutral gauge boson vertices, induced by the Fermion-Loop contributions, are separately gauge invariant .
Instead of explicitly including the loop vertices, we follow a Modified Fermion-Loop approach. Namely, we neglect the separately gauge invariant neutral boson vertices, and include only the part of the $`WW\gamma `$ loop function necessary to renormalize the bare $`WW\gamma `$ vertex and to insure the $`U(1)`$ gauge invariance. Our procedure is as follows: besides running couplings, we use bosonic propagators
$`P_w^{\mu \nu }(s)=\left(sM_w^2(s)\right)^1\left(g_{\mu \nu }{\displaystyle \frac{p_\mu p_\nu }{M_w^2(s)}}\right)`$
$`P_z^{\mu \nu }(s)=\left(sM_z^2(s)\right)^1\left(g_{\mu \nu }{\displaystyle \frac{p_\mu p_\nu }{M_z^2(s)}}\right)`$
with running boson masses defined as
$`M_w^2(s)=\mu _w{\displaystyle \frac{g^2(s)}{g^2(\mu _w)}}g^2(s)[T_W(s)T_W(\mu _w)]`$
$`M_z^2(s)=\mu _z{\displaystyle \frac{g^2(s)}{c_\theta ^2(s)}}{\displaystyle \frac{c_\theta ^2(\mu _z)}{g^2(\mu _z)}}{\displaystyle \frac{g^2(s)}{c_\theta ^2(s)}}[T_Z(s)T_Z(\mu _z)].`$
$`T_{W,Z}(s)`$ are contributions due to the top quark, $`\mu _{w,z}`$ the complex poles of the propagators (one can take, for instance, $`\mu _{w,z}=M_{w,z}^2i\mathrm{\Gamma }_{w,z}M_{w,z}`$) and
$$s_\theta ^2(s)=\frac{e^2(s)}{g^2(s)},c_\theta ^2(s)=1s_\theta ^2(s).$$
The leading contributions are in the real part of the running couplings therefore we take only the real part of them. This also means that one can replace, in the above formulae, $`g^2(\mu _{w,z})g^2(M_{w,z}^2)`$, $`c_\theta ^2(\mu _z)c_\theta ^2(M__Z^2)`$ and also $`T_{W,Z}(\mu _{w,z})T_{W,Z}(M_{w,z}^2)`$.
When the $`WW\gamma `$ coupling is present, we introduce, in addition, the following effective three gauge boson vertex
with $`s=p^2,s^+=p_+^2,s^{}=p_{}^2`$ and
$`V_{\mu \nu \rho }`$ $`=`$ $`g_{\mu \nu }(pp_+)_\rho +g_{\nu \rho }(p_+p_{})_\mu (1+\delta _V)+g_{\rho \mu }(p_{}p)_\nu `$
$`+`$ $`{\displaystyle \frac{(p_+p_{})_\mu }{s^{}s^+}}\left[\left({\displaystyle \frac{g(s^{})}{g(s^+)}}1\right)p_{+\nu }p_{+\rho }\left({\displaystyle \frac{g(s^+)}{g(s^{})}}1\right)p_\nu p_\rho \right]`$
$`\delta _V`$ $`=`$ $`{\displaystyle \frac{1}{g(s^+)g(s^{})(s^{}s^+)}}[g^2(s^+)g^2(s^{})[T_W(s^{})T_W(s^+)]`$ (74)
$`+`$ $`[g(s^+)g(s^{})][s^{}g(s^+)+s^+g(s^{})]].`$
It is the easy to see that, with the above choice for $`V_{\mu \nu \rho }`$, the $`U(1)`$ gauge invariance - namely current conservation - is preserved, even in presence of complex masses and running couplings, also with massive final state fermions.
By looking at Eq.(74), one can notice at least two effective ways to preserve $`U(1)`$. One can either compute $`g(s)`$ at a fixed scale (for example always with $`s=M__W^2`$), while keeping only the running of $`e(s)`$, or let all the couplings run at the proper scale.<sup>16</sup><sup>16</sup>16Note, however, that in the complete formulation of the EFL-scheme there is no ambiguity and all scales are automatically fixed.
With the first choice the modification of the three gauge boson vertex is kept minimal (but the leading running effects included). With the second choice everything runs at the proper scale, but a heavier modification of the Feynman rules is required. At this point one should not forget that our approach is an effective one, the goodness of which can be judged only by comparing with the exact calculation of Ref. . We found that the second choice gives a better agreement for leptonic single-$`W`$ final states, while the first one is closer to the exact result in the hadronic case, which is phenomenologically more relevant. Therefore, we adopted this first option as our default implementation in NEXTCALIBUR. The results of the EFL-scheme are then reproduced at $`2\%`$ accuracy for both leptonic and hadronic single-$`W`$ final states.
We want to stress once more that the outlined solution is flexible enough to deal with any four-fermion final state, whenever small scales dominate. For example, once the given formulae are implemented in the Monte Carlo, the correct running of $`\alpha _{\mathrm{QED}}`$ is taken into account also for $`s`$-channel processes as $`Z\gamma ^{}`$ production.
Also naive QCD corrections can be easily included, without breaking $`U(1)`$ gauge invariance, by the usual recipe of rescaling the total $`W`$-width and the cross-section.
In fact, in our approach, $`\mathrm{\Gamma }__W`$ can be generic, and the above procedure respects current conservation, provided the same $`W`$-width is used everywhere<sup>17</sup><sup>17</sup>17Note, however, that in a complete EFL-scheme the relevant objects are the complex poles and QCD corrections should be computed accordingly, see Sect. 6.3.1.
#### Improving the treatment of the QED radiation
Once the matrix element calculation is fixed one can add externally the QED leading logarithmic effects in the Structure Function method . Such a strategy is implemented in most of the programs used for the analysis of the LEP 2 data and accurately reproduces the inclusive four-fermion cross-sections, at least for $`s`$-channel dominated processes. In principle both initial and final state radiation (ISR and FSR) can be treated in this way, as it has been explicitly done originally for Bhabha scattering . Here only the implementation of ISR in NEXTCALIBUR is discussed. There are two issues to be discussed. One is the choice of scale $`q^2`$ in the leading logarithm $`L=\mathrm{ln}(q^2/m_e^2)`$. Another is the unfolding of this leading logarithm in terms of an emitted photon. For the latter issue a particular form of $`p_t`$ dependent Structure Functions is implemented. These are derived, at the first leading logarithmic order, for small values of $`p_t`$. In practice, we replace the quantity
$$\mathrm{ln}(\frac{q^2}{m_e^2})\mathrm{by}\frac{1}{1c_i+2\frac{m_e^2}{q^2}}$$
in the strictly collinear Structure Function for the $`i^{th}`$ incoming particle, by explicitly generating $`c_1`$ and $`c_2`$, the cosines (in the laboratory frame) of the emitted photons with respect to the incoming particles. Once $`c_{1,2}`$ are generated, together with the energy fractions $`x_{1,2}`$, and the azimuthal angles $`\varphi _{1,2}`$, the momenta of two ISR photons are known. The four-fermion event is then generated in the c.m.s. of the incoming particles after QED radiation, and then boosted back to the laboratory frame.
We also take into account non leading terms with the substitution
$$\mathrm{ln}(\frac{q^2}{m_e^2})1\frac{1}{1c_i+2\frac{m_e^2}{q^2}}2\frac{m_e^2}{q^2}\frac{1}{(1c_i+2\frac{m_e^2}{q^2})^2}.$$
The above choice ensures that the residue of the soft-photon pole gets proportional to $`\mathrm{ln}(\frac{q^2}{m_e^2})1`$, after integration over $`c_i`$.
As to the scale $`q^2`$, $`s`$ should be taken for $`s`$-channel dominated processes, while, when a process is dominated by small $`t`$ exchanges and $`t`$ is much smaller than $`s`$, the scale is related to $`t`$. This is e.g. . the case in small angle Bhabha scattering and the proper scale is chosen as the one which reproduces roughly the exact first order QED correction, which is known for Bhabha scattering. A similar procedure now also exists for the multi-peripheral two photon process , since an exact first order calculation is also available . In these $`t`$-channel dominated processes it is important to know whether a cross-section with angular cuts is wanted, since then the $`t`$-related scale will increase and the QED corrections as well. When no exact first order calculations are available the scale occurring in the first order soft corrections is also used as guideline to guess $`q^2`$ .
In NEXTCALIBUR the choice of the scale is performed automatically by the program, event by event, according to the selected final state (see Tab.(30)).
#### Numerical results
In Tabs.(3132) we show single-$`W`$ numbers produced with the Modified Fermion-Loop approach, as discussed in the previous section. Comparisons are made with the EFL calculation of Ref. . The results of EFL are reproduced within $`2\%`$ accuracy for both leptonic and hadronic single-$`W`$ final states.
It should also be noted that, when neglecting Fermion-Loop corrections, one can directly compare NEXTCALIBUR with other massive Monte Carlo’s and one finds excellent agreement for single-$`W`$ production in the whole phase space.
In Figs. 5556 we show the $`\mathrm{cos}\theta _\gamma `$ and $`E_\gamma `$ distributions for the most energetic photon in the process $`e^+e^{}e^{}\overline{\nu }_eu\overline{d}(\gamma )`$. We used $`\sqrt{s}=200`$ GeV, $`|\mathrm{cos}\theta _e|>0.997`$ and $`M(u\overline{d})>45`$ GeV. Only ISR photons are taken into account, according to the scheme given in Tab.(30).
Note, however, that the recipe of using Structure Functions with a proper choice of scales is not enough to determine without ambiguity the pattern of the radiation in $`t`$-channel dominated processes. The reason is that, when $`|t|=q^2m_e^2`$, the Leading Order Structure Function approach fails and one has to introduce a minimum value for $`|t|`$, below which only non-radiative events from the corresponding leg are generated. Since Structure Functions behave like $`\delta `$ functions for vanishing $`q^2`$, this is automatically achieved by introducing a minimum value $`|t_{min}|`$, such that, for events with $`|t|<|t_{min}|`$, the scale in the corresponding Structure Function is always set equal to $`|t_{min}|`$<sup>18</sup><sup>18</sup>18Note, however, that this behavior is known and could be implemented. It is enough to consider the standard YFS infrared emission factor $`\widehat{}`$, see e.g. Eq.(3) of Ref. .. We observed deviations at the order of $`0.5\%`$ by varying $`|t_{min}|`$ from $`2.71828m_e^2`$ to $`100m_e^2`$. The default value of $`|t_{min}|`$ in NEXTCALIBUR is taken to be the latter.
In table (33) we also give the cross sections (in pb) corresponding to the above distributions. tot refers to radiative plus non radiative events (within the specified separation cuts for the generated photons), nrad to non-radiative events, srad to single-radiative events and drad to double radiative events.
Finally, in order to quantify the effects due to ISR scales and running of $`\alpha _{\mathrm{QED}}`$, we show, in tables (34) and (35), the cross sections obtained by using both ISR scales = $`s`$ and switching on and off the Modified Fermion loop corrections.
#### Single-$`W`$ with GRACE
#### Authors
#### Introduction
The single-$`W`$ production processes present an opportunity to study the anomalous triple-gauge-couplings (hereafter TGC) at LEP 2 experiments. In order to proceed to the precise measurement of TGC, the inclusion of an initial state radiative correction (ISR) in any generator is an inevitable step. As a tool for the ISR the structure function (SF) and the parton shower methods are widely used for the $`e^+e^{}`$ annihilation processes. Since the main contribution for the single-$`W`$ production processes comes, however, from the non-annihilation type diagrams, the universal factorization method used for the annihilation processes is, obviously, inappropriate. The main problem lies in the determination of the energy scale of the factorization. According to the study of the two photon process, SF and QED parton shower (QEDPS) methods were shown to reproduce the exact $`O(\alpha )`$ results precisely even for the non-annihilation processes, when the appropriate energy scale is used in those algorithms.
Here, we propose a general method to determine the energy scale to be used in SF and QEDPS. The numerical results of testing SF and QEDPS for $`e^{}e^+e^{}\overline{\nu }_eu\overline{d}`$ and $`e^{}e^+e^{}\overline{\nu }_e\mu ^+\nu _\mu `$ are given. The systematic errors are also discussed.
#### 6.2.1 Energy Scale Determination in QED corrections
Single-$`W`$ is not dominated by annihilation and, therefore, standard methods as $`s`$-channel structure functions fail to reproduce the correct result.
The factorization theorem for the QED radiative corrections in the LL approximation is valid independently of the structure of the matrix element of the kernel process. Hence structure functions (hereafter SF) and QEDPS must be applicable to any $`e^+e^{}`$ scattering processes. However, the choice of the energy scale in SF and QEDPS is not a trivial issue. For simple processes like $`e^+e^{}`$ annihilation into fermion pairs and two-photon process (with only the multi-peripheral diagrams considered so far), the evolution energy scale could be determined in terms of the exact perturbative calculations. However, for more complicated processes, this is not always possible. Hence a way to find a suitable energy scale without knowing the exact loop calculations should be established.
First we look at the general consequence of the soft photon approximation.
The soft photon cross-section is given, in some approximation, by the Born cross-section multiplied by the following correction factor :
$`{\displaystyle \frac{d\sigma _{soft}(s)}{d\mathrm{\Omega }}}`$ $`=`$ $`{\displaystyle \frac{d\sigma _0(s)}{d\mathrm{\Omega }}}\left|\mathrm{exp}\left[{\displaystyle \frac{\alpha }{\pi }}\mathrm{ln}\left({\displaystyle \frac{E}{k_c}}\right){\displaystyle \underset{i,j}{}}{\displaystyle \frac{e_ie_j\eta _i\eta _j}{\beta _{ij}}}\mathrm{ln}\left({\displaystyle \frac{1+\beta _{ij}}{1\beta _{ij}}}\right)\right]\right|^2,`$ (75)
$`\beta _{ij}`$ $`=`$ $`\left(1{\displaystyle \frac{m_i^2m_j^2}{(p_ip_j)^2}}\right)^{\frac{1}{2}},`$ (76)
where $`m_j`$ ($`p_j`$) are the mass(momentum) of $`j`$-th charged particle, $`k_c`$ is the maximum energy of the soft photon (the boundary between soft- and hard-photons), $`E`$ is the beam energy, and $`e_j`$ the electric charge in unit of the $`e^+`$ charge. The factor $`\eta _j`$ is $`1`$ for the initial particles and $`+1`$ for the final particles. The indices ($`i,j`$) run over all the charged particles in the initial and final states.
The part proportional to $`\mathrm{ln}(E/k_c)`$ that is shown explicitly in Eq.(75) is exact and not only LL-approximated. However, the single-logarithmic part is omitted, so that the formula is not a complete LL-approximation, but it is enough to guess the energy scale appearing in SF and QEDPS. For the two-photon process, $`e^{}(p_{})+e^+(p_+)e^{}(q_{})+e^+(q_+)+\mu ^{}(k_{})+\mu ^+(k_+)`$, it was shown in Ref. that the soft-photon factor in Eq.(75) with a ($`p_{}q_{}`$)-term gives a good numerical approximation to the exact $`𝒪\left(\alpha \right)`$ correction.
This implies that one is able to make and educated guess about the possible evolution energy scale in SF from Eq.(75) without an explicit loop calculation.
However, one may question why the energy scale $`s=(p_{}+p_+)^2`$ does not appear in the soft-photon correction, even if they are included in Eq.(75). When applying SF to the two-photon process we have ignored those terms which come from the photon connecting different charged lines. This is because the contributions from the box diagrams, with photon exchange between the $`e^+`$ and $`e^{}`$ lines, is known to be small. Fortunately, the infrared part of the loop correction is already included in Eq.(75) and there is no need to know the full form of the loop diagram. For the two-photon processes we look at those two terms where, for example, ($`p_{}p_+`$)-terms and ($`q_{}p_+`$)-terms are present; here, the momentum of $`e^{}`$ is almost the same, before and after the scattering($`p_{}q_{}`$). The difference only appears in $`\eta _j\eta _k=+1`$ for a ($`p_{}p_+`$)-term and in $`\eta _j\eta _k=1`$ for a ($`q_{}p_+`$)-term. Then these terms compensate each other after summing them up for the forward scattering, which is the dominant kinematical region of this process. This is why the energy scale $`s=(p_{}+p_+)^2`$ does not appear in the soft-photon correction, despite its presence in Eq.(75).
When experimental cuts are imposed, for example the final $`e^{}`$ is tagged in a large angle, this cancellation is not perfect but only partial and the energy scale $`s`$ must appear in the soft-photon correction. In this case the annihilation-type diagrams will also contribute to the matrix elements. Then the usual SF and QEDPS formulation for the annihilation processes are justified and can be used for ISR with the energy scale $`s`$. One can check which energy scale is dominant under the given experimental cuts by numerically integrating the soft-photon cross-section given by Eq.(75) over the allowed kinematical region. Thus, in order to determine the energy scale it is sufficient to know the infrared behavior of the radiative process using the soft-photon factor.
Next, we determine the energy-scale of the QED radiative corrections to the single-$`W`$ production process,
$`e^{}(p_{})+e^+(p_+)`$ $``$ $`e^{}(q_{})+\overline{\nu }_e(q_\nu )+u(k_u)+\overline{d}(k_d).`$ (77)
The soft-photon correction factor shown in Eq.(75) is numerically integrated with the Born matrix element of the process (77), with $`t`$-channel diagrams only and without any cut on the final fermions. The results are shown in Table36.
One can see that the main contribution comes from an electron-line ($`p_{}q_{}`$-term) and a positron-line ($`p_+k_uk_d`$-term), while all the other contributions are negligibly small. As in the case of the two-photon processes, the energy scale $`s`$ does not appear in the soft-photon correction. Applying SF or QEDPS for the electron and positron charged-lines individually and with an energy scale given by their momentum-transfer squared might be legitimate, according to the above results.
#### 6.2.2 Structure Function Method
The corrected cross-section is given by
$`\sigma _{total}(s)`$ $`=`$ $`{\displaystyle 𝑑x_I𝑑x_F𝑑x_{I+}𝑑x_u𝑑x_dD_e^{}(x_I,t_{})D_e^{}(x_F,t_{})}`$ (78)
$`D_{e^+}(x_{I+},p_{Tud}^2)D_u(x_u,m_{ud}^2)D_d(x_d,m_{ud}^2)\sigma _0(\widehat{s}),`$
using the structure function ($`D_f`$) with an energy scale $`t_{}=(p_{}q_{})^2`$, $`p_{Tud}^2`$ i.e. the transverse-momentum squared of the $`u`$-$`\overline{d}`$ system and $`m_{ud}^2=(k_u+k_d)^2`$.
The energy-scale determination for the positron line is rather ambiguous. The $`p_{Tu+d}`$ is distributing around $`M__W/3`$, then the difference between these two energy scales does not give a significant effect on the correction factor. After(before) the photon radiation the initial(final) momenta $`p_\pm `$ ($`q_\pm `$) become $`\widehat{p}_\pm `$ ($`\widehat{q}_\pm `$) defined by:
$`\widehat{p}_{}`$ $`=`$ $`x_Ip_{},\widehat{q}_{}={\displaystyle \frac{1}{x_F}}q_{},\mathrm{}`$ (79)
Then the c.m.s. energy squared $`s`$ is scaled as $`\widehat{s}=x_Ix_{I+}s`$.
#### 6.2.3 Parton Shower Method
Instead of the analytic formula of the structure-function approach, a Monte Carlo method based on the parton shower algorithm in QED (QEDPS) can be used to solve the Altarelli-Parisi equation in the LL approximation. The detailed QEDPS-algorithm can be found in Ref. for the $`e^+e^{}`$ annihilation processes, in Ref. for the Bhabha process, and in Ref. for the two-photon process. In QEDPS we use the same energy scale as in the SF method. One difference between SF and QEDPS is that the ad hoc replacement of the perturbative expansion coefficient $`L(=\mathrm{ln}(Q^2/m_f^2))`$ by $`L1`$, which was realized by hand for SF, does not apply for QEDPS. Another significant difference between these two methods is that QEDPS can give a correct treatment of the transverse momentum of emitted photons by imposing the exact kinematics at the $`ee\gamma `$ splitting. Note that it does not affect the total cross sections too much when the final $`e^{}`$ are unconstrained. However, the finite recoiling of the final $`e^\pm `$ may result into a large effect on the tagged cross-sections.
As a consequence of the exact kinematics at the $`ee\gamma `$ splitting, the $`e^\pm `$ are no more on-shell after photon emission. On the other hand the matrix element of the hard scattering process must be calculated with the on-shell external particles. A trick to map the off-shell four-momenta of the initial $`e^\pm `$ to those at on-shell is needed. The following method is used in the calculations: First $`\widehat{s}=(\widehat{p}_{}+\widehat{p}_+)^2`$ is calculated, where $`\widehat{p}_\pm `$ are the four-momenta of the initial $`e^\pm `$ after the photon emission by QEDPS. $`\widehat{s}`$ is mainly positive even for the off-shell $`e^\pm `$. (When $`\widehat{s}`$ is negative, that event is discarded.) Subsequently, all four-momenta are generated in the rest-frame of the initial $`e^\pm `$ after the photon emission. Four-momenta of the hard scattering in their rest-frame are $`\stackrel{~}{p}_\pm `$, where $`\stackrel{~}{p}_\pm ^2=m_e^2`$ (on-shell) and $`\widehat{s}=(\stackrel{~}{p}_{}+\stackrel{~}{p}_+)^2`$.
Finally, all four-momenta are rotated and boosted to match the three-momenta of $`\stackrel{~}{p}_\pm `$ with those of $`\widehat{p}_\pm `$. This method respects the direction of the final $`e^\pm `$ rather than the c.m.s. energy of the collision. The total energy is not conserved because of the virtuality of the initial $`e^\pm `$. The violation of energy-conservation is of the order of $`10^6`$ GeV or less. The probability to violate it by more than $`1`$ MeV is $`10^4`$.
#### Numerical Calculations, the total cross-sections
Total and differential cross-sections of the semi-leptonic process $`e^{}e^+e^{}\overline{\nu }_eu\overline{d}`$ and of the leptonic one, $`e^{}e^+e^{}\overline{\nu }_e\mu ^+\nu _\mu `$, are calculated with the radiative correction by using SF or QEDPS. Fortran codes to calculate amplitudes of the above processes are produced using $`\mathrm{𝙶𝚁𝙰𝙲𝙴}`$ system. All fermion-masses are kept finite in calculations. Numerical integrations of the matrix element squared in the four-body phase space are done using $`\mathrm{𝙱𝙰𝚂𝙴𝚂}`$. For the study of the radiative correction for the single-$`W`$ productions, only $`t`$-channel diagrams(non-annihilation diagrams) are taken into account.
For the total energy of the emitted photons, both methods must give the same spectrum, when the same energy scale are used. That is confirmed by the results shown in Fig. 57 at the c.m.s. energy of $`200`$GeV for the semi-leptonic process.
Total cross-sections as a function of the c.m.s. energies at LEP 2 with and without experimental cuts are shown in Fig. 58. The experimental cuts applied here are $`M_{q\overline{q}}>45\mathrm{GeV}`$ and $`E_l>20\mathrm{GeV}`$.
The effect of the QED radiative corrections on the total cross-sections are obtained to be $`7`$ to $`10\%`$ on LEP 2 energies. If one uses the wrong energy scale $`s`$ in SF, the ISR effect is overestimated of about $`4\%`$ as shown in Fig. 59 both with and without cuts. For the fully extrapolated case the SF-algorithm with a correct energy scale is consistent with QEDPS within $`0.2\%`$. It may reflect the difference between $`L`$ and $`L1`$, as mentioned in Sect. 6.2.1. On the other hand, with the experimental cuts the SF-method at the correct energy-scale gives a deviation of around $`1\%`$ from QEDPS.
#### Numerical Calculations, the hard photon spectrum
Energy and angular distributions of the hard photon from QEDPS are compared with those from the calculations of the exact matrix elements. The cross-sections of the process $`e^{}e^+e^{}\overline{\nu }_eu\overline{d}\gamma `$ are calculated based on the exact amplitudes generated by $`\mathrm{𝙶𝚁𝙰𝙲𝙴}`$ and integrated numerically in five-body phase space using $`\mathrm{𝙱𝙰𝚂𝙴𝚂}`$. To compare the distributions, the soft-photon correction for the radiative process must be included. For this purpose QEDPS is implemented into the calculation of the process $`e^{}\overline{\nu }_eu\overline{d}\gamma `$ with a careful treatment aimed to avoid a double-counting of the radiation effect. The definition of the hard photon is $`E_\gamma >1\mathrm{GeV}`$ with an opening angle between the photon and the nearest final-state charged particles that is greater than $`5^{}`$. The distributions of the hard photons are in good agreement as shown in Fig. 60. The total cross-section of the hard photon emission is consistent at the $`2\%`$ level. On the other hand, if the soft-photon correction is not implemented on the radiative process, we end up with an over-estimate of $`30\%`$.
### 6.3 Technical precision in single-$`W`$
An old comparison for single-$`W`$ has been extended to cover
1. $`e^+e^{}q\overline{q}e\nu (\gamma )`$, $`|\mathrm{cos}\theta _e|>0.997`$, either $`M(q\overline{q})>45`$GeV or $`E_{q_1},E_{q_2}>15`$GeV, inclusive cross-section accuracy $`2\%`$, photon energy and polar angle ($`|\mathrm{cos}\theta _\gamma |<0.997(0.9995)`$) spectrum
2. $`e^+e^{}e\nu e\nu (\gamma )`$, $`|\mathrm{cos}\theta _e|>0.997`$, $`E_e>15`$GeV, $`|\mathrm{cos}\theta _e|<0.7(0.95)`$, inclusive cross-section accuracy $`5\%`$, photon energy and polar angle ($`|\mathrm{cos}\theta _\gamma |<0.997(0.9995)`$) spectrum.
3. $`e^+e^{}e\nu \mu \nu (\gamma )`$ and $`e^+e^{}e\nu \tau \nu (\gamma )`$, $`|\mathrm{cos}\theta _e|>0.997`$, $`E_{\mu /\tau }>15`$GeV, $`|\mathrm{cos}\theta _{\mu /\tau }|<0.95`$, inclusive cross-section accuracy $`5\%`$, photon energy and polar angle ($`|\mathrm{cos}\theta _\gamma |<0.997(0.9995)`$) spectrum.
With this comparison we want to check a) technical precision at the Born level, b) the correct inclusion of QED radiation, c) QCD corrections, especially in the low-mass region.
The first answer is that technical precision is not a problem anymore, all codes agree on single-$`W`$ cross-sections and distributions, even for $`\theta _e<0.1^{}`$, even for leptonic final states. On $`\sigma _{\mathrm{Born}}`$ the technical accuracy is $`0.1\%`$, the same for $`d\sigma /d\theta _e`$ for $`\theta _e0`$. Not only invariant-mass cuts, but also energy-cuts have been tested as shown in Tabs.(3738) and in Fig. 61.
#### 6.3.1 QCD corrections
QCD corrections are usually implemented in their naive form, a recipe where the total $`W`$-width is corrected by a factor
$$\mathrm{\Gamma }__W=\mathrm{\Gamma }__W^{\mathrm{EW}}\left(1+\frac{2}{3}\frac{\alpha _s(M__W^2)}{\pi }\right),$$
(80)
and the cross-section gets multiplied by $`1+\alpha _s(M__W^2)/\pi `$. In all those approaches where the Fermion-Loop is included or simulated, one should pay particular attention to QCD, for instance in WTO QCD corrections are incorporated in the evaluation of the complex poles by using the $`𝒪\left(\alpha \alpha __S\right)`$ vector-boson self-energies of Ref. (the location of the poles is gauge-invariant). Furthermore, the vertices are effectively corrected so that the relevant Ward identity remains satisfied. In a similar way WPHACT also includes QCD effects in the computation of the imaginary part of both the re-summed propagators and the vertices, to preserve gauge invariance.
To check the effect of QCD corrections we have compared WPHACT (IFL<sub>α</sub>) with WTO (EFL) for $`e\nu _eud`$ final states in LEP 2 configuration with and without QCD. The comparison is shown in Tab.(39) where the first error for WTO comes from a variation of the scale $`\mu `$ from $`\mu /2`$ to $`2\mu `$, where we adopt $`\mu =M__W`$ as the scale for light quarks and $`\mu =m_t`$ for the $`bb,bt`$ and $`tt`$ contributions. Therefore, QCD effects in single-$`W`$ are under control in those programs that implement them consistently with Fermion-Loop.
#### 6.3.2 Assessing the theoretical uncertainty in single-$`W`$
If we do not want to use the Fermion-Loop prediction then, by a careful examination of the most plausible re-scaling procedure, we end up with approximately $`1\%,2\%`$ and $`3\%`$ theoretical uncertainty to be assigned to the energy scale in the channels $`ud`$, $`\mu \nu _\mu `$ and $`e\nu _e`$ respectively. Therefore, a conservative estimate of the theoretical uncertainty would read as follows:
$`\pm 2÷3\%`$ from a tuned comparison among NEXTCALIBUR, WPHACT and WTO;
$`\pm 4\%`$: if one uses the wrong energy scale $`s`$ in SF, the ISR effect is, indeed, overestimated by approximately $`4\%`$ as shown in the subsequent analysis.
giving a conservative total upper bound of $`\pm 5\%`$, see Sub-Sect. 6.4 for a more complete discussion.
One should stress that most of the theorists were interested in gauge invariance issues due to unstable particle for CC20. The experimentalists, however, were asking from the beginning for ISR $`p_t`$ effects, comparison with QEDPSt, SF and YFS. Unfortunately, only few groups have been working on these issues.
In the previous section few recipes have been introduced to improve upon QED ISR; they are all equivalent insofar as they translate into different choices for the scale in the leading-logarithms of the structure functions. However this problem has not yet received its final solution and a full $`𝒪\left(\alpha \right)`$ calculation would be needed.
There is, however, an additional complication in the use of QED structure functions originating from mass effects. The single-$`W`$ is $`st`$\- channels and the $`t`$-channel parts look as in Fig. 62.
The corresponding cross-section is proportional to
$$𝑑\mathrm{\Phi }_3\frac{1}{\widehat{Q}^4}\widehat{L}_{\mu \nu }\widehat{W}_{\mu \nu },\widehat{Q}=\widehat{p}_{}q_{},\widehat{L}_{\mu \nu }=\frac{1}{2}\widehat{Q}^2\delta _{\mu \nu }+\widehat{p}__\mu \widehat{q}__\nu +\widehat{q}_\mu \widehat{p}_\nu $$
(82)
$`{\displaystyle 𝑑\mathrm{\Phi }_3\widehat{W}_{\mu \nu }}`$ $`=`$ $`\widehat{W}_1(\delta _{\mu \nu }+{\displaystyle \frac{\widehat{Q}_\mu \widehat{Q}_\nu }{\widehat{Q}^2}}){\displaystyle \frac{\widehat{Q}^2}{(\widehat{p}_+\widehat{Q})^2}}\widehat{W}_2𝒫_\mu 𝒫_\nu `$
$`𝒫^\mu `$ $`=`$ $`\widehat{p}_+^\mu {\displaystyle \frac{\widehat{p}_+\widehat{Q}}{\widehat{Q}^2}}\widehat{Q}^\mu `$
where $`\widehat{p}`$ and $`\widehat{q}`$ denote emission of soft and collinear photons. Usually, $`p_{}^2=\widehat{p}_{}^{}{}_{}{}^{2}=0`$ and $`q_{}^2=\widehat{q}_{}^{}{}_{}{}^{2}=0`$, and one writes $`\widehat{p}_{}=x_{\mathrm{in}}p_{}`$ and $`\widehat{q}_{}=q_{}/x_{\mathrm{out}}`$ with the kernel cross-section to be weighted with structure functions. Here, however, masses matter should not be neglected and the electrons are in a virtual state, i.e. off their mass-shell. A possible choice is to write
$`(\widehat{p}_{})^2`$ $`=`$ $`m_e^2+{\displaystyle \frac{1}{2}}(1\beta )(1x_{\mathrm{in}})sx_{\mathrm{in}}m_e^2`$
The important facts are that $`\widehat{Q}_\mu \widehat{L}_{\mu \nu }=0`$ owing to gauge invariance. Instead we get
$`\widehat{Q}_\mu \widehat{L}_{\mu \nu }`$ $`=`$ $`4(1x_{\mathrm{in}})m_e^2q_\nu +4(1{\displaystyle \frac{1}{x_{\mathrm{in}}}})m_e^2p_\nu .`$
Even if we insist in putting $`\widehat{p}_{}=x_{\mathrm{in}}p_{}`$ and $`p_{}^2=m_e^2`$, gauge invariance is violated by terms of $`𝒪\left(m_e^2/s\right)`$. The effect of constant terms on the $`Q^2`$-integrated photon flux-function can be as large as $`6\%`$. Gauge-invariance violation affects this term, resulting in some intrinsic theoretical uncertainty, although we expect that the effect will be strongly decreased after convolution with SF peaking at $`x_{\mathrm{in}/\mathrm{out}}=1`$. Alternatively one may adopt formulations where the electron remains on-shell after emission but at the price of having collinear photons of non-zero virtuality, $`(\widehat{p}p)^20`$.
It is worth noticing that, the rescaled incoming four-momenta are implemented in SWAP as $`\widehat{p}_\pm =(xE,0,0,\pm \sqrt{x^2E^2m_e^2})`$, by interpreting $`x`$ as the energy fraction after photon radiation, as motivated in Ref. . If required, $`p_{}/p_L`$ effects can be implemented in the treatment of ISR, by means of either $`p_{}`$-dependent SF or a QED Parton Shower algorithm . Therefore, in practice SWAP adopts a formulation that preserves on-shell incoming electrons. Furthermore, in NEXTCALIBUR it is possible to have both on-shell initial state particles and on shell generated photons but at the price of loosing part of the information on the direction of the initial states after radiation.
A final set of comments is needed to quantify the theoretical accuracy of single-$`W`$ production.
#### GRACE
The method to apply the QED radiative correction on the non-annihilation processes are established. The conventional method, SF with energy scale $`s`$ gives about $`4\%`$ overestimation for the QED radiative effect on the LEP 2 energies. If one wants to look at the hard photon spectrum, the soft-photon correction on these radiative processes are needed.
#### SWAP
The difference shown in Fig. 53 between the predictions given by the two set of $`Q^2`$ scales of Eq.(71) and Eq.(72) is at the per mille level, and therefore the simple naive scales of Eq.(72) are a good ansatz for the energy scale of QED radiation, which could be corroborated by the comparison with the results of other groups. QED corrections missing in the present approach are beyond the LL approximation. The present study shows that the choice $`Q_\pm ^2=s`$ as scale in the IS QED SF(s) can lead to over-estimate the effect of LL photonic corrections by a factor of $`1.5`$, implying an under-estimate of the QED corrected cross-section of about $`4\%`$. Also the choice of fixing the scale to $`Q_\pm ^2=|q_\gamma ^{}^2|`$ for both the IS SF(s), as recently suggested , leads to an under-estimate of the photon correction of about $`4\%`$. Since these effects are not negligible in the light of the expected theoretical accuracy, it is recommended to adopt the $`Q^2`$-scales as given in eq. (71) or eq. (72), which are motivated by the arguments sketched above.
Further, the effect of rescaling $`\alpha _{\mathrm{QED}}`$ from the high-energy value $`\alpha _{G_F}`$ to $`\alpha (q_\gamma ^{}^2)`$ amounts to a negative correction of about $`5÷6\%`$, to be taken into account carefully.
#### WPHACT
From the cross-sections of Tabs.(2627) one deduces that the difference between IFL and IFL<sub>α</sub> is of the order of $`6\%`$, for both $`e^{}e^+e^{}e^+\nu _e\overline{\nu }_e`$ and $`e^{}e^+e^{}e^+\nu _\mu \overline{\nu }_\mu `$. The discrepancy between IFL<sub>α</sub> and IFL<sub>α1</sub> predictions is always of the order of $`3\%`$ and one has, therefore, to attribute an estimate of $`3÷4`$ % error to the IFL<sub>α</sub> calculations for these processes.
Considering all processes together one can conclude that the implementation of a proper running $`\alpha _{\mathrm{QED}}`$ reduces the theoretical uncertainty by about one half with respect to fixed-width or imaginary fermion-loop alone. In some cases this uncertainty is further reduced to less than one percent, but only a comparison with complete EFL calculations, as a reference point, may assess whether this is the case. If no comparisons are available for the process and cuts at hand our study points towards a $`3\%`$ uncertainty for the calculations using the running $`\alpha _{\mathrm{QED}}`$, for both single-$`W`$ and single-$`Z`$ processes. Of course, one should add the uncertainty due to the fact that ISR for annihilation processes is not suited for $`t`$-channel contributions. Some obvious improvements on this point will soon be implemented in WPHACT, however a more careful study both for the theoretical uncertainty of these solutions and for a better treatment of $`t`$-channel ISR is still needed.
#### WTO
Bosonic corrections are still missing and, very often, our experience has shown, especially at LEP 1, that bosonic corrections may become sizeable . A large part of the bosonic corrections, as e.g. the leading-logarithmic corrections, factorize and can be treated by a convolution. Nevertheless the remaining bosonic corrections can still be non-negligible, i.e. , of the order of a few percent at LEP 2 . For the Born cross-sections $`1÷2\%`$ should, therefore, be understood as the present limit for the theoretical uncertainty. This will have to be improved, soon or later, since bosonic corrections are even larger at higher energies and the single-$`W`$ cross-section will cross over the $`WW`$ one at $`500`$ GeV.
### 6.4 Summary and conclusions
A fairly large amount of work has been done in the last years on the topic of single-$`W`$. In the previous sections we presented the most recent theoretical developments in single-$`W`$ and their implementation in the generators. There are common problematic situations with more or less equivalent solutions. One has to assign an error band to the cross-section for our partial knowledge of ISR, with or without $`p_t`$, and for the uncertainty in the scale of the running couplings. As for the energy scale in couplings we have an exact calculation based on the EFL-scheme which, at the Born-level (no QED) is known to be at the $`1\%`$ level of accuracy. EFL-scheme, however, is implemented only in one generator while the other offer a wide range of approximations based on the idea of re-scaling the cross-section. Furthermore, no program includes $`𝒪\left(\alpha \right)`$ electroweak corrections, not even in Weizsäcker-Williams approximation (for the subprocess $`\mathrm{e}\gamma \mathrm{W}\nu _\mathrm{e}`$), the counterpart of DPA in CC03.
A description of single-$`W`$ processes by means of the EFL-scheme is mandatory from, at least, two points of view. EFL is the only known field-theoretically consistent scheme that preserves gauge invariance in processes including unstable vector-bosons coupled to e.m. currents. Furthermore, single-$`W`$ production is a process that depends on several scales, the single-resonant $`s`$-channel exchange of $`W`$-bosons, the exchange of $`W`$-bosons in $`t`$-channel, the small scattering angle peak of outgoing electrons. A correct treatment of the multi-scale problem can only be achieved when we include radiative corrections in the calculation, not only one-loop terms but also the re-summation of leading higher-order terms. Recent months have shown that this project can be brought to a very satisfactory level by identifying the correct approximation, process-by-process. In particular, the $`WW`$ configuration, dominated by double-resonant terms, can be treated within DPA. As a consequence, the theoretical uncertainty associated with the determination of the $`WW`$ cross-section is sizably decreased. In principle, the same procedure applies to the determination of the $`ZZ`$ cross-section, where one develops a NC02-DPA instead of the CC03-DPA one.
We have found that all the modifications introduced via the EFL-scheme are relevant: running of the couplings, $`\rho `$-factors and vertices, not only the change $`\alpha _{\mathrm{QED}}(\mathrm{fixed})\alpha _{\mathrm{QED}}(\mathrm{running})`$. Therefore, a naive rescaling cannot reproduce the EFL answers for all situations, all kinematical cuts. The high-energy Input Parameter Set used in all calculations that are presently available – we quote, among the various schemes, the Fixed-Width scheme, the Overall scheme and the IFL one – is based on $`G_F,M__W`$ and $`M__Z`$ with $`\alpha _{\mathrm{QED}}(\mathrm{fixed})=1/131.95798`$. It allows for the inclusion of part of higher order effects in the Born cross-sections but, it fails to give a correct and accurate description of the $`q^20`$ dominated processes. A naive, overall, rescaling would lower the single-$`W`$ cross-section of about $`7\%`$. We have found, with the EFL calculation, that this decrease is process and cut dependent. Moreover, the effect is larger in the first bin for $`\theta _e`$$`0.0^{}÷0.01^{}`$ – in the distribution $`d\sigma /d\theta _e`$ and tends to become less pronounced for larger scattering angles of the electron. However, the first bin represents almost $`50\%`$ of the total single-$`W`$ cross-section, so that, in general, the compensations that occur among several effects never bring the EFL/FW ratio to one. We obtain a maximal decrease of about $`7\%`$ in the result but, on average, the effect is smaller. We have also found that the effect is rather sensitive to the relative weight of multi-peripheral contributions.
Finally, the effect of the QED radiative corrections on the total cross-sections are between $`7\%`$ and $`10\%`$ at LEP 2 energies. grc4f and SWAP have estimated that if one uses the wrong energy scale $`s`$ in the structure functions, the ISR effect is overestimated of about $`4\%`$, as shown in Figs. 5359, both with and without cuts. For the no-cut case SF with a correct energy scale is consistent with QEDPS around $`0.2\%`$. On the other hand with the experimental cuts, SF with correct energy-scale gives around $`1\%`$ deviation from QEDPS.
At the same time SWAP estimates that the effects due to two different scales (eq. (71) and eq. (72) are in good agreement and both predict a lowering of the Born cross-section of about $`8\%`$, almost independent of the c.m.s. LEP 2 energy. SWAP results show a good agreement with those of grc4f when both are referred to $`s`$-channel SF.
Although we register substantial improvements upon the standard treatment of QED ISR, the problem is not yet fully solved for processes where the non-annihilation component is relevant. A solution of it should rely on the complete calculation of the $`𝒪\left(\alpha \right)`$ correction, therefore the the basic YFS approach or any equivalent one augmented by virtual corrections.
At the moment, a total upper bound of $`\pm 5\%`$ theoretical uncertainty should still be assigned to the single-$`W`$ cross-section. In particular, the difference between annihilation-like QED radiation and the optimized scales amounts to a $`4\%`$, which is conservatively used (by the LEP EWWG) in the global estimate of theoretical uncertainty. Alternatively one should use the differences between different implementations of ISR in the $`t`$-channel as a basis for the systematic uncertainty. However, we are not yet ready to formulate a strict and definitive statement along these lines. Furthermore, there seems to be and indication of some numerical difference arising from different QED treatments in GRACE and in SWAP. At present no direct comparison has been attempted to understand the origin. We could say that QED radiation in single-$`W`$ is understood at a level better than $`4\%`$ but we are presently unable to quantify this assertion.
In this sense the current $`5\%`$ should be considered as a good estimate of the global upper limit for theoretical uncertainty. The origins of this upper bound are as follows. QED effects are bounded by a $`4\%`$, saturated only by those programs that do not improve upon the scale. Effects due to running couplings and vertices are bounded by a $`2\%÷3\%`$, saturated by those programs that do not implement an exact massive FL-scheme. To lower this estimate is presently possible only in a multi-step procedure where program A is used vs. B to assess the effect of EFL/$`\alpha _{\mathrm{QED}}`$, then A is used vs. C to assess the effect of QED ISR and finally A is corrected to take into account the missing pieces and assign an uncertainty. This procedure should be performed within the experimental community, using the individual estimates of theoretical uncertainties as declared by the programs in this Section.
We expect an improvement upon this estimate when more implementation of the Fermion-Loop scheme will be available. Presently, the results with a rescaling of $`\alpha _{\mathrm{QED}}`$ for the $`t`$-channel photon show an agreement with EFL predictions that is between $`1\%`$ and $`2\%`$. Note that in $`e\nu e\nu `$ EFL is not yet implemented and there we use the estimate by WPHACT of roughly $`3\%`$. All programs that implement the correct running of $`\alpha _{\mathrm{QED}}`$ should be able to reach this level of accuracy, but not all of them have this implementation.
All program that still implement $`s`$-channel structure functions saturate the $`5\%`$ level of theoretical accuracy. Further and more complete studies are needed for QED corrections and ad hoc solutions, like fudge factors should be avoided.
As stated above, the present level of global theoretical uncertainty, $`5\%`$, comes from different sources and different effects. Some of them have been fully understood from a theoretical point of view but, sometimes, not yet implemented in most of the programs. There are remaining problems that have not yet received a satisfactory solution and some of the programs implement educated guesses. In general we should say that single-$`W`$ remains, to a large extent, the land of fudge factors. As for individual programs, the following collaborations (listed in alphabetic order) have agreed to quantify their performances:
NEXTCALIBUR tries to include all leading higher order effects. At present, by comparing with the Exact-Fermion Loop and varying the internal parameters of the program, we can assign a conservative $`3\%`$ uncertainty coming from our Modified fermion loop approach. On the other hand, our solution to the $`t`$-channel ISR problem represents the best we have so far at the theoretical level. Therefore, the final precision of $`\pm 5\%`$ on single-$`W`$ has to be considered as a safe estimate of the accuracy reached by the program, at least in absence of large angle hard photons.
SWAP includes exact tree-level matrix elements with finite fermion masses and anomalous trilinear gauge couplings, the effect of vacuum polarization, higher-order leading QED corrections according to the treatment for the energy scale as given by the (equivalent) choices of eq. (71) and eq. (72). Since, apart from the effect of the running of $`\alpha _{\mathrm{QED}}`$, other one-loop fermionic and bosonic corrections are still missing in SWAP, its theoretical uncertainty is at the level of $`23\%`$<sup>19</sup><sup>19</sup>19to be compared with the estimated upper bound of $`5\%`$, depending on the channel and/or event selection considered.
WPHACT can be used for single-$`W`$ in its version IFL$`\alpha `$. This is at present the best choice: all other schemes have been employed for studies and comparisons but are not recommended. As already explained, the theoretical uncertainty due to non implementation of the complete EFL amounts to $`34\%`$ for CC20/Mix56. This, together with the non correct QED radiation for $`t`$-channel, leads to an estimate of $`56\%`$ accuracy in actual WPHACT single-$`W`$ predictions.
WTO can only be used as a benchmark for the determination of the scales in the coupling constants. In its default WTO saturates the upper bound of $`5\%`$ of accuracy. Ideally, the difference between any program using some approximation and WTO should be considered as systematic uncertainty for the scale determination (in couplings) of that program. In practice EFL, the right approach, is only implemented in WTO and a cross-check is needed before being able to apply the previous rule. The correct treatment of QED radiation is still missing, it is a choice of the author to avoid ad hoc solutions and a consistent upgrading is currently under study. Furthermore, Fermion-Loop (as DPA) implies certain characteristics and programs that implement a incomplete-FL that does not reflect at least a large fraction of them should refrain from using the label FL.
The collaborations each take responsibility for the above statements that range from conservative to more optimistic ones.
### 6.5 Outlook
A substantial amount of work was done in the last two years on the topic of single-$`W`$ production. This has triggered theoretical developments which can be used also in other areas, e.g. massive Fermion-Loop scheme, QED radiation in processes dominated by $`t`$-channel diagrams. One of the main results of the theoretical activity has been to upgrade programs that where available prior to the workshop and did not provide a satisfactory simulation of the process. They might have given a numerically more or less correct cross section, but this was mostly an accident.
Some work has not yet been done, e.g. low–invariant-mass $`e\nu +`$hadrons final states (searches), DPA-equivalent set of radiative corrections (high-luminosity LC). Finally, we still do not have a complete, solution to the ISR problem, although there has been considerable progress in the treatment of QED radiation, in particular in the determination of the radiation scale. Going beyond the present level of theoretical accuracy would require a complete $`𝒪\left(\alpha \right)`$ calculation therefore contributing to improve the present level of theoretical accuracy.
### Appendix: the diagrams
## 7 The NC02 cross-section, $`\sigma _{ZZ}`$
The cross-section for $`e^+e^{}ZZ`$ is defined starting from the NC02 process, very much as $`e^+e^{}W^+W^{}`$ in terms of CC03, hence we sum over all channels, $`e^+e^{}Z(f_1\overline{f}_1)+Z(f_2\overline{f}_2)`$ including four neutrinos in the final state. The electroweak corrections to $`eeZZ`$ were calculated in . Actually there the weak corrections have been discussed separately, but unfortunately in the $`\alpha `$ scheme.
As usual, it is left for the experimenters to evaluate the background, i.e. to define a neutral current observable cross-section as follows:
$$\sigma _{\mathrm{NC}}=\sigma _{\mathrm{NC02}}\left(1+\delta _{\mathrm{NC}}^{\mathrm{DPA}}\right)+\left[\sigma _{4\mathrm{f}}\sigma _{\mathrm{NC02}}\right].$$
(83)
The theoretical prediction, therefore, should concentrate on $`\sigma _{\mathrm{NC02}}`$, with or without $`𝒪\left(\alpha \right)`$ radiative correction in DPA-approximation. In particular, the background should account for the Mixed processes (Mix43). There is some important remark to be made. When dealing with $`u\overline{u}u\overline{u}`$ etc, i.e. with channels containing identical particles, we have to evaluate the unphysical sum of the two diagrams corresponding to $`e^+e^{}Z(u_1\overline{u}_1)+Z(u_2\overline{u}_2)`$, tacitly assuming that there are two $`u`$ quarks, $`u`$ of type 1 and $`u`$ of type 2. Since the interferences between the crossings are not double-resonant, it is customary to consider them as background and to define the $`ZZ`$ signal, i.e. $`\sigma _{\mathrm{NC02}}`$, from the absolute squares of the double-resonant diagrams only. This is a matter of definition, i.e. , we could define the $`ZZ`$ signal to contain all crossings in case of four identical flavors in the final state. That one chooses the first option is largely based on the drawback that, with the latter,
$$\sigma (e^+e^{}ZZ)\times \mathrm{BR}^2(Zu\overline{u}),$$
(84)
is no longer $`\sigma _{ZZu\overline{u}u\overline{u}}`$. It is certainly true that the cross-section containing all crossings would be more physical but, for the time being this is the convention. Furthermore, one should remember that the $`\mathrm{e}^+\mathrm{e}^{}\gamma ^{}\gamma ^{},\gamma ^{}\mathrm{Z}`$ background is quite large (see e.g. Ref. ).
Bearing this in mind, we should stress that the terminology $`\sigma ^{\mathrm{NC02}}`$ is, sometimes, unfaithful, simply because this is not what experiments use in their analysis. A common procedure is to use EXCALIBUR and to restrict it to the complete set of double-resonant diagrams. In other words, experiments measure data in some window of invariant mass and extrapolate with some coefficient, evaluated by MC, to what one finally calls the NC02-total cross-section but it represents, instead, the sum of all double-resonant $`Z`$ diagrams (for some channel $`4`$ instead of $`2`$).
However, by definition, we select NC02 to be $`e^+e^{}ZZ`$, two diagrams ($`t`$ and $`u`$ channel), with all $`Z`$ decay modes allowed for both $`Z`$-bosons. If one computes everything as production $``$ decay then, as long as one remembers to include factors $`1/2`$, everything is reasonable. The conclusion is based on the following observation. When all diagrams are taken into account we find
$`\sigma (e^+e^{}\overline{u}u\overline{c}c)`$ $`=`$ $`208.9\mathrm{fb},\sigma (e^+e^{}\overline{u}u\overline{s}s)=204.4\mathrm{fb},`$
$`\sigma (e^+e^{}\overline{d}d\overline{s}s)`$ $`=`$ $`182.6\mathrm{fb},\sigma (e^+e^{}\overline{u}u\overline{d}d)=1.980\mathrm{pb},`$
$`\sigma (e^+e^{}\overline{u}u\overline{u}u)`$ $`=`$ $`101.4\mathrm{fb},\sigma (e^+e^{}\overline{d}d\overline{d}d)=87.88\mathrm{fb},`$ (85)
and, as a consequence,
$$R_{uucc/uuuu}=2.06,R_{ddss/dddd}=2.08.$$
(86)
In other words, even if we define on-shell and compute off-shell the same result, within few percents, is obtained.
The relative significance of the $`ZZ`$ cross-section is considerably less than the one attributed to the $`WW`$ cross-section. Its is smaller and with much larger experimental errors, even at the level of projected ones. As a consequence the NC02 process has received less attention that the CC03 one and, so far, we have no published result on $`𝒪\left(\alpha \right)`$ DPA calculations for it although, in principle, there is no major obstacle to it.
### 7.1 Description of programs and results
#### YFSZZ
#### Authors
#### General Description
The program evaluates the NC02 double resonant process $`e^+e^{}ZZ4\mathrm{f}`$ in the presence of multiple photon radiation using Monte Carlo event generator techniques. The theoretical formulation is based, in the leading pole approximation (LPA), on $`𝒪\left(\alpha ^2\right)`$ LL YFS exponentiation for the production process, with the possibility of anomalous gauge couplings if the user so desires. The Monte Carlo algorithm used to realize the YFS exponentiation is based on the YFS2 algorithm presented in Ref. and in Ref. . In this way, we achieve an event-by-event realization of our calculation in which arbitrary detector cuts are possible and in which infrared singularities are cancelled to all orders in $`\alpha `$. A detailed description of our work can be found in Ref. .
#### Features of the program
The code is a complete Monte Carlo event generator and gives for each event the final particle four-momenta for the entire $`4f+n\gamma `$ final state over the entire phase space for each final state particle. The events may be weighted or unweighted, as it is more or less convenient for the user accordingly. The code features the realization of the LPA for the NC02 process that is the analog of that given in Ref. for the CC03 process of production and decay of $`WW`$ pairs. A technical precision check on the program at the level of 2 per mille for the total cross-section has been done by comparison with the results in Ref. . The accuracy of the combined result from YFSZZ 1.02 and KoralW 1.42, when the combination is taken in analogy with that presented in Ref. for YFSWW3 1.14 and KoralW 1.42, is expected to be at the level of $`2\%`$ for the total cross-section, due to the missing $`𝒪\left(\alpha \right)`$ pure weak corrections in YFSZZ 1.02 (we do not expect the other effects missing from our calculation such as non-universal QED corrections to enter at this level), when all tests are finished. These tests are currently in progress.
The operation of the code is entirely analogous to that of the MC YFS2 in Refs. . A crude distribution based on the primitive Born level distribution and the most dominant part of the YFS form factors that can be treated analytically is used to generate a background population of events. The weight for these events is then computed by standard rejection techniques involving the ratio of the complete distribution and the crude distribution. As the user wishes, these weights may be either used directly with the events, which have the four-momenta of all final state particles available, or they may be accepted/rejected against a maximal weight WTMAX to produce unweighted events via again standard MC methods. Standard final statistics of the run are provided, such as statistical error analysis, total cross-sections, etc. The total phase space for the process is always active in the code.
#### Description of output and availability
The program prints certain control outputs. The most important output of the program is the series of Monte Carlo events. The total cross-section in $`fb`$ is available for arbitrary cuts in the same standard way as it is for YFS2, i.e. the user may impose arbitrary detector cuts by the usual rejection methods.
The program is available from the authors via e-mail. The program is currently posted on WWW at http://enigma.phys.utk.edu as well as on anonymous ftp at enigma.phys.utk.edu in the form of a tar.gz file in the /pub/YFSZZ/ directory together with all relevant papers and documentation in postscript.
#### ZZTO
#### Author
#### Description.
ZZTO is a newly created code for computing $`\sigma ^{\mathrm{NC02}}`$ which, at the moment, has universal Initial State QED, Final State QED, Final State QCD, is fully massive with $`b`$ and $`c`$ quarks running masses. Fermion-Loop is also implemented. ZZTO is missing non-universal QED ISR and purely weak effects (in DPA); however, it is under construction with the final goal of including those effects. The code sums over all $`ZZ`$ decay modes, even $`\nu \nu \nu \nu `$. However, single channels are available, i.e. $`qqqq,qq\nu \nu ,qq\mathrm{ll},\mathrm{ll}\nu \nu ,\mathrm{llll},\nu \nu \nu \nu `$. Therefor, inside ZZTO we have the exact matrix element for $`e^+e^{}ZZ4\mathrm{f}`$ with massive fermions and running masses for the $`b,c`$-quarks. Cuts are only implemented on the $`Z`$ invariant masses, therefore we can apply final state QCD correction factors beyond the usual naive correction. In other words, the total hadronic decay rate of each $`Z`$-boson is split into the sum of the vector current induced rate, $`\mathrm{\Gamma }^V`$, and of the axial decay rate, $`\mathrm{\Gamma }^A`$, which receive different QCD corrections evaluated at the scale equal to the virtuality of the $`q\overline{q}`$-pair. Non-factorizable QCD corrections are neglected. Final state QED corrections are also included, again evaluated at the virtuality of the pair, i.e. with $`\alpha _{\mathrm{QED}}(M_{\mathrm{pair}}^2)`$. Initial state QED corrections include, so far, only the universal part of the structure functions evaluated at the scale $`s`$.
To implement the Fermion-Loop scheme we had to incorporate QCD corrections in the evaluation of the complex pole $`p__Z`$ and of the $`\rho `$-parameter associated to the $`Z`$-propagator. This we have done by taking into account also the massive top quark, while the light quarks, including the $`b`$ one, are treated as massless. QCD is exactly implemented by using the $`𝒪\left(\alpha \alpha __S\right)`$ vector-boson self-energies of Ref. with $`M__Z`$ as the scale for light quarks and $`m_t`$ for the $`bb,bt`$ and $`tt`$ contributions. For $`M__W=80.350\mathrm{GeV},M__Z=91.1888\mathrm{GeV}`$ and $`\alpha __S(M__Z^2)=0.120`$ we find a QCD effect illustrated in Tab.(40).
The program ZZTO is currently posted on WWW at http://www.to.infn.it/giampier/zzto.
#### Distributions.
The $`ZZ`$-signal is basically defined through invariant masses, for instance $`e^+e^{}q\overline{q}l^+l^{}(\gamma )`$, $`q`$-flavour blind or heavy $`q`$-flavors, $`l=e/\mu /\tau `$, $`|\mathrm{cos}\theta _{l_1}|<0.985`$, no cut on the second lepton (only one lepton tagged), $`M(q\overline{q})>10(45)`$GeV.
Here, we do not discuss invariant mass distributions in terms of the full processes but only in terms of the signal NC02. The angular cuts are there only because of detector holes at the beam pipe. Since for NC02 there are no poles at edge of phase space, we could leave these cuts out for simplicity. Furthermore, we analyze only $`e^+e^{}q\overline{q}l^+l^{}(\gamma )`$ where the definition of invariant masses is free of ambiguities. ZZTO includes final state radiations in two different options. In a first case ZZTO implements the exact, factorizable, $`𝒪\left(\alpha \right)`$ corrections for some extrapolated setup were one can only cut on the $`Z`$-virtuality, see Ref. . In the second one, hard and collinear photons are included, within a cone of angular resolution $`\delta 1`$, according to the formalism of Ref. . Moreover, soft photons are exponentiated.
Therefore, we can define invariant mass distributions according to the following choices: a) $`M(l^+l^{}\gamma )`$ or $`M(\overline{q}q\gamma )`$ where $`M`$ represents the virtuality of the decaying $`Z`$-boson and b) $`m(l^+l^{})`$ where $`m`$ is the $`l^+l^{}`$ invariant mass and hard photons are included whenever the angle between the photon and the nearest charged final-state fermion is less than $`\delta 1`$. Above $`\delta `$ photons are not included in the mass calculation. Gluons are always included in $`𝒪\left(\alpha __S\right)`$ with a fully extrapolated setup, i.e. the $`M`$-variable for $`\overline{q}q`$ final states is always understood as $`Z(M)\overline{q}q+\gamma +g`$.
In Fig. 71 we show the $`M`$-distribution for $`e^+e^{}+`$hadrons and for $`\overline{b}b(\overline{c}c)+`$leptons at one energy, $`\sqrt{s}=188.6`$GeV. There is no appreciable difference with $`\mu ^+\mu ^{}+`$hadrons due to the fact that the FSR correction factor is approximately $`3/4Q_f^2\alpha /\pi `$ since we cut on the $`Z`$-virtuality and not on the $`\overline{f}f`$ invariant mass.
In Fig. 72 we show $`e^+e^{}+`$hadrons and compare $`M`$ and $`m`$ distributions for the $`e^+e^{}(\gamma )`$ pair. The latter includes collinear photons within a cone of half-opening angle $`\delta =5^{}`$. In the same figure we also compare the $`m(\overline{f}f)`$ distributions for $`e^+e^{}+`$hadrons and $`\mu ^+\mu ^{}+`$hadrons. Since the cut is on the invariant mass of the pair one starts appreciating differences between different flavors.
All distributions are computed by ZZTO in the Fermion-Loop mode. The largest effects in the theoretical uncertainty are associated to the fact that non-factorizable QED corrections are neglected, although one can show that they vanish in the limit of on-shell $`Z`$-bosons, $`M^2M__Z^2\mathrm{\Gamma }_ZM__Z`$. They also vanish for a fully extrapolated setup, i.e. after integrating over the full range of the two $`Z`$ virtualities, which is not the case for distributions.
#### GENTLE
#### Authors
The NC cross-sections in package 4fan include now besides the NC32 class also the NC02 process (NC08 is unchanged); also some new options introduced, ICHNNL=0,1: switching between NC02 and NC32 classes (for IPROC=2); Note, that the treatment of NC08 sub-family is not changed compared to the version v.2.10. It remains accessible only via NCqed branch of the package.
### 7.2 Comparisons for the NC02 cross-section
In this Section we will compare the NC02 cross-section between YFSZZ, GENTLE and the newly created code ZZTO. First, the comparison between YFSZZ and ZZTO. Here, $`\sqrt{s}=188.6GeV`$ and QCD is not included. The result is shown in Tab.(41). From Tab.(41) we see a remarkable agreement, further quantified in Tab.(42). Furthermore, for $`\sigma _{ZZ}`$ with Born+ISR+QCD the uncertainty related to the IPS (Input Parameter Set) is approximately $`1\%`$. This does not mean that the total, true, theoretical uncertainty is $`1\%`$. The $`ZZ`$ line-shape, as predicted by ZZTO and including QCD corrections is shown in Tab.(44) where the results refer to three schemes, $`\alpha ,G_F`$ and Fermion-Loop.
Finally, in Fig. 73 we present the NC02 line-shape for a wide range of energy, comparing the $`\alpha `$-scheme with the $`G_F`$-scheme and the Fermion-loop one. Missing an implementation of the Fermion-Loop scheme in other codes, our recommendation is to use the $`G_F`$-scheme since it allows us to include part of higher order effects in the Born cross-sections.
In Tab.(45) we show the $`\sigma _{ZZ}`$ cross-section as predicted from GENTLE. Tab.(45) is produced with the following GENTLE/4fan flag settings:
IPROC,IINPT,IONSHL,IBORNF,IBCKGR,ICHNNL = 2 2 1 1 0 0
IGAMZS,IGAMWS,IGAMW,IDCS,IANO,IBIN = 0 0 0 0 0 0
ICONVL,IZERO,IQEDHS,ITNONU,IZETTA = x x 3 0 1
ICOLMB,IFUDGF,IIFSR,IIQCD = 0 0 1 0
IMAP,IRMAX,IRSTP,IMMIN,IMMAX = 1 0 1 1 1
and with the following NCqed branch settings:
IPROC ,IINPT ,IONSHL,IBORNF,IBCKGR,ICHNNL= 3 2 1 1 1 2
IGAMZS,IGAMWS,IGAMW ,IDCS ,IANO ,IBIN = 0 0 0 0 0 0
ICONVL,IZERO ,IQEDHS,ITNONU,IZETTA = 0 1 x x 1
ICOLMB,IFUDGF,IIFSR ,IIQCD = 2 1 1 0
IMAP ,IRMAX ,IRSTP ,IMMIN ,IMMAX = 1 0 1 1 1
The Table deserves an extended comment. Its upper part is obtained with the aid of the standard GENTLE approach to ISR: the band of theoretical uncertainties is produced by choosing standard structure functions (SF) for the minimum and flux functions (FF) for the maximum with a reasonable choice in between for the preferred one. For the maximum, we include LLA second order corrections and exclude the lowest order constant term (option IZERO=0). The band has a typical width of about $`3÷4\%`$. This approach finds its roots in the treatment of the CC03 cross-section where we used the so-called current-splitting technique, the precision of which is difficult to evaluate since it takes into account only a part of diagrams. We emphasize again that nowadays, after the advent of DPA calculations, the theoretical uncertainties in the CC-sector are reduced.
For NC-processes, the ISR is well defined and no current-splitting is required. In paper we provided the complete lowest order ISR QED corrections (option ITNONU=1). In our complete calculations the constant term is full reproduced and there are no justifications to exclude it. This is why in the lower part of the Table we always use IZERO =1. For the theoretical uncertainties, we vary then over three working options IQEDHS,ITNONU=00,10,11 and select preferred, min and max out of them. As seen from the lower part of the Table, the theoretical uncertainty derived in such a way is about twice as narrow as compared to the upper part. It is important to emphasize that the two bands overlap, although there is a systematic shift towards slightly higher cross-sections.
This shift is due to the constant term. If we had chosen IZERO=1 for the upper part, its band would totally contain the band for the lower part. We tend to consider the lower part to be a more correct treatment of the ISR for the case of NC-processes.
### 7.3 Summary and conclusions
Three different programs have produced numbers for the NC02 cross-section showing remarkable agreement over a wide energy range. ZZTO has produced results with two different renormalization schemes, $`G_F`$ and $`\alpha `$, showing differences of the order of a percent. GENTLE confirms the finding with nearly the same shifts as ZZTO between the two schemes. It looks plausible to have a $`\pm 2\%`$ of theoretical uncertainty assigned to the NC02 cross-section. There is an indications, coming from the Fermion-Loop analysis of ZZTO, that show smaller deviations with respect to the $`G_F`$-scheme and the Fermion-Loop is usually accurate at the $`1÷2\%`$ level.
At the moment the estimated theoretical uncertainty comes from the comparisons between GENTLE, YFSZZ and ZZTO and it is roughly about $`2\%`$. The size of the uncertainty is confirmed by an internal estimate of GENTLE, as given in Tab.(45). With the complete lowest order ISR QED included GENTLE gives a total cross-section at $`\sqrt{s}=188.6`$GeV of $`627.37_{2.14}^{+0.15}`$fb where ZZTO gives $`621.22`$fb, i.e. GENTLE predicts a $`0.4\%`$ uncertainty with GENTLE and ZZTO differing by roughly $`1\%`$. Furthermore, GENTLE predicts a $`+0.6\%`$ shift due to the constant term in ISR and both programs predict a $`0.8\%`$ shift from the $`G_F`$-scheme to the $`\alpha `$-scheme.
Given the experimental uncertainty on the cross-section a difference below $`2\%`$ is reasonable and, most likely, do not require the implementation of missing effects which are beyond the reach of the experiments. Nevertheless, work is in progress for ZZTO towards a complete DPA calculation for NC02.
## 8 Conclusions and outlook
An extensive collection of theoretical predictions for observables in $`e^+e^{}`$ interactions at LEP 2 energies had been presented in the 1996 CERN Report of the Workshop on Physics at LEP2. However, an update with improved theoretical prescriptions is needed in order to match the precision achieved by now in the experimental analyses.
The aim of the four-fermion contribution to this workshop effort is twofold. We have summarized the most recent theoretical developments concerning $`e^+e^{}`$ annihilation into four-fermions at LEP 2 energies. Furthermore, applications to the four most important classes of processes have been discussed in detail. In decreasing order of importance they are the $`WW`$-signal, the inclusion of an extra photon in the final state, the single-$`W`$ production and the $`ZZ`$-signal.
To gauge the priorities of this Report one should remember that the experimental situation is rather different for $`WW`$ when compared to the other processes. For $`W`$-pairs, LEP (ADLO) is able to test the theory to below $`1\%`$, i.e. , below the old uncertainty of $`\pm 2\%`$ established in 1995. Thus the CC03-DPA, including non-leading electroweak corrections, constitutes a very important theoretical development. However, ADLO cannot test single-$`W`$ or $`ZZ`$-signal to an equivalent level, since their total cross-section is of the order of $`1pb`$ or less, $`20`$ times smaller than that of $`W`$-pair production <sup>20</sup><sup>20</sup>20For $`ZZ`$ with 1997+1998+1999 data, the present analyses and global LEP combination method give an average measurement with $`7\%`$ accuracy. At the end of LEP, we may reach better than $`5\%`$..
The authors of the four-fermion report agree on the following conclusions from this study:
* There is a nice global agreement between the new DPA predictions for CC03, which are $`2\%÷3\%`$ lower than the old approach<sup>21</sup><sup>21</sup>21see Sect. 4 for a proper definition of the old approach..
* The Monte Carlo programs RacoonWW and YFSWW3 agree within $`0.3\%`$ at $`\sqrt{s}=200`$GeV. The present estimated theoretical uncertainty of these programs is $`0.4\%`$, $`0.5\%`$, and $`0.7\%`$ for $`\sqrt{s}=200\mathrm{GeV}`$, $`180\mathrm{GeV}`$, and $`170\mathrm{GeV}`$, respectively.
* There is a general satisfaction with the progress induced by new DPA calculations. Nevertheless, the theoretical uncertainty could probably be improved somewhat in the future.
* More work will be needed to reduce the uncertainty for $`4\mathrm{f}+\gamma `$ and of parton shower with $`p_t`$.
* In single-$`W`$ production most of the theorists were interested in gauge-invariance issues due to unstable particle. The experimentalists were asking for ISR and $`p_t`$ effects, comparisons including parton shower, structure functions and exponentiation. Unfortunately, only few groups have been working on these issues. Their work represents an important result of this Report.
* In single-$`W`$ production we have a (global) $`2\%÷3\%`$ theoretical uncertainty associated with the scale of the $`t`$-channel photon, with a projected $`1\%`$ uncertainty when the implementation of the Fermion-Loop scheme will receive more cross-checks.
* For simple processes like $`e^+e^{}`$ annihilation and two-photon collision, the evolution of the energy scale in the structure function or in the parton-shower algorithms can be determined by the exact perturbative calculations. However, this is not available for more complicated processes. When no exact first order calculations are available then one resorts to the scale occurring in the first order soft corrections. Therefore, at the moment, we may apply a very conservative (global) upper bound of $`4\%`$ theoretical uncertainty for ISR in single-$`W`$ production. Here we repeat one of the conclusions of Sect. 6, we understand the implementation of QED radiation in the MC better than before, Structure Functions at the scale $`s`$ are obviously wrong, but we are presently unable to precisely quantify the improvement upon the quoted – global – upper bound. Single programs may claim to have more stringent internal estimates. In conclusion, the current upper bound on the global estimate of the theoretical uncertainty is $`5\%`$ for single-$`W`$. A detailed explanation of this bound is given in Sub-Sect. 6.4.
* Compared to the experimental uncertainty on the NC02 $`ZZ`$ cross-section a difference of about $`1\%`$ between theoretical predictions is acceptable. The global estimate of theoretical uncertainty is $`2\%`$, again acceptable. However, it would be nice to improve upon the existing calculations.
These points are discussed in more detail in the following.
The new DPA predictions for CC03 are $`2\%÷3\%`$ lower than in the old approach. The new Monte Carlo programs RacoonWW and YFSWW3 agree within $`0.3\%`$ at $`\sqrt{s}=200\mathrm{GeV}`$, i.e. at a level that is consistent with the accuracy of the DPA. The theoretical uncertainty of these programs for the CC03 $`W`$-pair cross section, which we estimate to be below $`0.4÷0.5\%`$ for $`\sqrt{s}=180`$$`210\mathrm{GeV}`$, should be compared with the current experimental precision of $`\pm 0.9\%`$ with all ADLO data at $`183`$$`202`$GeV combined. It should be mentioned as well that RacoonWW and the semi-analytical BBC calculations agree very well where they should, i.e. above $`185\mathrm{GeV}`$.
Turning to distributions, the deviations seem to become somewhat larger for large $`\mathrm{W}^{}`$ production angles, although compatible with the statistical accuracy. The invariant mass distributions agree within roughly $`1\%`$ with a distortion of the distributions that is mainly due to radiation off the final state and the $`W`$ bosons. We expect that the present uncertainty of the CC03 $`W`$-pair cross-section can be reduced somewhat when the sources of the differences between RacoonWW and YFSWW and the leading higher-order corrections will be better analyzed. To go below the level of a few per-mille of accuracy would require the complete calculation of one-loop radiative corrections in four-fermion production for all 4f final states, a program that does not seem feasible in a foreseeable future.
The presence of real photons can also change the quantitative agreement of DPA calculations. For integrated quantities the differences between alternative approaches are expected to be of the order of the accuracy of the DPA while for more exclusive observables larger differences can be expected. A comparison between RacoonWW and YFSWW3 for various distributions in the semi-leptonic channel $`\mathrm{e}^+\mathrm{e}^{}\mathrm{u}\overline{\mathrm{d}}\mu ^{}\overline{\nu }_\mu `$ and with a specified set of separation and recombination cuts reflects, however, for observables inclusive in the photon the same global difference as the total cross-section.
The technical precision for $`e^+e^{}4\mathrm{f}+\gamma `$ has reached high standards as shown by the comparisons among PHEGAS/HELAC, RacoonWW and WRAP, but at the moment we are unable to present any overall statement on the theoretical uncertainty process by process. This is true in particular for the single-$`W`$ configuration. Furthermore, no detailed comparison has been performed including parton shower and hadronization.
In general more work will be needed to establish the uncertainty for $`4\mathrm{f}+\gamma `$. This should be done process by process, with the target of achieving the required accuracy. At the moment we can fix an upper bound of $`2.5\%`$ based on missing non-logarithmic corrections.
In single-$`W`$ production most of our activity was centered around gauge-invariance issues due to unstable particle. Although, no coordinate effort has been performed, at the moment, to study the theoretical uncertainty induced by ISR $`p_t`$ effects, comparison with parton shower, structure functions and exponentiation the interested reader can find in the Report details on QED corrections as they stand now. Few programs, noticeably GRACE and SWAP, have produced a preliminary internal estimate of the uncertainty associated with the treatment of QED radiation; the net effect of QED is between $`8\%`$ and $`10\%`$ in the LEP 2 energy range, with $`s`$-channel structure functions over-estimating the effect by $`4\%`$. Furthermore, structure functions with a modified scale seems to agree with parton shower at the level of $`1\%`$ when experimental cuts are included or even better for a fully extrapolated setup.
As far as the scale of the electromagnetic coupling is concerned we find that the results with a rescaling of $`\alpha _{\mathrm{QED}}`$ for the $`t`$-channel photon that has been implemented in NEXTCALIBUR, SWAP and WPHACT show an agreement with WTO predictions that is roughly around $`2\%`$.
For single-$`W`$, therefore, we register a conservative, overall, upper bound of $`\pm 5\%`$ for the theoretical uncertainty. Single programs may claim better internal estimates but this does not transform, yet, into a global one<sup>22</sup><sup>22</sup>22 We recall that, at the moment an uncertainty associated to QED ISR is quoted, by the LEP EWWG, that follows from taking the average of the Born result with the one corrected via $`s`$-channel structure functions, where $`\mathrm{SF}(t,p_{tW}^2)>\mathrm{SF}(s)`$ by $`+5\%`$ at $`200`$GeV and Born $`>\mathrm{SF}(s)`$ by $`+12\%`$. Note, however, that this is not a real estimate of uncertainty but just a pragmatic way of determining the effects of ISR..
Implementation of the EFL-scheme in single-$`W`$ (in addition to WTO) will give a more solid basis to the estimate of $`1÷2\%`$ for the uncertainty associated with the scale of the e.m. coupling.
The next, obvious, step is represented by the evaluation of missing $`𝒪\left(\alpha \right)`$ electroweak effects, e.g. in Weizsäcker-Williams approximation (for the sub-process $`\mathrm{e}\gamma \mathrm{W}\nu _\mathrm{e}`$), the analogous of DPA for CC03.
A better understanding of QED ISR and of all radiative corrections in single-$`W`$ production is certainly needed in order to reduce the corresponding uncertainty, hopefully around $`1\%`$. This, however, requires to go beyond the present approximations, not an easy task and with a considerably large experimental error. Since DPA cannot be applied to single-$`W`$ production one has to follow some alternative path, like including radiative corrections in (improved) Weizsäcker-Williams approximation, or WWA. It is expected that already the normal WWA ( i.e. logarithmic terms only), with a typical Born-accuracy of $`5\%`$, will yield results accurate at the level of $`5\%\times \alpha /\pi `$. For the moment this is not strictly needed but one should consider that single-$`W`$ will be one of the major processes at LC.
For the NC02 cross-section we have a $`1\%`$ variation, obtained by changing the Input Parameter Set in GENTLE and in ZZTO and by varying from the standard GENTLE approach for ISR to the complete lowest order corrections. We estimate the real uncertainty to be $`2\%`$. However, given the experimental uncertainty a theoretical uncertainty in this order is acceptable and does not seem to require the implementation of missing effects. Furthermore, ZZTO which is not yet a DPA calculation agrees rather well with YFSZZ, roughly below the typical DPA accuracy of $`0.5\%`$, and the latter features the realization of the LPA for the NC02 process. The implementation of a DPA calculation, in more than one code, in the NC02 $`Z`$-pair cross-section will bring the corresponding accuracy at the level of $`0.5\%`$, similar to the CC03 case.
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# RELATIVISTIC CORRECTIONS TO THE SUNYAEV-ZEL’DOVICH EFFECT FOR CLUSTERS OF GALAXIES. V. MULTIPLE SCATTERING
## 1 INTRODUCTION
Compton scattering of the cosmic microwave background (CMB) radiation by hot intracluster gas — the Sunyaev-Zel’dovich effect (Zel’dovich & Sunyaev 1969; Sunyaev & Zel’dovich 1972, 1980a, 1980b, 1981) — provides a useful method to measure the Hubble constant $`H_0`$ (Gunn 1978; Silk & White 1978; Birkinshaw 1979; Cavaliere, Danese, & De Zotti 1979; Birkinshaw, Hughes, & Arnaud 1991; Birkinshaw & Hughes 1994; Myers et al. 1995; Herbig et al. 1995; Jones 1995; Markevitch et al. 1996; Holzapfel et al. 1997; Furuzawa et al. 1998; Komatsu et al. 1999). The original Sunyaev-Zel’dovich formula has been derived from a kinetic equation for the photon distribution function taking into account the Compton scattering by electrons: the Kompaneets equation (Kompaneets 1957; Weymann 1965). The original Kompaneets equation has been derived with a nonrelativistic approximation for the electron. However, recent X-ray observations have revealed the existence of many high-temperature galaxy clusters (David et al. 1993; Arnaud et al. 1994; Markevitch et al. 1994; Markevitch et al. 1996; Holzapfel et al. 1997; Mushotzky & Scharf 1997; Markevitch 1998). In particular, Tucker et al. (1998) reported the discovery of a galaxy cluster with the electron temperature $`k_BT_e=17.4\pm 2.5`$ keV. Rephaeli and his collaborator (Rephaeli 1995; Rephaeli & Yankovitch 1997) have emphasized the need to take into account the relativistic corrections to the Sunyaev-Zel’dovich effect for clusters of galaxies.
In recent years remarkable progress has been achieved in the theoretical studies of the relativistic corrections to the Sunyaev-Zel’dovich effects for clusters of galaxies. Stebbins (1997) generalized the Kompaneets equation. Itoh, Kohyama, & Nozawa (1998) have adopted a relativistically covariant formalism to describe the Compton scattering process (Berestetskii, Lifshitz, & Pitaevskii 1982; Buchler & Yueh 1976), thereby obtaining higher-order relativistic corrections to the thermal Sunyaev-Zel’dovich effect in the form of the Fokker-Planck expansion. In their derivation, the scheme to conserve the photon number at every stage of the expansion which has been proposed by Challinor & Lasenby (1998) played an essential role. The results of Challinor & Lasenby (1998) are in agreement with those of Itoh, Kohyama, & Nozawa (1998). The latter results include higher-order expansions. Itoh, Kohyama, & Nozawa (1998) have also calculated the collision integral of the Boltzmann equation numerically and have compared the results with those obtained by the Fokker-Planck expansion method. They have confirmed that the Fokker-Planck expansion method gives an excellent result for $`k_BT_e15`$keV, where $`T_e`$ is the electron temperature. For $`k_BT_e15`$keV, however, the Fokker-Planck expansion results show nonnegligible deviations from the results obtained by the numerical integration of the collision term of the Boltzmann equation.
Nozawa, Itoh, & Kohyama (1998b) have extended their method to the case where the galaxy cluster is moving with a peculiar velocity with respect to CMB. They have thereby obtained the relativistic corrections to the kinematical Sunyaev-Zel’dovich effect. Challinor & Lasenby (1999) have confirmed the correctness of the result obtained by Nozawa, Itoh, & Kohyama (1998b). Sazonov & Sunyaev (1998a, b) have calculated the kinematical Sunyaev-Zel’dovich effect by a different method. Their results are in agreement with those of Nozawa, Itoh, & Kohyama (1998b). The latter authors have given the results of the higher-order expansions.
Itoh, Nozawa, & Kohyama (2000) have also applied their method to the calculation of the relativistic corrections to the polarization Sunyaev-Zel’dovich effect (Sunyaev & Zel’dovich 1980b, 1981). They have thereby confirmed the result of Challinor, Ford, & Lasenby (1999) which has been obtained with a completely different method. Recent works on the polarization Sunyaev-Zel’dovich effect include Audit & Simons (1999), Hansen & Lilje (1999), and Sazonov & Sunyaev (1999).
As stated above, Itoh, Kohyama, & Nozawa (1998) have carried out the numerical integration of the collision term of the Boltzmann equation. This method produces the exact results without the power series expansion approximation. Sazonov & Sunyaev (1998a, b) have reported the results of the Monte Carlo calculations on the relativistic corrections to the Sunyaev-Zel’dovich effect. In Sazonov & Sunyaev (1998b), a numerical table which summarizes the results of the Monte Carlo calculations has been presented. This table is of great value when one wishes to calculate the relativistic corrections to the Sunyaev-Zel’dovich effect for galaxy clusters of extremely high temperatures. Accurate analytic fitting formulae would be still more convenient to use for the observers who wish to analyze the galaxy clusters with extremely high temperatures. For this purpose, Nozawa et al. (2000) have presented an accurate analytical fitting formula of 0.1% accuracy for the numerical results for the relativistic corrections to the thermal Sunyaev-Zel’dovich effect for clusters of galaxies. For the analyses of the galaxy clusters with extremely high temperatures, the results of the calculation of the relativistic thermal bremsstrahlung Gaunt factor (Nozawa, Itoh, & Kohyama 1998a) and their accurate analytic fitting formulae (Itoh et al. 2000) will be useful.
In this series of papers devoted to the study of the relativisitc corrections to the Sunyaev-Zel’dovich effect for clusters of galaxies, we have so far restricted ourselves to the case of single Compton scattering. This is justified because the optical depth for the Compton scattering of the CMB photon inside the galaxy clusters is generally about $`10^2`$ or smaller (Birkinshaw 1999). Nevertheless, it would be desirable to evaluate the effects of the multiple Compton scattering of the CMB photon inside the galaxy clusters accurately, as we have already developed the method to calculate the relativistic corrections to the Sunyaev-Zel’dovich effect for the galaxy clusters with high accuracy. The multiple scattering effects have been already considered by many authors (see Birkinshaw 1999 for references). In this paper we wish to evaluate the multiple scattering effects in the same theoretical framework of this series of papers.
The present paper is organized as follows. In $`\mathrm{\S }`$ 2 we give the method of the calculation and the results. In $`\mathrm{\S }`$ 3 we give discussion of the results and concluding remarks.
## 2 MULTIPLE SCATTERING CONTRIBUTION
In the present paper, we would like to derive the analytic expression for the multiple scattering contribution in the Sunyaev-Zeldovich effect for clusters of galaxies. As a reference system, we choose the system that is fixed to the center of mass of the galaxy cluster. The galaxy cluster is assumed to be fixed to the cosmic microwave background (CMB). Following Itoh, Kohyama & Nozawa (1998), we start with the Fokker-Plank expansion for the time evolution equation of the CMB photon distribution function $`n(\omega )`$:
$`{\displaystyle \frac{n(\omega )}{t}}`$ $`=`$ $`2\left[{\displaystyle \frac{n}{x}}+n(1+n)\right]I_1`$ (2.1)
$`+`$ $`2\left[{\displaystyle \frac{^2n}{x^2}}+2(1+n){\displaystyle \frac{n}{x}}+n(1+n)\right]I_2`$
$`+`$ $`2\left[{\displaystyle \frac{^3n}{x^3}}+3(1+n){\displaystyle \frac{^2n}{x^2}}+3(1+n){\displaystyle \frac{n}{x}}+n(1+n)\right]I_3`$
$`+`$ $`2\left[{\displaystyle \frac{^4n}{x^4}}+4(1+n){\displaystyle \frac{^3n}{x^3}}+6(1+n){\displaystyle \frac{^2n}{x^2}}+4(1+n){\displaystyle \frac{n}{x}}+n(1+n)\right]I_4`$
$`+`$ $`\mathrm{},`$
where
$`x`$ $``$ $`{\displaystyle \frac{\omega }{k_BT_e}},`$ (2.2)
$`\mathrm{\Delta }x`$ $``$ $`{\displaystyle \frac{\omega ^{}\omega }{k_BT_e}},`$ (2.3)
$`I_k`$ $``$ $`{\displaystyle \frac{1}{k!}}{\displaystyle \frac{d^3p}{(2\pi )^3}d^3p^{}d^3k^{}Wf(E)(\mathrm{\Delta }x)^k}.`$ (2.4)
In equation (2.4), $`W`$ is the transition probability of the Compton scattering, $`f(E)`$ is the relativistic Maxwellian distribution function for electrons with temperature $`T_e`$. We have integrated equation (2.4) analytically with power series expansions of the integrand. The expansion parameter is
$$\theta _e\frac{k_BT_e}{mc^2}.$$
(2.5)
The explicit forms for $`I_k`$ are given in Itoh, Kohyama & Nozawa (1998).
We first assume the initial photon distribution of the CMB radiation to be Planckian with temperature $`T_0`$:
$`n(X)`$ $`=`$ $`n_0(X){\displaystyle \frac{1}{e^X1}},`$ (2.6)
where
$`X`$ $``$ $`{\displaystyle \frac{\omega }{k_BT_0}}.`$ (2.7)
Assuming $`T_0/T_e1`$, one obtains the following expression for the fractional distortion of the photon spectrum derived by Itoh, Kohyama & Nozawa (1998):
$`{\displaystyle \frac{\mathrm{\Delta }n(X)}{n_0(X)}}`$ $`=`$ $`{\displaystyle \frac{y\theta _eXe^X}{e^X1}}\left[Y_0+\theta _eY_1+\theta _e^2Y_2+\theta _e^3Y_3+\theta _e^4Y_4\right],`$ (2.8)
$`y`$ $``$ $`\sigma _T{\displaystyle _0^{\mathrm{}}}𝑑\mathrm{}_1N_e(\mathrm{}_1),`$ (2.9)
where $`\sigma _T`$ is the Thomson scattering cross section, $`N_e`$ is the electron number density, and the integral is over the photon path length in the cluster. The explicit forms for $`Y_0,Y_1,Y_2,Y_3`$ and $`Y_4`$ are given in Itoh, Kohyama & Nozawa (1998). Equation (2.8) is the single scattering contribution, i.e. the first-order term in $`y`$. If the cluster of galaxies is optically thin, i.e. $`y1`$, the single scattering approximation is a good approximation. In fact, the approximation is valid for most of the clusters. However, it is extremely important to calculate the next-order contribution in order to obtain more accurate theoretical prediction for the future observation of the Sunyaev-Zeldovich effect for clusters of galaxies.
We now calculate the multiple scattering contribution. Since $`y1`$ is realized for most of clusters of galaxies, the second-order contribution is considered to be sufficient. We now assume that the initial photon distribution has a first-order perturbation. Namely,
$`n(X)`$ $`=`$ $`n_1(X)n_0(X)+\mathrm{\Delta }n(X),`$ (2.10)
$`=n_0(X)\left\{1+{\displaystyle \frac{\mathrm{\Delta }n(X)}{n_0(X)}}\right\},`$
where the second term in equation (2.10) is given by equation (2.8). Inserting equation (2.10) into RHS of equation (2.1), and performing the standard calculation, we obtain the following expression for the fractional distortion of the photon distribution function including the second-order contribution:
$`{\displaystyle \frac{\mathrm{\Delta }n(X)}{n_0(X)}}`$ $`=`$ $`{\displaystyle \frac{y\theta _eXe^X}{e^X1}}\left[Y_0+\theta _eY_1+\theta _e^2Y_2+\theta _e^3Y_3+\theta _e^4Y_4\right],`$ (2.11)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{y^2\theta _e^2Xe^X}{e^X1}}\left[Z_0+\theta _eZ_1+\theta _e^2Z_2\right],`$
$`Z_0`$ $`=`$ $`16+34\stackrel{~}{X}12\stackrel{~}{X}^2+\stackrel{~}{X}^3+\stackrel{~}{S}^2\left(6+2\stackrel{~}{X}\right),`$ (2.12)
$`Z_1`$ $`=`$ $`80+590\stackrel{~}{X}{\displaystyle \frac{3492}{5}}\stackrel{~}{X}^2+{\displaystyle \frac{1271}{5}}\stackrel{~}{X}^3{\displaystyle \frac{168}{5}}\stackrel{~}{X}^4+{\displaystyle \frac{7}{5}}\stackrel{~}{X}^5`$ (2.13)
$`+`$ $`\stackrel{~}{S}^2\left({\displaystyle \frac{1746}{5}}+{\displaystyle \frac{2542}{5}}\stackrel{~}{X}{\displaystyle \frac{924}{5}}\stackrel{~}{X}^2+{\displaystyle \frac{91}{5}}\stackrel{~}{X}^3\right)`$
$`+`$ $`\stackrel{~}{S}^4\left({\displaystyle \frac{168}{5}}+{\displaystyle \frac{119}{10}}\stackrel{~}{X}\right)`$
$`Z_2`$ $`=`$ $`160+4792\stackrel{~}{X}{\displaystyle \frac{357144}{25}}\stackrel{~}{X}^2+{\displaystyle \frac{312912}{25}}\stackrel{~}{X}^3{\displaystyle \frac{110196}{25}}\stackrel{~}{X}^4`$ (2.14)
$`+{\displaystyle \frac{34873}{50}}\stackrel{~}{X}^5{\displaystyle \frac{734}{15}}\stackrel{~}{X}^6+{\displaystyle \frac{367}{300}}\stackrel{~}{X}^7`$
$`+`$ $`\stackrel{~}{S}^2({\displaystyle \frac{178572}{25}}+{\displaystyle \frac{625824}{25}}\stackrel{~}{X}{\displaystyle \frac{606078}{25}}\stackrel{~}{X}^2+{\displaystyle \frac{453349}{50}}\stackrel{~}{X}^3`$
$`{\displaystyle \frac{20919}{15}}\stackrel{~}{X}^4+{\displaystyle \frac{367}{5}}\stackrel{~}{X}^5)`$
$`+`$ $`\stackrel{~}{S}^4\left({\displaystyle \frac{110196}{25}}+{\displaystyle \frac{592841}{100}}\stackrel{~}{X}2202\stackrel{~}{X}^2+{\displaystyle \frac{5872}{25}}\stackrel{~}{X}^3\right)`$
$`+`$ $`\stackrel{~}{S}^6\left({\displaystyle \frac{6239}{30}}+{\displaystyle \frac{11377}{150}}\stackrel{~}{X}\right),`$
where
$`\stackrel{~}{X}`$ $``$ $`X\mathrm{coth}\left({\displaystyle \frac{X}{2}}\right),`$ (2.15)
$`\stackrel{~}{S}`$ $``$ $`{\displaystyle \frac{X}{\mathrm{sinh}\left({\displaystyle \frac{X}{2}}\right)}}.`$ (2.16)
In equation (2.11), the first term corresponds to the first-odrder contribution and the second term corresponds to the second-order contribution of the multiple scattering. In deriving equation (2.11), we have used the following identity relation for $`y`$:
$`\sigma _T{\displaystyle _0^{\mathrm{}}}𝑑\mathrm{}_1N_e(\mathrm{}_1)\sigma _T{\displaystyle _0^\mathrm{}_1}𝑑\mathrm{}_2N_e(\mathrm{}_2)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\sigma _T{\displaystyle _0^{\mathrm{}}}𝑑\mathrm{}_1N_e(\mathrm{}_1)\right)^2`$ (2.17)
$`=`$ $`{\displaystyle \frac{1}{2}}y^2.`$
We have also neglected terms higher than $`O(\theta _e^4)`$ in the $`y^2`$ contributions in equation (2.11). It is important to note that equation (2.11) satisfies the photon number conservation. With equation (2.11), we define the distortion of the spectral intensity as follows:
$`\mathrm{\Delta }I`$ $`=`$ $`{\displaystyle \frac{X^3}{e^X1}}{\displaystyle \frac{\mathrm{\Delta }n(X)}{n_0(X)}}=\mathrm{\Delta }I_1+\mathrm{\Delta }I_2.`$ (2.18)
The first term $`\mathrm{\Delta }I_1`$ contains a factor $`y`$, whereas the second term $`\mathrm{\Delta }I_2`$ contains a factor $`y^2`$. In Figure 1 we show $`\mathrm{\Delta }I_2/y^2`$ as a function of $`X`$ for the case $`k_BT_e`$ = 10keV. It is clear that the magnitude of $`\mathrm{\Delta }I_2/y^2`$ has a maximum value at $`X5`$ for $`k_BT_e`$=10keV.
In order to estimate the relative importance of the multiple scattering contribution, we now define the following ratio:
$`\mathrm{\Gamma }`$ $``$ $`{\displaystyle \frac{\mathrm{\Delta }I_2/y^2}{\mathrm{\Delta }I_1/y}}.`$ (2.19)
In Figure 2 we show $`\mathrm{\Gamma }`$ for $`X=5`$ as a function of the electron temperature $`T_e`$. It is clear that $`\mathrm{\Gamma }`$ increases with a negative sign as the temperature of the cluster of galaxies increases, i.e., $`\mathrm{\Gamma }0.3`$ at $`k_BT_e=15`$keV. However, the multiple scattering contribution is small because of a further factor $`y`$. Namely, for the cluster of galaxies of $`k_BT_e=15`$keV, we have
$`{\displaystyle \frac{\mathrm{\Delta }I_2}{\mathrm{\Delta }I_1}}`$ $`=`$ $`y\mathrm{\Gamma }0.3y0.3\%,`$ (2.20)
where we used a typical value $`y0.01`$ of the galaxy clusters. Therefore the maximum effect of the multiple scattering contribution is $`0.3\%`$ of the single scattering contribution for the observed high-temperature galaxy clusters.
In the Rayleigh–Jeans limit where $`X0`$, equation (2.11) is further simplified:
$`{\displaystyle \frac{\mathrm{\Delta }n(X)}{n_0(X)}}`$ $`=`$ $`2y\theta _e\left(1{\displaystyle \frac{17}{10}}\theta _e+{\displaystyle \frac{123}{40}}\theta _e^2{\displaystyle \frac{1989}{280}}\theta _e^3+{\displaystyle \frac{14403}{640}}\theta _e^4\right)`$ (2.21)
$`+2y^2\theta _e^2\left(1{\displaystyle \frac{17}{5}}\theta _e+{\displaystyle \frac{226}{25}}\theta _e^2\right).`$
With equation (2.21) we have the multiple scattering contribution for $`k_BT_e=`$15keV and $`y=0.01`$ as follows:
$`{\displaystyle \frac{\mathrm{\Delta }I_2}{\mathrm{\Delta }I_1}}`$ $``$ $`y\theta _e0.03\%.`$ (2.22)
In the Rayleigh–Jeans region the multiple scattering contribution is safely neglected.
## 3 DISCUSSION AND CONCLUDING REMARKS
From the results presented in the previous section it is clear that the multiple scattering contribution $`\mathrm{\Delta }I_2`$ is very small compared with the single scattering contribution $`\mathrm{\Delta }I_1`$. For high-temperature galaxy clusters of $`k_BT_e15`$keV, we obtain the ratio $`\mathrm{\Delta }I_2/\mathrm{\Delta }I_10.3\%`$ at $`X=5`$. In the Rayleigh–Jeans region we have $`\mathrm{\Delta }I_2/\mathrm{\Delta }I_10.03\%`$. Therefore it is concluded that the multiple scattering contribution to the thermal Sunyaev-Zel’dovich effect for galaxy clusters can be safely neglected. The reader is therefore referred to the previous four papers in this series of papers which deal with the single scattering contribution in detail.
This work is financially supported in part by the Grant-in-Aid of Japanese Ministry of Education, Science, Sports, and Culture under the contract #10640289.
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# 1 Introduction
## 1 Introduction
The purpose of this paper is to count BPS monopoles and dyons in supersymmetric Yang-Mills theories. Such computations have been performed by many authors using moduli space dynamics of monopoles . Recently , however, it was realized that the monopole dynamics in a generic vacuum is qualitatively different from the old moduli space dynamics of Manton employed in most such endeavors, where the low energy dynamics of monopoles were considered when only one adjoint Higgs field is turned on, while supersymmetric Yang-Mills field theories come with 2 or 6 such scalars. This restriction disallows static interaction between monopoles , so that all interaction comes from nontrivial coefficients of kinetic terms. In a generic vacuum with more than one adjoint Higgs turned on, monopoles of the same type still have no static force among them, but dynamics of monopoles of distinct type could have a static potential. In this article, we solve various index problems as a first step towards counting all BPS states.
In the old monopole dynamics, the Lagrangian one finds is a pure sigma model with the moduli space as the target manifold. Classically, one solves for geodesic trajectories to find classical orbits of monopoles. Quantum mechanically, the Hamiltonian is a square of a supercharge which can be regarded as a Dirac operator acting either on the spinor bundle or on the Clifford bundle over the moduli space,
$$Qi\gamma ^m_m.$$
(1)
Supersymmetric bound states, for example, would be found as normalizable spinors or forms on the moduli space that are also zero modes of this Dirac operator.
The new supersymmetric low energy dynamics is obtained by augmenting old moduli space dynamics with a set of supersymmetric potential terms, and was written explicitly in Ref.. In principle, the question of BPS states must be reconsidered in the new dynamics. (One example of states that cannot be probed in the old formalism is the now well-known 1/4 BPS states .) In this new setting, the supercharges of the low energy dynamics will again be interpreted as Dirac operators on the moduli space, which is now twisted by some triholomorphic Killing vector field, say $`K`$,
$$Q\gamma ^m(i_mK_m).$$
(2)
The corresponding bosonic potential is precisely half the squared norm of $`K`$ . We count the number of normalizable states annihilated by a Dirac operator, weighted by $`\pm 1`$ for the chiral and the antichiral states respectively. The resulting integer is the index of the Dirac operator. In fact this index would be infinite or ill defined without further restriction of the domain of the Dirac operator. We restrict the problem to each charge eigensector before counting the zero modes of the Dirac operator. Thus, effectively, we will be computing an equivariant index with $`L^2`$ condition. This gives information on the existence and the degeneracy of dyonic bound states for each electric charge sector.
Understanding monopole dynamics in generic vacua is particularly significant in the context of $`N=2`$ Yang-Mills theories, because many (1/2) BPS dyons exist only when both adjoint Higgs are turned on. By ignoring the potential term in the low energy dynamics, one would in effect be searching for a bound state in the vacua where the state cannot exist as a supersymmetric one-particle state. In the language of Seiberg-Witten , one would be looking for it on one side of a marginal stability domain wall in the vacuum moduli space, while the bound state in question exists only on the other side of the domain wall. For example, most dyons that become massless at special hypersurfaces in the Seiberg-Witten moduli space, are of this type and cannot be probed by old moduli space dynamics.
In Section 2, we review the supersymmetric quantum mechanics with potential for the case of 4 and 8 supercharges. We isolate various involutions, with respect to which the index is defined. Section 3 introduces the explicit form of the moduli space metric that governs the dynamics of BPS monopoles. The only known moduli space for arbitrarily many monopoles is the case of all distinct monopoles, and this is the case for which we will compute the index explicitly. Section 4 recalls a recent explicit computation of two-monopole bound states in $`SU(3)`$ theories and presents the resulting value of the indices. In section 5, we finally delve into the computation of the index by using a Fredholm deformation of the Dirac operator in question. The computation reduces to that of a superharmonic oscillator in 4 dimensions, whose index is computed explicitly. Section 6 translates the results to degeneracies of various dyonic and purely magnetic bound states and checks its consistency with anticipated nonperturbative physics. We close with a summary.
## 2 Supersymmetric Sigma Model with Potential
In this section, we briefly review the supersymmetric sigma-model quantum mechanics with potentials. These quantum mechanics have been introduced as the low energy dynamics of monopoles in pure $`N=2`$ Yang-Mills field theory and in $`N=4`$ Yang-Mills field theory . They also appeared in other systems such as the dynamics of instanton solitons .
### 2.1 Quantum Mechanics with 4 Real SUSY
The SUSY dynamics we consider is a sigma model with potential, whose Lagrangian is written compactly as
$$=\frac{1}{2}\left(g_{mn}\dot{z}^m\dot{z}^n+ig_{mn}\lambda ^mD_t\lambda ^ng^{mn}G_mG_ni_mG_n\lambda ^m\lambda ^n\right),$$
(3)
where $`D_t\lambda ^m=\dot{\lambda }^m+\mathrm{\Gamma }_{np}^m\dot{z}^n\lambda ^p`$. The target manifold must be hyperKähler, which means that there are three covariantly constant complex structures $`J^{(s)}`$ satisfying quaternionic algebra,
$$J^{(s)}J^{(t)}=\delta ^{st}+ϵ^{stu}J^{(u)},$$
(4)
and the Killing vector field $`G`$ should be triholomorphic.
$`_Gg=0,_GJ^{(s)}=0.`$ (5)
Introducing vielbein $`e_m^E`$ and defining $`\lambda ^E=\lambda ^me_m^E`$ which commute with all bosonic variables, the canonical commutators are
$`[z^m,p_n]`$ $`=`$ $`i\delta _n^m,`$
$`\{\lambda ^E,\lambda ^F\}`$ $`=`$ $`\delta ^{EF}.`$ (6)
We can realize this algebra on spinors on the moduli space by letting $`\lambda ^E=\gamma ^E/\sqrt{2}`$, where $`\gamma ^E`$ are gamma matrices. (Since the moduli space is hyperKähler an equivalent quantization is obtained using holomorphic differential forms.) The supercovariant momentum operator, defined by
$$\pi _m=p_m\frac{i}{4}\omega _{mEF}[\lambda ^E,\lambda ^F],$$
(7)
where $`\omega _{mE}^F`$ is the spin connection, then becomes the covariant derivative acting on spinors $`\pi _m=i_m`$. Note that
$`[\pi _m,\lambda ^n]`$ $`=`$ $`i\mathrm{\Gamma }_{mp}^n\lambda ^p,`$
$`[\pi _m,\pi _n]`$ $`=`$ $`{\displaystyle \frac{1}{2}}R_{mnpq}\lambda ^p\lambda ^q.`$ (8)
The supersymmetry charges take the form
$`Q`$ $`=`$ $`\lambda ^m(\pi _mG_m),`$
$`Q^{(s)}`$ $`=`$ $`\lambda ^mJ_m^{(s)n}(\pi _nG_n),`$ (9)
which obey
$`\{Q,Q\}`$ $`=`$ $`2(𝒵),`$
$`\{Q^{(s)},Q^{(t)}\}`$ $`=`$ $`2\delta _{st}(𝒵),`$
$`\{Q,Q^{(s)}\}`$ $`=`$ $`0.`$ (10)
Here the Hamiltonian $``$ and the central charge $`𝒵`$ are given by
$`={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{\sqrt{g}}}\pi _m\sqrt{g}g^{mn}\pi _n+G_mG^m+i_mG_n\lambda ^m\lambda ^n\right),`$ (11)
$`𝒵=G^m\pi _m{\displaystyle \frac{i}{2}}(_mG_n)\lambda ^m\lambda ^n.`$ (12)
Note that the operator $`i𝒵`$ is the Lie derivative $`_G`$ acting on spinors (see e.g., )
$$_GG^m_m+\frac{1}{8}_mG_n[\gamma ^m,\gamma ^n].$$
(13)
The SUSY quantum mechanics comes with a natural $`Z_2`$ grading defined by the operator,
$$\tau _2=2^{1/2}\lambda ^E=\gamma ^E,$$
(14)
which anticommutes with the Dirac operator,
$$D=\sqrt{2}Q=\gamma ^m(i_mG_m).$$
(15)
This pair defines the Witten index that counts the difference between the number of bosonic states and the number of fermionic states annihilated by the supercharge. In fact, the index is defined in each superselection sector with fixed $`𝒵`$, and effectively counts the difference in the numbers of BPS states of given central charges. The index will be denoted collectively by $`_2`$. See Section 5 for detailed computation of $`_2`$.
### 2.2 Quantum Mechanics with 4 Complex SUSY
When the number of supercharges and the number of fermions double, we obtain the following form of sigma model with potential,
$``$ $`=`$ $`{\displaystyle \frac{1}{2}}(g_{mn}\dot{z}^m\dot{z}^n+ig_{mn}\overline{\psi }^m\gamma ^0D_t\psi ^n+{\displaystyle \frac{1}{6}}R_{mnpq}\overline{\psi }^m\psi ^p\overline{\psi }^n\psi ^q`$ (16)
$`g_{mn}G_I^mG_I^ni_mG_{In}\overline{\psi }^m(\mathrm{\Omega }^I\psi )^n),`$
where $`\psi ^m`$ is a two component Majorana spinor, $`\gamma ^0=\sigma _2,\gamma ^1=i\sigma _1,\gamma ^2=i\sigma _3`$, $`\overline{\psi }=\psi ^T\gamma ^0`$. The operator $`\mathrm{\Omega }_I`$’s are defined respectively by $`\mathrm{\Omega }_4=\delta _n^m\gamma _{\alpha \beta }^1`$, $`\mathrm{\Omega }_5=\delta _n^m\gamma _{\alpha \beta }^2`$ and $`\mathrm{\Omega }_s=iJ_n^{(s)m}\delta _{\alpha \beta }`$ for $`s=1,2,3`$. The supersymmetry algebra again requires the manifold to be hyperKähler. As in the previous subsection, the $`G^I`$’s must be triholomorphic Killing vector fields.
When quantized, the spinors $`\psi ^E=e_m^E\psi ^m`$ with vielbein $`e_m^E`$, commute with all the bosonic dynamical variables, especially with $`p`$’s that are canonical momenta of the coordinate $`z`$’s. The remaining fundamental commutation relations are
$`[z^m,p_n]=i\delta _n^m,`$
$`\{\psi _\alpha ^E,\psi _\beta ^F\}=\delta ^{EF}\delta _{\alpha \beta }.`$ (17)
Define supercovariant momenta by
$`\pi _mp_m{\displaystyle \frac{i}{2}}\omega _{EFm}\overline{\psi }^E\gamma ^0\psi ^F,`$ (18)
where $`\omega _{EFm}`$ is the spin connection. The N=4 SUSY generators in real spinors can be written as,
$`Q_\alpha =\psi _\alpha ^m\pi _m(\gamma ^0\mathrm{\Omega }^I\psi )^mG_m^I,`$ (19)
$`Q_\alpha ^{(s)}=(J^{(s)}\psi )_\alpha ^m\pi _m(\gamma ^0J^{(s)}\mathrm{\Omega }^I\psi )^mG_m^I.`$ (20)
These charges satisfy the $`N=4`$ complex superalgebra:
$`\{Q_\alpha ,Q_\beta \}=\{Q_\alpha ^{(s)},Q_\beta ^{(s)}\}=2\delta _{\alpha \beta }2(\gamma ^0\gamma ^1)_{\alpha \beta }𝒵_42(\gamma ^0\gamma ^2)_{\alpha \beta }𝒵_5,`$ (21)
$`\{Q_\alpha ,Q_\beta ^{(s)}\}=2\gamma _{\alpha \beta }^0𝒵_s,\{Q_\alpha ^{(1)},Q_\beta ^{(2)}\}=2\gamma _{\alpha \beta }^0𝒵_3,`$ (22)
$`\{Q_\alpha ^{(2)},Q_\beta ^{(3)}\}=2\gamma _{\alpha \beta }^0𝒵_1,\{Q_\alpha ^{(3)},Q_\beta ^{(1)}\}=2\gamma _{\alpha \beta }^0𝒵_2,`$ (23)
where $``$ is the Hamiltonian, and the $`𝒵_I`$’s are central charges,
$`𝒵_I=G_I^m\pi _m{\displaystyle \frac{i}{2}}_mG_n^I\overline{\psi }^m\gamma ^0\psi ^n.`$ (24)
The sigma-model without the potential possesses an $`SO(5)`$ R-symmetry which is explicitly broken by the $`G^I`$’s. The $`G^I`$’s transform as 5 of $`SO(5)_R`$.
The complex form of the supercharges is often useful. To this end, we introduce $`\phi ^m\frac{1}{\sqrt{2}}(\psi _1^mi\psi _2^m)`$ and define $`Q\frac{1}{\sqrt{2}}(Q_1iQ_2)`$. The supercharges in (19) can be rewritten as
$`Q=\phi ^m\pi _m\phi ^m(G_m^4iG_m^5)i{\displaystyle \underset{s=1}{\overset{3}{}}}G_m^s(J^{(s)}\phi )^m,`$ (25)
$`Q^{}=\phi ^m\pi _m\phi ^m(G_m^4+iG_m^5)+i{\displaystyle \underset{s=1}{\overset{3}{}}}G_m^s(J^{(s)}\phi ^{})^m.`$ (26)
The charges $`Q^{(s)}`$ and $`Q_{}^{(s)}{}_{}{}^{}`$ are analogously defined from (20). The positive definite nature of the Hamiltonian can be seen easily in the anticommutator
$$\{Q,Q^{}\}=\{Q^{(s)},Q_{}^{(s)}{}_{}{}^{}\}=2,$$
(27)
while the central charges appear in other parts of the superalgebra. For instance, we have
$`\{Q,Q\}=𝒵_4+i𝒵_5,`$
$`\{Q^{},Q^{}\}=𝒵_4i𝒵_5.`$ (28)
Once we adopt this complex notation, it is natural to introduce an equivalent geometrical notation. Defining the vacuum state $`|0`$ to be annihilated by $`\phi ^m`$’s, and using the 1-1 correspondence,
$$(\phi ^{m_1}\phi ^{m_2}\mathrm{}\phi ^{m_k})|0dz^{m_1}dz^{m_2}\mathrm{}dz^{m_k},$$
(29)
we can reinterpret $`\phi ^m`$ as the exterior product with $`dz^m`$, and $`\phi _m^{}=g_{mn}\phi ^n`$ as the contraction with $`/z^m`$. The supercharge operators can be rewritten as,
$`Q`$ $`=`$ $`id\iota _{G^4iG^5}+i\iota _{J^{(1)}(G^1)}^{}+i\iota _{J^{(2)}(G^2)}^{}+i\iota _{J^{(3)}(G^3)}^{},`$
$`Q^{}`$ $`=`$ $`id^{}\iota _{G^4+iG^5}^{}i\iota _{J^{(1)}(G^1)}i\iota _{J^{(2)}(G^2)}i\iota _{J^{(3)}(G^3)},`$ (30)
where $`\iota _K`$ is the contraction with the vector field $`K`$, and its conjugate $`\iota _K^{}`$ is the exterior product by the 1-form obtained from $`K`$ by lowering its indices.
The SUSY quantum mechanics admit a canonical $`Z_2`$ grading, which in the geometrical notation of (29) is defined on $`k`$-forms by
$$\tau _4(1)^k.$$
(31)
or equivalently by
$$\tau _42\psi _1^E\psi _2^E=(\phi ^E\phi ^E\phi ^E\phi ^E).$$
(32)
The involution $`\tau _4`$ anticommutes with all supercharges and determines the usual Witten index, $`_4`$.
In some special limits, however, there could be an additional $`Z_2`$ grading. Suppose that we have only one nonzero $`G^I`$, say $`G^5`$. The operators
$$\tau _\pm (\sqrt{i}\phi ^E\pm \sqrt{i}\phi ^E).$$
(33)
anticommutes with the Dirac operators defined as,
$$D_\pm iQ\pm Q^{}=(i\phi ^m\pm \phi ^m)(\pi _mG_m^5)=d\iota _{G^5}\pm i(d^{}\iota _{G^5}^{}),$$
(34)
the square of which is
$$D_\pm ^2=\pm 2i(𝒵_5).$$
(35)
So the $`Z_2`$ gradings define an analog of the signature index for each choice of sign and for each charge-eigensector. A given state with nonzero $`𝒵_5`$ can be annihilated by one of $`D_\pm `$ at most, and in fact must break at least half of the supercharges. The corresponding indices will be denoted by $`_s^\pm `$.
In Section 5, we will compute both $`_4`$ and $`_s^\pm `$ in such a special limit with only one of five $`G^I`$’s present, which we can take to be $`G^5`$ without loss of generality. For $`_4`$, we may take any one of $`D_\pm `$ as the Dirac operator, since $`\tau _4`$ anticommutes with both. A standard index theorem will then allow us to deduce $`_4`$ in more general setting.
## 3 Moduli Spaces
Moduli space dynamics of monopoles decompose into the interacting relative part and the non-interacting “center of mass” part. The latter corresponds to a 4-dimensional flat metric of the form,
$$g_{cm}=Ad\stackrel{}{X}^2+Bd\xi _T^2,$$
(36)
where $`\stackrel{}{X}`$ is a three-vector. Since we are interested in establishing existence of bound states, this part of the dynamics will be ignored for the most part.
The free center-of-mass sector generates two kinds of quantum numbers, nevertheless. One is the overall, conserved $`U(1)`$ charge, and the other is a supermultiplet structure generated by the fermionic partners of $`\stackrel{}{X}`$ and $`\xi _T`$. The resulting degeneracies, 4 and 16 for $`N=4`$ real and complex supersymmetric quantum mechanics respectively, correspond to the smallest possible BPS multiplet of the underlying SUSY Yang-Mills field theories with 8 and 16 supercharges, respectively.
### 3.1 Distinct Monopoles
A simple case of this dynamics involves a collection of distinct monopoles in $`SU(n)`$ gauge theories. The interacting part of the moduli space metric is a simple generalization of four-dimensional Taub-NUT metric . Without loss of generality, consider a collection of $`k+1`$ distinct monopoles, whose magnetic charges are given by an irreducible (sub)set of simple roots, $`\beta _1,\mathrm{},\beta _{k+1}`$. The simple roots satisfy relations $`\beta _a^2=1`$, $`\beta _a\beta _{a+1}=1/2`$, and $`\beta _a\beta _{a+b}=0`$ for $`b>1`$. The relative part of the corresponding metric is
$$g=C_{ab}d\stackrel{}{r}_ad\stackrel{}{r}_b+\frac{4\pi ^2}{e^4}(C^1)_{ab}(d\psi _a+\mathrm{cos}\theta _ad\varphi _a)(d\psi _b+\mathrm{cos}\theta _bd\varphi _b),$$
(37)
where the matrix $`C`$ for the relative moduli space is<sup>1</sup><sup>1</sup>1The coupling constant $`e`$ will be assumed to be positive without loss of generality.
$$C_{ab}=\mu _{ab}+\frac{2\pi }{e^2}\delta _{ab}\frac{1}{r_a}.$$
(38)
The 3-vector $`\stackrel{}{r}_a`$ is the relative position between the $`a^{th}`$ and $`(a+1)^{th}`$ monopoles,
$$\stackrel{}{r}_a=\stackrel{}{x}_{a+1}\stackrel{}{x}_a,$$
(39)
while the angles $`\psi _a`$ of period $`4\pi `$ are related to the $`U(1)`$ phases of each monopole, $`\xi _a`$’s (of period $`2\pi `$), by
$`2{\displaystyle \frac{}{\psi _a}}`$ $`=`$ $`{\displaystyle \frac{}{\xi _{a+1}}}{\displaystyle \frac{}{\xi _a}}`$
$`\left({\displaystyle \underset{a=1}{\overset{k+1}{}}}m_a\right){\displaystyle \frac{}{\xi _T}}`$ $`=`$ $`{\displaystyle \underset{a=1}{\overset{k+1}{}}}m_a{\displaystyle \frac{}{\xi _a}}.`$ (40)
where $`\xi _T`$ is a coordinate that appears in free center-of-mass part of the dynamics and $`m_a`$ is the mass of the $`a^{th}`$ monopole.
For a generic reduced mass matrix $`\mu `$, the triholomorphic Killing vector fields of this geometry are exhausted by
$$K_a=\frac{}{\psi _a},$$
(41)
so the vector fields $`G`$ and $`G^I`$ are linear combinations of $`K_a`$’s with constant coefficients;
$`G`$ $`=`$ $`e{\displaystyle \underset{c}{}}a_cK_c,`$
$`G^I`$ $`=`$ $`e{\displaystyle \underset{c}{}}a_c^IK_c.`$ (42)
The electric charges are measured by the charge operators,
$$i_{K_a},$$
(43)
whose (half-)integer eigenvalues will be denoted by $`q_a`$. In terms of the simple roots $`\beta _a`$, the electric charge of a dyonic state with charge $`q_a`$’s is
$$\begin{array}{c}e(+q_1+q_2+q_3+\mathrm{}+q_k+n/2)\beta _1+\hfill \\ e(q_1+q_2+q_3+\mathrm{}+q_k+n/2)\beta _2+\hfill \\ e(q_1q_2+q_3+\mathrm{}+q_k+n/2)\beta _3+\hfill \\ e(q_1q_2q_3+\mathrm{}+q_k+n/2)\beta _4+\hfill \\ \mathrm{}\hfill \\ e(q_1q_2q_3\mathrm{}q_k+n/2)\beta _{k+1}.\hfill \end{array}$$
(44)
where the integer $`n`$ comes from quantization of an overall $`U(1)`$ angle and should be even or odd when $`2_aq_a`$ is even or odd, respectively.
### 3.2 Unit Noncommutative Instanton
A simple deformation of the above moduli space appeared in another context recently, where one considers low energy dynamics of an instanton soliton in the 5-dimensional noncommutative $`U(k+1)`$ Yang-Mills theory . This happens because an instanton in $`S^1\times R^3`$ can be regarded as a collection of $`k+1`$ distinct monopoles of the underlying Yang-Mills theory . When we compactify the theory on a circle of radius $`R`$, the nontrivial part of the moduli space of a single instanton soliton is given by the metric
$$g=\frac{4\pi ^2R}{\stackrel{~}{e}^2}\left(\stackrel{~}{C}_{ab}d\stackrel{}{r}_ad\stackrel{}{r}_b+(\stackrel{~}{C}^1)_{ab}(d\psi _a+\mathrm{cos}\theta _ad\varphi _a)(d\psi _b+\mathrm{cos}\theta _bd\varphi _b)\right),$$
(45)
where $`\stackrel{~}{e}`$ is the 5-dimensional Yang-Mills coupling. The matrix $`\stackrel{~}{C}`$ for the relative moduli space is
$$\stackrel{~}{C}_{ab}=\nu _{ab}+\delta _{ab}\frac{1}{r_a}+\frac{1}{|\stackrel{}{r}_a2\pi \stackrel{}{\zeta }/R|},$$
(46)
where $`\stackrel{}{\zeta }`$ encodes the noncommutativity. The matrix $`\nu `$ is determined by the Wilson line along $`S^1`$ that breaks the gauge symmetry to $`U(1)^n`$. When the supersymmetry of the underlying field theory is maximal with 16 supercharges, the low energy dynamics of the instanton is given by our SUSY quantum mechanics with 4 complex supercharges. When the field theory comes with 8 supercharges, the instanton dynamics is described by the SUSY quantum mechanics with 4 real supercharges.
## 4 Bound States of Two Distinct Monopoles
For a pair of two distinct and interacting monopoles, the dynamics have been solved for supersymmetric ground states in each charge eigensector. The geometry is that of a Taub-NUT manifold which comes with a single triholomorphic Killing vector field $`K_1`$. Accordingly, there is only one conserved $`U(1)`$ charge, $`q_1`$, which labels superselection sectors.
In the pure $`N=2`$ Yang-Mills case, define $`\stackrel{~}{a}_14\pi ^2a_1/e^3\mu `$ where $`\mu `$ is the reduced mass and $`a_1`$ is defined by $`G=ea_1K_1`$. The normalizable wavefunctions had been constructed by Pope in another context , and the number of dyonic bound states of charge $`q_1`$ was found to be
$$2|q_1|,$$
(47)
if $`0<q_1<\stackrel{~}{a}`$ or $`\stackrel{~}{a}_1<q_1<0`$, and
$$0,$$
(48)
otherwise. For each $`q_1`$, the solutions belong to the same chirality spinors, and thus contribute to the Witten index equally. Thus the Witten index $`_2`$ in each charge eigensectors are
$$_2=\left\{\begin{array}{cc}2|q_1|& 0<|q_1|<|\stackrel{~}{a}_1|\text{and}0<q\stackrel{~}{a}_1\hfill \\ & \\ 0& \mathrm{otherwise}\hfill \end{array}\right\}.$$
(49)
For a pair of distinct monopoles in $`N=4`$ Yang-Mills , the five $`G^I`$’s must be proportional to the single triholomorphic vector field $`K_1=/\psi _1`$. We may rotate them into a single triholomorphic vector field, say $`G^{I=5}=ea_1K_1`$, upon which we can define $`\stackrel{~}{a}_1`$ similarly as above, $`\stackrel{~}{a}_14\pi ^2a_1/e^3\mu `$. The degeneracy is found to be
$$1,$$
(50)
for purely magnetic state ($`q_1=0`$), while for dyons
$$8|q_1|,$$
(51)
when $`0<|q_1|<|\stackrel{~}{a}_1|`$, and zero otherwise. All solutions are self-dual differential forms, when we take the convention that the curvature tensor of the moduli space is self-dual.
When the central charge $`𝒵_5=ea_1q_1`$ of the state is positive (negative), the bound state is annihilated by $`D_+`$ ($`D_{}`$) only, while for $`𝒵_5=0`$, the state is annihilated by both. For given $`q_1`$, we find
$$_s^+=\left\{\begin{array}{cc}1& q_1=0\hfill \\ 8|q_1|& 0<|q_1|<|\stackrel{~}{a}_1|\mathrm{and}q_1a_1>0\hfill \\ 0& \mathrm{otherwise}\hfill \end{array}\right\},$$
(52)
and
$$_s^{}=\left\{\begin{array}{cc}1& q_1=0\hfill \\ 8|q_1|& 0<|q_1|<|\stackrel{~}{a}_1|\mathrm{and}q_1a_1<0\hfill \\ 0& \mathrm{otherwise}\hfill \end{array}\right\}.$$
(53)
The Witten index $`_4`$ counts the number of even forms minus the number of the odd forms. Of solutions with $`q_10`$, half are even and the other half are odd, so we find that
$$_4=\left\{\begin{array}{cc}1& q_1=0\hfill \\ & \\ 0& q_10\hfill \end{array}\right\}$$
(54)
regardless of $`a_1`$.
## 5 Index Computation
We would like to put a lower bound on the number of bound states in the above SUSY dynamics by computing indices. The index problems can be quite involved, given that the quantum mechanics involve many degrees of freedom with complicated interaction terms. However, the problem can be simplified by utilizing the invariance of the index under certain deformations. In this section we will use the invariance of the index under Fredholm deformation to simplify our index computations. Before proceeding with the computation, however, we need to restrict to the regime where a massgap exists.
### 5.1 Massgap
When restricted to specific charge eigensectors, the operators above may exhibit two drastically different behavior. For small charges, the sector has a massgap; the continuum part of the spectrum is bounded below by a positive gap. For large charges, the massgap disappears. This is the reason why there is an upper bound on the electric charge $`q_1`$ of bound states of two monopoles. In the two-body problems, the condition for the massgap to exist in a sector with electric charge $`q_1`$ is
$$|q_1|<\frac{4\pi ^2}{e^3}\frac{|a_1|}{\mu },$$
(55)
where the bosonic potential is generated by a single triholomorphic vector field $`G=ea_1K_1`$. When we consider many distinct monopoles, the condition for the massgap to exist is equally simple:
$$|q_c|<|\stackrel{~}{a}_c|,$$
(56)
where
$$\stackrel{~}{a}_c\frac{4\pi ^2}{e^3}\underset{b=1}{\overset{k}{}}(\mu ^1)_{cb}a_b,$$
(57)
with $`G=ea_cK_c`$. In the quantum mechanics with four complex supersymmetries, the same holds true provided that only one $`G^I`$, say, $`G^5=ea_cK_c`$ is turned on. We will compute the indices, $`_2`$, $`_4`$, $`_s^\pm `$ assuming that all of these conditions hold.<sup>2</sup><sup>2</sup>2If five $`G^I=e_ca_c^IK_c`$’s are involved, the massgap condition generalizes to
$$(q_c)^2<\underset{I=1}{\overset{5}{}}(\stackrel{~}{a}_c^I)^2,$$
(58) where $`\stackrel{~}{a}_c^I`$ are defined similarly as above for each $`G^I`$. However, the Indices $`_s^\pm `$ are not well-defined unless all $`G^I`$’s are proportional to each other. We will discuss such generic cases in Section 6.
### 5.2 Index Generalities
First we recall basic definitions.
###### Definition 5.1
A bounded linear operator $`L:E_1E_2`$ between two Hilbert spaces is Fredholm if there exists a bounded operator $`P:E_2E_1`$ such that $`PLI_1`$, and $`LPI_2`$ are compact operators. Here $`I_j`$ denotes the identity map on $`E_j`$.
The operator $`P`$ in the above definition is called a parametrix. We will be interested in the case where $`L`$ is a Dirac operator. In this case, although Dirac operators are unbounded on $`L_2`$, we may trivially make $`L`$ bounded by taking $`E_1`$ to be the closure of the domain of $`L`$ with respect to the norm (graph norm)
$$f_{graph}^2f^2+Lf^2,$$
where unsubscripted norms denote $`L_2`$ norms.
If $`L`$ is a Dirac operator on a compact manifold, then it is well known to be Fredholm. In the compact case, one takes, for example, $`P`$ to be the Green’s operator, $`𝒢`$ defined to be the unique operator satisfying:
(i) $`𝒢`$ annihilates the kernel of $`L^{}`$.
(ii) The range of $`𝒢`$ is orthogonal to the kernel of $`L`$.
(iii) $`L𝒢f=f`$ for $`f`$ the image of $`L`$.
Then $`𝒢`$ is bounded by $`1+\lambda _1^{1/2}`$, where $`\lambda _1`$ is the first nonzero eigenvalue of $`L^{}L.`$ Also, $`𝒢L=I_1\mathrm{\Pi }_1,`$ and $`L𝒢=I_2\mathrm{\Pi }_2`$, where $`\mathrm{\Pi }_1`$ and $`\mathrm{\Pi }_2`$ denote the orthogonal projections onto the kernels of $`L`$ and $`L^{}`$ respectively (and are finite rank and thus compact operators). Hence $`P=𝒢`$ satisfies all the conditions of the definition.
In the case of a Dirac operator on a noncompact manifold the preceding construction of a Greens operator may fail to yield Fredholmness for several reasons. The kernel of $`L`$ or $`L^{}`$ may fail to be finite dimensional, making one of the projections not a compact operator. Also, if there is no gap in the spectrum, $`𝒢`$ will fail to be bounded. These deficits are all avoided, however, under the assumption that the essential spectrum of $`L^{}L`$ is bounded away from zero. (We recall that the essential spectrum includes the continuous spectrum and any eigenvalue of infinite multiplicity.) Then the kernels are finite dimensional and $`𝒢`$ is again bounded by $`1+\lambda _1^{1/2}`$, where $`\lambda _1`$ is the smallest nonzero element of the spectrum of $`L^{}L`$. It is well known that the essential spectrum is bounded away from zero whenever $`L^{}L`$ has the form $`\mathrm{\Delta }+V`$, where for two positive constants $`c`$ and $`R`$, $`V`$ satisfies $`V(x)>c`$ for $`x`$ outside a fixed compact set. All the operators we consider in this paper have this form.
A basic result in index theory (eg p.122), is the following.
###### Proposition 5.2
Let $`L_t`$, $`t[0,1]`$ be a continuous family of Fredholm operators. Then $`index(L_0)=index(L_1)`$.
Thus one can sometimes deform an index problem to a more tractable index computation. To avoid potential confusion, we recall the notion of continuity assumed in the above proposition. $`L_t`$ is a continuous family of operators if for each $`s`$ and for every $`ϵ>0`$, there exists $`\delta >0`$ so that $`L_tFL_sF/F<ϵ`$ for all nonzero $`F`$ if $`|ts|<\delta `$. In particular, we note that we require the $`\delta `$ to be independent of $`F`$. Hence, for example, the super harmonic oscillator in one variable $`\psi _1\frac{d}{dx}+\psi _2x`$, cannot be continuously deformed to $`\psi _1\frac{d}{dx}+\psi _2`$ by scaling away the interaction term. If one uses the graph norm for $`\psi _1\frac{d}{dx}+\psi _2`$, then the oscillator is unbounded and hence clearly cannot be deformed to a bounded operator. If one instead uses the graph norm for the oscillator, it is easy to see that all frequencies give equivalent norms and by construction, multiplication by $`x`$ (as a map to $`E_2`$ is continuous in each of these norms. Hence the deformation
$$L_t=\psi _1\frac{d}{dx}+\psi _2((1t)x+t)$$
is continuous. The limit operator, however, is not Fredholm as a map from $`E_1E_2`$ even though it is easy to show that it is Fredholm if the oscillator graph norm is replaced by the $`\psi _1\frac{d}{dx}+\psi _2`$ graph norm.
In analyzing continuous families of operators $`L_t`$ it is often useful to utilize also families of parametrices $`P_t`$. If, however, we choose $`P_t`$ to be the Greens operator $`𝒢_t`$ of $`L_t`$ then we will be plagued by the possibility that if eigenvalues converge to zero, $`𝒢_t`$ will become unbounded. Hence, for no other reason than to avoid such problems of bounding $`P_t`$, it is useful to define a modified Greens operator
$$𝒢_{L_t,ϵ}=𝒢_t(I\mathrm{\Pi }_{L_t,ϵ}),$$
where $`\mathrm{\Pi }_{L_t,e}`$ denotes the projection onto the $`\lambda ϵ`$ eigenspaces of $`L_t^{}L_t`$. This operator is bounded by $`1+ϵ^{1/2}`$ and is a parametrix as long as $`ϵ`$ lies below the essential spectrum.
### 5.3 Deforming the Index
For several of the index problems we will be considering, it seems likely that one can simply deform the given operator into a standard superharmonic oscillator and then immediately deduce the index. There are some minor issues fitting such a deformation into a continuous family. We will not treat those here because for one of our index computations \- that of the noncommutative instanton - there is no single model operator to which to deform. Instead we will use the deformation invariance of the index to localize all the problems to an elementary computation around the zeros of our triholomorphic vector field $`G`$.
The case of interest to us then is $`D`$ a Dirac operator anticommuting with an involution $`\tau `$, $`L`$ the restriction of $`D`$ to the $`+1`$ eigenspace of $`\tau `$, and $`E_1`$ and $`E_2`$ the spaces of sections of the associated bundles with finite graph and $`L_2`$ norms respectively. In this context, Fredholmness follows from the conditions in the preceding sections guaranteeing a mass gap (i.e., bounding the essential spectrum of $`D^{}D`$ away from $`0`$.)
The deformations we will consider involve replacing $`G`$ by $`TG`$ for some $`T`$ large. This fits into the above framework without modification since $`T1`$ ensures preservation of the mass gap. Moreover, scaling $`G`$ is clearly continuous because the norm of $`G`$ is a bounded in the given metric. Recall that $`G`$ enters the Dirac operators in the form of operators,
$$\lambda ^mG_m$$
(59)
or
$$(\sqrt{i}\varphi ^m\pm \sqrt{i}\varphi ^m)G_m$$
(60)
which are Clifford multiplications by $`G`$. Denote these operators by $`\widehat{G}`$. We see that the sup norm of the difference between two Dirac operators (associated to $`TG`$ and $`SG`$) is bounded by $`(T\widehat{G}S\widehat{G})f|(TS)|\times |G|_{sup}\times f,`$ which clearly gives the desired inequalities for the continuity of the deformation. We note that even had the metric allowed for unbounded $`|G|`$, we still would have $`\widehat{G}`$ bounded as an operator from $`E_1`$ (equipped with the graph norm) to $`E_2`$, as in the oscillator example of the previous section.
In addition, we will modify the metric on certain compact subsets. This modification may change the actual domain and range of our operator. For example $`\tau `$, and hence its eigenspaces may vary with the metric. Nonetheless, we may choose quasiisometries between them. Thus if we have Fredholm operators $`D_T:E_1(T)E_2(T)`$ and quasiisometries $`h_i(T):E_iE_i(T)`$ then the index of $`h_2(T)^1D_Th_1(T)`$ is $`T`$ independent by the proposition and is equal to the index of $`D_T`$ since the index is unchanged under composition with bounded operators with bounded inverse. We note, although we will not need it here, that the condition that $`h_i`$ be quasiisometries may be relaxed to the condition that the eigenvalues of $`h_i`$ and $`h_i^1`$ grow at most polynomially (subexponentially even) in distance from some choice of origin. This is an easy consequence of the fact that the Fredholm estimate implies exponential decay of the elements in the $`L_2`$ kernel of $`D_T`$. (See ). As we will use these decay properties, let us recall them in a crude form now.
Suppose we have $`N`$ points $`y_i,i=1,\mathrm{},N`$ and a Hamiltonian of the form $`H=\mathrm{\Delta }+4T^2V,`$ with $`V(x)1`$ if $`|xy_i|>1`$, $`i=1,\mathrm{},N`$. Suppose also that $`Hf=\lambda _0^2f`$ for some small constant $`\lambda _0`$, and $`fL_2`$, say with $`L_2`$ norm $`1.`$ Let
$$|x|_m:=min_{1iN}|xy_i|.$$
Then $`e^{(T^2\lambda _0^2)^{1/2}|x|_m}fL_2,`$ and the $`L_2`$ norm of $`|e^{(T^2\lambda _0^2)^{1/2}|x|_m}f|`$ restricted to the exterior of the balls of radius $`R>1`$ about the $`y_i`$ is finite and bounded by $`4e^{(T^2\lambda _0^2)^{1/2}R}`$. (This is not sharp. See for sharper statements.) In particular, we observe that the $`L_2`$ norm of $`f`$ restricted to the complement of the balls of radius $`2R`$ about the $`y_i`$ satisfies
$$f_{|B_{2R}^c}^24e^{2(T^2\lambda _0^2)^{1/2}R}.$$
(61)
Hence, $`f`$ is concentrated near the zeroes of $`V`$.
Let $`D_T`$ denote our Dirac operator with $`G`$ replaced by $`2TG`$ and the metric modified to be Euclidean in a ball of radius $`10R`$ some $`R>>1`$ about each zero of $`|G|^2`$. This metric modification allows us to compare $`D_T`$ to a model Dirac operator which agrees with $`D_T`$ near the zeros of $`G`$ and has known index. Assume, as we may by replacing $`G`$ initially by a suitable multiple, that $`D_T^{}D_T`$ has the form $`\mathrm{\Delta }+4T^2V`$, with $`V(x)>1`$ in the complement of the balls of radius 1 about each zero of $`G`$. As in the previous section, for $`ϵ`$ below the continuous spectrum of $`D_T^2`$, $`\mathrm{\Pi }_{D_T,ϵ}`$ denotes the projection onto the $`\lambda ϵ`$ eigenspaces of $`D_T^2`$. Then
$$\mathrm{Index}D_T^+=𝑑x\mathrm{tr}\tau \mathrm{\Pi }_{T,ϵ}(x,x).$$
(62)
Using (61) we see that for $`d_ϵ:=\mathrm{rank}\mathrm{\Pi }_{D_T,ϵ}`$,
$$𝑑x\mathrm{tr}\tau \mathrm{\Pi }_{D_T,ϵ}(x,x)=_{|x|_m<2c}𝑑x\mathrm{tr}\tau \mathrm{\Pi }_{D_T,ϵ}(x,x)+O(d_ϵe^{2c(T^2ϵ)^{1/2}}),$$
(63)
for some choice of $`c`$. Hence it suffices to bound $`d_ϵ`$ independently of $`T`$ large and to compute the integral of $`\mathrm{tr}\tau \mathrm{\Pi }_{T,ϵ}(x,x)`$ over $`|x|_m<2c`$ in the large $`T`$ limit.
First we estimate $`d_ϵ`$. Let $`D_T^2f=\lambda ^2f,`$ for $`\lambda ϵ`$, and $`f_{L_2}=1.`$ Let $`Q_T`$ denote the Green’s operator for the super harmonic oscillator (SHO) which agrees with $`D_T^2`$ in a neighborhood of radius $`4R`$ about the zeros of $`G`$. Let $`\rho _R`$ denote a cutoff function supported on a ball of radius $`2R`$ where $`D_T^2`$ is the SHO and identically one on a ball of radius $`R`$. Then $`(\rho _RfQ_TD_T^2\rho _Rf)`$ is in the kernel of the SHO. Denote its norm by $`a`$ and introduce a unit vector $`v`$ in the kernel of SHO such that
$$(\rho _RfQ_TD_T^2\rho _Rf)=av.$$
Observe that $`Q_T`$ has sup norm $`T^1.`$ Now consider the equality
$$D_T^2\rho _Rf=\lambda ^2\rho _Rf+[\mathrm{\Delta },\rho _R]f.$$
By our assumptions, the right hand side is $`O(\lambda ^2)+O(e^{TR})`$ (not sharp). Hence $`\rho _Rfav=O(\lambda ^2)`$ for $`\lambda >O(e^{TR/2}).`$ Setting, for example, $`ϵ=1/T,`$ we have then
$$fav=O(1/T^2).$$
(64)
Moreover, such an inequality is true for any vector in the image of $`\mathrm{\Pi }_{D_T,1/T}`$. We conclude then that rank $`\mathrm{\Pi }_{D_T,1/T}`$ is no larger than the dimension of the kernel of the SHO (times the number of zeros of $`G`$). This bounds $`d_{1/T}`$ and completes our demonstration that it suffices to compute the trace over a bounded region.
In the following, for simplicity of notation we will consider the case of a single zero for $`G`$, but the general case follows similarly with only notational complications.
Let $`S_T^+`$ and $`\tau _E`$ denote the Euclidean Dirac operator and involution which agree with $`D_T^+`$ and $`\tau `$ near the zeros of $`G`$. Let $`F_T`$ denote the Greens operators for $`S_T^+`$. To define a parametrix for $`D_T`$, introduce $`𝒢_T`$, the Greens operator for $`D_T^+`$, and let $`P_T`$ be the modified Greens operator
$$P_T:=𝒢_T(I\mathrm{\Pi }_{D_T,1/T}).$$
Define
$$I_1:=\mathrm{Index}(D_T^+)\mathrm{Index}(S_T^+)=Tr([D_T^+P_TS_T^+F_T][P_TD_T^+F_TS_T^+]).$$
Then we wish to show that the integer $`I_1=0`$. By (63), the above traces can be approximated for $`R,T`$ large as
$$I_1=\mathrm{Tr}\rho _R([D_T^+P_TS_T^+F_T][P_TD_T^+F_TS_T^+])+O(d_{1/T}e^{RT/2}).$$
On the support of $`\rho _R,`$ $`D_T^+=S_T^+,`$ hence we have
$$I_1=\mathrm{Tr}\rho _R[D_T^+,P_TF_T]+O(d_{1/T}e^{RT/2}).$$
Using the cyclic property of the trace, we rewrite the first term on the right hand side of the above formula as
$$\mathrm{Tr}\rho _R[D_T^+,P_TF_T]=Tr[D^+,\rho _R](P_TF_T)+𝑑x_iV(x)^i.$$
where $`V(x)`$ is the vector with
$$V_i(x):=\mathrm{Tr}\gamma _i(\rho _R(x)(P_TF_T)(x,x)).$$
The integral vanishes by Stoke’s theorem, leaving
$$I_1=\mathrm{tr}[D^+,\rho _R](x)(P_TF_T)(x,x)𝑑x.$$
We estimate the last term by converting it back into an expression involving the exponentially decaying projection operators. Write
$`[D_T^+,\rho _R](P_TF_T)`$
$`=`$ $`[D_T^+,\rho _R](F_TS_T^++\mathrm{\Pi }_{S_T^+,1/T})(P_TF_T)`$
$`=`$ $`[D_T^+,\rho _R]F_T(S_T^+P_TS_T^+F_T)+[D_T^+,\rho _R](\mathrm{\Pi }_{S_T^+,1/T})(P_TF_T).`$ (65)
The last term is $`O(Te^{RT})`$ because
$$[D_T^+,\rho _R](\mathrm{\Pi }_{S_T^+,1/T})=O(e^{RT})$$
by (61) and because $`(P_TF_T)`$ has sup norm $`<T`$ by construction. We separate the first term into two additional terms
$`[D_T^+,\rho _R]F_T(S_T^+P_TS_T^+F_T)`$
$`=`$ $`[D_T^+,\rho _R]F_T(D_T^+P_TS_T^+F_T)+[D^+,\rho _R]F_T(S_T^+D_T^+)P_T.`$ (66)
The first term is again exponentially decreasing because $`(D_T^+P_TS_T^+F_T)`$ is a difference of exponentially decaying projection operators and $`F_T`$ is uniformly bounded. We can compute $`F_T`$ explicitly and it is $`O(e^{T\delta })`$, where $`\delta `$ is the distance between the support of $`[D_T^+,\rho _R]`$ and the support of $`(S_T^+D_T^+)`$. Hence all the terms are exponentially decreasing, and we deduce that $`I_1`$ is exponentially decreasing. On the other hand, it is the difference between two integers and must therefore vanish.
We summarize our results.
$$\mathrm{Index}(D_T^+)=\mathrm{Index}(S_T^+).$$
When there is more than one zero of $`G`$, a minor variation of the same argument yields
$$\mathrm{Index}(D_T^+)=\underset{i}{}\mathrm{Index}(S_T(i)^+),$$
where $`S_T(i)`$ is the local model for $`D_T`$ at the $`i^{th}`$ zero of $`G`$. In order to extend our argument to this case, we must replace expressions of the form $`F_TP_T`$ and $`\mathrm{\Pi }_{S_T^+,1/T}P_T`$ in the previous expression by $`F_T\rho _{nT}P_T`$ and $`\mathrm{\Pi }_{S_T^+,1/T}\rho _{nT}P_T`$ for some large $`n`$, because $`F_T`$ and $`\mathrm{\Pi }_{S_T^+,1/T}`$ need not extend naturally to the full moduli space in the many zero case. This will introduce new error terms of the form $`\rho _{nT}P_T\rho _Ttobeestimated.`$ Our decay estimates can once again be used to show these terms are also exponentially decaying.
### 5.4 Computing the Deformed Index
We now use the deformation arguments of the preceding section to complete the index computations. We consider first the case of the quantum mechanics with 4 complex supersymmetries on a moduli space of dimension 4k and compute $`_s^+.`$ The other cases are very similar and follow with minor modifications. In the $`_s^+`$ case, we have reduced the problem to computing the index of the operator $`S_{1/e}:=d\iota _G+i(d^{}\iota _G^{})`$ acting on selfdual forms ($`\tau _+=1`$) on $`C^{2k}`$. Separating variables, we see that the index of $`S_1`$ is the product of the indices of the $`D_c`$, $`c=1,\mathrm{},k,`$ where where $`D_c:=d\iota _{a_cK_c}+i(d^{}\iota _{a_cK_c}^{})`$ (no sum over $`c`$) acting on selfdual forms on $`C^2`$. Using the deformation invariance of the index again, we may assume $`a_c=2`$.
This latter index is easy to calculate exactly as follows. We compute
$`D_cD_c^{}+D_c^{}D_c`$
$`=\mathrm{\Delta }+|2K_c|^2\{d,\iota _{2K_c}^{}\}+i\{\iota _{2K_c},d\}i\{d^{},\iota _{2K_c}^{}\}\{d^{},\iota _{2K_c}\}.`$ (67)
Let $`z_1`$ and $`z_2`$ denote complex coordinates on $`C^2`$. Then $`2K_c=i_{j=1}^2(z_c\frac{}{z_c}\overline{z}_c\frac{}{\overline{z}_c}).`$ Hence, $`|2K_c|^2=|z_1|^2+|z_2|^2.`$ In the coordinate frame, on $`(p,q)`$ forms we have
$$i\{\iota _{2K_c},d\}=i\{d^{},\iota _{2K_c}^{}\}=i_{2K_c}=(pq)2K_c/i,$$
and
$`\{d,\iota _{2K_c}^{}\}=idz_1d\overline{z}_1+idz_2d\overline{z}_2,`$
$`\{d^{},\iota _{2K_c}\}=idz_1^{}d\overline{z}_1^{}+idz_2^{}d\overline{z}_2^{}.`$ (68)
Hence, we have
$`D_cD_c^{}+D_c^{}D_c`$
$`=`$ $`\mathrm{\Delta }+|z|^22(pq)4K_c/i\{d,\iota _{2K_c}^{}\}\{d^{},\iota _{2K_c}\}.`$ (69)
The functions
$$f(a,b,c,d):=(_{z_1}\overline{z_1})^a(_{\overline{z}_1}z_1)^b(_{z_2}\overline{z_2})^c(_{\overline{z}_2}z_2)^de^{|z|^2/2},$$
(70)
with $`a,b,c,d0,`$ span the eigenspace of $`\mathrm{\Delta }+|z|^2`$ with the eigenvalue $`2(a+b+c+d)+4.`$ We compute commutators to obtain
$$2K_cf(a,b,c,d)=i(a+bc+d)f(a,b,c,d).$$
Therefore, on the algebraic span of $`f(a,b,c,d)`$ we have
$`D_cD_c^{}+D_c^{}D_c=4(a+c)+42(pq)\{d,\iota _{2K_c}^{}\}\{d^{},\iota _{2K_c}\}.`$ (71)
This vanishes if and only if $`a=c=0`$ and the differential form coefficient of $`f(a,b,c,d)`$ takes one of the following forms:
$$dz_1dz_2,$$
(72)
or
$$1+(idz_1d\overline{z}_1+idz_2d\overline{z}_2)/2dz_1d\overline{z}_1dz_2d\overline{z}_2,$$
or
$$(dz_1+idz_2d\overline{z}_2dz_1/2),$$
or
$$(dz_2+idz_1d\overline{z}_1dz_2/2).$$
We have, therefore, an infinite dimensional kernel to our operator before taking into account the constraint on charge. We now recall that we wish to restrict to the space
$$2q_c=_{2K_c}/i=(pq)+2K_c/i.$$
With the above normalization of $`K_c`$, $`q_c`$’s are integers or half-integers. On $`f(0,b,0,d)`$ this imposes the constraint
$$2q_c=(pq)+(b+d).$$
The index of $`D_c`$ is thus given by the number of ways to choose nonnegative integers $`b`$ and $`d`$ so that
$$b+d=2q_c(pq)$$
with $`pq`$ = $`2`$ or $`0`$ plus twice the number of ways to choose nonnegative integers $`b`$ and $`d`$ so that
$$b+d=2q_c1.$$
There are $`8q_c`$ such solutions for positive $`q_c`$, one such solution for $`q_c=0`$, and none for negative $`q_c`$. All of these solutions are self-dual, so the index of $`D_c`$ is
$$\mathrm{Index}(D_c)=\left\{\begin{array}{cc}8q_c\hfill & q_c>0\hfill \\ 1\hfill & q_c=0\hfill \\ 0\hfill & q_c<0\hfill \end{array}\right\}.$$
Note that this result assumes a positive coefficient of $`K_c`$. For a negative coefficient, the computation proceeds exactly as above, provided that we make the following exchanges of coordinates,
$`z_1`$ $``$ $`\overline{z}_1`$
$`z_2`$ $``$ $`\overline{z}_2`$ (73)
This maps $`K_c`$ to $`K_c`$, and flips the sign of $`q_c`$ in the charge constraint above. In other words, the sign condition in the index formula is really on $`a_cq_c`$ for each $`j`$. Thus the index is
$$\mathrm{Index}(D_c)=\left\{\begin{array}{cc}8|q_c|\hfill & a_cq_c>0\hfill \\ 1\hfill & a_cq_c=0\hfill \\ 0\hfill & a_cq_c<0\hfill \end{array}\right\}.$$
for each $`c`$.
Thus, whenever there exist a massgap, the index $`_s^+`$ is
$$_s^+=\left(\underset{c}{}\left\{\begin{array}{cc}8|q_c|\hfill & a_cq_c>0\hfill \\ 1\hfill & a_cq_c=0\hfill \\ 0\hfill & a_cq_c<0\hfill \end{array}\right\}\right).$$
where the sum is over the zeros of the potential. Note that the index is nonvanishing only if all $`a_cq_c`$ (no summation) are nonnegative. The states in the kernel of the Dirac operator must be annihilated by $`𝒵`$ as well, and the central charge $`𝒵`$ of the states
$$e\underset{c}{}a_cq_c>0$$
(74)
equals the energy.
Computation of $`_s^{}`$, appropriate for those states with positive central charge, proceeds similarly. In fact, this problem can be mapped to that of $`_s^+`$ by
$`\phi `$ $``$ $`\phi ^{}`$
$`\phi ^{}`$ $``$ $`\phi `$
$`K_c`$ $``$ $`K_c`$ (75)
The net effect is to flip the sign condition on the charges $`q_c`$, so
$$_s^{}=\left(\underset{c}{}\left\{\begin{array}{cc}8|q_c|\hfill & a_cq_c<0\hfill \\ 1\hfill & a_cq_c=0\hfill \\ 0\hfill & a_cq_c>0\hfill \end{array}\right\}\right).$$
whenever a massgap exists. The sum is over zeros of the potential. The energy of the contributing states is $`e_ca_cq_c>0`$.
We consider next the same set of operators but now restricted to the $`+1`$ eigenspace of $`\tau _4`$. We see that the terms with $`p+q`$ even are in the $`+1`$ eigenspace of $`\tau _4`$, and the terms with $`p+q`$ odd are in the $`1`$ eigenspace of $`\tau _4`$. This leads to zero index for all nonzero $`q_c`$. When $`q_c=0`$, we get a solution with $`pq=0`$. There is only one of these. The index of $`\tau _4`$ is then
$$_4=\left(\underset{c}{}\left\{\begin{array}{cc}1\hfill & q_c=0\hfill \\ 0\hfill & q_c0\hfill \end{array}\right\}\right)$$
where the sum is over the zeros of the potential.
Finally, we consider the minor modifications necessary to compute $`_2`$. Once again a separation of variables allows us to reduce the index of the Euclidean operator to a product of indices of operators $`B_c`$, $`l=1,\mathrm{},k`$ on $`C^2`$. In coordinates, $`B_c`$ has the form
$$B_c=\underset{j=1}{\overset{2}{}}[\lambda _{2j1}(\frac{}{x_j}+iy_j)+\lambda _{2j}(\frac{}{y_j}ix_j)],$$
acting on the $`+1`$ eigenspace of $`4\lambda _1\lambda _2\lambda _3\lambda _4`$. This choice of the Dirac operator corresponds to positive coefficients $`a_c=2`$ for all $`c=1,\mathrm{},k`$. Then in a covariant constant frame,
$$2B_cB_c^{}+2B_c^{}B_c=\mathrm{\Delta }+|z|^2+4iK4i\underset{j}{}\lambda _{2j1}\lambda _{2j}.$$
Acting on the algebraic span of $`f(a,b,c,d)`$,
$$2B_cB_c^{}+2B_c^{}B_c=4(a+c)+4+4i\underset{j}{}\lambda _{2j1}\lambda _{2j}.$$
This has infinite dimensional kernel spanned by the product of $`f(0,b,0,d)`$ and a covariant constant spinor in the (one dimensional) intersection of the $`1/2`$ eigenspaces of $`i\lambda _1\lambda _2`$ and $`i\lambda _3\lambda _4.`$
The charge constraint in a covariant constant frame takes the form
$$2q_c=2K_c/ii\lambda _1\lambda _2i\lambda _3\lambda _4.$$
Acting on the above basis elements of the kernel of $`2B_cB_c^{}+2B_c^{}B_c`$ this reduces to
$$2q_c=b+d+1.$$
Counting as before this yields a $`2q_c`$ dimensional kernel which lies entirely in the $`+1`$ eigenspace of $`4\lambda _1\lambda _2\lambda _3\lambda _4`$. Hence, the index of $`B_c`$ is
$$\mathrm{Index}(B_c)=\left\{\begin{array}{cc}2|q_c|\hfill & a_cq_c>0\hfill \\ 0\hfill & a_cq_c0\hfill \end{array}\right\}.$$
Thus, whenever there exist a massgap, the index $`_2`$ is
$$_2=\left(\underset{c}{}\left\{\begin{array}{cc}2|q_c|\hfill & a_cq_c>0\hfill \\ 0\hfill & a_cq_c0\hfill \end{array}\right\}\right).$$
Finally, we conclude the index computations by noting that they reproduce the four-dimensional results summarized in the previous section. In fact, the wavefunctions found in Ref. and in Ref. can be seen easily to reduce to the superharmonic oscillator wavefunctions above in the limit of $`\stackrel{~}{a}q`$.
## 6 BPS Bound States
The above index computations count differences in the number of ground states with respect to $`Z_2`$ involutions $`\tau _\pm ,\tau _4,\tau _2`$,
$$\mathrm{Index}=n_+n_{}$$
(76)
where $`n_\pm `$ are the number of ground states with $`\tau `$ eigenvalue $`\pm 1`$. We are actually interested in the sum $`n_++n_{}`$ instead, for which one needs a more refined understanding of the dynamics. For the case of $`\tau _2`$ and $`\tau _\pm `$, we anticipate $`n_{}`$ vanishes by itself. Such a vanishing theorem is shown rigorously for the simplest cases in Appendix. We will assume in this section that $`n_{}=0`$ holds true for $`\tau _2`$ and $`\tau _\pm `$ in all cases, and compare the results to what are expected on physical grounds.
### 6.1 $`N=4`$ Yang-Mills Theories
The supersymmetric quantum mechanics with four complex supercharges describe dynamics of monopoles in $`N=4`$ Yang-Mills theories. Recent studies of D-branes indicates the following three possibilities for dyonic bound states of monopoles.
* The state is 1/2 BPS in the Yang-Mills field theory. These states would be annihilated by all supercharges of the low energy monopole dynamics, which is possible only if the central charges in the relative part of the dynamics is absent. This is guaranteed when all relative electric charge $`q_a`$’s vanish. In particular, this includes purely magnetic bound states.
* The state is 1/4 BPS in the Yang-Mills field theory. These states would be annihilated by half of the supercharges of the low energy monopole dynamics and not by the other half. This is possible only if at least one central charge is nonzero.
* The state is non-BPS.
The index computation of the previous section tells us something about 1/2 BPS and 1/4 BPS states, where we counted indices $`_4`$ and $`_s^\pm `$ in the special limit where only one $`G^I`$, say $`G^5`$, is turned on. Equivalently, we considered vacua where two Higgs fields are turned on.
Of the three indices, only $`_4`$ is robust against turning on more than one $`G^I`$’s. The Dirac operator $`iQ\pm Q^{}`$ would no longer anticommute with $`\tau _\pm `$ but does anticommutes with $`\tau _4`$. Only $`_4`$ is a well-defined index in such generic vacua. Turning on additional $`G^I`$ always increases the massgap, and is a Fredholm deformation that preserves $`_4`$. Thus our result shows that, in generic vacua,
$$_4=1,$$
(77)
when $`q_a0`$, and zero otherwise. Since the central charge of the state that contributes to the index is zero, the state must be annihilated by all supercharges of the quantum mechanics and is a 1/2 BPS in $`N=4`$ Yang-Mills theory.<sup>3</sup><sup>3</sup>3One might think that existence of this bound state is obvious since the potentials are all attractive and also there exists a classical BPS monopole of the same magnetic charge. However, none of these guarantee the existence of BPS bound state at quantum level. In fact, the same set of facts are true for a pair of distinct monopoles in $`N=2`$ $`SU(3)`$ Yang-Mills theory but we know that such a purely magnetic bound state does not exist as a BPS state . This is consistent with the existence of a unique magnetic 1/2 BPS bound state of monopoles in generic Coulomb vacua, which is expected from the $`SL(2,Z)`$ electromagnetic duality. One of the generators of $`SL(2,Z)`$ maps massive charged vector multiplets to purely magnetic bound states in 1-1 fashion. After taking into account the automatic degeneracy $`16`$ from the free center-of-mass fermions, the total degeneracy of these bound states is alway 16, which fits the $`N=4`$ vector multiplet nicely. This purely magnetic bound state was previously constructed by Gibbons in special vacua where all $`G^I`$’s vanishes.
Existence of 1/4 BPS states are more sensitive to the vacuum choice and the electric charges. The existence criteria were first found by Bergman , where he constructed these dyons as string webs ending on D3-branes. The first necessary condition is that the string web should be planar, which is equivalent to the condition that, effectively, only one linearly independent $`G^I`$ is present. This allows us to assume without loss of generality that only $`G^5`$ is turned on, as far as counting 1/4 BPS states are concerned. Thus, the computation of $`_s^\pm `$ in the previous section is directly applicable.
Secondly, at each junction of the string web, the string tensions must balance against each other, which in the present language of low energy dynamics translates to the condition that the effective potential in the charge-eigensector is nonrepulsive along all asymptotic directions ;
$$|q_c||\stackrel{~}{a}_c|.$$
(78)
This second condition may indicate the existence of a minimal energy bound state, however, does not guarantee that the state would preserve some supersymmetry.
Finally, a minimal energy configuration is supersymmetric when the orientation of string segments are consistent with each other. Say, if one fundamental string segment is directed to one particular direction, then another fundamental string in the same web must be directed the same way. The second string can point toward the opposite direction and still balance the tension, but such a combination breaks all supersymmetry. This orientation condition on the string web, is nothing but the condition that the product $`a_cq_c`$’s (no summation) are all of same sign. See figure 1. Thus, a 1/4 BPS dyon may exist only when $`|q_c||\stackrel{~}{a}_c|`$ for all $`c`$ and $`a_cq_c`$ are all of same sign, at least one of which is nonzero.
The indices $`_s^\pm `$ were computed with the massgap condition $`|q_c|<|\stackrel{~}{a}_c|`$ to begin with, and yielded nonzero value only when all $`a_cq_c`$ were of the same sign; For positive signs of $`a_cq_c`$, $`_s^+0`$, while for negative $`a_cq_c`$’s, we have $`_s^{}0`$. The result is clearly consistent with the existence criteria set by the string-web construction, and furthermore gives us extra information beyond the string web picture. The index indicates that the degeneracy of such a 1/4 BPS state is
$$16\times \underset{c}{}\mathrm{Max}\{8|q_c|,1\}.$$
(79)
The factor 16 arises from the free center-of-mass fermions.
In the two monopole bound states, the number $`8|q|`$ is accounted for by four angular momentum multiplets of $`j=|q|,|q|1/2,|q|1/2,|q|1`$ .<sup>4</sup><sup>4</sup>4The first three suffices for $`|q|=1/2`$. The top angular momentum $`|q|`$ in the relative part of the wavefunction has a well-known classical origin: when an electrically charged particle moves around a magnetic object, the conserved angular momentum is shifted by a factor of $`eg/4\pi `$. While fermions can and do contribute, the number of fermions scales with the number of monopoles, and not with the charge $`q_a`$. In fact, it is most likely that the top angular momentum of such a dyonic bound state wavefunction is
$$j_{top}=\underset{c}{}|q_c|,$$
(80)
for large charges, so that the highest spin of the dyon would be
$$1+j_{top}=1+\underset{c}{}|q_c|,$$
(81)
after taking into account the universal vector multiplet structure from the free center-of-mass part. On the other hand, a 1/4 BPS supermultiplet with the highest spin $`j_{top}+1`$ has the total degeneracy of
$$16\times 8\underset{c}{}|q_c|,$$
(82)
which is much less than the number of states we found above unless all but one $`q_c`$ vanishes. Thus, this implies that there are many 1/4 BPS, thus degenerate, supermultiplets of dyons for a given set of electromagnetic charges. This is probably the least understood of our results. While one would expect to find degenerate states within a supermultiplet, there is no natural symmetry that accounts for the existence of many supermultiplets of the same electromagnetic charges and of the same energy.
For large electric charges $`q_a`$, the number of dyon supermultiplets scales as, at least,
$$\left(\underset{c}{}\mathrm{Max}\{8|q_c|,1\}\right)/\left(8\underset{c}{}|q_c|\right).$$
(83)
Proliferation of dyonic states of a given charge was anticipated by Kol some time ago in the context of string webs in 5 dimensions . Because our computation was performed for a collection of distinct monopoles, which put some constraint on the possible topology of the related string web, it is not immediately clear to us whether we can make any sensible statement in the regime where Kol’s prediction is applicable.<sup>5</sup><sup>5</sup>5Kol anticipated exponentially large numbers of states, in fact, which is much more than our powerlike result. Nevertheless, it is tantalizing that we found the number of states increasing much faster than would have been expected from supersymmetry alone. It is not clear to us why this happens and what interpretation this may have in the Yang-Mills field theory.
In the regime where $`|q_c||\stackrel{~}{a}_c|`$ for some $`q_c`$, we cannot rely on the current index computation. On the other hand, since even a single repulsive direction, i.e., $`|q_c|>|\stackrel{~}{a}_c|`$ for some $`c`$, prohibits a bound state (supersymmetric or not), the unresolved question boils down to the marginal case, where $`|q_c|`$ equals $`|\stackrel{~}{a}_c|`$ for some $`c`$’s while the others satisfy $`|q_c|<|\stackrel{~}{a}_c|`$. The only state that must exist for sure is the purely magnetic bound state ($`q_a=0`$), which was constructed by Gibbons when $`a_c0`$ and which is necessary for $`SL(2,Z)`$ invariance. The explicit construction of two-monopole bound states in Ref. seem to indicate that no dyonic bound state may form along such marginal directions, but this remains to be seen for multi-monopole cases.
### 6.2 $`N=2`$ Yang-Mills Theories
In $`N=2`$ Yang-Mills theories, a state could be either BPS or non-BPS. There is no such thing as a 1/4 BPS state. Dyons that would have been 1/4 BPS when embedded in $`N=4`$ theories, are realized as either 1/2 BPS or non-BPS depending on the sign of the electric charges. The index of this Dirac operator was nonzero only when
$`0<q_c<\stackrel{~}{a}_c`$ for all $`c`$
$`\text{or}0>q_c>\stackrel{~}{a}_c`$ for all $`c`$ (84)
which gives us a possible criterion for BPS dyon to exist. This condition is similar to the condition for BPS dyons or monopoles to exist in $`N=4`$ Yang-Mills theories but differs in two aspects. The first is that given a set of $`a_c`$’s, all of which are positive (negative), the electric charge $`q_c`$’s must be all positive (negative).
The second difference from $`N=4`$ case is that a purely magnetic bound state of monopoles does not seem to exist as a BPS state, even though there exists a classical BPS solution of such a charge. This feature was noted previously in Ref. , where bound states of a pair of distinct monopoles were counted explicitly. In fact, the index indicates that all relative $`q_a`$ must be nonvanishing for a BPS state to exist. Assuming the vanishing theorem, the number of BPS dyonic bound state under the above condition is
$$4\times \underset{c}{}2|q_c|,$$
(85)
The overall factor $`4`$ is from the quantization of the free center-of-mass fermions.
For large electric charges we again observe the proliferation of supermultiplets. The top angular momentum and thus the size of the largest supermultiplet can grow only linearly with $`|q_c|`$ which means that the number of supermultiplets with the same electric charges scales at least as
$$\left(\underset{c}{}2|q_c|\right)/\left(2\underset{c}{}|q_c|\right)$$
(86)
for large $`q_c`$’s. Again it is not clear to us what implication this has in the Yang-Mills field theory.
In the regime where $`|q_c||\stackrel{~}{a}_c|`$ for some $`q_c`$, again we cannot rely on the current index computation. For the same reason as in $`N=4`$ Yang-Mills theory, no bound state can exist if even a single repulsive direction ($`|q_c|>|\stackrel{~}{a}_c|`$ for some $`c`$) exists, so the unresolved question boils down again to the marginal case, where $`|q_c|`$ equals $`|\stackrel{~}{a}_c|`$ for some $`c`$’s while the others satisfy $`|q_c|<|\stackrel{~}{a}_c|`$. Extrapolating from the explicit construction of two-monopole bound states in Ref. , we suspect that no bound state may form along such marginal directions.
### 6.3 Ground States of a Noncommutative Instanton Soliton
Supersymmetric ground states and excited BPS states of an instanton soliton in $`S^1\times R^{3+1}`$ can be counted similarly as above. The only difference as far as computing the index goes, is that the potential has many zeros. For a single instanton in noncommutative $`U(n)`$ theory, there are precisely $`n`$ zeros of the bosonic potential, and near each of these zeros, the Dirac operator can be deformed to that of a superharmonic oscillator. One crucial difference in interpreting the result in physical terms comes from identification of the conserved charges. Of $`n`$ possible conserved $`U(1)`$ charges, $`n1`$ relative charges are again electric charges. However, the overall conserved $`U(1)`$ does not correspond to a gauge symmetry, and comes from translation of the instanton along $`S^1`$. This last $`U(1)`$ charge is just the Kaluza-Klein momentum along $`S^1`$.
Of particular interest are the quantum ground state of the instanton with no $`U(1)`$ charges excited. In the maximally supersymmetric $`U(n)`$ Yang-Mills theory, the index tells us that there are $`n`$ distinct BPS supermultiplets of ground states. This result was anticipated in Ref. . With half as much supersymmetry, however, the index is consistent with no supersymmetric quantum ground state exist at all.
Quantum states of instanton soliton in $`R^4`$ was previously studied in the commutative setting . In particular, the absence of a quantum ground state of instanton soliton in the nonmaximal supersymmetric Yang-Mills theories, has been observed from string-web construction.
## 7 Conclusion
By computing indices and assuming vanishing theorems, we counted supersymmetric bound states of arbitrarily many distinct monopoles in $`N=2`$ pure Yang-Mills theories and also in $`N=4`$ Yang-Mills theories. The relevant low energy dynamics are supersymmetric sigma-models with potential(s), where the supercharges preserved by supersymmetric bound states can be interpreted as Dirac operators twisted by triholomorphic Killing vector fields. An obvious generalization of this computation is to include hypermultiplets in $`N=2`$ Yang-Mills theories, but it goes beyond the scope of this paper.
Counting of 1/2 BPS states in $`N=4`$ Yang-Mills yielded a result consistent with electromagnetic duality of the theory. In particular, the necessary purely magnetic bound states of distinct monopoles are all accounted for in $`SU(n)`$ theories. While this result is not surprising, it is still significant in that this was shown for the first time in all generic Coulomb vacua of the Yang-Mills theory. In contrast, distinct $`N=2`$ monopoles do not seem to bind at all unless all possible relative charges are turned on.
Existence criteria for 1/4 BPS states, previously found in the context of string-webs, are also faithfully reflected in the index formulae. On the other hand, the degeneracy of most 1/4 BPS dyons is shown to be much larger that one would have expected from a single 1/4 BPS supermultiplet with a physically reasonable angular momentum. $`N=2`$ Dyons of the same electromagnetic charges as 1/4 BPS dyons of $`N=4`$ theories, could be BPS or non-BPS, depending on the signs of the electric charge. We also counted the degeneracy of such $`N=2`$ BPS dyons, which shows similar proliferation of supermultiplets. This phenomenon is not understood at the moment. It should be also interesting to see how the degeneracy behaves when both electric and magnetic charges are large.
We are grateful to Jerome Gauntlett for stimulating discussions. We thank Aspen Center for Physics and also the organizers of the workshop, ”The Geometry and Physics of Monopoles”, where this work was initiated. The work of M.S. is supported by NSF grant DMS-9870161.
## Appendix
Here, we prove a vanishing theorem for $`\tau _2`$ and $`\tau _\pm `$ on four-dimensional moduli space. Let us consider $`\tau _2`$ first. Because of the triholomorphic Killing conditions on $`G`$, $`dG`$ is self-dual and does not couple to antichiral spinors. Then the Dirac operator is a simple Laplacian;
$$D_mD^m$$
(87)
when acting on antichiral spinors. Using the standard trick of sandwiching this operator by a hypothetical zero mode $`\mathrm{\Psi }`$ and its complex conjugate $`\mathrm{\Psi }^{}`$, we find
$$0=\underset{m}{}dz^m\mathrm{\Psi }^{}_mD^m\mathrm{\Psi }=\underset{m}{}dz^mg^{mn}(D_n\mathrm{\Psi })^{}(D_m\mathrm{\Psi })$$
(88)
where the possible boundary term vanishes by itself since the massgap forces $`\mathrm{\Psi }`$ to be exponentially small at large distances. Therefore,
$$0=D_m\mathrm{\Psi }=(_miG_m)\mathrm{\Psi }$$
(89)
everywhere. This modified connection is still unitary as $`G_m`$ is real. Hence, $`\mathrm{\Psi }`$ is covariant constant with respect to metric compatible connection and is therefore of constant norm. Such an $`f`$ cannot be normalizable on an infinite volume space unless it is identically zero, which proves the vanishing theorem in four dimensions for $`\tau _2`$.
The case of $`\tau _\pm `$ can be handled similarly. Let us recall that differential forms can be thought of as a tensor product of two spinors. With an appropriate sign convention, we can identify various sectors of the former with those of the latter as follows
* selfdual even form $``$ \[chiral spinor\]$``$\[chiral spinor\]
* selfdual odd form $``$ \[chiral spinor\]$``$\[antichiral spinor\]
* antiselfdual odd form $``$ \[antichiral spinor\]$``$\[chiral spinor\]
* antiselfdual even form $``$ \[antichiral spinor\]$``$\[antichiral spinor\]
On antiselfdual even form, then, the Clifford action of $`dG`$ is trivial. Since the self-dual curvature does not couple to antiselfdual forms either, its action is also trivial. Thus, the square of the Dirac operator becomes a simple Laplacian again,
$$D_\pm ^2=D_mD^m$$
(90)
with $`D_m=_miG_m^5`$. By the same logic as in the spinor case, therefore, no antiselfdual even-form solution can exist. Finally this also shows that antiselfdual odd-form solution does not exist; Unless the central charge vanishes, a solution generates other solutions via the action of broken supercharges. The broken supercharges are the linear combination of $`Q`$ and $`Q^{}`$ orthogonal to $`D_\pm `$, so that it flips $`\tau _4`$ while preserving $`\tau _\pm `$. Thus, the number of odd-form solutions equals the number of even-form solutions, in each $`\tau _\pm `$ eigensectors, whenever the central charge is nonzero. This proves the vanishing theorem for $`\tau _\pm `$ for sectors with nonzero $`U(1)`$ charges $`q_c`$.
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# 1 Introduction
## 1 Introduction
Renormalization group short-distance QCD effects play an important role in the phenomenology of non-leptonic weak transitions of $`K`$-, $`D`$\- and $`B`$-mesons. An essential ingredient in any renormalization group analysis is the anomalous dimension matrix (ADM), which describes the mixing of the relevant local four-quark operators under renormalization .
The operators considered in the present paper have the form
$$\overline{\mathrm{\Psi }}_1^\alpha \mathrm{\Gamma }_A^k\mathrm{\Psi }_2^\alpha \overline{\mathrm{\Psi }}_3^\beta \mathrm{\Gamma }_B^k\mathrm{\Psi }_4^\beta ,\overline{\mathrm{\Psi }}_1^\alpha \mathrm{\Gamma }_A^k\mathrm{\Psi }_2^\beta \overline{\mathrm{\Psi }}_3^\beta \mathrm{\Gamma }_B^k\mathrm{\Psi }_4^\alpha ,$$
(1.1)
where $`\alpha `$, $`\beta `$ are colour indices and $`\mathrm{\Gamma }_{A,B}^k`$ are generic Dirac matrices given explicitly below. The subscripts $`i`$ in $`\mathrm{\Psi }_i`$ are flavour indices. In the case of FCNC transitions with $`\mathrm{\Delta }F=2`$, such as neutral meson mixing, one has
$$\mathrm{\Psi }_1=\mathrm{\Psi }_3,\mathrm{\Psi }_2=\mathrm{\Psi }_4.$$
(1.2)
Known examples are the operators $`(\overline{s}d)_{VA}(\overline{s}d)_{VA}`$ and $`(\overline{b}d)_{VA}(\overline{b}d)_{VA}`$ relevant in the Standard Model (SM) to $`K^0`$$`\overline{K}^0`$ and $`B_d^0`$$`\overline{B}_d^0`$ mixing, respectively.
Four-quark operators that occur in the SM calculations of flavour-changing processes do not form a complete set of all the dimension-six four-quark operators. Other operators need to be considered in many extensions of the SM, e.g. in the Supersymmetric Standard Model (SSM) (see e.g. ref. ). For instance, the SSM and SM predictions for $`K^0`$$`\overline{K}^0`$ and $`B_d^0`$$`\overline{B}_d^0`$ mixing can have similar precision only if the two-loop ADM for all the $`\mathrm{\Delta }F=2`$ operators is known.
The main purpose of the present paper is a calculation of the two-loop ADM for all the dimension-six flavour-changing four-quark operators in the NDR–$`\overline{\mathrm{MS}}`$ scheme ($`\overline{\mathrm{MS}}`$ scheme with fully anticommuting $`\gamma _5`$). Our main findings are the NDR–$`\overline{\mathrm{MS}}`$ anomalous dimensions of the operators with Dirac structures (cf. eq. (1.1)):
$$\mathrm{\Gamma }_A^k\mathrm{\Gamma }_B^k=(1\pm \gamma _5)(1\pm \gamma _5)\mathrm{and}\mathrm{\Gamma }_A^k\mathrm{\Gamma }_B^k=[\sigma _{\mu \nu }(1\pm \gamma _5)][\sigma ^{\mu \nu }(1\pm \gamma _5)].$$
(1.3)
For these operators, our two-loop results differ from the NDR–$`\overline{\mathrm{MS}}`$ ones of Ciuchini et al. , but are compatible with their RI-scheme ADM. For all the other operators, no new calculation is actually necessary — all the two-loop results can be extracted from the existing Standard Model ones.
Our paper is organized as follows. In section 2, we perform a direct calculation of the NDR–$`\overline{\mathrm{MS}}`$-scheme ADM of $`\mathrm{\Delta }F=2`$ operators. This is a relatively straightforward computation, since all the methods are already known from similar SM calculations (see e.g. refs. ). The only novelty here is the introduction of evanescent operators that vanish by the Fierz identities.
In section 3, we compute the NDR–$`\overline{\mathrm{MS}}`$ ADM for such $`\mathrm{\Delta }F=1`$ operators, to which only the current–current diagrams are relevant. Some of the $`\mathrm{\Delta }F=1`$ results have never been published before. The ones that are not new agree with the old SM calculations. The subject of section 4 are $`\mathrm{\Delta }F=1`$ operators containing one quark–antiquark pair of the same flavour. We identify the operators to which the so-called penguin diagrams are relevant, and give the corresponding anomalous dimensions.
In section 5, we derive the matrix $`\mathrm{\Delta }\widehat{r}`$ necessary for transforming the Wilson coefficients from the NDR–$`\overline{\mathrm{MS}}`$ to the RI scheme (originally called the MOM scheme) that is more useful for non-perturbative calculations of hadronic matrix elements .
Section 6 is devoted to performing a consistency check of our $`\mathrm{\Delta }F=1`$ and $`\mathrm{\Delta }F=2`$ results. The current–current ADM of $`\mathrm{\Delta }F=1`$ operators is transformed there to such an operator basis, in which the $`\mathrm{\Delta }F=2`$ results can be easily read off. This calculation serves also as a preparation for the comparison with Ciuchini et al. . Comparison with this article and other existing literature is the subject of section 7. We conclude in section 8.
In appendix A, we list the evanescent operators relevant to the $`\mathrm{\Delta }F=2`$ calculation. In appendix B, an analogous list for the $`\mathrm{\Delta }F=1`$ case is presented. Appendix C contains additional evanescent operators that become important only when one wants to derive the $`\mathrm{\Delta }F=2`$ results from the $`\mathrm{\Delta }F=1`$ ones, as in section 6. Appendix D is devoted to recalling and generalizing the notion of “Greek projections”. Appendix E contains a list of separate contributions from different diagrams to the one- and two-loop ADMs for $`\mathrm{\Delta }F=1`$ operators with Dirac structures (1.3). Finally, in appendix F, we outline our determination of two-loop mixing via penguin diagrams that involves beyond-SM operators.
## 2 Direct calculation of the ADM in the $`\mathrm{\Delta }F=2`$ case
For definiteness, we shall consider here operators responsible for the $`K^0`$$`\overline{K}^0`$ mixing. There are 8 such operators of dimension 6. They can be split into 5 separate sectors, according to the chirality of the quark fields they contain. The operators belonging to the first three sectors (VLL, LR and SLL) read
$`Q_1^{\mathrm{VLL}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu P_Ld^\alpha )(\overline{s}^\beta \gamma ^\mu P_Ld^\beta ),`$
$`Q_1^{\mathrm{LR}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu P_Ld^\alpha )(\overline{s}^\beta \gamma ^\mu P_Rd^\beta ),`$
$`Q_2^{\mathrm{LR}}`$ $`=`$ $`(\overline{s}^\alpha P_Ld^\alpha )(\overline{s}^\beta P_Rd^\beta ),`$
$`Q_1^{\mathrm{SLL}}`$ $`=`$ $`(\overline{s}^\alpha P_Ld^\alpha )(\overline{s}^\beta P_Ld^\beta ),`$
$`Q_2^{\mathrm{SLL}}`$ $`=`$ $`(\overline{s}^\alpha \sigma _{\mu \nu }P_Ld^\alpha )(\overline{s}^\beta \sigma ^{\mu \nu }P_Ld^\beta ),`$ (2.1)
where $`\sigma _{\mu \nu }=\frac{1}{2}[\gamma _\mu ,\gamma _\nu ]`$ and $`P_{L,R}=\frac{1}{2}(1\gamma _5)`$. The operators belonging to the two remaining sectors (VRR and SRR) are obtained from $`Q_1^{\mathrm{VLL}}`$ and $`Q_i^{\mathrm{SLL}}`$ by interchanging $`P_L`$ and $`P_R`$. Since QCD preserves chirality, there is no mixing between different sectors. Moreover, the ADMs in the VRR and SRR sectors are the same as in the VLL and SLL sectors, respectively. In the following, we shall consider only the VLL, LR and SLL sectors.
In dimensional regularization, the four-quark operators from eq. (2.1) mix at one loop into the evanescent operators listed in appendix A. Specifying these evanescent operators is necessary to make precise the definition of the NDR–$`\overline{\mathrm{MS}}`$ scheme in the effective theory . An important novelty in the present case (when compared to $`\mathrm{\Delta }F=1`$ calculations) is the necessity of introducing evanescent operators that vanish in 4 dimensions by the Fierz identities. The Fierz identities cannot be analytically continued to $`D`$ dimensions. Therefore, they have to be treated in dimensional regularization in the same manner as the identity
$$\gamma _\mu \gamma _\nu \gamma _\rho =g_{\mu \nu }\gamma _\rho +g_{\nu \rho }\gamma _\mu g_{\mu \rho }\gamma _\nu +iϵ_{\alpha \mu \nu \rho }\gamma ^\alpha \gamma _5,$$
(2.2)
i.e. appropriate evanescent operators have to be introduced.
As an example, consider the operators $`Q_1^{\mathrm{SLL}}`$ and $`Q_2^{\mathrm{SLL}}`$. When these operators are inserted into one- and two-loop diagrams, the operators
$`\stackrel{~}{Q}_1^{\mathrm{SLL}}`$ $`=`$ $`(\overline{s}^\alpha P_Ld^\beta )(\overline{s}^\beta P_Ld^\alpha ),`$ (2.3)
$`\stackrel{~}{Q}_2^{\mathrm{SLL}}`$ $`=`$ $`(\overline{s}^\alpha \sigma _{\mu \nu }P_Ld^\beta )(\overline{s}^\beta \sigma ^{\mu \nu }P_Ld^\alpha )`$ (2.4)
are generated. In 4 dimensions these operators can be expressed through $`Q_1^{\mathrm{SLL}}`$ and $`Q_2^{\mathrm{SLL}}`$ by using the Fierz identities
$`(P_L)_{ij}(P_L)_{kl}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(P_L)_{il}(P_L)_{kj}{\displaystyle \frac{1}{8}}(\sigma _{\mu \nu }P_L)_{il}(\sigma ^{\mu \nu }P_L)_{kj},`$
$`(\sigma _{\mu \nu }P_L)_{ij}(\sigma ^{\mu \nu }P_L)_{kl}`$ $`=`$ $`6(P_L)_{il}(P_L)_{kj}{\displaystyle \frac{1}{2}}(\sigma _{\mu \nu }P_L)_{il}(\sigma ^{\mu \nu }P_L)_{kj},`$ (2.5)
which give
$`\stackrel{~}{Q}_1^{\mathrm{SLL}}`$ $`\underset{D=4}{=}`$ $`{\displaystyle \frac{1}{2}}Q_1^{\mathrm{SLL}}+{\displaystyle \frac{1}{8}}Q_2^{\mathrm{SLL}},`$ (2.6)
$`\stackrel{~}{Q}_2^{\mathrm{SLL}}`$ $`\underset{D=4}{=}`$ $`6Q_1^{\mathrm{SLL}}+{\displaystyle \frac{1}{2}}Q_2^{\mathrm{SLL}}.`$ (2.7)
These relations can be used in the calculation of one-loop ADM. In the case of two-loop calculations, in the NDR–$`\overline{\mathrm{MS}}`$ scheme, where Dirac algebra has to be performed in $`D4`$ dimensions, these relations have to be generalized to
$`\stackrel{~}{Q}_1^{\mathrm{SLL}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}Q_1^{\mathrm{SLL}}+{\displaystyle \frac{1}{8}}Q_2^{\mathrm{SLL}}+E_1^{\mathrm{SLL}},`$ (2.8)
$`\stackrel{~}{Q}_2^{\mathrm{SLL}}`$ $`=`$ $`6Q_1^{\mathrm{SLL}}+{\displaystyle \frac{1}{2}}Q_2^{\mathrm{SLL}}+E_2^{\mathrm{SLL}}.`$ (2.9)
Here, $`E_1^{\mathrm{SLL}}`$ and $`E_2^{\mathrm{SLL}}`$ are the evanescent operators that vanish in 4 dimensions by Fierz identities. They are simply defined by (2.8) and (2.9) and are given in appendix A.
The effective Lagrangian can be written separately for each sector. It takes the form
$$_{eff}=\frac{G_F^2M_W^2}{4\pi ^2}(V_{ts}^{}V_{td})^2Z_q^2\underset{i}{}C_i(\mu )\left[Q_i+(\text{counterterms})_i\right],$$
(2.10)
where $`Z_q`$ is the quark wave-function renormalization constant.
The coefficients $`C_i(\mu )`$ satisfy the Renormalization Group Equation (RGE)
$$\mu \frac{d}{d\mu }\stackrel{}{C}(\mu )=\widehat{\gamma }(\mu )^T\stackrel{}{C}(\mu )$$
(2.11)
governed by the ADM $`\widehat{\gamma }(\mu )`$ that has the following perturbative expansion:
$$\widehat{\gamma }(\mu )=\frac{\alpha _s(\mu )}{4\pi }\widehat{\gamma }^{(0)}+\frac{\alpha _s^2(\mu )}{(4\pi )^2}\widehat{\gamma }^{(1)}+𝒪(\alpha _s^3).$$
(2.12)
The ADM in the MS or $`\overline{\mathrm{MS}}`$ scheme is found from one- and two-loop counterterms in the effective theory, according to the following relations (equivalent to eqs. (4.26)–(4.37) of ref. ):
$`\widehat{\gamma }^{(0)}`$ $`=`$ $`2\widehat{a}^{11},`$ (2.13)
$`\widehat{\gamma }^{(1)}`$ $`=`$ $`4\widehat{a}^{12}2\widehat{b}\widehat{c}.`$ (2.14)
The matrices $`\widehat{a}^{11}`$, $`\widehat{a}^{12}`$ and $`\widehat{b}`$ in the above equations parametrize the MS-scheme counterterms in eq. (2.10) (for $`D=42ϵ`$)
$`(\text{counterterms})_i`$ $`=`$ $`{\displaystyle \frac{\alpha _s}{4\pi ϵ}}\left[{\displaystyle \underset{k}{}}a_{ik}^{11}Q_k+{\displaystyle \underset{k}{}}b_{ik}E_k\right]+{\displaystyle \frac{\alpha _s^2}{(4\pi )^2}}{\displaystyle \underset{k}{}}\left({\displaystyle \frac{1}{ϵ^2}}a_{ik}^{22}+{\displaystyle \frac{1}{ϵ}}a_{ik}^{12}\right)Q_k`$ (2.15)
$`+`$ $`(\text{two-loop evanescent counterterms})+𝒪(\alpha _s^3).`$
The matrix $`\widehat{c}`$ is recovered from one-loop matrix elements of the evanescent operators. Let us denote by $`E_k_{1\mathrm{l}\mathrm{o}\mathrm{o}\mathrm{p}}`$ the one-loop $`K^0`$$`\overline{K^0}`$ amplitude with an insertion of some evanescent operator $`E_k`$. The pole part of such an amplitude is proportional to some linear combination of tree-level matrix elements of evanescent operators. The remaining part in the limit $`D4`$ can be expressed by tree-level matrix elements of the physical operators $`Q_i`$. The finite coefficients of these matrix elements define the matrix $`\widehat{c}`$ as follows:
$`E_k_{1\mathrm{l}\mathrm{o}\mathrm{o}\mathrm{p}}`$ $`=`$ $`{\displaystyle \frac{1}{ϵ}}\left[{\displaystyle \underset{j}{}}d_{kj}E_j_{\mathrm{tree}}+{\displaystyle \underset{j}{}}e_{kj}F_j_{\mathrm{tree}}\right]{\displaystyle \underset{i}{}}c_{ki}Q_i_{\mathrm{tree}}+𝒪(ϵ).`$ (2.16)
Here, $`F_j`$ stand for such evanescent operators that are not necessary as counterterms for the one-loop Green functions with insertions of the physical operators $`Q_i`$. The matrices $`\widehat{c}`$ and $`\widehat{a}^{12}`$ depend on the structure of $`F_j`$, but $`\widehat{\gamma }^{(1)}`$ does not.
The matrices $`\widehat{\gamma }^{(0)}=2\widehat{a}^{11}`$$`\widehat{b}`$ and $`\widehat{c}`$ in each sector are found from the one-loop $`d\overline{s}s\overline{d}`$ diagrams presented in fig. 1 with insertions of the physical operators $`Q_i`$, as well as the evanescent operators $`E_k`$. We calculate only the “annihilation-type” diagrams, i.e. we drop all the diagrams where fermion lines connect the incoming and outgoing particles. Dropping such diagrams consistently at the tree level, at one loop and (later) at two loops does not alter the final results for the renormalization constants.
All the one- and two-loop diagrams considered in the present article are calculated using two different methods. In both of them, a covariant gauge-fixing term
$$_{gf}=\frac{1}{2\lambda }(^\mu G_\mu ^a)(^\nu G_\nu ^a)$$
(2.17)
is used, and the physical masses are set to zero. In the first method, the external quarks are assumed to have momentum $`\pm p`$. In the second method, the external momenta are set to zero, but a common mass parameter is introduced in all the propagator denominators as IR regulator . The two methods give the same results for the $`\overline{\mathrm{MS}}`$ renormalization constants. The ADMs calculated from these renormalization constants with the help of eqs. (2.13) and (2.14) are independent of the gauge-fixing parameter $`\lambda `$.
We begin with presenting the ADM in the SLL sector, because in this very sector our results are going to differ (at two loops) from those of ref. . The matrices $`\widehat{\gamma }^{(0)\mathrm{SLL}}`$ and $`\widehat{b}^{\mathrm{SLL}}`$ are found to be the following:
$`\widehat{\gamma }^{(0)\mathrm{SLL}}`$ $`=`$ $`\left(\begin{array}{ccc}6N+6+\frac{6}{N}& & \frac{1}{2}\frac{1}{N}\\ 24\frac{48}{N}& & 2N+6\frac{2}{N}\end{array}\right),`$ (2.20)
$`\widehat{b}^{\mathrm{SLL}}`$ $`=`$ $`\left(\begin{array}{cccc}0& \frac{1}{2}& 0& 0\\ 8& 8& \frac{1}{2N}& \frac{1}{2}\end{array}\right),`$ (2.23)
where $`N`$ stands for the number of colours.
In order to find the matrix $`\widehat{a}^{12}`$, we need to calculate two-loop diagrams obtained from the ones in fig. 1 by including one-loop corrections on the gluon lines or adding another gluon that couples to the open quark lines. Of course, one-loop diagrams with counterterm insertions need to be included, too. All the two-loop diagrams and the corresponding colour factors are the same as in fig. 2 and table 2 of ref. . However, in the present article, we also consider additional Dirac structures (1.3) in the four-quark vertices.
Inserting the calculated matrix $`\widehat{a}^{12}`$ into eq. (2.14), we obtain the two-loop ADM. Its entries are found to be the following:
$$\begin{array}{ccc}\hfill \gamma _{11}^{(1)\mathrm{SLL}}& =& \frac{203}{6}N^2+\frac{107}{3}N+\frac{136}{3}\frac{12}{N}\frac{107}{2N^2}+\frac{10}{3}Nf\frac{2}{3}f\frac{10}{3N}f,\hfill \\ \hfill \gamma _{12}^{(1)\mathrm{SLL}}& =& \frac{1}{36}N\frac{31}{9}+\frac{9}{N}\frac{4}{N^2}\frac{1}{18}f+\frac{1}{9N}f,\hfill \\ \hfill \gamma _{21}^{(1)\mathrm{SLL}}& =& \frac{364}{3}N\frac{704}{3}\frac{208}{N}\frac{320}{N^2}+\frac{136}{3}f+\frac{176}{3N}f,\hfill \\ \hfill \gamma _{22}^{(1)\mathrm{SLL}}& =& \frac{343}{18}N^2+21N\frac{188}{9}+\frac{44}{N}+\frac{21}{2N^2}\frac{26}{9}Nf6f+\frac{2}{9N}f,\hfill \end{array}$$
(2.24)
where $`f`$ stands for the number of active flavours. The above equation is one of the main results of the present paper.
Proceeding analogously in the VLL sector, we reproduce the well-known results for the one- and two-loop anomalous dimensions of the operator $`Q_1^{\mathrm{VLL}}`$ :
$$\begin{array}{ccc}\hfill \gamma ^{(0)\mathrm{VLL}}& =& 6\frac{6}{N},\hfill \\ \hfill \gamma ^{(1)\mathrm{VLL}}& =& \frac{19}{6}N\frac{22}{3}+\frac{39}{N}\frac{57}{2N^2}+\frac{2}{3}f\frac{2}{3N}f.\hfill \end{array}$$
(2.25)
The matrix $`\widehat{b}`$ in the VLL sector reads
$$\widehat{b}^{\mathrm{VLL}}=\left(\begin{array}{ccc}5& \frac{1}{2N}& \frac{1}{2}\end{array}\right).$$
(2.26)
Finally, our results for the LR sector read
$`\widehat{\gamma }^{(0)\mathrm{LR}}`$ $`=`$ $`\left(\begin{array}{ccc}\frac{6}{N}& & 12\\ 0& & 6N+\frac{6}{N}\end{array}\right),`$ (2.29)
$`\widehat{\gamma }^{(1)\mathrm{LR}}`$ $`=`$ $`\left(\begin{array}{ccc}\frac{137}{6}+\frac{15}{2N^2}\frac{22}{3N}f& & \frac{200}{3}N\frac{6}{N}\frac{44}{3}f\\ \frac{71}{4}N+\frac{9}{N}2f& & \frac{203}{6}N^2+\frac{479}{6}+\frac{15}{2N^2}+\frac{10}{3}Nf\frac{22}{3N}f\end{array}\right),`$ (2.32)
$`\widehat{b}^{\mathrm{LR}}`$ $`=`$ $`\left(\begin{array}{cccccc}0& 5& \frac{1}{2N}& \frac{1}{2}& 0& 0\\ 0& 0& 0& 0& \frac{1}{2N}& \frac{1}{2}\end{array}\right).`$ (2.35)
As mentioned in the introduction, all the comparisons with existing literature are relegated to section 7.
## 3 Current–current contributions to the ADM <br>of $`\mathrm{\Delta }F=1`$ operators
In the present section, we evaluate contributions from the current–current diagrams to the ADM of $`\mathrm{\Delta }F=1`$ operators. For this purpose, we choose the operators in such a manner that all the four flavours they contain are different: $`\overline{s}`$, $`d`$, $`\overline{u}`$, $`c`$. In such a case, the only possible diagrams are the current–current ones.
Twenty linearly independent operators can be built out of four different quark fields. They can be split into 8 separate sectors, between which there is no mixing. The operators belonging to the first four sectors (VLL, VLR, SLR and SLL) read
$`Q_1^{\mathrm{VLL}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu P_Ld^\beta )(\overline{u}^\beta \gamma ^\mu P_Lc^\alpha )=\stackrel{~}{Q}_{V_LV_L},`$
$`Q_2^{\mathrm{VLL}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu P_Ld^\alpha )(\overline{u}^\beta \gamma ^\mu P_Lc^\beta )=Q_{V_LV_L},`$
$`Q_1^{\mathrm{VLR}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu P_Ld^\beta )(\overline{u}^\beta \gamma ^\mu P_Rc^\alpha )=\stackrel{~}{Q}_{V_LV_R},`$
$`Q_2^{\mathrm{VLR}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu P_Ld^\alpha )(\overline{u}^\beta \gamma ^\mu P_Rc^\beta )=Q_{V_LV_R},`$
$`Q_1^{\mathrm{SLR}}`$ $`=`$ $`(\overline{s}^\alpha P_Ld^\beta )(\overline{u}^\beta P_Rc^\alpha )=\stackrel{~}{Q}_{LR},`$
$`Q_2^{\mathrm{SLR}}`$ $`=`$ $`(\overline{s}^\alpha P_Ld^\alpha )(\overline{u}^\beta P_Rc^\beta )=Q_{LR},`$
$`Q_1^{\mathrm{SLL}}`$ $`=`$ $`(\overline{s}^\alpha P_Ld^\beta )(\overline{u}^\beta P_Lc^\alpha )=\stackrel{~}{Q}_{LL},`$
$`Q_2^{\mathrm{SLL}}`$ $`=`$ $`(\overline{s}^\alpha P_Ld^\alpha )(\overline{u}^\beta P_Lc^\beta )=Q_{LL},`$
$`Q_3^{\mathrm{SLL}}`$ $`=`$ $`(\overline{s}^\alpha \sigma _{\mu \nu }P_Ld^\beta )(\overline{u}^\beta \sigma ^{\mu \nu }P_Lc^\alpha )=\stackrel{~}{Q}_{T_LT_L},`$
$`Q_4^{\mathrm{SLL}}`$ $`=`$ $`(\overline{s}^\alpha \sigma _{\mu \nu }P_Ld^\alpha )(\overline{u}^\beta \sigma ^{\mu \nu }P_Lc^\beta )=Q_{T_LT_L},`$ (3.1)
where on the r.h.s. we have shown the notation of ref. .
The operators belonging to the four remaining sectors (VRR, VRL, SRL and SRR) are obtained from the above by interchanging $`P_L`$ and $`P_R`$. Obviously, it is sufficient to calculate the ADMs only for the VLL, VLR, SLR and SLL sectors. The “mirror” operators in the VRR, VRL, SRL and SRR sectors will have exactly the same properties under QCD renormalization.
The evanescent operators for the VLL, VLR, SLR and SLL sectors are listed in appendix B. Calculation of the renormalization constants and the ADMs proceeds along the same lines as in the previous section. The relevant divergences in one- and two-loop diagrams in the cases of VLL, VLR and SLR sectors are given in refs. and . For completeness we give in appendix E the corresponding results for the SLL sector. These have not been published so far in the NDR–$`\overline{\mathrm{MS}}`$ scheme.
Our final results for the $`\mathrm{\Delta }F=1`$ ADMs are as follows:
$`\widehat{\gamma }^{(0)\mathrm{VLL}}`$ $`=`$ $`\left(\begin{array}{ccc}\frac{6}{N}& & 6\\ 6& & \frac{6}{N}\end{array}\right),`$ (3.4)
$`\widehat{\gamma }^{(1)\mathrm{VLL}}`$ $`=`$ $`\left(\begin{array}{ccc}\frac{22}{3}\frac{57}{2N^2}\frac{2}{3N}f& & \frac{19}{6}N+\frac{39}{N}+\frac{2}{3}f\\ \frac{19}{6}N+\frac{39}{N}+\frac{2}{3}f& & \frac{22}{3}\frac{57}{2N^2}\frac{2}{3N}f\end{array}\right),`$ (3.7)
$`\widehat{\gamma }^{(0)\mathrm{VLR}}`$ $`=`$ $`\left(\begin{array}{ccc}6N+\frac{6}{N}& & 0\\ 6& & \frac{6}{N}\end{array}\right),`$ (3.10)
$`\widehat{\gamma }^{(1)\mathrm{VLR}}`$ $`=`$ $`\left(\begin{array}{ccc}\frac{203}{6}N^2+\frac{479}{6}+\frac{15}{2N^2}+\frac{10}{3}Nf\frac{22}{3N}f& & \frac{71}{2}N\frac{18}{N}+4f\\ \frac{100}{3}N+\frac{3}{N}+\frac{22}{3}f& & \frac{137}{6}+\frac{15}{2N^2}\frac{22}{3N}f\end{array}\right),`$ (3.13)
$`\widehat{\gamma }^{(0)\mathrm{SLR}}`$ $`=`$ $`\left(\begin{array}{ccc}\frac{6}{N}& & 6\\ 0& & 6N+\frac{6}{N}\end{array}\right),`$ (3.16)
$`\widehat{\gamma }^{(1)\mathrm{SLR}}`$ $`=`$ $`\left(\begin{array}{ccc}\frac{137}{6}+\frac{15}{2N^2}\frac{22}{3N}f& & \frac{100}{3}N+\frac{3}{N}+\frac{22}{3}f\\ \frac{71}{2}N\frac{18}{N}+4f& & \frac{203}{6}N^2+\frac{479}{6}+\frac{15}{2N^2}+\frac{10}{3}Nf\frac{22}{3N}f\end{array}\right),`$ (3.19)
$`\widehat{\gamma }^{(0)\mathrm{SLL}}`$ $`=`$ $`\left(\begin{array}{cccc}\frac{6}{N}& 6& \frac{N}{2}\frac{1}{N}& \frac{1}{2}\\ 0& 6N+\frac{6}{N}& 1& \frac{1}{N}\\ \frac{48}{N}+24N& 24& \frac{2}{N}4N& 6\\ 48& \frac{48}{N}& 0& 2N\frac{2}{N}\end{array}\right),`$ (3.24)
$`\gamma _{11}^{(1)\mathrm{SLL}}`$ $`=`$ $`\begin{array}{c}\frac{N^2}{2}+\frac{148}{3}\frac{107}{2N^2}2Nf\frac{10}{3N}f,\hfill \end{array}`$ (3.26)
$`\gamma _{12}^{(1)\mathrm{SLL}}`$ $`=`$ $`\begin{array}{c}\frac{178}{3}N+\frac{64}{N}+\frac{16}{3}f,\hfill \end{array}`$ (3.28)
$`\gamma _{13}^{(1)\mathrm{SLL}}`$ $`=`$ $`\begin{array}{c}\frac{107}{36}N^2\frac{71}{18}\frac{4}{N^2}\frac{1}{18}Nf+\frac{f}{9N},\hfill \end{array}`$ (3.30)
$`\gamma _{14}^{(1)\mathrm{SLL}}`$ $`=`$ $`\begin{array}{c}\frac{109}{36}N+\frac{8}{N}\frac{f}{18},\hfill \end{array}`$ (3.32)
$`\gamma _{21}^{(1)\mathrm{SLL}}`$ $`=`$ $`\begin{array}{c}26N+\frac{104}{N},\hfill \end{array}`$ (3.34)
$`\gamma _{22}^{(1)\mathrm{SLL}}`$ $`=`$ $`\begin{array}{c}\frac{203}{6}N^2+\frac{28}{3}\frac{107}{2N^2}+\frac{10}{3}Nf\frac{10}{3N}f,\hfill \end{array}`$ (3.36)
$`\gamma _{23}^{(1)\mathrm{SLL}}`$ $`=`$ $`\begin{array}{c}\frac{89}{18}N+\frac{2}{N}\frac{1}{9}f,\hfill \end{array}`$ (3.38)
$`\gamma _{24}^{(1)\mathrm{SLL}}`$ $`=`$ $`\begin{array}{c}\frac{53}{18}\frac{4}{N^2}+\frac{1}{9N}f,\hfill \end{array}`$ (3.40)
$`\gamma _{31}^{(1)\mathrm{SLL}}`$ $`=`$ $`\begin{array}{c}\frac{676}{3}N^2\frac{1880}{3}\frac{320}{N^2}\frac{88}{3}Nf+\frac{176}{3N}f,\hfill \end{array}`$ (3.42)
$`\gamma _{32}^{(1)\mathrm{SLL}}`$ $`=`$ $`\begin{array}{c}\frac{820}{3}N+\frac{448}{N}\frac{88}{3}f,\hfill \end{array}`$ (3.44)
$`\gamma _{33}^{(1)\mathrm{SLL}}`$ $`=`$ $`\begin{array}{c}\frac{257}{18}N^2\frac{116}{9}+\frac{21}{2N^2}+\frac{22}{9}Nf+\frac{2}{9N}f,\hfill \end{array}`$ (3.46)
$`\gamma _{34}^{(1)\mathrm{SLL}}`$ $`=`$ $`\begin{array}{c}\frac{50}{3}N\frac{8}{3}f,\hfill \end{array}`$ (3.48)
$`\gamma _{41}^{(1)\mathrm{SLL}}`$ $`=`$ $`\begin{array}{c}\frac{488}{3}N+\frac{416}{N}\frac{176}{3}f,\hfill \end{array}`$ (3.50)
$`\gamma _{42}^{(1)\mathrm{SLL}}`$ $`=`$ $`\begin{array}{c}\frac{776}{3}\frac{320}{N^2}+\frac{176}{3N}f,\hfill \end{array}`$ (3.52)
$`\gamma _{43}^{(1)\mathrm{SLL}}`$ $`=`$ $`\begin{array}{c}\frac{22}{3}N\frac{40}{N}+\frac{8}{3}f,\hfill \end{array}`$ (3.54)
$`\gamma _{44}^{(1)\mathrm{SLL}}`$ $`=`$ $`\begin{array}{c}\frac{343}{18}N^2+\frac{28}{9}+\frac{21}{2N^2}\frac{26}{9}Nf+\frac{2}{9N}f.\hfill \end{array}`$ (3.56)
Equation (3.56) is one of the main results of this work.
The careful reader has already noticed that the following equalities hold up to $`𝒪`$$`(\alpha _s^2)`$:
$$\begin{array}{ccccc}\gamma _{11}^{\mathrm{VLL}}=\gamma _{22}^{\mathrm{VLL}},& & \gamma _{12}^{\mathrm{VLL}}=\gamma _{21}^{\mathrm{VLL}},& & \gamma _{11}^{\mathrm{VLR}}=\gamma _{22}^{\mathrm{SLR}},\\ \gamma _{22}^{\mathrm{VLR}}=\gamma _{11}^{\mathrm{SLR}},& & \gamma _{12}^{\mathrm{VLR}}=\gamma _{21}^{\mathrm{SLR}},& & \gamma _{21}^{\mathrm{VLR}}=\gamma _{12}^{\mathrm{SLR}}.\end{array}$$
(3.57)
At one loop, these equalities are a consequence of the Fierz identities
$`(\gamma _\mu P_L)_{ij}(\gamma ^\mu P_L)_{kl}`$ $`=`$ $`(\gamma _\mu P_L)_{il}(\gamma ^\mu P_L)_{kj},`$ (3.58)
$`(\gamma _\mu P_L)_{ij}(\gamma ^\mu P_R)_{kl}`$ $`=`$ $`2(P_R)_{il}(P_L)_{kj},`$ (3.59)
as well as the flavour- and chirality-blind character of QCD interactions. Since the Fierz identities are satisfied in four spacetime dimensions only, the relations (3.57) could be potentially broken at two loops in the NDR–$`\overline{\mathrm{MS}}`$ scheme. Surprisingly, they are not.<sup>1</sup><sup>1</sup>1 In section 4, where the penguin diagrams are considered, no invariance under Fierz rearrangement is observed at two loops for the operators with VLL Dirac structure. A detailed discussion of this fact can be found in ref. .
On the contrary, analogous relations are broken at two loops in the SLL sector. Because of the Fierz relations (2.5), the one-loop matrix $`\widehat{\gamma }^{(0)\mathrm{SLL}}`$ must satisfy the following identity (cf. eqs. (9) and (10) of ref. ):
$$\widehat{\gamma }^{(0)\mathrm{SLL}}=\widehat{}\widehat{\gamma }^{(0)\mathrm{SLL}}\widehat{}$$
(3.60)
with
$$\widehat{}=\left(\begin{array}{cccc}0& \frac{1}{2}& 0& \frac{1}{8}\\ \frac{1}{2}& 0& \frac{1}{8}& 0\\ 0& 6& 0& \frac{1}{2}\\ 6& 0& \frac{1}{2}& 0\end{array}\right).$$
(3.61)
No similar relation holds for $`\widehat{\gamma }^{(1)\mathrm{SLL}}`$ in the NDR–$`\overline{\mathrm{MS}}`$ scheme. As it has already been said, this is not surprising, because the Fierz identities are not true in $`D4`$ dimensions.
It is unclear to us whether the symmetries (3.57) for the VLL, VLR and SLR sectors are preserved at two loops in the NDR–$`\overline{\mathrm{MS}}`$ scheme only by coincidence, or if there is some reason beyond this. As we shall see in section 6, this question is related to the properties of one-loop matrix elements of certain evanescent operators.
## 4 Penguin contributions to the ADM <br>of $`\mathrm{\Delta }F=1`$ operators
In the present section, we shall describe additional contributions to the ADM of $`\mathrm{\Delta }F=1`$ operators that are due to penguin diagrams. Such contributions may arise only when the operators contain one quark-antiquark pair of the same flavour.
For definiteness, let us consider $`\mathrm{\Delta }S=1`$ operators. In the SM analysis of ref. , 10 such operators were considered<sup>2</sup><sup>2</sup>2 Our operators here differ from the ones in ref. by a global normalization factor of 4. Of course, it does not affect their ADM. The factor of 4 can be absorbed into the global normalization factor of the effective Lagrangian, as the first ratio on the r.h.s. of eq. (2.10). In this case, the Wilson coefficients of our operators are exactly the same as those in ref. .
$`Q_1`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu P_Lu^\beta )(\overline{u}^\beta \gamma ^\mu P_Ld^\alpha ),`$
$`Q_2`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu P_Lu^\alpha )(\overline{u}^\beta \gamma ^\mu P_Ld^\beta ),`$
$`Q_3`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu P_Ld^\alpha ){\displaystyle \underset{q}{}}(\overline{q}^\beta \gamma ^\mu P_Lq^\beta ),`$
$`Q_4`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu P_Ld^\beta ){\displaystyle \underset{q}{}}(\overline{q}^\beta \gamma ^\mu P_Lq^\alpha ),`$
$`Q_5`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu P_Ld^\alpha ){\displaystyle \underset{q}{}}(\overline{q}^\beta \gamma ^\mu P_Rq^\beta ),`$
$`Q_6`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu P_Ld^\beta ){\displaystyle \underset{q}{}}(\overline{q}^\beta \gamma ^\mu P_Rq^\alpha ),`$
$`Q_7`$ $`=`$ $`{\displaystyle \frac{3}{2}}(\overline{s}^\alpha \gamma _\mu P_Ld^\alpha ){\displaystyle \underset{q}{}}e_q(\overline{q}^\beta \gamma ^\mu P_Rq^\beta ),`$
$`Q_8`$ $`=`$ $`{\displaystyle \frac{3}{2}}(\overline{s}^\alpha \gamma _\mu P_Ld^\beta ){\displaystyle \underset{q}{}}e_q(\overline{q}^\beta \gamma ^\mu P_Rq^\alpha ),`$
$`Q_9`$ $`=`$ $`{\displaystyle \frac{3}{2}}(\overline{s}^\alpha \gamma _\mu P_Ld^\alpha ){\displaystyle \underset{q}{}}e_q(\overline{q}^\beta \gamma ^\mu P_Lq^\beta ),`$
$`Q_{10}`$ $`=`$ $`{\displaystyle \frac{3}{2}}(\overline{s}^\alpha \gamma _\mu P_Ld^\beta ){\displaystyle \underset{q}{}}e_q(\overline{q}^\beta \gamma ^\mu P_Lq^\alpha ).`$ (4.1)
Their one- and two-loop ADMs, including current–current and penguin diagrams, can be found in appendices A and B of ref. . They were also obtained in ref. . The same results hold for the mirror copies of the SM operators, i.e. for the operators obtained from the ones in eq. (4.1) by $`P_LP_R`$ interchange.
Beyond SM, new linearly independent operators appear. Their Dirac structures are as in eq. (3.1). Our aim is to find a minimal set of linearly independent new operators. In the process of identifying these operators, we shall use four-dimensional Dirac algebra, including the Fierz relations (2.5), (3.58) and (3.59). It turns out that only 3 additional operators (and their mirror copies) undergo mixing via penguin diagrams into other four-quark operators in eq. (4.1). These are
$`Q_{11}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu P_Ld^\alpha )\left[(\overline{d}^\beta \gamma ^\mu P_Ld^\beta )+(\overline{s}^\beta \gamma ^\mu P_Ls^\beta )\right],`$
$`Q_{12}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu P_Ld^\beta )\left[(\overline{d}^\beta \gamma ^\mu P_Rd^\alpha )+(\overline{s}^\beta \gamma ^\mu P_Rs^\alpha )\right],`$
$`Q_{13}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu P_Ld^\alpha )\left[(\overline{d}^\beta \gamma ^\mu P_Rd^\beta )+(\overline{s}^\beta \gamma ^\mu P_Rs^\beta )\right].`$ (4.2)
The remaining elements of the operator basis can be chosen in such a manner that massless penguin diagrams with their insertions vanish. The first three of the remaining operators have the structure of $`Q_{11}`$, …, $`Q_{13}`$, but with a relative minus sign between the two terms. The next two have the structure of $`Q_5`$ and $`Q_6`$, but the sum over flavour-conserving currents is replaced by a difference between the analogous $`u`$-quark and $`c`$-quark currents. Their mirror copies have to be included, as well. Further operators have the SLL and SRR Dirac structures as in eq. (1.3), or they have the form
$`(\overline{s}^\alpha P_{L,R}d^\beta )(\overline{q}^\beta P_{R,L}q^\alpha ),`$ (4.3)
$`(\overline{s}^\alpha P_{L,R}d^\alpha )(\overline{q}^\beta P_{R,L}q^\beta )`$ (4.4)
where $`q`$ has flavour different from $`s`$ and $`d`$. It is straightforward to convince oneself that we have not missed any linearly independent $`\mathrm{\Delta }S=1`$ operator in the above considerations.
Massless penguin diagrams with insertions of the operators (1.3), (4.3) and (4.4) vanish, because
$$\mathrm{Tr}(S_{\mathrm{odd}}P_{L,R})=0\mathrm{and}P_{L,R}S_{\mathrm{odd}}P_{L,R}=0,$$
(4.5)
where $`S_{\mathrm{odd}}`$ is a product of an odd number of Dirac $`\gamma `$-matrices. For dimensional reasons, only massless penguin diagrams can cause mixing into other four-quark operators. This means that all the $`\mathrm{\Delta }S=1`$ operators, except for $`Q_1`$, …, $`Q_{13}`$ and their mirror copies, mix only due to current–current diagrams, i.e. their ADMs are identical to the ones we have already calculated in sections 2 and 3.
At the two-loop level, a complication arises because generally the Fierz relations could be broken in $`D4`$ dimensions. Consequently, our use of these relations in the identification of linearly independent operators could be put in question. However, as we have already discussed in section 2 and will elaborate in section 5, this complication can be avoided by introducing appropriate evanescent operators that vanish in four dimensions by Fierz identities. This allows us to restrict the basis of new physical operators (undergoing penguin mixing) to the one in eq. (4.2), even at the two-loop level.
The introduction of evanescent operators that vanish in four dimensions by Fierz identities turns out to have no effect on the two-loop ADM in the case of the operators with VLR and SLR structures, because the Fierz identity (3.59) remains valid at two loops in the NDR–$`\overline{\mathrm{MS}}`$ scheme, even if the penguin insertions are considered . On the other hand, as pointed out in ref. , the Fierz identity (3.58) is broken at two loops in the NDR–$`\overline{\mathrm{MS}}`$ scheme through penguin diagrams, although it remains valid for current–current diagrams. As a result, the mixing of the operator
$$Q_{11}^{}=(\overline{s}^\alpha \gamma _\mu P_Ld^\beta )\left[(\overline{d}^\beta \gamma ^\mu P_Ld^\alpha )+(\overline{s}^\beta \gamma ^\mu P_Ls^\alpha )\right]$$
(4.6)
with the operators in eq. (4.1), through penguin diagrams, differs from the one of $`Q_{11}`$ at the two-loop level. This can be easily verified by using the results of ref. . As $`Q_{11}^{}=Q_{11}`$ in $`D=4`$ dimensions due to the Fierz identity (3.58), $`Q_{11}^{}`$ was not included in the basis (4.2). By working with $`Q_{11}`$ and the evanescent operator $`Q_{11}^{}Q_{11}`$, the explicit appearance of $`Q_{11}^{}`$ can be avoided at any number of loops, so that the basis (4.2) remains unchanged.
The above discussion implies that the only additional ADMs we need to find in the present section are:
* The 3$`\times `$3 matrix $`\widehat{\gamma }_{cc}`$ describing the mixing of $`Q_{11}`$, …, $`Q_{13}`$ among themselves.
* The 3$`\times `$4 matrix $`\widehat{\gamma }_p`$ describing the mixing of $`Q_{11}`$, …, $`Q_{13}`$ into $`Q_3`$, …, $`Q_6`$ via penguin diagrams. (Only $`Q_3`$, …, $`Q_6`$ are generated by massless QCD penguin diagrams with four-quark operator insertions.)
The matrix $`\widehat{\gamma }_{cc}`$ is given by current–current diagrams only. It takes the form
$$\widehat{\gamma }_{cc}=\left(\begin{array}{cc}\gamma _{\mathrm{\Delta }F=2}^{\mathrm{VLL}}& 0\\ 0& \widehat{\gamma }_{\mathrm{\Delta }F=1}^{\mathrm{VLR}}\end{array}\right)$$
(4.7)
with $`\gamma _{\mathrm{\Delta }F=2}^{\mathrm{VLL}}`$ and $`\widehat{\gamma }_{\mathrm{\Delta }F=1}^{\mathrm{VLR}}`$ taken from eqs. (2.25), (3.4) and (3.7).
The matrix $`\widehat{\gamma }_p=\widehat{\gamma }_p^{(0)}+{\displaystyle \frac{\alpha _s}{4\pi }}\widehat{\gamma }_p^{(1)}+\mathrm{}`$ that originates from penguin diagrams can be easily extracted from sections 3.2 and 5.3 of ref. . We find<sup>3</sup><sup>3</sup>3 The four integers marked in red in eq. (4.14) have been corrected with respect to the first (v1) arXiv version of the current paper. They are now in agreement with ref. where a mistake in our original determination of $`\widehat{\gamma }_p^{(1)}`$ was pointed out. In the current version of the article, we include an extra appendix F where our extraction of $`\widehat{\gamma }_p^{(0)}`$ and $`\widehat{\gamma }_p^{(1)}`$ from the results of ref. is outlined.
$`\widehat{\gamma }_p^{(0)}`$ $`=`$ $`\begin{array}{c}(\frac{4}{3},\frac{4}{3},0)^T\times (\frac{1}{N},1,\frac{1}{N},1),\end{array}`$ (4.9)
$`\widehat{\gamma }_p^{(1)T}`$ $`=`$ $`\left(\begin{array}{ccccc}6N\frac{64}{27}\frac{\textcolor[rgb]{1,0,0}{4}}{3N}+\frac{172}{27N^2}& & \frac{112}{27}\frac{356}{27N^2}& & 6N+\frac{40}{3N}\\ \frac{352}{27}N\frac{\textcolor[rgb]{1,0,0}{14}}{3}\frac{460}{27N}& & \frac{32}{27}N+\frac{500}{27N}& & \frac{22}{3}\\ 6N\frac{244}{27}+\frac{\textcolor[rgb]{1,0,0}{32}}{3N}\frac{188}{27N^2}& & \frac{140}{27}+\frac{148}{27N^2}& & 6N+\frac{4}{3N}\\ \frac{172}{27}N\frac{\textcolor[rgb]{1,0,0}{14}}{3}+\frac{260}{27N}& & \frac{220}{27}N\frac{508}{27N}& & \frac{22}{3}\end{array}\right).`$ (4.14)
The above discussion changes very little in the case of $`\mathrm{\Delta }F=1`$ operators, in which $`F`$ is the up-type flavour. Similarly to the $`\mathrm{\Delta }S=1`$ case, all the contributions from penguin diagrams can be easily extracted from ref. .
## 5 Transformation of the Wilson coefficients <br>to the RI scheme
The ADMs calculated in the present work are given in the NDR–$`\overline{\mathrm{MS}}`$ scheme that is most convenient for perturbative calculations. However, after the Wilson coefficients are evolved with the help of RGE (2.11) down to a low energy scale, it might be necessary to transform them to another scheme that is more appropriate for non-perturbative calculations of hadronic matrix elements . One such scheme is the so-called Regularization-Independent (RI) scheme (originally called the MOM scheme) used in ref. . Below, we shall give relations between the NDR–$`\overline{\mathrm{MS}}`$-renormalized and RI-renormalized Wilson coefficients of all the operators considered in sections 2 and 3.
For completeness, we begin with the definition of the RI scheme. For the massless quark propagator, the renormalization condition can be written as
$$\frac{i}{4}\left[\gamma ^\rho \frac{}{p^\rho }S(p)_R^1\right]_{p^2=\mu ^2}=1,$$
(5.1)
where $`\mu `$ is the subtraction scale. A simple one-loop calculation is necessary to verify that the renormalized inverse propagator in the RI scheme reads
$$S(p)_R^1=ip\text{/}\left[1\frac{\alpha _s}{4\pi }C_F\lambda \left(\frac{1}{2}\mathrm{ln}\frac{p^2}{\mu ^2}\right)\right]+𝒪(\alpha _s^2),$$
(5.2)
where $`C_F=\frac{N^21}{2N}`$ and $`\lambda `$ is the gauge-fixing parameter (cf. eq. (2.17)). In dimensional regularization, the corresponding quark wave-function renormalization constant reads
$$Z_q^{\mathrm{RI}}=1\frac{\alpha _s}{4\pi }C_F\lambda \left(\frac{1}{ϵ}\gamma +\mathrm{ln}(4\pi )+\frac{1}{2}\right),$$
(5.3)
provided the subtraction scale $`\mu `$ is identified with the standard $`\overline{\mathrm{MS}}`$ renormalization scale.
Conditions similar to eq. (5.1) are imposed on renormalized matrix elements of the operators (2.1) and (3.1) among four external quarks with the same momentum $`p`$. The quarks are assumed to be massless here. For the $`\mathrm{\Delta }F=2`$ operators, such matrix elements have the following form
$`Q_1^{\mathrm{VLL}}_R`$ $`=`$ $`A_{11}^{\mathrm{VLL}}(p^2)Q_1^{\mathrm{VLL}}_{\mathrm{tree}}+B_{11}^{\mathrm{VLL}}(p^2)p^\mu p_\nu (\overline{s}^\alpha \gamma _\mu P_Ld^\alpha )(\overline{s}^\beta \gamma ^\nu P_Ld^\beta )_{\mathrm{tree}},`$
$`Q_1^{\mathrm{LR}}_R`$ $`=`$ $`A_{11}^{\mathrm{LR}}(p^2)Q_1^{\mathrm{LR}}_{\mathrm{tree}}+B_{11}^{\mathrm{LR}}(p^2)p^\mu p_\nu (\overline{s}^\alpha \gamma _\mu P_Ld^\alpha )(\overline{s}^\beta \gamma ^\nu P_Rd^\beta )_{\mathrm{tree}}`$
$`+`$ $`A_{12}^{\mathrm{LR}}(p^2)Q_2^{\mathrm{LR}}_{\mathrm{tree}}+B_{12}^{\mathrm{LR}}(p^2)p^\mu p_\nu (\overline{s}^\alpha \sigma _{\mu \rho }P_Ld^\alpha )(\overline{s}^\beta \sigma ^{\nu \rho }P_Rd^\beta )_{\mathrm{tree}},`$
$`Q_2^{\mathrm{LR}}_R`$ $`=`$ $`A_{21}^{\mathrm{LR}}(p^2)Q_1^{\mathrm{LR}}_{\mathrm{tree}}+B_{21}^{\mathrm{LR}}(p^2)p^\mu p_\nu (\overline{s}^\alpha \gamma _\mu P_Ld^\alpha )(\overline{s}^\beta \gamma ^\nu P_Rd^\beta )_{\mathrm{tree}}`$
$`+`$ $`A_{22}^{\mathrm{LR}}(p^2)Q_2^{\mathrm{LR}}_{\mathrm{tree}}+B_{22}^{\mathrm{LR}}(p^2)p^\mu p_\nu (\overline{s}^\alpha \sigma _{\mu \rho }P_Ld^\alpha )(\overline{s}^\beta \sigma ^{\nu \rho }P_Rd^\beta )_{\mathrm{tree}},`$
$`Q_1^{\mathrm{SLL}}_R`$ $`=`$ $`A_{11}^{\mathrm{SLL}}(p^2)Q_1^{\mathrm{SLL}}_{\mathrm{tree}}+A_{12}^{\mathrm{SLL}}(p^2)Q_2^{\mathrm{SLL}}_{\mathrm{tree}},`$
$`Q_2^{\mathrm{SLL}}_R`$ $`=`$ $`A_{21}^{\mathrm{SLL}}(p^2)Q_1^{\mathrm{SLL}}_{\mathrm{tree}}+A_{22}^{\mathrm{SLL}}(p^2)Q_2^{\mathrm{SLL}}_{\mathrm{tree}}.`$ (5.4)
The formfactors $`B_{ij}(p^2)`$ originate from UV-finite parts of Feynman diagrams and are scheme-independent. Note that in all the matrix elements multiplied by $`B_{ij}(p^2)`$, only colour-singlet quark currents occur. Colour-octet currents are removed from these terms with the help of the following Fierz identities (which are independent from the ones in eqs. (2.5), (3.58) and (3.59)):
$`(p\text{/}P_L)_{ij}(p\text{/}P_L)_{kl}`$ $`=`$ $`(p\text{/}P_L)_{il}(p\text{/}P_L)_{kj}{\displaystyle \frac{1}{2}}p^2(\gamma _\mu P_L)_{il}(\gamma ^\mu P_L)_{kj},`$ (5.5)
$`(p\text{/}P_L)_{ij}(p\text{/}P_R)_{kl}`$ $`=`$ $`{\displaystyle \frac{1}{2}}p^\mu p_\nu (\sigma _{\mu \rho }P_R)_{il}(\sigma ^{\nu \rho }P_L)_{kj}+{\displaystyle \frac{1}{2}}p^2(P_R)_{il}(P_L)_{kj},`$ (5.6)
$`p^\mu p_\nu (\sigma _{\mu \rho }P_L)_{ij}(\sigma ^{\nu \rho }P_R)_{kl}`$ $`=`$ $`2(p\text{/}P_R)_{il}(p\text{/}P_L)_{kj}{\displaystyle \frac{1}{2}}p^2(\gamma _\mu P_R)_{il}(\gamma ^\mu P_L)_{kj}.`$ (5.7)
No $`B_{ij}`$ formfactors occur in the SLL sector thanks to the four-dimensional identity
$$p^\mu p_\nu (\sigma _{\mu \rho }P_L)_{ij}(\sigma ^{\nu \rho }P_L)_{kl}=\frac{1}{4}p^2(\sigma _{\mu \nu }P_L)_{ij}(\sigma ^{\mu \nu }P_L)_{kl}.$$
(5.8)
The RI renormalization condition reads
$$A_{ij}(\mu ^2)\omega \mu ^2B_{ij}(\mu ^2)=\delta _{ij},$$
(5.9)
with
$$\omega =\{\begin{array}{cc}\frac{1}{4}& \mathrm{for}B_{11}^{\mathrm{VLL}},B_{11}^{\mathrm{LR}}\mathrm{and}B_{21}^{\mathrm{LR}},\hfill \\ 0& \mathrm{otherwise}.\hfill \end{array}$$
(5.10)
The renormalization condition (5.9) can be equivalently written as
$$A_{ij}^{\mathrm{effective}}(p^2=\mu ^2)=\delta _{ij},$$
(5.11)
with $`A_{ij}^{\mathrm{effective}}`$ obtained from eqs. (5.4) by making the following ad hoc replacements
$`(p\text{/}P_L)(p\text{/}P_{L,R})`$ $``$ $`{\displaystyle \frac{1}{4}}p^2(\gamma _\mu P_L)(\gamma ^\mu P_{L,R}),`$
$`p^\mu p_\nu (\sigma _{\mu \rho }P_{L,R})(\sigma ^{\nu \rho }P_{R,L})`$ $``$ $`0.`$ (5.12)
In the case of $`\mathrm{\Delta }F=1`$ operators, the general structure of one-loop matrix elements is similar to that in eq. (5.4), but the number of formfactors is larger, because operators with colour-octet currents are now linearly independent. The matrix elements can be written as
$$Q_i_R=A_{ij}^{\mathrm{effective}}(p^2)Q_j_{\mathrm{tree}}+N_i,$$
(5.13)
where $`N_i`$ vanish under the replacements (5.12). The RI renormalization condition then has the same form as in eq. (5.11).
In each of the sectors, the RI-renormalized Wilson coefficients can be obtained from the NDR-$`\overline{\mathrm{MS}}`$–renormalized ones with the help of the following relation
$$\stackrel{}{C}^{\mathrm{RI}}(\mu )=\left(1\frac{\alpha _s(\mu )}{4\pi }\mathrm{\Delta }\widehat{r}_{\overline{\mathrm{MS}}RI}^T(\mu )\right)\stackrel{}{C}^{\overline{\mathrm{MS}}}(\mu )+𝒪\left(\alpha _s^2\right),$$
(5.14)
where
$$\left[\mathrm{\Delta }r_{\overline{\mathrm{MS}}RI}(\mu )\right]_{ij}=\frac{4\pi }{\alpha _s(\mu )}\left[A_{ij}^{\mathrm{RI}}(p^2)A_{ij}^{\overline{\mathrm{MS}}}(p^2)\right].$$
(5.15)
The above relations can be easily derived from the fact that the renormalized matrix element of the whole effective Hamiltonian is scheme-independent, i.e.
$$\stackrel{}{C}^{\mathrm{RI}}(\mu )\stackrel{}{Q}(\mu ,p^2)^{\mathrm{RI}}=\stackrel{}{C}^{\overline{\mathrm{MS}}}(\mu )\stackrel{}{Q}(\mu ,p^2)^{\overline{\mathrm{MS}}}.$$
(5.16)
Again, the RI subtraction scale and the standard $`\overline{\mathrm{MS}}`$ renormalization scale have been tacitly identified. The external states must be the same in eq. (5.16). Consequently, the RI-scheme renormalization constant (5.3) must be used for external quark lines in $`A_{ij}^{\overline{\mathrm{MS}}}(p^2)`$ that enters into eq. (5.15).
The dependence on $`p^2`$ and the explicit dependence on $`\mu `$ cancels out in $`\mathrm{\Delta }\widehat{r}_{\overline{\mathrm{MS}}RI}`$ (5.15). However, one should not forget that this matrix depends on the gauge-fixing parameter $`\lambda `$ that is, in turn, $`\mu `$-dependent.
Once the RI renormalization conditions have been specified, finding the explicit form of $`\mathrm{\Delta }\widehat{r}_{\overline{\mathrm{MS}}RI}`$ is only a matter of a straightforward one-loop computation. Our results for the $`\mathrm{\Delta }F=2`$ operators are as follows:
$`\mathrm{\Delta }r_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{VLL}}`$ $`=`$ $`\begin{array}{c}7\frac{7}{N}12\mathrm{ln}2+\frac{12\mathrm{ln}2}{N}+\lambda \left(\frac{3}{2}\frac{3}{2N}4\mathrm{ln}2+\frac{4\mathrm{ln}2}{N}\right),\end{array}`$ (5.18)
$`\mathrm{\Delta }\widehat{r}_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{LR}}`$ $`=`$ $`\left(\begin{array}{cc}\frac{2}{N}+\frac{2\mathrm{ln}2}{N}+\lambda \left(\frac{1}{2N}+\frac{2\mathrm{ln}2}{N}\right)& 4+4\mathrm{ln}2+\lambda \left(1+4\mathrm{ln}2\right)\\ 1+\mathrm{ln}2\lambda \left(\frac{1}{2}\mathrm{ln}2\right)& 4N+\frac{2}{N}+\frac{2\mathrm{ln}2}{N}\lambda \left(\frac{3N}{2}\frac{1}{2N}\frac{2\mathrm{ln}2}{N}\right)\end{array}\right),`$ (5.21)
$`\left[\mathrm{\Delta }r_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{SLL}}\right]_{11}`$ $`=`$ $`\begin{array}{c}4N+7+\frac{5}{N}4\mathrm{ln}2+\frac{2\mathrm{ln}2}{N}+\lambda \left(\frac{1}{2}+\frac{1}{2N}\frac{3N}{2}+\frac{2\mathrm{ln}2}{N}\right),\end{array}`$ (5.23)
$`\left[\mathrm{\Delta }r_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{SLL}}\right]_{12}`$ $`=`$ $`\begin{array}{c}\frac{5}{12}\frac{13}{12N}\frac{2\mathrm{ln}2}{3}+\frac{5\mathrm{ln}2}{6N}+\lambda \left(\frac{5}{24}\frac{1}{6N}\frac{\mathrm{ln}2}{3}+\frac{\mathrm{ln}2}{6N}\right),\end{array}`$ (5.25)
$`\left[\mathrm{\Delta }r_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{SLL}}\right]_{21}`$ $`=`$ $`\begin{array}{c}4\frac{12}{N}32\mathrm{ln}2+\frac{40\mathrm{ln}2}{N}\lambda \left(2+\frac{8}{N}+16\mathrm{ln}2\frac{8\mathrm{ln}2}{N}\right),\end{array}`$ (5.27)
$`\left[\mathrm{\Delta }r_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{SLL}}\right]_{22}`$ $`=`$ $`\begin{array}{c}\frac{7}{3}\frac{5}{3N}\frac{28\mathrm{ln}2}{3}+\frac{26\mathrm{ln}2}{3N}+\lambda \left(\frac{N}{2}+\frac{7}{6}\frac{5}{6N}\frac{8\mathrm{ln}2}{3}+\frac{10\mathrm{ln}2}{3N}\right).\end{array}`$
In the $`\mathrm{\Delta }F=1`$ case, we find
$`\mathrm{\Delta }\widehat{r}_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{VLL}}`$ $`=`$ $`\left(\begin{array}{cc}\frac{7}{N}+\frac{12\mathrm{ln}2}{N}\lambda \left(\frac{3}{2N}\frac{4\mathrm{ln}2}{N}\right)& 712\mathrm{ln}2+\lambda \left(\frac{3}{2}4\mathrm{ln}2\right)\\ 712\mathrm{ln}2+\lambda \left(\frac{3}{2}4\mathrm{ln}2\right)& \frac{7}{N}+\frac{12\mathrm{ln}2}{N}\lambda \left(\frac{3}{2N}\frac{4\mathrm{ln}2}{N}\right)\end{array}\right),`$ (5.31)
$`\mathrm{\Delta }\widehat{r}_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{VLR}}`$ $`=`$ $`\left(\begin{array}{cc}4N+\frac{2}{N}+\frac{2\mathrm{ln}2}{N}\lambda \left(\frac{3N}{2}\frac{1}{2N}\frac{2\mathrm{ln}2}{N}\right)& 22\mathrm{ln}2+\lambda \left(12\mathrm{ln}2\right)\\ 22\mathrm{ln}2\lambda \left(\frac{1}{2}+2\mathrm{ln}2\right)& \frac{2}{N}+\frac{2\mathrm{ln}2}{N}+\lambda \left(\frac{1}{2N}+\frac{2\mathrm{ln}2}{N}\right)\end{array}\right),`$ (5.34)
$`\mathrm{\Delta }\widehat{r}_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{SLR}}`$ $`=`$ $`\left(\begin{array}{cc}\frac{2}{N}+\frac{2\mathrm{ln}2}{N}+\lambda \left(\frac{1}{2N}+\frac{2\mathrm{ln}2}{N}\right)& 22\mathrm{ln}2\lambda \left(\frac{1}{2}+2\mathrm{ln}2\right)\\ 22\mathrm{ln}2+\lambda \left(12\mathrm{ln}2\right)& 4N+\frac{2}{N}+\frac{2\mathrm{ln}2}{N}\lambda \left(\frac{3N}{2}\frac{1}{2N}\frac{2\mathrm{ln}2}{N}\right)\end{array}\right),`$ (5.37)
$`\left[\mathrm{\Delta }r_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{SLL}}\right]_{11}`$ $`=`$ $`\begin{array}{c}\frac{3N}{2}+\frac{5}{N}+\frac{2\mathrm{ln}2}{N}+\lambda \left(\frac{1}{2N}+\frac{2\mathrm{ln}2}{N}\right),\end{array}`$ (5.39)
$`\left[\mathrm{\Delta }r_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{SLL}}\right]_{12}`$ $`=`$ $`\begin{array}{c}\frac{7}{2}2\mathrm{ln}2\lambda \left(\frac{1}{2}+2\mathrm{ln}2\right),\end{array}`$ (5.41)
$`\left[\mathrm{\Delta }r_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{SLL}}\right]_{13}`$ $`=`$ $`\begin{array}{c}\frac{N}{2}\frac{13}{12N}+\frac{5\mathrm{ln}2}{6N}+\lambda \left(\frac{N}{8}\frac{1}{6N}+\frac{\mathrm{ln}2}{6N}\right),\end{array}`$ (5.43)
$`\left[\mathrm{\Delta }r_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{SLL}}\right]_{14}`$ $`=`$ $`\begin{array}{c}\frac{7}{12}\frac{5\mathrm{ln}2}{6}+\lambda \left(\frac{1}{24}\frac{\mathrm{ln}2}{6}\right),\end{array}`$ (5.45)
$`\left[\mathrm{\Delta }r_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{SLL}}\right]_{21}`$ $`=`$ $`\begin{array}{c}12\mathrm{ln}2+\lambda \left(12\mathrm{ln}2\right),\end{array}`$ (5.47)
$`\left[\mathrm{\Delta }r_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{SLL}}\right]_{22}`$ $`=`$ $`\begin{array}{c}4N+\frac{5}{N}+\frac{2\mathrm{ln}2}{N}\lambda \left(\frac{3N}{2}\frac{1}{2N}\frac{2\mathrm{ln}2}{N}\right),\end{array}`$ (5.49)
$`\left[\mathrm{\Delta }r_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{SLL}}\right]_{23}`$ $`=`$ $`\begin{array}{c}\frac{13}{12}\frac{5\mathrm{ln}2}{6}+\lambda \left(\frac{1}{6}\frac{\mathrm{ln}2}{6}\right),\end{array}`$ (5.51)
$`\left[\mathrm{\Delta }r_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{SLL}}\right]_{24}`$ $`=`$ $`\begin{array}{c}\frac{13}{12N}+\frac{5\mathrm{ln}2}{6N}\lambda \left(\frac{1}{6N}\frac{\mathrm{ln}2}{6N}\right),\end{array}`$ (5.53)
$`\left[\mathrm{\Delta }r_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{SLL}}\right]_{31}`$ $`=`$ $`\begin{array}{c}4N\frac{12}{N}+\frac{40\mathrm{ln}2}{N}+\lambda \left(6N\frac{8}{N}+\frac{8\mathrm{ln}2}{N}\right),\end{array}`$ (5.55)
$`\left[\mathrm{\Delta }r_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{SLL}}\right]_{32}`$ $`=`$ $`\begin{array}{c}840\mathrm{ln}2+\lambda \left(28\mathrm{ln}2\right),\end{array}`$ (5.57)
$`\left[\mathrm{\Delta }r_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{SLL}}\right]_{33}`$ $`=`$ $`\begin{array}{c}\frac{5N}{2}\frac{5}{3N}+\frac{26\mathrm{ln}2}{3N}\lambda \left(N+\frac{5}{6N}\frac{10\mathrm{ln}2}{3N}\right),\end{array}`$ (5.59)
$`\left[\mathrm{\Delta }r_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{SLL}}\right]_{34}`$ $`=`$ $`\begin{array}{c}\frac{25}{6}\frac{26\mathrm{ln}2}{3}+\lambda \left(\frac{11}{6}\frac{10\mathrm{ln}2}{3}\right),\end{array}`$ (5.61)
$`\left[\mathrm{\Delta }r_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{SLL}}\right]_{41}`$ $`=`$ $`\begin{array}{c}1240\mathrm{ln}2+\lambda \left(88\mathrm{ln}2\right),\end{array}`$ (5.63)
$`\left[\mathrm{\Delta }r_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{SLL}}\right]_{42}`$ $`=`$ $`\begin{array}{c}\frac{12}{N}+\frac{40\mathrm{ln}2}{N}\lambda \left(\frac{8}{N}\frac{8\mathrm{ln}2}{N}\right),\end{array}`$ (5.65)
$`\left[\mathrm{\Delta }r_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{SLL}}\right]_{43}`$ $`=`$ $`\begin{array}{c}\frac{5}{3}\frac{26\mathrm{ln}2}{3}+\lambda \left(\frac{1}{3}\frac{10\mathrm{ln}2}{3}\right),\end{array}`$ (5.67)
$`\left[\mathrm{\Delta }r_{\overline{\mathrm{MS}}\mathrm{RI}}^{\mathrm{SLL}}\right]_{44}`$ $`=`$ $`\begin{array}{c}\frac{5}{3N}+\frac{26\mathrm{ln}2}{3N}+\lambda \left(\frac{N}{2}\frac{5}{6N}+\frac{10\mathrm{ln}2}{3N}\right).\end{array}`$ (5.69)
In section 7, the above results will be used in performing the comparison with ref. .
## 6 Recovering the ADM of $`\mathrm{\Delta }F=2`$ operators <br>from $`\mathrm{\Delta }F=1`$ results
Let us now use our $`\mathrm{\Delta }F=1`$ anomalous dimensions from section 3 to find again the ADM of $`\mathrm{\Delta }F=2`$ operators. This will serve as a cross-check of our findings and as a preparation for the comparison with ref. in section 7.
Starting from eq. (3.1), we shall pass to another operator basis where the operators are either symmetric or antisymmetric under $`dc`$ interchange. Next, the flavours of both quarks and both antiquarks will be set equal. For definiteness, we shall do it first in the SLL sector. The superscript “SLL” will be understood for all the relevant quantities below, and we shall not write it explicitly.
In four spacetime dimensions, passing to the new operator basis would be equivalent to performing a simple linear transformation of the operators. In the framework of dimensional regularization, introducing additional evanescent operators becomes necessary. In the SLL sector, only two evanescent operators were needed in the $`\mathrm{\Delta }F=1`$ calculation (see appendix B). Now, we need to introduce six additional evanescent operators in this sector. They are defined in appendix C.
We begin with a redefinition of the physical operators $`Q_i(i=1,\mathrm{},4)`$ that amounts to adding to them appropriate linear combinations of the evanescent operators $`E_i`$:
$$Q_iQ_i+\underset{k=1}{\overset{8}{}}W_{ik}E_k\left[Q_i\right]_{\mathrm{new}}$$
(6.1)
where
$$\widehat{W}=\left(\begin{array}{cccccccc}0& 0& \frac{1}{2}& 0& \frac{1}{8}& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 6& 0& \frac{1}{2}& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0\end{array}\right).$$
(6.2)
The “new” operators read
$`\left[Q_1\right]_{\mathrm{new}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[Q_1\right]_F+{\displaystyle \frac{1}{8}}\left[Q_3\right]_F,\left[Q_2\right]_{\mathrm{new}}=Q_2,`$ (6.3)
$`\left[Q_3\right]_{\mathrm{new}}`$ $`=`$ $`6\left[Q_1\right]_F+{\displaystyle \frac{1}{2}}\left[Q_3\right]_F,\left[Q_4\right]_{\mathrm{new}}=Q_4,`$ (6.4)
where
$$\left[Q_1\right]_F=\left(\overline{s}^\alpha P_Lc^\alpha \right)\left(\overline{u}^\beta P_Ld^\beta \right),\left[Q_3\right]_F=\left(\overline{s}^\alpha \sigma _{\mu \nu }P_Lc^\alpha \right)\left(\overline{u}^\beta \sigma ^{\mu \nu }P_Ld^\beta \right).$$
(6.5)
In 4 spacetime dimensions, the transformation (6.1) would be equivalent to performing the Fierz rearrangement of $`Q_1`$ and $`Q_3`$, as $`E_k`$ would not contribute. Since the Fierz identities cannot be analytically continued to $`D`$ dimensions, the Fierz rearrangement must be understood in terms of the transformation (6.1), so long as the $`\overline{\mathrm{MS}}`$ scheme is used. The $`\overline{\mathrm{MS}}`$-renormalized one-loop matrix elements of $`Q_1`$ and $`Q_3`$ are affected by this transformation. This means that the renormalization scheme is changed. We pass from one version of the NDR–$`\overline{\mathrm{MS}}`$ scheme to another, even though the evanescent operators remain unchanged.
After the redefinition (6.1), we perform a simple linear transformation of the operators
$$\left[Q_i\right]_{\mathrm{new}}\underset{j=1}{\overset{4}{}}R_{ij}\left[Q_j\right]_{\mathrm{new}}$$
(6.6)
with
$$\widehat{R}=\left(\begin{array}{cccc}\frac{1}{4}& \frac{1}{2}& \frac{1}{16}& 0\\ 3& 0& \frac{1}{4}& \frac{1}{2}\\ \frac{1}{4}& \frac{1}{2}& \frac{1}{16}& 0\\ 3& 0& \frac{1}{4}& \frac{1}{2}\end{array}\right).$$
(6.7)
As one can easily check, our final operator basis is $`\{Q_1^+,Q_2^+,Q_1^{},Q_2^{}\}`$, where
$`Q_1^\pm `$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[(\overline{s}^\alpha P_Ld^\alpha )(\overline{u}^\beta P_Lc^\beta )\pm (\overline{s}^\alpha P_Lc^\alpha )(\overline{u}^\beta P_Ld^\beta )\right],`$
$`Q_2^\pm `$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[(\overline{s}^\alpha \sigma _{\mu \nu }P_Ld^\alpha )(\overline{u}^\beta \sigma ^{\mu \nu }P_Lc^\beta )\pm (\overline{s}^\alpha \sigma _{\mu \nu }P_Lc^\alpha )(\overline{u}^\beta \sigma ^{\mu \nu }P_Ld^\beta )\right].`$ (6.8)
The ADM transforms as follows:
$`\widehat{\gamma }^{(0)}`$ $``$ $`\widehat{R}\widehat{\gamma }^{(0)}\widehat{R}^1,`$ (6.9)
$`\widehat{\gamma }^{(1)}`$ $``$ $`\widehat{R}\left\{\widehat{\gamma }^{(1)}+[\mathrm{\Delta }\widehat{r},\widehat{\gamma }^{(0)}]+2\beta _0\mathrm{\Delta }\widehat{r}\right\}\widehat{R}^1,`$ (6.10)
where $`\beta _0=\frac{11}{3}N\frac{2}{3}f`$. The matrix $`\mathrm{\Delta }\widehat{r}`$ reflects in the usual manner change of the renormalization scheme that follows from eq. (6.1). The explicit form of $`\mathrm{\Delta }\widehat{r}`$ is
$$\mathrm{\Delta }\widehat{r}=\widehat{W}\widehat{c},$$
(6.11)
provided $`\widehat{W}\widehat{e}=0`$. The matrices $`\widehat{c}`$ and $`\widehat{e}`$ are found from one-loop matrix elements of evanescent operators, as in eq. (2.16). The product $`\widehat{W}\widehat{e}`$ indeed vanishes in our case, and
$$\widehat{c}=\left(\begin{array}{cccc}& & & \\ & & & \\ \frac{3}{4}N\frac{5}{N}& \frac{17}{4}& \frac{1}{16}N\frac{1}{4N}& \frac{3}{16}\\ & & & \\ 7N\frac{28}{N}& 21& \frac{7}{4}N+\frac{5}{N}& \frac{13}{4}\\ & & & \\ & & & \\ & & & \end{array}\right).$$
(6.12)
Here, stars denote non-vanishing elements of $`\widehat{c}`$ that are irrelevant for us, since they do not affect the matrix
$$\mathrm{\Delta }\widehat{r}=\widehat{W}\widehat{c}=\left(\begin{array}{cccc}\frac{1}{2}N+\frac{1}{N}& \frac{1}{2}& \frac{1}{4}N\frac{3}{4N}& \frac{1}{2}\\ 0& 0& 0& 0\\ 8N+\frac{44}{N}& 36& \frac{1}{2}N\frac{1}{N}& \frac{1}{2}\\ 0& 0& 0& 0\end{array}\right).$$
(6.13)
After the transformation (6.9, 6.10), the ADM in the basis $`\{Q_1^+,Q_2^+,Q_1^{},Q_2^{}\}`$ is found to have the form
$$\widehat{\gamma }_{4\times 4}=\left(\begin{array}{cc}\widehat{\gamma }_{2\times 2}^+& 0_{2\times 2}\\ 0_{2\times 2}& \widehat{\gamma }_{2\times 2}^{}\end{array}\right),$$
(6.14)
where $`\widehat{\gamma }^\pm =\widehat{\gamma }^{(0)\pm }+{\displaystyle \frac{\alpha _s}{4\pi }}\widehat{\gamma }^{(1)\pm }+\mathrm{}`$,
$$\widehat{\gamma }^{(0)\pm }=\left(\begin{array}{ccc}6N\pm 6+\frac{6}{N}& & \pm \frac{1}{2}\frac{1}{N}\\ 24\frac{48}{N}& & 2N\pm 6\frac{2}{N}\end{array}\right),$$
(6.15)
and
$$\begin{array}{ccc}\hfill \gamma _{11}^{(1)\pm }& =& \frac{203}{6}N^2\pm \frac{107}{3}N+\frac{136}{3}\frac{12}{N}\frac{107}{2N^2}+\frac{10}{3}Nf\frac{2}{3}f\frac{10}{3N}f,\hfill \\ \hfill \gamma _{12}^{(1)\pm }& =& \frac{1}{36}N\frac{31}{9}\pm \frac{9}{N}\frac{4}{N^2}\frac{1}{18}f+\frac{1}{9N}f,\hfill \\ \hfill \gamma _{21}^{(1)\pm }& =& \frac{364}{3}N\frac{704}{3}\frac{208}{N}\frac{320}{N^2}\pm \frac{136}{3}f+\frac{176}{3N}f,\hfill \\ \hfill \gamma _{22}^{(1)\pm }& =& \frac{343}{18}N^2\pm 21N\frac{188}{9}\pm \frac{44}{N}+\frac{21}{2N^2}\frac{26}{9}Nf6f+\frac{2}{9N}f.\hfill \end{array}$$
(6.16)
One can easily verify that the matrix $`\widehat{\gamma }^+`$ is equal to the one we have already found in eqs. (2.20) and (2.24). It must be so, because the operators $`Q_i^+`$ from eq. (6.8) reduce to $`Q_i^{\mathrm{SLL}}`$ from eq. (2.1) when the flavour replacements $`cd`$ and $`\overline{u}\overline{s}`$ are made. Moreover, the evanescent operators listed in appendices B and C can be linearly combined to the ones that are either symmetric or antisymmetric under $`dc`$ interchange. When the flavour replacements $`cd`$ and $`\overline{u}\overline{s}`$ are made, the antisymmetric operators vanish, while the symmetric ones become equal to those in appendix A. Thus, we have shown how to extract the $`\mathrm{\Delta }F=2`$ results from the $`\mathrm{\Delta }F=1`$ ones.
Let us now briefly describe the analogous transformations in the VLL and LR$``$VLR$``$SLR sectors. All the necessary evanescent operators are given in appendices B and C. The relevant matrices $`\widehat{W}`$ and $`\widehat{R}`$ are the following:
$`\widehat{W}^{\mathrm{VLL}}`$ $`=`$ $`\left(\begin{array}{cccccc}0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 0\end{array}\right),`$ (6.19)
$`\widehat{R}^{\mathrm{VLL}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right),`$ (6.22)
$`\widehat{W}_{4\times 12}^{\mathrm{LR}}`$ $``$ $`\left(\begin{array}{cc}\widehat{W}_{2\times 6}^{\mathrm{VLR}}& 0_{2\times 6}\\ 0_{2\times 6}& \widehat{W}_{2\times 6}^{\mathrm{SLR}}\end{array}\right)=\left(\begin{array}{cc}2\widehat{W}^{\mathrm{VLL}}& 0_{2\times 6}\\ 0_{2\times 6}& \frac{1}{2}\widehat{W}^{\mathrm{VLL}}\end{array}\right),`$ (6.27)
$`\widehat{R}^{\mathrm{LR}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\begin{array}{cccc}0& 1& 2& 0\\ \frac{1}{2}& 0& 0& 1\\ 0& 1& 2& 0\\ \frac{1}{2}& 0& 0& 1\end{array}\right).`$ (6.32)
Consequently, the final operator bases are $`\{Q_1^{VLL+},Q_1^{VLL}\}`$ and $`\{Q_1^{LR+},Q_2^{LR+},Q_1^{LR},Q_2^{LR}\}`$, where
$`Q_1^{VLL\pm }`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[(\overline{s}^\alpha \gamma _\mu P_Ld^\alpha )(\overline{u}^\beta \gamma ^\mu P_Lc^\beta )\pm (\overline{s}^\alpha \gamma _\mu P_Lc^\alpha )(\overline{u}^\beta \gamma ^\mu P_Ld^\beta )\right],`$
$`Q_1^{LR\pm }`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[(\overline{s}^\alpha \gamma _\mu P_Ld^\alpha )(\overline{u}^\beta \gamma ^\mu P_Rc^\beta )\pm (\overline{s}^\alpha \gamma _\mu P_Rc^\alpha )(\overline{u}^\beta \gamma ^\mu P_Ld^\beta )\right],`$
$`Q_2^{LR\pm }`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[(\overline{s}^\alpha P_Ld^\alpha )(\overline{u}^\beta P_Rc^\beta )\pm (\overline{s}^\alpha P_Rc^\alpha )(\overline{u}^\beta P_Ld^\beta )\right].`$ (6.33)
An important simplification in the present case is that the one-loop matrix elements of the evanescent operators $`E_3^{\mathrm{VLL}}`$, $`E_3^{\mathrm{VLR}}`$ and $`E_3^{\mathrm{SLR}}`$ from appendix C vanish in the limit $`D4`$, after subtraction of the MS-counterterms proportional to evanescent operators only. This means that the third rows of $`\widehat{c}^{\mathrm{VLL}}`$, $`\widehat{c}^{\mathrm{VLR}}`$ and $`\widehat{c}^{\mathrm{SLR}}`$ vanish (cf. eq. (2.16)). Consequently, $`\mathrm{\Delta }\widehat{r}^{\mathrm{VLL}}=\widehat{W}^{\mathrm{VLL}}\widehat{c}^{\mathrm{VLL}}=0`$ and $`\mathrm{\Delta }\widehat{r}^{\mathrm{LR}}=\widehat{W}^{\mathrm{LR}}\widehat{c}^{\mathrm{LR}}=0`$. This is why the two-loop $`\mathrm{\Delta }F=1`$ matrices of the VLL, VLR and SLR sectors exhibited Fierz symmetry in eq. (3.57). The transformations of the two-loop ADMs in the VLL and LR sectors thus look as if we worked in 4 dimensions, i.e. they reduce to simple multiplications by the corresponding $`\widehat{R}`$-matrices and their inversions. The final results are
$$\widehat{\gamma }_{2\times 2}^{\mathrm{VLL}}=\left(\begin{array}{cc}\widehat{\gamma }^{VLL+}& 0\\ 0& \widehat{\gamma }^{VLL}\end{array}\right),\widehat{\gamma }_{4\times 4}^{\mathrm{LR}}=\left(\begin{array}{cc}\widehat{\gamma }_{2\times 2}^{LR+}& 0_{2\times 2}\\ 0_{2\times 2}& \widehat{\gamma }_{2\times 2}^{LR}\end{array}\right),$$
(6.34)
where
$`\gamma ^{(0)VLL\pm }`$ $`=`$ $`\pm 6{\displaystyle \frac{6}{N}},`$
$`\gamma ^{(1)VLL\pm }`$ $`=`$ $`{\displaystyle \frac{19}{6}}N{\displaystyle \frac{22}{3}}\pm {\displaystyle \frac{39}{N}}{\displaystyle \frac{57}{2N^2}}\pm {\displaystyle \frac{2}{3}}f{\displaystyle \frac{2}{3N}}f,`$
$`\widehat{\gamma }^{(0)LR\pm }`$ $`=`$ $`\left(\begin{array}{ccc}\frac{6}{N}& & \pm 12\\ 0& & 6N+\frac{6}{N}\end{array}\right),`$ (6.37)
$`\widehat{\gamma }^{(1)LR\pm }`$ $`=`$ $`\left(\begin{array}{ccc}\frac{137}{6}+\frac{15}{2N^2}\frac{22}{3N}f& & \pm \frac{200}{3}N\frac{6}{N}\frac{44}{3}f\\ \pm \frac{71}{4}N\pm \frac{9}{N}2f& & \frac{203}{6}N^2+\frac{479}{6}+\frac{15}{2N^2}+\frac{10}{3}Nf\frac{22}{3N}f\end{array}\right).`$ (6.40)
One can see that $`\gamma ^{VLL+}`$ and $`\widehat{\gamma }^{LR+}`$ are identical to our $`\mathrm{\Delta }F=2`$ results in eqs. (2.25), (2.29) and (2.32).
## 7 Comparison with previous ADM calculations
In the present section, we compare our findings from sections 2, 3 and 4 with the previously published results for anomalous-dimension matrices.
### 7.1 One-loop results
As far as the one-loop QCD ADMs of four-quark operators are concerned, the historical order of their evaluation was as follows:
* Current–current contributions to the one-loop ADM of $`\mathrm{\Delta }F=1`$ operators belonging to the VLL and VLR sectors were originally calculated in refs. . These results were also immediately applicable to the SLR sector, because the Fierz rearrangement has a trivial effect at one loop. For the same reason, one-loop anomalous dimensions of the $`\mathrm{\Delta }F=2`$ operators belonging to the VLL and LR sectors could have been immediately read off from these articles. Thus, after 1974, the only unpublished one-loop current–current anomalous dimensions were those of the SLL sector, both in the $`\mathrm{\Delta }F=1`$ and $`\mathrm{\Delta }F=2`$ cases.
* One-loop penguin contributions to the ADM of the Standard Model operators were originally evaluated in refs. . As we have shown in section 4, penguin contributions to the ADM of other (beyond–SM) flavour-changing dimension-six operators can be easily extracted from the SM calculations, both at one and at two loops.
* To our knowledge, the first published results for $`\gamma ^{(0)\mathrm{SLL}}`$ occur in refs. and , for the $`\mathrm{\Delta }F=2`$ and $`\mathrm{\Delta }F=1`$ cases, respectively.
The one-loop ADMs given in the present article agree with all the papers quoted above. However, in order to perform comparisons, one often needs to make simple linear transformations, because different operator bases are used by different authors. For instance, the results for $`\widehat{\gamma }^{(0)\mathrm{SLL}}`$ in ref. are given in the basis $`\{Q_1^{\mathrm{SLL}},\stackrel{~}{Q}_1^{\mathrm{SLL}}\}`$. In order to compare them with our eq. (2.20), one should use the relation (2.6). Similarly, eqs. (6.7) and (6.9) need to be used for comparing our $`\widehat{\gamma }^{(0)\mathrm{SLL}}`$ in eq. (3.24) with the corresponding results in ref. .
### 7.2 Two-loop results
The history of previous two-loop computations is as follows:
* The current–current anomalous dimensions of the $`\mathrm{\Delta }F=1`$ operators belonging to the VLL sector were originally calculated in ref. (in the DRED–$`\overline{\mathrm{MS}}`$ scheme), and confirmed in ref. (where the NDR–$`\overline{\mathrm{MS}}`$ and HV–$`\overline{\mathrm{MS}}`$ results were also given).
* The remaining elements of the two-loop QCD ADM for $`\mathrm{\Delta }F=1`$ operators relevant in the SM were calculated in refs. . New results in these papers were the current–current contributions in the VLR sector, as well as all the penguin contributions. The SLR sector results in the $`\mathrm{\Delta }F=1`$ case, as well as the $`\mathrm{\Delta }F=2`$ results for the VLL and LR sectors could be easily derived from them with the help of Fierz identities, because the NDR–$`\overline{\mathrm{MS}}`$-renormalized one-loop matrix elements remain invariant under Fierz transformations, except for the current–current ones in the SLL sector, and the penguin ones in the VLL sector. Therefore, in the early 1990’s, the only unknown two-loop anomalous dimensions were those of the SLL sector.
* The first calculation of the two-loop ADM in the SLL sector was performed by Ciuchini et al. , in both the $`\mathrm{\Delta }F=1`$ and $`\mathrm{\Delta }F=2`$ cases. The ADM was calculated there in the so-called “FRI” renormalization scheme. The transformation rules were given to the LRI scheme (Landau-gauge RI-scheme) and to the NDR–$`\overline{\mathrm{MS}}`$ scheme. Current–current anomalous dimensions for the remaining sectors were recalculated as well.
* Penguin contributions to the ADM of non-SM operators are considered for the first time in the present article.
All the two-loop results presented here agree with the previous calculations mentioned above, except for the NDR–$`\overline{\mathrm{MS}}`$ ones for the SLL sector found in ref. . Below, we explain the reason for this disagreement.
### 7.3 Comparison with ref.
In ref. , the two-loop ADM for $`\mathrm{\Delta }F=1`$ operators of the SLL sector was given in the basis defined in eq. (13) of that paper, which is equivalent to our eq. (6.8). It was presented in the so-called “FRI” renormalization scheme, and the transformation rules to the NDR–$`\overline{\mathrm{MS}}`$ scheme were appended. Applying these transformation rules to their “FRI”-scheme ADM, one obtains results that differ from our eq. (6.16). In particular, a mixing between $`Q_i^{}`$ and $`Q_i^+`$ occurs, which is absent in our result (6.16). We could obtain their result if we ignored the transformation (6.1) and, consequently, used $`\mathrm{\Delta }\widehat{r}=0`$ in our eq. (6.10). However, the final results would then correspond to the basis
$`Q_{}^{}{}_{1}{}^{\pm }`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[(\overline{s}^\alpha P_Ld^\alpha )(\overline{u}^\beta P_Lc^\beta ){\displaystyle \frac{1}{2}}(\overline{s}^\alpha P_Ld^\beta )(\overline{u}^\beta P_Lc^\alpha )\pm {\displaystyle \frac{1}{8}}(\overline{s}^\alpha \sigma _{\mu \nu }P_Ld^\beta )(\overline{u}^\beta \sigma ^{\mu \nu }P_Lc^\alpha )\right],`$
$`Q_{}^{}{}_{2}{}^{\pm }`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[(\overline{s}^\alpha \sigma _{\mu \nu }P_Ld^\alpha )(\overline{u}^\beta \sigma ^{\mu \nu }P_Lc^\beta )\pm 6(\overline{s}^\alpha P_Ld^\beta )(\overline{u}^\beta P_Lc^\alpha )\pm {\displaystyle \frac{1}{2}}(\overline{s}^\alpha \sigma _{\mu \nu }P_Ld^\beta )(\overline{u}^\beta \sigma ^{\mu \nu }P_Lc^\alpha )\right],`$
rather than the one in eq. (6.8). In 4 spacetime dimensions, the operators (6.8) and (LABEL:unfierzed) are identical, thanks to the Fierz identities (2.5). However, in $`D`$ dimensions they are not. Consequently, their NDR–$`\overline{\mathrm{MS}}`$-renormalized matrix elements differ at one loop, and it is not surprising that the two-loop ADM depends on which of the two bases is used.
We informed the authors of ref. about our findings prior to publication of the present article. They responded that although their NDR–$`\overline{\mathrm{MS}}`$ results had been claimed to correspond to the basis (6.8), the NDR–$`\overline{\mathrm{MS}}`$ renormalization conditions had been actually imposed in the basis (LABEL:unfierzed). However, they had forgotten to mention this in their article. Unfortunately, such a mistake in the presentation has the same effect on the final result as a mistake in the calculation that amounts to missing $`\mathrm{\Delta }\widehat{r}0`$ in eq. (6.10).
As far as the two-loop ADM for $`\mathrm{\Delta }F=2`$ operators of the SLL sector is concerned, the situation is as follows. If we made the flavour replacements $`cd`$ and $`\overline{u}\overline{s}`$ in the basis (LABEL:unfierzed), but did not change anything in the ADM, we could interpret this ADM as the one for $`\mathrm{\Delta }F=2`$ operators, as the authors of ref. did. However, it would correspond to quite non-standard conventions for the treatment of the evanescent operators obtained from $`Q_{}^{}{}_{1}{}^{}`$ and $`Q_{}^{}{}_{2}{}^{}`$ after the flavour replacements. One would need to assume that the finite one-loop matrix elements of these evanescent operators are not renormalized away, contrary to the usual procedure for any evanescent operator . Such non-standard conventions make the RGE evolution more complicated, because one has to deal with a $`4\times 4`$ instead of a $`2\times 2`$ ADM in the NDR–$`\overline{\mathrm{MS}}`$ RGE for the SLL sector, in the $`\mathrm{\Delta }F=2`$ case. The calculation of the one-loop matrix elements becomes more involved, as well.
In the $`\mathrm{\Delta }F=2`$ case, no calculation is necessary to convince oneself that the results of ref. cannot correspond to the NDR–$`\overline{\mathrm{MS}}`$ renormalization conditions imposed in the basis (6.8) (their eq. (13)). Once the $`cd`$ and $`\overline{u}\overline{s}`$ replacements have been made, the operators $`Q_i^{}`$ in eq. (6.8) vanish identically in $`D`$ dimensions. Therefore, they cannot mix into the $`Q_i^+`$ operators, independently of what the treatment of evanescent operators is. On the other hand, mixing of $`Q_i^{}`$ into $`Q_i^+`$ was claimed to be found in the NDR–$`\overline{\mathrm{MS}}`$ scheme in ref. . Therefore, an inconsistency is clearly seen.
In the remainder of this section, we shall verify that our NDR–$`\overline{\mathrm{MS}}`$ results are compatible with the LRI ones of ref. . By differentiating eq. (5.14) with respect to $`\mu `$, one obtains
$`\widehat{\gamma }_{\mathrm{RI}}^T(\mu )\stackrel{}{C}^{\mathrm{RI}}(\mu )`$ $`=`$ $`[{\displaystyle \frac{\beta _0\alpha _s^2(\mu )}{8\pi ^2}}\mathrm{\Delta }\widehat{r}_{\overline{\mathrm{MS}}RI}^T(\mu )+{\displaystyle \frac{\beta _\lambda ^0\alpha _s(\mu )}{8\pi ^2}}\lambda (\mu )\left({\displaystyle \frac{}{\lambda }}\mathrm{\Delta }\widehat{r}_{\overline{\mathrm{MS}}RI}^T(\mu )\right)`$ (7.2)
$`+(1{\displaystyle \frac{\alpha _s(\mu )}{4\pi }}\mathrm{\Delta }\widehat{r}_{\overline{\mathrm{MS}}RI}^T(\mu ))\widehat{\gamma }_{\overline{\mathrm{MS}}}^T(\mu )]\stackrel{}{C}^{\overline{\mathrm{MS}}}(\mu )+𝒪(\alpha _s^3),`$
where we have used the RGE (2.11),
$$\mu \frac{d}{d\mu }\alpha _s(\mu )=\frac{\beta _0\alpha _s(\mu )^2}{2\pi }+𝒪(\alpha _s^3)\mathrm{and}\mu \frac{d}{d\mu }\lambda (\mu )=\frac{\beta _\lambda ^0\alpha _s(\mu )}{2\pi }\lambda (\mu )+𝒪(\alpha _s^2).$$
(7.3)
We have also used the fact that the dependence of $`\mathrm{\Delta }\widehat{r}_{\overline{\mathrm{MS}}RI}`$ on $`\mu `$ originates solely from its dependence on the gauge-fixing parameter $`\lambda (\mu )`$.
Next, we use eq. (5.14) again to express $`\stackrel{}{C}^{\overline{\mathrm{MS}}}(\mu )`$ by $`\stackrel{}{C}^{\mathrm{RI}}(\mu )`$ in eq. (7.2). Then, the first two terms of the perturbative expansion (2.12) of $`\widehat{\gamma }_{\mathrm{RI}}`$ can be easily read off
$`\widehat{\gamma }_{\mathrm{RI}}^{(0)}`$ $`=`$ $`\widehat{\gamma }_{\overline{\mathrm{MS}}}^{(0)},`$ (7.4)
$`\widehat{\gamma }_{\mathrm{RI}}^{(1)}`$ $`=`$ $`\widehat{\gamma }_{\overline{\mathrm{MS}}}^{(1)}+[\mathrm{\Delta }\widehat{r}_{\overline{\mathrm{MS}}RI},\widehat{\gamma }_{\overline{\mathrm{MS}}}^{(0)}]+2\left(\beta _0+\beta _\lambda ^0\lambda {\displaystyle \frac{}{\lambda }}\right)\mathrm{\Delta }\widehat{r}_{\overline{\mathrm{MS}}RI}.`$ (7.5)
Armed with our explicit expressions for $`\mathrm{\Delta }\widehat{r}_{\overline{\mathrm{MS}}RI}`$ given in section 5 and with the values of
$$\beta _0=\frac{11}{3}N\frac{2}{3}f\mathrm{and}\beta _\lambda ^0=\left(\frac{\lambda }{2}\frac{13}{6}\right)N+\frac{2}{3}f,$$
(7.6)
we can easily calculate the RI-scheme ADM from our $`\overline{\mathrm{MS}}`$ results, for arbitrary $`\lambda `$. Setting then $`\lambda 0`$, we recover all the LRI-scheme anomalous dimensions given in ref. .
As far as the “FRI”-scheme ADMs of ref. are concerned, we can confirm them as well. However, it should be emphasized that the “FRI” scheme is not equivalent to the RI scheme considered in section 5 for any choice of $`\lambda `$. The “FRI” scheme cannot be defined beyond perturbation theory, because different external momenta are chosen in different diagrams when the renormalization conditions are specified. Therefore, in our opinion, the main advantage of the RI scheme is lost.
## 8 Conclusions
In the present paper, we have calculated the two-loop QCD anomalous dimensions matrix (ADM) $`(\widehat{\gamma }^{(1)})_{NDR}`$ in the NDR–$`\overline{\mathrm{MS}}`$ scheme for all the four-fermion dimension-six flavour-changing operators that are relevant to both the Standard Model and its extensions.
The $`\mathrm{\Delta }F=2`$ two-loop results can be found in eqs. (2.24), (2.25) and (2.32). While the matrices in eqs. (2.25) and (2.32) could be extracted from the already published results, the two-loop NDR–$`\overline{\mathrm{MS}}`$ ADM (2.24) for the SLL operators defined in eq. (2.1) is correctly calculated for the first time here.
The $`\mathrm{\Delta }F=1`$ two-loop results for operators containing four different quark flavours can be found in eqs. (3.7), (3.13), (3.19) and (3.56). While the matrices in eqs. (3.7), (3.13) and (3.19) could be extracted from the already published results, the two-loop NDR–$`\overline{\mathrm{MS}}`$ ADM (3.56) for the SLL operators defined in eq. (3.1) is correctly calculated for the first time here.
Penguin contributions to the ADM of non-SM operators have been considered for the first time here. These contributions can be easily extracted from the existing SM calculations. We have identified the relevant non-SM operators in the $`\mathrm{\Delta }S=1`$ case, and presented the corresponding ADM explicitly in eqs. (4.9) and (4.14).
We have demonstrated that the main findings of our paper, given in eqs. (2.24) and (3.56), are compatible with each other, i.e. we have shown how to properly transform the ADMs from the $`\mathrm{\Delta }F=1`$ to the $`\mathrm{\Delta }F=2`$ case. In this context, we have pointed out that in the process of this transformation it is necessary to introduce additional evanescent operators that vanish in four spacetime dimensions because of the Fierz identities.
We have also given the rules that allow transforming our NDR–$`\overline{\mathrm{MS}}`$ ADMs to the corresponding results in the RI scheme, for arbitrary gauge-fixing parameter $`\lambda `$. They can be found in the end of section 5.
The $`\mathrm{\Delta }F=1`$ two-loop ADMs for all the operators defined in eq. (3.1) were previously presented in ref. , in the $`Q_i^\pm `$ basis. In the case of VLL, VLR and SLR operators, there is full agreement between their and our results. The case of SLL operators is more subtle. We can confirm their LRI-scheme results (RI scheme with $`\lambda =0`$). However, their NDR–$`\overline{\mathrm{MS}}`$ ADM is compatible with ours only after correcting their eq. (13), i.e. after changing the definitions of their SLL operators to the ones given in eq. (LABEL:unfierzed).
After such a correction in eq. (13) of ref. , also their $`\mathrm{\Delta }F=2`$ NDR–$`\overline{\mathrm{MS}}`$ results are compatible with ours, provided they are understood in terms of quite non-standard conventions for the treatment of evanescent operators. In their conventions, the two-loop $`\mathrm{\Delta }F=2`$ NDR–$`\overline{\mathrm{MS}}`$ ADM is a 4$`\times `$4 rather than 2$`\times `$2 matrix, which makes the RGE evolution and calculating low-energy matrix elements unnecessarily complicated. Consequently, the results presented here should be more useful for phenomenological applications.
Acknowledgements
We thank the authors of ref. for extensive discussions concerning their paper. Furthermore, we would like to thank Christoph Bobeth and Gerhard Buchalla for carefully reading the manuscript. A.B. and J.U. acknowledge support from the German Bundesministerium für Bildung und Forschung under the contract O5HT9WOA0. M.M. has been supported in part by the Polish Committee for Scientific Research under grant 2 P03B 014 14, 1998-2000.
Appendix A
Here, we specify the evanescent operators that are necessary as counterterms for one-loop diagrams with insertions of the $`\mathrm{\Delta }F=2`$ operators (2.1).
$`E_1^{\mathrm{VLL}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu P_Ld^\beta )(\overline{s}^\beta \gamma ^\mu P_Ld^\alpha )Q_1^{\mathrm{VLL}},`$
$`E_2^{\mathrm{VLL}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu \gamma _\nu \gamma _\rho P_Ld^\alpha )(\overline{s}^\beta \gamma ^\mu \gamma ^\nu \gamma ^\rho P_Ld^\beta )+(16+4ϵ)Q_1^{\mathrm{VLL}},`$
$`E_3^{\mathrm{VLL}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu \gamma _\nu \gamma _\rho P_Ld^\beta )(\overline{s}^\beta \gamma ^\mu \gamma ^\nu \gamma ^\rho P_Ld^\alpha )+(16+4ϵ)Q_1^{\mathrm{VLL}},`$
$`E_1^{\mathrm{LR}}`$ $`=`$ $`(\overline{s}^\alpha P_Ld^\beta )(\overline{s}^\beta P_Rd^\alpha )+{\displaystyle \frac{1}{2}}Q_1^{\mathrm{LR}},`$
$`E_2^{\mathrm{LR}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu P_Ld^\beta )(\overline{s}^\beta \gamma ^\mu P_Rd^\alpha )+2Q_2^{\mathrm{LR}},`$
$`E_3^{\mathrm{LR}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu \gamma _\nu \gamma _\rho P_Ld^\alpha )(\overline{s}^\beta \gamma ^\mu \gamma ^\nu \gamma ^\rho P_Rd^\beta )+(44ϵ)Q_1^{\mathrm{LR}},`$
$`E_4^{\mathrm{LR}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu \gamma _\nu \gamma _\rho P_Ld^\beta )(\overline{s}^\beta \gamma ^\mu \gamma ^\nu \gamma ^\rho P_Rd^\alpha )+(8+8ϵ)Q_2^{\mathrm{LR}},`$
$`E_5^{\mathrm{LR}}`$ $`=`$ $`(\overline{s}^\alpha \sigma _{\mu \nu }P_Ld^\alpha )(\overline{s}^\beta \sigma ^{\mu \nu }P_Rd^\beta )6ϵQ_2^{\mathrm{LR}},`$
$`E_6^{\mathrm{LR}}`$ $`=`$ $`(\overline{s}^\alpha \sigma _{\mu \nu }P_Ld^\beta )(\overline{s}^\beta \sigma ^{\mu \nu }P_Rd^\alpha )+3ϵQ_1^{\mathrm{LR}},`$
$`E_1^{\mathrm{SLL}}`$ $`=`$ $`(\overline{s}^\alpha P_Ld^\beta )(\overline{s}^\beta P_Ld^\alpha )+{\displaystyle \frac{1}{2}}Q_1^{\mathrm{SLL}}{\displaystyle \frac{1}{8}}Q_2^{\mathrm{SLL}},`$
$`E_2^{\mathrm{SLL}}`$ $`=`$ $`(\overline{s}^\alpha \sigma _{\mu \nu }P_Ld^\beta )(\overline{s}^\beta \sigma ^{\mu \nu }P_Ld^\alpha )6Q_1^{\mathrm{SLL}}{\displaystyle \frac{1}{2}}Q_2^{\mathrm{SLL}},`$
$`E_3^{\mathrm{SLL}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu \gamma _\nu \gamma _\rho \gamma _\sigma P_Ld^\alpha )(\overline{s}^\beta \gamma ^\mu \gamma ^\nu \gamma ^\rho \gamma ^\sigma P_Ld^\beta )+(64+96ϵ)Q_1^{\mathrm{SLL}}+(16+8ϵ)Q_2^{\mathrm{SLL}},`$
$`E_4^{\mathrm{SLL}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu \gamma _\nu \gamma _\rho \gamma _\sigma P_Ld^\beta )(\overline{s}^\beta \gamma ^\mu \gamma ^\nu \gamma ^\rho \gamma ^\sigma P_Ld^\alpha )64Q_1^{\mathrm{SLL}}+(16+16ϵ)Q_2^{\mathrm{SLL}}.`$
The evanescent operators for the VRR and SRR sectors, i.e. $`E_k^{VRR}`$ and $`E_k^{SRR}`$ are obtained by replacing $`L`$ by $`R`$ in the definitions of $`E_k^{\mathrm{VLL}}`$ and $`E_k^{\mathrm{SLL}}`$.
The operators $`E_1^{\mathrm{VLL}}`$, $`E_1^{\mathrm{LR}}`$, $`E_2^{\mathrm{LR}}`$, $`E_1^{\mathrm{SLL}}`$ and $`E_2^{\mathrm{SLL}}`$ vanish in four spacetime dimensions because of the Fierz identities (3.58), (3.59) and (2.5). The operators $`E_2^{\mathrm{VLL}}`$, $`E_3^{\mathrm{VLL}}`$, $`E_3^{\mathrm{LR}}`$, $`E_4^{\mathrm{LR}}`$, $`E_3^{\mathrm{SLL}}`$ and $`E_4^{\mathrm{SLL}}`$ vanish by the four-dimensional identity (2.2). Finally, $`E_5^{\mathrm{LR}}`$ and $`E_6^{\mathrm{LR}}`$ vanish in four dimensions, because they become full contractions of self-dual and self-antidual antisymmetric tensors.
The evanescent operators listed here would look somewhat simpler if we removed from them all the terms proportional to $`ϵ`$. It would be equivalent to changing one version of the $`\overline{\mathrm{MS}}`$ scheme to another. Keeping the terms proportional to $`ϵ`$ in the above equations makes our NDR–$`\overline{\mathrm{MS}}`$ scheme equivalent to the one where the so-called “Greek projections” are used (see appendix D).
Appendix B
Here, we specify the evanescent operators that are necessary as counterterms for one-loop diagrams with insertions of the $`\mathrm{\Delta }F=1`$ operators (3.1).
$`E_1^{\mathrm{VLL}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu \gamma _\nu \gamma _\rho P_Ld^\beta )(\overline{u}^\beta \gamma ^\mu \gamma ^\nu \gamma ^\rho P_Lc^\alpha )+(16+4ϵ)Q_1^{\mathrm{VLL}},`$
$`E_2^{\mathrm{VLL}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu \gamma _\nu \gamma _\rho P_Ld^\alpha )(\overline{u}^\beta \gamma ^\mu \gamma ^\nu \gamma ^\rho P_Lc^\beta )+(16+4ϵ)Q_2^{\mathrm{VLL}},`$
$`E_1^{\mathrm{VLR}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu \gamma _\nu \gamma _\rho P_Ld^\beta )(\overline{u}^\beta \gamma ^\mu \gamma ^\nu \gamma ^\rho P_Rc^\alpha )+(44ϵ)Q_1^{\mathrm{VLR}},`$
$`E_2^{\mathrm{VLR}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu \gamma _\nu \gamma _\rho P_Ld^\alpha )(\overline{u}^\beta \gamma ^\mu \gamma ^\nu \gamma ^\rho P_Rc^\beta )+(44ϵ)Q_2^{\mathrm{VLR}},`$
$`E_1^{\mathrm{SLR}}`$ $`=`$ $`(\overline{s}^\alpha \sigma _{\mu \nu }P_Ld^\beta )(\overline{u}^\beta \sigma ^{\mu \nu }P_Rc^\alpha )6ϵQ_1^{\mathrm{SLR}},`$
$`E_2^{\mathrm{SLR}}`$ $`=`$ $`(\overline{s}^\alpha \sigma _{\mu \nu }P_Ld^\alpha )(\overline{u}^\beta \sigma ^{\mu \nu }P_Rc^\beta )6ϵQ_2^{\mathrm{SLR}},`$
$`E_1^{\mathrm{SLL}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu \gamma _\nu \gamma _\rho \gamma _\sigma P_Ld^\beta )(\overline{u}^\beta \gamma ^\mu \gamma ^\nu \gamma ^\rho \gamma ^\sigma P_Lc^\alpha )+(64+96ϵ)Q_1^{\mathrm{SLL}}+(16+8ϵ)Q_3^{\mathrm{SLL}},`$
$`E_2^{\mathrm{SLL}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu \gamma _\nu \gamma _\rho \gamma _\sigma P_Ld^\alpha )(\overline{u}^\beta \gamma ^\mu \gamma ^\nu \gamma ^\rho \gamma ^\sigma P_Lc^\beta )+(64+96ϵ)Q_2^{\mathrm{SLL}}+(16+8ϵ)Q_4^{\mathrm{SLL}}.`$
The remaining evanescent operators (for the VRR, VRL, SRL and SRR sectors) are obtained by interchanging $`L`$ and $`R`$ above.
Appendix C
This appendix contains definitions of the “additional” evanescent operators that are not necessary as one-loop counterterms in the $`\mathrm{\Delta }F=1`$ effective Lagrangian in section 3. However, they have to be included before performing transformation to the “plus–minus” basis in section 6.
$`E_3^{\mathrm{VLL}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu P_Lc^\alpha )(\overline{u}^\beta \gamma ^\mu P_Ld^\beta )Q_1^{\mathrm{VLL}},`$
$`E_4^{\mathrm{VLL}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu P_Lc^\beta )(\overline{u}^\beta \gamma ^\mu P_Ld^\alpha )Q_2^{\mathrm{VLL}},`$
$`E_5^{\mathrm{VLL}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu \gamma _\nu \gamma _\rho P_Lc^\alpha )(\overline{u}^\beta \gamma ^\mu \gamma ^\nu \gamma ^\rho P_Ld^\beta )+(16+4ϵ)Q_1^{\mathrm{VLL}},`$
$`E_6^{\mathrm{VLL}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu \gamma _\nu \gamma _\rho P_Lc^\beta )(\overline{u}^\beta \gamma ^\mu \gamma ^\nu \gamma ^\rho P_Ld^\alpha )+(16+4ϵ)Q_2^{\mathrm{VLL}},`$
$`E_3^{\mathrm{VLR}}`$ $`=`$ $`(\overline{s}^\alpha P_Rc^\alpha )(\overline{u}^\beta P_Ld^\beta )+{\displaystyle \frac{1}{2}}Q_1^{\mathrm{VLR}},`$
$`E_4^{\mathrm{VLR}}`$ $`=`$ $`(\overline{s}^\alpha P_Rc^\beta )(\overline{u}^\beta P_Ld^\alpha )+{\displaystyle \frac{1}{2}}Q_2^{\mathrm{VLR}},`$
$`E_5^{\mathrm{VLR}}`$ $`=`$ $`(\overline{s}^\alpha \sigma _{\mu \nu }P_Rc^\alpha )(\overline{u}^\beta \sigma ^{\mu \nu }P_Ld^\beta )+3ϵQ_1^{\mathrm{VLR}},`$
$`E_6^{\mathrm{VLR}}`$ $`=`$ $`(\overline{s}^\alpha \sigma _{\mu \nu }P_Rc^\beta )(\overline{u}^\beta \sigma ^{\mu \nu }P_Ld^\alpha )+3ϵQ_2^{\mathrm{VLR}},`$
$`E_3^{\mathrm{SLR}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu P_Rc^\alpha )(\overline{u}^\beta \gamma ^\mu P_Ld^\beta )+2Q_1^{\mathrm{SLR}},`$
$`E_4^{\mathrm{SLR}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu P_Rc^\beta )(\overline{u}^\beta \gamma ^\mu P_Ld^\alpha )+2Q_2^{\mathrm{SLR}},`$
$`E_5^{\mathrm{SLR}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu \gamma _\nu \gamma _\rho P_Rc^\alpha )(\overline{u}^\beta \gamma ^\mu \gamma ^\nu \gamma ^\rho P_Ld^\beta )+(8+8ϵ)Q_1^{\mathrm{SLR}},`$
$`E_6^{\mathrm{SLR}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu \gamma _\nu \gamma _\rho P_Rc^\beta )(\overline{u}^\beta \gamma ^\mu \gamma ^\nu \gamma ^\rho P_Ld^\alpha )+(8+8ϵ)Q_2^{\mathrm{SLR}},`$
$`E_3^{\mathrm{SLL}}`$ $`=`$ $`(\overline{s}^\alpha P_Lc^\alpha )(\overline{u}^\beta P_Ld^\beta )+{\displaystyle \frac{1}{2}}Q_1^{\mathrm{SLL}}{\displaystyle \frac{1}{8}}Q_3^{\mathrm{SLL}},`$
$`E_4^{\mathrm{SLL}}`$ $`=`$ $`(\overline{s}^\alpha P_Lc^\beta )(\overline{u}^\beta P_Ld^\alpha )+{\displaystyle \frac{1}{2}}Q_2^{\mathrm{SLL}}{\displaystyle \frac{1}{8}}Q_4^{\mathrm{SLL}},`$
$`E_5^{\mathrm{SLL}}`$ $`=`$ $`(\overline{s}^\alpha \sigma _{\mu \nu }P_Lc^\alpha )(\overline{u}^\beta \sigma ^{\mu \nu }P_Ld^\beta )6Q_1^{\mathrm{SLL}}{\displaystyle \frac{1}{2}}Q_3^{\mathrm{SLL}},`$
$`E_6^{\mathrm{SLL}}`$ $`=`$ $`(\overline{s}^\alpha \sigma _{\mu \nu }P_Lc^\beta )(\overline{u}^\beta \sigma ^{\mu \nu }P_Ld^\alpha )6Q_2^{\mathrm{SLL}}{\displaystyle \frac{1}{2}}Q_4^{\mathrm{SLL}},`$
$`E_7^{\mathrm{SLL}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu \gamma _\nu \gamma _\rho \gamma _\sigma P_Lc^\alpha )(\overline{u}^\beta \gamma ^\mu \gamma ^\nu \gamma ^\rho \gamma ^\sigma P_Ld^\beta )64Q_1^{\mathrm{SLL}}+(16+16ϵ)Q_3^{\mathrm{SLL}},`$
$`E_8^{\mathrm{SLL}}`$ $`=`$ $`(\overline{s}^\alpha \gamma _\mu \gamma _\nu \gamma _\rho \gamma _\sigma P_Lc^\beta )(\overline{u}^\beta \gamma ^\mu \gamma ^\nu \gamma ^\rho \gamma ^\sigma P_Ld^\alpha )64Q_2^{\mathrm{SLL}}+(16+16ϵ)Q_4^{\mathrm{SLL}}.`$
The remaining evanescent operators (for the VRR, VRL, SRL and SRR sectors) are obtained by interchanging $`L`$ and $`R`$ above.
Appendix D
In the present appendix, the notion of “Greek projections” is recalled and generalized to the case of SLL-sector operators. Let us denote the Dirac structure of the operator in eq. (1.1) by $`\mathrm{\Gamma }_A\mathrm{\Gamma }_B`$. The insertion of this operator in one- and two-loop diagrams results in new Dirac structures like
$$\mathrm{\Gamma }_n\mathrm{\Gamma }_A\mathrm{\Gamma }^n\mathrm{\Gamma }_B,$$
(D.1)
where $`\mathrm{\Gamma }_n=\gamma _{\mu _1}\gamma _{\mu _2}\mathrm{}\gamma _{\mu _n}`$. Several examples of such structures occur in appendices A–C. It has been suggested in ref. to project them onto physical operators as follows. One defines the projection $`G`$ so that the following equality is satisfied
$$G\left[\mathrm{\Gamma }_n\mathrm{\Gamma }_A\mathrm{\Gamma }^n\mathrm{\Gamma }_B\right]=\xi G\left[\mathrm{\Gamma }^A\mathrm{\Gamma }^B\right].$$
(D.2)
In the case of $`\mathrm{\Gamma }_A=\mathrm{\Gamma }_B=\gamma _\alpha P_L`$, performing the projection $`G`$ amounts to replacing $``$ by $`\gamma _\tau `$ on both sides of the above equation and contracting the indices using $`D`$-dimensional Dirac algebra. In this manner, the coefficient $`\xi `$ is determined. One finds for instance:
$$G\left[(\overline{s}^\alpha \gamma _\mu \gamma _\nu \gamma _\rho P_Ld^\beta )(\overline{u}^\beta \gamma ^\mu \gamma ^\nu \gamma ^\rho P_Lc^\alpha )\right]=(164ϵ)G\left[Q_1^{\mathrm{VLL}}\right]+𝒪(ϵ^2)$$
(D.3)
with $`Q_1^{\mathrm{VLL}}`$ as defined in eq. (3.1).
It has been pointed out in ref. that for a proper treatment of counterterms in two-loop calculations, one has to use eq. (D.3) only as a prescription for defining an evanescent operator. In the case at hand, this is the operator $`E_1^{\mathrm{VLL}}`$ of appendix B. As discussed in ref. , in the case of VLR and SLR operators, the analogous projections are performed by replacing $``$ by 1 and $`\gamma _\tau `$, respectively. Examples of the corresponding evanescent operators can be found in appendices A–C.
The projections in the SLL sector are slightly more involved. In the case of the insertion of $`Q_1^{\mathrm{SLL}}`$ or $`Q_3^{\mathrm{SLL}}`$, the r.h.s. of eq. (D.2) has to be generalized to a linear combination of these two operators. The same applies to the pair $`(Q_2^{\mathrm{SLL}},Q_4^{\mathrm{SLL}})`$. The projection $`G`$ is now performed by replacing $``$ by $`\gamma _\alpha \gamma _\beta `$. After the projection, one finds linear combinations of $`g_{\alpha \beta }`$ and $`\gamma _\alpha \gamma _\beta `$ on both sides of the equation. This allows extracting the coefficients in question. One finds for instance
$$G\left[(\overline{s}^\alpha \gamma _\mu \gamma _\nu \gamma _\rho \gamma _\sigma P_Ld^\beta )(\overline{u}^\beta \gamma ^\mu \gamma ^\nu \gamma ^\rho \gamma ^\sigma P_Lc^\alpha )\right]=(6496ϵ)G\left[Q_1^{\mathrm{SLL}}\right]+(168ϵ)G\left[Q_3^{\mathrm{SLL}}\right]+𝒪(ϵ^2).$$
(D.4)
The corresponding evanescent operator is $`E_1^{\mathrm{SLL}}`$ in appendix B. An alternative approach to projections can be found in ref. .
Appendix E
In this appendix, the $`1/ϵ`$ and $`1/ϵ^2`$ poles in the one- and two-loop diagrams are given for the $`\mathrm{\Delta }F=1`$ calculation in the SLL sector. Analogous results for the remaining sectors can be found in refs. and . The gauge-fixing parameter $`\lambda `$ is set to unity here, i.e. the Feynman–’t Hooft gauge is used.
Each insertion results in a linear combination of $`Q_1^{\mathrm{SLL}}`$, …, $`Q_4^{\mathrm{SLL}}`$, after subtracting the evanescent counterterms (see appendix B) or, alternatively, after performing the “Greek projections” (see appendix D). Table 1 gives the singularities (without colour factors) in the coefficients of the resulting operators, for each diagram separately. The numbering of the diagrams and values of the colour factors are exactly as in figs. 1, 2 and tables 1, 2 of ref. . The multiplicity factors of the diagrams are included.
In the two-loop case, the singularities include one-loop diagrams with counterterm insertions. The counterterms proportional to evanescent operators are multiplied by an additional factor $`1/2`$, and, at the same time, the term $`2\widehat{b}\widehat{c}`$ in eq. (2.14) is ignored. Correctness of such a trick has been justified in refs. .
The singularities from table 1 apply for the pair $`(Q_2^{\mathrm{SLL}},Q_4^{\mathrm{SLL}})`$, too. After including colour factors and summing the diagrams, the $`1/ϵ`$ singularities build a $`4\times 4`$ matrix in the basis $`\{Q_1^{\mathrm{SLL}},Q_2^{\mathrm{SLL}},Q_3^{\mathrm{SLL}},Q_4^{\mathrm{SLL}}\}`$
$$\widehat{B}=\frac{\alpha _s}{4\pi }\widehat{B}_1+\left(\frac{\alpha _s}{4\pi }\right)^2\widehat{B}_2+𝒪(\alpha _s^3),$$
(E.5)
from which the anomalous-dimension matrix can be obtained by means of
$`\gamma _{ij}^{(0)\mathrm{SLL}}`$ $`=`$ $`2\left[2a_1\delta _{ij}+(\widehat{B}_1)_{ij}\right],`$ (E.6)
$`\gamma _{ij}^{(1)\mathrm{SLL}}`$ $`=`$ $`4\left[2a_2\delta _{ij}+(\widehat{B}_2)_{ij}\right].`$ (E.7)
Here, $`a_1`$ and $`a_2`$ originate from $`1/ϵ`$ singularities in the quark field renormalization constants. They read
$$a_1=C_F,a_2=C_F\left[\frac{3}{4}C_F\frac{17}{4}N+\frac{1}{2}f\right].$$
(E.8)
Remembering the trick applied to evanescent operators here, it is easy to verify that eqs. (E.6) and (E.7) are equivalent to eqs. (2.13) and (2.14) from section 2.
Appendix F
In the present appendix, we outline details of our determination of the penguin-diagram generated ADMs in eqs. (4.9) and (4.14). Let us begin with the SM QCD-penguin operators $`Q_3`$, …$`Q_6`$ that are listed in eq. (4.1). Their $`4\times 4`$ ADM can be split into contributions from the current-current and penguin diagrams: $`\widehat{\gamma }^{\mathrm{SM}}=\widehat{\gamma }_{cc}^{\mathrm{SM}}+\widehat{\gamma }_p^{\mathrm{SM}}`$. The structure of $`Q_3`$, …$`Q_6`$ implies that
$$\widehat{\gamma }_{cc}^{\mathrm{SM}}=\left(\begin{array}{cccc}\gamma _{22}^{\mathrm{VLL}}& \gamma _{21}^{\mathrm{VLL}}& 0& 0\\ \gamma _{12}^{\mathrm{VLL}}& \gamma _{11}^{\mathrm{VLL}}& 0& 0\\ 0& 0& \gamma _{22}^{\mathrm{VLR}}& \gamma _{21}^{\mathrm{VLR}}\\ 0& 0& \gamma _{12}^{\mathrm{VLR}}& \gamma _{11}^{\mathrm{VLR}}\end{array}\right),$$
(F.9)
where $`\widehat{\gamma }^{\mathrm{VLL}}`$ and $`\widehat{\gamma }^{\mathrm{VLR}}`$ up to the two-loop level are given in eqs. (3.2)–(3.5). The matrix $`\widehat{\gamma }_p^{\mathrm{SM}}`$ is most easily found by subtracting $`\widehat{\gamma }_{cc}^{\mathrm{SM}}`$ from the full $`\widehat{\gamma }^{\mathrm{SM}}`$. One- and two-loop contributions to the latter matrix in the NDR–$`\overline{\mathrm{MS}}`$ are listed in table 5 of ref. .<sup>4</sup><sup>4</sup>4 They coincide with results of independent calculations in refs. and . This way one finds
$`\widehat{\gamma }_p^{\mathrm{SM}(0)}`$ $`=`$ $`\begin{array}{c}(\frac{4}{3},\frac{2f}{3},0,\frac{2f}{3})^T\times (\frac{1}{N},1,\frac{1}{N},1),\end{array}`$ (F.11)
$`\widehat{\gamma }_p^{\mathrm{SM}(1)}`$ $`=`$ $`\left(\begin{array}{ccccccc}\frac{64}{27}+\frac{172}{27N^2}& & \frac{460}{27N}+\frac{352N}{27}& & \frac{244}{27}\frac{188}{27N^2}& & \frac{260}{27N}+\frac{172N}{27}\\ \frac{4}{3N}+6N& & \frac{14}{3}& & \frac{32}{3N}6N& & \frac{14}{3}\\ 0& & 0& & 0& & 0\\ 0& & 0& & 0& & 0\end{array}\right)`$ (F.16)
$`+`$ $`\left(\begin{array}{ccccccc}\frac{8}{3N}+3N& & \frac{1}{3}& & \frac{10}{3N}3N& & \frac{1}{3}\\ \frac{20}{27}+\frac{74}{27N^2}& & \frac{164}{27N}+\frac{110N}{27}& & \frac{56}{27}+\frac{2}{27N^2}& & \frac{20}{27N}+\frac{74N}{27}\\ \frac{20}{3N}3N& & \frac{11}{3}& & \frac{2}{3N}+3N& & \frac{11}{3}\\ \frac{56}{27}\frac{178}{27N^2}& & \frac{250}{27N}\frac{16N}{27}& & \frac{70}{27}+\frac{74}{27N^2}& & \frac{254}{27N}+\frac{110N}{27}\end{array}\right)f.`$ (F.21)
Let us now extend the considered operator set to $`\{Q_3,Q_4,Q_5,Q_6,Q_{11},Q_{11}^{},Q_{13},Q_{12}\}`$, where the extra operators have been defined in eqs. (4.2) and (4.6). At this point, we ignore the fact that $`Q_{11}`$ and $`Q_{11}^{}`$ are related by a Fierz relation in $`D=4`$, i.e. we treat both of them as independent normal (non-evanescent) operators.
The four extra operators $`\{Q_{11},Q_{11}^{},Q_{13},Q_{12}\}`$ can be obtained from $`\{Q_3,Q_4,Q_5,Q_6\}`$ by skipping the $`u`$-, $`c`$\- and $`b`$-quarks in the sum over flavours. Consequently, the full $`8\times 8`$ ADM up to two loops takes the form
$$\widehat{\gamma }_{8\times 8}^{}=\left(\begin{array}{cc}\widehat{\gamma }^{\mathrm{SM}}& 0_{4\times 4}\\ \widehat{\gamma }_{p,f=2}^{\mathrm{SM}}& \widehat{\gamma }_{cc}^{\mathrm{SM}}\end{array}\right).$$
(F.22)
We have tacitly assumed here that each of our operators is accompanied by a corresponding one-loop evanescent operator containing triple products of Dirac matrices inside the quark currents, defined in full analogy to $`E_1^{\mathrm{VLL}}`$, $`E_2^{\mathrm{VLL}}`$, $`E_1^{\mathrm{VLR}}`$ and $`E_2^{\mathrm{VLR}}`$ in appendix B, including the coefficients at the $`𝒪(ϵ)`$ terms there.
In the next step, we transform our $`8\times 8`$ ADM to the basis $`\{Q_3,Q_4,Q_5,Q_6,Q_{11},Q_{12},Q_{13},E_{11}\}`$, where $`E_{11}=Q_{11}^{}Q_{11}`$ is still treated as a normal (non-evanescent) operator. The transformed ADM reads
$$\widehat{\gamma }_{8\times 8}^{}=\widehat{R}\widehat{\gamma }_{8\times 8}^{}\widehat{R}^1\text{with}\widehat{R}=\left(\begin{array}{cc}1_{4\times 4}& 0_{4\times 4}\\ 0_{4\times 4}& X_{4\times 4}\end{array}\right)\text{and}X_{4\times 4}=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 1& 0\\ 1& 1& 0& 0\end{array}\right).$$
(F.23)
Finally, we depart from the $`\mathrm{MS}`$ scheme for $`E_{11}`$ (still thinking of it as of a normal operator though) by introducing finite terms in the one-loop renormalization constants that correspond to its mixing via penguin diagrams into $`Q_3`$, …$`Q_6`$. This amounts to replacing $`a_{ik}^{11}`$ in eq. (2.15) by $`a_{ik}^{11}+ϵa_{ik}^{01}`$ when the index $`i`$ corresponds to $`E_{11}`$, and the index $`k`$ corresponds to $`Q_3`$, …$`Q_6`$. We adjust
$$a_{E_{11},Q_3}^{01}=a_{E_{11},Q_5}^{01}=\frac{2}{3N}\text{and}a_{E_{11},Q_4}^{01}=a_{E_{11},Q_6}^{01}=\frac{2}{3}$$
(F.24)
to make the renormalized one-loop penguin matrix element of $`E_{11}`$ vanish. The one-loop ADM remains intact, while the resulting transformation of the two-loop ADM reads (c.f. eq. (6.10))
$$\widehat{\gamma }_{8\times 8}^{\prime \prime (1)}=\widehat{\gamma }^{}{}_{8\times 8}{}^{(1)}+[\mathrm{\Delta }\widehat{r},\widehat{\gamma }_{8\times 8}^{(0)}]+2\beta _0\mathrm{\Delta }\widehat{r}\text{with}(\mathrm{\Delta }r)_{ij}=\delta _{i8}[\frac{2}{3N}(\delta _{j1}+\delta _{j3})+\frac{2}{3}(\delta _{j2}+\delta _{j4})].$$
(F.25)
At this point, our non-$`\mathrm{MS}`$ scheme with $`E_{11}`$ treated as a normal operator becomes equivalent to the $`\mathrm{MS}`$ scheme with $`E_{11}`$ treated as an evanescent operator.<sup>5</sup><sup>5</sup>5 Strictly speaking, this is true only up to finite renormalization in the one-loop mixing of evanescent operators among themselves, which should be absent in the $`\mathrm{MS}`$ scheme. However, such a finite renormalization has no effect on the ADM of the physical operators up to two loops. We explicitly verify that the 8th row of $`\widehat{\gamma }_{8\times 8}^{\prime \prime }`$ up to two loops has only a single non-vanishing entry that corresponds to the mixing of $`E_{11}`$ with itself. Consequently, the Wilson coefficient of $`E_{11}`$ has no effect on the RG evolution of the Wilson coefficients of normal operators, as it should be for any evanescent operator in the $`\mathrm{MS}`$ scheme. Let us note that it would not be the case if the transformation (F.25) was not performed.
To find the actual ADM for the normal operators only (now with $`E_{11}`$ treated as an evanescent operator, and in the $`\mathrm{MS}`$ scheme), we remove the 8th row and 8th column from the matrix $`\widehat{\gamma }_{8\times 8}^{\prime \prime }`$. The resulting $`7\times 7`$ matrix reads
$$\widehat{\gamma }_{7\times 7}^{}=\left(\begin{array}{cc}\widehat{\gamma }^{\mathrm{SM}}& 0_{4\times 3}\\ \widehat{\gamma }_p& \widehat{\gamma }_{cc}\end{array}\right),$$
(F.26)
where the one- and two-loop contributions to $`\widehat{\gamma }^{\mathrm{SM}}`$, $`\widehat{\gamma }_{cc}`$ and $`\widehat{\gamma }_p`$ coincide with what has been already given in eqs. (F.9)-(F.21), (4.7) and (4.9)-(4.14), respectively. Our final results for $`\widehat{\gamma }_p`$ in eqs. (4.9)-(4.14) have actually been extracted from eq. (F.26). As far as $`\widehat{\gamma }_{cc}`$ in eq. (4.7) is concerned, eq. (F.26) gives us a nice confirmation of the result that has in practice been determined using a much simpler method.
The transformation (F.25) was missed in the original version of our paper in 2000. The mistake was pointed out in ref. . The current appendix was added in 2024, simultaneously with implementing the proper correction in eq. (4.14).
The reader might wonder why our finite subtraction in eq. (F.24) was restricted to the penguin matrix elements only, i.e. why no similar operation was necessary for the one-loop current-current matrix element of $`E_{11}`$. It was the case because such a matrix element turns to vanish after subtracting the evanescent counterterms only, and passing to $`D=4`$, as already discussed in section 6, below eq. (6.33).
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# I Introduction
## I Introduction
Testing the QCD dynamics of heavy quarks in various conditions provides us with a qualitative and quantitative knowledge that allows us to distinguish fine complex effects caused by the electroweak nature of CP-violation or physics beyond the Standard Model. The list of hadrons containing the heavy quarks as available to the experimental observations and measurements, was recently enriched by a new member, the long-lived $`B_c`$-meson, in addition to the heavy quarkonia $`\overline{b}b`$ and $`\overline{c}c`$ as well as the mesons and baryons with a single heavy quark. This success of CDF Collaboration in the first observation of $`B_c`$-meson was based on the progress of experimental technique in the reconstruction of rare processes with heavy quarks by use of vertex detectors. This experience supports a hope to observe other rare long-lived doubly heavy hadrons. i.e. the baryons containing two heavy quarks. As expected they have production rate and lifetime similar to the $`B_c`$-meson ones.
In the present paper we investigate the two-point sum rules of NRQCD . The light quark-doubly heavy diquark structure of baryon leads to the definite expressions for baryonic currents written in terms of nonrelativistic heavy quark fields. To relate the nonrelativstic heavy quark correlators to the full QCD ones we need to take into account the hard gluon corrections by means of solving the renormalization group equation known up to the two-loop accuracy.
The convergency of sum rules results are essentially improved by account for a nonzero light quark mass. As was mentioned in the sum rules stability can be achieved by destroying the baryon-diquark factorization in the correlators. The convergency was obtained due to taking into account the nonperturbative interactions caused by higher dimension operators in contrast to ref., where a signigicant instability of results were observed in full QCD sum rules with no product of quark and gluon condensates. We show that for the strange $`\mathrm{\Omega }_{QQ^{}}`$ baryons this factorization is broken already in the perturbative limit which allows us to introduce a new criterion for the determination of baryon masses since we observe the stability of sum rules for the masses obtained from both correlators standing in front of two independent Lorentz structures for the spinor field of $`\mathrm{\Omega }_{QQ^{}}`$.
Moreover our choice of baryonic current is convenient to take into account the $`\alpha _s/v`$ coulomb-like corrections inside the doubly heavy diquark.
In Section II we describe the scheme of calculation. There we define the currents and represent the spectral densities in the NRQCD sum rules for various operators included into the consideration. The numerical estimates in comparison with the values obtained in potential models are given in Section III. The results are summarized in Conclusion.
## II NRQCD sum rules for doubly heavy baryons
### A Description of the method
In order to determine the masses and coupling constants of baryons in sum rules we consider the two point correlators of interpolating baryon currents. The quantum numbers of doubly heavy diquark in the ground states are given by its spin and parity, so that $`j_d^P=1^+`$ or $`j_d^P=0^+`$ (if the identical heavy quarks form the diquark then the scalar state $`j_d^P=0^+`$ is forbidden). Adding the light quark to form the baryon, we obtain the pair of degenerate states $`j^P=\frac{1}{2}^+`$ and $`j^P=\frac{3}{2}^+`$ for the baryonsThe superscript $``$ denotes various electric charges depending on the flavor of the light quark. $`\mathrm{\Xi }_{cc}^{}`$, $`\mathrm{\Xi }_{bc}^{}`$, $`\mathrm{\Xi }_{bb}^{}`$ and $`\mathrm{\Xi }_{cc}^{}`$, $`\mathrm{\Xi }_{bc}^{}`$, $`\mathrm{\Xi }_{bb}^{}`$ with the vector diquark, and $`j^P=\frac{1}{2}^+`$ for the $`\mathrm{\Xi }_{bc}^{}`$ baryons with the scalar diquark. Unlike the case of baryons with a single heavy quark , there is the only independent current component for each ground state. We find
$`J_{\mathrm{\Xi }_{QQ^{}}^{}}`$ $`=`$ $`[Q^{iT}C\tau \gamma _5Q^j]q^k\epsilon _{ijk},`$ (1)
$`J_{\mathrm{\Xi }_{QQ}^{}}`$ $`=`$ $`[Q^{iT}C\tau 𝜸^mQ^j]𝜸_m\gamma _5q^k\epsilon _{ijk},`$ (2)
$`J_{\mathrm{\Xi }_{QQ}^{}}^n`$ $`=`$ $`[Q^{iT}C\tau 𝜸^nQ^j]q^k\epsilon _{ijk}+{\displaystyle \frac{1}{3}}𝜸^n[Q^{iT}C𝜸^mQ^j]𝜸_mq^k\epsilon _{ijk},`$ (3)
where $`J_{\mathrm{\Xi }_{QQ}^{}}^n`$ satisfies the spin-3/2 condition $`𝜸_nJ_{\mathrm{\Xi }_{QQ}^{}}^n=0`$. The flavor matrix $`\tau `$ is antisymmetric for $`\mathrm{\Xi }_{bc}^{}`$ and symmetric for $`\mathrm{\Xi }_{QQ}^{}`$ and $`\mathrm{\Xi }_{QQ}^{}`$. Here $`T`$ means transposition, $`C`$ is the charge conjugation matrix.
The matrix structure of correlator for two baryonic currents with the spin of $`1/2`$ has the form
$$\mathrm{\Pi }(w)=id^4xe^{ipx}0|TJ(x),\overline{J}(0)|0=\text{ / }vF_1(w)+F_2(w),$$
(4)
where $`w`$ is defined by $`p^2=(+w)^2`$, $`=m_Q+m_Q^{}+m_s`$, $`m_{Q,Q^{}}`$ are the heavy quark masses and $`m_s`$ is the strange quark mass. The appropriate definitions of scalar formfactors for the 3/2-spin baryon are given by the following:
$$\mathrm{\Pi }_{\mu \nu }(w)=id^4xe^{ipx}0|T\{J_\mu (x),\overline{J}_\nu (0)\}|0=g_{\mu \nu }[\text{ / }v\stackrel{~}{F}_1(w)+\stackrel{~}{F}_2(w)]+\mathrm{},$$
(5)
where we do not concern for distinct Lorentz structures. The scalar correlators $`F`$ can be evaluated in a deep euclidean region by employing the Operator Product Expansion (OPE) in the framework of NRQCD,
$$F_{1,2}(w)=\underset{d}{}C_d^{(1,2)}(w)O_d,$$
(6)
where $`O_d`$ denotes the local operator with a given dimension $`d`$: $`O_0=\widehat{1}`$, $`O_3=\overline{q}q`$, $`O_4=\frac{\alpha _s}{\pi }G^2`$, …, and the functions $`C_d(w)`$ are the corresponding Wilson coefficients of OPE. For the contribution of quark condensate operator we explore the following OPE up to the terms of the fourth order in $`x`$ (the derivation is presented in Appendix):
$`0|Ts_i^a(x)\overline{s}_j^b(0)|0`$ $`=`$ $`{\displaystyle \frac{1}{12}}\delta ^{ab}\delta _{ij}\overline{s}s`$ (8)
$`\left[1+{\displaystyle \frac{x^2(m_0^22m_s^2)}{16}}+{\displaystyle \frac{x^4(\pi ^2\frac{a_s}{\pi }G^2\frac{3}{2}m_s^2(m_0^2m_s^2))}{288}}\right]`$
$`+`$ $`im_s\delta ^{ab}x_\mu \gamma _{ij}^\mu \overline{q}q\left[{\displaystyle \frac{1}{48}}+{\displaystyle \frac{x^2}{24^2}}\left({\displaystyle \frac{3m_0^2}{4}}m_s^2\right)\right].`$ (9)
Note that at $`m_s0`$ the expansion of quark condensate gives contributions in both correlators in contrast with the sum rules for $`\mathrm{\Xi }_{QQ^{}}`$ , where putting $`m_s=0`$ and neglecting the higher condensates, the authors found the factorization of diquark correlator in $`F_2`$ and full baryonic correlator in $`F_1`$. This fact was the physical reason for the divergency of SR method.
We write down the Wilson coefficient in front of unity and quark-gluon operators by making use of the dispersion relation over $`w`$,
$$C_d(w)=\frac{1}{\pi }_0^{\mathrm{}}\frac{\rho _d(\omega )d\omega }{\omega w},$$
(10)
where $`\rho `$ denotes the imaginary part in the physical region of NRQCD.
To relate the NRQCD correlators to the real hadrons, we use the dispersion representation for the two point function with the physical spectral density given by the appropriate resonance and continuum part. The coupling constants of baryons are defined by the following expressions:
$`0|J(x)|\mathrm{\Xi }(\mathrm{\Omega })_{QQ}^{}(p)`$ $`=`$ $`iZ_{\mathrm{\Xi }(\mathrm{\Omega })_{QQ}^{}}u(v,M_{\mathrm{\Xi }(\mathrm{\Omega })})e^{ipx},`$ (11)
$`0|J^m(x)|\mathrm{\Xi }(\mathrm{\Omega })_{QQ}^{}(p,\lambda )`$ $`=`$ $`iZ_{\mathrm{\Xi }(\mathrm{\Omega })_{QQ}^{}}u^m(v,M_{\mathrm{\Xi }(\mathrm{\Omega })})e^{ipx},`$ (12)
where the spinor field with the four-velocity $`v`$ and mass $`M`$ (the mass of baryon) satisfies the equation $`\text{ / }vu(v,M)=u(v,M)`$, and $`u^m(v,M)`$ denotes the transversal spinor.
Then we use the nonrelativistic expressions for the physical spectral functions
$$\rho _{1,2}^{phys}(\omega )=\frac{M}{2}|Z|^2\delta (\overline{\mathrm{\Lambda }}\omega ),$$
(13)
where we have performed the substitution $`\delta (p^2M^2)\frac{1}{2}\delta (\overline{\mathrm{\Lambda }}w)`$, here $`\overline{\mathrm{\Lambda }}`$ is the binding energy of baryon and $`M=+\overline{\mathrm{\Lambda }}`$. The nonrelativistic dispersion relation for the hadronic part of sum rules has the form
$$\frac{\rho _{1,2}^{phys}d\omega }{\omega w}=\frac{1}{2}\frac{|Z|^2}{\overline{\mathrm{\Lambda }}w}.$$
(14)
We suppose that the continuum densities starting from the threshold $`\omega _{cont}`$, is modelled by the NRQCD expressions. Then, in the sum rules equalizing the correlators calculated in NRQCD and given by the physical states, the integration above $`\omega _{cont}`$ cancel each other in two sides of sum rules. Further, we write down the correlators in the deep underthreshold point of $`w=+t`$ with $`t0`$.
Now the sum rules in the scheme of moments with respect to $`t`$ can be written down as follows:
$$\frac{1}{\pi }_0^{\omega _{cont}}\frac{\rho _{1,2}d\omega }{(\omega +)^n}=\frac{M}{2}\frac{|Z|^2}{M^n},$$
(15)
where $`\rho _j`$ contains the contributions given by various operators in OPE for the corresponding scalar functions $`F_j`$. Introducing the following notation for the $`n`$-th moment of two point correlation function:
$$_n=\frac{1}{\pi }_0^{\omega _{cont}}\frac{\rho (\omega )d\omega }{(\omega +)^{n+1}},$$
(16)
we can obtain the estimates of baryon mass $`M`$, for example, as the following:
$$M[n]=\frac{_n}{_{n+1}},$$
(17)
and the coupling is determined by the expression
$$|Z[n]|^2=\frac{2}{M}_nM^{n+1},$$
(18)
where we see the dependence of sum rule results on the scheme parameter.
### B Calculating the spectral densities
In this subsection we present analytical expression for the perturbative spectral functions in the NRQCD approximation. The evaluation of spectral densities involves the standard use of Cutkosky rules , with some modification motivated by NRQCD
$`\mathrm{heavy}\mathrm{quark}:`$ $`{\displaystyle \frac{1}{p_0(m+\frac{𝐩^2}{2m})}}2\pi i\delta (p_0(m+{\displaystyle \frac{𝐩^2}{2m}})),`$ (19)
$`\mathrm{light}\mathrm{quark}:`$ $`{\displaystyle \frac{1}{p^2m^2}}2\pi i\delta (p^2m^2).`$ (20)
We derive the spin symmetry relations for all the spectral densities due to the fact that in the leading order of the heavy quark effective theory the spins of heavy quarks are decoupled, so
$$\rho _{1,\mathrm{\Omega }_{QQ}^{}}=3\rho _{1,\mathrm{\Omega }_{QQ^{}}^{{}_{}{}^{}}}=3\rho _{1,\mathrm{\Omega }_{QQ}^{}},$$
(21)
$$\rho _{2,\mathrm{\Omega }_{QQ}^{}}=3\rho _{2,\mathrm{\Omega }_{QQ^{}}^{{}_{}{}^{}}}=3\rho _{2,\mathrm{\Omega }_{QQ}^{}},$$
(22)
and we have the following relation for the baryon couplings in NRQCD:
$$|Z_\mathrm{\Omega }|^2=3|Z_\mathrm{\Omega }^{}|^2=3|Z_\mathrm{\Omega }^{}|^2.$$
(23)
Using the smallness of the strange quark mass we use the following expansions in $`m_s`$ for the perturbative spectral densities standing in front of unity operator ($`m_{QQ^{}}=m_Qm_Q^{}/(m_Q+m_Q^{})`$ is the reduced diquark mass, $`_{diq}=m_Q+m_Q^{}`$):
$$\rho _{1,\mathrm{\Omega }_{QQ^{}}^{{}_{}{}^{}}}(\omega )=\frac{\sqrt{2}(m_{QQ^{}}\omega )^{3/2}}{15015\pi ^3(_{diq}+\omega )^3}(\eta _{1,0}(\omega )+m_s\eta _{1,1}(\omega )+m_s^2\eta _{1,2}(\omega )),$$
(24)
where we have found
$`\eta _{1,0}(\omega )`$ $`=`$ $`16\omega ^2(429_{diq}^3+715_{diq}^2\omega +403_{diq}\omega ^2+77\omega ^3),`$ (25)
$`\eta _{1,1}(\omega )`$ $`=`$ $`104\omega (231_{diq}^3+297_{diq}^2\omega +121_{diq}\omega ^2+15\omega ^3),`$ (26)
$`\eta _{1,2}(\omega )`$ $`=`$ $`{\displaystyle \frac{10}{(_{diq}+\omega )^2}}(3003_{diq}^5+9009_{diq}^4\omega +9438_{diq}^3\omega ^2`$ (28)
$`+4290_{diq}^2\omega ^3+871_{diq}\omega ^4+77\omega ^5).`$
The first term of this expansion reproduces the result obtained in . A new feature is the appearance of nonzero perturbative $`\rho _{2,\mathrm{\Omega }_{QQ^{}}^{{}_{}{}^{}}}`$ density which is proportional to $`m_s`$,
$$\rho _{2,\mathrm{\Omega }_{QQ^{}}^{{}_{}{}^{}}}(\omega )=\frac{2\sqrt{2}\omega (m_{QQ^{}}\omega )^{3/2}m_s}{105\pi ^3(_{diq}+\omega )^2}(\eta _{2,0}+m_s\eta _{2,1}+m_s^2\eta _{2,2}),$$
(29)
and
$`\eta _{2,0}`$ $`=`$ $`42\omega (_{diq}^2+48_{diq}\omega +14\omega ^2),`$ (30)
$`\eta _{2,1}`$ $`=`$ $`3(35_{diq}^2+28_{diq}\omega +5\omega ^2),`$ (31)
$`\eta _{2,2}`$ $`=`$ $`{\displaystyle \frac{1}{(_{diq}+\omega )^2}}(105_{diq}^3+315_{diq}^2\omega +279_{diq}\omega ^2+77\omega ^3).`$ (32)
The account for the coulomb-like interaction leads to the finite renormalization of the diquark spectral densities before the integration over the diquark invariant mass by the introduction of Sommerfeld factor C, so that
$$\rho _{diquark}^𝐂=\rho _{diquark}^{bare}𝐂$$
(33)
with
$$𝐂=\frac{2\pi \alpha _s}{3v_{QQ^{}}}\left[1\mathrm{exp}\left(\frac{2\pi \alpha _s}{3v_{QQ^{}}}\right)\right]^1,$$
(34)
where $`v_{12}`$ denotes the relative velocity of heavy quarks inside the diquark, and we have taken into account the color anti-triplet structure of diquark. The relative velocity is given by the following expression:
$$v_{QQ^{}}=\sqrt{1\frac{4m_Qm_Q^{}}{Q^2(m_Qm_Q^{})^2}},$$
(35)
where $`Q^2`$ is the square of heavy diquark four-momentum. In NRQCD we take the limit of low velocities, so that denoting the diquark invariant mass squared by $`Q^2=(_{diq}+ϵ)^2`$, we find
$$𝐂=\frac{2\pi \alpha _s}{3v_{QQ^{}}},v_{QQ^{}}^2=\frac{ϵ}{2m_{QQ^{}}},$$
at $`ϵm_{QQ^{}}`$.
The corrected spectral densities are equal to
$$\rho _1^𝐂(\omega )=\frac{m_{QQ^{}}^2\alpha _s\omega (2_{diq}+\omega )}{6\pi ^2(_{diq}+\omega )^3}(\eta _{1,0}^𝐂+m_s\eta _{1,1}^𝐂+m_s^2\eta _{1,2}^𝐂),$$
(36)
where
$`\eta _{1,0}^𝐂`$ $`=`$ $`(2_{diq}+\omega )^2\omega ^2,`$ (37)
$`\eta _{1,1}^𝐂`$ $`=`$ $`{\displaystyle \frac{3(2_{diq}+\omega )\omega }{(_{diq}+\omega )}}(4_{diq}^3+6_{diq}^2\omega +4_{diq}\omega ^2+\omega ^3),`$ (38)
$`\eta _{1,2}^𝐂`$ $`=`$ $`{\displaystyle \frac{1}{(_{diq}+\omega )^2}}(12_{diq}^4+24_{diq}^3\omega +32_{diq}^2\omega ^2+20_{diq}\omega ^3+5\omega ^4).`$ (39)
We see that the first term again reproduces the result of . For the $`\rho _{2,\mathrm{\Omega }_{QQ^{}}^{{}_{}{}^{}}}^𝐂`$ we find
$$\rho _2^𝐂=\frac{m_sm_{QQ^{}}^2(2_{diq}+\omega )\omega \alpha _s}{2\pi (_{diq}+\omega )^2}(\eta _{2,0}^𝐂+m_s\eta _{2,1}^𝐂+m_s^2\eta _{2,2}^𝐂),$$
(40)
$`\eta _{2,0}^𝐂`$ $`=`$ $`(2_{diq}+\omega )\omega ,`$ (41)
$`\eta _{2,1}^𝐂`$ $`=`$ $`{\displaystyle \frac{2}{_{diq}+\omega }}(2_{diq}^2+2_{diq}\omega +\omega ^2),`$ (42)
$`\eta _{2,0}^𝐂`$ $`=`$ $`{\displaystyle \frac{2}{(_{diq}+\omega )^2}}(2_{diq}^2+2_{diq}\omega +\omega ^2).`$ (43)
The use of these expansions numerically leads to very small deviations from the exact integral representations of spectral densities (about 0.5%), but they are more convenient in calculations.
The contribution to the moments of the spectral densities determined by the light quark condensate can be calculated by the exploration of (9)
$`_{\overline{q}q}^{(1)}(n)`$ $`=`$ $`{\displaystyle \frac{(n+1)!}{n!}}𝒫_1^{diq}(n+1)+{\displaystyle \frac{(n+3)!}{n!}}𝒫_3^{diq}(n+3)`$ (44)
$`_{\overline{q}q}^{(2)}(n)`$ $`=`$ $`𝒫_0^{diq}(n){\displaystyle \frac{(n+2)!}{n!}}𝒫_2^{diq}(n+2)+{\displaystyle \frac{(n+4)!}{n!}}𝒫_4^{diq}(n+4),`$ (45)
where we have introduced the coefficients of expansion in $`x`$ by $`𝒫_i`$ (see (9) and Appendix). The n-th moment of two-point correlator function of diquark is denoted by $`^{diq}(n)`$. Then the diquark spectral density takes the following form:
$$\rho _{diq}=\frac{12\sqrt{2}m_{QQ^{}}^{3/2}\sqrt{\omega }}{\pi },$$
(46)
which has to be multiplied by the Sommerfeld factor $`𝐂`$, wherein the variable $`ϵ`$ is substituted by $`\omega `$, since in this case there is no integration over the quark-diquark invariant mass. This corrected density is
$$\rho _{diq}^𝐂=\frac{48\pi \alpha _sm_{QQ^{}}^2}{3},$$
(47)
and it is independent of $`\omega `$.
The corrections due to the gluon condensate are given by the density
$$\rho _1^{G^2}(\omega )=\frac{(m_Q^2+m_Q^{}^2+11m_Qm_Q^{})m_{QQ^{}}^{5/2}\sqrt{\omega }}{212^{10}\sqrt{2}\pi m_Q^2m_Q^{}^2(_{diq}+\omega )^2}(\eta _{1,0}^{G^2}+m_s\eta _{1,1}^{G^2}+m_s^2\eta _{1,2}^{G^2}),$$
(48)
with
$`\eta _{1,0}^{G^2}`$ $`=`$ $`84_{diq}^3+196_{diq}^2\omega +133_{diq}\omega ^2+11\omega ^3,`$ (49)
$`\eta _{1,1}^{G^2}`$ $`=`$ $`{\displaystyle \frac{2(210_{diq}^3+70_{diq}^2\omega +21_{diq}\omega ^2+3\omega ^3)}{_{diq}+\omega }},`$ (50)
$`\eta _{1,2}^{G^2}`$ $`=`$ $`{\displaystyle \frac{2(210_{diq}^3+70_{diq}^2\omega +21_{diq}\omega ^2+3\omega ^3)}{(_{diq}+\omega )^2}},`$ (51)
where we again make the expansion in $`m_s`$. In the case of nonzero quark mass we get the nonzero density proportional to $`m_s`$,
$$\rho _2^{G^2}(\omega )=\frac{m_s(m_Q^2+m_Q^{}^2+11m_Qm_Q^{})m_{QQ^{}}^{5/2}\sqrt{\omega }}{32^9\sqrt{2}\pi m_Q^2m_Q^{}^2(_{diq}+\omega )}(\eta _{2,0}^{G^2}+m_s\eta _{2,1}^{G^2}),$$
(52)
with
$`\eta _{2,0}^{G^2}`$ $`=`$ $`(9_{diq}+\omega ),`$ (53)
$`\eta _{2,1}^{G^2}`$ $`=`$ $`{\displaystyle \frac{9_{diq}+\omega }{_{diq}+\omega }}.`$ (54)
For the product of condensates $`\overline{q}q\frac{\alpha _s}{\pi }G^2`$, wherein the gluon fields are connected to the heavy quarks, it is convenient to compute the contribution to the two-point correlation function itself. We have found
$$F_2^{\overline{q}qG^2}(\omega )=\frac{m_{QQ^{}}^{5/2}(m_Q^2+m_Q^{}^2+11m_Qm_Q^{})}{2^9\sqrt{2}\pi m_Qm_Q^{}(\omega )^{5/2}},$$
(55)
and we put $`F_1^{\overline{q}qG^2}(\omega )=0`$, since we have restricted the dimension of condensate operators, while the corresponding term in $`F_1`$ would appear in the fifth order in expansion (9). The result is given in the form, which allows the analytical continuation over $`\omega =+w`$.
### C Matching with full QCD
To connect the NRQCD sum rules to the baryonic couplings in full QCD we have to use the relation
$$J^{QCD}=𝒦_J(\alpha _s,\mu _{\mathrm{soft}},\mu _{\mathrm{hard}})J^{NRQCD},$$
where the coefficient $`𝒦_J(\alpha _s,\mu _{\mathrm{soft}},\mu _{\mathrm{hard}})`$ depends on the soft normalization scale $`\mu _{\mathrm{soft}}`$. The $`𝒦`$-factor obeys the matching condition at the hard scale $`\mu _{\mathrm{hard}}=_{diq}`$ and is determined by the anomalous dimensions of effective baryonic currents which are independent of the diquark spin in the leading order. They are known up to the two loop accuracy . In our consideration we use the one-loop accuracy, since the subleading corrections in the first $`\alpha _s`$ order are not available yet. So,
$`\gamma `$ $`=`$ $`{\displaystyle \frac{d\mathrm{ln}𝒦_J(\alpha _s,\mu _{\mathrm{soft}},\mu _{\mathrm{hard}})}{d\mathrm{ln}(\mu )}}={\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{\alpha _s}{4\pi }}\right)^m\gamma ^{(m)},`$ (56)
$`\gamma ^{(1)}`$ $`=`$ $`\left(2C_B(3a3)+3C_F(a2)\right),`$ (57)
where $`C_F=(N_c^21)/2N_c`$, $`C_B=(N_c+1)/2N_c`$ for $`N_c=3`$ and $`a`$ is the gauge parameter. In Feynman gauge $`a=1`$, and we get $`\gamma ^{(1)}=4`$. So, in the leading logarithmic approximation and to the one-loop accuracy we find
$$𝒦_J(\alpha _s,\mu _{\mathrm{soft}},\mu _{\mathrm{hard}})=\left(\frac{\alpha _s(\mu _{\mathrm{hard}})}{\alpha _s(\mu _{\mathrm{soft}})}\right)^{\frac{\gamma ^{(1)}}{2\beta _0}},$$
(58)
where $`\beta _0=112/3N_F=9`$. Further, we determine the soft normalization scale for the NRQCD estimates by the average momentum transfer inside the doubly heavy diquark, so that $`\mu _{\mathrm{soft}}^2=2m_{QQ^{}}T_{diq}`$, where $`T_{diq}`$ is the kinetic energy in the system of two heavy quarks, which is phenomenologically independent of the heavy quark flavours and approximately equal to 0.2 GeV . Then, the coefficients $`𝒦_J`$ are equal to
$$𝒦_{\mathrm{\Omega }_{cc}}1.95,𝒦_{\mathrm{\Omega }_{bc}}1.52,𝒦_{\mathrm{\Omega }_{bb}}1.30,$$
(59)
with the characteristic uncertainty about $`10\%`$ basically due to the variation of hard and soft scale points $`\mu _{\mathrm{hard},\mathrm{soft}}`$. Note that the values of $`𝒦_J`$ do not change the estimates of baryon masses, but they are essential in the evaluation of baryon couplings.
## III Numerical results
Evaluating the two-point sum rules, we explore the scheme of moments. We point out the well-known fact that an essential part of uncertainties is caused by the variation of heavy quark masses. Indeed, the results of sum rules for the systems containing two heavy quarks strongly depend on the choice of masses, and this fact allows us to pin down the values of masses with a high precision up to 20 MeV , so that $`m_b=4.60\pm 0.02`$GeV, $`m_c=1.40\pm 0.03`$GeV, which are extracted from the two-point sum rules for the families of $`\mathrm{{\rm Y}}`$ and $`\psi `$. To get these values we have use the correlators evaluated up to the same accuracy in $`\alpha _s`$, i.e. we have put the quark loop with the appropriate Sommerfeld factor. Then, the stability criterion for the leptonic constant of heavy quarkonium strictly fixes the heavy quark masses, which are close to the results of . The same sum rules are also explored to estimate the couplings determining the coulomb-like interactions inside the heavy quarkonia
$$\alpha _s(b\overline{b})=0.37,\alpha _s(c\overline{b})=0.45,\alpha _s(c\overline{c})=0.60,$$
(60)
since they fix the absolute normalization of corresponding leptonic constants known experimentally.
Since the squared size of diquark is two times larger than that of the meson the effective coulumb constants have to be rescaled according to the equation of evolution in QCD. We use the one-loop evolution equation
$$\alpha _s(QQ^{})=\frac{\alpha _s(Q\overline{Q}^{})}{1\frac{9}{4\pi }\alpha _s(Q\overline{Q}^{})\mathrm{ln}2}.$$
So,
$$\alpha _s(bb)=0.45,\alpha _s(bc)=0.58,\alpha _s(cc)=0.85.$$
(61)
The values of condensates are taken in the ranges $`\overline{q}q=(0.26÷0.27\text{GeV})^3`$, $`m_0^2=0.75÷0.85\text{GeV}^2`$, $`\frac{\alpha _s}{\pi }G^2=(1.7÷1.8)10^2\text{GeV}^4`$.
The main source of uncertainties in the ratios of the baryonic couplings is the ratio of the condensates of the strange quark and light quark. We use $`\overline{s}s/\overline{q}q=0.8\pm 0.2`$ that corresponds to the variations of the sum $`(m_u+m_d)[1\mathrm{GeV}]=12÷14`$ MeV .
So, we have described the set of parameters entering the scheme of calculations. In Figs. 1-3 we present the results of the two-point sum rules for the masses of $`\mathrm{\Xi }_{bc}`$ and $`\mathrm{\Omega }_{bc}`$ (the figures for the other baryons are similar). For the $`\mathrm{\Omega }_{bc}`$-baryons one can observe the stability of mass with respect to the changing of the moment numbers in both correlators. We suppose it is connected with the destroying of diquark-$`\mathrm{\Omega }`$baryon factorization in the perturbative limit in contrast to the $`\mathrm{\Xi }`$-baryons. The stability regions for $`F_1`$ and $`F_2`$ are not coincide because the contributions of higher dimension operators become valuable at the different numbers of moments. However, the quantity $`1/2(M_1[n]+M_2[n])`$ has the larger stability region, and we explore this fact to determine the $`\mathrm{\Omega }`$ baryons masses as well as that of $`\mathrm{\Xi }`$ baryons. The theoretical uncertainties in the $`\mathrm{\Omega }`$-baryon masses are mainly determined by the difference between the values of baryon masses at the regions of stability.
Then, we investigate the difference between the masses $`1/2((M_{1,\mathrm{\Omega }}+M_{2,\mathrm{\Omega }})(M_{1,\mathrm{\Xi }}+M_{2,\mathrm{\Xi }}))`$ shown in Fig. 4. In our scheme of baryon masses determination this quantity has the meaning of the difference between the $`\mathrm{\Omega }`$ and $`\mathrm{\Xi }`$ baryon masses. It has the large region of stability and is determined with a good precision. We obtain
$$\mathrm{\Delta }M=M_{\mathrm{\Omega }_{bb}}M_{\mathrm{\Xi }_{bb}}=M_{\mathrm{\Omega }_{cc}}M_{\mathrm{\Xi }_{cc}}=M_{\mathrm{\Omega }_{bc}}M_{\mathrm{\Xi }_{bc}}=100\pm 10\mathrm{MeV}.$$
The uncertainty in the $`\mathrm{\Xi }`$-baryons masses are determined through the uncertainty in the $`\mathrm{\Omega }`$-baryons masses and that of in $`\mathrm{\Delta }M`$. So, for the masses we find the following results:
$$\begin{array}{cccccccc}M_{\mathrm{\Omega }_{bc}}\hfill & =& \hfill 6.89\pm 0.05& \hfill \text{GeV},& M_{\mathrm{\Xi }_{bc}}\hfill & =& \hfill 6.79\pm 0.06& \hfill \text{GeV},\\ M_{\mathrm{\Omega }_{bb}}\hfill & =& \hfill 10.09\pm 0.05& \hfill \text{GeV},& M_{\mathrm{\Xi }_{bb}}\hfill & =& \hfill 10.00\pm 0.06& \hfill \text{GeV},\\ M_{\mathrm{\Omega }_{cc}}\hfill & =& \hfill 3.65\pm 0.05& \hfill \text{GeV},& M_{\mathrm{\Xi }_{cc}}\hfill & =& \hfill 3.55\pm 0.06& \hfill \text{GeV}.\end{array}$$
(62)
The obtained values are in agreement with the calculations in the framework of nonrelativistic potential models , though the models based on the calculation of three body systems with the pair interactions give slightly higher values of masses. In the other method of baryon mass determination was used, since the quantities $`M_{1,\mathrm{\Xi }}`$ and $`M_{2,\mathrm{\Xi }}`$ separately have no good stability in the sum rules. So, the difference of $`M_1M_2`$ close to zero was stable. The use of $`\frac{1}{2}(M_1+M_2)`$ stability criterion results in the $`\mathrm{\Xi }_{QQ^{}}`$ masses coinciding with those of up to 10 MeV. Figs.5, 6 show the dependence of baryon couplings calculated in the moment scheme of NRQCD sum rules. Numerically, we find
$$\begin{array}{cccccccc}|Z_{\mathrm{\Omega }_{cc}}|^2\hfill & =& \hfill (10.0\pm 1.2)10^3& \hfill \text{GeV}^6,& |Z_{\mathrm{\Xi }_{cc}}|^2\hfill & =& \hfill (7.2\pm 0.8)10^3& \hfill \text{GeV}^6,\\ |Z_{\mathrm{\Omega }_{bc}}|^2\hfill & =& \hfill (15.6\pm 1.6)10^3& \hfill \text{GeV}^6,& |Z_{\mathrm{\Xi }_{bc}}|^2\hfill & =& \hfill (11.6\pm 1.0)10^3& \hfill \text{GeV}^6,\\ |Z_{\mathrm{\Omega }_{bb}}|^2\hfill & =& \hfill (6.0\pm 0.8)10^2& \hfill \text{GeV}^6,& |Z_{\mathrm{\Xi }_{bb}}|^2\hfill & =& \hfill (4.2\pm 0.6)10^2& \hfill \text{GeV}^6.\end{array}$$
(63)
In Fig. 7 we present the sum rules results for the ratio of baryonic constants $`|Z_{\mathrm{\Omega }_{bc}}|^2/|Z_{\mathrm{\Xi }_{bc}}|^2`$. We have also found
$$|Z_{\mathrm{\Omega }_{bc}}|^2/|Z_{\mathrm{\Xi }_{bc}}|^2=|Z_{\mathrm{\Omega }_{cc}}|^2/|Z_{\mathrm{\Xi }_{cc}}|^2=|Z_{\mathrm{\Omega }_{bb}}|^2/|Z_{\mathrm{\Xi }_{bb}}|^2=1.3\pm 0.2.$$
The uncertainty of this result as was mentioned above is mainly connected with the pourly known ratio of $`\overline{s}s/\overline{q}q=0.8\pm 0.2`$.
For the sake of comparison, we derive the relation between the baryon coupling and the wave function of doubly heavy baryon evaluated in the framework of potential model, where we have used the approximation of quark-diquark factorization. So, we find
$$|Z^{\mathrm{PM}}|=2\sqrt{3}|\mathrm{\Psi }_d(0)\mathrm{\Psi }_{l,s}(0)|,$$
(64)
where $`\mathrm{\Psi }_d(0)`$ and $`\mathrm{\Psi }_{l,s}(0)`$ denote the wave functions at the origin for the doubly heavy diquark and light (strange) quark-diquark systems, respectively. In the approximation used, the values of $`\mathrm{\Psi }(0)`$ were calculated in in the potential by Buchmüller–Tye , so that
$`\sqrt{4\pi }|\mathrm{\Psi }_l(0)|`$ $`=`$ $`0.53\mathrm{GeV}^{3/2},`$ (65)
$`\sqrt{4\pi }|\mathrm{\Psi }_s(0)|`$ $`=`$ $`0.64\mathrm{GeV}^{3/2},`$ (66)
$`\sqrt{4\pi }|\mathrm{\Psi }_{cc}(0)|`$ $`=`$ $`0.53\mathrm{GeV}^{3/2},`$ (67)
$`\sqrt{4\pi }|\mathrm{\Psi }_{bc}(0)|`$ $`=`$ $`0.73\mathrm{GeV}^{3/2},`$ (68)
$`\sqrt{4\pi }|\mathrm{\Psi }_{bb}(0)|`$ $`=`$ $`1.35\mathrm{GeV}^{3/2}.`$ (69)
In the static limit of potential models, these parameters result in the estimates
$`|Z_{\mathrm{\Omega }_{cc}}^{\mathrm{PM}}|^2=8.810^3\mathrm{GeV}^6,|Z_{\mathrm{\Xi }_{cc}}^{\mathrm{PM}}|^2`$ $`=`$ $`6.010^3\mathrm{GeV}^6,`$ (70)
$`|Z_{\mathrm{\Omega }_{bc}}^{\mathrm{PM}}|^2=1.610^2\mathrm{GeV}^6,|Z_{\mathrm{\Xi }_{bc}}^{\mathrm{PM}}|^2`$ $`=`$ $`1.110^2\mathrm{GeV}^6,`$ (71)
$`|Z_{\mathrm{\Omega }_{bb}}^{\mathrm{PM}}|^2=5.610^2\mathrm{GeV}^6,|Z_{\mathrm{\Xi }_{bb}}^{\mathrm{PM}}|^2`$ $`=`$ $`3.910^2\mathrm{GeV}^6.`$ (72)
The estimates in the potential model (71) are close to the values obtained in the sum rules of NRQCD (63). We also see that the SU(3)-flavor splitting for the baryonic constants $`|Z_\mathrm{\Omega }|^2/|Z_\mathrm{\Xi }|^2`$ is determined by the ratio $`|\mathrm{\Psi }_s(0)|^2/|\mathrm{\Psi }_l(0)|^2=1.45`$ which is in agreement with the sum rules result. The values obtained in the NRQCD sum rules have to be multiplied by the Wilson coefficients coming from the expansion of full QCD operators in terms of NRQCD fields, as they have been estimated by use of corresponding anomalous dimensions. This procedure results in the final estimates
$`|Z_{\mathrm{\Omega }_{cc}}|^2`$ $`=`$ $`(38\pm 5)10^3\text{GeV}^6,|Z_{\mathrm{\Xi }_{cc}}|^2=(27\pm 3)10^3\text{GeV}^6,`$ (73)
$`|Z_{\mathrm{\Omega }_{bc}}|^2`$ $`=`$ $`(36\pm 4)10^3\text{GeV}^6,|Z_{\mathrm{\Xi }_{bc}}|^2=(27\pm 3)10^3\text{GeV}^6,`$ (74)
$`|Z_{\mathrm{\Omega }_{bb}}|^2`$ $`=`$ $`(10\pm 1)10^2\text{GeV}^6,|Z_{\mathrm{\Xi }_{bb}}|^2=(70\pm 8)10^3\text{GeV}^6.`$ (75)
## IV Conclusion
In this paper the NRQCD sum rules applied to the doubly heavy baryons have been considered. The nonrelativistic approximation for the heavy quark fields allows us to fix the structure of baryonic currents (the light quark-doubly heavy diquark) and to take into account the coulomb-like interactions inside the doubly heavy diquark. The presence of both the nonzero mass of light quark and the contribution of nonperturbative terms of the quark, gluon, mixed condensates and the product of condensates destroys the factorization of the correlators. This fact provides the convergency of sum rules for each correlator and allows us to obtain the reliable results for the masses and baryonic constants, which agree with the estimates in the framework of potential models. We also have calculated the mass splitting of $`\mathrm{\Omega }`$ and $`\mathrm{\Xi }`$ doubly heavy baryons and the ratio of baryonic constants $`|Z_\mathrm{\Omega }|^2/|Z_\mathrm{\Xi }|^2`$.
The authors are grateful to prof. A.K.Likhoded for stimulating discussions.
This work is in part supported by the Russian Foundation for Basic Research, grants 99-02-16558 and 00-15-96645.
## V Appendix
Here the derivation of expansion (9) is briefly presented. The calculations are done in the technique of fixed point gauge , so we write down the expansion of quark field:
$$q(x)=q(0)+x^\alpha D_\alpha q(0)+\frac{1}{2}x^\alpha x^\beta D_\alpha D_\beta q(0)+\mathrm{},$$
and in the evaluation of $`0|Tq_i^a(x)\overline{q}_j^b(0)|0`$, where $`i`$ and$`j`$ are the spinor indices, $`a,b`$ are the color indices, we have to know how to get the vacuum average of type $`0|D_\alpha \mathrm{}D_\omega q(0)\overline{q}(0)|0`$. The main formulae are the followings:
the definitions of condensates
$$q_i^a(0)\overline{q}_j^b(0)_0=\frac{1}{12}\delta ^{ab}\delta _{ij}\overline{q}q,$$
$$G_{\alpha \beta }^aG_{\alpha ^{}\beta ^{}}^a^{}=\frac{\delta ^{aa^{}}}{96}(g_{\alpha \alpha ^{}}g_{\beta \beta ^{}}g_{\alpha \beta ^{}}g_{\alpha ^{}\beta })G^2,$$
$$\overline{q}igG_{\alpha \beta }^at^a\sigma _{\alpha \beta }q_0=m_0^2\overline{q}q,$$
the commutator of covariant derivatives
$$[D_\alpha ,D_\beta ]=igG_{\alpha \beta }^at^a,$$
and the equation of motion for the spinor field
$$\text{ / }Dq=im_qq.$$
Form the last two equations we derive the so-called quadratic Dirac equation,
$$D^2q=m_q^2q+\frac{\sigma _{\alpha \beta }}{2}igG_{\alpha \beta }^at^aq.$$
Now it is an easy challenge to obtain the first term in expansion (9).
Since the tensor $`x_\alpha \mathrm{}x_\omega `$ is the symmetric one, we may perform the symmetrization
$$D_\alpha \mathrm{}D_\omega \{D_\alpha ,\mathrm{},D_\omega \}_+,$$
to find the n-th term of expansion for $`\overline{q}(x)q(0)`$, which equals
$$\frac{1}{n!}x_\alpha \mathrm{}x_\omega \overline{q}(0)D_\alpha \mathrm{}D_\omega q(0)=\frac{1}{n!}x_\alpha \mathrm{}x_\omega \overline{q}(0)\{D_\alpha ,\mathrm{},D_\omega \}_+q(0).$$
Note, the tensor $`\overline{q}(0)\{D_\alpha \mathrm{}D_\omega \}_+q(0)`$ is also symmetric one.
The second term of expansion is derived from
$$\{D_\alpha ,D_\beta \}_+q_\rho ^i(0)\overline{q}_\eta ^j(0)=2!𝒫_2g_{\alpha \beta }\delta ^{ij}\delta _{\rho \eta }\overline{q}q,$$
and the coefficient $`𝒫_2`$ is determined by contracting the indices $`\alpha ,\beta `$ and using the quadratic Dirac equation,
$$𝒫_2=(m_0^22m_q^2)/192.$$
The third term can be derived from the following structure:
$$\{D_\alpha ,D_\beta ,D_\delta \}_+q_\rho ^i(0)\overline{q}_\eta (0)=3!𝒫_3\delta ^{ij}((\gamma _\alpha )_{\rho \eta }g_{\beta \delta }+(\gamma _\beta )_{\rho \eta }g_{\alpha \delta }+(\gamma _\delta )_{\rho \eta }g_{\alpha \beta })\overline{q}q.$$
Then, contracting $`\alpha `$ and $`\beta `$ and using of the equation of motion, the quadratic Dirac equation and the commutation relation, we obtain
$$𝒫_3=im_q(3m_0^2/4m_q^2)/576.$$
This includes the evaluation of vacuum averages
$$D^2D_\alpha q(0)\overline{q}(0),D_\beta D_\alpha D_\beta q(0)\overline{q}(0)\text{and}D_\alpha D^2q(0)\overline{q}(0).$$
Considering the structure
$$\{D_\alpha ,D_\beta ,D_\delta ,D_\xi \}_+q_\rho ^i(0)\overline{q}_\eta (0)=4!𝒫_4\delta ^{ij}\delta _{\rho \eta }(g_{\alpha \beta }g_{\delta \xi }+g_{\alpha \delta }g_{\beta \xi }+g_{\alpha \xi }g_{\delta \beta })\overline{q}q$$
contracted over any pair of indices, we derive
$$𝒫_4=(\pi ^2\alpha _s/\pi G^2+3/2m_q^2(m_q^2m_0^2))/3456.$$
Here we evaluated the following types of vacuum expectations:
$$D^2D^2q(0)\overline{q}(0),D_\alpha D_\beta D_\alpha D_\beta q(0)\overline{q}(0),D_\alpha D^2D_\alpha q(0)\overline{q}(0).$$
Then the OPE for the quark condensate can be expressed in terms of $`𝒫_i`$ by
$$q_\rho ^i(x)\overline{q}_\eta ^j(0)=\delta ^{ij}\overline{q}q(𝒫_0\delta _{\rho \eta }+𝒫_1x_\alpha \gamma _{\rho \eta }^\alpha +𝒫_2\delta _{\rho \eta }x^2+𝒫_3x_\alpha \gamma _{\rho \eta }^\alpha x^2+𝒫_4\delta _{\rho \eta }x^4),$$
with $`𝒫_0=1/12`$, and $`𝒫_1=im_q/48`$.
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# Event reconstruction in high resolution Compton telescopes
## 1 Compton telescopes for $`\gamma `$-ray astrophysics
Looking beyond the INTErnational Gamma-Ray Astrophysics Laboratory (INTEGRAL), the next generation soft $`\gamma `$-ray ($``$0.2-20 MeV) observatory will require high angular and spectral resolution imaging to significantly improve sensitivity to astrophysical sources of nuclear line emission. Building upon the success of COMPTEL/CGRO (Schönfelder et al. (1993)), and the high spectral resolution of the upcoming SPI/INTEGRAL (Vedrenne et al. (1998); Lichti et al. (1996)), a number of researchers (Johnson et al. (1996); Jean et al. (1996); Boggs (1998)) have discussed the merits of a high spectral/angular resolution germanium Compton telescope (GCT); the ability to achieve high sensitivity to point sources while maintaining a large field-of-view make a high resolution Compton telescope an attractive option for the next soft $`\gamma `$-ray observatory.
The development of Compton telescopes began in the 1970’s, with work done at the Max Planck Institut (Schönfelder et al. (1973)), University of California, Riverside (Herzo et al. (1975)), and the University of New Hampshire (Lockwood et al. (1979)), culminating in the design and flight of COMPTEL/CGRO. These historical Compton telescopes consist of two scintillation detector planes – a low atomic number ‘converter’ and a high atomic number ‘absorber.’ The model interaction of a Compton telescope is a single Compton scatter in the converter plane, followed by photoelectric absorption of the scattered photon in the absorber. By measuring the position and energy of the interactions, the event can be reconstructed to determine the initial photon direction to within an annulus on the sky.
A handful of groups are actively developing imaging germanium detectors (GeDs) partly in anticipation of a GCT (Luke et al. (1994); Kroeger et al. (1996)). The goal of these researchers is to develop large area detectors with (sub)millimeter spatial resolution, while maintaining the high spectral resolution ($`E/\delta E500`$ at 1 MeV) characteristic of GeDs. The use of high spectral/spatial resolution GeDs as converter and absorber planes would significantly improve the performance of a Compton telescope, but will add a number of complications to the event reconstruction. Most significantly, with the moderate atomic number $`(Z=32)`$ of germanium, photons will predominantly undergo multiple Compton scatters before being photoabsorbed in the instrument. Furthermore, with interaction timing capabilities of $``$10 ns, the interaction order will not be determined unambiguously by timing alone. Compton Kinematic Discrimination (CKD) is proposed here to overcome these complications, an extension of a method first discussed in context of liquid xenon time projection chambers (Aprile et al. (1993)). The ability of this technique to allow proper event reconstruction is investigated in detail.
Due to their relatively low efficiency (typically $`1\%`$), Compton telescopes rely on efficient background suppression to maintain their sensitivity. In addition to interaction ordering, techniques are presented using CKD, in combination with other tests and restrictions, to suppress the dominant background components.
The goal of this work is to outline a complete set of event reconstruction techniques for GCTs, taking into account realistic detector/instrument performance and uncertainties. Examples of the techniques are presented for a GCT configuration outlined in Appendix A; however, full analysis of this configuration will be presented in a second paper dedicated to the optimization and performance of several GCT configurations. The full analysis of a GCT configuration is complicated, requiring a detailed study of the tradeoffs between efficiency, angular and spectral resolution; therefore, this paper focuses only on the detailed discussion of the event reconstruction techniques which will be used in future work dedicated to analyzing GCT performance.
## 2 Principles of Compton imaging
The principle of Compton imaging of $`\gamma `$-ray photons is illustrated in Figure 1. (See von Ballmoos, Diehl and Schönfelder (1989) for an excellent review of historical Compton telescope configurations.) An incoming photon of energy $`E`$ and direction $`\widehat{p}`$ undergoes a Compton scatter at an angle $`\varphi _1`$ at the position $`r_1`$ within a detector, creating a recoil electron of energy of $`E_1`$ which is quickly absorbed and measured by the detector itself. The scattered photon then deposits the rest of its energy in the instrument in a series of one or more interactions of energies $`E_i`$ at the positions $`r_i`$, until eventually photoabsorbed. Here the total photon energy after each scatter $`i`$, normalized to the electron mass, is defined as
$$W_i=\frac{1}{m_ec^2}\underset{j=i+1}{\overset{N}{}}E_j,$$
(1)
where $`W_0=E/m_ec^2`$, and $`N`$ is the total number of interactions. The initial photon direction is related to scatter direction vector $`r_1^{}=r_2r_1`$ ($`\widehat{r}_1^{}`$ after normalization), and the scattered photon energies $`W_i`$ by the Compton formula
$$\widehat{r}_1^{}\widehat{p}=\mathrm{cos}\varphi _1=1+\frac{1}{W_0}\frac{1}{W_1}.$$
(2)
Given the measured scatter direction $`\widehat{r}_1^{}`$ and the angle $`\mathrm{cos}\varphi _1`$ implied from the energy depositions, the equation for $`\widehat{p}`$ is not unique (if the electron recoil direction could be measured, it would solve this ambiguity); therefore, the initial direction of the photon cannot be determined directly, but it can be limited to an annulus of directions $`\widehat{p}^{}`$ which satisfy the equation
$$\widehat{r}_1^{}\widehat{p}^{}=\mathrm{cos}\varphi _1.$$
(3)
There are two uncertainties in determining the event annulus: the uncertainty in $`\varphi _1`$ due to the finite energy resolution of the detectors, here labelled $`\delta \varphi _{1,E}`$, and the uncertainty in $`r_1^{}`$ determined by the spatial resolution of the detectors. Both of these uncertainties add to determine the uncertainty (effective width) of the event annulus $`\delta \varphi _1`$. From Equation 2, the derivation of $`\delta \varphi _{1,E}`$ is straightforward and yields
$$\delta \varphi _{i,E}=\frac{1}{\mathrm{sin}\varphi _i}[(\frac{\delta W_{i1}^2}{W_{i1}^4})+\delta W_i^2((\frac{1}{W_i^2}\frac{1}{W_{i1}^2})^2\frac{1}{W_{i1}^4})]^{1/2},$$
(4)
where,
$$\delta W_i=\frac{1}{m_ec^2}[\underset{j=i+1}{\overset{N}{}}\delta E_j^2]^{1/2}.$$
(5)
In order to simplify the analysis, it is convenient to transform the uncertainty $`\delta r_1^{}`$ into an effective uncertainty in $`\varphi _1`$, defined as $`\delta \varphi _{1,r}`$, such that
$$\delta \varphi _1=\sqrt{\delta \varphi _{1,E}^2+\delta \varphi _{1,r}^2}.$$
(6)
The angular resolution $`\delta \varphi _{i,r}`$ is the effective ‘wiggle’ of $`\widehat{r}_i^{}`$ around its measured direction due to the uncertainties in the spatial measurements. The spatial uncertainties are defined as $`\delta x_i^{}=\sqrt{\delta x_i^2+\delta x_{i+1}^2}`$, $`\delta y_i^{}=\sqrt{\delta y_i^2+\delta y_{i+1}^2}`$, $`\delta z_i^{}=\sqrt{\delta z_i^2+\delta z_{i+1}^2}`$. It is simplest to analyze the situation for each axis separately as shown in Figure 2. The uncertainty in the direction of $`\widehat{r}_i^{}`$ due to the uncertainty $`\delta x_i^{}`$ is given by
$$\delta \varphi _{i,x}\mathrm{tan}(\delta \varphi _{i,x})=(\frac{\delta x_i^{}}{r_i^{}})\sqrt{1(\widehat{r}_i^{}\widehat{x})^2}.$$
(7)
Likewise for the other axis,
$$\delta \varphi _{i,y}(\frac{\delta y_i^{}}{r_i^{}})\sqrt{1(\widehat{r}_i^{}\widehat{y})^2},$$
$$\delta \varphi _{i,z}(\frac{\delta z_i^{}}{r_i^{}})\sqrt{1(\widehat{r}_i^{}\widehat{z})^2},$$
which combine to yield the total uncertainty $`\delta \varphi _{i,r}`$ given by
$$\delta \varphi _{i,r}=\sqrt{\delta \varphi _{i,x}^2+\delta \varphi _{i,y}^2+\delta \varphi _{i,z}^2}.$$
(8)
For detectors with a given energy resolution, in order to optimize the performance of a Compton telescope one would require that $`\delta \varphi _{i,r}\delta \varphi _{i,E}`$ in the energy range of interest. To first order, this implies that the spatial resolution in relation to the scale size of the instrument must be comparable to or less than the energy resolution, i.e. $`\delta r_1^{}/r_1^{}\delta E_1^{}/E_1^{}`$.
## 3 Complications of germanium Compton telescopes
Finite detector thresholds, energy resolutions, and spatial resolutions produce systematic biases in the imaging capabilities of Compton telescopes. These limitations have been discussed in detail elsewhere in context of two-layer, low-Z converter and high-Z absorber, scintillation detector designs (von Ballmoos et al. (1989)), and the conclusions can be directly applied to GCTs. However, GCT designs will introduce additional complications which significantly alter the event reconstruction techniques. Historical Compton telescope configurations make two assumptions about the events which do not generally hold in GCTs: (i) the events are a single Compton scatter in the converter, followed by photoelectric absorption in the absorber, and (ii) the time-of-flight (TOF) between the photon interactions is measured to determine their order.
The distributions of number and type of interaction sites in a GCT for normally incident, fully-absorbed photons ranging from 0.2-10 MeV are shown in Figure 3, for the instrument configuration discussed in Appendix A. Here we distinguish three event types: a single photoelectric absorption, one or more Compton scatters followed by a single photoabsorption, and one or more pair productions. Compton scatters followed by pair production could potentially be reconstructed; however, here we include these events with other pair productions. These distributions account for the finite spatial resolution of the detectors, so that interactions occurring too closely together are not resolved. From these distributions it is clear that events with $``$8 or more interaction sites can be immediately rejected as probable pair production events, with little effect on the Compton photopeak efficiency. For incident photon energies above 0.5 MeV, 3-7 interaction site Compton scatter events dominate the photopeak.
To accurately reconstruct a Compton scatter event, the first and second interaction sites must be spatially resolved, and their order correctly determined. The need to determine the proper ordering of three or more (3+) interaction sites is complicated by the timing capabilities of GeDs. In the scintillation detectors of COMPTEL/CGRO the interaction timing can be performed to $``$0.25 ns (Schönfelder et al. (1993)), which is adequate to determine the TOF between two interactions in the separate detector planes. With the slower rise time of GeDs one can reasonably expect event timing to $``$10 ns, which is inadequate for TOF measurement in reasonably-sized instruments. While Pulse Shape Discrimination methods have been proposed to push the interaction timing in GeDs to $``$1 ns (Boggs (1998)), even this timing would be unreliable for determining TOF among three or more interaction sites. A method of reliably determining the photon interaction order without timing information must be developed.
## 4 Multiple Compton scatter events: Compton kinematic discrimination
One method has been suggested to overcome these complications in the context of liquid xenon time projection chambers (Aprile et al. (1993)). Here, this method is formalized as Compton Kinematic Discrimination (CKD) and examined in more detail. This technique allows the order of the photon interactions to be determined with high probability, as well as providing the basis of a powerful tool for background suppression in GCTs.
CKD takes advantage of redundant measurement information in an event to determine the most likely interaction sequence. A photon of initial energy $`E`$ (using the notation in Section 2) interacts in the instrument at $`N`$ sites, depositing an energy of $`E_i`$ at each location $`r_i`$. It is assumed that the interactions $`1,\mathrm{},N1`$ are Compton scatters, and interaction $`N`$ is the final photoabsorption. Given the correct ordering of the interactions, there are two independent ways of measuring $`N2`$ of the scattering angles, $`\mathrm{cos}\varphi _2,\mathrm{},\mathrm{cos}\varphi _{N1}`$.
Geometrical measurement of $`\mathrm{cos}\varphi _i`$. From simple vector analysis, given the correct ordering of the interaction sites one can derive the scatter angles
$$\mathrm{cos}\varphi _i=\widehat{r}_i^{}\widehat{r}_{i1}^{},i=2,\mathrm{},N1,$$
(9)
where the uncertainties in the scattering angles, $`\delta \mathrm{cos}\varphi _i`$, can be estimated from the spatial uncertainty in the scattering angles (Equation 8), yielding
$$\delta (\mathrm{cos}\varphi _i)=\delta \varphi _{i,r}\mathrm{sin}\varphi _i.$$
(10)
Compton kinematics measurement of $`\mathrm{cos}\varphi _i^{}`$. Given the correct ordering, the measured values of $`W_i`$ can be derived, which were defined earlier as the energy of the photon after each scattering $`i`$, in units of $`m_ec^2`$. The Compton scatter formula (Equation 2) gives:
$$\mathrm{cos}\varphi _i^{}=1+\frac{1}{W_{i1}}\frac{1}{W_i},i=1,\mathrm{},N1,$$
(11)
$$\delta (\mathrm{cos}\varphi _i^{})=[(\frac{\delta W_{i1}^2}{W_{i1}^4})+\delta W_i^2((\frac{1}{W_i^2}\frac{1}{W_{i1}^2})^2\frac{1}{W_{i1}^4})]^{1/2}.$$
(12)
Given the $`N2`$ independent measurements of $`\mathrm{cos}\varphi _2,\mathrm{},\mathrm{cos}\varphi _{N1}`$, a trial ordering of the interaction sites can be tested for consistency. If the assumed ordering is incorrect $`\mathrm{cos}\varphi _i`$ will not equal $`\mathrm{cos}\varphi _i^{}`$ in general. Every possible permutation of orderings can be tested to determine the one most consistent with $`\mathrm{cos}\varphi _i=\mathrm{cos}\varphi _i^{}`$. Given a trial ordering, the two angle cosines for sites $`i=2,\mathrm{},N2`$ are relabelled for convenience $`\eta _i=\mathrm{cos}\varphi _i`$, $`\eta _i^{}=\mathrm{cos}\varphi _i^{}`$.
As a first test, trial orderings that produce values of $`\eta _i^{}1`$ are ruled out, since $`\mathrm{cos}\varphi _i^{}<1`$ for any scattering angle $`\varphi _i^{}`$. This condition will eliminate many orderings which cannot physically be due to multiple Compton scatters followed by photoabsorption. Next a least-squares statistic measuring the agreement between the redundant scatter angle measurements is defined:
$$\chi ^2=\frac{1}{N2}\underset{i=2}{\overset{N1}{}}\frac{(\eta _i\eta _i^{})^2}{(\delta \eta _i^2+\delta \eta _i^2)}.$$
(13)
In general, $`\chi ^2`$ will be minimized when the interactions are properly ordered (i.e. the order in which they occurred). Therefore, all possible permutations can be tested for their value of $`\chi ^2`$, and the ordering corresponding to the minimum value, $`\chi _{min}^2`$, is taken as the most likely ordering.
This consistency statistic also provides a powerful tool for rejecting background events. If the event is truly a multiple Compton scatter event followed by a photoabsorption then $`\chi _{min}^21`$. By setting a maximum acceptable level for $`\chi _{min}^2`$, events that do not fit this scenario can be rejected. Such events include partially-deposited photons which scatter out of the instrument (Compton continuum), photon interactions with spatially unresolved interaction sites, events with interactions below the detector threshold, pair-production events, and similarly $`\beta ^+`$ decays. These events frequently have $`\chi _{min}^21`$, allowing a strong rejection statistic that is not very sensitive on the level set on $`\chi _{min}^2`$. Here, $`\chi _{min}^2`$ has been treated as a normal least-squares statistic with $`N2`$ degrees of freedom, and events are rejected which have probabilities of $`\chi _{min}^2<5\%`$. Variations in the level between $`1\%`$ and $`10\%`$ do not strongly affect CKD rejection capabilities. For example, varying this level from $`5\%`$ to $`1\%`$ shifted the CKD efficiency curves in Figure 4 by 1-2%.
The fraction of 3+ site photopeak events which have the first and second interaction sites spatially resolved – and hence could be imaged to the proper direction – is shown in Figure 4 as a function of energy, for the instrument model discussed in Appendix A. Roughly $`90\%`$ of all events from 0.2-20 MeV have their first and second sites spatially resolved from each other. (Some of these events do not have their second, third, etc., interactions spatially resolved from each other, and will be rejected by the limits on $`\chi _{min}^2`$.) This figure also shows the fraction of the photopeak events which CKD properly orders (correctly reconstructed), as well as the fractions improperly ordered (hence incorrectly imaged to off-source background), and the fraction completely rejected. For energies below $``$10 MeV, CKD allows proper reconstruction (hence imaging) of $`6070\%`$ of the photopeak events, while rejecting $`1020\%`$. The remaining $`1020\%`$ are incorrectly imaged into the off-source background. For comparison, if the order of the interaction sites were randomly chosen $`<15\%`$ would be correctly imaged, while the remaining $`>85\%`$ would be incorrectly imaged into the background.
## 5 Single Compton scatter events: single scatter discrimination
CKD will only work for $`N>2`$ since there are no independent scattering angle measurements for a single Compton scatter followed by a photoabsorption. It turns out that the ordering of two-site photopeak events can still be determined with a high probability; however, the ability to reject background events is lost. As is discussed further in Section 6, the loss of background rejection, coupled with low peak-to-Compton ratios, and a larger fraction of backscatter events mean that the inclusion of two-site events will likely hurt the sensitivity of a GCT; however, discussion of event ordering is still included for completeness.
Given a two-site event, the first test one can perform is to determine whether both possible orderings of the interaction sites are energetically compatible with a single Compton scatter, i.e. are compatible with the requirement that $`\mathrm{cos}\varphi _1<1`$ (Equation 2). In Figure 5, the fraction of spatially-resolved, photopeak events with unique orderings are plotted versus energy. Also plotted are the fraction with ambiguous orderings. At energies below $``$0.4 MeV the majority of resolved photopeak events have a unique ordering, while at higher energies most events are ambiguous.
As an empirical test of the ambiguous events, the relative magnitude of the energy lost in the initial scatter $`(E_1)`$ compared to the photoabsorption $`(E_2)`$ can be compared. The fraction of resolved two-site photopeak events which have ambiguous orderings with $`E_1>E_2`$ is plotted in Figure 5. At higher energies, nearly all of the resolved photopeak events with ambiguous interaction orders have $`E_1>E_2`$, which can be used to determine the most likely interaction order. This empirical result can be easily understood, in hindsight, by the fact that photons which deposit most of their energy in the initial Compton scatter are much more likely to be photoabsorbed in the second interaction.
Therefore, a simple Single Scatter Discrimination (SSD) technique to determine the most likely interaction ordering of two-site events follows. First one determines whether a physically unique ordering exists; if not, the larger energy deposit is assumed to be the initial Compton scatter. Only at the lowest energies are two-site events possible in which neither ordering is acceptable (unresolved events, Compton continuum), and some background rejection is possible.
The fraction of two-site photopeak events which are spatially resolved is shown in Figure 6 as a function of energy, for the instrument model in Appendix A. Roughly $`80\%`$ of all events from 0.2-20 MeV are resolved. This number is about $`10\%`$ lower than the 3+ site events, due to the smaller path lengths of the lower energy scattered photons in two-site photopeak events. Also shown in Figure 6 is the fraction of events which SSD has properly reconstructed. SSD allows proper reconstruction (hence imaging) of $`6080\%`$ of the photopeak events, while improperly imaging the remaining $`2040\%`$ into the off-source background. Only a relatively small number of low energy events can be rejected outright. SSD is least effective around 0.5 MeV, where the unique/ambiguous ordering signatures are not as clear. For comparison, if the order of the interaction sites were randomly chosen $`40\%`$ would be properly imaged, while the remaining $`60\%`$ would be improperly imaged into the background.
## 6 Full reconstruction with background rejection
Given these two methods of determining the photon interaction order in GCTs, CKD and SSD, a full event reconstruction technique incorporating other background rejection techniques can be developed. The high spectral and spatial resolution of a GCT make several powerful background rejection techniques possible. The predominance of multiple scattering events, while initially a complication, dramatically helps in the overall background rejection.
The dominant sources of background in GCTs are expected to be diffuse cosmic $`\gamma `$-ray emission, induced satellite $`\gamma `$-ray emission, and induced $`\beta ^+`$, $`\beta ^{}`$ radioactivities in the GeDs themselves (Jean et al. (1996); Graham et al. (1997); Gehrels (1985); Dean et al. (1991); Naya et al. (1996)). Source/background photons which scatter out of the instrument before depositing all of their energy, and hence are improperly imaged, must also be included in these calculations.
Restrictions on the acceptable events can have a dramatic effect on sensitivity of a GCT. Specifically, several factors can affect the angular resolution of the instrument as well as the background rates – such as the inclusion of backscatter events, limits on the accepted scatter angles, and the minimum acceptable lever arm – and must be included in any discussion of full reconstruction and background rejection.
### 6.1 Shield veto
Placing an active shield below the bottom GCT detector plane could be useful for rejecting background photons from below the instrument (induced satellite $`\gamma `$-ray emission), as well as helping to reject Compton continuum and $`\beta `$decay events in the instrument. However, many of these events can be distinguished and rejected using CKD and other tests/restrictions outlined below; therefore, the usefulness of including a shield in the GCT design must be studied in detail for a given telescope configuration.
### 6.2 Restrictions on number of interaction sites: pair production/$`\beta ^+`$ decays
While it is obvious that any single-site interactions should be rejected in a Compton telescope, events with $``$8 or more interaction sites should also be rejected since these are very likely due to pair production events, as is evident from Figure 3, or similarly $`\beta ^+`$ decays. (See Section 6.8)
### 6.3 CKD $`\chi _{min}^2`$ test: Compton continuum, unresolved interactions, etc.
Using the tests outlined in Sections 4 and 5, the most likely ordering of the interaction sites can be determined, and for 3+ site events many of the unresolved and Compton continuum events, as well as pair production and $`\beta `$decays, rejected.
Shown in Figure 7 is the peak-to-Compton ratio for 3+ site events, here defined as the ratio of the properly imaged photopeak events to the corresponding integrated Compton continuum (photons which scatter out of the instrument before depositing all of their energy). This standard measure for $`\gamma `$-ray spectroscopy instruments has an altered meaning here, since the Compton continuum events will be incorrectly imaged, and thus will appear as off-source background. The peak-to-Compton ratio is shown both before and after rejection of the continuum events with the CKD statistic. CKD rejection of the Compton continuum events increases the peak-to-Compton ratio by factors of 3-6, an important improvement for low background instruments. By rejection of events appearing to originate from below the instrument (Section 6.4), as well as backscattered interactions (Section 6.5), this ratio can be increased by further factors of 2-4.
Also shown in Figure 7 is the photopeak-to-Compton ratio for two-site events using SSD, which is significantly lower than the corresponding ratio for 3+ site events. In fact, this ratio drops significantly below unity, which means that more background than signal is being created in the instrument for two-site events. This result questions whether two-site events should be included in actual observational analysis given the accompanying increase in unrejectable background. This conclusion is further supported by the fact that the majority of two-site events are backscatters (Section 6.5), which will significantly degrade the angular resolution. Even though inclusion of two-site events is unlikely to improve the overall sensitivity for GCTs, detailed background analysis for specific instrument configurations is required to determine the overall effects.
### 6.4 Effective TOF
Once the most likely order of interactions and the initial scatter angle are determined, it is possible to determine whether the incident photon scattered upwards or downwards in the instrument, as well as whether the initial scatter was forward or backwards. Thus, events which appear to be photons originating from below the instrument can be rejected, which include the induced satellite $`\gamma `$-ray emission, many Compton continuum events which were not rejected by CKD, photons which scatter in the passive satellite material before interacting in the detectors, and many of the pair production and $`\beta `$decay events. The simulation results for the configuration in Appendix A show that $`95\%`$ of photons originating from below the instrument are rejected.
### 6.5 Backscatters
Once the most likely ordering of interaction sites has been determined, this information can also be used to accept/reject backscattered source photons. The fraction of photopeak events which backscatter during the initial interaction is not strongly energy dependent, $`6070\%`$ for two site events, and $`30\%`$ for 3+ site events. These events can significantly increase the effective area at lower energies, where two-site events are most common, at the expense of degrading the angular resolution due to larger uncertainties in $`\delta \varphi _{1,E}`$ for backscatters events (Equation 4). It is unlikely that the overall sensitivity will improve by including backscattered events given the increased background rates and degraded angular resolution; however, the effects on sensitivity will depend on the exact instrument configuration and observational goals.
### 6.6 ‘Standard $`\varphi `$ restriction’
Restrictions can also be set on the scattering angles accepted for forward-scattering photons entering the front of the instrument. These limits can be used to restrict the instrument FOV to improve imaging capabilities and background, such as the ‘standard $`\varphi `$ restriction’ (Schönfelder et al. (1982)). These restricitions will have to be reanalyzed in detail for specific GCT configurations.
### 6.7 Nonlocalized $`\beta ^{}`$ decays
In a nonlocalized $`\beta ^{}`$ decay, the daughter nuclide is produced in an excited state which quickly decays on timescales relative to the detector collection time, emitting a photon with energy characteristic to the daughter nuclide. Therefore, the event consists of the intial $`\beta ^{}`$ decay site, plus the interaction sites of the emitted photon. Such an event can be rejected if the characteristic photon energy can be detected in any combination of the interaction site energies. The seven dominant $`\beta ^{}`$ isotopes and characteristic photon energies for natural Ge are given in Table 1 (Naya et al. (1996); Gehrels (1985)). In general, if the coincident $`\gamma `$-ray is fully deposited, then the event will have $`\chi _{min}^21`$, with the $`\beta ^{}`$ electron interaction ordered as the initial ‘scatter’ site. In these cases, $`W_1`$ will have the characteristic photon energy specific to that decay. The rejection of all events with $`W_1`$ equal to one of the characteristic energies in Table 1 can dramatically decrease the $`\beta ^{}`$ decay background, with only a small effect, typically $`3\%`$ drop in photopeak efficiencies for true photon events. More $`\beta ^{}`$ decays can be rejected if every possible combination of interaction sites is tested for the decay photon energies, at the expense, however, of more rejection of photopeak events, typically $`1520\%`$. If the coincident $`\gamma `$-ray is only partially deposited, the event will likely be rejected by the limits set on $`\chi _{min}^2`$.
$`\beta ^{}`$ decay background events were simulated for the instrument discussed in Appendix A. The results of these background calculations will be presented in a separate paper, but here we make preliminary use of these simulations to demonstrate the background rejection capabilities. After initial rejection of two site ($`34.4\%`$) and 8+ site ($`2.9\%`$) events, $`62.7\%`$ of the nonlocalized $`\beta ^{}`$ events remain. Applying the CKD test, requiring a $`5\%`$ probability of $`\chi _{min}^2`$, brings the remaining number of decays down to 17.9%. After screening the interactions for characteristic $`\beta ^{}`$ decay energies, this number is reduced to $`9.3\%(15.1\%)`$. Finally, after rejecting events which appear to originate from below the instrument, or which appear to be backscatter events, the final number of unrejected $`\beta ^{}`$ decays comes to $`4.2\%(6.8\%)`$, a factor of 20 (15) reduction in this background component. (First numbers give results when all combinations of interaction sites are searched for $`\beta ^{}`$ decay energies, while the numbers in parenthesis are results when only $`W_1`$ is tested. Typical errors $`0.2\%`$.)
### 6.8 Positron signatures
A further test for rejecting the pair production/$`\beta ^+`$ background events that survive the other tests/restrictions outlined above is to search for positron annihilation signatures in the interaction energies. By analyzing all combinations of the interaction sites to see if the energies sum to $`m_ec^2=0.511MeV`$, events with a positron annihilation signature can be rejected. This test typically reduces the non-pair production photopeak events by $`2\%`$.
$`\beta ^+`$ background events were simulated for the instrument discussed in Appendix A. After initial rejection of two site events ($`23.0\%`$) and 8+ site events ($`3.0\%`$), $`74.0\%`$ of the events remain. After the CKD test, $`16.7\%`$ remain. After screening the interactions for 0.511 MeV positron annihilation signatures, the number is reduced to $`5.2\%`$. Finally, after rejecting events which appear to originate from below the instrument, or which appear to be backscatter events, the final number of unrejected $`\beta ^+`$ events is $`1.9\%`$, a factor of 50 reduction in this background component. Similar reductions occur when these tests are applied to pair production events. (Typical errors $`0.1\%`$.)
### 6.9 Minimum lever arm
In general, a minimum acceptable distance between the first and second interaction sites – the lever arm – must be set. Figure 8 shows the fraction of 0.5 and 2.0 MeV photopeak events with lever arms above a given level, for the instrument configuration in Appendix A. Similar to the case of backscattered events, a smaller minimum lever arm means a higher effective area at the expense of poorer angular resolution. The exact lever arm chosen will depend on the instrument configuration and observational goals.
## 7 Conclusions
Event reconstruction in future high resolution Compton telescopes will present a number of complications compared to historical configurations. The initial complication of multiple-scattering of photons in GCTs, however, turns out to be an advantage: the application of CKD to 3+ site events, combined with the high spectral and spatial resolution of GeDs, allows extremely efficient background suppression, crucial for Compton telescope performance. This paper has outlined a set of tests and restrictions, accounting for realistic instrument/detector performance, to reconstruct photopeak events in GCTs while rejecting a large fraction of the background events. Table 2 presents the fraction of events, photon and background, that remain after each rejection technique is subsequently applied. (The numbers in Table 2 assume only $`W_1`$ is tested for $`\beta ^{}`$ decay energies.) Development of these event reconstruction techniques allows realistic evaluation of the performace and sensitivity of GCT designs. Our next goal is to simulate the efficiency, resolution, background and sensitivity of several Compton telescope configurations, utilizing the event reconstruction techniques developed here to realistically determine the performace of these instruments. CKD rejection has been shown to be the most efficient background rejection technique; however, the addition of effective TOF, backscatter, nonlocalized $`\beta ^{}`$ decay, and positron signature tests dramatically improve background rejection capabilities. We anticipate that use of these techniques will achieve overall sensitivity improvements in GCTs by factors of $`510`$.
###### Acknowledgements.
S. E. Boggs would like to thank the Millikan Postdoctoral Fellowship Program, CIT Deparment of Physics, for support.
## Appendix A Example GCT configuration
While it is not the intention of this paper to fully characterize the performance of a specific GCT, it is useful to have a telescope model for which the results of the event reconstruction can be presented. The telescope configuration modeled in this study is presented in Figure 1. The instrument consists of five planar arrays of 15 mm thick germanium, each of area $`100cm\times 100cm`$. In reality each array would consist of separate smaller detectors $`(5cm\times 5cm)`$ tiled to form the entire plane; however, the simulation performed here modeled each plane as a solid detector for simplicity. The five planar arrays are spaced 20 cm apart.
This configuration differs from historical Compton telescope configurations which generally consist of two detector planes separated by 100-150 cm. This separation distance is determined by the spatial resolution in z and the desired angular resolution. As will be discussed in a second paper, the configuration modeled here significantly improves the effective area of the telescope by letting each plane act as converter, and permitting a much wider range of scatter angles to produce good events. Allowing large-angle scatters also significantly increases the instrument FOV, and limits the effects of point spread function smearing for sources at large off-axis angles. The potential drawbacks of this configuration are increased background and degraded angular resolution.
The instrument was simulated using CERN’s GEANT Monte Carlo code. The Monte Carlo simulation produces a file of interaction locations and energy depositions for each photon/$`\beta `$decay event. Before performing event reconstruction on the interactions, the simulated events are modified to reflect realistic measurement uncertainties of an instrument: for each interaction, a random Gaussian-distributed uncertainty is added to the energy and position of each interaction. All interaction locations which lie within twice the instrumental spatial resolution of each other are combined into a single interaction site, to accurately reflect the resolving power of the detectors. Finally, interaction sites with energy deposits below the assumed detector threshold of 10 keV are ignored.
Two components are assumed to add in quadrature to determine the energy resolution: (i) a constant electronic noise, $`W_e=1.0keVFWHM`$, and (ii) the intrinsic resolution $`W_i`$ determined by the germanium Fano factor, $`F=0.13`$, and average free electron-hole pair energy, $`\epsilon =2.98eV`$, giving $`W_i=2.35\sqrt{F\epsilon E}FWHM`$. This corresponds to a resolution $`1.8keVFWHM`$ at 1 MeV, which is optimistic but not unrealistic. It is assumed that charge trapping and ballistic deficit do not significantly alter this energy resolution. The two components as well as the total energy resolution are shown in Figure A1.
It is assumed that two components add in quadrature to determine the $`1D`$ spatial resolutions, $`\delta x,\delta y,\delta z`$, of the detectors: (i) the range of the recoil electrons in the detector, and (ii) the positioning limits of the detector due to physical segmentation and/or signal analysis. Calculated electron ranges in germanium for different energies (Mukoyama (1976)) are used here as the 1-D FWHM positional uncertainties, $`\delta x_e,\delta y_e,\delta z_e`$. Methods to determine the event position by physically segmenting the GeD contacts into cross strips or pixels (Luke et al. (1994); Kroeger et al. (1996)), as well as using advanced signal processing to interpolate to even better positions (Boggs (1998); Luke et al. (1994)), are currently active fields of research – so this component of the spatial resolution remains speculative for now. Here it is assumed that signal processing will allow positional resolutions of $`0.5mmFWHM`$ at 100 keV, and that the discrimination capabilities go as the signal-to-noise ratio of the induced detector signal to electronic noise, i.e. as the inverse power of the interaction energy. It is also assumed that there is $`1mm`$ physical segmentation of the detector contacts in x, y, so that this component never exceeds this value. The z uncertainty, however, is not constrained by any such segmentation at the lowest energies. Therefore, the signal processing uncertainty is given by $`\delta x_s,\delta y_s,\delta z_s0.50(E/100keV)^1mmFWHM`$, maximizing at 1 mm in x, y below 50 keV, and approaching, but never maximizing at 15 mm in z at low energies. The two components as well as the total spatial resolution are shown in Figure A1.
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# Applying MDL to Learning Best Model GranularityParts of this work were presented at IJCAI, 1989, and at ESANN, 1994.
## 1 Introduction
It is commonly accepted that all learning involves compression of experimental data in a compact ‘theory’ , ‘hypothesis’, or ‘model’ of the phenomenon under investigation. In the last two authors analysed the theory of such approaches related to shortest effective description length (Kolmogorov complexity). The question arises whether these theoretical insights can be directly applied to real world problems. Selecting models on the basis of compression properties ignores the ‘meaning’ of the models. Therefore we should aim at optimizing a model parameter that has no direct semantics, such as the precision at which we represent the other parameters: too high precision causes accidental noise to be modeled as well, too low precision may cause models that should be distinct to be confused. In two quite different experimental settings the theoretically predicted values are shown to coincide with the best values found experimentally.
In general, the performance of a model for a given data sample depends critically on what we may call the “degree of discretization” or the “granularity” of the model: the choice of precision of the parameters, the number of nodes in the hidden layer of a neural network, and so on. The granularity is often determined ad hoc. Here we give a theoretical method to determine optimal granularity based on an application of the Minimum Description Length (MDL) principle in supervised learning. The cost of describing a set of data with respect to a particular model is the sum of the lengths of the model description and the description of the data-to-model error. According to MDL the best model is one with minimum cost, that is, the model that explains the data in the most concise way.
The first two authors, in , carried out an experiment in on-line learning to recognize isolated alphanumerical characters written in one subject’s handwriting, irrespective of size and orientation. Some novel features here are the use of multiple prototypes per character, and the use of the MDL principle to choose the optimal feature extraction interval. It is satisfactory that in this case the general learning theory can without much ado be applied to obtain the best sampling rate. We believe that the same method is applicable in a wide range of problems. To obtain evidence for this assertion, in the third author with G. te Brake and J. Kok applied the same method to modeling a robot-arm.
The genesis of this work is not rooted in traditional approaches to artificial intelligence (AI), but rather on new exciting general learning theories which have developed out from the computational complexity theory , statistics and descriptional (Kolmogorov) complexity . These new theories have received great attention in theoretical computer science and statistics, . On the other hand, the design of real learning systems seems to be dominated by ad hoc trial-and-error methods. Applications of these recent theoretical results to real world learning system design are scarce and far between. One exception is the elegant paper by Quinlan and Rivest, .
In a companion paper we develop the theory and mathematical validation of the MDL principle based on ultimate data compression up to the Kolmogorov complexity. Our purpose here is trying to bring theory and practice together by testing the theory on simple real applications. We give a brief accessible exposition of the theoretical background of MDL and the mathematical proof that it works. The principle is then applied to two distinct practical issues: that of on-line character recognition and of robot modeling. In both cases the issue is the supervised learning of best model granularity for classification or extrapolation on unseen examples. The systems and experiments are quite simple, and are intended to just be demonstrations that the theoretical approach works in practice. Both applications are in topical areas and the results give evidence that the theoretical approach can be extended to more demanding settings.
Contents: From a theoretical point of view we explain the general modeling principle of MDL and its relation to Bayesianism. We show that the MDL theory is solidly based on a provably ideal method of inference using Kolmogorov complexity and demonstrate that it is valid under the assumption that the set of data is “typical” (in a precise rigorous sense) for the targeted model. We then apply the theory to two experimental tasks that both concern supervised learning of the best model granularity but are otherwise quite dissimilar: In the first task we want to obtain the best model sampling rate in the design of an on-line hand-written character learning system, and in the second task our goal is to determine the best number of nodes in the hidden layer of a three-layer feedforward neural network modeling a robot arm with two degrees of freedom. It turns out that the theoretically predicted optimal precision coincides with the best experimentally determined precision. We conclude with a discussion concerning the equivalence of code-length based methods with probability-based methods, that is, of MDL and Bayesianism.
### 1.1 Introduction to Learning On-Line Handwritten Characters
One of the important aspects of AI research is machine cognition of various aspects of natural human languages. An enormous effort has been invested in the problem of recognizing isolated handwritten characters or character strings, . Recognizing isolated hand-written characters has applications, for example, in signature recognition and Chinese character input. The alphanumerical character learning experiment reported here is a pilot project which ultimately aims at providing a practicable method to learn Chinese characters. This problem knows many practical difficulties—both quantitatively and qualitatively.
There are several thousand independent Chinese characters. No current key-board input method is natural enough for casual users. Some of these methods require the user to memorize a separate code for each of seven thousand characters. Some methods require the user to know ping ying—the phonological representation of mandarin Chinese in latin characters. But then the translation into computer representation of characters is not easy because there are too many homonyms. Similarly, sound recognition techniques do not help much either because almost every commonly used Chinese character has more than one commonly used homonym. For non-professional casual users, hand-written input which is mechanically scanned and processed, seems to be a quite reasonable way of entering data into the computer.
A variety of approaches and algorithms have been proposed to achieve a high recognition rate. The recognition process is usually divided into two steps:
1. feature extraction from the sample characters, and
2. classification of unknown characters.
The latter often uses either deterministic or statistical inference based on the sample data, and various known mathematical and statistical approaches can be used. Our contribution is on that level. This leaves the technical problem of feature extraction, whose purpose is to capture the essence from the raw data. The state of the art is more like art than science. Here we use existing methods which we explain now.
#### 1.1.1 Feature Extraction
The most common significant features extracted are center of gravity, moments, distribution of points, character loci, planar curve transformation, coordinates, slopes and curvatures at certain points along the character curve. The obvious difficulty of the recognition task is the variability involved in handwritten characters. Not only does the shape of the characters depend on the writing style which varies from person to person, but even for the same person trying to write consistently writing style changes with time.
One way to deal with this problem is the idea of ‘elastic matching’ . Roughly speaking, the elastic matching method takes the coordinates or slopes of certain points approximately equally spaced along the curve of the character drawing as feature to establish the character prototypes. To classify an unknown character drawing, the machine compares the drawing with all the prototypes in its knowledge base according to some distance function and the entered character is classified as the character represented by the closest prototype. When an unknown character is compared to a prototype, the comparison of the features is not only made strictly between the corresponding points with the prototype but also with the points adjacent to the corresponding point in the prototype. The method we use for feature extraction is a new modification of the standard elastic matching method.
#### 1.1.2 Classification
Each implemented algorithm for character recognition embodies a model or a family of models for character recognition. One problem of the extant research is the lack of a common basis for evaluation and comparison among various techniques. This is especially true for on-line character recognition due to the lack of common standard and limited raw data source. A model from a family of models induced by a particular method of feature extraction is usually specified by a set of parameters. Varying the parameters gives a class of models with similar characteristics.
Consider the above mentioned elastic matching. It uses certain points along the character curve as features. The interval size used to extract these points along the curve is a parameter. How to determine the value of this parameter which gives optimal recognition?
Practically speaking, we can set the interval size to different values and experiment on a given sample set of data to see which value gives the best performance. However, since the experiment is based on one particular set of data, we do not know if this interval size value gives a similar optimal performance for all possible observations from the same data source. A theory is needed to guide the parameter selection in order to obtain the best model from the given class of models.
#### 1.1.3 Model Selection
Suppose we have models $`M_1,M_2,\mathrm{}`$. Let $`H_i`$ be the hypothesis ‘$`M_i`$ gives the best recognition rate’. Our problem in selecting the best model consists in finding the most likely hypothesis. We use the Minimum Description Length principle (referred as MDL hereafter) for this purpose. MDL finds its root in the well-known Bayesian inference and not so well-known Kolmogorov complexity.
Below we give the classic “Bayes’s rule.” According to Bayes’s rule, a specific hypothesis is preferred if the probability of that hypothesis takes maximum value for a given set of data and a given prior probability distribution over the set of hypotheses. This happens for the hypothesis under which the product of the conditional probability of the data for the given hypothesis and the prior probability of the hypothesis is maximal. When we take the negative logarithm of Bayes’s formula, then this maximal probability is achieved by the hypothesis under which the sum of the following two terms is minimized: the description length of the error of the data for the given hypothesis and the description length of the model (the hypothesis). Therefore, finding a maximum value of the conditional probability of a given set of hypotheses and data becomes minimizing the combined complexity or description length of the error and the model for a given set of candidate models.
To quantify this idea, the two description lengths are expressed in terms of the coding length of the model (set of prototypes) and the coding length of the error (combined length of all data failed to be described by the model). The trade-off between simplicity and complexity of both quantities is as follows.
1. If a model is too simple, in the sense of having too short an encoding, it may fail to capture the essence of the mechanism generating the data, resulting in increased error coding lengths.
2. If a model is too complicated, in the sense of having too long code length (like when it consists of a table of all data), it may contain a lot of redundance from the data and become too sensitive to minor irregularities to give accurate predictions of the future data.
The MDL principle states that among the given set of models, the one with the minimum combined description lengths of both the model and the error for given set of data is the best approximation of the mechanism behind data and can be used to predict the future data with best accuracy.
The objective of this work is to implement a small system which learns to recognize handwritten alphanumericals based on both elastic matching and statistical inference. MDL is used to guide the model selection, specifically the selection of the interval of feature extraction. The result is then tested experimentally to validate the application of the theory.
### 1.2 Introduction to Modeling a Robot Arm
We consider the problem of modeling a robot arm consisting of two joints and two stiff limbs connected as: joint, limb, joint, limb. The entire arm moves in a fixed two-dimensional plane. The first joint is fixed at the origin. The position of the other end of the arm is determined by the lengths of the two limbs, together with the angle of rotation in the first joint of the first limb, and the angle of rotation in the second joint of the second limb with respect to the first limb. The mobile end of the arm thus has two degrees of freedom given by the two angles. In this problem is modeled using a Bayesian framework to obtain a backpropagation network model. We use MDL to obtain the best number of nodes in the hidden layer of a three layer feedforward network model. The method is essentially the same as in the character recognition experiment. Just as before it is validated on a test set on unseen data.
## 2 Theoretic Preliminaries
We first explain the idea of Bayesian reasoning and give “ideal MDL” as a noncomputable but provably good approach to learning. Then, we will dilute the approach to obtain a feasible modification of it in the form of the real MDL. For another viewpoint of the relation between Bayesianism and MDL see . In the later sections we apply MDL to learning best model granularity.
### 2.1 Bayesianism
Bayesianism is an induction principle with a faultless derivation, yet allows us to estimate the relative likelihood of different possible hypotheses—which is hard or impossible with the commonly used Pearson-Neyman testing. With the latter tests we accept or reject a zero hypothesis with a given confidence. If we reject the zero hypothesis, then this does not mean that we do accept the alternative hypothesis. We cannot even use the same data to test the alternative hypthesis. (Or a subhypothesis of the alternative hypothesis—because note that all hypotheses different from the zero hypothesis must be taken together to form the alternative hypothesis.) In fact, this type of testing does not establish the relative likelihood between competing hypotheses at all.
###### Definition 1
Consider a discrete sample space $`\mathrm{\Omega }`$. Let $`D,H_1,H_2,\mathrm{}`$ be a countable set of events (subsets) of $`\mathrm{\Omega }`$. $`𝐇=\{H_1,H_2,\mathrm{}\}`$ is called hypotheses space . The hypotheses $`H_i`$ are exhaustive (at least one is true). From the definition of conditional probability, that is, $`P(A|B)=P(AB)/P(B)`$, it is easy to derive Bayes’s formula (rewrite $`P(AB)`$ in two different ways):
$$P(H_i|D)=\frac{P(D|H_i)P(H_i)}{P(D)}.$$
(1)
If the hypotheses are mutually exclusive ($`H_iH_j=\mathrm{}`$ for all $`i,j`$), then
$$P(D)=\underset{i}{}P(D|H_i)P(H_i).$$
Despite the fact that Bayes’s rule is just a rewriting of the definition of conditional probability, its interpretation and applications are most profound and have caused bitter controversy over the past two centuries. In Equation 1, the $`H_i`$’s represent the possible alternative hypotheses concerning the phenomenon we wish to discover. The term $`D`$ represents the empirically or otherwise known data concerning this phenomenon. The term $`P(D)`$, the probability of data $`D`$, may be considered as a normalizing factor so that $`_iP(H_i|D)=1`$. The term $`P(H_i)`$ is called the a priori probability or $`initial`$ probability of hypothesis $`H_i`$, that is, it is the probability of $`H_i`$ being true before we see any data. The term $`P(H_i|D)`$ is called a posteriori or inferred probability In model selection we want to select the hypothesis (model) with the maximum a posteriori probability (MAP). <sup>1</sup><sup>1</sup>1If we want to predict then we determine the expected a posteriori probability by integrating over hypotheses rather than choosing one hypothesis which maximises the posterior.
The most interesting term is the prior probability $`P(H_i)`$. In the context of machine learning, $`P(H_i)`$ is often considered as the learner’s initial degree of belief in hypothesis $`H_i`$. In essence Bayes’s rule is a mapping from a priori probability $`P(H_i)`$ to a posteriori probability $`P(H_i|D)`$ determined by data $`D`$. In general, the problem is not so much that in the limit the inferred hypothesis would not concentrate on the true hypothesis, but that the inferred probability gives as much information as possible about the possible hypotheses from only a limited number of data. In fact, the continuous acrimonious debate between the Bayesian and non-Bayesian opinions centered on the prior probability. The controversy is caused by the fact that Bayesian theory does not say how to initially derive the prior probabilities for the hypotheses. Rather, Bayes’s rule only tells how they are to be updated. In the real world problems, the prior proabilities may be unknown, uncomputable, or even conceivably non-existent. (What is the prior probability of use of a word in written English? There are many different sources of different social backgrounds living in different ages.) This problem would be solved if we can find a single probability distribution to use as the prior distribution in each different case, with approximately the same result as if we had used the real distribution. Surprisingly, this turns out to be possible up to some mild restrictions.
### 2.2 Kolmogorov Complexity
So as not to divert from the main thrust of the paper, we recapitulate the basic formal definitions and notations in Appendixes A, B. Here we give an informal overview.
Universal description length: For us, descriptions are finite binary strings. Since we want to be able determine where a description ends, we require that the set of descriptions is a prefix code: no description is a proper initial segment (proper prefix) of another description. Intuitively, the prefix Kolmogorov complexity of a finite object $`x`$ conditional $`y`$ is the length $`K(x|y)`$ in bits of the shortest effective description of $`x`$ using $`y`$ as input. Thus, for every fixed $`y`$ the set of such shortest effective descriptions is required to be a prefix code. We define $`K(x)=K(x|ϵ)`$ where $`ϵ`$ means “zero input”. Shortest effective descriptions are “effective” in the sense that we can compute the described objects from them. Unfortunately, , there is no general method to compute the length of a shortest description (the prefix Kolmogorov complexity) from the object being described. This obviously impedes actual use. Instead, one needs to consider computable approximations to shortest descriptions, for example by restricting the allowable approximation time. This course is followed in one sense or another in the practical incarnations such as MDL. There one often uses simply the Shannon-Fano code , which assigns prefix code length $`l_x:=\mathrm{log}P(x)`$ to $`x`$ irrespective of the regularities in $`x`$. If $`P(x)=2^{l_x}`$ for every $`x\{0,1\}^n`$, then the code word length of an all-zero $`x`$ equals the code word length of a truly irregular $`x`$. While the Shannon-Fano code gives an expected code word length close to the entropy, it does not distinguish the regular elements of a probability ensemble from the random ones, by compressing individual regular objects more than the irregular ones. The prefix code consisting of shortest prefix-free programs with the prefix Kolmogorov complexities as the code word length set does both: for every computable distribution $`P`$ the $`P`$-expected code-word length (prefix Kolmogorov complexity) is close to the entropy of $`P`$ as well as that every individual element is compressed as much as is possible using an effective code.
Universal probability distribution: Just as the Kolmogorov complexity measures the shortest effective description length of an object, the algorithmic universal probability $`𝐦(x|y)`$ of $`x`$ conditional $`y`$ measures the greatest effective probability of $`x`$ conditional $`y`$. It turns out that we can set $`𝐦(x|y)=2^{K(x|y)}`$, (2) in the Appendix B. For precise definitions of the notion of “greatest effective probability” the reader is referred to this appendix as wel. It expresses a property of the probability of every individual $`x`$, rather than entropy which measures an “average” or “expectation” over the entire ensemble of elements but does not tell what happens to the individual elements. <sup>2</sup><sup>2</sup>2 As an aside, for every fixed conditional $`y`$ the entropy $`_x𝐦(x|y)\mathrm{log}𝐦(x|y)=\mathrm{}`$. We will use the algorithmic universal probability as a universal prior in Bayes’s rule to analyze ideal MDL.
Individual randomness: The common meaning of a “random object” is an outcome of a random source. Such outcomes have expected properties but particular outcomes may or may not possess these expected properties. In contrast, we use the notion of randomness of individual objects. This elusive notion’s long history goes back to the initial attempts by von Mises, , to formulate the principles of application of the calculus of probabilities to real-world phenomena. Classical probability theory cannot even express the notion of “randomness of individual objects.” Following almost half a century of unsuccessful attempts, the theory of Kolmogorov complexity, , and Martin-Löf tests for randomness, , finally succeeded in formally expressing the novel notion of individual randomness in a correct manner, see . Every individually random object possesses individually all effectively testable properties that are only expected for outcomes of the random source concerned. It is “typical” or “in general position” in that it will satisfy all effective tests for randomness— known and unknown alike. A major result states that an object $`x`$ is individually random with respect to a conditional probability distribution $`P(|y)`$ iff $`\mathrm{log}(𝐦(x|y)/P(x|y))`$ is close to zero. In particular this means that $`x`$ is “typical” or “in general position” with respect to conditional distribution $`P(|y)`$ iff the real probability $`P(x|y)`$ is close to the algorithmic universal probability $`𝐦(x|y)=2^{K(x|y)}`$. That is, the prefix Kolmogorov complexity $`K(x|y)`$ is close to the Shannon-Fano code length of $`x`$ as element of the a set with probability distribution $`P(|y)`$.
For example, if $`H`$ is the hypothesis that we deal with a fair coin and the data sample $`D`$ is a hundred outcomes ‘heads’ in a row, then $`D`$ isn’t typical for $`H`$. But if $`D`$ is a truly random individual sequence with respect to $`H`$ (a notion that has a precise formal and quantitative meaning), then $`D`$ is typical for $`H`$. The probability of atypical sequences is very small and goes to zero when the data sample grows unboundedly.
Prediction and Model Selection: It has been shown by Solomonoff that the continuous variant of $`𝐦`$ has astonishing performance in predicting sequences where the probability of the next element is computable from the initial segment. We now come to the punch line: Bayes’s rule using the algorithmic universal prior distribution, suggested by Solomonoff already in , yields Occam’s Razor principle and is rigorously shown to work correctly in the companion paper . Namely, it is shown that this implies that data compression is almost always the best strategy, both in hypothesis identification and prediction.
### 2.3 Minimum Description Length Principle
Scientists formulate their theories in two steps. Firstly, a scientist, based on scientific observations, formulates alternative hypotheses (there can be an infinity of alternatives), and secondly a definite hypothesis is selected. The second step is the subject of inference in statistics. Historically this was done by many different principles, like Fisher’s Maximum Likelihood principle, various ways of using Bayesian formula (with different prior distributions). Among the most dominant ones is the ‘common sense’ idea of applying Occam’s razor principle of choosing the simplest consistent theory. But what is “simple”? We equate “simplicity” with “shortness of binary description,” thus reducing the razor to objective data compression.
However, no single principle is both theoretically sound and practically satisfiable in all situations. Fisher’s principle ignores the prior probability distribution (of hypotheses). To apply Bayes’s rule is difficult because we usually do not know the actual prior probability distribution. (What is the prior distribution of words in written English, where there are many sources of many ages and social classes?) No single principle turned out to be satisfiable in all situations. Philosophically speaking, relative shortness achievable by ultimate data compression presents an ideal way of solving induction problems. However, due to the non-computability of the Kolmogorov complexity and the associated algorithmic universal prior function, such a theory cannot be directly used. Some approximation is needed in the real world applications.
Rissanen follows Solomonoff’s idea, but substitutes a ‘good’ computable approximation to $`𝐦(x)`$ to obtain the so-called Minimum Description Length principle. Rissanen not only gives the principle, more importantly he also gives the detailed formulas on how to use this principle. This made it possible to use the MDL principle. The basic form of the MDL principle can be intuitively stated as follows:
> Minimum Description Length Principle. The best theory to explain a set of data is the one which minimizes the sum of
> $``$ the length, in bits, of the description of the theory;
> $``$ the length, in bits, of data when encoded with the help of the theory.
A survey of the development of the MDL principle in statistical inference and its applications is given in . In the relationship between the Bayesian approach and the minimum description length approach is established. The general modeling principle MDL is sharpened and clarified, abstracted as the ideal MDL principle and defined from Bayes’s rule by means of Kolmogorov complexity. The argument runs as follows:
Given a data sample and a family of models (hypotheses) one wants to select the model that produced the data. A priori it is possible that the data is atypical for the model that actually produced it. Meaningful induction is possible only by ignoring this possibility. Strictly speaking, selection of a “true” model is improper usage, “modeling the data” irrespective of truth and falsehood of the resulting model is more appropriate. In fact, given data sample and model class the truth about the models is impossible to ascertain and modeling as well as possible is all we can hope for. Thus, one wants to select a model for which the data is typical. The best models make the two-part description of the data using the model as concise as possible. The simplest one is best in accordance with Occam’s razor principle since it summarizes the relevant properties of the data as concisely as possible. In probabilistic data or data subject to noise this involves separating regularities (structure) in the data from random effects.
From Bayes’s Formula 1, we must choose the hypothesis $`H`$ that maximizes the posterior $`P(H|D)`$. Taking the negative logarithm on both sides of Equation 1:
$$\mathrm{log}P(H|D)=\mathrm{log}P(D|H)\mathrm{log}P(H)+\mathrm{log}P(D).$$
Here, $`\mathrm{log}P(D)`$ is a constant and can be ignored because we just want to optimize the left-hand side of the equation over $`H`$. Maximizing the $`P(H|D)`$’s over all possible $`H`$’s is equivalent to minimizing $`\mathrm{log}P(H|D)`$, that is, minimizing
$$\mathrm{log}P(D|H)\mathrm{log}P(H).$$
To obtain the ideal MDL principle it suffices to replace the terms in the sum by $`K(D|H)`$ and $`K(H)`$, respectively. In view of (2) in the Appendix B this is justified provided $`\mathrm{log}P(D|H)\stackrel{+}{=}\mathrm{log}𝐦(D|H)`$ and also $`\mathrm{log}P(H)\stackrel{+}{=}\mathrm{log}𝐦(H)`$. In we show that the basic condition under which this substitution is justified is encapsulated as the Fundamental Inequality, which in broad terms states that the substitution is valid when the data are random, relative to every contemplated hypothesis and also these hypotheses are random relative to the (universal) prior. Basically, the ideal MDL principle states that the prior probability associated with the model should be given by the algorithmic universal probability, and the sum of the log universal probability of the model plus the log of the probability of the data given the model should be minimized. For technical reasons the latter probability $`P(D|H)`$ must be computable.
It is important to note that using the algorithmic universal prior we compress every model $`H`$ to its prefix Kolmogorov complexity $`K(H)=\mathrm{log}𝐦(H)`$. Applying the ideal MDL principle then compresses de description of the data encoded using the model, $`D|H`$, to its prefix Kolmogorov complexity $`K(D|H)=\mathrm{log}𝐦(D|H)`$ as well for the model $`H`$ minimizing the sum of the two complexities. Roughly speaking, the MDL selection assumes that the data set is “typical” for the selected model $`H`$. Thus, MDL aims at selecting a model for which the data are “typical”, even if there happened to be a different “true” model that inappropriately generated “atypical” data. In this manner application of MDL is resilient to overfitting the model.
### 2.4 Ideal MDL versus Real MDL
Using the algorithmic universal prior, the ideal MDL principle is valid for a set of data samples of Lebesgue measure one, the “random”, “typical” outcomes, for every contemplated hypothesis. For these “typical” outcomes we have $`K(D|H)\stackrel{+}{=}\mathrm{log}P(D|H)`$ which means that the classic Shannon-Fano code length reaches the prefix Kolmogorov complexity on these data samples. The Shannon-Fano code that assigns code words of length $`\stackrel{+}{=}\mathrm{log}P()`$ to elements randomly drawn according to a probability density $`P()`$ is in fact used in the applied statistical version of MDL. Thus, under the assumption that the data sample is typical for the contemplated hypotheses, the ideal MDL principle and the applied statistical one coincide, and moreover, both are valid for a set of data samples of Lebesgue measure one . The latter result has been obtained in the statistical theory using probabilistic arguments .
The term $`\mathrm{log}P(D|H)`$, also known as the self-information in information theory and the negative log likelihood in statistics, can now be regarded as the number of bits it takes to redescribe or encode $`D`$ with an ideal code relative to $`H`$. In different applications, the hypothesis $`H`$ can mean many different things, such as decision trees, finite automata, Boolean formulas, or a polynomial.
###### Example 1
In general statistical applications, one assumes that $`H`$ is some model $`H(\theta )`$ with a set of parameters $`\theta =\{\theta _1,\mathrm{},\theta _k\}`$ of precision $`c`$, where the number $`k`$ may vary and influence the descriptional complexity of $`H(\theta )`$. For example, if we want to determine the distribution of the length of beans, then $`H`$ is a normal distribution $`N(\mu ,\sigma )`$ with parameters median $`\mu `$ and variation $`\sigma `$. So essentially we have to determine the correct hypothesis described by identifying the type of distribution (normal) and the correct parameter vector $`(\mu ,\sigma )`$.
In such cases, we minimize
$$\mathrm{log}P(D|\theta )\mathrm{log}P(\theta ).$$
###### Example 2
Let’s consider the fitting of a ‘best’ polynomial on $`n`$ given sample points in the 2-dimensional plane. This question is notoriously underdefined, since both a 1st degree polynomial with $`\chi ^2`$ best fit, and a $`(n1)`$th degree polynomial with perfect fit are arguably the right solutions. But with the MDL principle we can find an objective ‘best’ polynomial among the polynomials of all degrees.
For each fixed $`k`$, $`k=0,\mathrm{},n1`$, let $`f_k`$ be the best polynomial of degree $`k`$, fitted on points $`(x_i,y_i)`$ ($`1in`$), which minimizes the error
$$error(f_k)=\underset{i=1}{\overset{n}{}}(f_k(x_i)y_i)^2.$$
Assume each coefficient takes $`c`$ bits. So $`f_k`$ is encoded in $`c(k+1)`$ bits. Let us interpret the $`y_i`$’s as measurements for argument $`x_i`$ of some true polynomial to be determined. Assume that the measurement process involves errors. Such errors are accounted for by the commonly used Gaussian (normal) distribution of the error on $`y_i`$’s. Thus, given that $`f`$ is the true polynomial,
$$\mathrm{Pr}(y_1,\mathrm{},y_n|f,x_1,\mathrm{},x_n)=\mathrm{exp}(O((f(x_i)y_i)^2)).$$
The negative logarithm of above is $`c^{}error(f)`$ for some computable $`c^{}`$. The MDL principle tells us to choose $`f=f_m`$, with $`m\{0,\mathrm{},n1\}`$, which minimizes $`c(m+1)+c^{}error(f_m)`$.
In the original Solomonoff approach a hypothesis $`H`$ is a Turing machine. In general we must avoid such a too general approach in order to keep things computable. In different applications, $`H`$ can mean many different things. For example, if we infer decision trees, then $`H`$ is a decision tree. In case of learning Boolean formulas, then $`H`$ may be a Boolean formula. If we are fitting a polynomial curve to a set of data, then $`H`$ may be a polynomial of some degree. In Experiment 1 below, $`H`$ will be the model for a particular character. Each such $`H`$ can be encoded by a binary string from a prefix-free set, where a set of codes is prefix-free if no code in the set is a prefix of another.
## 3 Experiment 1: On-Line Handwritten Characters
### 3.1 Model Development
#### 3.1.1 Basic Assumptions
When an alphanumeral character is drawn on a planar surface, it can be viewed as a composite planar curve, the shape of which is completely determined by the coordinates of the sequence of points along the curve. The order of the sequence is determined by on-line processing the data from the scanning machinery at the time of writing the character. Since the shape tends to vary from person to person and from time to time, so do the coordinates of the point sequence. A key assumption in our treatment is that for a particular person writing consistently the shape of the curve tends to converge to an average shape, in the sense that the means of corresponding coordinates of the sampled point sequences converge.
That is, we assume:
* Each shape of an example curve for a particular character contains a set distinguished feature points.
* For each such point, the average of the instances in the different examples converges to a mean.
* Moreover, there is a fixed probability distribution (possibly unknown) for each such point which is symmetric about the mean value, and the variance is assumed to be the same for all the character drawings.
Essentially, we only assume that one person’s hand-writing has a fixed associated probability distribution, which does not change.
#### 3.1.2 Feature Space, Feature Extraction and Prototypes
A Kurta ISONE digitizer tablet with 200/inch resolution in both horizontal and vertical directions was used as the scanner to obtain and send the coordinates of the sequential points of the character curve on the tablet to the microprocessor of a IBM PS/2 model 30 computer. The system was implemented using programming language C. The coordinates were normalized on a 30 by 30 grid in horizontal and vertical directions.
The character drawn on the tablet is processed on-line. The sequence of the coordinates in order of time of entry is stored in the form of a linked list. This list is preprocessed in order to remove the repeating points due to hesitation at the time of writing, and to fill in the gaps between sampled points resulted from the sampling rate limit of the tablet.
The latter needs some explanation: the digitizer has a maximum sampling rate of 100 points/second. If a person writes a character in 0.2 seconds, only 20 points on the character curve will be sampled, leaving gaps between those points.
The preprocessing procedure ensures that in the resulting linked list any pair of consecutive points on the curve has at least one component of coordinates (stored as integers between 0 and 30) differing by 1 and no coordinate component differing by more than 1. For a preprocessed list $`((x_1,y_1),\mathrm{},(x_n,y_n))`$ therefore we have that for all $`i`$ ($`1i<n`$)
$`|x_ix_{i+1}|+|y_iy_{i+1}|`$ $``$ $`1`$
$`|x_ix_{i+1}|,|y_iy_{i+1}|`$ $``$ $`1`$
The preprocessed curve coordinate list is then sent to the feature extraction process. So far the coordinates are still integers in the range of 0 to 30.
The coordinates of certain points along the character curves are taken as relevant features. Feature extraction is done as follows. A character may consist of more than one stroke (a stroke is the trace from a pen drop-down to pen lift-up), the starting and ending points of every stroke are mandatorily taken as features. In between, feature points are taken at a fixed interval, say, one point for every $`n`$ points along the preprocessed curve, where $`n`$ is called feature extraction interval. This is to ensure that the feature points are roughly equally spaced. Actually the Euclidean length between any two points on the stroke curve, excluding the last point of a stroke, varies from $`n`$ to $`\sqrt{2}n`$ (for the diagonal).
The sequence of the feature point coordinates extracted from a given character drawing constitute a feature vector. (If the character drawing contains more than one stroke, its feature vector consists of the concatenation of the feature vectors of the individual strokes in time-order.) The dimension of the feature vector is the number of entries in it—or rather twice that number since each entry has two coordinate components. Obviously the dimension of the feature vector is also a random variable since the shape and the total number of points on the character curve varies from time to time. The dimension of the feature vector is largely determined by the feature extraction interval.
The extracted feature vector of a character is viewed as a prototype of a character, and is stored in the knowledge base of the system as such.
#### 3.1.3 Comparison between Feature Vectors
Before the system is employed to recognize characters, it must first learn them. It is trained with examples of the character drawings from the same data source which it is supposed to recognize afterwards. Here the ‘same data source’ means the same person writing consistently. The basic technique used in both training and recognition is the comparison or matching between prototypes or feature vectors. To compare any two prototypes or feature vectors of equal dimension, we can simply take the Euclidean distance between the two vectors. Mathematically this means subtracting each component of one vector from its corresponding component in the other feature vector, summing up the square of the differences and taking the square root of the sum. If the two prototypes are $`\chi =((x_1,y_1)\mathrm{},(x_n,y_n))`$ and $`\chi ^{}=((x_1^{},y_1^{})\mathrm{},(x_n^{},y_n^{}))`$, then the distance between them is
$$\sqrt{\underset{i=1}{\overset{n}{}}(x_ix_i^{})^2+(y_iy_i^{})^2}.$$
The knowledge base of the system is a collection of feature vectors stored in the form of a linked list. Each such feature vector is an example of a particular character and is called a prototype for that character. Each newly entered character drawing in the form of a feature vector is compared to the prototypes in the knowledge base. But we do not (cannot) assume that all feature vectors extracted from examples of the same character will have the same dimension.
Therefore, the comparison technique used in our system follows the spirit of this mathematical definition but is more elastic. As a consequence, corresponding feature points may be located in different places in the feature vectors. This problem is solved by so-called elastic matching which compares a newly sampled feature vector with the set of stored feature vectors, the prototypes. The elasticity is reflected in two aspects:
is a constant integer $`T_d`$ such that the new feature vector is compared with all stored feature vector of which the dimension is not more than $`T_d`$ different. That is, if the new feature vector has $`n`$ feature points, it will be compared (matched) with all the prototypes with a number of feature points in the range of $`[nT_d,n+T_d]`$.
is an integer constant $`N_e`$ such that the $`i`$th feature point of the new feature vector is compared with the feature points with index ranging from $`iN_e`$ to $`i+N_e`$ of each prototype satisfying the dimension tolerance. The least Euclidean distance found this way is considered to be the ‘true’ difference $`d_i`$ between the two vectors at $`i`$th feature point.
###### Definition 2
If the dimension of the new feature vector $`x`$ is $`n`$, then the elastic distance $`\delta (x,x^{})`$ between $`x`$ and a prototype $`x^{}`$ is defined as
$$\delta (x,x^{})=\sqrt{\underset{i=1}{\overset{n}{}}d_i^2}$$
if $`x^{}`$ is within the dimension tolerance $`T_d`$ of $`x`$, and $`\delta (x,x^{})=\mathrm{}`$ otherwise.
For our particular problem, experimental evidence indicates that it suffices to set both $`T_d`$ and $`N_e`$ to 1. In our experiment we used elastic distance between a new feature vector $`x`$ and a prototype $`x^{}`$ as computed above with $`T_d=N_e=1`$.
#### 3.1.4 Knowledge Base and Learning
The knowledge base is constructed in the learning phase of the system by sampling feature vectors of handwritten characters while telling the system which character it is an example of. Our system uses the following Learning Algorithm to establish the knowledge base.
Initialize the knowledge base $`S`$ to the empty set $`\mathrm{}`$. (The elements of $`S`$ will be triples $`(x,\chi ,c)`$ with $`x`$ a preprocessed feature vector (a prototype), $`\chi `$ is the character value of which $`x`$ is a prototype, and $`c`$ is a counter.) Assign values to weights $`\alpha ,\beta `$ (used later to combine prototypes) so that $`\alpha +\beta =1`$.
Sample a new example of a character feature vector, say $`x`$ after preprocessing, together with its character value, say $`\chi `$. (Actually, the user draws a new example handwritten character on the tablet and indicates the character value—which character the drawing represents—to the system.)
Check $`S`$ whether or not any prototypes exist for character $`\chi `$.
If there is no prototype for $`\chi `$ in $`S`$, then store $`x`$ in $`S`$ as a prototype for $`\chi `$ by setting
$$S:=S\{(x,\chi ,1)\}.$$
If $`P_x=\{y,\mathrm{},z\}`$ is a nonempty list of prototypes for $`\chi `$ in $`S`$, then determine elastic distances $`\delta (x,y)\mathrm{}\delta (x,z)`$. Let $`P_x^{min}P_x`$ be the set of prototypes in $`P_x`$ such that for all $`x^{}P_x^{min}`$ we have
$$\delta (x,x^{})=\delta _{min}\stackrel{\mathrm{def}}{=}\mathrm{min}\{\delta (x,y):yP_x\}.$$
Now $`x^{}P_x^{min}`$ may or may not be one of the prototypes with the character value $`\chi `$.
Step 2.1. If $`x^{}P_x^{min}`$ and $`(x^{},\chi ,m)S`$ (the minimum distance $`\delta _{min}`$ is between $`x`$ and one of the prototypes for $`\chi `$ in $`S`$), then
$$x^{}:=\alpha x^{}+\beta x;m:=m+1.$$
(The new prototype is combined with the existing prototype by taking the weighted average of every coordinate to produce a modified prototype for that character. Moreover, we add one to the counter associated with this prototype)
Step 2.2. If Step 2.1 is not applicable, then for no $`x^{}P_x^{min}`$ we have $`(x^{},\chi ,)S`$ (the minimum distance is between $`x`$ and prototype(s) in $`S`$ with character value $`\chi `$). Then we distinguish two cases.
Case 1. There is $`(x^{},\chi ,m)S`$ such that $`\delta (x,x^{})\delta _{min}(m+1)/m`$. (Here $`m`$ is number of character drawings which have consecutively be combined to form the current $`x^{}`$ prototype.) Then set
$$x^{}:=\alpha x^{}+\beta x;m:=m+1.$$
It is expected that the modified prototype will have minimal distance $`\delta _{min}`$ next time when a similar drawing of the same character value arrives.
Case 2. The condition in Case 1 is not satisfied. Then the new prototype will be saved in the knowledge base as a new prototype for the character by setting
$$S:=S\{(x,\chi ,1)\}.$$
Notice that more than one prototype for a single character is allowed in the knowledge base.
#### 3.1.5 Recognition of an Unknown Character Drawing
The recognition algorithm is simple. Assume the Learning Algorithm above has been executed, and a knowledge base $`S`$ with prototypes of all possible characters has been constructed. When an new character drawing is presented to the system, it is compared to all the prototypes in the knowledge base with dimension variation within the range specified by the dimension tolerance. The character of the prototype which has minimum $`\delta `$-distance from the presented feature vector is considered to be the character value of that feature vector. The rationale here is that the prototypes are considered as the ‘mean’ values of the feature vectors of the characters, and the variances of the distribution are assumed to be the same for all prototypes. Formally, the Recognition Algorithm is as follows.
Sample a new example of a character feature vector, say $`x`$ after preprocessing. (Actually, the user draws a new example handwritten character on the tablet which is preprocessed to form feature vector $`x`$.)
If $`S=\{(x_1,\chi _1,),\mathrm{},(x_n,\chi _n,)\}`$ is the knowledge base, then determine elastic distances $`\delta (x,x_1)\mathrm{}\delta (x,x_n)`$. If $`\delta (x,x_i)`$ is the minimal distance in this set, with $`i`$ is least in case more than one prototype induces minimum distance, then set $`\chi _i`$ is the character value for $`x`$ and
$$\text{Recognize character }\chi _i.$$
This concludes the main procedure of training and classification. A few remarks are in order to explain differences with the original elastic matching method.
1. This process differs from the original elastic matching method in the the way of prototype construction. More than one prototype are allowed for a single character. By our procedure in the Learning Algorithm, a prototype is the statistical mean of a number of positive examples of the character.
2. Every prototype is a feature vector which in turn is a point in the feature space of its dimension. Since the classification is based on statistical inference, the rate of correct classification depends not only on how well the prototypes in the knowledge base are constructed, but also on the variability of the handwriting of the subject. Even though more than one prototype is allowed for any character in the knowledge base, too many prototypes may result in an overly dense feature space. When the $`\delta `$-distance between two points (two prototypes in the knowledge base) in the feature space is comparable to the variability of the subjects handwriting, then the rate of correct classification may drop considerably.
3. The prototypes in the knowledge base constitute the model for the system. How well the prototypes are constructed will essentially determine the rate of correct classification and therefore the performance of the model. For the scheme described above, the prototypes are constructed by extracting points at a constant interval. Generally speaking, more points in the prototypes gives a more detailed image of the character drawing but may also insert random ‘noise’ in the model.
Application of MDL to guide the selection of ‘best’ feature extraction interval is the main thrust of this work, to which we proceed now.
### 3.2 Implemented Description Lengths and Minimization
The expression in MDL consists of two terms: the model and error coding lengths. The coding efficiency for both of these two terms must be comparable, otherwise minimizing the resulted expression of total description length will give either too complicated or too simple models. For this particular problem, the coding lengths are determined by practical programming considerations.
A set of 186 character drawings, exactly 3 for each of the 62 alphanumeral characters, were processed to feature vectors and presented to the Learning Algorithm, to form the raw database. The character drawings were stored in standardized integer coordinate system standardized from 0 to 30 in both $`x`$ and $`y`$ axis. After preprocessing as above, they were then input to the Learning Algorithm to establish a knowledge base: the collection of prototypes with normalized real coordinates, based on a selected feature extraction interval.
Subsequent to the construction of the knowledge base, the system was tested by having it classify the same set of character drawings using the Recognition Algorithm. This procedure served to establish the error code length and the model code length which are defined as follows.
###### Definition 3
The error code length or exception complexity is the sum of the total number of points for all the incorrectly classified character drawings. This represents the description of the data given the hypothesis.
The model code length or model complexity is the total number of points in all the prototypes in the machine’s knowledge base multiplied by 2. This represents the hypothesis.
The total code length is the sum of the error code length and the model code length.
###### Remark 1
The factor of 2 in the model code length is due to the fact that the prototype coordinates are stored as real numbers which takes twice as much memory (in programming language C) as the character drawing coordinates which are represented in integer form. One might wonder why the prototype coordinates are real instead of integer numbers. The reason is to facilitate the elastic matching to give small resolution for comparisons of classification.
Thus, both the model and error code lengths are directly related to the feature extraction interval. The smaller this interval, the more complex the model, but the smaller the error code length. The effect is reversed if the feature extraction interval goes toward larger values. Since the total code length is the sum of the two code lengths, there should be a value of feature extraction interval which minimizes the total code length.
This feature extraction interval is considered to be the ‘best’ one in the spirit of MDL. The corresponding model, the knowledge base, is considered to be optimal in the sense that it contains enough essence from the raw data but eliminates most redundancy due to noise from the raw data. This optimal feature extraction interval can be experimentally determined by carrying out the above described build-and-test (building the knowledge base and then test it based on the same set of characters on which it was built) for a number of different feature extraction intervals.
The actual optimization process was implemented on the actual system we constructed, and available to the user. For our particular set of characters and trial, the results of classifying by the Recognition Algorithm the same set of 186 character drawings used by the Learning Algorithm to establish the knowledge base, is given in Figure 1. Three quantities are depicted: the model code length, the error code length, and the total code length, versus different feature extraction intervals (FEATURE EXTRACTION INTERVAL in the figure).
For larger feature extraction intervals, the model complexity is small but most of the character drawings are misclassified, giving the very large error code length and hence the very large
total code length. On the other hand, when the feature extraction interval is at its low extremal value, all training characters get correctly classified which gives zero error coding length. But now the model complexity reaches its largest value, resulting also in a large total code length again. The minimum total code length occurred in our experiment at an extraction interval of 8, which gives 98.2 percent correct classification. Figure 2 illustrates the fraction of correctly classified character drawings for the training data.
### 3.3 Validation of the Model
Whether the ‘optimal’ model, determined by choosing the interval yielding minimal total code length for the training data, really performs better than models in the same class using different feature extraction intervals, can be tested by classification of new data—new character drawings.
We have executed such a test by having the set of 62 characters drawn anew by the same person who provided the raw data base to build the knowledge base. After preprocessing, the feature vectors resulting from these data were entered in the Recognition Algorithm. The new data are considered to be from the same source as the previous data set.
This new data set was classified by the system using the knowledge bases built by the Learning Algorithm from the training data set of 186 character drawings, based on different feature extraction intervals. The test results are plotted in Figure 3 in terms of the fraction of correct classification (CORRECT RATIO) versus feature extraction interval (FEATURE EXTRACTION INTERVAL) It is interesting to see that a 100% correct classification occurred at feature extraction intervals 5, 6 and 7. These values of feature extraction intervals are close to the optimal value 8 resulting from MDL considerations. Furthermore, at the lower feature extraction intervals, the correct classification rate drops, indicating the disturbance caused by too much redundance in the model. The recommended working feature extraction interval is thus either 7 or 8 for this particular type of character drawings.
## 4 Experiment 2: Modeling a Robot Arm
In the second experment, the problem is to model a two-jointed robot arm described in the introduction. A mathematical description is as follows. Let $`r_1`$ and $`r_2`$ be the lengths of the two limbs constituting the arm. One end of the limb of length $`r_1`$ is located in the joint at the origin $`(0,0)`$ of the two-dimensional plane in which the arm moves. The angle the limb makes with the horizontal axis is $`\theta _1`$. The angle the limb of length $`r_2`$, the second limb, makes with the first limb (in the second joint) is $`\theta _2`$. Then the relationship between the coordinates $`(y_1,y_2)`$ of the free end of the second limb (the hand so to speak) and the variables $`\theta _1,\theta _2`$ is given by
$`y_1`$ $`=`$ $`r_1\mathrm{cos}(\theta _1)+r_2\mathrm{cos}(\theta _1+\theta _2)`$
$`y_2`$ $`=`$ $`r_1\mathrm{sin}(\theta _1)+r_2\mathrm{sin}(\theta _1+\theta _2).`$
The goal is to construct a feedforward neural network that correctly associates the $`(y_1,y_2)`$ coordinates to the $`(\theta _1,\theta _2)`$ coordinates.
As in we set $`r_1=2`$ and $`r_2=1.3`$. The setup is similar to the character recognition experiment except that the data are not real-world but computer generated. We generated random examples of the relation between $`y_1,y_2`$ and $`\theta _1,\theta _2`$ as in the above formula and gaussian noise of magnitude 0.05 was added to the outputs.
Since we want the learned model to extrapolate from the training examples rather than interpolate between them, the training sets consist of random examples taken from two limited and separate areas of the domain. In the training data the first angle $`\theta _1`$ was in between 90 and 150 degrees or between 180 and 240 degrees, and the second angle was in between 30 and 150 degrees. To test extrapolation capabality of the learned model we used a unseen test set in which $`\theta _1`$ ranges between 0 and 270 and $`\theta _2`$ between 0 and 180 degrees.
### 4.1 Model Features
The model class consists of three-layer feedforward networks. The first layer is the input layer consisting of two input nodes with as input the real values of the two angles $`\theta _1,\theta _2`$. Both nodes in the input layer are connected with every node in the second layer—the hidden layer—of which the number of nodes is to be determined. Every node in the second layer is connected to both nodes in the third (output) layer, yielding the two real-valued output values $`y_1,y_2`$. There are no other connections between pairs of nodes. The second layer nodes have sigmoidial transfer functions and in the third layer output nodes have linear transfer functions. Thus, the only unknowns in the network are the number of nodes in the hidden layer, the weights on the connections and the biases of the nodes in the hidden layer. For every number of $`k`$ nodes ($`k=2,3,\mathrm{},15`$) in the hidden layer we learned the weights and biases of the network using standard methods by repeatingly presenting the training set. After that, the learned models are evaluated expermentally as to their prediction errors on the unseen test set.
During the experiments we noticed that if we used a test set from the same domain as the training set—thus testing interpolation rather than extrapolation—then the increase of error with increasing number of nodes in the hidden layer (after the optimal number) was small. For the unseen test set described earlier—testing extrapolation or generalization—the increase of the error after the optimal network size was more steep. Below we used the latter “generalization” test set.
### 4.2 Determining Size of Hidden Layer by MDL
We verify the contention that in this experimental setting the hypothesis selected by the MDL principle using the training data set can be expected to be a good predictor for the classification of unseen data from a test set.
Neural networks can be coded in the following way. Both the topology, the biases of the nodes, and the weights on the links are coded. Assume that the network contains $`k`$ nodes. The code starts with the number $`k`$. Next a list of $`k`$ bias values is encoded using $`l`$ bits for each bias value. We need $`k\times (k1)`$ bits to describe which pairs of nodes are connected by directed arcs (possibly in two ways). The weight for each link is given using a precision of $`l`$ bits. Concatenating all these descriptions in a binary string we can only retrieve the network if we can parse the constituent parts. Keeping the above order of the constituents we can do that if we know $`k`$. Therefore, we start the encoding with a prefix-free code for $`k`$ in $`\mathrm{log}k+2\mathrm{log}\mathrm{log}k`$ bits.<sup>3</sup><sup>3</sup>3This is standard in prefix-free coding, see Appendix A. The total description now takes at most $`\mathrm{log}k+2\mathrm{log}\mathrm{log}k+k\times l+k(k1)+m\times l`$ bits, where $`m`$ is the number of directed edges (links).
For three-layer feedforward networks that constitute our models, with two input nodes and two output nodes and $`k`$ nodes in the hidden layer, the topology is fixed. As already stated, we have to choose only the weights on the links and the biases of the hidden nodes. This gives descriptions of length $`\mathrm{log}k+2\mathrm{log}\mathrm{log}k+5k\times l`$ bits. For the range of $`k,l`$ we consider the logarithmic terms can be ignored. Thus, the model cost is set at $`5kl`$ bits, and with precision $`l=16`$ the model cost is linear in $`k`$ at $`80k`$ bits.
The encoding of the output data for the neural network, to determine the error cost in the MDL setting, depends on whether they are given as integers or reals. For integers one takes the 16 bits that the MaxInt format requires, and for real numbers usually twice as many, that is, 32 bits. For reals such an encoding introduces a new problem: when is the output correct? We consider it correct if the real distance between the output vector and the target vector is under a small fixed real value. In the MDL code every example consists of two reals which are encoded as 64 bits. Thus, the erroroneous examples, those exceeding the small fixed error cut-off level that we set, are encoded in 64 bits each. We ignore the amount with which the misclassified output real value differs from the target real value, it may be large or small. The total error is encoded as an explicit list of the misclassified examples.
Because MDL selects the model that minimizes sum of model length and total error length, it is important how large a training set we choose. The coding length of the models is the same for every fixed $`k`$ and training set size, but the total error length depends on this. For a small training set the number of erroneous (misclassified) examples may be very small compared to the model code length, and the difference between simple models with small $`k`$ and complex models is large. With large training sets the opposite happens. This is exactly right: with a small number of examples the simpler models are encouraged. How complex a model can be must be justified by the size of the training set. Intuitively, with increasing training set size, eventually the smallest model that has low error on this set can be expected to stabilize and to have low prediction error.
In the following experiment we used a random training set of 100 examples, and for every $`k`$ ($`2k15`$) a network with $`k`$ nodes in the hidden layer was trained with $`10^6`$ training cycles. Figure 4 shows the results in terms of MDL: the model code length, the error code length, and the total description length, as a function of the number of nodes in the hidden layer. The optmimum of the total code length is reached for seven hidden nodes, that is, MDL predicts that seven hidden nodes (the granularity of the hidden layer so to speak) give the best model. This is only one node away from the optimal network size determined experimentally below.
### 4.3 Validation of the Model
To determine the best number of hidden nodes we used thirty different random training sets of 100 examples each. For every $`k`$ ($`2k15`$) the network was trained using $`10^6`$ training cycles. Other more sophisticated stop criteria could have been used, but some checking showed that in general the performance of the network after $`10^6`$ training cycles was close to optimal. The error per example is the real distance between the output vector and the target vector. In figure LABEL:err:draw the average squared error on the training set and the average squared prediction error on the unseen test set are displayed as a function of the number of nodes in the hidden layer. The optimal network in the sense of having best extrapolation and generalization properties in modeling the unseen examples in the test set most correctly, is a network with 8 hidden nodes.
As expected, the error on the training set keeps on decreasing with an increasing number of nodes in the hidden layer, that is, when the model becomes increasingly complex and is capable of modeling more and more detail. When we look at the prediction error of new examples that were not in the training data (figure LABEL:err:draw), we see that the average squared prediction error first decreases when the model complexity increases, but that there is an optimum of minimum error after which the error starts to increase again. Experimentally, the best number of hidden nodes for this problem with a training set size of 100 examples is 8, that is, one more than predicted by the simplified application of MDL above.
## 5 Discussion
We applied the theoretical Minimum Description Length principle to two different experimental tasks aimed at learning the best model discretization. The first application was learning to recognize isolated handwritten characters on-line using elastic matching and some statistical technique. A ‘model’ is a collection of prototypes built from raw training character drawings by on-line taking points on the curves of the executed character drawing at a constant feature extraction interval, and by combining closely related character drawings. Some novel features here are the use of multiple prototypes per character, and the use of the MDL principle to choose the optimal feature extraction interval.
The model is optimized in the spirit of MDL by minimizing the total code length, which is the sum of the model and error-to-model code lengths, against different feature extraction intervals. The resulting model is optimal according to the theory. It is then validated by testing using a different set of character drawings from the same source. We believe that the result of this small test gives evidence that MDL may be a good tool in the area of handwritten character recognition.
The second application was modeling a robot arm by a three layer feedforward neural network, where the precision parameter to be learned is the number of nodes in the hidden layer. The MDL predicted number of nodes was validated by extensive testing of the model with respect to extrapolation and generalization capabilities using unseen examples from a test set.
The optimal granularity of the models was predicted for sensible values–only marginally different from the experimentally determined optimal ones. This shows that this rigorous and not ad hoc form of “Occam’s Razor” is quite succesfull. Comparison of the performance of the—admittedly limited— experiments on the robot arm problem with that of other principles, such as NIC and AIC, indicated that MDL’s performance was better or competitive .
A similar theory and practice validation in case of the Bayesian framework for model comparison was given by Mackay . This paper inspired us to use the robot arm problem in the MDL setting. We note that the Bayesian framework is genuinely different as is rigorously demonstrated in our companion paper . It is well known that prefix code length is equivalent to negative log probability through the Shannon-Fano code , and therefore with every such code there corresponds an equivalent probability. Thus, the MDL coding approach can in principle be translated back into a Bayesian approach where the model code gives the prior. The analysis we have given in shows that the data-to-model error may not correspond the the conditional data-to-model probability if the data are “atypical” for the contemplated prior. Moreover, to the authors coding of large data is more natural than reasoning about possibly nonexisting probabilities.
### 5.1 Directions for Future Work
In general the MDL method appears to be well suited for supervised learning of best model discretization parameters for classification problems in which error coding is straightforward. Applying the MDL method is simple, and it is computationally not expensive.
The central point is that using MDL the optimal granularity of the model parameters can be computed automatically rather than tuned manually. This approach constitutes a rational and feasibly computable approach for feature selection as opposed to customary rather ad hoc approaches. The purpose of presenting the theory outline and the example applications is to stimulate re-use in different areas of pattern recognition, classification, and image understanding (region segmentation, color clustering segmentation, and so on).
## Acknowledgement
We are grateful to Les Valiant for many discussions on machine learning and the suggestion for this research, to Guido te Brake and Joost Kok for executing the robot arm experiment, and to the referees for their insightful comments.
## Appendix A Appendix: Kolmogorov Complexity
The Kolmogorov complexity of a finite object $`x`$ is the length of the shortest effective binary description of $`x`$. We give a brief outline of definitions and properties. For more details see . Let $`x,y,z𝒩`$, where $`𝒩`$ denotes the natural numbers and we identify $`𝒩`$ and $`\{0,1\}^{}`$ according to the correspondence
$$(0,ϵ),(1,0),(2,1),(3,00),(4,01),\mathrm{}$$
Here $`ϵ`$ denotes the empty word ‘’ with no letters. The length $`l(x)`$ of $`x`$ is the number of bits in the binary string $`x`$. For example, $`l(010)=3`$ and $`l(ϵ)=0`$.
The emphasis is on binary sequences only for convenience; observations in any alphabet can be so encoded in a way that is ‘theory neutral’.
A binary string $`x`$ is a proper prefix of a binary string $`y`$ if we can write $`x=yz`$ for $`zϵ`$. A set $`\{x,y,\mathrm{}\}\{0,1\}^{}`$ is prefix-free if for any pair of distinct elements in the set neither is a proper prefix of the other. A prefix-free set is also called a prefix code. Each binary string $`x=x_1x_2\mathrm{}x_n`$ has a special type of prefix code, called a self-delimiting code,
$$\overline{x}=x_1x_1x_2x_2\mathrm{}x_n\neg x_n,$$
where $`\neg x_n=0`$ if $`x_n=1`$ and $`\neg x_n=1`$ otherwise. This code is self-delimiting because we can determine where the code word $`\overline{x}`$ ends by reading it from left to right without backing up. Using this code we define the standard self-delimiting code for $`x`$ to be $`x^{}=\overline{l(x)}x`$. It is easy to check that $`l(\overline{x})=2n`$ and $`l(x^{})=n+2\mathrm{log}n`$.
We develop the theory using Turing machines, but we can as well use the set of LISP programs or the set of FORTRAN programs.
Let $`T_1,T_2,\mathrm{}`$ be a standard enumeration of all Turing machines, and let $`\varphi _1,\varphi _2,\mathrm{}`$ be the enumeration of corresponding functions which are computed by the respective Turing machines. That is, $`T_i`$ computes $`\varphi _i`$. These functions are the partial recursive functions or computable functions. The Kolmogorov complexity $`C(x)`$ of $`x`$ is the length of the shortest binary program from which $`x`$ is computed. Formally, we define this as follows.
###### Definition 4
The Kolmogorov complexity of $`x`$ given $`y`$ (for free on a special input tape) is
$$C(x|y)=\underset{p,i}{\mathrm{min}}\{l(i^{}p):\varphi _i(p,y)=x,p\{0,1\}^{},i𝒩\}.$$
Define $`C(x)=C(x|ϵ)`$.
Though defined in terms of a particular machine model, the Kolmogorov complexity is machine-independent up to an additive constant and acquires an asymptotically universal and absolute character through Church’s thesis, from the ability of universal machines to simulate one another and execute any effective process. The Kolmogorov complexity of an object can be viewed as an absolute and objective quantification of the amount of information in it. This leads to a theory of absolute information contents of individual objects in contrast to classic information theory which deals with average information to communicate objects produced by a random source .
For technical reasons we also need a variant of complexity, so-called prefix kolmogorov complexity, which is associated with Turing machines for which the set of programs resulting in a halting computation is prefix free. We can realize this by equiping the Turing machine with a one-way input tape, a separate work tape, and a one-way output tape. Such Turing machines are called prefix machines since the halting programs for anyone of them form a prefix free set. Taking the universal prefix machine $`U`$ we can define the prefix complexity analogously with the plain Kolmogorov complexity. If $`x^{}`$ is the first shortest program for $`x`$ then the set $`\{x^{}:U(x^{})=x,x\{0,1\}^{}\}`$ is a prefix code. That is, each $`x^{}`$ is a code word for some $`x`$, and if $`x^{}`$ and $`y^{}`$ are code words for $`x`$ and $`y`$ with $`xy`$ then $`x^{}`$ is not a prefix of $`x`$.
Let $``$ be a standard invertible effective one-one encoding from $`𝒩\times 𝒩`$ to prefix-free recursive subset of $`𝒩`$. For example, we can set $`x,y=x^{}y^{}`$. We insist on prefix-freeness and recursiveness because we want a universal Turing machine to be able to read an image under $``$ from left to right and determine where it ends.
###### Definition 5
The prefix Kolmogorov complexity of $`x`$ given $`y`$ (for free) is
$$K(x|y)=\underset{p,i}{\mathrm{min}}\{l(p,i):\varphi _i(p,y)=x,p\{0,1\}^{},i𝒩\}.$$
Define $`K(x)=K(x|ϵ)`$.
The nice thing about $`K(x)`$ is that we can interpret $`2^{K(x)}`$ as a probability distribution since $`K(x)`$ is the length of a shortest prefix-free program for $`x`$. By the fundamental Kraft’s inequality, see for example , we know that if $`l_1,l_2,\mathrm{}`$ are the code-word lengths of a prefix code, then $`_x2^{l_x}1`$. This leads to the notion of algorithmic universal distribution—a rigorous form of Occam’s razor–below.
## Appendix B Appendix: Universal Distribution
A Turing machine $`T`$ computes a function on the natural numbers. However, we can also consider the computation of real valued functions. For this purpose we consider both the argument of $`\varphi `$ and the value of $`\varphi `$ as a pair of natural numbers according to the standard pairing function $``$. We define a function from $`𝒩`$ to the reals $``$ by a Turing machine $`T`$ computing a function $`\varphi `$ as follows. Interprete the computation $`\varphi (x,t)=p,q`$ to mean that the quotient $`p/q`$ is the rational valued $`t`$th approxmation of $`f(x)`$.
###### Definition 6
A function $`f:𝒩`$ is enumerable if there is a Turing machine $`T`$ computing a total function $`\varphi `$ such that $`\varphi (x,t+1)\varphi (x,t)`$ and $`lim_t\mathrm{}\varphi (x,t)=f(x)`$. This means that $`f`$ can be computably approximated from below. If $`f`$ can also be computably approximated from above then we call $`f`$ recursive.
A function $`P:𝒩[0,1]`$ is a probability distribution if $`_{x𝒩}P(x)1`$. (The inequality is a technical convenience. We can consider the surplus probability to be concentrated on the undefined element $`u𝒩`$).
Consider the family $`𝒫`$ of enumerable probability distributions on the sample space $`𝒩`$ (equivalently, $`\{0,1\}^{}`$). It is known, , that $`𝒫`$ contains an element m that multiplicatively dominates all elements of $`𝒫`$. That is, for each $`P𝒫`$ there is a constant $`c`$ such that $`c\text{m}(x)>P(x)`$ for all $`x𝒩`$. We call $`𝐦`$ an algorithmic universal distribution or shortly universal distribution.
The family $`𝒫`$ contains all distributions with computable parameters which have a name, or in which we could conceivably be interested, or which have ever been considered. The dominating property means that m assigns at least as much probability to each object as any other distribution in the family $`𝒫`$ does. In this sense it is a universal a priori by accounting for maximal ignorance. It turns out that if the true a priori distribution in Bayes’s rule is recursive, then using the single distribution m, or its continuous analogue the measure M on the sample space $`\{0,1\}^{\mathrm{}}`$ (for prediction as in ) is provably as good as using the true a priori distribution.
We also know, , that we can choose
$$\mathrm{log}\text{m}(x)=K(x)$$
(2)
That means that m assigns high probability to simple objects and low probability to complex or random objects. For example, for $`x=00\mathrm{}0`$ ($`n`$ 0’s) we have $`K(x)\stackrel{+}{=}K(n)\stackrel{+}{<}\mathrm{log}n+2\mathrm{log}\mathrm{log}n`$ since the program
$$\text{print }n\text{\_times a ‘‘0’’}$$
prints $`x`$. (The additional $`2\mathrm{log}\mathrm{log}n`$ term is the penalty term for a self-delimiting encoding.) Then, $`1/(n\mathrm{log}^2n)=O(\text{m}(x))`$. But if we flip a coin to obtain a string $`y`$ of $`n`$ bits, then with overwhelming probability $`K(y)\stackrel{+}{>}n`$ (because $`y`$ does not contain effective regularities which allow compression), and hence $`\text{m}(y)=O(1/2^n)`$.
The algorithmic universal distribution has many astonishing properties . One of these, of interest to the AI community, is that it gives a rigorous meaning to Occam’s Razor by assigning high probability to the “simple,” “regular,” objects and low probability to the “complex,” “irregular”, ones. For a popular account see . A celebrated result states that an object $`x`$ is individually random with respect to a conditional probability distribution $`P(|y)`$ iff $`\mathrm{log}(𝐦(x|y)/P(x|y))\stackrel{+}{=}0`$. Here the implied constant in the $`\stackrel{+}{=}`$ notation is in fact related to $`K(P(|y))`$—the length of the shortest program that computes the probability $`P(x|y)`$ on input $`x`$. In particular this means that for $`x`$ is “typical” or “in general position” with respect to conditional distribution $`P(|y)`$ iff the real probability $`P(x|y)`$ is close to the algorithmic universal probability $`𝐦(x|y)=2^{K(x|y)}`$.
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# Interference at quantum transitions: lasing without inversion and resonant four-wave mixing in strong fields at Doppler-broadened transitions
## I INTRODUCTION
Many concepts of quantum optics were originated proceeding from the assumed equality of probabilities of induced transitions accompanied by an absorption and emission of photons predicted by A. Einstein. Requirement of population inversion for lasing is direct consequence from this equality.
At the presence of several resonant electromagnetic fields, probability amplitudes of a coupled quantum states contain several oscillating components at close frequencies. Therefore alongside with squared modules of appropriate components cross terms indicating an interference of quantum transitions appear while calculating transitions probabilities. The coherent nonlinear optical phenomena stipulated by the indicated evolution of quantum states, driven by several fields, were called as nonlinear interference effects ($`NIE`$. In quantum optics $`NIE`$ may result in different coupling of a radiation with atoms in absorbing and emitting states controlled by the auxiliary fields . Various appearances of these effects are feasible. Soon after discovery of lasers Rautian and Sobelman showed feasibility of amplification without inversion ($`AWI`$) in two-level systems. The features of $`AWI`$ in optical three-level systems were explored in . Studies of $`NIE`$ in absorption/gain spectra including experiments on generation of an optical radiation in three-level systems, so that generation was possible only at the expense of nonlinear interference effects, drew much attention in 60th and 70th . (Review of relevant optical experiments and of earlier papers on $`AWI`$ in microwave range see in .) Coherent population trapping ($`CPT`$) is one of the appearances of $`NIE`$ for the states with negligible relaxation rates and Doppler effects. In 80th – 90th studies of coherent interference processes at quantum transitions have been attracting much interest again in the context of $`AWI`$, electromagnetically - induced transparency ($`EIT`$), $`CPT`$, enhanced nonlinear optical frequency conversion and other manifestations of these effects .
In classical terms an emission and absorption of a radiation are stipulated by forced oscillations of bound charges and depend on phase difference between radiation and induced oscillations. However, a radiation may simultaneously drive several coherent interfering oscillations of a various origin. Depending on a relation of their phases and amplitudes the interference can be either constructive or destructive, full or partial. Thus the matching components of an optical response can either amplify or suppress each other. On the other hand, the macroscopic response of a medium can be thought as result of quantum transitions, at which the photons can
11th International Vavilov Conference on Nonlinear Optics, 24-28 June 1997, Novosibirsk, Russia SPIE Vol. 3485
simultaneously contribute in several quantum pathways. By applying semi-classical approach a deep analogy with many well known effects of classical physics can be used to interpret and foresee the relevant quantum optics effects. Thus, leaving aside classifications of involved elementary quantum processes (introduced and valid for weak fields in the limits of perturbation theory), it is possible to predict and to explain wide range of optical processes, stipulated by quantum interference, some of them are quite unusual.
The objective of the paper is to consider various appearance of interference effects in resonance nonlinear - optical processes with the aid of the outlined approach in a context of some recent experiments. The amplitude and phase relations of interfering intra-atomic oscillations depend on configuration and on relaxation characteristics of the coupled transitions, on type of nonlinear-optical process, on intensities and frequency detunings of the radiations from resonances. Due to the difference in Doppler shifts the contributions to the macroscopic polarization of atoms at various velocities in gases may interfere in a different way too. The interference appears differently in an absorption, refraction and in different four-wave mixing ($`FWM`$) processes.
As an illustration of $`NIE`$ the following results will be presented:
1. The possibility of an amplification of a radiation without inversion of saturated populations on resonant transition is investigated. Influence of the growth of intensity of an amplified radiation on inversionless amplification in various open and closed transition configurations is analyzed. The conditions are formulated and with the concrete examples is shown, that by proper change of incoherent excitation rate of levels and of auxiliary radiation intensity the index of an inversionless amplification does not decrease with growth of intensity of an amplified radiation. The elements of the theory of such lasing without inversion are presented.
2. It is shown, that due to Maxwell velocity distribution of atoms and corresponding inhomogeneous broadening of the coupled transitions, incoherent excitation of the intermediate levels may drastically change both spectral properties and a magnitude (by orders of magnitude) of nonlinear susceptibilities for resonant $`FWM`$ processes. As the consequence, important power saturation effects appear. These features must be taken into account for explanation of the experiments and optimization of frequency-conversion. Resonant $`FWM`$ coupling of two strong and two weak radiations is considered. Formulas for both cases of coupling, one is relevant to coherent population trapping, another – when each level is coupled to only one driving field are derived. The outcomes are applied to numerical analysis and to discussion of recent experiments .
## II ABSORPTION AND REFRACTION INDICES FOR A STRONG RADIATION AT THE PRESENCE OF OTHER STRONG RADIATION, COUPLED TO AN ADJACENT TRANSITION
First, consider interaction of two strong laser fields at the three-level system. Possible configurations of such systems are shown in FIG.1: folded $`V`$ and $`\mathrm{\Lambda }`$, and cascade - $`H`$. In further we shall investigate spectral features of a gas material for a radiation $`E_4`$ at frequency $`\omega _4`$, tunable in the vicinity of a transition $`lm`$. It’s intensity is not supposed weak. Depending on the configuration of transitions under consideration one of the auxiliary strong fields $`E_1`$, $`E_3`$ or $`E_2`$ with frequencies $`\omega _1`$, $`\omega _3`$, $`\omega _2`$, resonant to adjacent transitions shown in the figure is turned on. All radiations are supposed to be uniformly polarized co- or counter-propagating travelling wave: $`E_j(z,t)=E_j\mathrm{exp}\{[i(\omega _jtk_jz)]\}+k.c.,`$ where $`k_j`$ \- can take both positive and negative values, $`j=1,2,3,4`$. Incoherent excitation of the levels with Maxwell’s velocity distribution, all possible population and coherence relaxation channels are accounted for.
It is necessary to distinguish the open and closed energy-level configurations. In open one (lowest level is not ground), the rate of incoherent excitation of the levels by an external source practically does not depend on the rate of induced transitions between considered levels. In the closed one (lowest level is ground one), the excitation rate for atoms at different levels and velocities depends on the value and velocity distribution of the other populations, which are dependent on the intensity of the driving fields.
### A General equations for absorption end refraction indices
Power dependent susceptibility $`\chi _4`$, responsible for absorption and refraction, can be found from the equation:
$$P^{NL}(\omega _4)=N\chi _4E_4,$$
(1)
where polarization $`P^{NL}(\omega _4)`$ is convenient to calculate with aid of density matrix $`\rho _{ij}`$:
$$P=N\rho _{ij}d_{ji}+c.c..$$
(2)
Taking into account above discussed relaxation and incoherent excitation processes, density matrix equations for a mixture of pure quantum mechanical ensembles in the interaction representation can be written in general form as:
$`L_{nn}\rho _{nn}=q_ni[V,\rho ]_{nn}+\gamma _{mn}\rho _{mm},L_{lm}\rho _{lm}=L_4\rho _4=i[V,\rho ]_{lm},`$ (3)
$`L_{ij}=d/dt+\mathrm{\Gamma }_{ij},V_{lm}=G_{lm}\mathrm{exp}\{i[\mathrm{\Omega }_4tkz]\},G_{lm}=𝐄_\mathrm{𝟒}𝐝_{\mathrm{𝐥𝐦}}/2\mathrm{},`$
where $`\mathrm{\Omega }_4=\omega _4\omega _{mn}`$ \- frequency detuning from resonance; $`\mathrm{\Gamma }_{mn}`$ \- homogeneous half-widths of transitions, in absence of collisions $`\mathrm{\Gamma }_{mn}=(\mathrm{\Gamma }_m+\mathrm{\Gamma }_n)/2`$; $`\mathrm{\Gamma }_n=_j\gamma _{nj}`$ \- inverse lifetimes of levels; $`\gamma _{mn}`$ \- rate of relaxation from the level $`m`$ to $`n`$, $`q_n=_jw_{nj}r_j`$ \- rate of incoherent excitation to a state $`n`$ from underlying levels. For open configurations $`q_i`$ \- is mainly determined by the population of the ground state and practically does not depend on the driving fields.
In a steady-state regime a set of density-matrix equation may be reduced to the set of algebraic equations . Below we present only results of calculations. Despite of essential distinctions in manifestations of $`NIE`$ in different open and closed configuration, formulas for absorption/gain ($`\alpha `$) and resonant part of refractive ($`\delta n`$) indices and also for power dependent populations of the levels can be presented uniformly for all configurations shown on FIG.1:
$`\alpha _4/\alpha _{04}=Re\{\chi _4/\chi _4^0\},\delta n_4/2\delta n_{04}=Im\{\chi _4/\chi _4^0\},\alpha _i/\alpha _{0i}=Re\{\chi _i/\chi _i^0\},\delta n_i/2\delta n_{0i}=Im\{\chi _i/\chi _i^0\},`$
$`{\displaystyle \frac{\chi _4}{\chi _4^0}}={\displaystyle \frac{\mathrm{\Gamma }_4}{P_4}}{\displaystyle \frac{\mathrm{\Delta }r_4(1+u_2)\mathrm{\Delta }r_ig_2}{\mathrm{\Delta }n_4(1+g_1+u_2)}},{\displaystyle \frac{\chi _i}{\chi _i^0}}={\displaystyle \frac{\mathrm{\Gamma }_i}{P_i}}{\displaystyle \frac{\mathrm{\Delta }r_i(1+g_1^{})\mathrm{\Delta }r_4u_1^{}}{\mathrm{\Delta }n_4(1+g_1^{}+u_2^{})}},`$ (4)
Here and further index $`i`$ specifies transition, resonant to the auxiliary radiation (see FIG.1), $`\chi `$ is susceptibility, $`\alpha _0,\delta n_0,\chi _0`$ are corresponding maximum resonant values at zero field intensities, $`P_j=\mathrm{\Gamma }_j+i\mathrm{\Omega }_j`$ (for example: $`P_{lm}=P_4=\mathrm{\Gamma }_4+i\mathrm{\Omega }_4`$, $`P_{lm}=P_{ml}^{}`$, $`P_{lf}=P_{42}=\mathrm{\Gamma }_{lf}+i(\mathrm{\Omega }_4+\mathrm{\Omega }_2)`$ etc.). If the atom moves with speed $`v`$, Doppler shift of resonances must be taken into account by substitution $`\mathrm{\Omega }_j`$ for $`\mathrm{\Omega }_j^{^{}}=\mathrm{\Omega }_jk_jv`$. In further strokes will be omitted, but it is supposed, that the Doppler shift in the formulas is taken into account. $`\mathrm{\Delta }r_4=r_lr_m`$ is power dependent population difference; $`\mathrm{\Delta }n_4=n_ln_m`$, $`n_i`$ \- population of the level in absence of driving fields, which is described by the formula: $`n_i=(q_i/\mathrm{\Gamma }_i)+(\gamma _{ki}/\mathrm{\Gamma }_i)(q_k/\mathrm{\Gamma }_k)`$.
$`g_1=|G_i|^2/P_4P_{4i},g_2=|G_i|^2/P_i^{}P_{4i},u_1=|G_4|^2/P_4P_{4i},u_2=|G_4|^2/P_i^{}P_{4i}.`$ (5)
$`G_j`$ are coupling Rabi frequencies. Formulas for populations differences can be presented uniformly too:
$`\mathrm{\Delta }r_4=(\mathrm{\Delta }n_4X_2\mathrm{\Delta }n_iX_3)/(X_1X_2X_3X_4),\mathrm{\Delta }r_i=(\mathrm{\Delta }n_iX_1\mathrm{\Delta }n_4X_4)/(X_1X_2X_3X_4).`$ (6)
$`X_2=1+Re\{a_{24}\text{æ}_4{\displaystyle \frac{\mathrm{\Gamma }_4}{P_4}}{\displaystyle \frac{g_2}{1+g_1+u_2}}+a_{2i}\text{æ}_i{\displaystyle \frac{\mathrm{\Gamma }_i}{P_i}}{\displaystyle \frac{1+g_1^{}}{1+g_1^{}+u_2^{}}}\},`$
$`X_3=Re\{a_{34}\text{æ}_4{\displaystyle \frac{\mathrm{\Gamma }_4}{P_4}}{\displaystyle \frac{g_2}{1+g_1+u_2}}+a_{3i}\text{æ}_i{\displaystyle \frac{\mathrm{\Gamma }_i}{P_i}}{\displaystyle \frac{1+g_1^{}}{1+g_1^{}+u_2^{}}}\},`$
$`X_1=1+Re\{a_{14}\text{æ}_4{\displaystyle \frac{\mathrm{\Gamma }_4}{P_4}}{\displaystyle \frac{1+u_2}{1+g_1+u_2}}+a_{1i}\text{æ}_i{\displaystyle \frac{\mathrm{\Gamma }_i}{P_i}}{\displaystyle \frac{u_1^{}}{1+g_1^{}+u_2^{}}}\},`$
$`X_4=Re\{a_{44}\text{æ}_4{\displaystyle \frac{\mathrm{\Gamma }_4}{P_4}}{\displaystyle \frac{1+u_2}{1+g_1+u_2}}+a_{4i}\text{æ}_i{\displaystyle \frac{\mathrm{\Gamma }_i}{P_i}}{\displaystyle \frac{u_1^{}}{1+g_1^{}+u_2^{}}}\},`$
$`a_{14}=a_{34},a_{1i}=a_{3i},a_{24}=a_{44},a_{2i}=a_{4i}.`$ (7)
Sign minus in (A),(6) concerns to folded $`V`$ ($`E_4,E_1`$) and $`\mathrm{\Lambda }`$ ($`E_4,E_3`$) schemes, plus - to cascade $`H`$ ($`E_4,E_2`$) scheme. Beside that in the ladder $`H`$-scheme, $`P_i`$ must be substituted for $`P_i^{}`$, and vice a versa: $`P_i^{}`$ for $`P_i`$. $`\text{æ}_4`$ and $`\text{æ}_i`$ \- are saturation parameters accordingly for transitions $`4`$ and $`i`$. For open configurations
$$\text{æ}_4=2|G_4|^2(\mathrm{\Gamma }_l+\mathrm{\Gamma }_m\gamma _4)/(\mathrm{\Gamma }_l\mathrm{\Gamma }_m\mathrm{\Gamma }_4),$$
(8)
whereas $`\text{æ}_i`$ and parameters $`a_{ij}`$ depending only on relaxation constants are defined below for each configuration.
1. $`V`$ \- scheme (fields $`E_4`$, $`E_1`$; $`i=1`$)
OPEN CONFIGURATION
$`\text{æ}_i=\text{æ}_1={\displaystyle \frac{2|G_1|^2(\mathrm{\Gamma }_g+\mathrm{\Gamma }_l\gamma _1)}{\mathrm{\Gamma }_g\mathrm{\Gamma }_l\mathrm{\Gamma }_1}},a_{2i}=a_{21}=1,a_{14}=1,a_{3i}=a_{31}={\displaystyle \frac{\mathrm{\Gamma }_g\gamma _1}{\mathrm{\Gamma }_g+\mathrm{\Gamma }_l\gamma _1}},a_{44}={\displaystyle \frac{\mathrm{\Gamma }_m\gamma _4}{\mathrm{\Gamma }_l+\mathrm{\Gamma }_m\gamma _4}}.`$ (9)
CLOSED CONFIGURATION
$`\text{æ}_4=4|G_4|^2/\mathrm{\Gamma }_m\mathrm{\Gamma }_4,\text{æ}_i=\text{æ}_1=4|G_1|^2/\mathrm{\Gamma }_g\mathrm{\Gamma }_1,`$
$`a_{3i}=a_{31}=0.5\mathrm{\Delta }n_4,a_{44}=0.5\mathrm{\Delta }n_1,a_{2i}=a_{21}=0.5[1+\mathrm{\Delta }n_1],a_{14}=0.5[1+\mathrm{\Delta }n_4].`$ (10)
2. $`\mathrm{\Lambda }`$ \- scheme (fields $`E_4`$, $`E_3`$, $`i=3`$)
OPEN CONFIGURATION
$`\text{æ}_3=2|G_3|^2(\mathrm{\Gamma }_m+\mathrm{\Gamma }_n\gamma _3)/\mathrm{\Gamma }_m\mathrm{\Gamma }_n\mathrm{\Gamma }_3,a_{2i}=a_{23}=1,a_{14}=1,`$
$`a_{3i}=a_{33}=\mathrm{\Gamma }_n(\mathrm{\Gamma }_l\gamma _4)/\mathrm{\Gamma }_l(\mathrm{\Gamma }_m+\mathrm{\Gamma }_n\gamma _3),a_{44}=\mathrm{\Gamma }_l(\mathrm{\Gamma }_n\gamma _3)/\mathrm{\Gamma }_n(\mathrm{\Gamma }_m+\mathrm{\Gamma }_l\gamma _4).`$ (11)
CLOSED CONFIGURATION
$`\text{æ}_4={\displaystyle \frac{4|G_4|^2}{\mathrm{\Gamma }_m\mathrm{\Gamma }_4}},a_{3i}=a_{33}=1+\mathrm{\Delta }n_4(1+2\mathrm{\Delta }n_4){\displaystyle \frac{\mathrm{\Gamma }_m\gamma _3}{\mathrm{\Gamma }_m+\mathrm{\Gamma }_n\gamma _3}},a_{44}=0.5[1{\displaystyle \frac{\gamma _3}{\mathrm{\Gamma }_n}}+\mathrm{\Delta }n_3(1+{\displaystyle \frac{\gamma _3}{\mathrm{\Gamma }_n}})],`$
$`a_{2i}=a_{23}=1+\mathrm{\Delta }n_3(\mathrm{\Gamma }_n\mathrm{\Gamma }_m+\gamma _3)/(\mathrm{\Gamma }_n+\mathrm{\Gamma }_m\gamma _3),a_{14}=0.5[1+\mathrm{\Delta }n_4(1+\gamma _3/\mathrm{\Gamma }_n)].`$ (12)
H-scheme (fields $`E_4`$, $`E_2`$; $`i=2`$)
OPEN CONFIGURATION
$`\text{æ}_2={\displaystyle \frac{2\left|G_2\right|^2\left(\mathrm{\Gamma }_f+\mathrm{\Gamma }_m\gamma _2\right)}{\mathrm{\Gamma }_f\mathrm{\Gamma }_m\mathrm{\Gamma }_2}};a_{14}=1;a_{2i}=a_{22}=1;a_{3i}=a_{32}={\displaystyle \frac{\mathrm{\Gamma }_l\gamma _4}{\mathrm{\Gamma }_l}}{\displaystyle \frac{\mathrm{\Gamma }_f\gamma _2}{\mathrm{\Gamma }_m+\mathrm{\Gamma }_f\gamma _2}};a_{44}={\displaystyle \frac{\mathrm{\Gamma }_l}{\mathrm{\Gamma }_l+\mathrm{\Gamma }_m\gamma _4}}.`$ (13)
CLOSED CONFIGURATION
$`\text{æ}_4=4|G_4|^2/\mathrm{\Gamma }_m\mathrm{\Gamma }_4,a_{3i}=a_{32}=(1+2\mathrm{\Delta }n_4)(\mathrm{\Gamma }_f\gamma _2)/(\mathrm{\Gamma }_m+\mathrm{\Gamma }_f\gamma _2)\mathrm{\Delta }n_4,a_{44}=0.5(1\mathrm{\Delta }n_2),`$
$`a_{14}=0.5(1+\mathrm{\Delta }n_4),a_{2i}=a_{22}=1+\mathrm{\Delta }n_2(\mathrm{\Gamma }_m\mathrm{\Gamma }_f+\gamma _2)/(\mathrm{\Gamma }_m+\mathrm{\Gamma }_f\gamma _2).`$ (14)
$`NIE`$ are associated with coherence at two-photon transitions and disappear at $`|P_{4i}|\mathrm{}`$. At $`G_4=0`$ formulas (A), (6) converge into those, similar to discussed and analyzed in . Following , range of parameters where amplification is not accompanied by the inversion of power-dependent populations can be easily found from (A), (6). For the resonant coupling in $`V`$ and $`\mathrm{\Lambda }`$ schemes conditions for AWI at the transitions 4 and $`i`$ take the form, correspondingly:
$$\mathrm{\Delta }r_4/\mathrm{\Delta }r_i<g_2/(1+u_2);\mathrm{\Delta }r_i/\mathrm{\Delta }r_4<u_1/(1+g_1).$$
(15)
Similar formulas for $`H`$ schemes show, that on the contrary to the previous configurations, inversion of populations on the adjacent transition is required for $`AWI`$ in center of the resonance, or amplification under certain conditions arises in wings of a resonance. (More detail formulas and analysis are given in .
Threshold and output power of lasing without inversion can be found from the equation:
$$\alpha _4=T,$$
(16)
Where $`T`$ is loss of a radiation from a laser cavity per one pass, scaled to the unit of length of the amplifying medium. Thus, the derived expressions determine conditions and characteristics of inversionless generation too.
### B Numerical analysis
Below we shall apply the derived expressions for numerical analysis of $`NIE`$ in open ($`Ne`$) and closed ($`Na`$) configurations of transitions. The same transitions of $`Ne`$ were considered in for illustration of possible AWI of weak probe field. The formulas for velocity averaged absorption index were derived and difference in $`NIE`$ for backward and forward waves in inhmogeneously broadened transitions was analyzed in for the cases of weak probe field and coupling Rabi frequency of driving field not exceeding Doppler width of the transition. Therefore in further main attention will be given to effects accompanying increase of intensity of an amplifying radiation.
FIG. 2a shows that inhomogeneous broadening does not destroy macroscopic coherence effects and $`NIE`$. The relative change of absorption index by an auxiliary field appears even larger, than for a homogeneously broadened transition. Coherent coupling gives rise to amplification at $`\omega _4`$ and establishes populations so that there is no inversion neither on one- nor on two-photon transitions FIG. 2b. With the growth of $`E_4`$ populations of $`m`$ and $`g`$ levels aim to equal magnitude, which indicates appearance of coherent population trapping. In the latter case a small modulation of velocity distribution on the level $`m`$ appears. Like in , analysis shows that in order to attain $`AWI`$ certain ratio between initial (in the absence of the radiations) populations must be fulfilled. FIG. 2c shows that absorption (gain) strongly depends on the intensities of both driving and probe fields. As it was outlined, the open and closed systems differ both in possible magnitudes of relaxation parameters and in dependence of incoherent excitation on an intensity of the driving fields. FIG. 4 considers the case, when 36% of a ground state atoms are initially excited by an incoherent pump to a level $`m`$, that still may correspond to strong absorption at the transition $`ml`$. By that strong driving field $`E_1`$ may produce $`AWI`$ for co-propagating shorter-wavelength weak radiation at $`\omega _4`$, which makes approximately 50% from an initial absorption (FIG. 4a). The amplification happens in absence of saturated population inversion for all transitions (FIG. 4b). It is essential that a population of a top level $`m`$ depends on the strength of $`E_1`$ even at zero intensities of a probe radiation $`E_4`$.
Curve 2 (FIG. 4a) shows, that the AWI strongly depend on intensity of an amplified radiation, that is accompanied for the given configuration by noticeable change of populations of levels $`m`$ and $`l`$. (FIG. 4b) displays energy and velocity distribution of the atoms corresponding to appearance of transparency at $`\omega _4`$. It is interesting to note that the distribution sharply varies with increase of intensity of $`E_1`$.
FIG. 4 shows, that $`AWI`$ may be maintained in a certain level with growth of intensity of an amplified radiation by changing incoherent excitation rate and strength of an auxiliary radiation.
## III INTERFERENCE EFFECTS IN RESONANT $`FWM`$ AT DOPPLER BROADENED TRANSITIONS
The use of resonant $`FWM`$ in gases for frequency conversion allows one to decrease the required power of fundamental radiations down to the magnitudes characteristic for cw lasers . $`FWM`$ concerns to so-called coherent nonlinear - optical processes, depending on phase-mismatch. As it was outline above, at resonant coupling, various coherent component stipulated by correlated quantum transitions and giving contribution to the process of radiation conversion can interfere. Constructive interference gives rise to enhancement of appropriate components of nonlinear polarization and destructive on the contrary - to elimination. Studies of appearances of quantum interference at $`FWM`$ continue to attract significant interest in the context of possible use for increase of conversion efficiency. The values of interfering components depend on energy level populations, which, in turn, depend on intensity of radiations and on processes of incoherent excitation and relaxation. Constructive or destructive character of an interference depends on a relation of phases coherent component and, therefore, on detunings from resonance and on type of nonlinear - optical process. In gaseous media inhomogeneous Doppler broadening of transitions is characteristic of typical experimental conditions. Depending on energy level, value and sign of Doppler shift, contributions of atoms, moving at various speeds, to macroscopic nonlinear polarization can both enhance and suppress each other. The above listed effects appear in a different way in absorption, refraction and in $`FWM`$ macroscopic polarization. It turns out that even small incoherent excitation of the levels at Doppler broadened transitions may drastically change spectral properties and by orders of magnitude value of the nonlinear susceptibility . The choice of optimal conditions of conversion is essentially determined not only by influence of interference processes on nonlinear susceptibilities, but also on indices of an absorption (amplification) and refraction for coupled waves propagating through a resonance medium. Velocity selective population transfer and other effects of resonant coupling with strong fields give rise to specific power saturation effects in $`FWM`$. In experimental features were observed, which did not find explanations in framework of before published lowest order perturbation theory.
This section is devoted to investigation of mutual influence of quantum interference, relaxation, incoherent pump of levels, Doppler broadening, effects of strong fields on $`FWM`$ nonlinear polarization, and absorption of coupled radiations in the context of experiments . We shall consider Raman like coupling (FIG.5) and $`FWM`$ processes of a type $`\omega _4=\omega _1\omega _2+\omega _3`$ as well as inverse process $`\omega _3=\omega _4\omega _1+\omega _2`$. That allows to compare influence of the above discussed elementary effects on conversion processes in different conditions. The elementary quantum mechanical processes, determining $`FWM`$, essentially depend on that whether the atom at an energy level couples with two strong fields or with one strong and one weak field.
In the range of negligibly small change of the strong radiations both due to absorption and conversion and assuming exact phase-matching, quantum efficiency of conversion $`E_3`$ in $`E_4`$ $`\eta _{q4}(\omega _2)=(k_3/k_4)|E_4(z)/E_3|^2`$ can be presented in the form:
$`\eta _{q4}=k_3k_42\pi N\chi _4^{(3)}E_1E_2^{}/(\mathrm{\Delta }\alpha /2)^2exp\{\alpha _4z\}(exp\{(\mathrm{\Delta }\alpha /2)z\}1)^2,`$ (17)
where $`N`$ – atomic number density, $`\chi _4^{(3)}`$ – nonlinear susceptibility, $`\mathrm{\Delta }\alpha =\alpha _1+\alpha _2+\alpha _3\alpha _4\alpha _3\alpha _4`$, $`\alpha _j`$ – absorption index for corresponding radiation. Quantum efficiency for conversion of $`E_4`$ in $`E_3`$ is written in symmetrical form by substitution of $`\chi _4^{(3)},\alpha _4`$ for $`\chi _3^{(3)},\alpha _3`$. For small length $`z<<min\{(\alpha _{4,3})^1,(\mathrm{\Delta }\alpha /2)^1\}`$ spectral features of conversion is determined only by nonlinear susceptibility and by intensities of the strong fields.
$$\eta _{q4,3}=k_3k_42\pi N\chi _{4,3}^{(3)}E_1E_2^2z^2$$
(18)
### A $`FWM`$ in two strong and one weak fields at the conditions of maximum coherence and coherent population trapping
In this section expressions for nonlinear polarization at frequencies $`\omega _3`$ and $`\omega _4`$ will be presented for cases, when the fields $`E_1`$ and $`E_2`$ are strong. Their frequencies $`\omega _1`$ and $`\omega _2`$ are close accordingly to the transition frequencies $`\omega _{lg}`$ and $`\omega _{ng}`$. Radiations $`E_3`$ and $`E_4`$ with frequencies $`\omega _3`$ and $`\omega _4`$, close to the transition frequencies $`\omega _{nm}`$ and $`\omega _{lm}`$ \- are supposed nonperturbatively weak. $`NIE`$ in two strong fields may give rise to such population transfer between levels $`l`$ and $`n`$ that population of the intermediate level $`g`$ change negligibly. Such behavior is similar to $`CPT`$ (for review see ). In more general sense a term $`CPT`$ is often applied to the processes whereas contribution of $`NIE`$ in population transfer is crucial. Such type of coupling will be considered in this section.
Nonlinear $`FWM`$ susceptibilities $`\stackrel{~}{\chi }_{3,4}^{(3)}`$ are determined by the equations:
$$P^{NL}(\omega _3)=N\stackrel{~}{\chi }_3^{(3)}E_1^{}E_2E_4,P^{NL}(\omega _4)=N\stackrel{~}{\chi }_4^{(3)}E_1E_2^{}E_3.$$
(19)
Expressions for $`\stackrel{~}{\chi }_{3,4}`$ as well as for absorption (refractive) susceptibilities $`\chi _i`$, derived with density matrix in similar way as that in Section II, are given by :
$`\stackrel{~}{\chi }_3={\displaystyle \frac{iK}{d_3(1+q_2)}}\left[{\displaystyle \frac{R_1^{}}{P_1^{}}}\left({\displaystyle \frac{1}{P_{41}}}+{\displaystyle \frac{1}{P_{12}^{}}}\right)+{\displaystyle \frac{R_2}{P_2P_{12}^{}}}+{\displaystyle \frac{R_4}{P_4P_{41}}}\right],`$
$`\stackrel{~}{\chi }_4={\displaystyle \frac{iK}{d_4(1+q_1)}}\left[{\displaystyle \frac{R_1}{P_1P_{12}}}+{\displaystyle \frac{R_2^{}}{P_2^{}}}\left({\displaystyle \frac{1}{P_{32}}}+{\displaystyle \frac{1}{P_{12}}}\right)+{\displaystyle \frac{R_3}{P_3P_{32}}}\right],`$ (20)
$`\chi _i/\chi _i^0=\mathrm{\Gamma }_iR_i/P_i\mathrm{\Delta }n_i,R_1=[(1+g_2^{})\mathrm{\Delta }r_1u_1^{}\mathrm{\Delta }r_2]/(1+g_2^{}+u_2),R_2=[(1+u_2^{})\mathrm{\Delta }r_2g_3\mathrm{\Delta }r_1]/(1+g_2+u_2^{}),`$
$`R_3={\displaystyle \frac{\mathrm{\Delta }r_3(1+q_1)u_3R_2^{}(1q_3)+q_1u_2R_1}{1+q_1+u_4}},R_4={\displaystyle \frac{\mathrm{\Delta }r_4(1+q_2)g_1R_1^{}(1q_4)+q_2g_2R_2}{1+q_2+g_4}}.`$
$`g_1={\displaystyle \frac{|G_1|^2}{P_{41}P_1^{}}},g_2={\displaystyle \frac{|G_1|^2}{P_{12}^{}P_2}},g_3={\displaystyle \frac{|G_1|^2}{P_{12}^{}P_1^{}}},g_4={\displaystyle \frac{|G_1|^2}{P_{41}P_4}},u_1={\displaystyle \frac{|G_2|^2}{P_{12}^{}P_2}},u_2={\displaystyle \frac{|G_2|^2}{P_{12}P_1}},u_3={\displaystyle \frac{|G_2|^2}{P_{32}P_2^{}}},u_4={\displaystyle \frac{|G_2|^2}{P_{32}P_3}},`$
$`q_1={\displaystyle \frac{|G_1|^2}{P_{32}d_4}},q_2={\displaystyle \frac{|G_2|^2}{P_{41}d_3}},q_3={\displaystyle \frac{|G_1|^2}{P_{12}d_4}},q_4={\displaystyle \frac{|G_2|^2}{P_{12}^{}d_3}},\mathrm{\Delta }r_1={\displaystyle \frac{\mathrm{\Delta }n_1X_2\mathrm{\Delta }n_2X_3}{X_1X_2X_3X_4}},\mathrm{\Delta }r_2={\displaystyle \frac{\mathrm{\Delta }n_2X_1\mathrm{\Delta }n_1X_4}{X_1X_2X_3X_4}},`$
$`d_3=\mathrm{\Gamma }_3+i(\mathrm{\Omega }_4\mathrm{\Omega }_1+\mathrm{\Omega }_2),d_4=\mathrm{\Gamma }_4+i(\mathrm{\Omega }_1\mathrm{\Omega }_2+\mathrm{\Omega }_3),`$
$`r_j=n_j+\mathrm{\Delta }r_2(b_1^j\text{æ}_1F_2+b_2^j\text{æ}_2F_3)\mathrm{\Delta }r_1(b_1^j\text{æ}_1F_1+b_2^j\text{æ}_2F_4),`$
$`X_1=1+a_1\text{æ}_1F_1a_2\text{æ}_2F_4,X_2=1+a_3\text{æ}_2F_3a_4\text{æ}_1F_2,`$
$`X_3=a_2\text{æ}_2F_3a_1\text{æ}_1F_2,X_4=a_4\text{æ}_1F_1a_3\text{æ}_2F_4.`$ (21)
$`F_1=Re\{{\displaystyle \frac{\mathrm{\Gamma }_1}{P_1}}{\displaystyle \frac{1+g_2^{}}{1+g_2^{}+u_2}}\},F_2=Re\{{\displaystyle \frac{\mathrm{\Gamma }_1}{P_1}}{\displaystyle \frac{u_1^{}}{1+g_2^{}+u_2}}\},F_3=Re\{{\displaystyle \frac{\mathrm{\Gamma }_2}{P_2^{}}}{\displaystyle \frac{1+u_2}{1+g_2^{}+u_2}}\},F_4=Re\{{\displaystyle \frac{\mathrm{\Gamma }_2}{P_2^{}}}{\displaystyle \frac{g_3^{}}{1+g_2^{}+u_2}}\}.`$
Here and further $`\mathrm{\Delta }r_1=r_lr_g,\mathrm{\Delta }r_2=r_nr_g,\mathrm{\Delta }r_3=r_nr_m,\mathrm{\Delta }r_4=r_lr_m,\mathrm{\Delta }n_4=n_ln_m`$, etc., $`r_j`$ and $`n_j`$ are population of levels accordingly power dependent and in absence of the fields; $`\text{æ}_1`$, $`\text{æ}_2`$ – saturation parameters; $`a_k`$, $`b_k^j`$ \- coefficients, determined by relaxation properties of the transitions, which are different for open and closed transition configurations and will be defined below; $`K`$ is constant.
OPEN CONFIGURATIONS
For the open system the parameters in (A)are:
$`\text{æ}_1={\displaystyle \frac{2|G_1|^2(\mathrm{\Gamma }_g+\mathrm{\Gamma }_l\gamma _1)}{\mathrm{\Gamma }_g\mathrm{\Gamma }_l\mathrm{\Gamma }_1}},\text{æ}_2={\displaystyle \frac{2|G_2|^2(\mathrm{\Gamma }_g+\mathrm{\Gamma }_n\gamma _2)}{\mathrm{\Gamma }_g\mathrm{\Gamma }_n\mathrm{\Gamma }_2}},`$
$`a_2={\displaystyle \frac{\mathrm{\Gamma }_n}{\mathrm{\Gamma }_l}}{\displaystyle \frac{\mathrm{\Gamma }_l\gamma _1}{\mathrm{\Gamma }_g+\mathrm{\Gamma }_n\gamma _2}},a_4={\displaystyle \frac{\mathrm{\Gamma }_l}{\mathrm{\Gamma }_n}}{\displaystyle \frac{\mathrm{\Gamma }_n\gamma _2}{\mathrm{\Gamma }_l+\mathrm{\Gamma }_g\gamma _1}},a_1=a_3=1,`$
$`b_1^g={\displaystyle \frac{\mathrm{\Gamma }_l}{\mathrm{\Gamma }_l+\mathrm{\Gamma }_g\gamma _1}},b_2^g={\displaystyle \frac{\mathrm{\Gamma }_n}{\mathrm{\Gamma }_g+\mathrm{\Gamma }_n\gamma _2}},b_1^n={\displaystyle \frac{\gamma _2}{\mathrm{\Gamma }_n}}b_1^g,`$
$`b_2^n={\displaystyle \frac{\mathrm{\Gamma }_g\gamma _2}{\mathrm{\Gamma }_g+\mathrm{\Gamma }_n\gamma _2}},b_1^l={\displaystyle \frac{\mathrm{\Gamma }_g\gamma _1}{\mathrm{\Gamma }_l+\mathrm{\Gamma }_g\gamma _1}},b_2^l={\displaystyle \frac{\gamma _1}{\mathrm{\Gamma }_l}}b_2^g,b_i^m=0.`$
CLOSED CONFIGURATION
The corresponding parameters in (A) take the values:
$`\text{æ}_1=4|G_1|^2/\mathrm{\Gamma }_g\mathrm{\Gamma }_1,a_1=0.5[1+\mathrm{\Delta }n_1(1+\gamma _2/\mathrm{\Gamma }_n)],a_2=1+\mathrm{\Delta }n_1(1+2\mathrm{\Delta }n_1)(\mathrm{\Gamma }_g\gamma _2)/(\mathrm{\Gamma }_g+\mathrm{\Gamma }_n\gamma _2),`$
$`a_3=1+\mathrm{\Delta }n_2[12(\mathrm{\Gamma }_g\gamma _2)/(\mathrm{\Gamma }_g+\mathrm{\Gamma }_n\gamma _2)],a_4=0.5[1(\gamma _2/\mathrm{\Gamma }_n)+\mathrm{\Delta }n_2[1+(\gamma _2/\mathrm{\Gamma }_n)],`$
$`b_1^l/n_l=b_1^m/n_m=(\mathrm{\Gamma }_n+\gamma _2)/2\mathrm{\Gamma }_n,b_2^l/n_l=b_2^l/n_l=(\mathrm{\Gamma }_n\mathrm{\Gamma }_g+\gamma _2)/(\mathrm{\Gamma }_g+\mathrm{\Gamma }_n\gamma _2),`$
$`b_1^n=0.5n_n(1n_n)\gamma _2/2\mathrm{\Gamma }_n,b_2^n=(2n_n1)(\mathrm{\Gamma }_g\gamma _2)/(\mathrm{\Gamma }_g+\mathrm{\Gamma }_n\gamma _2)n_n,`$
$`b_1^g=n_g(\mathrm{\Gamma }_n+\gamma _2)/2\mathrm{\Gamma }_n0.5,b_2^g=1n_g(12n_g)(\mathrm{\Gamma }_g\gamma _2)/(\mathrm{\Gamma }_g+\mathrm{\Gamma }_n\gamma _2),`$
$`n_l=(1+w_m/\mathrm{\Gamma }_m+w_g/\mathrm{\Gamma }_g+w_n^{}/\mathrm{\Gamma }_n)^1,n_m=w_mn_l/\mathrm{\Gamma }_m,n_g=w_gn_l/\mathrm{\Gamma }_g,`$
$`n_n=w_n^{}n_l/\mathrm{\Gamma }_n,w_n^{}=w_n+w_g\gamma _{gn}/\mathrm{\Gamma }_g+w_m\gamma _{mn}/\mathrm{\Gamma }_m`$
### B $`FWM`$ of two strong and two weak radiations under condition of perturbation of each resonant energy level only by one strong radiation
In the features coming out from increase of intensity of a radiation at frequency $`\omega _3`$ were investigated too. Consider a case, where the fields $`E_1`$ and $`E_3`$ are strong, and $`E_2`$ and $`E_4`$ \- weak. With the aid of solution of a set of equations for off- and diagonal elements of density matrix up to the first order of perturbation theory in respect of the weak fields equations for the susceptibilities can be presented as :
$`\stackrel{~}{\chi }_2={\displaystyle \frac{iK}{d_2(1+v_5^{}+g_5^{})}}\left[\left({\displaystyle \frac{\mathrm{\Delta }r_1}{P_1P_{41}^{}}}+{\displaystyle \frac{\mathrm{\Delta }r_3}{P_3P_{43}^{}}}\right)+{\displaystyle \frac{R_4^{}}{P_4^{}}}\left({\displaystyle \frac{1}{P_{41}^{}}}+{\displaystyle \frac{1}{P_{43}^{}}}\right)\right],v_5={\displaystyle \frac{|G_3|^2}{P_{41}d_2^{}}},g_5={\displaystyle \frac{|G_1|^2}{P_{43}d_2^{}}},v_1={\displaystyle \frac{|G_3|^2}{P_{43}P_3^{}}},`$
$`\stackrel{~}{\chi }_4={\displaystyle \frac{iK}{d_4(1+v_7^{}+g_7^{})}}\left[\left({\displaystyle \frac{\mathrm{\Delta }r_1}{P_1P_{12}}}+{\displaystyle \frac{\mathrm{\Delta }r_3}{P_3P_{32}}}\right)+{\displaystyle \frac{R_2^{}}{P_2^{}}}\left({\displaystyle \frac{1}{P_{12}}}+{\displaystyle \frac{1}{P_{32}}}\right)\right],v_7={\displaystyle \frac{|G_3|^2}{P_{12}^{}d_4^{}}},g_7={\displaystyle \frac{|G_1|^2}{P_{32}^{}d_4^{}}},v_2={\displaystyle \frac{|G_3|^2}{P_{32}^{}P_2}},`$ (22)
$`R_2={\displaystyle \frac{\mathrm{\Delta }r_2(1+g_7+v_7)v_3(1+v_7g_8)\mathrm{\Delta }r_3g_3(1+g_7v_8)\mathrm{\Delta }r_1}{(1+g_2+v_2)+[g_7+g_2(g_7v_8)+v_7+v_2(v_7g_8)]}},v_8={\displaystyle \frac{|G_3|^2}{P_{32}^{}d_4^{}}},g_8={\displaystyle \frac{|G_1|^2}{P_{12}^{}d_4^{}}},v_3={\displaystyle \frac{|G_3|^2}{P_{32}^{}P_3^{}}},`$
$`R_4={\displaystyle \frac{\mathrm{\Delta }r_4(1+v_5+g_5)g_1(1+g_5v_6)\mathrm{\Delta }r_1v_1(1+v_5g_6)\mathrm{\Delta }r_3}{(1+g_4+v_4)+[v_5+v_4(v_5g_6)+g_5+g_4(g_5v_6)]}},v_6={\displaystyle \frac{|G_3|^2}{P_{43}d_2^{}}},g_6={\displaystyle \frac{|G_1|^2}{P_{41}d_2^{}}},v_4={\displaystyle \frac{|G_3|^2}{P_{43}P_4}},`$
$`d_2=\mathrm{\Gamma }_{ng}+i(\mathrm{\Omega }_1+\mathrm{\Omega }_3\mathrm{\Omega }_4),\chi _i/\chi _i^0=\mathrm{\Gamma }_i\mathrm{\Delta }r_i/P_i\mathrm{\Delta }n_i,(i=1,3),\chi _i/\chi _i^0=\mathrm{\Gamma }_iR_i/P_i\mathrm{\Delta }n_i,(i=2,4).`$ (23)
The rest notations are the same as in Subsection III.A. Expressions for the populations are:
OPEN CONFIGURATION
$`\mathrm{\Delta }r_1=[(1+\text{æ}_3)\mathrm{\Delta }n_1+b_1\text{æ}_3\mathrm{\Delta }n_3]/[(1+\text{æ}_1)(1+\text{æ}_3)a_1\text{æ}_1b_1\text{æ}_3],`$
$`\mathrm{\Delta }r_3=[(1+\text{æ}_1)\mathrm{\Delta }n_3+a_1\text{æ}_1\mathrm{\Delta }n_1]/[(1+\text{æ}_1)(1+\text{æ}_3)a_1\text{æ}_1b_1\text{æ}_3],`$
$`\mathrm{\Delta }r_2=\mathrm{\Delta }n_2b_2\text{æ}_3\mathrm{\Delta }r_3a_2\text{æ}_1\mathrm{\Delta }r_1,\mathrm{\Delta }r_4=\mathrm{\Delta }n_4a_3\text{æ}_1\mathrm{\Delta }r_1b_3\text{æ}_3\mathrm{\Delta }r_3,`$
$`r_m=n_m+(1b_2)\text{æ}_3\mathrm{\Delta }r_3,r_g=n_g+(1a_3)\text{æ}_1\mathrm{\Delta }r_1,`$
$`r_n=n_nb_2\text{æ}_3\mathrm{\Delta }r_3+a_1\text{æ}_1\mathrm{\Delta }r_1,r_l=n_lb_1\text{æ}_3\mathrm{\Delta }r_3+a_3\text{æ}_1\mathrm{\Delta }r_1,`$
$`\text{æ}_1=\text{æ}_1^0{\displaystyle \frac{\mathrm{\Gamma }_{lg}^2}{|P_1|^2}},\text{æ}_1^0={\displaystyle \frac{2(\mathrm{\Gamma }_l+\mathrm{\Gamma }_g\gamma _{gl})}{\mathrm{\Gamma }_l\mathrm{\Gamma }_g\mathrm{\Gamma }_{lg}}}|G_1|^2,\text{æ}_3=\text{æ}_3^0{\displaystyle \frac{\mathrm{\Gamma }_{mn}^2}{|P_3|^2}},\text{æ}_3^0={\displaystyle \frac{2(\mathrm{\Gamma }_m+\mathrm{\Gamma }_n\gamma _{mn})}{\mathrm{\Gamma }_m\mathrm{\Gamma }_n\mathrm{\Gamma }_{mn}}}|G_3|^2,`$
$`a_1={\displaystyle \frac{\gamma _{gn}a_2}{\mathrm{\Gamma }_n\gamma _{gn}}}={\displaystyle \frac{\gamma _{gn}\mathrm{\Gamma }_la_3}{\mathrm{\Gamma }_n(\mathrm{\Gamma }_g\gamma _{gl})}}={\displaystyle \frac{\gamma _{gn}\mathrm{\Gamma }_l}{\mathrm{\Gamma }_n(\mathrm{\Gamma }_l+\mathrm{\Gamma }_g\gamma _{gl})}},`$
$`b_1={\displaystyle \frac{\gamma _{ml}\mathrm{\Gamma }_nb_2}{\mathrm{\Gamma }_l(\mathrm{\Gamma }_m\gamma _{mn})}}={\displaystyle \frac{\gamma _{ml}b_3}{\mathrm{\Gamma }_l(\mathrm{\Gamma }_l\gamma _{ml})}}={\displaystyle \frac{\gamma _{ml}\mathrm{\Gamma }_n}{\mathrm{\Gamma }_l(\mathrm{\Gamma }_m+\mathrm{\Gamma }_n\gamma _{mn})}}.`$
CLOSED CONFIGURATION
The populations of levels are described by the equations:
$`\mathrm{\Gamma }_mr_m=w_mr_l2\text{Re}\left\{iG_3^{}r_3\right\},\mathrm{\Gamma }_gr_g=w_gr_l2\text{Re}\left\{iG_1^{}r_1\right\},`$
$`\mathrm{\Gamma }_nr_n=w_nr_l+2\text{Re}\left\{iG_3^{}r_3\right\}+\gamma _{gn}r_g+\gamma _{mn}r_m,r_l=1r_mr_gr_n,`$
where $`r_1=iG_1\mathrm{\Delta }r_1/P_1`$, $`r_3=iG_3\mathrm{\Delta }r_3/P_3`$. The solution is
$`r_l=n_l(1+\text{æ}_3)(1+\text{æ}_1)/\beta ,r_g=(1+\text{æ}_3)[n_l(1+\text{æ}_1)\mathrm{\Delta }n_1]/\beta ,`$
$`r_n=\left\{n_m(1+\text{æ}_3)(1+\text{æ}_1)+[\mathrm{\Delta }n_3(1+\text{æ}_1)+\mathrm{\Delta }n_1\gamma _2\text{æ}_1/\mathrm{\Gamma }_n](1+b\text{æ}_3)\right\}/\beta `$
$`r_m=\left\{n_m(1+\text{æ}_3)(1+\text{æ}_1)+[\mathrm{\Delta }n_3(1+\text{æ}_1)+\mathrm{\Delta }n_1\gamma _2\text{æ}_1/\mathrm{\Gamma }_n]b\text{æ}_3\right\}/\beta ,`$
$`\mathrm{\Delta }r_1=r_lr_g=\mathrm{\Delta }n_1(1+\text{æ}_3)/\beta ,\mathrm{\Delta }r_3=r_nr_m=[\mathrm{\Delta }n_3(1+\text{æ}_1)+\mathrm{\Delta }n_1\gamma _2\text{æ}_1/\mathrm{\Gamma }_n]/\beta ,`$
where $`\beta =(1+\text{æ}_3)[1\mathrm{\Delta }n_3+2(n_l+n_m)\text{æ}_1]+(1+2b\text{æ}_3)[\mathrm{\Delta }n_3(1+\text{æ}_1)+\mathrm{\Delta }n_1\gamma _2\text{æ}_1/\mathrm{\Gamma }_n]`$, $`\mathrm{\Delta }n_1=n_ln_g`$, $`\mathrm{\Delta }n_3=n_nn_m`$, $`n_m=n_lw_m/\mathrm{\Gamma }_m`$, $`n_g=n_lw_g/\mathrm{\Gamma }_g`$, $`n_n=n_lw_{n}^{}{}_{}{}^{}/\mathrm{\Gamma }_n`$, $`n_l=(1+w_m/\mathrm{\Gamma }_m+w_g/\mathrm{\Gamma }_g+w_{n}^{}{}_{}{}^{}/\mathrm{\Gamma }_n)^1`$, $`w_{n}^{}{}_{}{}^{}=w_n+w_g\gamma _{gn}/\mathrm{\Gamma }_n+w_m\gamma _{mn}/\mathrm{\Gamma }_n`$, $`b=\mathrm{\Gamma }_n/(\mathrm{\Gamma }_m+\mathrm{\Gamma }_n\gamma _3)`$, $`\text{æ}_1=(2|G_1|^2/\mathrm{\Gamma }_1\mathrm{\Gamma }_g)(\mathrm{\Gamma }_1^2/|P_1|^2)`$, $`\text{æ}_3=(2|G_3|^2(\mathrm{\Gamma }_m+\mathrm{\Gamma }_n\gamma _3)/\mathrm{\Gamma }_m\mathrm{\Gamma }_n\mathrm{\Gamma }_3)(\mathrm{\Gamma }_3^2/|P_3|^2)`$. The remaining denotations are former.
### C Effect of Doppler broadening on resonant $`FWM`$
Formula for $`\stackrel{~}{\chi }_4^{(3)}`$ in lowest order of perturbation theory can be derived from (A), (22) at $`G_i0`$:
$`\stackrel{~}{\chi }_4^{(3)}(\omega _4=\omega _1\omega _2+\omega _3)={\displaystyle \frac{iK}{\mathrm{\Gamma }_{ml}+i(\mathrm{\Omega }_1^{^{}}\mathrm{\Omega }_2^{^{}}+\mathrm{\Omega }_3^{^{}})}}`$
$`\left\{{\displaystyle \frac{1}{\mathrm{\Gamma }_{gm}+i(\mathrm{\Omega }_3^{^{}}\mathrm{\Omega }_2^{^{}})}}[{\displaystyle \frac{n_gn_n}{\mathrm{\Gamma }_{ng}i\mathrm{\Omega }_2^{^{}}}}+{\displaystyle \frac{n_mn_n}{\mathrm{\Gamma }_{mn}+i\mathrm{\Omega }_3^{^{}}}}]+{\displaystyle \frac{1}{\mathrm{\Gamma }_{ln}+i(\mathrm{\Omega }_1^{^{}}\mathrm{\Omega }_2^{^{}})}}[{\displaystyle \frac{n_gn_n}{\mathrm{\Gamma }_{ng}i\mathrm{\Omega }_2^{^{}}}}+{\displaystyle \frac{n_gn_l}{\mathrm{\Gamma }_{lg}+i\mathrm{\Omega }_1^{^{}}}}]\right\},`$ (24)
where $`\mathrm{\Omega }_j^{^{}}=\mathrm{\Omega }_j𝐤_j𝐯`$, $`n_i=N_iexp\{(𝐯/\overline{v})^2\}/\sqrt{\pi }\overline{v}`$. As the function of $`v`$ all terms in (24), besides those proportional to $`n_gn_n`$, have all poles in one and the same complex half plane. Therefore at $`\mathrm{\Gamma }_i<<k_i\overline{v}`$ only terms, proportional to $`n_gn_n`$, do not vanish after averaging over Maxwell’s velocity distribution. Velocity averaged susceptibility is:
$`<\stackrel{~}{\chi }_4^{(3)}(\omega _4)>_v={\displaystyle \frac{iK\pi ^{1/2}\mathrm{exp}\{(\mathrm{\Omega }_2/k_2\overline{v})^2\}(N_gN_n)}{k_2\overline{v}[\stackrel{~}{\mathrm{\Gamma }}_1+i(\mathrm{\Omega }_1k_1\mathrm{\Omega }_2/k_2)][\stackrel{~}{\mathrm{\Gamma }}_3+i(\mathrm{\Omega }_3k_3\mathrm{\Omega }_2/k_2)]}},`$ (25)
$`\stackrel{~}{\mathrm{\Gamma }}_1=\mathrm{\Gamma }_{nl}+(k_1/k_21)\mathrm{\Gamma }_{ng},\stackrel{~}{\mathrm{\Gamma }}_3=\mathrm{\Gamma }_{gm}+(k_3/k_21)\mathrm{\Gamma }_{ng}.`$
As it is seen from (25), for the process $`\omega _4=\omega _1\omega _2+\omega _3`$ interference of contributions of atoms at different velocities to the velocity averaged nonperturbed $`FWM`$ nonlinear susceptibility $`<\stackrel{~}{\chi }_4>_v`$ leads to the fact, that in the lowest order on the small parameter $`\mathrm{\Gamma }_2/k_2\overline{v}`$, it is proportional to the velocity integrated difference between populations of the excited states $`N_gN_n`$. In the similar way, with aid of (22) one can find, that on the contrary, $`<\stackrel{~}{\chi }_2>_v`$ for the process $`\omega _2=\omega _1\omega _4+\omega _3`$ in the same approximation is determined by the population difference on transitions from the lowest level. For the resonant sum frequency $`FWM`$ $`\omega _4=\omega _1+\omega _2+\omega _3`$ in the cascade configuration of levels velocity averaged susceptibility occurs proportional to higher order of the small parameter $`\mathrm{\Gamma }/k\overline{v}`$ compared to Raman-type difference-frequency coupling . These features demonstrate great difference between resonant $`FWM`$ processes in homogeneously and inhomogeneously broadened transitions.
In strong electromagnetic fields above mentioned processes are accompanied by the velocity selective population transfer and by some other intensity dependent effects. In experiments on resonant $`cw`$ $`FWM`$ at Raman-like electronic molecular transitions of $`Na_2`$ have been carried out. Frequency tunable radiation at $`\omega _2`$ was generated at the same transition of $`Na_2`$ either in external dimer Raman laser or $`FWM`$ was performed inside the Raman laser cavity. Frequency $`\omega _2`$ was tuned by tuning $`\omega _1`$. Radiation at $`\omega _3`$ was provided from $`cw`$ dye laser. High conversion efficiency have been attained in single frequency nearly power saturation regime. Observed $`FWM`$ frequency tuning characteristics occured in disagreement with the predictions of lowest order perturbative theory. From that the authors derived the questions to be answered with the aid of an advanced nonperturbative theory. We shall use above presented expressions for the numerical analysis of the models with the parameters, close to that in the experiments, in order to explain main observed features.
The electronic - vibration-rotation transitions between $`X`$, $`A`$ and $`B`$ electronic levels of the dimer were used in the experiments, the lowest electronic level being ground one. Two $`FWM`$ processes were investigated: when frequency $`\omega _3`$ less than $`\omega _2`$ and, therefore frequency of a generated radiation $`\omega _4`$ was less than $`\omega _1`$ (down conversion) and opposite upconversion process. As it was discussed above, different appearances of interference processes at Doppler broadened transitions can be expected in those cases. Main observed experimental dependencies, which did not find explanations, can be summarized as follows.
According to (25), at $`\mathrm{\Omega }_1k_1\mathrm{\Omega }_2/k_2=0`$, in the lowest order of perturbation theory the maximum output of $`FWM`$ at $`\omega _4`$ as a function of $`\omega _3`$ corresponds to $`\mathrm{\Omega }_3=k_3\mathrm{\Omega }_2/k_2=0`$. Lineshape of the resonance is Lorentzian with the linewidth of the order of characteristic homogeneous widths of optical transitions. However, in the down conversion experiments the wide resonance of the order of Doppler width of transition $`ml`$ with the center being locked at $`\omega _3\omega _{mg}`$ was observed. It’s position practically did not vary at tuning $`\omega _2`$ within Doppler resonance of the transition $`gn`$ (at the expense of tuning of $`\omega _1`$, so that $`\mathrm{\Omega }_2=k_2\mathrm{\Omega }_1/k_1`$). In the upconversion experiments the resonance was tunable by tuning frequency $`\omega _2`$, but with the slope less than $`d\mathrm{\Omega }_3/d\mathrm{\Omega }_2=k_3/k_2`$. Width of the resonance also was commensurable with the Doppler width of the transition $`ml`$.
For numerical analysis we have used a model with the transitions parameters, close to those from the experiment.
1. Down conversion: $`\lambda _{ml}`$ = 598 nm, $`\lambda _{gl}`$ = 488 nm, $`\lambda _{mn}`$ = 655 nm, $`\lambda _{gn}`$ = 525 nm; $`k_1\overline{v}=6.94`$, $`k_2\overline{v}=6.45`$, $`k_3\overline{v}=5.17`$, $`k_4u=5.66`$, (in terms of $`10^9`$ $`c^1`$); $`\mathrm{\Gamma }_m=200`$, $`\mathrm{\Gamma }_n=30`$, $`\mathrm{\Gamma }_g=260`$, $`\gamma _{mn}=2`$, $`\gamma _{ml}=4`$, $`\gamma _{gn}=20`$, $`\gamma _{gl}=10`$, $`\mathrm{\Gamma }_{ln}=40`$, $`\mathrm{\Gamma }_{nm}=110`$, $`\mathrm{\Gamma }_{lm}=110`$, $`\mathrm{\Gamma }_{gm}=130`$, $`\mathrm{\Gamma }_{ng}=140`$, $`\mathrm{\Gamma }_{lg}=140`$, (in terms of $`10^6`$ $`c^1`$) $`N_l/N_n=30/2`$.
2. Up-conversion: $`\lambda _{ml}`$ = 473 nm, $`\lambda _{gl}`$ = 661 nm, $`\lambda _{mn}`$ = 514 nm, $`\lambda _{gn}`$ = 746 nm; $`k_1\overline{v}=5.12`$, $`k_2\overline{v}=4.54`$, $`k_3\overline{v}=6.59`$, $`k_4\overline{v}=7.16`$, (in terms of $`10^9`$ $`c^1`$); $`\mathrm{\Gamma }_m=200`$, $`\mathrm{\Gamma }_n=30`$, $`\mathrm{\Gamma }_g=260`$, $`\gamma _{mn}=2`$, $`\gamma _{ml}=4`$, $`\gamma _{gn}=20`$, $`\gamma _{gl}=10`$, $`\mathrm{\Gamma }_{ln}=40`$, $`\mathrm{\Gamma }_{nm}=110`$, $`\mathrm{\Gamma }_{lm}=110`$, $`\mathrm{\Gamma }_{gm}=130`$, $`\mathrm{\Gamma }_{ng}=140`$, $`\mathrm{\Gamma }_{lg}=140`$, (in terms of $`10^6`$ $`c^1`$), $`N_l/N_n=110/3.2`$, population of two upper levels being negligibly small.
First, with an aid of these models we shall illustrate a role of an interference at velocity averaging of nonlinear susceptibilities in weak fields. For down conversion in exact one- and multiphoton resonances and homogeneously broadened transitions computing gives the ratio of squared modulus of nonlinear susceptibilities $`|\chi _3^{(3)}/\chi _4^{(3)}|^2`$ = 2.5. For averaged values it yields $`|<\chi _3^{(3)}>_v/<\chi _4^{(3)}>_v|^2`$ = $`2.3110^2`$. If to change the population ratio for the inverse magnitude ($`N_n/N_l=30/2`$), we obtain: $`|\chi _3^{(3)}/\chi _4^{(3)}|^2`$ = 0.4. The difference between averaged values sharply decreases. Their ratio in this case yields: $`|<\chi _3^{(3)}>_v/<\chi _4^{(3)}>_v|^2`$ = 0.13.
For up-conversion similar computations give: $`|\chi _3^{(3)}/\chi _4^{(3)}|^2`$ = 2.7, $`|<\chi _3^{(3)}>_v/<\chi _4^{(3)}>_v|^2`$ = $`1.4510^3`$. At the inverse population ratio ($`N_n/N_l=110/3.2`$) we obtain: $`|\chi _3^{(3)}/\chi _4^{(3)}|^2`$ = 0.37. The difference between averaged values sharply decreases. Their ratio in this case is: $`|<\chi _3^{(3)}>_v/<\chi _4^{(3)}>_v|^2`$ = 0.24.
Thus, effect of inhomogeneous broadening of the resonant transitions on $`FWM`$ processes may be very strongly dependent on a specific process, as well as on distribution of the populations over levels and velocities. Therefore one can expect that velocity selective population transfer and other effects of strong fields may change conclusions of the lowest order perturbative theory.
For small number density and medium length $`FWM`$ conversion efficiency of weak radiation is proportional to a product of intensities of the strong radiations and squared modulus of velocity averaged nonlinear susceptibility (equation (18)). The later is intensity dependent too. Further we shall numerically analyze effects of the strong fields on conversion efficiency with an aid of the expressions of the subsections III.A and III.B. Intensities of the radiations will be characterized by the parameters $`S_1=|G_1|^2/\mathrm{\Gamma }_{gl}\mathrm{\Gamma }_{nl}`$ and $`S_2=|G_2|^2/\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{nl}`$, which are chosen in near saturation range like in the experiments.
Figures 6 show, that for the chosen parameters due to power broadening the resonance is much broader than homogeneous transition width and is commensurable with Doppler linewidth, which is of the order 60 in the used scale. For the parameters, corresponding to the up-conversion experiments (FIG. 6 a,b) in the range of small detunings $`\mathrm{\Omega }_1`$ ($`\mathrm{\Omega }_2=(k_2/k_1)\mathrm{\Omega }_1`$) (FIG. 6 a) the peak of the tuning curve is displaced very insignificantly (and even in the opposite side, depending on the value of $`\mathrm{\Omega }_1`$). At further increase of $`\mathrm{\Omega }_1`$ (FIG. 6 b) the maximum shifts with the increase of $`\mathrm{\Omega }_1`$, so that the slope $`\mathrm{\Omega }_3/\mathrm{\Omega }_1`$ is variable. A maximum of the slope corresponds to the detunings $`\mathrm{\Omega }_1`$ of about a half of Doppler width of the transition $`gl`$. Thus for the considered intensities the value of the slope reaches $`0.8`$, that makes $`0.5(k_3/k_1)`$.
FIG. 6 c is computed and drown for the parameters, corresponding to down conversion experiments. For the considered intensities the peak occurs locked to the center of transition $`ml`$ practically in all an interval of $`\mathrm{\Omega }_1`$ within the Doppler width of transition $`gl`$.
When the weak field detunings are fixed and driving fields frequencies $`\omega _1`$ and $`\omega _2`$ ($`\mathrm{\Omega }_2=(k_2/k_1)\mathrm{\Omega }_1`$) are tuned to the maximum, computer analysis of the slope $`d\mathrm{\Omega }_1/d\mathrm{\Omega }_3`$ in the range of $`S_1=65,S_2=2.33S_1`$ shows behavior similar to that observed in experiments ($`d\mathrm{\Omega }_1/d\mathrm{\Omega }_3k_1/2k_3`$). (Ratio $`S_1/S_2=2.33`$ is chosen according to the ratio of prodacts of experimental field intensities and Franck-Kondon factors.)
In the experiments the features, following increase of intensity of $`E_3`$, were observed too. In order to consider effects of this field and to understand whether $`CPT`$ play a decisive role in observed dependencies we have carried out numerical analysis of the up-conversion model with aid of formulas from the Subsection III.B (FIG. 7).
FIG. 7 a shows that at certain ratio of intensities even power narrowing of Doppler broadened $`FWM`$ resonance may happen. FIG. 7 b displays approximately constant slope $`d\mathrm{\Omega }_3/d\mathrm{\Omega }_1`$ ($`\mathrm{\Omega }_3`$ corresponds to the maximum output), which is about 0.75, in a quite wide interval of $`\mathrm{\Omega }_1`$. Only for $`\mathrm{\Omega }_1`$, larger than Doppler width, the slope starts decreasing. At larger intensities the slope varies more considerable, when $`\mathrm{\Omega }_1`$ is tuned within the Doppler line. So for FIG. 7 c the slope makes up 0.1 in the vicinity of $`y_1=20`$; 0.8 - in the range of $`y_1=40`$ and 0.6 - at $`y_1=60`$.
In an absorbing medium spectral properties of absorption indices may bring important effect on $`FWM`$. FIG. 8 a-c display corresponding lineshapes. FIG. 9 is computed with aid of formula (17) and shows that additional broadening of the tuning curve of $`FWM`$ output may arise from the propagation effect in an absorbing medium ($`Z=\alpha _{10}z`$, $`\mathrm{\Omega }_1=0,S_1=150,S_2=350`$ (upconversion)).
In conclusion, the theory of nonlinear interference processes at Doppler broadened quantum transitions in two strong resonant optical fields is developed. The derived formulas allow one account for such contributing processes relevant to experiments as various relaxation channels, incoherent excitation by an external source, population transfer and other coherent and incoherent effects accompanying coupling with strong optical fields in various open and closed configuration of transitions. Explicit formulae accounting for the effects of the strong fields are derived. Such appearance of quantum interference as amplification (and lasing) without inversion of power saturated populations and specific effects in resonant four wave mixing in gases are analyzed with the aid of numerical models. Parameters of the model are close to those of some recently carried out experiments. Crucial effect of Doppler-broadening on the contributions of the populations of different levels to four-wave mixing, which determine selection of optimal energy-level configuration and conditions for the experiments, is shown. Unexpected dependencies, observed in the experiments are explained.
ACKNOWLEDGMENTS
The author would like to thank S. A. Myslivets and V. M. Kuchin for assistance in calculations, A. A. Apolonskii, S. A. Babin, U. Hinze and B. Wellegehausen for useful discussion of their experiments prior publication. This work was supported by the Russian Foundation for Basic Research and by the Deutsche Forschungsgemeinschaft through collaborative Grant 96-02-00016 G. The author wishes to express his sincere thanks to L. J. F. Hermans from Hygens Laboratorium and the Netherlands Science Foundation (NWO) for support of this work.
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# Spin 1/2 Field Theory in the de Sitter space-time
## 1 Introduction
In general on curved space-time, no true spectral condition can be satisfied by Quantum Field Theory (QFT) and no unique vacum state exists. In the case of de Sitter spase-time, it has been discovered that the Hadamard condition selects an unique vacume state . The Hadamard condition is related to normal analyticity, that is to say, the two-point function is the boundary value of an analytic function. One can replace the usual spectral condition by a certain geodesic spectral condition (or KMS condition), and one can consider the generalized free fields on dS space-time. The generalized free fields can be defined entirely in terms of Wightman two-point function. In this context we present local free spinor fields (s=$`\frac{1}{2}`$) in 4-dimensional dS space-time based on analyticity in the complexified pseudo-Riemanian manifold. First we derive the dS-Dirac field equation as an eigenvalue equation for the Casimir operator and we find the solutions in terms of coordinate-independent dS plane-waves in tube domains. We define the two-point function $`W(z_1,z_2)`$ in terms of spinor dS plane-waves in their tube domains. Normal analyticity allows one to define Wigthman two-point function $`𝒲(x,y)`$ as the boundary value of $`W(z_1,z_2)`$ from the tube domains. Then the Hilbert space structure and the field operators $`\psi (f)`$ are defined. Finally, the unsmeared field operator $`\psi (x)`$ in terms of a coordinate-independent dS plane waves is also defined. This work is in the continuation of the previous ones concerning the scalar case .
## 2 Notation
De Sitter space-time is visualized as the hyperboloid with equation:
$$X_R=\{x^\alpha \mathrm{IR}^5:x.x=\eta _{\alpha \beta }x^\alpha x^\beta =(x^0)^2(x^1)^2(x^2)^2(x^3)^2(x^4)^2=R^2\}$$
$$\eta ^{\alpha \beta }=\text{diag}(1,1,1,1,1);\alpha ,\beta =0,1,\mathrm{},4.$$
(1)
Let us define the punctured half-con with: $`C=\{\xi ^\alpha \mathrm{IR}^5:\eta _{\alpha \beta }\xi ^\alpha \xi ^\beta =0\}`$. The kinematical group of the de Sitter space-time is $`G_R=SO_0(1,4)`$ and the double (and universal) covering group of $`G_R`$ is given in a quaternionic realisation by
$$Sp(2,2)=\{g=\left(\begin{array}{cccc}a& b\hfill & & \\ c& d\hfill & & \end{array}\right),\text{det}g=1,\gamma ^0\stackrel{~}{g}^t\gamma ^0=g^1;a,b,c,d\text{IHI}\},$$
(2)
$`g^t`$ denotes the $`2\times 2`$ transpose of $`g`$, $`\stackrel{~}{g}`$ denotes its quaternionic conjugate, and det$`g`$ is the determinant of $`g`$ viewed as a $`4\times 4`$ matrix with complex entries. Now we need five $`\gamma `$ matrices instead of the usual four ones in Minkowski space-time. They are defined by the Clifford algebra:
$$\{\gamma ^\alpha ,\gamma ^\beta \}=\gamma ^\alpha \gamma ^\beta +\gamma ^\beta \gamma ^\alpha =2\eta ^{\alpha \beta },\gamma ^\alpha =\gamma ^0\gamma ^\alpha \gamma ^0.$$
The quaternion representation of the $`\gamma `$ matrices is given by
$$\gamma ^0=\left(\begin{array}{cccc}\text{I1}& 0\hfill & & \\ 0& \text{I1}\hfill & & \end{array}\right),\gamma ^4=\left(\begin{array}{cccc}0& \text{I1}\hfill & & \\ \text{I1}& 0\hfill & & \end{array}\right)$$
$$\gamma ^1=\left(\begin{array}{cccc}0& i\sigma ^1\hfill & & \\ i\sigma ^1& 0\hfill & & \end{array}\right),\gamma ^2=\left(\begin{array}{cccc}0& i\sigma ^2\hfill & & \\ i\sigma ^2& 0\hfill & & \end{array}\right),\gamma ^3=\left(\begin{array}{cccc}0& i\sigma ^3\hfill & & \\ i\sigma ^3& 0\hfill & & \end{array}\right)$$
(3)
in terms of the $`2\times 2`$ unit I1 and Pauli matrices $`\sigma ^i`$. Casimir operator and infinitesimal generators read
$$Q=\frac{1}{2}L_{\alpha \beta }L^{\alpha \beta },L_{\alpha \beta }=M_{\alpha \beta }+S_{\alpha \beta }=i(x_\alpha \overline{}_\beta x_\beta \overline{}_\alpha )\frac{i}{4}[\gamma _\alpha ,\gamma _\beta ],$$
(4)
where $`\overline{}`$ is the tangential derivative: $`\overline{}_\beta =_\beta +H^2x_\beta x.`$.
## 3 dS-Dirac field equation and plane-waves solution
Starting from the Casimir operator and using the infinitesimal generators and the Casimir eigenvalue equation $`Q\psi (x)=(\nu ^2+\frac{3}{2})\psi (x)`$ give
$$Q\psi (x)=\{(\frac{1}{2}\gamma _\alpha \gamma _\beta M_{\alpha \beta }+2i)^2+\frac{3}{2}\}\psi (x)=(\nu ^2+\frac{3}{2})\psi (x),$$
(5)
where $`\psi (x)`$ is the 4-component spinor wave function and $`\nu \mathrm{IR}`$. A possible spinorial solution $`\psi (x)`$ to the above equation is afforded by the first-order equation :
$$(i\mathit{}\overline{\partial ̸}+2i\nu )\psi (x)=0,\text{(dS Dirac field equation)}$$
(6)
where $`\mathit{}=x.\gamma `$ in usual notations. For large $`R`$ behaviour $`(R\mathrm{})`$ we obtain the Dirac field equation in Minkowskian space. We know that any field quantity living on de Sitter space-time $`X_R`$ can be viewed as an homogenous function of the $`\mathrm{IR}^5`$-variable $`x^\alpha `$ with some arbitrarily chosen degree $`\sigma `$ . Let us choose a solution to the dS-Dirac equation of the type
$$\psi (x)=(\frac{1}{2}\gamma ^\alpha \gamma ^\beta M_{\alpha \beta }+i+\nu )\varphi (x)𝒰_T,$$
(7)
where $`\varphi (x)`$ is a scalar field with degree $`\sigma `$ and $`𝒰_T`$ is an arbitrary four-component spinor. T denotes the “orbital basis” of $`C^+`$ with respect to a subgroup $`L_e`$ of $`G_R`$ which is the stabilizer of an unit vector $`e`$ ( $`e^2=1`$) in $`\mathrm{IR}^5`$ . We have two solutions for $`\psi (x)`$ :
$$\psi _1^{\xi ,𝒱}(x)=(\frac{x.\xi }{R})^{2+i\nu }\frac{\mathit{}\mathit{\xi ̸}}{R\sqrt{2(\xi ^0+1)}}𝒰_T(\underset{+}{\overset{o}{\xi }})(\frac{x.\xi }{R})^{2+i\nu }𝒱(x,\xi ),$$
$$\psi _2^{\xi ,𝒰}(x)=(\frac{x.\xi }{R})^{2i\nu }\frac{\mathit{\xi ̸}\gamma ^4}{\sqrt{2(\xi ^0+1)}}𝒰_T(\underset{}{\overset{o}{\xi }})(\frac{x.\xi }{R})^{2i\nu }𝒰(\xi ).$$
(8)
We now consider the transformation of the Dirac free field $`\psi (x)`$ in such a way that the transformd $`\psi ^{}(x^{})`$ obyes the same dS-Dirac equation in the new frame . This lead to the simple relation $`\psi ^{}(x)=g\psi (\mathrm{\Lambda }^1(g)x)`$ where $`gSp(2,2)`$ is viewed as a $`4\times 4`$ matrix when acting on the 4-spinor $`\psi (x)`$ and $`\mathrm{\Lambda }SO(1,4)`$. We have
$$\psi ^{\xi ,𝒰_T}(x)=g\psi ^{\xi ,𝒰_T}(\mathrm{\Lambda }^1x)=\psi ^{\mathrm{\Lambda }\xi ,g𝒰_T}(x).$$
(9)
We see how the dS action is transfered onto the ”reciprocal space ” $`C^+`$ to which the parameter $`\xi `$ belongs, and the 4-component arbitary spinor $`𝒰_T`$. The plane-wave solutions to the free Dirac equation in Minkovskian space are given by the large-R behavior of the plane-wave solutions to the dS-Dirac field equation . For obtaining general field solutions, we consider the solution in the compolexed dS space-time $`X_R^{(c)}`$
$$X_R^{(c)}=\{z=x+iy\mathrm{lC}^5;\eta _{\alpha \beta }z^\alpha z^\beta =(z^0)^2\stackrel{}{z}.\stackrel{}{z}(z^4)^2=R^2\}$$
$$=\{(x,y)\mathrm{IR}^5\times \mathrm{IR}^5;x^2y^2=R^2,x.y=0\}.$$
(10)
Let $`T^\pm =\mathrm{IR}^5+iV^\pm `$ be the forward and backward tubes in $`\mathrm{lC}^5.V^+`$(resp. $`V^{})`$ stems from the causal structure on $`X_R`$,
$$V^\pm =\{x\mathrm{IR}^5;x^0\stackrel{>}{<}\sqrt{\stackrel{}{x}^2+(x^4)^2}\}.$$
(11)
We then introduce their respective intersections with $`X_R^{(c)},𝒯^\pm =T^\pm X_R^{(c)}`$, which will be called forward and backward tubes in the complex dS space-time $`X_R^{(c)}`$. Finally we define the set
$$𝒯_{12}=\{(z_1,z_2);z_1𝒯^+,z_2𝒯^{}\},$$
as a tube above $`X_R\times X_R`$ in $`X_R^{(c)}\times X_R^{(c)}`$. Details are given in . When $`z`$ varies in $`𝒯^+`$ (or $`𝒯^{}`$) and $`\xi `$ lies in the positive cone $`𝒞^+`$, the plane wave solutions $`\psi ^{\xi ,𝒱}(z)=(\frac{z.\xi }{R})^\sigma 𝒰(z,\xi ),\sigma \mathrm{lC}`$ are globally defined because the imaginary part of $`(z.\xi )`$ has a fixed sign. Now we can define the Wightman two-point function.
## 4 Two point function
Let us briefly recall the conditions we require on the Wightman two-point function
$$𝒲^\nu (x,y)=<\mathrm{\Omega },\psi (x)\overline{\psi }(y)\mathrm{\Omega }>,$$
where $`x,yX_R`$ and $`\overline{\psi }=\psi ^{}\gamma ^0\gamma ^4,`$ is the spinor field conjugate to $`\psi `$ . This function is $`4\times 4`$ matrix-valued in the present case, and has to satisfy the following requirements:
1. Positivity
for any test function $`f𝒟(X_R)`$ with values in $`\mathrm{lC}^4`$
$$_{X_R\times X_R}\overline{f}(x)𝒲(x,y)f(y)𝑑\sigma (x)𝑑\sigma (y)0,$$
(12)
where $`d\sigma (x)`$ denotes the dS-invariant measure on $`X_R`$.
2. Locality
for every space-like separated pair $`(x,y)`$, i.e. $`xy>R^2`$,
$$𝒲_{i\overline{j}}(x,y)=𝒲_{\overline{j}i}(y,x),$$
(13)
where $`𝒲_{\overline{j}i}(y,x)=<\mathrm{\Omega },\overline{\psi }_{\overline{j}}(y)\psi _i(x)\mathrm{\Omega }>`$.
3. Covariance
$$g𝒲(\mathrm{\Lambda }^1(g)x,\mathrm{\Lambda }^1(g)y)i(g^1)=𝒲(x,y),$$
(14)
where $`i(g^1)=\gamma ^4g^1\gamma ^4`$ ,
4. Normal analyticity
$`𝒲(x,y)`$ is the boundary value (in the sense of distributions) of a function $`W(z_1,z_2)`$ which is analytic in the domain $`𝒯_{12}`$.
The two-point $`W^\nu (z_1,z_2)`$ , lablled by the principal-series parameter $`\nu `$, is given by the following class of integral representations
$$W_{i\overline{j}}^\nu (z_1,z_2)=c_\nu _T(z_1.\xi )^{2i\nu }(z_2.\xi )^{2+i\nu }\underset{a=1,2}{}𝒰_i^a(\xi )\overline{𝒰}_{\overline{j}}^a(\xi )d\mu _T(\xi ),$$
and the two-point function $`W_{\overline{j}i}^\nu (z_2,z_1)`$ is given by the following class of integral representations
$$W_{\overline{j}i}^\nu (z_2,z_1)=H^2c_\nu _T(z_1.\xi )^{2+i\nu }(z_2.\xi )^{2i\nu }\underset{a=1,2}{}𝒱_i^a(z_1,\xi )\overline{𝒱}_{\overline{j}}^a(z_2,\xi )d\mu _T(\xi ).$$
$`d\mu _T(\xi )`$ is a measure invariant under $`L_e`$. We can write
$$W^\nu (z_1,z_2)=D(z_2)\gamma ^4𝒩(z_1,z_2),D(z_2)=\frac{1}{\nu +i}(i\mathit{}_2\overline{)\overline{}}_{z_2}+i+\nu ).$$
$`𝒩(z_1,z_2)`$ is a scalar two-point function
$$𝒩(z_1,z_2)=c_\nu _T(z_1.\xi )^{2i\nu }(z_2.\xi )^{1+i\nu }d\mu _T(\xi ).$$
(15)
In terms of the generalized Legendre function of the first kind we have
$$W_{i\overline{j}}^\nu (z_1,z_2)=C_\nu (D(z_2)\gamma ^4)_{i\overline{j}}P_{1+i\nu }^{(5)}(\frac{z_1.z_2}{R^2})=4C_\nu (D(z_2)\gamma ^4)_{i\overline{j}}(1\frac{z_1.z_2}{R^2})^{\frac{1}{2}}\text{P}_{i\nu }^1(\frac{z_1.z_2}{R^2})$$
$$W_{\overline{j}i}^\nu (z_2,z_1)=4H^2C_\nu (\mathit{}_1D(z_1)\mathit{}_2\gamma ^4)_{i\overline{j}}(1\frac{z_1.z_2}{R^2})^{\frac{1}{2}}\text{P}_{i\nu }^1(\frac{z_1.z_2}{R^2}).$$
(16)
The boundary value of $`W_{i\overline{j}}^\nu (z_1,z_2)`$ provides us with the following represntation for the Wightman two-point function :
$$𝒲(x,y)=c_\nu _T[(x.\xi )_+^{2i\nu }+e^{i\pi (2i\nu )}(x.\xi )_{}^{2i\nu }]$$
$$[(y.\xi )_+^{2+i\nu }+e^{i\pi (2+i\nu )}(y.\xi )_{}^{2+i\nu }]\mathit{\xi ̸}\gamma ^4d\mu _T.$$
(17)
This function satisfies the conditions of: a) positivity, b) locality, c) covariance, and d) normal analyticity . The existence of $`𝒲`$ allows one to make the QF formalism work . The spinor fields $`\psi (x)`$ are expected to be operator-valued distributions on $`X_R`$ acting on a Hilbert space $``$ . The Hilbert space $``$ of the representation can be described as the Hilbertian sum
$$=_0[\underset{n=1}{\overset{\mathrm{}}{}}A_1^{\scriptscriptstyle n}].$$
$`A`$ denotes the antisymmetrisation operation and $`_0=\{\lambda \mathrm{\Omega },\lambda \mathrm{lC}\}`$ where the vector $`\mathrm{\Omega }`$, cyclic for the polynomial algebra of field operators and invariant under the representation of $`G_R`$, is “the vacuum” . $`_1`$ is defined by the scalar product
$$(h_1,h_2)=_{X_R\times X_R}\overline{h}_1(x)𝒲(x,y)h_2(y)𝑑\sigma (x)𝑑\sigma (y)0,$$
(18)
where $`h𝒟(X_H)`$ with values in $`\mathrm{lC}^4`$. Each field operator $`\psi (f)`$ can be defined in terms of annihilation and creation operators by :
$$\psi (f)h^{(n)}(i_1,x_1;i_2,x_2;\mathrm{};i_n,x_n)=$$
$$\sqrt{n+1}_{X_R\times X_R}f^i(x)𝒲_{\overline{j}i}(y,x)h^{(n+1)}(\overline{j},y;i_1,x_1;\mathrm{};i_n,x_n)𝑑\sigma (x)𝑑\sigma (y)$$
$$+\frac{1}{\sqrt{n}}\underset{k=1}{\overset{n}{}}(1)^{k+1}f_{i_k}(x_k)h^{(n1)}(i_1,x_1;\mathrm{};\widehat{i_k},\widehat{x_k};\mathrm{};i_n,x_n),$$
(19)
where the $`i_k=1,2,3,4,`$ are spin indices. Here ‘$`\widehat{i_k}`$’ means omit it. By using the Fourier-Bros transformation on $`X_R`$ , we can write the unsmeared field operators $`\psi (x)`$
$$\psi (x)=_T\underset{a=1,2}{}\{a_a(\xi ,\nu )𝒰^a(\xi )[(x.\xi )_+^{2i\nu }+e^{i\pi (2i\nu )}(x.\xi )_{}^{2i\nu }]$$
$$+d_a^{}(\xi ,\nu )𝒱^a(x,\xi )[(x.\xi )_+^{2+i\nu }+e^{i\pi (2+i\nu )}(x.\xi )_{}^{2+i\nu }]\}d\mu _T(\xi ),$$
(20)
where $`a_a(\xi ,\nu )`$ and $`d_a(\xi ,\nu )`$ are defined by $`a_a(\xi ,\nu )\mathrm{\Omega }>=0=d_a(\xi ,\nu )\mathrm{\Omega }>`$. The field anticommutator is given by
$$\{\psi _i(x),\overline{\psi }_{\overline{j}}(y)\}=𝒲_{i\overline{j}}(x,y)+𝒲_{\overline{j}i}(y,x).$$
It can be esily checked that for space-like separated points (x,y) we have $`\{\psi _i(x),\overline{\psi }_{\overline{j}}(y)\}=0`$ . The integral representation for the two-point function is defined on the space which carries the principal series of the dS group $`Sp(2,2)/\text{ZZ}_2SO_0(1,4)`$. In the limit $`R\mathrm{}`$, we obtained
$$\underset{R\mathrm{}}{lim}\{\psi (x),\overline{\psi }(y)\}=$$
$$\frac{1}{2(2\pi )^3}\{e^{ik.(XY)}(\mathit{}\gamma ^4+m)+e^{ik.(XY)}(\mathit{}\gamma ^4m)\}\frac{d^3k}{k^0},$$
(21)
which is of the same form of the Minkowskian space.
## 5 Conclusion and outlook
Using the dS-plane waves in tube domains for spinor field, one can construct an analytic function $`W(z_1,z_2)`$ that its boundary value is the Wightman two-point function $`𝒲(x,y)`$. Then the dS-plane-waves allow us to construct the quantum field on dS space in the same way as the quantum field on Minkowski space. In the case of the massless spinor field, we must replace $`\nu `$ with $`0`$ in the Waightman two-point function as well as in the field operator $`\psi (x)`$ . In this case the corresponding UIR is known as the first term of the spinor discrete series of representation, which is written as $`\mathrm{\Pi }_{\frac{1}{2},\frac{1}{2}}^\pm `$ in . The $`\pm `$ define the helicity of the massless spinor field. We now intend to use these methods to construct a covariant quantum fields with spin-1 and spin-2. In the case of the massive field the procedure is the same as spin-$`\frac{1}{2}`$. In the case of the massless field, we must use the Gupta-Bleuler quantization for obtaining a fully covariant theory. At this moment, this kind of work is in progress on the ”massless” representations of dS group and the related QFT.
Acknowlegements: I am grateful to Professor J. P. Gazeau for giving inspiration to this investigation and for a fruitful discussion about it.
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# A Remark on BRST Singlets
## Abstract
Negative norm Hilbert space state vectors can be BRST invariant, we show in a simplified Y-M model that such states can be created by starting with gluons only.
Faddeev and Slavnov pointed out that BRST invariance alone is not sufficient for a Y-M type theory to be physically admissible since “physical” states that are annihilated by the generators of BRST transformation may have negative norm. They emphasized the need for additional investigation of the BRST-singlet sector. By using a second quantized F-P model I find that the transition amplitude from an initial BRST singlet of positive norm to a final state of negative norm is different from zero. There is no barrier to the production of states of negative norm.
The minimal Y-M Lagrangian for gluons, ghosts and antighost fields is
$$=_o(A)+_{GF}+_{FP}$$
(1)
$$_o=\frac{1}{4}F_{\mu \nu }^aF^{\mu \nu a}$$
(2)
$$_{GF}=\frac{1}{2}(_\mu A^{\mu a})^2$$
(3)
$$_{FP}=i^\mu \overline{C}^a(D_\mu C)^a$$
(4)
$$D_\mu =_\mu igA_\mu $$
(5)
The lowest order transition amplitude from an initial state with two gluons to a final ghost, anti-ghost state requires two types of Feynman diagrams. The one, with two gluon-ghost vertices gives the amplitude
$$\frac{i}{2}g^2\left[p_1\epsilon (k_1)p_2\epsilon (k_2)C_{a_1b_1f}C_{a_2b_2f}+p_1\epsilon (k_2)p_2\epsilon (k_1)C_{a_2b_1f}C_{a_1b_2f}\right]$$
(6)
The other, which has one gluon-ghost and a three gluon vertex, gives
$$\frac{ig^2}{2k_1k_2}\left[\epsilon (k_1)\epsilon (k_2)(k_1k_2)p_12p_1\epsilon (k_1)k_1\epsilon (k_2)+2p_1\epsilon (k_2)k_2\epsilon (k_1)\right]C_{a_1a_2f}C_{b_1b_2f}$$
(7)
where $`C_{abc}`$ are the structure constants of the group $`SU(n)`$. Repeated indices are summed. The incident gluon momenta, polarization and color indices respectively are
$`k_1,\epsilon (k_1),a_1;k_2,\epsilon (k_2),a_2`$. The ghost, antighost momenta, and color indices are $`p_1,b_1;p_2,b_2`$. A common feature of ghost production amplitudes is their gauge dependence, this will be seen later to follow directly from BRST invariance. The replacement $`\epsilon (k)\epsilon (k)+\lambda k`$ changes the sum of (6) and (7), whereas ghost free amplitudes are invariant under this residual gauge transformation.
From the anticommutation relations of ghost and antighost fields it follows that the final state in this process
$$\mathrm{\Psi }=\overline{c}_{b_1}^+(p_1)c_{b_2}^+(p_2)0$$
(8)
has zero norm if $`p_1p_2`$. Whereas for smeared ghost, antighost states
$$\mathrm{\Psi }^{^{}}=\underset{¯}{d}p_1\underset{¯}{d}p_2f(p_1)g(p_2)\overline{c}_{b_1}^+(p_1)c_{b_2}^+(p_2)0$$
(9)
the norm is
$$\mathrm{\Psi }^{^{}}\mathrm{\Psi }^{^{}}=f(p_1)g(p_2)\underset{¯}{d}p_1\underset{¯}{d}p_2^2\delta _{b_1b_2}$$
(10)
A negative norm results when $`b_1=b_2`$ for an appreciable overlap of the two momentum distribution $`f`$ and $`g`$. Therefore the sum of (6) and (7), with $`p_1p_2`$, is the amplitudes of a zero norm vector in Fock space.
The amplitude in Eq.(7) is ambiguous for $`p_1p_2`$ since then $`k_1k_2p_2`$. We avoid this by including a gluon in the final state with momenta, polarization and color index $`k_3,\epsilon (k_3),a_3`$ respectively. To obtain all amplitudes of order $`g^3`$ requires the quartic and cubic gluon-gluon vertices as well as a gluon-ghost interaction. For $`b_1=b_2=b`$ the amplitude which involves the quartic vertex vanishes. For $`p_1=p_2=p`$ the diagram that has three gluon-ghost vertices gives
$$\frac{g^3}{4}[\frac{1}{pk_1pk_3}p\epsilon (k_1)p\epsilon (k_3)(k_1+k_3)\epsilon (k_2)C_{a_1bf}C_{a_2fg}C_{a_3bg}$$
$$+\frac{1}{pk_2pk_3}p\epsilon (k_2)p\epsilon (k_3)(k_2+k_3)\epsilon (k_1)C_{a_2bf}C_{a_1fg}C_{a_3bg}$$
$$\frac{1}{pk_1pk_2}p\epsilon (k_1)p\epsilon (k_2)(k_1k_2)\epsilon (k_3)C_{a_1bf}C_{a_3fg}C_{a_2bg}]$$
(11)
where all repeated indices except $`b`$ are summed. In order to show that the product of the three structure constant in Eq.(11) does not vanish identically it is convenient to sum over $`b`$, and we obtain $`C_{a_1a_2a_3}`$.
For the same process, the amplitude that contains a gluon-ghost and the cubic gluon interaction gives
$$\frac{g^3}{2}\frac{p\epsilon (k_1)}{pk_1k_2k_3}[\epsilon (k_2)\epsilon (k_3)(k_2+k_3)p2p\epsilon (k_3)k_3\epsilon (k_2)$$
$$2p\epsilon (k_2)k_2\epsilon (k_3)]C_{a_1br}C_{rbl}C_{la_2a_3}$$
$$+\frac{g^3}{2}\frac{p\epsilon (k_2)}{pk_2k_1k_3}[\epsilon (k_1)\epsilon (k_3)(k_1+k_3)p2p\epsilon (k_3)k_3\epsilon (k_1)$$
$$2p\epsilon (k_1)k_1\epsilon (k_3)]C_{a_2br}C_{rbl}C_{la_1a_3}$$
$$+\frac{g^3}{2}\frac{p\epsilon (k_3)}{pk_3k_2k_1}[\epsilon (k_2)\epsilon (k_1)(k_2k_1)p+2p\epsilon (k_1)k_1\epsilon (k_2)$$
$$2p\epsilon (k_2)k_2\epsilon (k_1)]C_{a_3br}C_{rbl}C_{la_2a_1},$$
(12)
and since
$$\underset{b,r}{}C_{mbr}C_{mbr}\alpha \delta _{mn}$$
(13)
it follows that the products of the structure constants is not zero for every $`a_1,a_2,a_3`$ and $`b`$. Eqs.(12) and (13) give the transition amplitude to a ghost, antighost state of negative norm. Just as in the previous example the answer is gauge dependent.
In covariant gauges, the gauge independence of transition amplitudes between physical states is known to follow from BRST symmetry, in contrast to this the transition amplitude from physical states to ghost states depends on the gauge. The gauge dependence in Eqs.(6), (7), (11) and (12) is a non-perturbative result arising the BRST invariance of physical wave functions. For instance, for an infinitesimal BRST transformation $`\delta `$,
$$\delta oT\left(\overline{C}^{a_1}(x_1)A^{\mu _2,a_2}(x_2)_{\mu _3}A^{\mu _3,a_3}(x_3)_{\mu _4}A^{\mu _4,a_4}(x_4)\right)o=o$$
(14)
from which it follows that
$$oT\left(\delta \overline{C}^{a_1}(x_1)A^{\mu _2,a_2}(x_2)_{\mu _3}A^{\mu _3,a_3}(x_3)_{\mu _4}A^{\mu _4,a_4}(x_4)\right)o$$
$$=oT\left(\overline{C}^{a_1}(x_1)\delta A^{\mu _2,a_2}(x_2)_{\mu _3}A^{\mu _3,a_3}(x_3)_{\mu _4}A^{\mu _4,a_4}(x_4)\right)o$$
(15)
where
$$\delta \overline{C}^a=i^\mu A_\mu ^a$$
$$\delta A_\mu ^a=\left(D_\mu C\right)^a$$
(16)
are the standard BRST variations. Letting $`x_1^o,x_2^o+\mathrm{}`$ and $`x_3^o,x_4^o\mathrm{}`$, we use the LSZ scattering formalism. The lhs of Eq.(15) is proportional to the gluon-gluon scattering amplitude. The rhs is proportional to the gluon-gluon $``$ ghost-antighost amplitude. The l.h.s. is not zero since only three gluons are contracted with their momenta. This shows again that the gauge dependence of what was found in Eqs.(6) and (7) is not an artifact of the weak coupling approximation.
Covariant quantization of vector fields is known to require an enlarged Hilbert space, containing states of positive, zero and negative norm. Hence the time evolution of a state of positive norm will in general yield a superposition of all three states. Electrodynamics avoids the conceptual problem of negative norm states in two ways, by imposing the Gupta-Bleuler constraint on physical states or by showing that time-like and longitudinal photon contributions to the cross-section cancel each other.
In both cases that were dealt with above the initial state has positive norm, therefore the final state will be a Hilbert space superposition of positive and zero norm states in the first case. Even though the latter are not observable they can not be ignored completely unless they are orthogonal to physical states of positive norm. In the second example we find that production of ghost-antighost states of negative norm is not prohibited.
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# Aspects of nonlocality in atom-photon interactions in a cavity
## I Introduction
Nonlocality at the quantum level manifests itself in various kinds of phenomena. The study of this so far, has been predominantly confined to the study of interactions amongst similar kinds of particles, for example, photon-photon interactions or the interaction of subatomic particles among themselves. With the advance of technology over the last several years, it has now become conceivable to investigate a new kind of nonlocality in a controllable fashion, viz., the nonlocality generated through the interaction of distinct entities, like atoms and photons inside a high quality cavity.
The mathematical framework for demonstrating the violation of local realism in quantum mechanics was first provided by Bell through his famous inequalities. This work was subsequently generalized and also extended to consider the interaction of more than two particles. A different kind of proof of nonlocality without the use of inequalities, also exists . Furthermore, it has been shown that quantum nonlocality continues to persist even for the case of a large number of particles, or large quantum numbers. This has raised certain questions regarding the issue of the macroscopic or classical limit of quantum mechanics in examples where both the number of particles, and their quantum number is made arbitrarily large. The phenomenon of environment induced decoherence is of direct relevance here. It would be interesting if decoherence could be experimentally controlled and its effect on nonlocality be quantitatively monitored in particular examples of study.
Experimental proposals of demonstrating nonlocality have mostly been concerned with spin-$`1/2`$ particles, photons, or mesons. In recent times, several schemes involving two-level Rydberg atoms have been proposed. In such schemes two-level rydberg atoms are considered in analogy to the spin system in Bell’s original reasoning. The role of the polarization axis of the Stern-Gerlach apparatus used for spin-$`1/2`$ systems is played here by the phase of an auxiliary electromagnetic field. The primary advantages of the experiments using atoms, compared to those with photons or spin half particles are two-fold. First, the realization of spacelike separation for Rydberg atoms is easier because of their smaller velocities than photons or electrons. Secondly, the efficiencies of detectors used for the former is much larger in general in comparison with the detectors used for elementary particles. In addition, the interaction of large-sized atoms with the environment can be significant enough to be monitored in certain cases. In fact, in the experimental schemes which involve the interaction of Rydberg atoms with photons in a microwave cavity, dissipation through the loss of cavity photons always occurs. The effect of this is manifested in the form of loss of coherence in the atom-photon interactions. Thus, this is a natural arena to study the effects of decoherence on quantum nonlocality in a quantitative manner.
In this paper we propose a realistic experiment to test Bell’s inequality for real micromasers/microlasers by taking experimentally attainable data in the presence of both atomic decay and cavity dissipation under the influence of their respective reservoirs. The approach followed by us enables us to analyze both micromaser as well as microlaser dynamics within the framework of a single formalism in presence of decoherence. Our aim is to study the effects of decoherence on the magnitude of violation of Bell’s inequality in an experimentally controllable fashion. As an interesting sidelight, we are also able to demonstrate the effect of decoherence on multiparticle correlations that are a natural outcrop of experimental schemes using microwave cavity that we use for our analysis. In the next section we describe the experimental arrangement and the relevant Bell-type inequality (BI). We wish to emphasize that although the BI used by us is the same as in , the cavity dynamics considered by us (the steady-state dynamics is discussed in Section III) differs crucially form the one used in . We are hence able to take into account micromaser dissipation in a more realistic manner, and also analyse the microlaser by incorporating atomic decay. In Section IV we present and discuss our results. Section V is reserved for some concluding remarks.
## II Violations of Bell’s inequality in a microcavity
We consider the following experimental scenario. A two-level atom initially in its upper excited state $`|e>`$ traverses a high-Q single mode cavity. The cavity is in a steady state and tuned to a single mode resonant with the transition $`|e>|g>`$. The emerging single-atom wavefunction is a superposition of the upper $`|e>`$ and lower $`|g>`$ state, and it leaves an imprint on the photonic wavefunction in the cavity. After leaving the cavity, the atom passes through an electromagnetic field which gives it a $`\pi /2`$ pulse the phase of which can be varied for different atoms. The atom then reaches the detector, placed at a sufficient distance, capable of detecting the atom only in the upper or lower state. Thus, the role of the $`\pi /2`$ pulse may be considered as a component of the detection mechanism in the experiment. During the whole process, dissipation takes place, and is taken into account. Next, this process is repeated for a similar second atom. The important difference is that the second atom interacts with a photonic wavefunction which has been modified due to the passage of the first atom. There is no direct interaction between the two atoms, although secondary correlations develop between them. In other words, though there is no spatial overlap between the two atoms, the entanglement of their wavefunctions with the cavity photons can be used to formulate local-realist bounds on the detection probabilities for the two atoms.
The interplay of the atomic statistics with the photonic statistics plays a crucial role in the investigation of nonlocality here. The initial state of the cavity is built up by the passage of a large number of atoms, but only one at a time, through it. The pump parameter and the atom-photon interaction time are key inputs for the profile of the resultant photonic wavefunction which in turn governs the nature of entanglement between two successive experimental atoms detected in their upper or lower states by the detector. As stated earlier, dissipation due to the interaction of the pumping atoms with their reservoir, as well as the loss of cavity photons can be controlled, and their effects on the statistics of detected atoms can be studied. The formalism used by us has another generic feature. The effects of decoherence on nonlocality can be studied in context of the micromaser, as well as the microlaser, its optical counterpart.
It is easy to obtain a Bell-type inequality suitable for the scenario considered by us in analogy to Bell’s original reasoning. Two level Rydberg atoms are analogous to spin-$`1/2`$ systems, and the phase of the electromagnetic field plays the role of the polarization axis of the Stern-Gerlach apparatus used for spin-$`1/2`$ systems. In fact, several local realist bounds have earlier been derived to tailor such a situation. Let us very briefly describe one such derivation which we shall use in the present analysis. Assigning the value $`+1`$ for the atom detected in the upper state $`|e>`$, and $`1`$ for the lower state $`|g>`$, one can in any local realist theory define functions $`f(\varphi _1)=\pm 1;g(\varphi _2)=\pm 1`$ describing the outcome of measurement on the atom 1 and 2 when the phase of the electromagnetic field giving $`\pi /2`$ pulse to the atoms is set to be $`\varphi _1`$ and $`\varphi _2`$ for the respective atoms. The ensemble average for double click events is therefore defined as
$$E^\lambda (\varphi _1,\varphi _2)=𝑑\lambda f(\varphi _1)g(\varphi _2)$$
(1)
where $`\lambda `$ is a suitable probability measure on the space of all possible states.
The quantum mechanical expectation value for double click events is calculated from the probabilities of all possible double-click sequences. This is given by
$`E(\varphi _1,\varphi _2)`$ $`=`$ $`P_{ee}(\varphi _1,\varphi _2)+P_{gg}(\varphi _1,\varphi _2)`$ (2)
$``$ $`P_{eg}(\varphi _1,\varphi _2)P_{ge}(\varphi _1,\varphi _2)`$ (3)
where $`P_{eg}(\varphi _1,\varphi _2)`$ stands for the probability that the first atom is found to be in state $`|e>`$ after traversing the $`\pi /2`$ pulse with phase $`\varphi _1`$, and the second atom is found to be in state $`|g>`$ with the phase of the $`\pi /2`$ pulse being $`\varphi _2`$ for its case. Defining $`E_0=E^\lambda (\varphi _1=\varphi _2)`$ and $`M_0=P_{ee}(\varphi _1=\varphi _2)+P_{gg}(\varphi _1=\varphi _2)`$, and assuming perfect detections, it follows that $`E_0=2M_01`$. Further, it is easy to see that $`f(\varphi )=+g(\varphi )`$ with probability $`M_0`$, and $`f(\varphi )=g(\varphi )`$ with probability $`(1M_0)`$. Hence, $`E^\lambda (\varphi _1,\varphi _2)`$ can be written as
$`E^\lambda (\varphi _1,\varphi _2)=E_0{\displaystyle 𝑑\lambda f(\varphi _1)f(\varphi _2)}`$ (4)
Now, one can define a Bell sum
$`B`$ $``$ $`|E^\lambda (\varphi _1,\varphi _2)E^\lambda (\varphi _1,\varphi _3)|`$ (5)
$`+`$ $`sign(E_0)[E^\lambda (\varphi _2,\varphi _3)E_0]`$ (6)
It follows immediately from (2-4) that $`B0`$. In the next section we shall calculate this Bell sum $`B`$ and see how it evolves for various values of the cavity parameters. It is convenient to set the values of the phases $`\varphi _1=0`$, $`\varphi _2=\pi /3`$, and $`\varphi _3=2\pi /3`$, as for these values the Bell-type inequality is always violated, i.e., $`B>0`$ for the case of an idealised micromaser.
## III Steady-state micromaser/microlaser photon statistics
In realistic situations, one must consider the interacting systems (atoms as well as cavity field) coupled to their respective reservoirs. The couplings are governed by their equations of motion,
$`\dot{\rho }|_{atomreservoir}`$ $`=`$ $`\gamma (1+\overline{n}_{th})(s^+s^{}\rho 2s^{}\rho s^++\rho s^+s^{})`$ (7)
$``$ $`\gamma \overline{n}_{th}(s^{}s^+\rho 2s^+\rho s^{}+\rho s^{}s^+)`$ (8)
for the atom and
$`\dot{\rho }|_{fieldreservoir}`$ $`=`$ $`\kappa (1+\overline{n}_{th})(a^{}a\rho 2a\rho a^{}+\rho a^{}a)`$ (9)
$``$ $`\kappa \overline{n}_{th}(aa^{}\rho 2a^{}\rho a+\rho aa^{})`$ (10)
for the field. $`\rho `$ is the reduced density operator obtained after tracing over the reservoir. $`\gamma `$ and $`\kappa `$ are the decay constants for the atom and the field respectively. $`\overline{n}_{th}`$ is the average photon number representing the reservoir. $`s^+`$ and $`s^{}`$ are the usual Pauli operators for the pseudo-spin representation of the two-level model of the atom. $`a(a^{})`$ is the photon annihilation (creation) operator.
The dynamics we are interested in, involves two-level atoms steamed into a single mode cavity in such a way that there is at most one atom in the cavity at any time. Thus we have sequences of events (atom-field interactions) taking place randomly, but with each event of a fixed duration $`\tau `$. This interaction is governed by
$$\dot{\rho }|_{atomfield}=i[H,\rho ]$$
(11)
where $`H=g(s^+a+s^{}a^{})`$ is the well known Jaynes-Cummings Hamiltonian with $`g`$ being the coupling constant.
Thus, we have to solve the equation of motion
$$\dot{\rho }=\dot{\rho }|_{atomreservoir}+\dot{\rho }|_{fieldreservoir}+\dot{\rho }|_{atomfield}$$
(12)
where the terms on the r.h.s. are the r.h.s’ of (5), (6) and (7) respectively. For the duration between two events, we have to solve the equation of motion (6) only.
The steady-state photon statistics of the cavity field undergoing such dynamics has been derived in. In the cavity photon number representation, the probabilities $`P_n=<n|\rho _f^{(ss)}|n>`$ are given by
$$P_n=P_0\underset{m=1}{\overset{n}{}}v_m$$
(13)
$`P_0`$ is obtained from the normalisation $`_{n=0}^{\mathrm{}}P_n=1`$. The $`v_n`$ are given by the continued fractions
$$v_n=f_3^{(n)}/(f_2^{(n)}+f_1^{(n)}v_{n+1})$$
(14)
with $`f_1^{(n)}=(Z_n+C_n)/\kappa `$, $`f_2^{(n)}=2N+(Y_n+B_n)/\kappa `$ and $`f_3^{(n)}=(X_n+A_n)/\kappa `$. $`\kappa `$ is the cavity bandwidth and $`N=R/2\kappa `$ is the number of atoms passing through the cavity in a photon lifetime. $`A_n=2n\kappa \overline{n}_{th}`$, $`B_n=2\kappa (n+\overline{n}_{th}+2n\overline{n}_{th})`$ and $`C_n=2(n+1)(\overline{n}_{th}+1)\kappa `$. $`X_n`$, $`Y_n`$ and $`Z_n`$ are given by
$$X_n=R\mathrm{sin}^2(g\sqrt{n}\tau )\mathrm{exp}\{[\gamma +(2n1)\kappa ]\tau \}$$
(15)
$`Y_n={\displaystyle \frac{1}{2}}R(\{2\mathrm{cos}^2[g\sqrt{n+1}\tau ]{\displaystyle \frac{1}{2}}(\gamma /\kappa +2n+1)`$ (16)
$`+F_1(n1)\}\mathrm{exp}\{[\gamma +(2n+1)\kappa ]\tau \}+[{\displaystyle \frac{1}{2}}(\gamma /\kappa +2n+1)`$ (17)
$`F_2(n1)]\mathrm{exp}\{[\gamma +(2n1)\kappa ]\tau \}),`$ (18)
and
$`Z_n`$ $`=`$ $`{\displaystyle \frac{1}{2}}R([{\displaystyle \frac{1}{2}}(\gamma /\kappa +2n+3)+F_2(n)]\mathrm{exp}\{[\gamma +(2n+1)\kappa ]\tau \}`$ (19)
$`[{\displaystyle \frac{1}{2}}(\gamma /\kappa +2n+3)+F_1(n)]\mathrm{exp}\{[\gamma +(2n+3)\kappa ]\tau \})`$ (20)
$`\gamma `$ represents reservoir induced spontaneous emission from the upper to the lower masing level. The functions $`F_i`$ are
$`F_i(n)={\displaystyle \frac{\kappa /4g}{(\sqrt{n+2}\sqrt{n+1})^2}}`$ (21)
$`[{\displaystyle \frac{\gamma }{\kappa }}(\sqrt{n+2}\sqrt{n+1})\mathrm{sin}(2g\sqrt{m}\tau ){\displaystyle \frac{\gamma }{g}}\mathrm{cos}(2g\sqrt{m}\tau )`$ (22)
$`[2n+3+2\sqrt{(n+1)(n+2)}](\sqrt{n+2}\sqrt{n+1})\mathrm{sin}(2g\sqrt{m}\tau )]`$ (23)
$`+{\displaystyle \frac{\kappa /4g}{(\sqrt{n+2}+\sqrt{n+1})^2}}`$ (24)
$`[\pm {\displaystyle \frac{\gamma }{\kappa }}(\sqrt{n+2}+\sqrt{n+1})\mathrm{sin}(2g\sqrt{m}\tau ){\displaystyle \frac{\gamma }{g}}\mathrm{cos}(2g\sqrt{m}\tau )`$ (25)
$`[2n+32\sqrt{(n+1)(n+2)}](\sqrt{n+2}+\sqrt{n+1})\mathrm{sin}(2g\sqrt{m}\tau )]`$ (26)
where $`m=n+2`$ and $`n+1`$ for $`i=1`$ and $`2`$, respectively, with the upper sign for $`i=1`$.
The experimental atoms on which we plan to test the Bell’s inequality (BI), encounter this steady state radiation field $`\rho _f^{(ss)}`$. The atom-field interaction is, as mentioned earlier governed by the Jaynes-Cummings Hamiltonian $`H`$. After interaction with the cavity, the first experimental atom emerges in a superposition of the upper ($`|e>`$) and the lower ($`|g>`$) states and experiences a $`\pi /2`$ pulse with phase $`\varphi _1`$. The probability of detection of the atom in the upper state $`|e>`$ and the lower state $`|g>`$ can be written respectively as
$`P_e=\mathrm{Tr}._\mathrm{f}𝒫_e`$ (27)
$`P_g=\mathrm{Tr}._\mathrm{f}𝒫_g`$ (28)
with
$`𝒫_e={\displaystyle \frac{1}{2}}[𝒜\rho _f^{(ss)}𝒜^{}+𝒟\rho _f^{(ss)}𝒟^{}\{e^{i\varphi _1}𝒜\rho _f^{(ss)}𝒟^{}+e^{i\varphi _1}𝒟\rho _f^{(ss)}𝒜^{}\}]`$ (29)
$`𝒫_g={\displaystyle \frac{1}{2}}[𝒜\rho _f^{(ss)}𝒜^{}+𝒟\rho _f^{(ss)}𝒟^{}+\{e^{i\varphi _1}𝒜\rho _f^{(ss)}𝒟^{}+e^{i\varphi _1}𝒟\rho _f^{(ss)}𝒜^{}\}]`$ (30)
where trace is taken over the cavity field and the operators $`𝒜`$ and $`𝒟`$ are given by
$`𝒜=\mathrm{cos}(gt\sqrt{a^{}a+1})`$ (31)
$`𝒟=ia^{}{\displaystyle \frac{\mathrm{sin}(gt\sqrt{a^{}a+1})}{\sqrt{a^{}a+1}}}`$ (32)
After the passage of the first atom through the cavity and its detection in, for example, the state $`|e>`$, the second atom encounters the cavity field with density operator $`\rho _f^{(2)}`$ given by
$$\rho _f^{(2)}=[𝒜\rho _f^{(ss)}𝒜^{}+𝒟\rho _f^{(ss)}𝒟^{}\{e^{i\varphi _1}𝒜\rho _f^{(ss)}𝒟^{}+e^{i\varphi _1}𝒟\rho _f^{(ss)}𝒜^{}\}]$$
(33)
(since $`\mathrm{Tr}._\mathrm{f}𝒫_e=1/2`$). The phase of the $`\pi /2`$ pulse is set to $`\varphi _2`$ for the second atom. $`𝒫_e`$ for the second atom is given by
$$𝒫_e^{(2)}=\frac{1}{2}[𝒜\rho _f^{(2)}𝒜^{}+𝒟\rho _f^{(2)}𝒟^{}\{e^{i\varphi _1}𝒜\rho _f^{(2)}𝒟^{}+e^{i\varphi _1}𝒟\rho _f^{(2)}𝒜^{}\}]$$
(34)
The conditional probability $`P_{ee}(\varphi _1,\varphi _2)`$ is thus given by
$$P_{ee}(\varphi _1,\varphi _2)=\mathrm{Tr}._\mathrm{f}𝒫_e^{(2)}$$
(35)
$`P_{gg}`$ is obtained similarly, and using the relations $`E(\varphi _1,\varphi _2)=2P(\varphi _1,\varphi _2)1`$ and $`P(\varphi _1,\varphi _2)=P_{ee}(\varphi _1,\varphi _2)+P_{gg}(\varphi _1,\varphi _2)`$ one obtains the values of $`E`$ in the Bell sum (4). We perform this exercise numerically with the values of the phases in the Bell sum (4) set to $`\varphi _1=0`$, $`\varphi _2=\pi /3`$, and $`\varphi _3=2\pi /3`$.
It will be appropriate to mention here that although decoherence effects are inherent in the build up of the cavity field to its steady state since it encounters a large number of atoms over the time required for the steady state to be reached, the dynamics of a single experimental atom interacting with this field for a short duration will be unaffected by the decoherence effects, as has been observed in the cavity-QED of Jaynes-Cummings interaction. It was shown there that typical durations of atom-field interactions in realistic environments can be as large as $`t10/g`$ up to which decoherence effects are insignificant. We again stress here that even though the individual atom-field interaction time in the dynamics of the cavity field, pumped repeatedly with atoms, is uniformly of this order, the dissipative forces there play a crucial role due to the number of atoms ($`1`$) and the time ($`10/g`$) involved in reaching the steady state. The same argument applies to the case of atomic damping, as well. Indeed, we keep the above argument in mind in our choice of parameters while calculating the probabilities in the Bell sum (4).
The micromaser model considered by us differs crucially from that in because in the latter model cavity dissipation is neglected whenever an atom is present in the cavity. Our model takes into account dissipation even during the short atom-field interaction times. This a priori small effect on cavity photon statistics gets magnified due to the requirement of passage of a large number of atoms through the cavity for it to reach its steady state. Thus, as expected, the resultant steady-state photon statistics in our model is clearly different from the one used in . Let us again emphasize, that although the Bell’s inequality (BI) we propose to test is the same as in , our cavity dynamics are fundamentally different: (a) our micromaser model (discussed in detail in ) is more realistic, and (b) we are also able to analyse the microlaser by incorporating atomic decay.
## IV Effects of cavity dissipation and atomic decay on the Bell sum
In the previous section we have presented steady-state cavity dynamics describing both the micromaser as well as the microlaser within a unified framework. However, their distinctive features are manifested in the choice of parameters which we use below to study the violation of BI in both separately. Our choice of different sets of experimentally attainable parameters for both the micromaser and the microlaser is motivated by the following considerations. Recall that the experimental atoms are pumped into the cavity with steady-state photon statistics at such a rate that at most one atom is present in the cavity at any time. So, there are two situations: (a) atom present in the cavity, and (b) cavity empty of atom. Whatever be the situation, cavity dissipation, i.e., the interaction of the cavity mode with its reservoir continues uninterrupted. This process is of crucial importance in micromaser dynamics, and has been discussed in detail in . However, while considering micromaser dynamics, one can safely set atomic decay to zero. This is because the Rydberg levels involved in the micromaser have a lifetime of about $`0.2s`$, whereas the atomic flight time through the cavity (atom-field interaction time) is typically $`35\mu s`$. We present the variation of the Bell sum $`B`$ with respect to the parameter $`D`$ (which is nothing but the Rabi angle $`\varphi `$ scaled by the pump rate $`N`$, i.e., $`D=\varphi \sqrt{N}`$) for the case of the micromaser in Figure 1.
Our results show clearly that the Bell sum reflecting the degree of nonlocality exhibited in the atom-atom secondary correlations depends heavily on cavity dissipation. In particular, it is seen that the value of Bell sum $`B`$ decreases with the increase of pump rate $`N`$ for a large range of interaction times $`\tau `$. This can be understood from the way dissipative effects creep into the dynamics through two parameters $`N`$ and $`\tau `$. The genesis of atom-photon and the resultant atom-atom entanglement competes with decoherence in an interesting fashion over time. For shorter values of single atom interaction times we find that the correlations build up sharply with $`\tau `$, and the peak value of $`B`$ signifying maximum violation of BI is larger for higher values of $`N`$, (as a magnification of Figure 1 reveals). The maximum violation (the value of $`B`$ at its second peak) is $`0.5812`$ for $`N=20`$ (full line), $`0.6079`$ for $`N=50`$ (broken) and $`0.6241`$ for $`N=100`$ (dotted) (See Figure 2 where we magnify the curves in the region $`D5`$ of Figure 1 for a clear display of these peaks). This feature is a curious example of multiparticle induced nonlocality. It is analogous to the enhancement of nonlocality for multiparticle systems, and is in conformity with the mathematical demonstration of larger violation of BI with increase in the number of particles involved. For short interaction times, naturally the effects of decoherence are too small to affect the correlations. One noticeable feature in Figure 1 is the structures for low values of $`N`$ (full line). This originates from “trapped state” dynamics of photonic statistics where it has been shown that dissipative effects gradually wash out such states, as can be seen from the broken and dotted lines.
Let us now consider the violation of BI in the microlaser. It is well known that atomic decay is an important factor in the microlaser where atomic levels at optical frequencies are involved. Although decay is unimportant for the dynamics of a single atom interacting with the field up to a certain interaction time, its accumulated effect for a large number of atoms is crucial for the evolution of the microlaser field to a steady state. Furthermore, the interaction time of the atom with the $`\pi /2`$ pulse (in this case it is $`gt=\pi /2`$) is far less compared to even individual atom-field interaction times of interest in microlaser dynamics. For this reason the effect of atomic decay can be neglected during interaction of the atom with the $`\pi /2`$ pulse. Our results for the microlaser are shown in Figure 3.
The effect of decoherence (cavity leakage rate $`\kappa /g`$) on the violation of BI is clearly seen in the microlaser. One can check, similar to the case of the micromaser that the value of Bell sum $`B`$ decreases with the increase of pump rate $`N`$ for a large range of interaction times $`\tau `$. But in case of the microlaser, atomic damping $`\gamma `$ is a dominating factor, in contrast to the cavity photon loss in the micromaser. This makes the Bell sum fall off rapidly for large values of $`\tau `$. The second peak in $`B`$ however, is a consequence of the Jaynes-Cummings dynamics and survives such dissipation.
## V Conclusions
To summarize, we have shown that a demonstration of nonlocality, encompassing several of its varied aspects in atom-photon interactions in cavities and the effects of decoherence on it, can be possible in experimental set-ups already in operating conditions for the micromaser, as well as for the microlaser. The formalism chosen enables us to consider the steady-state dynamics of both the micromaser as well as the microlaser in the presence of atomic decay and cavity dissipation within a unified framework. Their distinguishing features are brought about by the different values of the parameters chosen to analyze the violation of BI in them separately. Certain notable features, such as enhancement of nonlocality for increased number of atoms, when decoherence effects are small, can be observed. We have seen how such features can be quantitatively monitored by control of decoherence. In an actual experiment, certain points have to be borne in mind. A few atoms may go undetected between two detector clicks. However, the steady state nature of the cavity field contributes to making the effect of this on the Bell sum insignificant compared to the effect of decoherence which we have probed in detail. Finally, the observed magnitude of violation of BI would be brought down by finite detector efficiency. Nevertheless, our selection of the particular type of BI, and the phases of $`\varphi `$, ensure that this BI is always violated for the range of parameters chosen irrespective of the efficiency factor of the detector, which can in any case be accounted for easily by the introduction of appropriate scaling factors in the expressions of the various probabilities appearing in the Bell sum.
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# Experimental evidences of Luttinger liquid behavior in the crossed multi-wall carbon nanotubes
\[
## Abstract
Luttinger liquid behavior was observed in a crossed junction formed with two metallic multi-wall carbon nanotubes whose differential conductance vanished with the power of bias voltage and temperature. With applying constant voltage or current to one of the two carbon nanotubes in a crossed geometry, the electrical transport properties of the other carbon nanotube were affected significantly, implying there exists strong correlation between the carbon nanotubes. Such characteristic features are in good agreement with the theoretical predictions for the crossed two Luttinger liquids.
\]
The non-Fermi liquid behavior of low dimensional systems is one of the most challenging subjects in recent theoretical and experimental studies . In one dimension, electron-electron interaction is known to invalidate the Fermi liquid description of metal and gives a new physical state, Luttinger liquid (LL), which is characterized by the spin-charge separation and the absence of single-particle excitation at low energy, etc. Until now, many experiments, mostly with the compound semiconductors , have been performed regarding the LL behavior of low dimensional conductor. The experimental results seem to support non-Fermi liquid behavior of one-dimensional conductor but conclusive evidence of the LL behavior is still to come.
Carbon nanotube (CNT) is known to be an ideal system to test the LL behavior of one-dimensional system . Previous experiment on single-wall CNT by Bockrath et al. has shown that the tunneling density of states vanishes both with temperature and bias voltage in the power-law functional form. However, the electrical transport properties of single-wall CNT in their measurements were dominated by the Coulomb blockade effect at low bias region, which might mask the LL behavior. Recently, Komnik and Egger have proposed an elegant way to verify the LL behavior of one-dimensional conductor. They have shown that if two LLs are contacted in a point-like manner, the electrical transport through one LL is perturbed significantly by the bias voltage applied to the other LL. A previous experiment on the crossed single-wall CNT did not show the LL behavior theoretically predicted , probably because of the large Coulomb blockade effect in each CNT. To test these theoretical predictions and detect possible deviation from the Fermi liquid theory, we have fabricated a cross junction formed with two metallic multi-wall CNTs. With single-wall CNTs, it is not easy to form a cross junction with enough electrostatic interaction between two CNTs which is essential to observe LL behavior in the crossed geometry . Due to the relative easiness of forming low-ohmic cross junction and weak Coulomb blockade effect compared to the single-wall counterpart, the multi-wall CNT has great advantages to test the LL behavior of CNT. Further, it is important to note that multi-wall CNT, as well as single-wall CNT, is predicted to show LL behavior in the limit of small number of conducting channels .
In this Letter, we report experimental evidences of LL behavior of multi-wall CNTs in the crossed geometry. Each CNT in the sample showed vanishing tunneling density of states with bias voltage in the power-law functional form, which is one of the signatures of LL behavior. The differential conductance curves at different temperatures are collapsed well into a single scaling curve. We have also measured the differential conductance of one CNT with applying constant voltage to the other CNT and observed that the differential conductance of one CNT is strongly influenced by the voltage applied to the other CNT . Furthermore, it is found that the off-diagonal differential resistance curve in the cross junction shows highly non-linear behavior along with negative differential resistance, which is another indication of the existence of strong correlation between two CNTs. All the experimental results support the LL behavior of multi-wall CNTs.
The multi-wall CNT used in this measurement was synthesized by arc discharge method. To select single CNT we have dispersed ultrasonically the CNT in chloroform for about half an hour and then dropped a droplet of the dispersed solution on the Si substrate with 500 nm-thick thermally-grown $`\text{SiO}_2`$ layer. The multi-wall CNTs in the crossed form were searched by scanning electron microscope (SEM). The patterns for electrical leads were generated using e-beam lithography technique onto the selected CNTs and then 20 nm of Ti and 50 nm of Au were deposited successively on the contact area by thermal evaporation. Shown in Fig. 1 is the SEM photograph of the crossed CNTs with the electric leads labeled. The atomic force microscope study has shown that the diameter of the CNT was in the range of 25 - 30 nm. In order to form low-ohmic contacts between the CNT and the Ti/Au electrodes, we have performed rapid thermal annealing at 800 C for 30 sec . The contact resistances are in the range of 5 k$`\mathrm{\Omega }`$ \- 18 k$`\mathrm{\Omega }`$ at room temperature and become 10 k$`\mathrm{\Omega }`$ \- 60 k$`\mathrm{\Omega }`$ at 4.2 K. The cross junction has junction resistance of 5.4 k$`\mathrm{\Omega }`$ at room temperature and 16.8 k$`\mathrm{\Omega }`$ at 4.2 K. The four-terminal resistance of each CNT increases monotonically with lowering temperature and depends sensitively on the bias current level, implying non-ohmic current-voltage characteristics of the CNT .
We have measured the current-voltage ($`I`$-$`V`$) characteristics both in two- and four-terminal measurement configurations. The differential conductance-voltage curve ($`dI/dV`$-$`V`$) was then obtained by numerically differentiating the $`I`$-$`V`$ characteristics. Insets of Fig. 2 display the four-terminal $`dI/dV`$-$`V`$ curves of each CNT as a function of temperature. Subtracting the contact resistance, both two- and four-terminal measurements give nearly identical $`dI/dV`$-$`V`$ curves. As shown in the insets of Fig. 2, the differential conductance, which is proportional to the density of states near the Fermi level, vanishes at low bias as the temperature is lowered. We have found that the density of states vanishes with the bias voltage in the power-law functional form, $`GV^\alpha `$, with $`\alpha =0.3`$ for the nanotube horizontally placed in Fig. 1 (from now on we call it CNT-1) and $`\alpha =0.9`$ for the nanotube vertically placed (CNT-2). By using the relation between the exponent and the Luttinger parameter for an end-contacted LL ,
$$\alpha =\left(\overline{g}^11\right)/4$$
(1)
where $`\overline{g}`$ is the effective Luttinger parameter for crossed LL given by $`\overline{g}=2g`$, we get $`g=0.23`$ for CNT-1 and $`g=0.11`$ for CNT-2.
For a LL, the temperature dependence of low-bias conductance is also expected to show the power-law functional form, $`G`$($`V0)T^\alpha `$. One way to exhibit this behavior is to plot the scaled differential conductance, $`T^\alpha dI/dV`$, as a function of the scaled voltage, $`eV/k_BT`$ . As shown in the main panels of Fig. 2, the scaled differential conductance curves for different temperatures fall well into a single scaling curve, except for high bias voltage where the differential conductance becomes saturated. Two CNTs showed similar scaling behavior but with different exponent $`\alpha `$. Such scaling behavior of the differential conductance curves is an indication of probable LL behavior of the two CNTs in the sample. The exponent $`\alpha `$ depends on the sample geometry.
The LL behavior of CNTs was further verified by the two-terminal differential conductance curves of CNT in a crossed geometry proposed by Komnik and Egger . We have measured the differential conductance of the CNT-1, $`dI_{34}/dV_{34}`$, with applying bias voltage $`V_{56}`$ to the CNT-2. Fig. 3 (a) shows the measured differential conductance curves with varying bias voltage $`V_{56}`$ from -24 mV to +21.6 mV with the increment of 2.4 mV. For $`V_{56}`$ close to zero, typical differential conductance curves with vanishing $`dI/dV`$ in a power of $`V`$ were shown. With increasing the magnitude of $`V_{56}`$, the zero-bias conductance increases rapidly and for $`|V_{56}|>12\text{mV}`$, the zero-bias differential conductance $`dI_{34}/dV_{34}`$ ($`V_{34}`$ = 0) switches from a dip to a peak. In addition the differential conductance curve begins to exhibit two separated dips at finite voltages. This characteristic feature agrees very well with the theoretical prediction and is one of experimental evidences of the strong correlation between the two CNTs. Such a strong dependence of the differential conductance $`dI_{34}/dV_{34}`$ on the bias $`V_{56}`$ may not be easily understood within the framework of the non-interacting electron picture . The dip separation increases monotonically with $`|V_{56}|`$. Following Komnik and Egger , the conductance dip should appear at $`|V_{56}|=|V_{34}|`$ for two identical LLs with $`g=1/4`$, a special point where an exact $`I`$-$`V`$ curve was calculated. We have plotted the dip position $`V_{34}`$ as a function of $`V_{56}`$ in Fig. 3 (b). As expected, the peak position $`V_{34}`$ increases or decreases linearly with $`V_{56}`$. Best fit gives the magnitude of the linear slopes close to 0.21, about 5 times smaller than that of the predicted ones for the two identical LLs. In our case, two CNTs have different Luttinger parameters, $`g`$ = 0.23 for CNT-1 and $`g`$ = 0.11 for CNT-2, respectively, which might be the origin of the discrepancy.
We have also measured the current-voltage characteristics with applying current to one CNT and measuring voltage drop on the other CNT in the crossed geometry. Fig. 4 shows the voltage-current characteristics, $`V_{34}`$-$`I_{56}`$, and the off-diagonal differential resistance (ODR)-current curve, $`dV_{56}/dI_{34}`$-$`I_{34}`$, measured at $`T`$ = 50 mK. A noticeable feature is the highly-non-linear behavior of the ODR together with the existence of the negative differential resistance (NDR) at low bias region. This can be interpreted as another indication of the existence of strong correlation between two CNTs which are considered to be LLs. The interchange of the voltage and the current leads gives almost identical voltage-current characteristics. The NDR and asymmetry in the ODR can be understood in a simple argument. The ODR can be written by
$$R_{ab}=\frac{G_{ab}}{G_{aa}^2G_{ab}^2},$$
(2)
where $`G_{aa}`$ and $`G_{ab}`$ are the diagonal and the off-diagonal conductances given in Ref.. Here the indices $`a,b`$ denote the current and the voltage probes in a given measurement configuration, respectively. It can be shown that for $`g<1/2`$, $`G_{ab}`$ is non-zero and bound by $`G_0/2<G_{ab}<G_0/2`$, where $`G_0`$ is the unit conductance of the system (in our case $`G_0=4e^2/h`$), and can be asymmetric under the bias reversal. Then $`R_{ab}`$ can be negative and also be asymmetric under the bias reversal. These features are well shown in Fig. 4.
In summary, we have investigated electrical transport properties of the crossed multi-wall CNTs. Each nanotube showed the tunneling density of states vanishing with the power of bias and temperature at low energy limit, which is an evidence of the LL behavior. The differential conductance curves of one CNT were disturbed significantly by applying bias voltage to the other one in a crossed geometry. This characteristic feature is an indication that strong correlation exists between the two crossed CNTs. With increasing bias voltage, two dips appear in the differential conductance curves and the dip separation increases linearly with the bias voltage, which is consistent with the theoretical prediction for two crossed LLs. Furthermore, the off-diagonal differential resistance exhibited highly non-linear behavior with negative differential resistance. We conclude that all these experimental results support the LL behavior of multi-wall CNTs.
We thank G. Cuniberti, H.-W. Lee, and H. S. Sim for helpful discussions and comments. This work was supported by the MOST through Nano Structure Project, Korea Research Foundation Grant (KRF-99-015-DP0128), BK21 project, and also by KRISS (Project No. 00-0502-100).
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# Brane cosmological perturbations
## I Introduction
Recently has emerged the fascinating idea that extra dimensions (beyond our familiar time and three spatial dimensions) could be large today or at energy scales much below what was thought before . The reason why these extra dimensions could remain hidden is that matter fields would be confined to a three-dimenional brane (our visible ‘space’) whereas gravitational fields would live everywhere and thus also in the extra-dimensions.
Even non-compact extra-dimensions can be envisaged , like in the five-dimensional model of Randall and Sundrum , where the bulk is endowed with a negative cosmological constant. Remarkably, in this model, the gravity felt by observers on the brane will behave, at least at first approximation, like usual four-dimensional gravity, because of the existence of zero-mode gravitons effectively confined in the brane.
When one considers a brane universe in the context of cosmology, it appears that one does not recover the standard Friedmann equations easily . The reason is that the matter content of the brane enters quadratically in the equations governing the dynamics of the brane geometry. However, when one generalizes the idea of Randall and Sundrum to cosmology, i.e. when one assumes the existence of a negative cosmological constant in the bulk, adjusted so as to compensate the tension of the brane, one ends up with a cosmological evolution that is conventional at sufficiently low energies ( and ).
The question now is whether this conventional behaviour will still be recovered in brane cosmology, when one relaxes the assumption of homogeneity and isotropy in the brane, in other words when one considers inhomogeneous brane cosmology. Up to now, almost all works (see e.g. ) on the cosmological aspects of the brane scenario have dealt with a homogeneous and isotropic brane universe. However, it is well known that a lot of information is contained in the cosmological perturbations of our Universe, and it is therefore of the uppermost importance to analyse the behaviour of cosmological perturbations in the context of brane cosmology. It would be crucial to test whether brane cosmology is compatible with current observations of cosmological perturbations, either the anisotropies of the Cosmic Microwave Background or the large scale structure data. It would also be interesting to check if the usual mechanism of generation of cosmological perturbations via amplification of quantum fluctuations during an inflationary phase would still be valid in brane cosmology (see ).
Another motivation to study the perturbations in the cosmological context is the question of the stabilisation of the radion and its implication on the cosmological evolution in the brane . As shown by in a non-cosmological context, a perturbative approach may be necessary for a full understanding of the radion.
The purpose of the present paper is to develop a formalism that can describe the evolution of the cosmological perturbations in the context of brane cosmology. In order to do so, we will proceed in several steps. First, one must compute the perturbations of the bulk Einstein’s equations, which relate the perturbations of the metric to the perturbations of the bulk matter, if there is any. For this first step, we have computed the perturbed Einstein’s equations in the more general case of a warped spacetime with any number of dimensions. A similar calculation has been done recently in the case of maximally symmetric spacetimes.
The results valid for any dimension can be applied to the model we wish to focus on: a five-dimensional spacetime. It is then useful to distinguish several types of cosmological modes, following the familiar decomposition of the standard theory of cosmological perturbations (see e.g. or ).
At this stage, however, the perturbations of the matter in the brane have not yet been taken into account. They play a rôle only in the junctions conditions: the jump in the derivatives (with respect to the fifth dimension) of the bulk metric perturbations is indeed governed by the matter perturbations in the brane. This is in this very unusual way that matter perturbations in our apparent Universe, i.e. the three-brane, are connected with the geometry perturbations in our Universe, i.e. simply the particular value on the brane of the bulk metric perturbations. It is nevertheless interesting to notice that the present problem has some analogy with the question of the interaction between a domain wall and gravitational waves in a four-dimensional spacetime .
The plan of the paper will be the following. In Section 2, we obtain the perturbations of Einstein’s equation in the bulk, in the case of a general warped spacetime. In Section 3, we consider the five-dimensional model and recall what is known about the homogeneous and isotropic solutions in an empty bulk or with a negative cosmological constant. Section 4 is devoted to the bulk perturbations of the five-dimensional spacetime, whereas section 5 deals with the junction conditions that take into account the matter in the brane. Finally, we conclude in the last section.
## II Bulk perturbations in general warped spacetime
Although we will be ultimately interested in the perturbations of a five-dimensional spacetime, it is instructive to begin with a more general situation and to compute the perturbations of the Ricci tensor for a D-dimensional spacetime that can be considered as a warped product of a p-dimensional spacetime with a d-dimensional space. The metric has thus the particular form
$$\overline{g}_{AB}dx^Adx^B=\stackrel{~}{\gamma }_{ab}dx^adx^b+a^2\{x^c\}\gamma _{ij}dx^idx^j,$$
(1)
where $`A,B,\mathrm{}`$ denote global spacetime indices, $`a,b,\mathrm{}`$ indices of the p-dimensional spacetime with metric $`\stackrel{~}{\gamma }_{ab}`$ , and $`i,j,\mathrm{}`$ indices of the d-dimensional space with metric $`\gamma _{ij}`$.
Let us denote $`D_A`$ the global covariant derivative associated with the metric $`\overline{g}_{AB}`$, whereas $`\stackrel{~}{}_a`$ will stand for the covariant derivative associated with the metric $`\stackrel{~}{\gamma }_{ab}`$ and $`_i`$ for the covariant derivative associated with the metric $`\gamma _{ij}`$. The nonvanishing mixed Christoffel symbols (those with indices of both types $`a`$ and $`i`$) are
$$\mathrm{\Gamma }_{aj}^i=\frac{_aa}{a}\gamma _j^i,\mathrm{\Gamma }_{ij}^a=(a^aa)\gamma _{ij}.$$
(2)
All tensors components will be decomposed in several sets, depending on the number of p-indices and of d-indices. And one can then decompose the action of the global spacetime covariant derivative $`D_A`$ into the action of the covariant derivative $`\stackrel{~}{}_a`$ and that of the covariant derivative $`_i`$. For illustration, on a vector $`V^A=(V^a,V^i)`$, the action of the covariant derivative $`D_A`$ gives
$`D_aV^b`$ $`=`$ $`\stackrel{~}{}_aV^b,`$ (3)
$`D_aV^i`$ $`=`$ $`{\displaystyle \frac{1}{a}}\stackrel{~}{}_a(aV^i),`$ (4)
$`D_iV^a`$ $`=`$ $`_iV^a(a^aa)\gamma _{ij}V^j,`$ (5)
$`D_iV^j`$ $`=`$ $`_iV^j+{\displaystyle \frac{_ba}{a}}\gamma _j^iV^b.`$ (6)
This type of formulas can be generalized to any type of tensor by using the mixed Christoffel symbols given above.
Let us now consider a linear perturbation of the spacetime metric, so that the total metric reads
$$ds^2=(\overline{g}_{AB}+h_{AB})dx^Adx^Bg_{AB}dx^Adx^B.$$
(7)
Our purpose will now be to compute the linear perturbation of the Ricci tensor, $`\delta R_{AB}`$. Quite generally, the expression of $`\delta R_{AB}`$ in terms of the metric perturbations is found to be of the following form (see e.g. )
$$\delta R_{AB}=\frac{1}{2}D_AD_Bh\frac{1}{2}D^CD_Ch_{AB}+D^CD_{(B}h_{A)C},$$
(8)
where $`h`$ denotes the trace of the perturbation $`h_{AB}`$, namely
$$h\overline{g}^{AB}h_{AB}.$$
(9)
After some tedious but straightforward calculations, it is possible to express the various components of the linearized Ricci tensor more specifically in terms of the components $`h_{ab}`$, $`h_{ai}`$ and $`h_{ij}`$, and of the covariant derivatives $`\stackrel{~}{}_a`$ and $`_i`$. One finds
$`\delta R_{ab}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\stackrel{~}{}_a\stackrel{~}{}_bh{\displaystyle \frac{1}{2}}\stackrel{~}{}^c\stackrel{~}{}_ch_{ab}+\stackrel{~}{}^c\stackrel{~}{}_{(a}h_{b)c}{\displaystyle \frac{1}{2}}a^2^2h_{ab}+a^2^k\stackrel{~}{}_{(a}h_{b)k}`$ (11)
$`+{\displaystyle \frac{d}{2}}{\displaystyle \frac{\stackrel{~}{}^ca}{a}}\left(\stackrel{~}{}_ah_{bc}+\stackrel{~}{}_bh_{ac}\stackrel{~}{}_ch_{ab}\right)a^1\stackrel{~}{}_{(a}a\stackrel{~}{}_{b)}\widehat{h},`$
$`\delta R_{ai}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\stackrel{~}{}_a_ih{\displaystyle \frac{1}{2}}\stackrel{~}{}^b\stackrel{~}{}_bh_{ai}+{\displaystyle \frac{1}{2}}\stackrel{~}{}^b\stackrel{~}{}_ah_{bi}+{\displaystyle \frac{1}{2}}_i\stackrel{~}{}^bh_{ab}+{\displaystyle \frac{1}{2}}a^2\stackrel{~}{}_a^kh_{ik}{\displaystyle \frac{1}{2}}a^2^k_kh_{ai}`$ (14)
$`+{\displaystyle \frac{1}{2}}a^2^k_ih_{ak}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\stackrel{~}{}_aa}{a}}_i(h\widehat{h}){\displaystyle \frac{\stackrel{~}{}_aa}{a}}\stackrel{~}{}^bh_{bi}+{\displaystyle \frac{d}{2}}{\displaystyle \frac{\stackrel{~}{}^ba}{a}}\stackrel{~}{}_ah_{bi}+\left(1{\displaystyle \frac{d}{2}}\right){\displaystyle \frac{\stackrel{~}{}^ba}{a}}\stackrel{~}{}_bh_{ai}`$
$`\left(1{\displaystyle \frac{d}{2}}\right){\displaystyle \frac{\stackrel{~}{}^ba}{a}}_ih_{ab}{\displaystyle \frac{\stackrel{~}{}_a}{a^3}}^kh_{ik}(d1){\displaystyle \frac{\stackrel{~}{}^ba\stackrel{~}{}_aa}{a^2}}h_{bi}{\displaystyle \frac{\stackrel{~}{}^b\stackrel{~}{}_aa}{a}}h_{bi}`$
$`\delta R_{ij}`$ $`=`$ $`{\displaystyle \frac{1}{2}}_i_jh{\displaystyle \frac{1}{2}}\left(a\stackrel{~}{}^aa\right)\gamma _{ij}\stackrel{~}{}_ah{\displaystyle \frac{1}{2}}a^2^k_kh_{ij}+{\displaystyle \frac{1}{2}}a^2^k_ih_{jk}+{\displaystyle \frac{1}{2}}a^2^k_jh_{ik}`$ (18)
$`{\displaystyle \frac{1}{2}}\stackrel{~}{}^a\stackrel{~}{}_ah_{ij}+{\displaystyle \frac{1}{2}}\stackrel{~}{}^a_ih_{aj}+{\displaystyle \frac{1}{2}}\stackrel{~}{}^a_jh_{ai}+{\displaystyle \frac{\stackrel{~}{}^aa}{a}}\gamma _{ij}^kh_{ak}+\left(a\stackrel{~}{}^ba\right)\gamma _{ij}\stackrel{~}{}^ah_{ab}`$
$`+\left(2{\displaystyle \frac{d}{2}}\right){\displaystyle \frac{\stackrel{~}{}^aa}{a}}\stackrel{~}{}_ah_{ij}+\left({\displaystyle \frac{d}{2}}1\right){\displaystyle \frac{\stackrel{~}{}^ba}{a}}_ih_{bj}+\left({\displaystyle \frac{d}{2}}1\right){\displaystyle \frac{\stackrel{~}{}^ba}{a}}_jh_{bi}`$
$`+(d1)(\stackrel{~}{}^aa\stackrel{~}{}^ba)\gamma _{ij}h_{ab}+a(\stackrel{~}{}^a\stackrel{~}{}^ba)\gamma _{ij}h_{ab}2{\displaystyle \frac{\stackrel{~}{}^ba\stackrel{~}{}_ba}{a^2}}h_{ij},`$
with $`\widehat{h}a^2\gamma ^{ij}h_{ij}`$.
Gauge transformations.
It is important to notice that perturbations which differ quantitatively can in fact describe the same geometry simply because they correspond to different systems of coordinates. Gauge transformations, corresponding to infinitesimal changes of coordinates
$$x^Ax^A+\xi ^A,$$
(19)
induce the following transformations for the metric perturbations,
$$h_{AB}h_{AB}+D_A\xi _B+D_B\xi _A.$$
(20)
This general expression gives for the three types of metric components or, decomposing into the two subsystems of coordinates,
$`h_{ab}`$ $``$ $`h_{ab}+\stackrel{~}{\gamma }_{bc}\stackrel{~}{}_a\xi ^c+\stackrel{~}{\gamma }_{ac}\stackrel{~}{}_b\xi ^c,`$ (21)
$`h_{ij}`$ $``$ $`h_{ij}+a^2\left(\gamma _{jk}_i\xi ^k+\gamma _{ik}_j\xi ^k\right)+2a(_aa)\gamma _{ij}\xi ^a,`$ (22)
$`h_{ia}`$ $``$ $`h_{ia}+\stackrel{~}{\gamma }_{ab}_i\xi ^b+a^2\gamma _{ij}\stackrel{~}{}_a\xi ^j.`$ (23)
## III Five-dimensional background spacetime
In this section, we recall briefly the cosmological solutions that have been found in the case of five-dimensional spacetimes, with a metric of the form
$$ds^2=n^2(\tau ,y)d\tau ^2+a^2(\tau ,y)\gamma _{ij}dx^idx^j+b^2(\tau ,y)dy^2,$$
(24)
where the spatial three-surfaces, defined by $`\tau `$ and $`y`$ constant, are homogeneous and isotropic and $`\gamma _{ij}`$ is a maximally symmetric 3-dimensional metric ($`k=1,0,1`$ will parametrize the spatial curvature).
The five-dimensional Einstein equations take the usual form
$$G_{AB}R_{AB}\frac{1}{2}Rg_{AB}=\kappa ^2T_{AB},$$
(25)
where $`T_{AB}`$ is the five-dimensional energy-momentum tensor.
With the above metric, the non-vanishing components of the Einstein tensor $`G_{AB}`$ are found to be
$`G_{00}`$ $`=`$ $`3\left\{{\displaystyle \frac{\dot{a}}{a}}\left({\displaystyle \frac{\dot{a}}{a}}+{\displaystyle \frac{\dot{b}}{b}}\right){\displaystyle \frac{n^2}{b^2}}\left({\displaystyle \frac{a^{\prime \prime }}{a}}+{\displaystyle \frac{a^{}}{a}}\left({\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{b^{}}{b}}\right)\right)+k{\displaystyle \frac{n^2}{a^2}}\right\},`$ (26)
$`G_{ij}`$ $`=`$ $`{\displaystyle \frac{a^2}{b^2}}\gamma _{ij}\left\{{\displaystyle \frac{a^{}}{a}}\left({\displaystyle \frac{a^{}}{a}}+2{\displaystyle \frac{n^{}}{n}}\right){\displaystyle \frac{b^{}}{b}}\left({\displaystyle \frac{n^{}}{n}}+2{\displaystyle \frac{a^{}}{a}}\right)+2{\displaystyle \frac{a^{\prime \prime }}{a}}+{\displaystyle \frac{n^{\prime \prime }}{n}}\right\}`$ (28)
$`+{\displaystyle \frac{a^2}{n^2}}\gamma _{ij}\left\{{\displaystyle \frac{\dot{a}}{a}}\left({\displaystyle \frac{\dot{a}}{a}}+2{\displaystyle \frac{\dot{n}}{n}}\right)2{\displaystyle \frac{\ddot{a}}{a}}+{\displaystyle \frac{\dot{b}}{b}}\left(2{\displaystyle \frac{\dot{a}}{a}}+{\displaystyle \frac{\dot{n}}{n}}\right){\displaystyle \frac{\ddot{b}}{b}}\right\}k\gamma _{ij},`$
$`G_{05}`$ $`=`$ $`3\left({\displaystyle \frac{n^{}}{n}}{\displaystyle \frac{\dot{a}}{a}}+{\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{\dot{b}}{b}}{\displaystyle \frac{\dot{a}^{}}{a}}\right),`$ (29)
$`G_{55}`$ $`=`$ $`3\left\{{\displaystyle \frac{a^{}}{a}}\left({\displaystyle \frac{a^{}}{a}}+{\displaystyle \frac{n^{}}{n}}\right){\displaystyle \frac{b^2}{n^2}}\left({\displaystyle \frac{\dot{a}}{a}}\left({\displaystyle \frac{\dot{a}}{a}}{\displaystyle \frac{\dot{n}}{n}}\right)+{\displaystyle \frac{\ddot{a}}{a}}\right)k{\displaystyle \frac{b^2}{a^2}}\right\},`$ (30)
where a prime stands for a derivative with respect to $`y`$, and a dot for a derivative with respect to $`\tau `$. The stress-energy-momentum tensor can be decomposed into two parts,
$$T_B^A=\stackrel{ˇ}{T}_B^A|_{_{\mathrm{bulk}}}+T_B^A|_{_{\mathrm{brane}}},$$
(31)
where $`\stackrel{ˇ}{T}_B^A|_{_{\mathrm{bulk}}}`$ is the energy momentum tensor of the bulk matter and $`T_B^A|_{_{\mathrm{brane}}}`$ corresponds to the matter content in the brane $`(y=0)`$. Since we consider in this section only strictly homogeneous and isotropic geometries inside the brane, the latter will be necessary of the form
$$T_B^A|_{_{\mathrm{brane}}}=\frac{\delta (y)}{b}\text{diag}(\rho ,p,p,p,0),$$
(32)
where the energy density $`\rho `$ and pressure $`p`$ are functions only of time.
In the case of a general stress-energy tensor for the matter on the brane, and assuming a time-independent metric along the fifth dimension, explicit solutions for the whole metric have been found in two simple cases: the case where the bulk is empty and the case where the bulk contains a cosmological constant , i.e.
$$\stackrel{ˇ}{T}_B^A|_{_{\mathrm{bulk}}}=\text{diag}(\rho _B,\rho _B,\rho _B,\rho _B,\rho _B),$$
(33)
with $`\rho _B=const`$.
In the case of a negative bulk cosmological constant, the most realistic, the explicit expressions for the metric components (with $`b=1`$ and $`n(t,y=0)=1`$) are given by
$`a(t,y)`$ $`=`$ $`\{{\displaystyle \frac{1}{2}}(1+{\displaystyle \frac{\kappa ^2\rho ^2}{6\rho _B}})a_0^2+{\displaystyle \frac{3𝒞}{\kappa ^2\rho _Ba_0^2}}+[{\displaystyle \frac{1}{2}}(1{\displaystyle \frac{\kappa ^2\rho ^2}{6\rho _B}})a_0^2{\displaystyle \frac{3𝒞}{\kappa ^2\rho _Ba_0^2}}]\mathrm{cosh}(\mu y)`$ (35)
$`{\displaystyle \frac{\kappa \rho }{\sqrt{6\rho _B}}}a_0^2\mathrm{sinh}(\mu |y|)\}^{1/2},`$
with $`\mu =\sqrt{\frac{2\kappa ^2}{3}\rho _B},`$ and
$$n(t,y)=\frac{\dot{a}(t,y)}{\dot{a}_0(t)},$$
(36)
where $`a_0`$ is the scale factor in the brane (i.e. in $`y=0`$). The two functions $`a_0(t)`$ and $`\rho (t)`$, which appear in the above solution, are obtained by solving
$$\frac{\dot{a}_0^2}{a_0^2}=\frac{\kappa ^2}{6}\rho _B+\frac{\kappa ^4}{36}\rho ^2+\frac{𝒞}{a_0^4}\frac{k}{a_0^2},$$
(37)
where $`𝒞`$ is a constant (of integration), and
$$\dot{\rho }+3\frac{\dot{a}_0}{a_0}(\rho +p)=0.$$
(38)
The last equation corresponds to the ordinary energy conservation law. But (37), analogous to the first Friedman equation, does not give the usual cosmological evolution because of the quadratic term $`\rho ^2`$. It is however possible to recover the standard cosmological evolution, at least at sufficiently late times, if one decomposes $`\rho `$ into a tension $`\sigma `$ and an energy density of ordinary matter living in the brane and if one assumes a negative cosmological constant in the bulk that will compensate the $`\sigma ^2`$ term (see ).
## IV Perturbations of a five-dimensional spacetime
After having described our reference spacetime in the previous section, we will now study linear perturbations about it. For this task, the calculations of Section 2, specialized to the case of five-dimensional spacetime, are going to be very useful. As it is clear from the previous section, the metric $`\stackrel{~}{\gamma }_{ab}`$ will be defined by
$$\stackrel{~}{\gamma }_{ab}dx^adx^b=n(t,y)^2dt^2+b(t,y)^2dy^2.$$
(39)
The second metric $`\gamma _{ij}`$ is simply the metric covering the three ordinary spatial dimensions. In a cosmological context, there are three choices for $`\gamma _{ij}`$ depending on the spatial curvature of spacetime. For simplicity, it will be assumed that our Universe is spatially flat, which means that
$$\gamma _{ij}=\delta _{ij}.$$
(40)
From now on, we also choose to work in a Gaussian normal (GN) system of coordinates adapted to the three-brane, in which the brane is localized in $`y=0`$ and the metric has the form
$$ds^2=g_{\mu \nu }dx^\mu dx^\nu +dy^2,$$
(41)
where the greek indices refer to ordinary spacetime coordinates, i.e. the time coordinate $`\tau `$ and the three ordinary spatial dimensions $`x^i`$. For the background, this choice of gauge corresponds simply to set
$$b(t,y)=1.$$
(42)
In the case of the perturbed spacetime, described by the metric (7), the choice of a GN system of coordinates will impose
$$h_{55}=h_{5\mu }=0.$$
(43)
Taking this into account, the linearized metric (7) can now be written, quite generally, in the form
$$ds^2=n^2(1+2A)dt^2+2B_idx^idt+a^2\left(\delta _{ij}+\widehat{h}_{ij}\right)dx^idx^j+dy^2.$$
(44)
Following Bardeen (see also ), the linearized quantities specified above can be decomposed further, into so-called scalar, vector and tensor quantities, according to the expressions,
$$B_i=_iB+\overline{B}_i,$$
(45)
where $`\overline{B}_i`$ satisfies $`_i\overline{B}^i=0`$, and
$$\widehat{h}_{ij}=2C\delta _{ij}+2_i_jE+2_{(i}E_{j)}+E_{ij},$$
(46)
where $`E_{ij}`$ is transverse traceless, i.e. $`_iE^{ij}=0`$, $`E_i^i=0`$. The quantities $`A`$, $`C`$, $`B`$ and $`E`$ are usually refered to as ‘scalar’ perturbations, $`\overline{B}_i`$ and $`E_i`$ as ‘vector’ perturbations, and $`E_{ij}`$ as ‘tensor’ perturbations.
### A Gauge transformation
As explained earlier, the perturbations defined above can be quantitatively different but describe the same geometry if they are written in different coordinate systems. In order to distinguish gauge effects and physical degrees of freedom, it is useful to write down the effect of a coordinate change (or gauge transformation) on all the metric perturbations defined above. These transformations follow directly from the general expressions given in (23).
We will parametrize the infinitesimal coordinate transformation by the vector $`\xi ^A=(\xi ^0,\xi ^i,\xi ^5)`$. Let us first consider the components $`h_{55}`$, $`h_{05}`$ and $`h_{i5}`$, which vanish in a GN coordinate system. They transform according to the laws,
$`h_{55}`$ $``$ $`h_{55}+2\xi _{}^{5}{}_{}{}^{},`$ (47)
$`h_{05}`$ $``$ $`h_{05}+\dot{\xi }^5n^2\xi _{}^{0}{}_{}{}^{},`$ (48)
$`h_{i5}`$ $``$ $`h_{i5}+_i\xi ^5+a^2\delta _{ij}\xi _{}^{j}{}_{}{}^{}.`$ (49)
In order to bring any coordinate system into a GN coordinate system, it is clear from the above relations that once one has used $`\xi ^5`$ to adjust the position of the brane at $`y=0`$, appropriate choices of $`\xi ^0`$ and of $`\xi ^i`$ will be required in order to make $`h_{05}`$ and $`h_{i5}`$ vanish. Note however that the GN gauge adapted to the brane is not completely fixed: there is some residual gauge freedom associated with parameters $`\xi ^0`$ and $`\xi ^i`$ that depend only on the four ordinary spacetime coordinates. This residual gauge freedom can be interpreted as possible redefinitions of the coodinates inside the brane worldsheet.
Let us now consider the transformations for the other components of the metric perturbations. To do that, it is convenient to decompose the spatial vector $`x^i`$ into
$$\xi ^i=^i\xi +\overline{\xi }^i$$
(50)
(where $`^i=\delta ^{ij}_j`$), such that $`\overline{\xi }^i`$ is transverse, i.e. $`_i\overline{\xi }^i=0`$. With this decomposition, one can see that the three scalar parameters $`\xi ^0`$, $`\xi `$ and $`\xi ^5`$ will induce transformations in the subset of scalar perturbations, whereas $`\overline{\xi }^i`$ will act in the ‘vector’ subspace. The tensor perturbations $`\overline{E}_{ij}`$ will be left untouched by the gauge transformations. Specializing the expressions (23) to our particular case here, one finds the following transformations:
#### 1 Scalar gauge transformations
$`A`$ $``$ $`A+\dot{\xi }^0+{\displaystyle \frac{\dot{n}}{n}}\xi ^0+{\displaystyle \frac{n^{}}{n}}\xi ^5,`$ (51)
$`B`$ $``$ $`Bn^2\xi ^0+a^2\dot{\xi },`$ (52)
$`C`$ $``$ $`C+{\displaystyle \frac{\dot{a}}{a}}\xi ^0+{\displaystyle \frac{a^{}}{a}}\xi ^5,`$ (53)
$`E`$ $``$ $`E+\xi ,`$ (54)
#### 2 Vector gauge transformations
$`\overline{B}_i`$ $``$ $`\overline{B}_i+a^2\dot{\overline{\xi }}_i,`$ (55)
$`E_i`$ $``$ $`E_i+\overline{\xi }_i.`$ (56)
In the above transformations, we have considered the most general gauge transformation. But if one considers only gauge transformations inside the subset of GN coordinate systems, which will be the case in the following, then $`\xi ^5`$ disappears from (54) and the parameters $`\xi ^0`$, $`\xi `$ and $`\overline{\xi }`$ cannot depend on the fifth coordinate.
Finally let us remark that one can use the remaining gauge freedom within the GN subset to impose some additional gauge conditions. For example, for scalar quantities, a choice which is familiar in the standard theory of cosmological perturbations is to impose
$$B=0,E=0,$$
(57)
which corresponds to the so-called longitudinal (or Newtonian) gauge. However, these conditions can be imposed only on one hypersurface, $`y=0`$ say, but not everywhere in the bulk, because $`B`$ and $`E`$ a priori depend on the coordinate $`y`$.
### B Perturbed Einstein equations in the bulk
Let us now compute the equations governing the perturbations in the bulk, ignoring for the moment the presence of the brane. The five dimensional Einstein’s equations (25), in the bulk, can be rewritten
$$R_{AB}=\kappa ^2\left(\stackrel{ˇ}{T}_{AB}\frac{1}{3}\stackrel{ˇ}{T}g_{AB}\right).$$
(58)
Therefore, the perturbed Einstein equations in the bulk have the form
$$\delta R_{AB}=\kappa ^2\left(\delta \stackrel{ˇ}{T}_{AB}\frac{1}{3}\stackrel{ˇ}{T}h_{AB}\frac{1}{3}g_{AB}\delta \stackrel{ˇ}{T}\right).$$
(59)
Since the background solutions have been found explicity, as recalled in Section 3, the cases of an empty bulk or of a cosmological constant are of particular interest, but of course, the present formalism applies to any bulk energy-momentum tensor. In the case of an empty bulk, the perturbed Einstein’s equations are simply
$$\delta R_{AB}=0,$$
(60)
whereas in the case of a bulk with a cosmological constant $`\mathrm{\Lambda }=\kappa ^2\rho _B`$, the perturbed Einstein’s equations read
$$\delta R_{AB}=\frac{2}{3}\mathrm{\Lambda }h_{AB}.$$
(61)
The remaining task consists in computing explicitly the components of the perturbed Ricci tensor. Using the expressions (11-18), one obtains the following expressions, which can conveniently be separated in scalar, vector and tensor parts.
#### 1 Scalar components
$`\delta R_{00}^S`$ $`=`$ $`n^2\left[A^{\prime \prime }+\left(3{\displaystyle \frac{a^{}}{a}}+2{\displaystyle \frac{n^{}}{n}}\right)A^{}+3{\displaystyle \frac{n^{}}{n}}C^{}+\left(6{\displaystyle \frac{a^{}n^{}}{an}}+2{\displaystyle \frac{n^{\prime \prime }}{n}}\right)A\right]`$ (64)
$`3\ddot{C}+3{\displaystyle \frac{\dot{a}}{a}}\dot{A}+\left(3{\displaystyle \frac{\dot{n}}{n}}6{\displaystyle \frac{\dot{a}}{a}}\right)\dot{C}+{\displaystyle \frac{n^2}{a^2}}\mathrm{\Delta }A`$
$`\mathrm{\Delta }\ddot{E}+\left({\displaystyle \frac{\dot{n}}{n}}2{\displaystyle \frac{\dot{a}}{a}}\right)\mathrm{\Delta }\dot{E}+nn^{}\mathrm{\Delta }E^{}+a^2\mathrm{\Delta }\dot{B}{\displaystyle \frac{\dot{n}}{a^2n}}\mathrm{\Delta }B,`$
$`\delta R_{0i}^S`$ $`=`$ $`_i\left\{2\dot{C}+2{\displaystyle \frac{\dot{a}}{a}}A{\displaystyle \frac{1}{2}}B^{\prime \prime }+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{n^{}}{n}}{\displaystyle \frac{a^{}}{a}}\right)B^{}+\left[{\displaystyle \frac{1}{n^2}}\left({\displaystyle \frac{\ddot{a}}{a}}{\displaystyle \frac{\dot{n}\dot{a}}{na}}+2{\displaystyle \frac{\dot{a}^2}{a^2}}\right)2{\displaystyle \frac{a^{}n^{}}{an}}\right]B\right\},`$ (65)
$`\delta R_{ij}^S`$ $`=`$ $`a^2\{(4{\displaystyle \frac{\dot{a}^2}{a^2n^2}}2{\displaystyle \frac{\dot{a}\dot{n}}{an^3}}+2{\displaystyle \frac{\ddot{a}}{an^2}})A+2({\displaystyle \frac{a^{}n^{}}{an}}+2{\displaystyle \frac{a^2}{a^2}}+{\displaystyle \frac{a^{\prime \prime }}{a}}+{\displaystyle \frac{\dot{a}\dot{n}}{an^3}}2{\displaystyle \frac{\dot{a}^2}{a^2n^2}}{\displaystyle \frac{\ddot{a}}{an^2}})C`$ (70)
$`+{\displaystyle \frac{a^{}}{a}}A^{}+\left(6{\displaystyle \frac{a^{}}{a}}+{\displaystyle \frac{n^{}}{n}}\right)C^{}+{\displaystyle \frac{\dot{a}}{an^2}}\dot{A}`$
$`+C^{\prime \prime }+({\displaystyle \frac{\dot{n}}{n^3}}6{\displaystyle \frac{\dot{a}}{an^2}})\dot{C}{\displaystyle \frac{1}{n^2}}\ddot{C}+{\displaystyle \frac{\dot{a}}{n^2a^3}}\mathrm{\Delta }B{\displaystyle \frac{\dot{a}}{n^2a}}\mathrm{\Delta }\dot{E}+{\displaystyle \frac{a^{}}{a}}\mathrm{\Delta }E^{}\}\delta _{ij}\mathrm{\Delta }C\delta _{ij}`$
$`+_i_j[AC+{\displaystyle \frac{a^2}{n^2}}\ddot{E}+({\displaystyle \frac{3a\dot{a}}{n^2}}{\displaystyle \frac{a^2\dot{n}}{n^3}})\dot{E}a^2E^{\prime \prime }(3aa^{}+{\displaystyle \frac{n^{}a^2}{n}})E^{}`$
$`+2({\displaystyle \frac{a\ddot{a}}{n^2}}+2{\displaystyle \frac{\dot{a}^2}{n^2}}{\displaystyle \frac{\dot{n}a\dot{a}}{n^3}}aa^{\prime \prime }2a^2{\displaystyle \frac{n^{}aa^{}}{n}})E{\displaystyle \frac{1}{n^2}}\dot{B}+({\displaystyle \frac{\dot{n}}{n^3}}{\displaystyle \frac{\dot{a}}{an^2}})B]`$
$`\delta R_{i5}^S`$ $`=`$ $`_i[A^{}2C^{}+({\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{n^{}}{n}})A{\displaystyle \frac{1}{2n^2}}\dot{B}^{}+{\displaystyle \frac{a^{}}{an^2}}\dot{B}+({\displaystyle \frac{\dot{n}}{2n^3}}{\displaystyle \frac{3\dot{a}}{2n^2a}})B^{}`$ (72)
$`+({\displaystyle \frac{\dot{a}^{}}{an^2}}+2{\displaystyle \frac{\dot{a}a^{}}{n^2a^2}}{\displaystyle \frac{\dot{n}a^{}}{n^3a}})B],`$
$`\delta R_{55}^S`$ $`=`$ $`6{\displaystyle \frac{a^{}}{a}}C^{}2{\displaystyle \frac{n^{}}{n}}A^{}A^{\prime \prime }3C^{\prime \prime }\mathrm{\Delta }E^{\prime \prime }2{\displaystyle \frac{a^{}}{a}}\mathrm{\Delta }E^{},`$ (73)
$`\delta R_{05}^S`$ $`=`$ $`3\left[{\displaystyle \frac{n^{}}{n}}\dot{C}\dot{C}^{}+{\displaystyle \frac{\dot{a}}{a}}A^{}{\displaystyle \frac{\dot{a}}{a}}C^{}{\displaystyle \frac{a^{}}{a}}\dot{C}\right]+{\displaystyle \frac{1}{2}}a^2\mathrm{\Delta }B^{}{\displaystyle \frac{n^{}}{a^2n}}\mathrm{\Delta }B`$ (75)
$`\mathrm{\Delta }\dot{E}^{}+\left({\displaystyle \frac{n^{}}{n}}{\displaystyle \frac{a^{}}{a}}\right)\mathrm{\Delta }\dot{E}{\displaystyle \frac{\dot{a}}{a}}\mathrm{\Delta }E^{}.`$
#### 2 Vector components
$`\delta R_{0i}^V`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Delta }\dot{E}_i{\displaystyle \frac{1}{2}}\overline{B}_i^{\prime \prime }+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{n^{}}{n}}{\displaystyle \frac{a^{}}{a}}\right)\overline{B}_i^{}+\left[{\displaystyle \frac{1}{n^2}}\left({\displaystyle \frac{\ddot{a}}{a}}{\displaystyle \frac{\dot{n}\dot{a}}{na}}+2{\displaystyle \frac{\dot{a}^2}{a^2}}\right)2{\displaystyle \frac{a^{}n^{}}{an}}\right]\overline{B}_i{\displaystyle \frac{1}{2a^2}}\mathrm{\Delta }\overline{B}_i,`$ (76)
$`\delta R_{i5}^V`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Delta }E_i^{}{\displaystyle \frac{1}{2n^2}}\dot{\overline{B}}_i^{}+{\displaystyle \frac{a^{}}{an^2}}\dot{\overline{B}}_i+\left({\displaystyle \frac{\dot{n}}{2n^3}}{\displaystyle \frac{3\dot{a}}{2n^2a}}\right)\overline{B}_i^{}+\left({\displaystyle \frac{\dot{a}^{}}{an^2}}+2{\displaystyle \frac{\dot{a}a^{}}{n^2a^2}}{\displaystyle \frac{\dot{n}a^{}}{n^3a}}\right)\overline{B}_i`$ (77)
$`\delta R_{ij}^V`$ $`=`$ $`{\displaystyle \frac{a^2}{n^2}}_{(i}\ddot{E}_{j)}+\left({\displaystyle \frac{3a\dot{a}}{n^2}}{\displaystyle \frac{a^2\dot{n}}{n^3}}\right)_{(i}\dot{E}_{j)}a^2_{(i}E_{j)}^{\prime \prime }\left(3aa^{}+{\displaystyle \frac{n^{}a^2}{n}}\right)_{(i}E_{j)}^{}`$ (80)
$`+2\left({\displaystyle \frac{a\ddot{a}}{n^2}}+2{\displaystyle \frac{\dot{a}^2}{n^2}}{\displaystyle \frac{\dot{n}a\dot{a}}{n^3}}aa^{\prime \prime }2a^2{\displaystyle \frac{n^{}aa^{}}{n}}\right)_{(i}E_{j)}`$
$`{\displaystyle \frac{1}{n^2}}_{(i}\dot{\overline{B}}_{j)}+\left({\displaystyle \frac{\dot{n}}{n^3}}{\displaystyle \frac{\dot{a}}{an^2}}\right)_{(i}\overline{B}_{j)}`$
#### 3 Tensor components
$`\delta R_{ij}^T`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Delta }E_{ij}+{\displaystyle \frac{a^2}{2n^2}}\ddot{E}_{ij}+\left({\displaystyle \frac{3a\dot{a}}{2n^2}}{\displaystyle \frac{a^2\dot{n}}{2n^3}}\right)\dot{E}_{ij}{\displaystyle \frac{1}{2}}a^2E_{ij}^{\prime \prime }\left({\displaystyle \frac{3}{2}}aa^{}+{\displaystyle \frac{n^{}a^2}{2n}}\right)E_{ij}^{}`$ (82)
$`+\left({\displaystyle \frac{a\ddot{a}}{n^2}}+2{\displaystyle \frac{\dot{a}^2}{n^2}}{\displaystyle \frac{\dot{n}a\dot{a}}{n^3}}aa^{\prime \prime }2a^2{\displaystyle \frac{n^{}aa^{}}{n}}\right)E_{ij}.`$
## V Junction conditions
In the previous section, we have considered only the perturbations in the bulk. It is now time to take into account the brane itself, and in particular the perturbations of the matter in the brane. The connection between the metric perturbations, living in the bulk, and the matter perturbations confined to the brane is made via the junction conditions .
The energy-momentum tensor describing matter content of the brane will be written in the form
$$T_\nu ^\mu |_{_{\mathrm{brane}}}=\delta (y)S_\nu ^\mu ,$$
(83)
and, in the background, it has the perfect fluid form
$$S_\nu ^\mu =(\rho +P)u^\mu u_\nu +Pg_\nu ^\mu .$$
(84)
The junction conditions are found to be given, in the five-dimensional case , by
$$[K_{\mu \nu }]=\kappa ^2\left(S_{\mu \nu }\frac{1}{3}Sg_{\mu \nu }\right),$$
(85)
where the brackets denote the jump across the brane, of the extrinsic curvature $`K_{\mu \nu }`$ and $`SS_\mu ^\mu `$. In a Gaussian normal coordinate system, the extrinsic curvature is given by the simple expression
$$K_{\mu \nu }=\frac{1}{2}_yg_{\mu \nu }.$$
(86)
In the present context, it is convenient to decompose the junction conditions (85) into a background part and a perturbed part. The use of the background junction conditions is necessary to find solutions of the Einstein equations when homogeneity and isotropy in the brane are assumed, such as those given in (35-36). The junction conditions for the background have been given in
$`{\displaystyle \frac{[a^{}]}{a_0b_0}}`$ $`=`$ $`{\displaystyle \frac{\kappa ^2}{3}}\rho ,`$ (87)
$`{\displaystyle \frac{[n^{}]}{n_0b_0}}`$ $`=`$ $`{\displaystyle \frac{\kappa ^2}{3}}\left(3p+2\rho \right),`$ (88)
where the subscript $`0`$ for $`a,b,n`$ means that these functions are taken in $`y=0`$.
Let us now consider the junction conditions for the perturbations, which can be written
$$[\delta K_{\mu \nu }]=\kappa ^2\left(\delta S_{\mu \nu }+\frac{1}{3}g_{\mu \nu }\delta S+\frac{1}{3}Sh_{\mu \nu }\right),$$
(89)
Using the expression (86) for the extrinsic curvature in the GN gauge as well as the $`yy`$ symmetry, one ends up with the following condition,
$$h_{\mu \nu |y=0^+}^{}=\kappa ^2\left(\delta S_{\mu \nu }+\frac{1}{3}g_{\mu \nu }\delta S+\frac{1}{3}Sh_{\mu \nu }\right).$$
(90)
It is then useful to decompose further these junction conditions by distinguishing space and time. Let us begin with the perturbations of the energy-momentum tensor in the brane. The perturbations for the unit four-velocity can be written
$$\delta u^\mu =\{n^1A,a^1v^i\},$$
(91)
where the expression for $`\delta u^0`$ follows automatically from the normalization condition satisfied by $`u^\mu `$. The perturbations of the energy-momentum tensor then have the following form:
$`\delta S_{00}`$ $`=`$ $`n^2\delta \rho +2\rho n^2A,`$ (92)
$`\delta S_{0i}`$ $`=`$ $`(\rho +P)nav_i\rho B_i,`$ (93)
$`\delta S_{ij}`$ $`=`$ $`a^2\delta P\delta _{ij}+Ph_{ij}+a^2\pi _{ij},`$ (94)
with $`v_i\delta _{ij}v^j`$ and where $`\pi _{ij}`$ is the anisotropic stress tensor. Once more it is possible to decompose the above expressions for the brane matter energy-momentum tensor into scalar, vector and tensor components, using
$$v_i=_iv+\overline{v}_i,$$
(95)
with $`_i\overline{v}^i=0`$ and
$$\pi _{ij}=\left(_i_j\frac{1}{3}\delta _{ij}\mathrm{\Delta }\right)\pi +2_{(i}\pi _{j)}+\overline{\pi }_{ij},$$
(96)
with as usual $`\pi _i`$ transverse and $`\overline{\pi }_{ij}`$ transverse traceless. Substituting the above expressions for the perturbed energy-momentum tensor and using the background junction conditions (87-88) yields the following conditions for the perturbations
$`A_{|y=0^+}^{}`$ $`=`$ $`{\displaystyle \frac{\kappa ^2}{6}}\left(2\delta \rho +3\delta P\right),`$ (97)
$`(n^2B_i)_{|y=0^+}^{}`$ $`=`$ $`\kappa ^2(\rho +P){\displaystyle \frac{a_0}{n_0}}v_i,`$ (98)
$`\widehat{h}_{ij|y=0^+}^{}`$ $`=`$ $`\kappa ^2\left({\displaystyle \frac{1}{3}}\delta \rho \delta _{ij}+\pi _{ij}\right).`$ (99)
These conditions can also be decomposed, in a straightforward manner, into scalar, vector and tensor junction conditions.
## VI Conclusion
In the present work, we have developed a formalism in order to deal with the evolution of cosmological perturbations in a brane universe. This formalism has the advantage to introduce quantities that are similar to the usual treatments of cosmological perturbations in standard cosmology, thus making easier in the future a comparison between brane cosmoloical perturbations and observable quantities, such as the large scale structure and the temperature anisotropies.
However, the equations governing the perturbations in the brane scenario are much more complicated than in standard cosmology. The main reason is that the equations for the metric perturbations, after a usual decomposition into ordinary (spatial) Fourier modes, will contain partial derivatives with respect to time and with respect to the fifth coordinate, in contrast with standard cosmology where one ends up with only ordinary differential equations with respect to time. Solving these equations appears to be quite a challenge.
Note: while the present work was being completed, two works on the same subject ( and ) have appeared on hep-th.
###### Acknowledgements.
I would like to thank P. Binétruy, C. Deffayet, B. Carter, N. Deruelle and T. Dolezel for very interesting discussions.
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# Evolution of interfaces and expansion in width Paper supported in part by KBN grant 2P03B 095 13
## 1 Introduction
Important aspect of dynamics of phase transitions in condensed matter is time evolution of an interface separating a retreating phase from the new one. Studied in the framework of Ginzburg-Landau type effective macroscopic models the interface can be regarded as a kind of smooth, asymmetric domain wall subject to a transverse force. The asymmetry and the force are due to a difference of potential energy across the interface. Pertinent evolution equations for order parameters typically are nonlinear partial differential equations. In general they imply rather nontrivial phase ordering dynamics, see, e.g., review article . Relativistic version of the problem, not considered here, also is interesting because of its connection with field-theoretical cosmology, .
Recently, evolution of ordinary domain walls has been studied with the help of Hilbert-Chapman-Enskog method applied in a suitably chosen comoving coordinate system , , – a systematic and consistent perturbative scheme has been developed. It yields the relevant solutions of the evolution equations in the form of expansion in a parameter $`l_0`$ which can be regarded as a measure of width of a static planar domain wall. Consecutive terms in this expansion contain extrinsic curvatures of a surface comoving with the wall, and also certain functions (below denoted by $`C_k`$) which can be regarded as fields defined on that surface and coupled to the extrinsic curvatures. In the present paper, which is a sequel to the previous paper , we show that that perturbative expansion can be generalised to the case of curved interfaces.
The Hilbert-Chapman-Enskog method and the comoving coordinates technique, which we have learned from , , respectively, have already been used in theoretical investigations of planar interfaces , and of curved ones in superconducting films . We apply these tools to curved interfaces in the three dimensional space, in a Ginzburg-Landau type model defined by formulas (1), (2) and Eq.(3) below. Interfaces in this model have been investigated for a long time, see, e.g., . Our main contribution consists in providing a systematic iterative scheme for generating the relevant solution in the form of a perturbative series. The role of small parameters is played by the ratios $`l_0/R_i,`$ where $`l_0`$ is the width and $`R_i`$ curvature radia of the interface. In spite of nonlinearity of the evolution equation the perturbative contributions can be generated in surprisingly simple manner. This is achieved by introducing the functions $`C_k`$ which saturate certain integrability conditions. As an application, we derive a formula for local velocity of the interface with curvature corrections included, and we discuss critical size of nucleating spherical droplets.
We consider a system described by a real, scalar, non-conserved order parameter $`\mathrm{\Phi }`$, with the free energy $`F`$ of the form
$$F=d^3x\left(\frac{1}{2}K\frac{\mathrm{\Phi }}{x^\alpha }\frac{\mathrm{\Phi }}{x^\alpha }+V(\mathrm{\Phi })\right),$$
(1)
where
$$V=A\mathrm{\Phi }^2+B\mathrm{\Phi }^3+C\mathrm{\Phi }^4.$$
(2)
Time evolution is governed by the dissipative nonlinear equation
$$\gamma \frac{\mathrm{\Phi }(\stackrel{}{x},t)}{t}=K\mathrm{\Delta }\mathrm{\Phi }V^{}(\mathrm{\Phi }).$$
(3)
Here $`(x^\alpha )_{\alpha =1,2,3}`$ are Cartesian coordinates in the space, $`V^{}`$ denotes the derivative $`dV/d\mathrm{\Phi }`$, and $`K,\gamma ,A,B,C`$ are positive constants. The free energy of the form (1) arises in, e.g., de Gennes-Landau description of nematic-isotropic transition in nematic liquid crystals in a single elastic constant approximation ($`L_2=0`$) . Then $`\mathrm{\Phi }=0`$ corresponds to the isotropic liquid phase, while in the nematic phase $`\mathrm{\Phi }0`$.
The concrete form (2) of the potential has the advantage that solution of Eq.(3) describing a planar interface has a simple, explicitly known form. It can be found in, e.g., . The planar interface plays important role in the perturbative scheme: the main idea is that there exist curved interfaces which do not differ much from the planar one if considered in appropriately chosen coordinate system (which in particular should comove with the interface). Therefore, one may hope that $`\mathrm{\Phi }(\stackrel{}{x},t)`$ for such curved interfaces can be calculated perturbatively, with the planar interface giving the zeroth order term. To this end, it is necessary to introduce the comoving coordinate system explicitly, and to give a prescription for the iterative computation of the perturbative corrections.
The plan of our paper is as follows. In Section 2 we recall the planar interface and the comoving coordinates. This Section contains preliminary material quoted here for convenience of the reader as well as in order to fix our notation. In Section 3 we describe the perturbative scheme for the curved interfaces. In Section 4 we present formulas for the velocity and the free energy of the curved interface. Section 5 is devoted to a discussion of spherical droplets of the stable phase which nucleate during the phase transition. Several remarks are collected in Section 6. In Appendix A we construct solutions of linear equations obeyed by corrections to transverse profile of the interface. Appendix B contains a brief discussion of stability of the interface.
## 2 The preliminaries
### 2.1 The homogeneous planar interface
Let us assume that the planar interface is perpendicular to the $`z`$ axis ($`zx^3`$) and homogeneous. Then $`\mathrm{\Phi }`$ depends only on $`z`$ and $`t`$, and Eq.(3) is reduced to
$$\gamma _t\mathrm{\Phi }=K_z^2\mathrm{\Phi }V^{}(\mathrm{\Phi }).$$
(4)
The interface type solution $`\mathrm{\Phi }(z,t)`$ interpolates smoothly between minima of $`V`$ when $`z`$ changes from $`\mathrm{}`$ to $`+\mathrm{}`$. $`V`$ given by formula (2) has two minima
$$\mathrm{\Phi }_{}=0,\mathrm{\Phi }_+=\sqrt{K/(8Cl_0^2)},$$
where $`l_0`$ is given by formula (9) below. The corresponding phases we shall call isotropic and ordered, respectively. Let us multiply Eq. (4) by $`_z\mathrm{\Phi }`$ and integrate over $`z`$. The resulting identity
$$\gamma _{\mathrm{}}^+\mathrm{}𝑑z_z\mathrm{\Phi }_t\mathrm{\Phi }=V(\mathrm{\Phi }_{})V(\mathrm{\Phi }_+),$$
implies that $`_t\mathrm{\Phi }0`$ if the minima are nondegenerate. Further assumption that the interface moves in a uniform manner with velocity $`v_0`$, that is that
$$\mathrm{\Phi }(z,t)=\mathrm{\Phi }(zz_0v_0t),$$
leads to the formula
$$\gamma v_0_{\mathrm{}}^+\mathrm{}𝑑Z\mathrm{\Phi }^2(Z)=V(\mathrm{\Phi }_+)V(\mathrm{\Phi }_{}),$$
(5)
where
$$Z=zz_0v_0t.$$
(6)
Hence the interface moves towards the region of higher potential $`V(\mathrm{\Phi }_\pm )`$, as expected.
The exact solution of Eq.(4) has the following form
$$\mathrm{\Phi }_0=\sqrt{\frac{K}{8l_0^2C}}\frac{1}{1+\mathrm{exp}(Z/2l_0)},$$
(7)
where
$$v_0=\sqrt{\frac{9K}{32\gamma ^2C}}\left(B\sqrt{9B^232AC}\right),$$
(8)
and
$$l_0^1=\frac{1}{2\sqrt{2KC}}\left(3B+\sqrt{9B^232AC}\right).$$
(9)
The constant $`z_0`$ can be regarded as the position of the interface at $`t=0,`$ and $`l_0`$ as its width. $`\mathrm{\Phi }_0`$ smoothly interpolates between the local minima of $`V`$: $`\mathrm{\Phi }_{}`$ for $`z\mathrm{}`$ and $`\mathrm{\Phi }_+`$ for $`z+\mathrm{}`$. The corresponding values of the potential are
$$V(\mathrm{\Phi }_{})=0,V(\mathrm{\Phi }_+)=\frac{K\gamma v_0}{96l_0^3C}.$$
At $`\mathrm{\Phi }_m=3B/4C\mathrm{\Phi }_+`$ the potential V has a local maximum if $`A>0`$. The substitutions $`ZZ`$ and $`v_0v_0`$ in formulas (7), (8) give another solution of Eq.(4), called anti-interface.
It is clear that solution (7) exists if
$$9B^232AC.$$
(10)
The parameter $`A`$ has the following dependence on the temperature $`T`$
$$A=a(TT_{}),$$
where $`a>0.`$ The constants $`B,C`$ and $`a`$, approximately do not depend on temperature. Condition (10) is satisfied if the temperature $`T`$ is from the interval $`(T_{},T_c)`$, where $`T_c`$ is determined from the equation $`9B^2=32aC(T_cT_{})`$. It is clear that $`T_c>T_{}`$. For temperatures in this interval one phase is stable and the other one is metastable.
The potential (2) can also lead to a static, symmetric domain wall. Namely, for the temperature $`T_0`$ such that $`B^2=4AC`$ the velocity $`v_0`$ vanishes, $`V(\mathrm{\Phi }_+)=V(\mathrm{\Phi }_{})=0`$, and the potential can be written in the following form
$$V=C\left[(\mathrm{\Phi }\mathrm{\Phi }_m)^2\mathrm{\Phi }_m^2\right]^2,$$
(11)
where for $`T=T_0`$
$$\mathrm{\Phi }_m=\frac{B}{4C}.$$
In this particular case there is the degenerate ground state given by $`\mathrm{\Phi }=\mathrm{\Phi }_\pm .`$ The potential (11) possesses the $`Z_2`$ symmetry
$$\mathrm{\Phi }2\mathrm{\Phi }_m\mathrm{\Phi },$$
and the interface becomes a static homogeneous, symmetric domain wall with the $`Z_2`$ topological charge.
### 2.2 The comoving coordinates
Here we quote the main definitions in order to introduce our notation. More detailed description of this change of coordinates, as well as a discussion of related mathematical questions, can be found in , , .
The comoving coordinates in the space, denoted by $`(\sigma ^\alpha )=(\sigma ^1,\sigma ^2,\sigma ^3=\xi )`$ where $`\alpha =1,2,3`$, are defined by the following formula
$$\stackrel{}{x}=\stackrel{}{X}(\sigma ^i,t)+\xi \stackrel{}{p}(\sigma ^i,t).$$
(12)
Here $`\stackrel{}{x}=(x^\alpha )`$, where $`x^\alpha `$ are the Cartesian coordinates in the space $`R^3`$. The points $`\stackrel{}{X}(\sigma ^i,t)`$ form a smooth surface $`S`$ which is parametrised by the two coordinates <sup>1</sup><sup>1</sup>1The Greek indices $`\alpha ,\beta ,\mathrm{}`$ have values 1,2,3 and they refer to the threedimensional space, while the Latin indices $`i,j,k,l,\mathrm{}`$ have values 1,2 and they refer to the inner coordinates $`\sigma ^1,\sigma ^2`$ on the surface $`S`$. $`\sigma ^1,\sigma ^2`$. In general, the interface moves in space, hence $`\stackrel{}{X}`$ depends on the time $`t`$.
The surface $`S`$ is fastened to the interface – the shape of it mimics the shape of the interface and they move together. We shall see that for consistency of the perturbative scheme $`\stackrel{}{X}(\sigma ^i,t)`$ has to obey certain equation from which one can determine evolution of the surface $`S`$. The coordinate $`\xi `$ parametrises the axis perpendicular to the interface at the point $`\stackrel{}{X}(\sigma ^i,t)`$. The vector $`\stackrel{}{p}(\sigma ^i,t)`$ is a unit normal to $`S`$ at this point, that is
$$\stackrel{}{p}^{\mathrm{\hspace{0.33em}2}}=1,\stackrel{}{p}\stackrel{}{X}_{,k}(\sigma ^i,t)=0,$$
where $`\stackrel{}{X}_{,k}=\stackrel{}{X}/\sigma ^k`$. The surface $`S`$ is characterized in particular by induced metric tensor on $`S`$
$$g_{ik}=\stackrel{}{X}_{,i}\stackrel{}{X}_{,k},$$
and the extrinsic curvature coefficients
$$K_{ij}=\stackrel{}{p}\stackrel{}{X}_{,ij}.$$
The matrix $`(g^{ik})`$ is by definition the inverse of the matrix $`(g_{kl})`$, i.e. $`g^{ik}g_{kl}=\delta _l^i`$.
The two by two matrix $`(K_{ik})`$ is symmetric. Two eigenvalues $`k_1,k_2`$ of the matrix $`(K_j^i)`$, where $`K_j^i=g^{il}K_{lj}`$, are called extrinsic curvatures of $`S`$ at the point $`\stackrel{}{X}`$. The main curvature radia are defined as $`R_i=1/k_i`$. Thus, by the definition
$$K_i^i=\frac{1}{R_1}+\frac{1}{R_2},det(K_j^i)=\frac{1}{R_1R_2}.$$
In general the curvature radia vary along $`S`$ and with time.
The coordinates $`(\sigma ^\alpha )`$ replace the Cartesian coordinates $`(x^\alpha )`$ in a vicinity of the interface. Components of metric tensor in the space transformed to the new coordinates are denoted by $`G_{\alpha \beta }`$. They are given by the following formulas
$$G_{33}=1,G_{3k}=G_{k3}=0,G_{ik}=N_i^lg_{lr}N_k^r,$$
where
$$N_i^l=\delta _i^l\xi K_i^l.$$
Dependence of $`G_{\alpha \beta }`$ on the transverse coordinate $`\xi `$ is explicit, and $`\sigma ^1,\sigma ^2`$ enter through the tensors $`g_{ik},K_r^l`$ which characterise geometry of the surface $`S`$.
Components $`G^{\alpha \beta }`$ of the inverse metric tensor have the form
$$G^{33}=1,G^{3k}=G^{k3}=0,G^{ik}=(N^1)_r^ig^{rl}(N^1)_l^k,$$
where
$$(N^1)_r^i=\frac{1}{N}\left[(1\xi K_l^l)\delta _r^i+\xi K_r^i\right],$$
and
$$N=det(N_k^i)=1\xi K_i^i+\frac{1}{2}\xi ^2(K_i^iK_l^lK_l^iK_i^l)=(1\frac{\xi }{R_1})(1\frac{\xi }{R_2}).$$
In order to transform Eq.(3) to the comoving coordinates we use the standard formula
$$\mathrm{\Delta }\mathrm{\Phi }=\frac{1}{\sqrt{G}}\frac{}{\sigma ^\alpha }\left(\sqrt{G}G^{\alpha \beta }\frac{\mathrm{\Phi }}{\sigma ^\beta }\right),$$
(13)
where $`G=det(G_{\alpha \beta }),\sqrt{G}=\sqrt{g}N,g=det(g_{ik})`$.
The time derivative in Eq.(3) is taken under the condition that all $`x^\alpha `$ are constant. It is convenient to use time derivative taken at constant $`\sigma ^\alpha `$. They are related by the formula
$$\frac{}{t}|_{x^\alpha }=\frac{}{t}|_{\sigma ^\alpha }+\frac{\sigma ^\beta }{t}|_{x^\alpha }\frac{}{\sigma ^\beta }.$$
(14)
Finally, let us introduce the dimensionless variables $`s`$ and $`\varphi `$ instead of, respectively, $`\xi `$ and $`\mathrm{\Phi }`$:
$$\xi =2l_0s,\mathrm{\Phi }(\xi ,\sigma ^i,t)=\sqrt{\frac{K}{8Cl_0^2}}\varphi (s,\sigma ^i,t).$$
(15)
The coordinate $`s`$ gives the distance from $`S`$ in the unit $`2l_0`$ related to the width of the planar interface.
Using formulas (13-15) we can write Eq.(3) in the following form, which is convenient for construction of the expansion in width:
$`{\displaystyle \frac{2l_0^2\gamma }{K}}{\displaystyle \frac{\varphi }{t}}|_{\sigma ^k}\overline{v}{\displaystyle \frac{\varphi }{s}}{\displaystyle \frac{2l_0^2\gamma }{K}}(N^1)_k^ig^{kr}\stackrel{}{X}_{,r}(\dot{\stackrel{}{X}}+2l_0s\dot{\stackrel{}{p}})\varphi _{,i}`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \frac{^2\varphi }{s^2}}+{\displaystyle \frac{1}{2N}}{\displaystyle \frac{N}{s}}{\displaystyle \frac{\varphi }{s}}+2l_0^2{\displaystyle \frac{1}{\sqrt{g}N}}{\displaystyle \frac{}{\sigma ^j}}\left(G^{jk}\sqrt{g}N\varphi _{,k}\right)`$ (16)
$`\alpha \varphi +(1+\alpha )\varphi ^2\varphi ^3,`$
where
$$\overline{v}=\frac{\gamma l_0}{K}\stackrel{}{p}\dot{\stackrel{}{X}}$$
is dimensionless transverse velocity of the surface $`S`$, the dot denotes the derivative $`/t|_{\sigma ^\alpha }`$, and
$$\alpha =\frac{4Al_0^2}{K}.$$
Formula (9) and the condition (10) imply that $`0\alpha 1`$ for temperatures in the range $`[T_{},T_c]`$.
The homogeneous planar interface (7) can be obtained from the evolution equation written in the form (16) in the following manner. As the surface $`S`$ we take a plane, hence $`K_j^i=0`$. Moreover, $`S`$ is assumed to move with constant velocity $`v_0`$, hence
$$\stackrel{}{p}\dot{\stackrel{}{X}}_0=v_0=\text{const}.$$
Finally,
$$\frac{\varphi }{t}|_{\sigma ^\alpha }=0$$
because we look at the interface from the comoving reference frame, and $`\varphi /\sigma ^i=0`$ because of the homogeneity. Then equation (16) is reduced to
$$\overline{v}_0\frac{\varphi }{s}=\frac{1}{2}\frac{^2\varphi }{s^2}\alpha \varphi +(1+\alpha )\varphi ^2\varphi ^3.$$
(17)
The solution previously given by formulas (6-9) now has the form
$$\varphi =\varphi _0(s),\overline{v}_0=\alpha \frac{1}{2},$$
(18)
where
$$\varphi _0(s)=\frac{\mathrm{exp}(ss_0)}{1+\mathrm{exp}(ss_0)}.$$
(19)
$`\varphi _0(s)`$ smoothly interpolates between 0 and 1. This corresponds to interpolation between the minima $`\mathrm{\Phi }_{},\mathrm{\Phi }_+`$ of the potential $`V`$ if $`\stackrel{}{p}`$ is directed from negative towards positive $`z`$’s. If we choose the opposite direction for $`\stackrel{}{p}`$ we obtain the anti-interface. The constant $`s_0`$ corresponds to $`z_0`$ from formula (6).
## 3 The expansion in width for curved interfaces
Let us begin from a brief description of ideas underlying the calculations presented below. The set of solutions of the nonlinear, partial differential equation (3) is very large. We are interested here only in a rather special subset of it, consisting of solutions which represent evolution of a smooth interface. Moreover, even within this subclass we concentrate on rather special interfaces, called by us the ‘basic’ ones. Their defining feature is that one can find a comoving coordinate system in which the order parameter of the interface is essentially given by $`\varphi _0(s)`$, formula (19), modified by small corrections which take into account the nonvanishing curvature.
By writing the evolution equation in the form (16) we have shown that $`l_0`$ can be regarded as a parameter analogous to a coupling constant – it appears in Eq.(16) only as a coefficient in several (but not all) terms. Therefore, one may hope that a systematic perturbative expansion in $`l_0`$ will turn out useful, as it is the case with other perturbative expansions so numerous in theoretical physics. The perturbative series can be constructed in the standard manner: the sought for solution $`\varphi `$ and the velocity $`\overline{v}`$ are written in the form
$$\varphi =\varphi _0+l_0\varphi _1+l_0^2\varphi _2+\mathrm{},\overline{v}=\overline{v}_0+l_0\overline{v}_1+l_0^2\overline{v}_2+\mathrm{},$$
(20)
and inserted in Eq.(16). Coefficients in front of successive powers of $`l_0`$ in this equation are equated to zero. Notice that after the rescaling $`\xi =2l_0s`$ the expansion parameter $`l_0`$ is present also in $`N`$ and $`(N^1)_k^i`$. In the zeroth order we obtain the following equation
$$\overline{v}_0\frac{\varphi _0}{s}=\frac{1}{2}\frac{^2\varphi _0}{s^2}\alpha \varphi _0+(1+\alpha )\varphi _0^2\varphi _0^3.$$
(21)
which formally coincides with Eq.(17). Therefore, we can immediately write the relevant solution
$$\overline{v}_0=\alpha \frac{1}{2},\varphi _0(s,\sigma ^i,t)=\frac{\mathrm{exp}(sC_0(\sigma ^i,t))}{1+\mathrm{exp}(sC_0(\sigma ^i,t))}.$$
(22)
There are however two differences between the planar solution (18), (19) and the solution (22). First, we do not assume homogeneity of the interface, therefore the constant $`s_0`$ from formula (19) is replaced by the function $`C_0(\sigma ^i,t)`$ of the indicated variables. Second, the surface $`S`$ is not fixed yet, while in the former case it was a plane.
It is convenient to rewrite Eq.(16) as equation for the corrections $`\delta \varphi ,\delta \overline{v}`$, which are defined by the formulas
$$\varphi =\varphi _0(s,\sigma ^i,t)+\delta \varphi ,\overline{v}=\overline{v}_0+\delta \overline{v}.$$
(23)
After taking onto account the fact that $`\varphi _0`$ obeys Eq.(21) we obtain equation of the form
$$\widehat{L}\delta \varphi =f,$$
(24)
with the linear operator $`\widehat{L}`$
$$\widehat{L}=\frac{1}{2}\frac{^2}{s^2}+(\alpha \frac{1}{2})\frac{}{s}\alpha +2(\alpha +1)\varphi _03\varphi _0^2,$$
and
$`f=\left({\displaystyle \frac{1}{2N}}{\displaystyle \frac{N}{s}}+\delta \overline{v}\right){\displaystyle \frac{\varphi _0}{s}}+{\displaystyle \frac{2\gamma l_0^2}{K}}\left({\displaystyle \frac{\delta \varphi }{t}}|_{\sigma ^\alpha }{\displaystyle \frac{C_0}{t}}|_{\sigma ^\alpha }{\displaystyle \frac{\varphi _0}{s}}\right)`$
$`{\displaystyle \frac{2l_0^2\gamma }{K}}(N^1)_k^ig^{kr}\stackrel{}{X}_{,r}(\dot{\stackrel{}{X}}+2l_0s\dot{\stackrel{}{p}})\left(\delta \varphi _{,i}C_{0,i}{\displaystyle \frac{\varphi _0}{s}}\right)`$
$`\left({\displaystyle \frac{1}{2N}}{\displaystyle \frac{N}{s}}+\delta \overline{v}\right){\displaystyle \frac{\delta \varphi }{s}}2l_0^2{\displaystyle \frac{1}{\sqrt{g}N}}{\displaystyle \frac{}{\sigma ^j}}\left(G^{jk}\sqrt{g}N\left(\delta \varphi _{,k}C_{0,k}{\displaystyle \frac{\varphi _0}{s}}\right)\right)`$
$`+(3\varphi _0\alpha 1)(\delta \varphi )^2+(\delta \varphi )^3,`$ (25)
where $`\overline{v}_0`$ and $`\varphi _0`$ are given by formulas (22). Now it is easy to see that order by order in $`l_0`$ we obtain enumerated by $`n=1,2,\mathrm{}`$ inhomogeneous linear differential equations of the form
$$\widehat{L}\varphi _n=f_n.$$
(26)
For example,
$$f_1=\left(K_i^i\overline{v}_1\right)\frac{\varphi _0}{s}.$$
(27)
In the whole perturbative scheme Eq.(21) is the only nonlinear equation for the contributions to the order parameter $`\varphi `$. We show in the Appendix that Eqs.(26) can easily be solved with the help of standard methods – one can construct the relevant Green’s function for $`\widehat{L}`$. It is remarkable that the same operator $`\widehat{L}`$ appears in all equations (26), and that the form of it does not depend on the surface $`S`$. For these reasons calculation of the corrections $`\varphi _n`$ is reduced to relatively simple task of finding $`f_n`$ and calculating the one dimensional integrals over $`s`$ shown in the Appendix.
The perturbative Ansatz (20) and Eqs.(26) are two parts of the Hilbert-Chapman-Enskog method. The third and most crucial part consists of integrability conditions for Eqs.(26) . Such conditions appear because the operator $`\widehat{L}^{}`$, the Hermitean conjugate of $`\widehat{L}`$, has a normalizable eigenstate with the eigenvalue equal to zero. Such eigenstate is called the zero mode. Let us first find the zero mode for the operator $`\widehat{L}`$: inserting $`\varphi _0`$ in Eq.(21) and differentiating this equation with respect to $`s`$ gives the following identity
$$\widehat{L}\psi _r=0,$$
(28)
where
$$\psi _r(s,\sigma ^i,t)=\frac{\varphi _0}{s}=\frac{\mathrm{exp}(sC_0)}{[1+\mathrm{exp}(sC_0)]^2}.$$
(29)
Notice that $`\psi _r`$ exponentially vanishes for $`s\pm \mathrm{}`$. Because the operator $`/s`$ is anti-Hermitean with respect to the scalar product $`<g_1|g_2>=_{\mathrm{}}^+\mathrm{}𝑑sg_1^{}g_2,`$ the operator $`\widehat{L}`$ is not Hermitean. Its Hermitean conjugate has the form
$$\widehat{L}^{}=\frac{1}{2}\frac{^2}{s^2}(\alpha \frac{1}{2})\frac{}{s}\alpha +2(\alpha +1)\varphi _03\varphi _0^2.$$
The operator $`\widehat{L}^{}`$ has a zero mode too, namely
$$\widehat{L}^{}\psi _l=0,$$
(30)
where <sup>2</sup><sup>2</sup>2The subscripts l and r stand for left and right, respectively. The point is that (30) can be written as $`\psi _l\widehat{L}=0`$.
$$\psi _l=\mathrm{exp}[(2\alpha 1)(sC_0)]\psi _r.$$
(31)
The function $`\psi _l`$ vanishes exponentially for $`s\pm \mathrm{}`$ because $`0<\alpha <1`$ for all temperatures in the range $`(T_{},T_c)`$. For $`\alpha =1/2`$, that is when the interface becomes the domain wall, the two zero modes coincide.
Let us multiply both sides of Eqs.(26) by $`\psi _l(s)`$ and take the integral $`_{\mathrm{}}^+\mathrm{}𝑑s`$. The l.h.s. of the resulting formula vanishes because of (30), hence
$$_{\mathrm{}}^+\mathrm{}𝑑s\psi _lf_n=0$$
(32)
for $`n=1,2,\mathrm{}`$ . It turns out that these conditions are nontrivial. In particular, they give evolution equation for the surface $`S`$.
It should be noted that the conditions (32) are in fact approximate, but the neglected terms are exponentially small. The point is that in order to obtain Eqs.(26) we use the expansions of the type
$$\frac{1}{12l_0s/R_i}=\underset{k=0}{\overset{\mathrm{}}{}}\left(\frac{2l_0s}{R_i}\right)^k$$
$`(i=1,2)`$, which are convergent for $`s<s_M`$ where $`s_M=\text{min}(R_1/2l_0,R_2/2l_0)`$. Therefore, when deriving conditions (32), the integration range should be restricted to $`|s|<s_M`$. Because of the exponential decrease of $`\psi _l`$ and $`f_n`$ at large $`|s|`$, this will give exponentially small corrections to these conditions. We assume that the ratios $`R_i/l_0`$ are so large that we may neglect those corrections.
Let us compute the basic interface up to the order $`l_0`$. For $`n=1`$, condition (32) gives
$$\overline{v}_1=\frac{1}{R_1}+\frac{1}{R_2}$$
(33)
because
$$a_0(\alpha )=_{\mathrm{}}^+\mathrm{}𝑑s\psi _l\psi _r=B(2\alpha +1,32\alpha )$$
does not vanish for $`\alpha `$ in the interval (0, 1). Here $`B`$ denotes the Euler beta function. The function $`a_0(\alpha )`$ has the symmetric "U" shape in the interval $`\alpha [0,1]`$, with the minimum equal to 1/6 at $`\alpha =1/2`$, and the upper ends reaching the value 1/3 for $`\alpha =0`$ and 1. The condition (33) implies that the surface $`S`$ obeys the following evolution equation
$$\frac{\gamma }{K}\dot{\stackrel{}{X}}\stackrel{}{p}=\frac{2\alpha 1}{2l_0}+K_i^i.$$
(34)
It formally coincides with the well-known Allen-Cahn equation .
Now Eq.(26) with $`n=1`$ is reduced to
$$\widehat{L}\varphi _1=0.$$
(35)
It has the following solution which vanishes at $`s\pm \mathrm{}`$
$$\varphi _1=C_1(\sigma ^i,t)\psi _r,$$
(36)
where $`C_1`$ is a smooth function of the indicated variables.
For $`n=2`$, we obtain from formula (25), after taking into account the results (33) and (36),
$`f_2=(2sK_j^iK_i^j\overline{v}_2)_s\varphi _0+(3\varphi _0\alpha 1)C_1^2(_s\varphi _0)^2`$
$`+{\displaystyle \frac{2\gamma }{K}}\left(_tC_0+g^{ik}\stackrel{}{X}_{,k}\dot{\stackrel{}{X}}C_{0,i}\right)_s\varphi _0+2{\displaystyle \frac{1}{\sqrt{g}}}{\displaystyle \frac{}{\sigma ^i}}\left(g^{ik}\sqrt{g}C_{0,k}_s\varphi _0\right).`$ (37)
Straightforward integration over $`s`$ as in (32) can be a little bit cumbersome. This calculation can be significantly simplified with the help of the following identity
$$2𝑑s\psi _l[3\varphi _0(\alpha +1)]_s\varphi _0\varphi _n=𝑑s_s\psi _lf_n,$$
(38)
where in the case at hand $`n=1`$. Identity (38) is obtained from Eq.(26) by differentiating both sides of it with respect to $`s`$, multiplying by $`\psi _l`$ and integrating over $`s`$, just as in the derivation of the integrability conditions (32). The integrability condition gives
$`a_1(\alpha )K_j^iK_i^j\overline{v}_2+{\displaystyle \frac{2\gamma }{K}}\left(_tC_0+g^{ik}\stackrel{}{X}_{,k}\dot{\stackrel{}{X}}C_{0,i}\right)`$
$`+2\mathrm{\Delta }_2C_0+(2\alpha 1)g^{ik}C_{0,i}C_{0,k}+2C_0K_j^iK_i^j=0,`$ (39)
where
$$\mathrm{\Delta }_2=\frac{1}{\sqrt{g}}\frac{}{\sigma ^i}\left(\sqrt{g}g^{ik}\frac{}{\sigma ^k}\right)$$
is the Laplacian on the surface $`S`$ and
$$a_1(\alpha )=a_0(\alpha )^1\frac{da_0(\alpha )}{d\alpha }.$$
For $`\alpha `$ from the interval the function $`a_1(\alpha )`$ is almost linear. In particular, $`a(0)=3,a(1/2)=0,a(1)=3.`$
The integrability condition (39) leaves the freedom of choosing whether we keep nonvanishing $`C_0`$ or $`\overline{v}_2`$. It is clear that we can not put to zero both of them unless $`K_j^iK_i^j=0`$ (then $`S`$ is a plane). The choice $`C_0=0`$ gives
$$\overline{v}_2=a_1(\alpha )K_j^iK_i^j.$$
(40)
This implies a correction to the Allen-Cahn equation of the form
$$\frac{\gamma l_0}{K}\delta _1(\dot{\stackrel{}{X}}\stackrel{}{p})=l_0^2\overline{v}_2,$$
where $`\delta _1(.)`$ denotes the first order correction to $`\dot{\stackrel{}{X}}\stackrel{}{p}`$. In consequence, there will be a correction to the solution $`\stackrel{}{X}`$ of the Allen-Cahn equation, and corrections to $`g_{ik}`$ and $`K_{ik}`$. These corrections have to be taken into account when calculating $`f_k`$ with $`k3.`$ It is clear that this version of the perturbative scheme is rather cumbersome.
On the other hand, if we put $`\overline{v}_2=0`$, then the evolution of the surface $`S`$ is still governed by the relatively simple Allen-Cahn equation (34). The integrability condition (39) is now saturated by the function $`C_0`$ – it has the form of evolution equation for $`C_0`$, namely
$`{\displaystyle \frac{\gamma }{K}}\left[_tC_0g^{ik}_{\sigma ^k}\stackrel{}{X}\dot{\stackrel{}{X}}C_{0,i}\right]\mathrm{\Delta }_2C_0`$ (41)
$`(\alpha {\displaystyle \frac{1}{2}})g^{ik}C_{0,i}C_{0,k}C_0K_j^iK_i^j={\displaystyle \frac{1}{2}}a_1(\alpha )K_j^iK_i^j.`$
Let us also check the third integrability condition. Formulas (20), (25) give
$`f_3=\overline{v}_3\psi _r+2K_i^i(3K_j^iK_i^j(K_i^i)^2)s^2\psi _r+{\displaystyle \frac{2\gamma }{K}}_t(C_1\psi _r)`$
$`+{\displaystyle \frac{2\gamma }{K}}g^{ir}\stackrel{}{X}_r\dot{\stackrel{}{X}}(C_1\psi _r)_{,i}+{\displaystyle \frac{4\gamma }{K}}C_{0,i}(g^{ir}\stackrel{}{X}_r\dot{\stackrel{}{p}}K^{ir}\stackrel{}{X}_r\dot{\stackrel{}{X}})s\psi _r`$
$`+2sK_j^iK_i^j_s\varphi _12\mathrm{\Delta }_2(C_1\psi _r)+4sK_i^i{\displaystyle \frac{1}{\sqrt{g}}}(\sqrt{g}g^{jk}C_{0,k}\psi _r)_{,j}`$
$`{\displaystyle \frac{4s}{\sqrt{g}}}(\sqrt{g}g^{jk}K_l^lC_{0,k}\psi _r)_{,j}+{\displaystyle \frac{8s}{\sqrt{g}}}(\sqrt{g}K^{jk}C_{0,k}\psi _r)_{,j}`$
$`+2(3\varphi _0\alpha 1)\varphi _1\varphi _2+\varphi _1^3,`$ (42)
where we have put $`\overline{v}_2=0`$. The term with $`\varphi _2`$ contains a new function $`C_2`$ (see the Appendix), however this function will not appear in the integrability condition because the integration over $`s`$ eliminates the term proportional to $`C_2`$. This can be seen from the identity (38): on the r.h.s. of it we have $`f_2`$ in which $`C_2`$ is not present. We see already from formula (42) that we can choose whether to keep nonvanishing $`\overline{v}_3`$ or $`C_1`$. For the same reason as in the case of $`n=2`$ we choose $`\overline{v}_3=0`$, and the integrability is saturated by $`C_1`$. The resulting evolution equation for $`C_1`$ has the following form
$`{\displaystyle \frac{\gamma }{K}}\left[_tC_1g^{ik}_{\sigma ^k}\stackrel{}{X}\dot{\stackrel{}{X}}_{\sigma ^i}C_1\right]\mathrm{\Delta }_2C_1K_j^iK_i^jC_1`$
$`(\alpha {\displaystyle \frac{1}{2}})g^{jk}_kC_0_jC_1=\left({\displaystyle \frac{1}{4}}a_2+a_1C_0+C_0^2\right)K_i^i[3K_j^kK_k^j(K_l^l)^2]`$
$`+2\left[g^{ij}_jK_l^l{\displaystyle \frac{\gamma }{K}}[g^{ir}(_r\dot{\stackrel{}{X}}\dot{\stackrel{}{p}})+K^{ir}(_r\stackrel{}{X}\dot{\stackrel{}{X}})]\right]_iC_0(C_0+{\displaystyle \frac{1}{2}}a_1)`$ (43)
$`{\displaystyle \frac{4}{\sqrt{g}}}_j(K^{jk}\sqrt{g}_kC_0)(C_0+{\displaystyle \frac{1}{2}}a_1)2K^{jk}_kC_0_jC_0\left[1+(2\alpha 1)(C_0+{\displaystyle \frac{1}{2}}a_1)\right],`$
where
$$a_2=\frac{a_0^{\prime \prime }}{a_0}.$$
It is easy to proceed to the second and higher orders. Using formulas from the Appendix one can write general solution $`\varphi _n`$ of Eq.(26). It contains the function $`C_n(\sigma ^i,t)`$ which obeys an equation analogous to (41) or (43). This equation follows from the integrability condition (32) with $`n`$ replaced by $`n+2`$ if we put $`\overline{v}_n=0`$. Due to identity (38), in the derivation of that equation we do not need explicit form of $`\varphi _{n+1}`$. In the present paper we will stop our considerations at the first order.
In order to obtain a concrete basic interface solution we have to specify initial data for equations (34), (41), (43). There is no restriction on the initial data, except the obvious requirement that perturbative corrections of given order should be small in comparison with the ones from preceding orders. In particular, $`l_0C_11,l_0^2C_21`$ and $`l_0/R_i1`$. If we consider the interface solution only up to the first order correction (36) then a simplification appears: without any loss in generality we may adopt the homogeneous initial data
$$C_0(\sigma ^i,t=t_0)=0,C_1(\sigma ^i,t=t_0)=0.$$
(44)
This can be justified as follows. Local deformation of $`S`$ by shifting a small piece of it along the direction $`\stackrel{}{p}`$ results in the corresponding shift of the coordinate $`s`$. Therefore for any given basic interface we can choose initial position of $`S`$ such that $`C_0(\sigma ^i,t=t_0)=0`$ in formula (22) for $`\varphi _0`$. This can be done at one instant of time, e.g., at the initial time. Values of $`C_0`$ and position of $`S`$ at later times are determined uniquely by Eqs.(34), (41), and in general $`C_0`$ does not vanish. Notice however that such a shift wiil influence terms of the order $`l_0^2`$ and higher in formula (25) for $`f`$ – due to the explicit presence of $`s`$ in $`N,(N^1)_k^i`$ $`f`$ is not invariant under the translations $`ss+C_0`$.
The r.h.s. of Eq.(41) vanishes for $`\alpha =1/2`$, that is in the domain wall case. In this case the initial condition (44) implies that $`C_0=0`$ for all times, and in consequence $`C_0`$ disappears from the first order perturbative solution.
As for $`C_1`$, the reason for the homogeneous initial condition is that the interface with the first order correction, that is
$$\varphi =\varphi _0(sC_0)+l_0C_1\psi _r(sC_0),$$
(45)
can be regarded as $`\varphi _0(sC_0+l_0C_1)`$ to the first order in $`l_0`$, so again we can cancel $`C_1`$ at the initial time $`t_0`$ by suitably correcting the initial position of the surface $`S`$. Let us stress again that this works only at the fixed time instant.
The initial data (44) imply that at the initial time the order parameter $`\varphi `$ is equal to $`\varphi _0(s)`$. Hence, the only freedom in choosing initial data for $`\varphi `$ we still have is position of the surface $`S`$. When the second and higher order corrections are included one has to allow for more general than (44) initial data for $`C_0`$ and $`C_1`$, nevertheless the initial form of $`\varphi `$ always is uniquely fixed by these data and the initial position of the surface $`S`$.
To recapitulate, the perturbative solution to the first order has the form (45) where $`\varphi _0,\psi _r`$ are given by formulas (22), (29), respectively. Formula (45) gives dependence on $`s`$ explicitly. The functions $`C_0,C_1`$ are to be determined from Eqs.(41),(43) with the initial data (44). One also has to solve Allen-Cahn equation (34). In certain cases these equations can be solved analytically, e. g., for spherical interface discussed in next Section. In general case one will be forced to use numerical methods. In comparison with the initial evolution equation (3) the gain is that the equations for $`S,C_0,C_1`$ involve only two spatial variables $`\sigma ^1,\sigma ^2`$. Such reduction of the number of independent variables is a great simplification for numerical calculations.
The perturbative solution obtained above can be used in calculations of physical characteristics of interfaces. In the following Section we obtain formulas for local velocity and surface tension of the interface. In the Section 5 we discuss evolution of a spherical interface.
## 4 Local velocity and surface tension of the interface
Let us apply the expansion in width in order to find local transverse velocity and surface tension of the interface. We shall use the first order solution (45). The velocity is obtained from the condition $`\varphi =\text{const}`$. It does not necessarily coincide with $`\stackrel{}{p}\dot{\stackrel{}{X}}`$ given by Allen-Cahn equation (34). Because we neglect terms of second and higher order in $`l_0`$, we may write $`\varphi `$ in the form $`\varphi _0(sC_0+l_0C_1)`$ from which we see that $`\varphi `$ is constant on surfaces given in the comoving coordinates by the condition $`sC_0+l_0C_1=s_0`$, where $`s_0`$ is a constant. It follows from formula (12) that in the laboratory Cartesian coordinate frame these surfaces are given by $`\stackrel{}{x}_0(\sigma ^i,t)`$, where
$$\stackrel{}{x}_0(\sigma ^i,t)=\stackrel{}{X}(\sigma ^i,t)+2l_0(s_0+C_0l_0C_1)\stackrel{}{p}(\sigma ^i,t).$$
The transverse velocity of the interface is equal to $`\dot{\stackrel{}{x}}_0\stackrel{}{p}.`$ In order to calculate it we take time derivative of $`\stackrel{}{x}_0`$, project it on $`\stackrel{}{p}`$, and use equations (34), (41), (43) with the initial data (44). The result can be written in the form
$$\frac{\gamma }{K}\stackrel{}{p}\dot{\stackrel{}{x}}_0=\frac{2\alpha 1}{2l_0}+\frac{1}{R_1}+\frac{1}{R_2}+l_0a_1\left(\frac{1}{R_1^2}+\frac{1}{R_2^2}\right)+l_0^2a_2\left(\frac{1}{R_1^3}+\frac{1}{R_2^3}\right).$$
(46)
The unit normal vector $`\stackrel{}{p}`$ is directed from the isotropic phase ($`\mathrm{\Phi }=\mathrm{\Phi }_0`$) to the ordered phase ($`\mathrm{\Phi }=\mathrm{\Phi }_+`$). In the next Section we shall use formula (46) in the case of spherical droplets.
The surface tension is another basic characteristics of the interface. It can be determined from a formula for free energy of the interface, which is defined as follows. The surface $`S`$ cuts the total volume of the sample into two regions denoted below by I and II. Let us imagine that in region I there is the homogeneous isotropic phase with constant free energy density equal to $`V(\mathrm{\Phi }_{})=0`$, and in region II the homogeneous ordered phase for which $`V(\mathrm{\Phi }_+)=K\gamma v_0/(96l_0^3C)`$. The normal vector $`\stackrel{}{p}`$ points to region II. The free energy of the interface is defined as the difference
$$F_i=FV(\mathrm{\Phi }_+)𝒱_{II},$$
where $`𝒱_{II}`$ denotes volume of the region II, and $`F`$ is the total free energy of the sample given by formula (1). We shall compute $`F_i`$ using the first order solution (45). Then, without any loss of generality we may put $`C_0=C_1=0`$ at the given time, as argued in the preceding Section, while the surface $`S`$ remains arbitrary. Therefore, we need only $`\varphi _0(s)=\mathrm{exp}(s)/(1+\mathrm{exp}(s))`$. Because the dependence on the coordinate $`s=\xi /2l_0`$ is explicit, we can integrate over $`s`$ in the formula (1) for the free energy $`F`$. The volume element and the gradient free energy are taken in the form
$$d^3x=\sqrt{G}d\xi d\sigma ^1d\sigma ^2,\frac{\varphi }{x^\alpha }\frac{\varphi }{x^\alpha }=G^{\alpha \beta }\frac{\varphi }{\sigma ^\alpha }\frac{\varphi }{\sigma ^\beta }.$$
Neglecting terms quadratic in $`l_0/R_i`$ we obtain the following formula for the free energy of the interface
$$F_i=_S\kappa 𝑑A,$$
where
$$\kappa =\frac{K^2}{96l_0^3C}\left[1(\frac{\pi ^2}{3}2)(12\alpha )(\frac{l_0}{R_1}+\frac{l_0}{R_2})\right]$$
(47)
can be regarded as the local surface tension of the interface at the point $`\sigma ^1,\sigma ^2`$. $`R_1,R_2`$ are the main curvature radia of the surface $`S`$ at that point, and $`dA=\sqrt{g}d\sigma ^1d\sigma ^2`$ is the surface element of $`S`$. Of course, this formula for $`F_i`$ can be trusted if $`l_0/R_i1`$.
For a spherical droplet of the ordered phase embedded in the isotropic phase $`R_1,R_2`$ are positive (the signs follow from formulas given in Section 2.2) and of course equal to the radius of the sphere. If $`\alpha <1/2`$, formula (46) implies that the droplet grows (if its radius is large enough) because $`\stackrel{}{p}\dot{\stackrel{}{x}}_0<0`$ and $`\stackrel{}{p}`$ is the inward normal. In this case the curvature correction diminishes the surface tension, and $`\kappa `$ increases as the droplet grows. In the reverse situation – the isotropic phase inside and the ordered one outside – the curvature increases the surface tension and $`\kappa `$ decreases as the droplet grows.
In the case of a growing droplet of radius $`R`$ of the isotropic phase in the ordered medium $`\kappa `$ has the same dependence on the curvature. Here $`\stackrel{}{p}`$ is the outward normal, $`\alpha >1/2,\stackrel{}{p}\dot{\stackrel{}{x}}_0>0`$, and $`R_1=R_2=R`$. Nevertheless the values of surface tension in both cases are different because $`l_0`$ present in formula (47) depends on $`\alpha `$, namely $`l_0(1+\alpha )`$, see formula (57) below.
Notice that the first order curvature correction to $`\kappa `$ vanishes in the domain wall case ($`\alpha =1/2`$).
## 5 Evolution of spherical interface
Let us now apply the formalism developed in Sections 3 and 4 to evolution of spherical droplets. We assume that $`\alpha 1/2`$ in order to exclude the relatively simpler case of the domain wall. The surface $`S`$ is parametrised by
$$\stackrel{}{X}_0=R(t)\stackrel{}{p}(\theta ,\psi ).$$
Here $`\theta ,\psi `$ are the spherical angles, and $`\stackrel{}{p}`$ is the inward (outward) normal to the sphere when $`\alpha <1/2`$ ($`\alpha >1/2`$). Thus, $`\stackrel{}{p}\dot{\stackrel{}{X}}_0=\dot{R},`$ with the upper sign corresponding to $`\alpha <1/2.`$ The Allen-Cahn equation has the form
$$\frac{\gamma }{2K}\dot{R}=\frac{1}{R_{}}\frac{1}{R},$$
(48)
where
$$R_{}=\frac{4l_0}{|12\alpha |}.$$
Integration of equation (48) yields the following formula
$$\frac{R(t)}{R_{}}+\mathrm{ln}|\frac{R(t)}{R_{}}1|=\frac{2K}{R_{}^2\gamma }t+\frac{R(0)}{R_{}}+\mathrm{ln}|\frac{R(0)}{R_{}}1|.$$
(49)
Evolution equation (41) for $`C_0`$ now has the from
$$\frac{\gamma }{2K}\dot{C}_0\frac{1}{R(t)^2}C_0=\frac{a_1}{2R(t)^2}.$$
(50)
It has the following solution
$$C_0(t)=\frac{a_1}{2}\frac{R_{}}{R(t)}\frac{R(0)R(t)}{R(0)R_{}}$$
(51)
which obeys the initial condition $`C_0(0)=0`$.
Evolution equation for $`C_1`$ is obtained from the general equation (43). For the spherical bubble it has the form
$$\frac{\gamma }{2K}\dot{C}_1\frac{1}{R^2}C_1=2(\frac{1}{4}a_2+a_1C_0+C_0^2)\frac{1}{R^3},$$
(52)
where as usual the upper sign corresponds to $`\alpha <1/2.`$ We do not know explicit solution of this equation.
There are two cases when the spherical droplets grow: a droplet of the ordered phase when $`\alpha <1/2`$, and a droplet of the disordered phase when $`\alpha >1/2`$. In both cases formula (46) for the radial velocity of expansion $`\dot{r}_0`$ gives
$$\frac{\gamma l_0}{2K}\dot{r}_0=\frac{|2\alpha 1|}{4}\frac{l_0}{R}+|a_1(\alpha )|(\frac{l_0}{R})^2a_2(\alpha )(\frac{l_0}{R})^3.$$
(53)
The expansion velocity is identical for all surfaces of constant $`\varphi `$.
It is clear that there is a minimal R, let us denote by $`R_{\text{min}}(\alpha )`$, such that $`\dot{r}_0>0`$. We have found numerically that
$$R_{\text{min}}(\alpha )=\frac{R_{}}{z(\alpha )},$$
(54)
where the function $`z(\alpha )`$ is symmetric with respect to $`\alpha =1/2`$ and it has values in the interval \[0.866, 1.049\]. For example, $`z(0)=0.866,z(0.1)=0.984,z(0.20)=1.045,z(0.3)=1.040z(0.4)=1.012,z(0.5)=1.0.`$ Notice that $`R_{\text{min}}(\alpha )`$ diverges when $`\alpha 1/2`$. Thus, in the Ginzburg-Landau model nucleation of expanding droplets is possible only if we heat the ordered phase to a temperature above $`T_0`$, or cool the isotropic phase below $`T_0`$.
For large time $`t`$, when the droplets are very large, the velocity $`\dot{r}(t)`$ becomes equal to the velocity of the planar interface
$$\dot{r}_{\mathrm{}}(\alpha )=\frac{K}{2\gamma l_0}|12\alpha |,$$
as expected. Notice that $`R_{\text{min}}(\alpha )`$ and $`\dot{r}_{\mathrm{}}(\alpha )`$ are not independent:
$$z(\alpha )\dot{r}_{\mathrm{}}(\alpha )R_{\text{min}}(\alpha )=\frac{2K}{\gamma }.$$
Parameter $`\alpha `$ is related to the temperature:
$$\alpha =\frac{2\theta }{12\theta +\sqrt{14\theta }},$$
where
$$\theta =\frac{8aC}{9B^2}(TT_{})$$
is a reduced temperature. $`T_0`$ and $`T_c`$ correspond to $`\theta =2/9`$ and 1/4, correspondingly. The interval $`\alpha [0,1/2]`$ corresponds to $`\theta [0,2/9]`$, and $`\alpha [1/2,1]`$ to $`\theta [2/9,1/4]`$. The temperature dependence of $`l_0`$ can be seen from formula
$$l_0=(1+\alpha )\frac{\sqrt{2KC}}{3B},$$
(55)
which follows from the definition (9) after some algebraic manipulations. $`R_{\text{min}}`$ is proportional to $`l_0/|12\alpha |`$, which can be written in the form
$$\frac{l_0}{|12\alpha |}=\frac{\sqrt{2KC}}{6B}\frac{4}{|3\sqrt{14\theta }1|}.$$
In the interval $`\theta [0,2/9]`$, which corresponds to $`T[T_{},T_0]`$, $`l_0/|12\alpha |`$ monotonically grows from $`\sqrt{2KC}/3B`$ to infinity.
Using formula (55) and the symmetry of $`z`$: $`z(1/2\delta )=z(1/2+\delta )`$, we obtain the following relation
$$\frac{R_{\text{min}}(1/2\delta )}{R_{\text{min}}(1/2+\delta )}=\frac{32\delta }{3+2\delta }<1,$$
(56)
where $`\delta (0,1/2)`$. Thus, the minimal size of the droplets of the isotropic phase which appear and grow when $`\alpha >1/2`$ is significantly larger than the size of the droplets of the ordered phase which can appear for $`\alpha <1/2`$.
The velocities $`\dot{r}_{\mathrm{}}`$ depend on temperature. In particular, comparing them for temperatures below and above $`T_0`$,
$$\frac{\dot{r}_{\mathrm{}}(1/2\delta )}{\dot{r}_{\mathrm{}}(1/2+\delta )}=\frac{3+2\delta }{32\delta }>1.$$
Hence, the droplets of the isotropic phase expand more slowly than the droplets with $`\varphi 1`$ inside.
Our main goal in this paper has been to develop the perturbative expansion for the curved interfaces. We plan to apply it to interfaces in liquid crystals in a subsequent work. Nevertheless, just in order to get an idea what our formulas predict, we have estimated $`l_0`$ and $`\dot{r}_{\mathrm{}}`$ for interfaces in nematic liquid crystal MBBA. The model defined by formulas (1), (2) and Eq.(3) can be related to de Gennes-Landau theory in a single elastic constant approximation ($`L_1=K,L_2=0`$). We take data found in : $`T_{}316K,a0.021J/(cm^3K),B0.07J/cm^3,C0.06J/cm^3`$ (after a change to our notation), and $`K610^{12}N`$, $`\gamma 5.210^2kg/(ms)`$. We have identified $`\gamma `$ with the rescaled rotational viscosity $`\gamma _1L_1/K_{11}`$ at a temperature close to $`T_{}`$. Then, $`T_0T_{}1K,T_cT_{}1.2K`$. The width $`l_0`$ and the velocity $`\dot{r}_{\mathrm{}}`$ of the planar interface are given by the following formulas:
$$l_040(1+\alpha )10^8cm,\dot{r}_{\mathrm{}}(\alpha )1.4\frac{|12\alpha |}{1+\alpha }\frac{cm}{s}.$$
Notice that even for rather small droplets with the radius of several hundred Ångström the ratio $`l_0/R`$ is rather small.
## 6 Remarks
1. We have shown how one can systematically compute curvature corrections to the transverse profile $`\varphi `$ and to the local velocity $`\dot{r}_0`$ of the interface. Due to the presence of functions $`C_k`$ evolution equation for the surface $`S`$ has the relatively simple form (34) and there are no curvature corrections to it. The formalism for interfaces is a generalisation of the one constructed for domain walls . The main new ingredients are the $`C_0`$ function and the $`(2\alpha 1)/l_0`$ term in Allen-Cahn equation (34). By including them we have significantly enlarged the range of physical applications of the perturbative scheme. This justifies the present publication.
2. The model we have considered is special in the sense that the exact planar interface solution $`\varphi _0`$ is known. Moreover, the solutions of the equations $`\widehat{L}\varphi _n=f_n`$ are given (almost) explicitly too, because the one dimensional integrations in formula (60) below can be easily calculated numerically. Here the crucial point is that we know explicitly the two linearly independent solutions $`\psi _r`$ and $`\psi _2`$ of the homogeneous equation $`\widehat{L}\psi =0`$. In other models, the analogs of $`\varphi _0,\psi _r,\psi _2`$ can be found at least numerically because the pertinent equations are relatively simple differential equations with the single independent variable $`\xi `$ (or $`s`$ after a rescaling). In our perturbative scheme we need to perform explicitly only integrals over $`s`$, as in integrability conditions (32) or in formula (60) for $`\stackrel{~}{\varphi }_n`$. Such integrals can easily be calculated numerically also in the case when only numerical solutions $`\varphi _0,\psi _r,\psi _2`$ are known.
3. In our approach, in order to describe the evolution of curved interface we introduce the surface $`S`$, and the functions $`C_k`$ which can be regarded as auxilliary fields defined on $`S`$ and coupled to extrinsic curvatures of it. The corresponding evolution equations, that is Allen-Cahn equation (34)for $`S`$ and equations (41), (43) and analogous ones for $`C_k`$, have one independent variable less than the original equation (3). This is significant simplification from the viewpoint of both computer simulations and analytical approaches. Therefore we think that our perturbative scheme is an interesting tool for investigations of dynamics of interfaces in Ginzburg-Landau effective models.
4. The perturbative solution we have presented above is based on the planar homogeneous interface $`\varphi _0(s)`$. Moreover, the dependence on the transverse coordinate $`s`$ is uniquely fixed by the perturbative scheme once the initial position of the surface $`S`$ and initial data for the functions $`C_k`$ are fixed. This means that in our scheme we obtain a special class of interfaces, distinguished by the particular form of the dependence on $`s`$. In other words, the transverse profiles of the interfaces provided by the perturbative solution are not arbitrary. Intuitively, the interfaces can be regarded as the planar interface folded to a required shape at the initial instant of time, and modified by necessary curvature corrections. Therefore, it seems appropriate to regard the curved interfaces obtained in our paper as the basic ones. More general interfaces could be obtained by choosing more general initial data and solving Eq.(3). For such generic interfaces no analytic perturbative approach is available.
5. One of advantages of a systematic perturbative approach is that one can make trustworthy estimates of neglected contributions and in consequence to check whether a given perturbative result is reliable. Formalism with that level of control can be used in order to make straightforward predictions of dynamical behaviour of an interface, but perhaps more important application is to "inverse problems", that is determination of parameters of Ginzburg-Landau effective theory. From dynamical behaviour of the interfaces one could reliably infer what values have parameters of the model. For example, for a given liquid crystal one could experimentally determine coefficients in front of higher powers of the order parameter like $`\mathrm{\Phi }^5,\mathrm{\Phi }^6`$ or terms of the type $`\mathrm{\Phi }_\alpha \mathrm{\Phi }_\alpha \mathrm{\Phi }`$ in the formula for free energy $`F`$. For a discussion of the form of $`F`$ for liquid crystals see, e.g., .
One could try to generalize our perturbative scheme for calculating the curvature corrections to interfaces coupled to a noise. In that case the r.h.s. of Eq.(3) would contain an external stochastic force which in particular would lead to fluctuations of the planar interface. In our formalism the dependence of, e.g., the surface tension $`\kappa `$ on curvature radia $`R_i`$ is explicit and it comes from purely geometric quantities like the metric tensor or Jacobian, while numerical coefficients in front of powers of $`l_0/R_i`$ are given by integrals over the $`s`$ coordinate and they are determined essentially by properties of the planar interface. Therefore, one may expect that if the stochastic force is present, values of these integrals would have to be averaged over the stochastic ensemble. This is another interesting direction in which one could continue the present work.
## 7 Appendix A. The equations $`\widehat{L}\varphi _n=f_n`$
In order to determine $`\varphi _n`$ we have to solve Eq.(26). Using standard methods it is not difficult to obtain appropriate solution.
Let us shift the variable $`s`$,
$$s=x+C_0,$$
in order to remove function $`C_0`$ from $`\varphi _0`$ present in the operator $`\widehat{L}`$. Then Eq.(26) acquires the following form
$$\stackrel{~}{L}\stackrel{~}{\varphi }_n(x,\sigma ^i,t)=\stackrel{~}{f}_n(x,\sigma ^i,t),$$
(57)
where
$$\stackrel{~}{L}=\frac{1}{2}\frac{^2}{x^2}+(\alpha \frac{1}{2})\frac{}{x}\alpha +2(\alpha +1)\stackrel{~}{\varphi }_0(x)3\stackrel{~}{\varphi }_0^2(x),$$
$$\stackrel{~}{f}_n(x,\sigma ^i,t)=f_n(x+C_0,\sigma ^i,t),\stackrel{~}{\varphi }_n(x,\sigma ^i,t)=\varphi _n(x+C_0,\sigma ^i,t),$$
and
$$\stackrel{~}{\varphi }_0(x)=\frac{\mathrm{exp}(x)}{1+\mathrm{exp}(x)}.$$
In the first step we find two linearly independent solutions of the homogeneous equation $`\stackrel{~}{L}\stackrel{~}{\psi }=0.`$ The zero mode $`\psi _r(x)`$ is one solution of this homogeneous equation, and the other one has the form
$$\psi _2(x)=\psi _r(x)h(x),$$
where $`\psi _r=d\stackrel{~}{\varphi }_0/dx`$ and
$`h(x)={\displaystyle \frac{1}{2\alpha +1}}\mathrm{exp}[(2\alpha +1)x]{\displaystyle \frac{1}{\alpha }}2\mathrm{exp}(2\alpha x)`$ (58)
$`+{\displaystyle \frac{6}{12\alpha }}[\mathrm{exp}[(12\alpha )x]1]+{\displaystyle \frac{1}{1\alpha }}2\mathrm{exp}[2(1\alpha )x]+{\displaystyle \frac{1}{32\alpha }}\mathrm{exp}[(32\alpha )x].`$
For $`\alpha 1/2`$ the first term in second line reduces to $`6x`$. Having those solutions one can construct the relevant Green’s function and the solution $`\stackrel{~}{\varphi }_n`$,
$`\stackrel{~}{\varphi }_n(x,\sigma ^i,t)=2\psi _r(x){\displaystyle _0^x}𝑑y\psi _l(y)h(y)f_n(y,\sigma ^i,t)`$
$`+2\psi _2(x){\displaystyle _{\mathrm{}}^x}𝑑y\psi _l(y)f_n(y,\sigma ^i,t)+C_n(\sigma ^i,t)\psi _r(x),`$ (59)
where
$$\psi _l(y)=\mathrm{exp}[(2\alpha 1)y]\psi _r(y).$$
The functions $`C_n,n=1,2,\mathrm{}`$, are utilised to saturate the integrability conditions (32).
The solution of Eq.(26) is given by the formula
$$\varphi _n(s,\sigma ^i,t)=\stackrel{~}{\varphi }_n(sC_0,\sigma ^i,t).$$
## 8 Appendix B. Stability of the interface
Significance of our theoretical analysis of the interfaces depends on their stability with respect to small perturbations. It is sufficient to check the stability of the planar interface because our perturbative solution is based on it. Mathematically, the stability is related to the sign of eigenvalues of certain operator, and it is a model dependent property. Considerations presented below apply to the model defined by formulas (1-3), of course. In this case the linearised evolution equation for small amplitude perturbations $`\delta \varphi `$ of the planar interface $`\varphi _0(s)`$, formula (19), in the comoving reference frame has the form
$$\frac{2l_0^2\gamma }{K}_t\delta \varphi =\widehat{L}\delta \varphi +2l_0^2(_{x^1}^2+_{x^2}^2)\delta \varphi ,$$
where $`\widehat{L}`$ has been given below formula (24), and $`x^1,x^2`$ are the two Cartesian coordinates in the plane of the interface. Because $`\widehat{L}`$ does not depend on $`x^1,x^2,`$ we may pass to Fourier modes $`\stackrel{~}{\delta \varphi }`$ of $`\delta \varphi `$ with respect to these coordinates. After the substitution
$$\stackrel{~}{\delta \varphi }=\mathrm{exp}[(\frac{1}{2}\alpha )s]\mathrm{\Psi }$$
we obtain the following equation
$$\frac{2l_0^2\gamma }{K}_t\mathrm{\Psi }=\widehat{N}\mathrm{\Psi },$$
where
$$\widehat{N}=\frac{1}{2}_s^2+\frac{1}{2}(\alpha +\frac{1}{2})^2+2l_0^2k^22(\alpha +1)\varphi _0+3\varphi _0^2,$$
with $`k^2=(k_1)^2+(k_2)^2,k_1,k_2`$ being the wave numbers Fourier conjugate with $`x^1,x^2`$.
Notice that
$$\mathrm{\Psi }_0=\mathrm{exp}[(\alpha \frac{1}{2})s]\psi _r,$$
where $`\psi _r`$ is the zero-mode given by formula (29), is eigenfunction of the Hermitean operator $`\widehat{N}`$. The corresponding eigenvalue is equal to zero. Because $`\mathrm{\Psi }_0`$ does not vanish at any finite $`s`$, it represents "the ground state" of a fictitious system with $`\widehat{N}`$ as "the Hamiltonian". Hence, all other eigenvalues of $`\widehat{N}`$ are strictly positive. Moreover, looking at "the potential" in "the Hamiltonian" $`\widehat{N}`$ one can see that the zero eigenvalue and the next one are separated by a finite gap. The eigenmode $`\mathrm{\Psi }_0`$ corresponds to a parallel shift of the interface as a whole in the $`x^3`$ direction. All other eigenmodes decay exponentially with a characteristic time equal to $`t_c=2l_0^2\gamma /(K\lambda )`$, where $`\lambda `$ denotes the corresponding eigenvalue. Thus, the planar interface in our model is stable. Consequently, also the curved interfaces are stable with respect to small perturbations provided that their curvature radia are large enough for validity of our perturbative expansion.
Let us point out that from point of view of applications in condensed matter physics even unstable interfaces can be interesting if unstable modes grow in time so slowly that the interface manages to travel across the sample before these modes become visible. The interface has finite normal velocity $`\stackrel{}{p}\dot{\stackrel{}{x}}_0`$, formula (46), and in any real experiment the sample occupies a finite volume.
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# Finite size effects on the phase diagram of a binary mixture confined between competing walls
## I Introduction
Although the finite size effects on phase transitions in thin films have been studied since a long time , only during the last decade it was discovered that in ferromagnetic Ising films with surface fields of different sign but of the same strength $`\pm H_1`$ (”competing walls”) novel types of phase transitions occur: namely, a phase transition occurs for zero bulk field $`H`$ from a state with an interface freely fluctuating in the center of the thin film to a state where the interface is bound either to the lower or to the upper wall confining the film, which then acquires a nonzero (positive or negative) magnetization. This interface localization-delocalization transition at $`T_c(D)`$ may be either second or first order , and for film thicknesses $`D\mathrm{}`$ the transition temperature $`T_c(D)`$ does not converge towards the bulk critical temperature $`T_{\text{cb}}`$ as usual , but rather towards the wetting transition temperature $`T_\text{w}(H_1)`$ .
Now it is well-known that, in general, wetting transitions can be second or first order . Thus it is plausible that also the interface localization-delocalization transition can be either second or first order. However, recently it was shown that also in cases where the wetting transition is first order, the transition $`T_c(D)`$ may be second order for small enough thickness $`(D<D_t)`$ and become first order only for $`D>D_t`$. Thus the transition at $`T_c(D_t)`$ is the finite-thickness analog of a wetting tricritical point .
All this work has only considered the case $`H=0`$, however. It is well known of course, that for $`D\mathrm{}`$, in the case of a first-order wetting transition at $`T_w(H_1<0)`$ there exists at $`T=T_{\text{pre}}(H,H_1)`$ a prewetting transition , where the distance of the interface bound to the wall discontinuously jumps from a smaller value to a larger value. While the analog of this prewetting transition in thin films has been studied occasionally for the case of ”capillary condensation” (where on both walls fields act that have the same sign) , it is only in the present work that the effect of prewetting phenomena on the interface localization-delocalization transition is considered . The physical systems that we have in mind are not magnets, of course, but rather binary (A,B) mixtures: as is well known, in an Ising model context the ”magnetization” simply translates into the relative concentration $`\varphi `$ if one component (A, say), and the field $`H`$ translates into the chemical potential difference $`\mathrm{\Delta }\mu `$ between the species (for simplicity we deal here with perfectly symmetric mixtures for which the bulk critical concentration is $`\varphi _{\text{cb}}=0.5`$).
However, one important aspect of binary mixtures is that physically it is a density of an extensive thermodynamic variable (namely $`\varphi `$) that is the fixed independent thermodynamic variable, rather than the intensive variable $`\mathrm{\Delta }\mu `$. As we shall see below, this fact has important consequences for the phase diagram of confined mixtures: the typical situation is that one encounters two successive lateral phase separation transitions!
In Section 2 we elaborate these ideas by a qualitative discussion of the phase diagrams {both in the space of intensive variables ($`\mathrm{\Delta }\mu ,T`$) and in the space ($`\varphi ,T`$)} and of the corresponding physical state of the confined mixture. Section 3 exemplifies these considerations by presenting a specific calculation for a symmetric polymer mixture, treated within a self-consistent field framework . There are numerous experimental studies of confined polymer mixtures and these systems might be convenient for testing our predictions. Finally section 4 summarizes our conclusions.
## II Qualitative phase diagrams of confined binary mixtures
We assume here a binary mixture confined by ”competing” walls in the sense that one wall attracts species A with the same strength as the opposite wall attracts species B, and consider the case of first-order wetting. Then the topology of the phase diagram can be estimated from the qualitative considerations as shown in the left part of Fig. 1: In the space ($`T,\mathrm{\Delta }\mu `$), bulk ($`D\mathrm{}`$) phase separation occurs for $`\mathrm{\Delta }\mu =0`$, and the walls are incompletely wet for $`T<T_w`$ but wet for $`T_w<T<T_{cb}`$. From the point $`T=T_w,\mathrm{\Delta }\mu =0`$ there extend two (first-order) prewetting lines, which end at the prewetting critical points $`T_{cp}`$. These prewetting transitions correspond to singularities of the surface free energies associated with the lower and upper wall confining the mixture (Fig. 2). Due to the special symmetries chosen for our model, both the wetting transitions for the lower and upper wall coincide, and the prewetting critical temperatures are also the same. There is a mirror symmetry of the phase diagrams around the line $`\mathrm{\Delta }\mu =0`$ (upper part) or $`\varphi =0.5`$ (lower part), respectively.
For finite thickness $`D`$ it may happen, as demonstrated by Monte Carlo simulations for Ising models with enhanced exchange interactions near the walls , that a tricritical point at $`D=D_t`$, $`T_c(D_t)`$, $`\mathrm{\Delta }\mu =0`$ (and $`\varphi =0.5`$) occurs. For $`D<D_t`$ then there exists a single critical point at $`T=T_c(D)`$, $`\mathrm{\Delta }\mu =0`$ (and $`\varphi =1/2`$), there is no remnant of the prewetting phenomena left, and the phase diagram both in the ($`T,\mathrm{\Delta }\mu `$) plane as well as in the ($`T,\varphi `$) plane looks qualitatively exactly like in the bulk three-dimensional case. Of course, we do expect a flatter shape near the critical point due to the occurence of the two-dimensional Ising exponents {$`\varphi _{\text{coex}}1/2(1T/T_c(D))^{1/8}`$ rather than $`\varphi _{\text{coex}}1/2(1T/T_{\text{cb}})^\beta `$ with $`\beta 0.325`$ in the bulk }. But the situation differs very much for $`D>D_t`$ (middle and right part of Fig. 1). In the semi-grand-canonical ensemble ($`\mathrm{\Delta }\mu `$ fixed), one experiences a single first order transition for $`\mathrm{\Delta }\mu _c(D)<\mathrm{\Delta }\mu <\mathrm{\Delta }\mu _c(D)`$ where $`\pm \mathrm{\Delta }\mu _c(D)`$ is the chemical potential reached for $`T=T_c(D)`$ along the remnants of the prewetting lines. In the canonical ensemble $`(\varphi `$ fixed), we encounter a single first-order phase transition only if $`\varphi =\varphi _{\text{trip}}=\frac{1}{2}`$, which corresponds to $`\mathrm{\Delta }\mu =0`$, while for $`\varphi _1<\varphi _{\text{trip}}`$ or for $`\varphi _{\text{trip}}<\varphi _2=1\varphi _1`$ one encounters two successive first-order transitions when one lowers the temperature. Only if one chooses $`\varphi =\varphi _{c1}`$, (the critical concentration corresponding to the critical point $`T_{c1}`$) or $`\varphi =\varphi _{c2}=1\varphi _{c1}`$ (corresponding to $`T_{c2}`$) one still encounters a (two-dimensional) critical behavior of phase separation, $`\varphi _{\text{coex}}\varphi _{c1}(1T/T_{c1})^{1/8}`$, and similarly $`\varphi _{\text{coex}}\varphi _{c2}(1T/T_{c2})^{1/8}`$. While $`\varphi _1`$ as well as $`\varphi _2`$ merge at $`\varphi _{\text{trip}}=1/2`$ as $`DD_t`$, $`\varphi _1`$ and $`\varphi _2`$ move outwards towards the prewetting critical concentration when $`D\mathrm{}`$. In this limit, the width of the two-phase coexistence regions for $`T>T_{\text{trip}}`$ must shrink and ultimately vanish, since the difference in concentration on both sides of the pseudo-prewetting first order lines with respect to the average concentration $`\varphi `$ of the film is an effect of order $`1/D`$, and therefore for $`D\mathrm{}`$ the prewetting transitions are lines in the $`(T,\varphi )`$ phase diagram and not split up in two-phase coexistence regions.
From this description it already is clear that the approach to the phase diagram of the bulk $`(D=\mathrm{})`$ as $`D\mathrm{}`$ is very nonuniform: for any finite $`D`$ the bulk transition is still rounded, and (for $`\varphi _1<\varphi _{\text{trip}}`$ or $`\varphi _{\text{trip}}>\varphi _1`$) the first transition is a lateral phase separation corresponding to the prewetting transition and the second transition is another lateral phase separation at $`T_{\text{trip}}`$ (Fig. 2). For $`D\mathrm{}`$ the phase diagram, hence, contains prewetting lines (as in the left part of Fig. 1) and a horizontal line at $`T_w`$ (to which the triple line in the middle part of Fig. 1 has converged). Of course, the pictures explaining the various phases in Fig. 2 are highly schematic, and in reality one expects the interface to turn around rather smoothly and avoid the 90 kinks. Such smooth interfaces (which also have an intrinsic width which need not be negligibly small in comparison with $`D`$) have in fact been observed in two-phase coexistence states associated with capillary condensation .
## III A quantitative example: a confined polymer mixture
Thin polymeric films confined between walls may have interesting applications and can also be studied conveniently by a variety of experimental tools . In fact, the ”soft mode” – phase with a single fluctuating interface in the middle of the film (as shown in the upper part of Fig. 2) has been experimentally observed , and we consider it likely that by fine-tuning of experimental control parameters it should also be possible to observe some of the transitions predicted in Figs. 1, 2. In fact, for polymers one need not have special interactions to get a first-order wetting transition as in the Ising model , rather one finds always first order wetting behavior except close to the critical point of the bulk .
As in our previous work we consider a situation where the wetting transition temperature $`T_w`$ lies in the strong segregation regime, for which the self-consistent field theory is accurate. The technical aspects of this approach are explained in detail elsewhere . Fig. 3 shows that for a typical choice of parameters indeed a phase diagram of the type of Fig. 1, right part, is reproduced. Note that the self-consistent field theory near the critical points $`T_{c1}(D)`$, $`T_{c2}(D)`$ implies mean-field like behavior, $`\varphi _{\text{coex}}\varphi _{c1}(1T/T_{c1}(D))^{1/2}`$, rather than yielding the expected two-dimensional Ising exponent . But for large molecular weight this Ising like critical region is expected to be rather narrow , and thus we consider Fig. 3 as a useful hint for the phase diagram to be searched for in the experiments.
## IV Conclusions
In this paper we have considered the problem of phase-separating binary mixtures confined between ”competing walls” and have shown by qualitative considerations (Fig.1.) and self–consistent field calculations (Fig.3.) that the phase diagram has either critical points and first order regions coexisting at a triple line or a single critical point resulting from the merging of these two critical points at the tricritical thickness $`D_t`$. In previous work treating the case $`D>D_t`$, only the case $`\mathrm{\Delta }\mu =0`$ in the semi-grandcanonical ensemble was studied , which in the $`(T,\varphi )`$ plane means that one cools the system at $`\varphi =\varphi _{\text{trip}}=0.5`$ and then a single first order transition (Fig. 2, left part) occurs: thus the existence of the two critical points was not previously discussed.
Of course, in reality one will have to abandon the special symmetry assumptions used in Figs. 1-3, allowing for asymmetric mixtures, differences in strength of the wall forces, etc, and thus the space of parameters to be considered gets much enlarged. However, as long as one works in the subspace where the wetting transition temperatures $`T_w`$ of both walls are the same, the phase diagrams still should have the topology of Fig. 1, only the mirror symmetry around $`\mathrm{\Delta }\mu =\mathrm{\Delta }\mu _{\text{coex}}(T)`$ or $`\varphi =\varphi _{\text{cb}}`$ is lost, and thus in general $`\varphi _{\text{trip}}`$ will differ from $`\varphi _{\text{cb}}`$. Also $`T_{c1}`$ and $`T_{c2}`$ will differ. Of course, in the most general case one must allow also for $`T_{w1}T_{w2}`$, different wetting transition temperatures of both walls. One can also consider first order wetting at one wall and second order wetting at the other. A description of the phase diagrams for these more complicated cases is a challenging task for future work.
Acknowledgements: This work was partially supported by the DFG under grant N$`^\text{o}`$ Bi314/17 and by the Volkswagenstiftung under grant N$`^\text{o}`$ I/74168.
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# Quantum Statistics of Hydrogen in Strong Magnetic Fields
## I Introduction
The recent discovery of magnetars has renewed interest in the behavior of charged particle systems in the presence of extremely strong external magnetic fields . In this new type of neutron stars, electrons and protons from decaying neutrons produce magnetic fields $`B`$ reaching up to $`10^{15}\mathrm{G}`$, much larger than those in neutron stars and white dwarfs, where $`B`$ is of order $`10^{10}10^{12}\mathrm{G}`$ and $`10^610^8\mathrm{G}`$, respectively.
Analytic treatments of the strong-field properties of an atomic system are difficult, even in the zero-temperature limit. The reason is the logarithmic asymptotic behavior of the ground state energy . In the weak-field limit, perturbative approaches yield well-known series expansions in powers of $`B^2`$. These are useful, however, only for $`BB_0`$, where $`B_0`$ is the atomic magnetic field strength $`B_0=e^3M^2/\mathrm{}^32.35\times 10^5\mathrm{T}=2.35\times 10^9\mathrm{G}`$.
So far, the most reliable values for strong uniform fields were obtained by numerical calculations . An analytic mapping procedure was introduced in Ref. to interpolate between the weak- and strong-field behavior, and a variational approximation was given in Ref. , both with quite good results.
In this note, we use an extension of the Feynman-Kleinert variational approach to find a single analytic approximation to the effective classical potential of the system for all temperatures and magnetic field strengths. From this, the quantum statistical partition function can be obtained by a simple configuration space integral over a classical-looking Boltzmann-factor. In the zero-temperature limit, the effective classical potential is the ground state energy of the system.
## II Effective Classical Potential
The Hamiltonian of the electron in a hydrogen atom in the presence of a uniform external magnetic field pointing along the positive $`z`$-axis is
$$H(𝐩,𝐱)=\frac{1}{2M}𝐩^2\frac{1}{2}\omega _cl_z(𝐩,𝐱)+\frac{1}{8}\omega _c^2𝐱^2\frac{e^2}{|𝐱|}.$$
(1)
Here we have used the symmetric gauge $`𝐀(𝐱)=(B/2)(y,x,0)`$, and denoted the $`z`$-component of the orbital angular momentum by $`l_z(𝐩,𝐱)=(𝐱\times 𝐩)_z`$. The quantum statistical partition function can always be expressed as a classical-looking configuration space integral
$$Z=\frac{d^3x_0}{\lambda _{\mathrm{th}}^3}\mathrm{exp}\left[\beta V_{\mathrm{eff}}(𝐱_0)\right],$$
(2)
where $`\lambda _{\mathrm{th}}=\sqrt{2\pi \mathrm{}^2\beta /M}`$ is the thermal wavelength, $`\beta =1/k_BT`$ is the inverse temperature, and $`V_{\mathrm{eff}}(𝐱_0)`$ is the effective classical potential $`V_{\mathrm{eff}}(𝐱_0)`$. Generalizing the development in Ref. , this is defined by the phase space path integral
$$\mathrm{exp}\left[\beta V_{\mathrm{eff}}(𝐱_0)\right]\lambda _{\mathrm{th}}^3d^3p_0𝒟^3x𝒟^3p\delta (𝐱_0\overline{𝐱(\tau )})\delta (𝐩_0\overline{𝐩(\tau )})e^{𝒜[𝐩,𝐱]/\mathrm{}},$$
(3)
where $`𝒜[𝐩,𝐱]`$ is the Euclidean action
$$𝒜[𝐩,𝐱]=_0^\mathrm{}\beta 𝑑\tau [i𝐩(\tau )\dot{𝐱}(\tau )+H(𝐩(\tau ),𝐱(\tau ))],$$
(4)
and $`\overline{𝐱(\tau )}=_0^\mathrm{}\beta 𝑑\tau 𝐱(\tau )/\mathrm{}\beta `$ and $`\overline{𝐩(\tau )}=_0^\mathrm{}\beta 𝑑\tau 𝐩(\tau )/\mathrm{}\beta `$ are the temporal averages of position and momentum. The special treatment of $`𝐱_0`$ and $`𝐩_0`$ is necessary, since the classical harmonic fluctuation widths $`𝐱^2^{\mathrm{cl}}`$ and $`𝐩^2^{\mathrm{cl}}`$ are proportional to the temperature $`T`$ (Dulong-Petit law). Thus they diverge for $`T\mathrm{}`$ and their fluctuations cannot be treated pertubatively. In contrast, the fluctuation widths $`(𝐱𝐱_0)^2`$, $`(𝐩𝐩_0)^2`$ around $`𝐱_0`$ and $`𝐩_0`$ go to zero for large $`T`$ and are limited down to $`T=0`$, thus allowing for a treatment by variational perturbation theory . For this we rewrite the action (4) as
$$𝒜[𝐩,𝐱]=𝒜_𝛀^{𝐩_0,𝐱_0}[𝐩,𝐱]+𝒜_{\mathrm{int}}[𝐩,𝐱],$$
(5)
with a harmonic trial action
$`𝒜_𝛀^{𝐩_0,𝐱_0}[𝐩,𝐱]={\displaystyle _0^\mathrm{}\beta }d\tau \{`$ $`i[𝐩(\tau )𝐩_0]\dot{𝐱}(\tau )+{\displaystyle \frac{1}{2M}}[𝐩(\tau )𝐩_0]^2+{\displaystyle \frac{1}{2}}\mathrm{\Omega }_1l_z(𝐩(\tau )𝐩_0,𝐱(\tau )𝐱_0)`$ (7)
$`+{\displaystyle \frac{1}{8}}M\mathrm{\Omega }_2^2[𝐱^{}(\tau )𝐱_0^{}]^2+{\displaystyle \frac{1}{2}}M\mathrm{\Omega }_{}^2[z(\tau )z_0]^2\},`$
in which $`𝐱^{}=(x,y)`$ denotes the transverse part of $`𝐱`$. The frequencies $`𝛀=(\mathrm{\Omega }_1,\mathrm{\Omega }_2,\mathrm{\Omega }_{})`$ are arbitrary for the moment. Inserting the decomposition (5) into (3), we expand the exponential of the interaction, $`\mathrm{exp}\left\{𝒜_{\mathrm{int}}[𝐩,𝐱]/\mathrm{}\right\}`$, yielding a series of expectation values of powers of the interaction
$$𝒜_{\mathrm{int}}^n[𝐩,𝐱]_𝛀^{𝐩_0,𝐱_0}=\frac{(2\pi \mathrm{})^3}{Z_𝛀^{𝐩_0,𝐱_0}}𝒟^3x𝒟^3p𝒜_{\mathrm{int}}^n[𝐩,𝐱]\delta (𝐱_0\overline{𝐱(\tau )})\delta (𝐩_0\overline{𝐩(\tau )})\mathrm{exp}\left\{\frac{1}{\mathrm{}}𝒜_𝛀^{𝐩_0,𝐱_0}[𝐩,𝐱]\right\}.$$
(8)
The path integral over the Boltzmann-factor involving the harmonic action (7) is exactly solvable and yields the restricted partition function
$$Z_𝛀^{𝐩_0,𝐱_0}=\frac{\mathrm{}\beta \mathrm{\Omega }_+/2}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_+/2}\frac{\mathrm{}\beta \mathrm{\Omega }_{}/2}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_{}/2}\frac{\mathrm{}\beta \mathrm{\Omega }_{}/2}{\mathrm{sinh}\mathrm{}\beta \mathrm{\Omega }_{}/2},$$
(9)
where $`\mathrm{\Omega }_\pm |\mathrm{\Omega }_1\pm \mathrm{\Omega }_2|/2`$. Rewriting the perturbation series as a cumulant expansion, evaluating the expectation values, and integrating out the momenta on the right-hand side of Eq. (3) leads to a series representation for the effective classical potential $`V_{\mathrm{eff}}(𝐱_0)`$. Since it is impossible to sum up the series, the perturbation expansion must be truncated, leading to an $`N`$th-order approximation $`W_𝛀^{(N)}(𝐱_0)`$ for the effective classical potential. Since the parameters $`𝛀`$ are arbitrary, $`W_𝛀^{(N)}(𝐱_0)`$ should depend minimally on $`𝛀`$. This determines the optimal values $`𝛀^{(N)}=(\mathrm{\Omega }_1^{(N)}(𝐱_0),\mathrm{\Omega }_2^{(N)}(𝐱_0),\mathrm{\Omega }_{}^{(N)}(𝐱_0))`$ of $`N`$th order. Reinserting these into $`W_𝛀^{(N)}(𝐱_0)`$ yields the optimal approximation $`W^{(N)}(𝐱_0)W_{𝛀^{(N)}}^{(N)}(𝐱_0)`$.
The first-order approximation to the effective classical potential is
$$W_𝛀^{(1)}(𝐱_0)=\frac{1}{\beta }\mathrm{ln}Z_𝛀^{𝐩_0,𝐱_0}+(\omega _c\mathrm{\Omega }_1)b_{}^2(𝐱_0)\frac{1}{4}\left(\mathrm{\Omega }_2^2\omega _c^2\right)a_{}^2(𝐱_0)\frac{1}{2}M\mathrm{\Omega }_{}^2a_{}^2(𝐱_0)\frac{e^2}{|𝐱|}_𝛀^{𝐩_0,𝐱_0},$$
(10)
where the quantities $`a_{}^2(𝐱_0)`$, $`b_{}^2(𝐱_0)`$, and $`a_{}^2(𝐱_0)`$ are the transverse and longitudinal fluctuation widths
$$a_{}^2(𝐱_0)=x^2(\tau )_𝛀^{𝐩_0,𝐱_0},a_{}^2(𝐱_0)=z^2(\tau )_𝛀^{𝐩_0,𝐱_0},b_{}^2(𝐱_0)=x(\tau )p_y(\tau )_𝛀^{𝐩_0,𝐱_0}.$$
(11)
The expectation value of the Coulomb potential on the right-hand side of Eq. (10) has the integral representation
$`{\displaystyle \frac{e^2}{|𝐱|}}_𝛀^{𝐩_0,𝐱_0}`$ $`=`$ $`e^2\sqrt{{\displaystyle \frac{2}{\pi }}a_{}^2(𝐱_0)}{\displaystyle \underset{0}{\overset{1}{}}}{\displaystyle \frac{d\xi }{a_{}^2(𝐱_0)+\xi ^2[a_{}^2(𝐱_0)a_{}^2(𝐱_0)]}}`$ (13)
$`\times \mathrm{exp}\left\{{\displaystyle \frac{\xi ^2}{2}}\left({\displaystyle \frac{x_0^2+y_0^2}{a_{}^2(𝐱_0)+\xi ^2[a_{}^2(𝐱_0)a_{}^2(𝐱_0)]}}+{\displaystyle \frac{z_0^2}{a_{}^2(𝐱_0)}}\right)\right\}.`$
The variational energy (10) is minimized at each $`𝐱_0`$, and the resulting $`W^{(N)}(𝐱_0)`$ is displayed for a low temperature and different magnetic fields in Fig. 1.
From now on we set $`\mathrm{}=e^2=k_B=c=M=1`$. Thus, energies are measured in units of $`ϵ_0=Me^4/\mathrm{}^22\mathrm{Ryd}27.21\mathrm{eV}`$, temperatures in $`ϵ_0/k_B3.16\times 10^5\mathrm{K}`$, distances in Bohr radii $`a_B=\mathrm{}^2/Me^20.53\times 10^8\mathrm{cm}`$, and magnetic field strengths in $`B_0=e^3M^2/\mathrm{}^32.35\times 10^5\mathrm{T}=2.35\times 10^9\mathrm{G}`$.
The curves $`W^{(1)}(𝐱_0)`$ are plotted to show their anisotropy with respect to the magnetic field direction. The anisotropy grows when lowering the temperature and increasing the field strength. Far away from the proton at the origin, the potential becomes isotropic, due to the decreasing influence of the Coulomb interaction. Analytically, this is seen by going to the limits $`\rho _0\mathrm{}`$ or $`z_0\mathrm{}`$, where the expectation value of the Coulomb potential (13) tends to zero, leaving an effective classical potential
$$W_𝛀^{(1)}(𝐱_0)\frac{1}{\beta }\mathrm{ln}Z_𝛀^{𝐩_0,𝐱_0}+(\omega _c\mathrm{\Omega }_1)b_{}^2\frac{1}{4}\left(\mathrm{\Omega }_2^2\omega _c^2\right)a_{}^2\frac{1}{2}M\mathrm{\Omega }_{}^2a_{}^2.$$
(14)
This is $`𝐱_0`$-independent, and optimization yields the constants $`\mathrm{\Omega }_1^{(1)}=\mathrm{\Omega }_2^{(1)}=\omega _c`$ and $`\mathrm{\Omega }_{}^{(1)}=0`$, with the asymptotic energy
$$W^{(1)}(𝐱_0)\frac{1}{\beta }\mathrm{ln}\frac{\mathrm{}\beta \omega _c/2}{\mathrm{sinh}\mathrm{}\beta \omega _c/2}.$$
(15)
The $`B=0`$ -curves are, of course, identical with those obtained from variational perturbation theory for the hydrogen atom .
For large temperatures, the anisotropy decreases since the violent thermal fluctuations have a smaller preference of the $`z`$-direction.
## III Zero-Temperature Limit
At zero-temperature, the first-order effective classical potential (10) at the origin yields an approximation for the ground state energy of the hydrogen atom in a uniform magnetic field: $`E_𝛀^{(1)}=lim_\beta \mathrm{}W_𝛀^{(1)}(\mathrm{𝟎})`$:
$$E_𝛀^{(1)}(B)=\frac{1}{4\mathrm{\Omega }_2}\left(\mathrm{\Omega }_2^2+\omega _c^2\right)+\frac{\mathrm{\Omega }_{}}{4}\frac{1}{|𝐱|}_𝛀^\mathrm{𝟎},$$
(16)
with the expectation value for the Coulomb potential
$$\frac{1}{|𝐱|}_𝛀^\mathrm{𝟎}=\frac{2}{\sqrt{\pi }}\times \{\begin{array}{cc}\sqrt{\frac{\mathrm{\Omega }_{}\mathrm{\Omega }_2}{\mathrm{\Omega }_{}\mathrm{\Omega }_2}}\mathrm{arctan}\sqrt{\frac{2\mathrm{\Omega }_{}}{\mathrm{\Omega }_2}1},& 2\mathrm{\Omega }_{}>\mathrm{\Omega }_2,\hfill \\ \sqrt{\mathrm{\Omega }_{}},& 2\mathrm{\Omega }_{}=\mathrm{\Omega }_2,\hfill \\ \frac{1}{2i}\sqrt{\frac{\mathrm{\Omega }_{}\mathrm{\Omega }_2}{\mathrm{\Omega }_{}\mathrm{\Omega }_2}}\mathrm{ln}\frac{1+i\sqrt{2\mathrm{\Omega }_{}/\mathrm{\Omega }_21}}{1i\sqrt{2\mathrm{\Omega }_{}/\mathrm{\Omega }_21}},& 2\mathrm{\Omega }_{}<\mathrm{\Omega }_2.\hfill \end{array}$$
(17)
Equations (16) and (17) are independent of the frequency parameter $`\mathrm{\Omega }_1`$, such that optimization of the ground state energy (16) is ensured by minimization. Reinserting the extremal $`\mathrm{\Omega }_2^{(1)}`$ and $`\mathrm{\Omega }_{}^{(1)}`$ into Eq. (16) yields the first-order approximation to the ground state energy $`E^{(1)}(B)`$. In the absence of the Coulomb interaction the optimization with respect to $`\mathrm{\Omega }_2`$ yields $`\mathrm{\Omega }_2^{(1)}=\omega _c`$, rendering the ground state energy $`E^{(1)}(B)=\omega _c/2`$, which is the zeroth Landau level in this special case. The trial frequency $`\mathrm{\Omega }_{}`$ must be set equal to zero to preserve translational symmetry along the $`z`$-axis.
In the opposite limit of a vanishing magnetic field, Eq. (16) coincides with the first-order variational result for the ground state energy of the hydrogen atom, whose optimization gave $`E^{(1)}(B=0)=4/3\pi 0.4244[2\mathrm{Ryd}]`$ obtained in Refs. . In Ref. , the $`B=0`$ -system was treated up to third order leading to the much more accurate result $`E^{(1)}(B=0)0.490[2\mathrm{Ryd}]`$, very near the exact value $`E^{\mathrm{ex}}(B=0)=0.5[2\mathrm{Ryd}]`$.
Let us investigate the asymptotics in the strong-field limit $`B\mathrm{}`$. The $`B`$-dependence of the binding energy
$$\epsilon (B)=\frac{B}{2}E$$
(18)
is plotted in Fig. 2, where it is compared with the results of Ref. , with satisfactory agreement. Our results are of similar accuracy as those of other first-order calculations, for example those from the operator optimization method in first order of Ref. . The advantage of variational perturbation theory is that it yields good results for all magnetic field strengths. From our experience with the fast convergence of the method \[9, Chaps. 5,9\], higher orders of variational perturbation theory will push the approximations rapidly towards the exact value.
### A Weak-Field Behavior
The calculations of the binding energy for weak magnetic fields show that the ratio $`\eta 2\mathrm{\Omega }_{}/\mathrm{\Omega }_2`$ is always smaller than one if $`B0`$. Setting $`\mathrm{\Omega }\mathrm{\Omega }_2`$, we rewrite the binding energy as a function of $`\mathrm{\Omega }`$ and $`\eta `$:
$$\epsilon _{\eta ,\mathrm{\Omega }}^{(1)}(B)\frac{B}{2}\frac{\mathrm{\Omega }}{4}\left(1+\frac{\eta }{2}\right)\frac{B^2}{4\mathrm{\Omega }}\sqrt{\frac{\eta \mathrm{\Omega }}{2\pi }}\frac{1}{\sqrt{1\eta }}\mathrm{ln}\frac{1\sqrt{1\eta }}{1+\sqrt{1\eta }}.$$
(19)
This is minimized with respect to the new variational parameters $`\eta `$ and $`\mathrm{\Omega }`$ by expanding $`\eta (B)`$ and $`\mathrm{\Omega }(B)`$ in powers of $`B^2`$ with unknown coefficients, and inserting these expansions into extremality equations. The expansion coefficients are then determined order by order. The optimal expansions are inserted into (19), yielding the optimized binding energy $`\epsilon ^{(1)}(B)`$ as a power series
$$\epsilon ^{(1)}(B)=\frac{B}{2}\underset{n=0}{\overset{\mathrm{}}{}}\epsilon _nB^{2n}.$$
(20)
The coefficients $`\epsilon _n`$ are listed in Table I and compared with the exact ones of Ref. . Of course, the higher-order coefficients of this first-order variational approximation become rapidly inaccurate, but the results can be improved, if desired, by going to higher orders in variational perturbation theory as in Ref. \[9, Chaps. 5,9\].
### B Strong-Field Behavior
In the discussion of the pure magnetic field we have mentioned that the variational calculation for the ground state energy, which is associated with the zeroth Landau level, yields a frequency $`\mathrm{\Omega }_2B`$, while $`\mathrm{\Omega }_{}=0`$. We therefore use the assumptions $`\mathrm{\Omega }_{}\mathrm{\Omega }_22\mathrm{\Omega }_{}`$ and $`\mathrm{\Omega }_{}B`$ for an analytic study of the strong-field behavior of the ground state energy (16). We expand the last expression of the expectation value (17) in terms of $`2\mathrm{\Omega }_{}/\mathrm{\Omega }_{}`$, and reinsert this expansion into (16). Then we omit all terms proportional to $`C/\mathrm{\Omega }_{}`$, where $`C`$ stands for any expression with a value much smaller than the field strength $`B`$. We thus obtain the strong-field approximation for the first-order binding energy (18)
$$\epsilon _{\mathrm{\Omega }_{},\mathrm{\Omega }_{}}^{(1)}=\frac{B}{2}\left(\frac{\mathrm{\Omega }_{}}{4}+\frac{B^2}{4\mathrm{\Omega }_{}}+\frac{\mathrm{\Omega }_{}}{4}+\sqrt{\frac{\mathrm{\Omega }_{}}{\pi }}\mathrm{ln}\frac{\mathrm{\Omega }_{}}{2\mathrm{\Omega }_{}}\right).$$
(21)
Determining $`\mathrm{\Omega }_{}`$, $`\mathrm{\Omega }_{}`$ by minimization, we obtain
$$\mathrm{\Omega }_{}B,\sqrt{\mathrm{\Omega }_{}^{(3)}}=\frac{2}{\sqrt{\pi }}\left(\mathrm{ln}B2\mathrm{l}\mathrm{n}\mathrm{l}\mathrm{n}B+\frac{2a}{\mathrm{ln}B}+\frac{a^2}{\mathrm{ln}^2B}+b\right)+𝒪(\mathrm{ln}^3B)$$
(22)
with abbreviations $`a=2\mathrm{ln}\mathrm{\hspace{0.17em}2}1.307`$ and $`b=\mathrm{ln}(\pi /2)21.548`$. Thus the optimized binding energy can be written up to the order $`\mathrm{ln}^2B`$:
$`\epsilon ^{(1)}(B)`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\left\{\mathrm{ln}^2B4\mathrm{ln}B\mathrm{lnln}B+4\mathrm{ln}^2\mathrm{ln}B4b\mathrm{lnln}B+2(b+2)\mathrm{ln}B+b^2{\displaystyle \frac{1}{\mathrm{ln}B}}\left[8\mathrm{ln}^2\mathrm{ln}B8b\mathrm{lnln}B+2b^2\right]\right\}`$ (24)
$`+𝒪(\mathrm{ln}^2B).`$
Note that the prefactor $`1/\pi `$ of the leading $`\mathrm{ln}^2B`$-term differs from a value $`1/2`$ obtained by Landau and Lifschitz . Our value is a consequence of the harmonic trial system. The calculation of higher orders in variational perturbation theory should drive our value towards $`1/2`$.
The convergence of the expansion (24) is quite slow. At a magnetic field strength $`B=10^5B_0`$, which corresponds to $`2.35\times 10^{10}\mathrm{T}=2.35\times 10^{14}\mathrm{G}`$, the contribution from the first six terms is $`22.87[2\mathrm{Ryd}]`$. The next three terms suppressed by a factor $`\mathrm{ln}^1B`$ contribute $`2.29[2\mathrm{Ryd}]`$, while an estimate for the $`\mathrm{ln}^2B`$-terms yields nearly $`0.3[2\mathrm{Ryd}]`$. Thus we find $`\epsilon ^{(1)}(10^5)=20.58\pm 0.3[2\mathrm{Ryd}]`$. This is in very good agreement with the value $`20.60[2\mathrm{Ryd}]`$ obtained from the full treatment described in Sec. III.
Table II lists the values of the first six terms of Eq. (24). This shows in particular the significance of the second term in (24), which is of the same order of the leading first term, but with an opposite sign. In Fig. 2, we have plotted the expression $`\epsilon _L(B)=(1/2)\mathrm{ln}^2B`$ of Landau and Lifschitz to illustrate that it gives far too large binding energies even at very large magnetic fields, e.g. at $`2000B_010^{12}\mathrm{G}`$. Obviously, the nonleading terms in Eq. (24) give important contributions to the asymptotic behavior even at such large magnetic fields. As an peculiar property of the asymptotic behavior, the absolute value of the difference between the Landau-expression $`\epsilon _L(B)`$ and our approximation (24) diverges with increasing magnetic field strengths $`B`$. Only the relative difference decreases.
## IV Summary
We have calculated the effective classical potential for the hydrogen atom in constant magnetic field, which governs the statistical mechanics of the system at all temperatures. At zero temperature, we find a rather accurate ground state energy which interpolates very well between weak and strong fields.
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# HST revisits the proper motion of PSR B0656+141footnote 11footnote 1Based on observations with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by AURA, Inc., under NASA contract NAS 5-26555.
## 1 Introduction
PSR B0656+14 is a middle-aged ($`100,000`$ yrs) isolated neutron star, detected as a radio (385 ms) as well as an X-ray pulsar by ROSAT (Finley et al., 1992), which confirmed the X-ray pulsation hinted in the discovery data of the EINSTEIN satellite (Cordova et al., 1989). A possible identification of the pulsar in gamma-rays was also claimed (Ramanamurthy et al., 1996) but so far only a marginal evidence of pulsations was found in the GRO/EGRET data. Although dispersion measure estimates put PSR B0656+14 at $``$ 760 pc, its true distance remains highly uncertain. Indeed, assuming for the neutron star a canonical radius of 10 km, blackbody fits to the soft thermal component of its X-ray spectrum (Finley et al., 1992) impose a significant downward revision of the nominal distance. PSR B0656+14 could thus join the thin crowd of pulsars detected within few hundred pc of the Sun. Since PSR B0656+14 is inside the putative supernova remnant known as Monogem Ring, the measurement of its proper motion was a natural way to investigate their possible association. Thus, 3 radio images of the field, covering roughly 4 years, were collected by Thompson and Cordova (1994) from the NRAO Very Large Array yielding a proper motion value of $`\mu _\alpha =64\pm 11`$ mas/yr; $`\mu _\delta =28\pm 4`$ mas/yr and ruling out the possible association. The proper motion value was later revised by Pavlov, Stringfellow and Cordova (1996) in $`\mu _\alpha =73\pm 20`$ mas/yr and $`\mu _\delta =26\pm 13`$ mas/yr (PA=$`109^{}\pm 10^{}`$), which, although rather uncertain, shall be taken as the radio “reference” value. In the optical, a likely counterpart to PSR B0656+14 ( m$`{}_{\mathrm{V}}{}^{}`$ 25) was proposed by Caraveo et al. (1994) on the basis of its positional coincidence with the radio pulsar. For the possible range of the pulsar distance ($`800÷250`$ pc), the optical emission of the candidate counterpart was found to be well above the predictions based on both the magnetospheric emission model of Pacini and Salvati (1987) and the extrapolation of the soft X-ray blackbody emission (Finley et al., 1992). Indeed, multicolor photometry, reported by Pavlov et al. (1997), gave support to this identification by showing a composite spectral shape, which could be interpreted as a magnetospheric component superimposed to a thermal one. To assess the reliability of the proposed optical identification, Mignani et al. (1997) tried to use proper motion as a distinctive character of the pulsar. The original ESO/3.6m and NTT observations of Caraveo et al. (1994) were compared to a newly acquired HST/WFPC2 one to search for the object’s displacement. At that time, the limited positional accuracy achievable with the ground based images hampered the astrometric measurements and no conclusive result could be obtained. Soon thereafter, optical pulsations at the radio/X-ray period were detected, at a 4 $`\sigma `$ level, from the proposed counterpart (Shearer et al., 1998), thus apparently closing the identification issue. However, the limited significance of this detection calls for a clear confirmation through new and independent observations. Thus, we started a successful programme of HST observations aimed at a solid detection of the candidate’s proper motion. In this paper, we report the outcome of this programme: the observations and data reduction are described in section 2, the results are discussed in section 3, while the comparison between optical and radio proper motion values is given in section 3.1. The implications of our measurement are summarized in section 4.
## 2 Observations and data reduction
The HST observations were performed on January 18th 1996, see Mignani et al. (1997), and on January 14th 2000, just after the third refurbishing mission. In both cases, to maximize the angular resolution, the target was centered in the Planetary Camera (PC), which has a pixel size of 45.5 mas. The images were taken with the 555W filter ($`\lambda =5252\AA `$; $`\mathrm{\Delta }\lambda =1222.5\AA `$). The pulsar was observed for three orbits in 1996 and for two orbits in 2000, which correspond to total exposure times of 6200 s and 4400 s, respectively. The observations were performed with the same telescope roll angle chosen in such a way that the PC X,Y axis were oriented along right ascension and declination.
After standard pipeline processing, the images have been cleaned from cosmic ray hits by combining repeated exposures through a median filter. Following the standard astrometric approach, e.g. Caraveo et al. (1996), the frames registration has been performed by computing the pixel coordinate transformation between a common set of reference objects identified in the two images and chosen to be neither extended objects nor relatively bright and saturated stars. The 6 reference objects that satisfy our criteria, together with the pulsar optical counterpart, are identifiable in the January 2000 image shown in Fig.1. Their pixel coordinates have then been computed in the two frames by 2-D gaussian fitting of PSF profiles. The procedure was repeated for different widths of the centering box until the coordinate values were shown to be stable. Finally, the centroids of the reference objects have been determined with an accuracy of $`0.05`$ px for both epochs. The coordinates of the PSR B0656+14 counterpart were computed following the same approach. Of course, owing to its lower S/N, its centering error was found slightly higher and more dependent on the size of the centering box. To be conservative, we have rounded up the associated error to 0.1 px. The quoted uncertainties also account for possible centering errors induced by the subtle exposure-to- exposure jitter ($`<0.04`$ pixels in both frames). All the measured centroids were subsequently corrected for the effects of the PC geometric distorsion (Gilmozzi et al., 1995) by applying the “metric” task of the STSDAS package. The coordinate transformation for the 2000 vs 1996 frame registration was finally computed through a linear fit (rotation plus offset) between the second and the first reference grid. The computed transformation turned out to be accurate to within 0.07 and 0.2 px along the X and Y axis, respectively. We note that the errors on the trasformation, certainly higher than the ones obtained e.g. by the De Luca et al. (2000a) in the case of the Vela pulsar, are independent of the alghoritm used to fit the centroids of the reference stars. This effect is probably to be ascribed to the much lower number of reference stars used to compute the coordinate transformation as well as to their relative distribution in the frame. Indeed, the error along Ra, where the reference stars are distributed on a larger area, is much smaller than the one in Dec, where the reference star are much closer to each other (see Fig.1).
## 3 Results
Having secured a single reference frame, we can compare the relative positions of the pulsar optical counterpart at the two epochs. The total offset turns out to be $`3.76\pm 0.16`$ px, where the error includes all the uncertainties of the astrometric steps (exposures alignement, object centering and frame registration). The displacement is obviously very significant and it represents a convincing proof of the object’s motion. After conversion from pixel to sky coordinates, the measured offset translates into an angular displacement of $`172\pm 8`$ mas over 4 years, corresponding to a proper motion in the plane of the sky of $`43\pm 2`$ mas/yr. We note that this result can not be improved by using the 1989/1991 images of Caraveo et al. (1994), which, although providing a $`3`$ times longer time span, would introduce an error larger than 0$`\stackrel{}{\mathrm{.}}`$2 on the corresponding target position, worsening significantly affecting the accuracy on the pulsar displacement. We note that the measured proper motion turns out to be almost totally in right ascension with the two components $`\mu _\alpha =42.7\pm 2`$ mas/yr, $`\mu _\delta =2.1\pm 3`$ mas/yr defining a position angle of $`93^{}\pm 4^{}`$. Since our WFPC2 observations have been taken virtually on the same day of the year (january 18 and 14, respectively), we can exclude that the measured displacement is affected by the object parallax, which, for a distance of e.g. 300 pc, would be $`3`$ mas (almost all in right ascension). Thus, the measured displacement is due to a genuine proper motion of the source.
### 3.1 Optical vs. Radio
We can now compare our optical proper motion ($`\mu _\alpha =42.7\pm 2`$ mas/yr, $`\mu _\delta =2.1\pm 3`$ mas/yr) to the revised radio one of Pavlov, Stringfellow and Cordova (1996) i.e. $`\mu _\alpha =73\pm 20`$ mas/yr and $`\mu _\delta =26\pm 13`$ mas/yr. Although certainly compatible within their quoted errors, the optical result appears much more accurate than the radio one. Since both measurements span about 4 years and use a comparable number of reference objects (Thompson and Cordova, 1994), the difference in accuracy is ascribable both to systematic errors that, in spite of the strengh of the signal from the radio pulsar, pleague the VLA data (see discussion in Pavlov, Stringfellow and Cordova 1996) and to the higher angular resolution of the WFPC2.
On the other hand, the reliability of HST astrometry has been recently assessed by DeLuca et al. (2000b) through comparing independent couples of WFPC2 observations of the Vela and Geminga pulsars, which yielded, for each object, fully consistent proper motion measurements.
## 4 Conclusions
Taking advantage of the outstanding performances of the Planetary Camera on HST, we have secured a firm measurement of the proper motion of the candidate optical counterpart of PSR B0656+14, which moves at a yearly rate of $`43\pm 2`$ mas. Even in the absence of any other supporting evidence, the presence of such a proper motion for a 25 mag object, coupled with a distance in the range $`250÷800`$ pc, would imply a transverse velocity anywhere between 50 to 160 km/sec. Such values are typical of just one class of astronomical object: isolated neutron stars. This was the case of Geminga, when the measurement of the proper motion of the optical counterpart (Bignami et al., 1993) clinched the identification of the faint G” with one of the brightest gamma-ray source in the sky (Bignami & Caraveo, 1996). Although similar to Geminga in many ways, PSR B0656+14 does not offer such a challenging situation: we face a bona fide radio pulsar, with a radio proper motion and a promising optical counterpart, the identification of which is supported both by multicolor photometry (Pavlov et al., 1997) and timing (Shearer et al., 1998). However, although both methods appear suggestive and definetely worth pursuing, the burden of the proof of PSR B0656+14 optical identification should rest on a firmer ground. Our highly significant measurement of a proper motion of the optical counterpart provides just such a straighforward evidence.
We wish to thank Giovanni F. Bignami for many stimulating discussions. ADL thanks the Space Telescope European Coordinating Facility for hospitality and support. This research is supported by the Italian Space Agency (ASI).
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# A New Analysis Method for Reconstructing the Arrival Direction of TeV Gamma-rays Using a Single Imaging Atmospheric Cherenkov Telescope
## 1 Introduction
The Whipple Collaboration operates a 10 m optical reflector for gamma-ray astronomy at the Fred Lawrence Whipple Observatory on Mt. Hopkins (elevation 2320 m) in southern Arizona. The reflector was originally constructed in 1968 weekes72 and numerous modifications have been made to improve the sensitivity and performance of the system. A camera consisting of photomultiplier tubes (PMTs) mounted in the focal plane of the reflector, detects the Cherenkov radiation produced by gamma-ray and cosmic-ray air showers from which an image of the Cherenkov light can be reconstructed. The camera is triggered when any two PMT signals are above a threshold within a short time coincidence.
In the past decade the camera has undergone significant expansion from a field of view (FOV) of $`3.0^{}`$ in 1989 to a FOV of $`4.8^{}`$ in 1999. From 1989 to 1996 the camera consisted of 109 PMTs, each viewing a circular field of $`0.259^{}`$ diameter, yielding a total FOV of $`3.0^{}`$. The trigger condition required any two of the inner 91 PMT signals to be above a threshold within a 15 ns time coincidence. During the spring of 1997 the camera was expanded to 151 PMTs yielding a total FOV of $`3.5^{}`$. The trigger condition required any two of the inner 91 PMT signals ($`3.0^{}`$ triggering FOV) to be above a threshold with a coincidence resolving time of 15 ns. From 1997 to 1999 the camera was further expanded to 331 PMTs yielding a total FOV of $`4.8^{}`$. The trigger was expanded to include all 331 PMT signals and the time coincidence was shortened to 10 ns due to the introduction of constant fraction discriminators. The layout of each camera is depicted in Figure 1. A full description of the reflector, which has not changed since its original construction, can be found in cawley90 .
Primary cosmic rays and gamma rays entering the atmosphere initiate showers of secondary particles which propagate down towards the ground. The trajectory of the shower continues along the path of the primary particle. If the optical reflector lies within the 300 m diameter Cherenkov light pool, it forms an image in the PMT camera. The appearance of this image depends upon a number of factors. The nature and energy of the incident particle, the arrival direction and the point of impact of the particle trajectory on the ground, all determine the initial shape and orientation of the image. This image is modified by the point spread function of the telescope, the addition of instrumental noise in the PMTs and subsequent electronics, the presence of bright stellar images in certain PMTs and by spurious signals from charged cosmic rays physically passing through the tubes. Monte Carlo studies have shown that gamma-ray induced showers give rise to more compact images than background hadronic showers and are preferentially oriented towards the source position in the image plane hillas85 . By making use of these differences, a gamma-ray signal can be extracted from the large background of hadronic showers and a gamma-ray map over the FOV can be obtained. The method of extracting a gamma-ray signal from the hadronic background can be found in fegan94 . The technique of obtaining a gamma-ray map of the FOV, by reconstructing the unique arrival direction of very high energy photons, is the topic of this paper.
Data from the Whipple Observatory’s 10 m high energy gamma-ray telescope, will be used to demonstrate methods of the reconstruction of the arrival direction of gamma-ray induced showers. These methods may be applied a) to a search for point sources located anywhere within the camera’s FOV, for example searches for counterparts to EGRET unidentified sources, gamma-ray bursts, described in connaughton98 , or sky surveys and b) analysis of extended sources of TeV gamma-rays such as supernova remnants, described in buckley98 and Galactic plane emission described in lebohec2000 .
## 2 Shower Image Processing and Characterization
Prior to analysis of the recorded images, two calibration operations must be performed: the subtraction of the pedestal analog-digital conversion (ADC) values and the normalization of the PMT gains, a process known as flat-fielding.
The pedestal of an ADC is the finite value which it outputs for zero input. This is usually set at 20 digital counts so that small negative fluctuations on the signal line, due to night sky noise variations, will not generate negative values in the ADC. The pedestal for each PMT is determined by artificially triggering the camera, thereby capturing ADC values in the absence of genuine input signals. The PMT pedestal and pedestal variance are calculated from the mean and variance of the pulse-height spectrum (PHS) generated from these injected events.
The relative PMT gains are determined by recording a thousand images using a fast Optitron Nitrogen Arc Lamp illuminating the focal plane through a diffuser. These nitrogen pulser images are used to determine the relative gains by comparing the relative mean signals seen by each PMT.
Fluctuations in the image usually arise from electronic noise and night-sky background variations. To reject these distortions a PMT is considered to be part of the image if it either has a signal above a certain threshold or is beside such a PMT and has a signal above a lower threshold. These two thresholds are defined as the picture and boundary thresholds, respectively. The picture threshold is the multiple of the root mean square (RMS) pedestal deviation which a PMTs signal must exceed to be considered part of the picture. The boundary threshold is the multiple of the RMS pedestal deviation which PMTs adjacent to the picture must exceed to be part of the boundary. The picture and boundary PMTs together make up the image; all others are zeroed. This image cleaning procedure is depicted in Figure 2. These thresholds were optimized using data taken on the Crab Nebula yielding picture threshold: 4.25, and boundary threshold: 2.25.
For the results reported in this paper, the data are obtained in an ON/OFF mode where the source position is tracked for 28 minutes (ON), followed by a 28 minute observation of a background region (OFF) covering the same path in elevation and azimuth. In an ON/OFF observation mode, differences in sky brightness between ON and OFF sky regions could introduce biases. These biases can severely affect the selection of pixels accredited to the boundary region and hence distort the Cherenkov image. For example, if the ON region is brighter than the OFF region, bias may arise as follows. Consider tubes which have a combination of small amounts of genuine signal coupled with some noise (e.g. boundary tubes). If the noise level is low, the boundary threshold is low and most of these tubes will pass the threshold test and be included as part of the image. However, if the noise level is high then the boundary threshold will also be high and the probability increases that a negative noise fluctuation will cancel the genuine signal component resulting in the tube being set to zero during image cleaning. Thus, the degree to which boundary tubes are set to zero depend on the noise level in the tubes. As a consequence a bright sky region will result in more boundary tubes being set to zero, making the image appear narrower and more gamma-ray like.
A software technique was developed to correct for the biases introduced by the differing sky brightness cawley93 . This technique, known as software padding works by adding software noise into the events for the darker sky region. Let $`P_{\mathrm{on}}(P_{\mathrm{off}}),\sigma _{\mathrm{on}}(\sigma _{\mathrm{off}})`$ be the ON(OFF) pedestal and pedestal deviation values for a particular pixel and $`C_{\mathrm{on}}(C_{\mathrm{off}}),\sigma _{C_{\mathrm{on}}}(\sigma _{C_{\mathrm{off}}})`$ be the component due to the Cherenkov signal and corresponding fluctuation. The total ON signal is
$$\mathrm{ON}=P_{\mathrm{on}}+\sigma _{\mathrm{on}}Gauss(0:1)+C_{\mathrm{on}}+\sigma _{C_{\mathrm{on}}}Gauss(0:1),$$
(1)
where $`Gauss(0:1)`$ is a random number drawn from a Gaussian distribution of zero mean and unit variance. The noise component due to the night sky light in the ON region is:
$$N_{\mathrm{on}}=\sigma _{\mathrm{on}}Gauss(0:1).$$
(2)
Similarly, for the OFF region:
$$N_{\mathrm{off}}=\sigma _{\mathrm{off}}Gauss(0:1).$$
(3)
Suppose we are working with pairs where $`N_{\mathrm{on}}`$ is larger than $`N_{\mathrm{off}}`$ then we wish to add noise, $`N_{\mathrm{add}}`$, in the OFF events such that
$$N_{\mathrm{on}}^2=N_{\mathrm{off}}^2+N_{\mathrm{add}}^2,$$
(4)
or
$$N_{\mathrm{add}}=\sqrt{N_{\mathrm{on}}^2N_{\mathrm{off}}^2}.$$
(5)
The total OFF signal is then
$`\mathrm{OFF}=P_{\mathrm{off}}+\sigma _{\mathrm{off}}Gauss(0:1)+N_{\mathrm{add}}Gauss(0:1)+`$
$`C_{\mathrm{off}}+\sigma _{C_{\mathrm{off}}}Gauss(0:1).`$ (6)
When the OFF region is brighter than the ON region then $`N_{\mathrm{add}}`$ is added to the ON pixels. This software padding procedure is very efficient in removing most of the biases induced by the sky brightness with only a modest reduction in sensitivity. Given the potential for large systematic errors without software padding, the modest reduction in sensitivity is a small price to pay.
An example of sky brightness is given in Figure 3 for the ON and OFF source regions of the supernova remnant G78.2+2.1 (see buckley98 for the results of TeV gamma-ray observations). For this object, located in a bright region of the galactic plane, the difference in sky brightness is quite large thus the application of software padding is essential. The results of the analysis of this extended object, with and without the application of software padding, will be presented later in this paper.
We characterize each Cherenkov image using a moment analysis reynolds93 . The roughly elliptical shape of the image is described by the length and width parameters. Its location and orientation within the FOV are given by the distance and $`\alpha `$ parameters, respectively. The asymmetry parameter, defined as the third moment of the light distribution, describes the skew of the image along its major axis. We also determine the two highest signals recorded by the PMTs (max1, max2) and the amount of light in the image (size). These parameters are defined in Table 1 and are depicted in Figure 4.
## 3 Analysis Techniques
Gamma-ray events give rise to shower images which are preferentially oriented towards the source position in the image plane. These images are narrow and compact in shape, elongating as the impact parameter increases. They generally have a cometary shape with their light distribution skewed towards their source location in the image plane. Hadronic events give rise to images that are, on average, broader (due to the emission angles of pions in nucleon collisions spreading the shower), and longer (since the nucleon component of the shower penetrates deeper into the atmosphere) and are randomly oriented within the FOV. Utilizing these differences, a gamma-ray signal can be extracted from the large background of hadronic showers.
### 3.1 Traditional Methods
The standard gamma-ray selection method utilized by the Whipple Collaboration is the Supercuts criteria (see Table 2; cf., reynolds93 ; catanese95 ; lebohec2000 ). These criteria were optimized on contemporaneous Crab Nebula data giving the best sensitivity to point sources positioned at the center of the FOV. In an effort to remove background events triggered by single muons and night sky fluctuations, Supercuts incorporates pre-selection cuts on the size and on max1 and max2. The cuts on width and length select the more compact gamma-ray images and the cut on distance selects images for which the pointing angle $`\alpha `$ is well defined. The changes to the upper bounds of the length and distance cuts for the various cameras reflects the increasing FOV which results in less image truncation thereby allowing an accurate reconstruction of more distant images. A final cut on $`\alpha `$ selects images which are aligned towards the source position, assumed to be at the center of the FOV. A gamma-ray signal is detected as a statistically significant excess of events, which pass the above criteria, between ON source and OFF source observations. In the case when the source is not positioned at the center of the FOV, this technique is insensitive to any excess.
### 3.2 Photon Arrival Direction
If the source position is ill-defined or if the source is extended, the traditional method becomes ineffective. As a result two different methods may be applied. In the first, the FOV is divided into a grid and each grid point is tested for an excess of events above background utilizing the Supercuts selection criteria given in Table 2 or using cuts which make use of asymmetry to break the pointing ambiguity and width/length to better localize the image. This method is described in detail in akerlof91 , fomin94 and buckley98 . In the second method, a unique arrival direction is determined thereby generating a TeV gamma-ray map of the FOV (see also lebohec98 for a description of an alternative method utilized by the Cherenkov Array at Thémis (CAT) group). The benefit of the second method is that it can be easily utilized in a search for extended or diffuse emission. This second method is described below.
By employing several features of gamma-ray induced showers, the arrival direction of a photon can be determined with the atmospheric Cherenkov technique, using a single telescope. Cherenkov light images of showers can be characterized as elongated ellipses with the major axis representing the projection of the shower trajectory along the image plane. The position of the source must lie along this axis near the tip of the light distribution corresponding to the initial interaction of the shower cascade. In addition, the position of the source must lie in the direction indicated by the asymmetry of the image. The elongation of a shower image and the angular distance between its centroid and the source position depend upon the impact parameter of the shower on the ground (see Figure 4). For small impact parameters, the image should have a form close to that of a circle and be positioned very near the source position in the focal plane fomin94 . For increasing impact parameter it should become elongated and have a form close to that of an ellipse and be positioned further from the source position in the focal plane as shown in Figure 4. The elongation of an image can be expressed as the ratio of its angular width and length. A simple form for the relationship between the elongation of an image and the angular distance between its centroid and the source position, disp, is
$$disp=\xi (1\frac{width}{length})$$
(7)
where $`\xi `$ is a scaling parameter. A combination of these features provides a unique arrival direction for each gamma-ray event.
The value of disp depends on an unknown scaling factor $`\xi `$. Generally, $`\xi `$ will depend on the height of the shower in the atmosphere, the elevation of the detector on the ground, the zenith angle of the observation, parameters of the model of the atmosphere and the energy of the primary particle. For observations at small zenith angles ($`<35^{}`$), and for the majority of events near the threshold of the instrument, the dependence on zenith angle and energy can be neglected. The value of $`\xi `$ will be further modified by the effect of the finite size of the camera in the focal plane giving rise to edge effects due to the truncation of images. Experimentally, $`\xi `$ is determined from data so that the calculated shower arrival directions line up with a known source position. For example, if a value of $`\xi `$ is chosen which is too small, then due to azimuthal symmetry, the arrival directions will circumvent the source location, similarly for values of $`\xi `$ too large. The optimum value of $`\xi `$ minimizes the spread in the angular distribution of events centered at the source position. The dependence of the angular spread of the arrival directions on the scaling factor $`\xi `$ is depicted in Figure 5. By minimizing the spread in the arrival directions the optimum angular resolution is achieved. This has the added advantage of increasing the signal to background ratio when determining an excess of events from a region.
Data from the Whipple 10 m gamma-ray telescope is used to optimize the value of $`\xi `$. We include data from each of the three cameras (see Table 3), taken on established sources, positioned at the center of the FOV which should appear point-like. First, events are selected as gamma-like based on the angular width and length criteria given in Table 2. No selection on distance or $`\alpha `$ is made. Secondly, the arrival direction of each event is determined utilizing Equation 7. Next, each arrival direction is corrected for the rotation of the FOV, due to the altitude-azimuth mount used by the Whipple 10 m telescope, by de-rotating the position to zero hour angle. These corrected arrival directions are binned on a two dimensional grid of bin size $`0.1^{}`$x$`0.1^{}`$. Figure 6 shows examples of the binned excess arrival directions, for data taken on Markarian 501 (see Table 3), projected along the x-axis, for two values of the scaling parameter $`\xi `$. The excess shows the difference between accumulated ON source events and corresponding OFF source control data.
The spread in arrival directions for each value of $`\xi `$ is determined by fitting a Gaussian function of standard deviation $`\sigma `$ to the binned excess. The results are depicted in Figure 7. The minimum spread occurs at a value of $`\xi =1.78^{},1.78^{},1.65^{}`$ for the $`3.0^{},3.5^{}`$ and $`4.8^{}`$ camera FOV respectively. The effect of the decreasing optimum value of $`\xi `$ with increasing FOV is the result of less image truncation imposed by the larger camera. This effect is also evident in the upper bound of the optimum length cut used to select gamma-ray images (see Table 2). The minimum spread is a measure of the optimum angular resolution, or point spread function (PSF) of the technique and corresponds to $`\sigma =0.11^{},0.13^{}`$ and $`0.12^{}`$ respectively.
## 4 Results
The determination of a photon’s arrival direction enables the search for emission from anywhere within the camera’s FOV. When the source position is known, a gamma-ray signal is detected as an excess of events, between corresponding ON and OFF source observations, originating from the source direction. Firstly, events are selected as gamma-like based on the angular width and length criteria given in Table 2 and their arrival directions determined as described above. Secondly, the radial distance from the arrival direction to the source position, $`\theta `$, is calculated and binned in $`0.02^{}`$ intervals. The results of data taken with the three cameras are shown in Figure 8. The area of exposure for each bin increases with greater radial distance, therefore each bin contents have been divided by the area of the annulus defined by the bin boundaries. The physical interpretation of such a distribution is the surface brightness of an extended object.
### 4.1 Aperture Selection Criteria
For each data set depicted in Figure 8, an excess gamma-ray signal is apparent in the ON source observation (solid line) when compared to the OFF source observation (dotted line) and the excess is confined to small values of $`\theta `$. The angular extent of the signal is similar for each camera, consistent with the results of the optimized angular resolution. The results of the OFF source observations illustrate a background which appears uniform for small values of $`\theta `$ with a gradual decline due to the finite size of the camera. In the $`3.0^{}`$ and $`4.8^{}`$ cameras, all PMTs enter into the trigger (all of the way out to the edge of the FOV), whereas in the $`3.5^{}`$ camera only the inner 91 out of the 151 PMTs enter into the trigger (see Figure 1). Thus in the $`3.0^{}`$ and $`4.8^{}`$ cameras there are a number of background cosmic-ray events which lie on the edge of the camera and which are distorted by image truncation. This truncation tends to produce images whose major axes make an angle of $`\alpha =90^{}`$ and results in a significant rise in the number of events collected at large values of $`\theta `$. This rise is not apparent in the results from the $`3.5^{}`$ camera. The selection of events consistent with the position of the source is accomplished by counting the number of events within a circular aperture of radius $`\theta _c`$ centered at the source position. The optimum value of $`\theta _c`$ is determined using data taken with the $`4.8^{}`$ camera. The results are shown in Figure 9 and yield an optimum selection criteria of $`\theta <0.22^{}`$. This selection criteria can be applied to the data for all camera configurations owing to the similar radial extent of the gamma-ray signal. When the angular extent of the source is of the order, or greater than the resolution of the technique, the size of the circular aperture can be adjusted to the size of the emission region. In doing so more background is included, hence sensitivity is diminished. In background dominated counting statistics, the sensitivity of a detector will be proportional to $`1/\sqrt{background}`$. If we neglect the gradual decline in the background with increasing angular offset, the sensitivity of an atmospheric Cherenkov telescope utilizing this method will scale as $`1/\theta _c`$.
TeV gamma-ray observations of the supernova remnant G78.2+2.1, reported by the Whipple Collaboration in buckley98 , is an example of an extended source. The expanding shell of the supernova explosion subtends an angle of approximately $`0.5^{}`$. The distribution of $`\theta `$, for a subset of the data reported in buckley98 , is given in Figure 10. In the search for a gamma-ray signal, a circular aperture was chosen to encompass the size of the emission region plus twice the angular resolution to account for smearing of the edges. As shown in Figure 10, no significant excess was obtained. Also shown is an example of a false signal resulting from biases due to sky brightness differences. Without the application of software padding, the brighter ON source region yields a greater number of background events selected as gamma-rays resulting in a significant excess between the ON and OFF source observations.
### 4.2 Gamma-ray Collection Efficiency
We compare the gamma-ray collection efficiency and background rejection of the aperture selection method with the traditional Supercuts criteria using data taken on the Crab Nebula, the standard candle for TeV gamma-ray astronomy. The results from the three camera configurations are given in Table 4. The data collected with the $`3.0^{}`$ camera was taken in January - February 1995. Like the traditional Supercuts criteria, the aperture selection includes events which are oriented towards the center of the FOV. However, the aperture selection has two further constraints in that events must be skewed towards the center of the FOV and must be elongated in proportion to the impact parameter on the ground. These additional criteria result in greater background rejection compared with Supercuts. Results from the $`3.0^{}`$ camera indicate a reduced gamma-ray collection efficiency due to image truncation imposed by the smaller FOV. Truncated images reduce the effectiveness of the asymmetry parameter by clipping the tails of the light distribution which defines the skewness of the image. Data collected with the $`3.5^{}`$ and $`4.8^{}`$ cameras were taken in January - February 1997 and November 1998 - January 1999 respectively. First note that for the data taken with the $`4.8^{}`$ FOV camera, a degradation of mirror reflectivity and lack of light cones resulted in an increased energy threshold and lower count rate on the Crab Nebula. Despite the increased threshold this data can still be used to judge the effectiveness of the extended camera. These results not only show the same improvement in background rejection over the traditional Supercuts criteria, but indicate similar gamma-ray collection efficiency. This results in a marginal increase in sensitivity of the aperture selection criteria over traditional Supercuts.
### 4.3 The Sky in TeV Gamma Rays
If the position of the source is unknown but assumed to be in the FOV, a two dimensional grid of bin size $`0.1^{}`$x$`0.1^{}`$ is constructed. Events are first selected as gamma-like based on the angular width and length criteria given in Table 2, and their arrival directions determined as described above. Next, the contents of each grid point within the optimum cut on the radial distance from the arrival direction is incremented by one. Figure 11 shows the accumulated grid contents for data taken with the $`3.5^{}`$ camera. Conceptually, this performs a type of image smoothing about the calculated photon arrival direction resulting in an oversampled gamma-ray map which is tuned to the search for point sources anywhere within the camera’s FOV. The statistical significance of excess events between ON and corresponding OFF source observations is calculated for each grid point using the method of li83 . The use of Poissonian statistics is valid in this case because grid points are incremented by unity. However, one must be careful in interpreting the significance. To be precise, the test statistic calculated at each grid point can only be interpreted as significance if there was a prior hypothesis for gamma-ray emission from a point source at that position. Otherwise, the number of independent grid points tested, or trials factor, must be taken into account thereby degrading the significance. For example, for a grid ($`3.9^{}`$ FOV divided into $`0.1^{}`$ bins) with 1521 grid points (not independent) we have effectively 324 trials and expect on average one $`3\sigma `$ excess.
An example of a TeV gamma-ray image of a region of sky (the supernova remnant G78.2+2.1), found to be devoid of a gamma-ray signal, is shown in Figure 12. The image shows several $`2\sigma `$ excesses and one $`3\sigma `$ excess. However, without a prior hypothesis for gamma-ray emission from a point source at these positions a trials factor for the number of independent grid points must applied, resulting in a reduction in significance. Also shown in Figure 12 is an example of a false signal resulting from biases due to sky brightness differences. The region of the image showing the greatest excesses appear well correlated with the sky brightness difference between the ON and OFF source regions as depicted in Figure 3.
To test the efficacy of the technique and measure the efficiency and sensitivity away from the center of the FOV, we analyzed data taken on the Crab Nebula which was offset in declination. The results are shown in Figures 13,14 and 15 for the $`3.0^{},3.5^{}`$ and $`4.8^{}`$ cameras respectively. The contours are proportional to the statistical significance of the excess between the ON and OFF source data. The peak excess and detected gamma-ray rates are given in Table 5.
The results show that the technique is sensitive to a gamma-ray signal offset from the center of the FOV. For the $`3.5^{}`$ and $`4.8^{}`$ cameras, the Crab Nebula was detected with a sensitivity of approximately $`3\sigma `$/hr for an offset of $`1.5^{}`$. The derived position of the Crab Nebula was within the $`0.1^{}`$ pointing accuracy of the Whipple 10 m telescope for offsets less than $`1.0^{}`$. We note that there appears to be a small systematic shift of the derived position of the Crab Nebula for the $`1.0^{}`$ and $`1.5^{}`$ offsets for the data taken with the $`3.5^{}`$ camera and for the $`1.5^{}`$ offsets for the data taken with the $`4.8^{}`$ camera. The shift is most likely due to the inclusion of truncated images as the source moves closer to the edge of the FOV. The spread of events about the derived source position was determined by fitting a Gaussian function to the unsmoothed binned excess. The results indicate that the PSF is approximately constant over the FOV of the camera and thus not dependent on the optical properties of the reflector which degrade with increasing angular offset.
As the gamma-ray source moves away from the center of the camera, less of the Cherenkov light pool falls on the detector at a given impact parameter. This has the effect of reducing the collection area for gamma-rays. The measurements of the gamma-ray rate from the Crab Nebula at increasing angular offset is a direct measure of this reduction as shown in Figure 16. The effective FOV, which we arbitrarily define as the radius corresponding to a 50% efficiency, appears to scale linearly with camera physical FOV, at least up to the range of offsets and camera sizes included in these results. For example, a camera with a physical FOV of $`6.0^{}`$ would have an effective FOV double that of a camera with a $`3.0^{}`$ FOV. Moreover, by doubling the diameter of the FOV a fourfold increase in sensitive area for sources offset from the center of the FOV is realized.
## 5 Conclusion
We have described a method of atmospheric Cherenkov imaging which reconstructs the unique arrival direction of TeV gamma-rays using a single telescope. This method is derived empirically making use of the Crab Nebula as a standard candle to optimize the angular resolution of the technique. This allows a selection of events based on the position of a source anywhere within the telescope’s FOV. We found that such a selection yields similar sensitivity and collection efficiency to the traditional Supercuts criteria utilized by the Whipple Collaboration. However, as demonstrated, the technique is easily applicable to sources offset from the center of the FOV or sources of extended emission.
This research is supported by grants from the U.S. Department of Energy. We are grateful to Dave Fegan and Trevor Weekes for their guidance and assistance. We thank the Whipple Gamma-ray Collaboration for the use of the data presented in this paper and acknowledge the technical assistance of K. Harris and E. Roache.
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# Slave-boson mean-field theory of spin- and orbital-ordered states in the degenerate Hubbard model*footnote **footnote *To be submitted to a special issue of ”Foundation of Physics” celebrating the 75-th birthday of Martin Gutzwiller
## Acknowledgements
It is my great pleasure to dedicate the present paper to Dr. Martin Gutzwiller on the occasion of his 75th birthday. This work has been supported partially by a Grant-in-Aid for Scientific Researach (B) from the Ministry of Education, Science, Sports and Culture of Japan.
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# Cylindrically symmetric dust spacetime
## Acknowledgements
We are grateful to W.B. Bonnor, A. Krasiński, P. Szekeres and J. Wainwright for their comments. The authors also thank financial support from the Basque Country University under project number UPV 172.310-G02/99. RV thanks the Spanish Secretaría de Estado de Universidades, Investigación y Desarrollo, Ministerio de Educación y Cultura, grant No. EX99 52155527.
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# Laplacian growth with separately controlled noise and anisotropy
## I Introduction
In a large class of pattern-forming systems the growth is controlled by a Laplacian field. In diffusion limited aggregation (DLA) this field is the probability density of aggregating particles . In viscous fingering it is pressure , and in crystal growth it can be either a diffusive or a thermal field . After the DLA model was introduced by Witten and Sander , it became standard to simulate the Laplacian field by random walkers, which after being released at the periphery of the system, diffuse towards growing cluster and freeze on it. To simulate DLA, several numerical techniques have been developed, of which most powerful are the off-lattice algorithms .
Applications to other Laplacian problems have been proposed based on the random walks idea. In particular, handling problems such as viscous fingering within this framework requires reducing noise of individual walkers as well as modeling surface tension. Reduction of noise was achieved by the method of multiple hits in which particles freeze on a particular site adjacent to the already grown cluster only after this site has been visited more than $`n_{\mathrm{min}}`$ times, where $`n_{\mathrm{min}}`$ is an acceptance threshold. The method of noise reduction has been introduced in the context of the on-lattice DLA models. More recently, this method was combined with the off-lattice technique and studied theoretically . In an advanced version of the multiple hits method the random walkers move off-lattice and sticking rules are defined by using a finite number $`m`$ of antennas attached to each particle, where, for instance, $`m=4`$ for the square lattice DLA. Each of $`m`$ antennas has a counter which scores the number of times $`n`$ random walkers arrive on it and is then used to set a threshold $`n_{\mathrm{min}}`$ for freezing. Having a finite number of antennas used per each particle makes the sticking rules in the mutiple hits models anisotropic. This anisotropy is essential because it is significantly amplified by the growth dynamics in the low noise regime of $`n_{\mathrm{min}}1`$. This built-in anisotropy of growth rules has been used to test universality of DLA and also to simulate dendritic crystal growth .
It has been proposed that surface tension can be modeled by introducing a probability $`t<1`$ of freezing upon each encounter with the cluster or by making freezing dependent on the local neighbors configuration . In the model with $`t1`$ each randomly walking particle freezes only after encountering the cluster about $`t^11`$ times. As a result, the freezing point is displaced from the point of first encounter by the distance $`dt^1`$ in the units of particle size. Effectively, in this model a finite length scale $`d`$ is introduced over which the harmonic (Laplacian) measure describing the probability of the first encounter is being probabilistically averaged. It has been conjectured in the works (and partially confirmed by various features observed in the growth patterns) that the length scale $`d`$ simulates capillary radius in the Laplacian problem with surface tension.
The field of Laplacian growth, despite being well developed by now, contains several long-standing unresolved problems. Firstly, the lattice simulation, although extremely efficient algorithmically, does not seem to be a natural starting point for analytic understanding of large scale phenomena, such as fingering, fractalization and scaling. Secondly, the methods of simulating Laplacian growth that have been used so far are not entirely free of problems, the most notable being an intrinsic anisotropy of growth rules. Already the original DLA rules use square lattice and thus are anisotropic. This anisotropy is weak and reveals itself only in very large DLA clusters . However, when the noise level is reduced using multiple hits, the underlying lattice anisotropy is amplified . The noise-reduced growth remains vulnerable to anisotropy even when off-lattice random walks are used, due to the anisotropy of freezing rules mentioned above.
It was demonstrated recently that Laplacian growth can be studied using an entirely different approach based on iterated conformal maps . The model uses analytic functions chosen in such a way that upon acting on a unit circle they produce bumps of prescribed size. Iterated $`n`$ times with the parameter defining the bump size chosen according to a certain rule, these maps produce a cluster of $`n`$ particles of nearly equal size. The conformal model of growth has become recently a subject of active work .
The goal of this article is to extend the conformal mapping methodology to the problems with reduced noise and growth anisotropy. Here one clear advantage is that the growth rules using conformal mapping are intrinsically isotropic. Because of that one can easily avoid problems pertinent to other models, in which growth anisotropy and reduced noise are intertwinned. The main idea behind the noise reduction method proposed in this work is to average the Laplacian measure over finite length which is larger than the particle size in the original model . For that we alter particle shape and use “flat” particles extended along the cluster boundary and thin in the growth direction. To compare our method of reducing noise to other techniques, we note that the positions of flat particles are chosen strictly according to Laplacian measure, like in the multiple hits method . The control over noise is achieved by suppressing noise at the length scales shorter than the particle larger dimension. This is in contrast with the multiple hits method, where noise is suppressed due to statistical averaging over many particle growth attempts uniformly over all length scales down to particle size. Because of the appearance of a new length scale our method somewhat resembles the surface tension models used in the DLA lattice growth .
One notable difference from previous models is in the dependence of the computation time on the achieved level of noise reduction. Reducing statistical fluctuations in the multiple hits models required increasing the number of random walkers used to grow the cluster inversely with the noise reduction parameter. In our method one can reduce noise arbitrarily without increasing computation length, simply by varying particle aspect ratio with particle areas kept fixed. We also demonstrate that growth anisotropy can be naturally incorporated in the conformal mapping method without affecting noise reduction.
Our plan in this article is as follows. We start with revisiting conformal mapping model. We discuss some issues ignored before and propose a generalization to the problems with reduced noise and anisotropy. In Sec. II we review the model, keeping focus on aspects that will be important in the rest of the article. In Sec. III we study the distribution of particle areas produced by growth rules employed in Ref. . We observe that these rules lead to occasional appearance of very large particles. To fix this problem, we evaluate particle areas at each growth step and apply an acceptance criterion for newly grown particles according to their area. In Sec. IV we describe a model with reduced noise. To suppress noise we use particles which are thin in the growth direction and smooth at corners. In Sec. V we show how these growth rules can be generalized for anisotropic growth. In Sec. VI we study scaling properties of all introduced varieties of the model and compare them with each other. We find that the fractal dimension estimated from cluster radius scaling is less sensitive to noise than to anisotropy. For isotropic growth, both with and without noise reduction, the dimension is very close to $`1.7`$. For anisotropic growth, reducing noise to the level at which anisotropy reveals itself strongly shifts fractal dimension to somewhat lower values. In this regime, the fractal dimension depends on symmetry, and is found to be $`1.62`$ for the four-fold symmetry and $`1.5`$ for the three-fold symmetry. Finally, in two appendices we discuss in detail the particles area distribution and possible improvements of our numerical procedure.
## II Conformal mapping model
We describe growing cluster by a sequence of domains $`𝒟_0𝒟_1𝒟_2\mathrm{}`$ corresponding to subsequent growth steps in time. In the canonical formulation , growth occurs due to particles diffusing from infinity one by one and freezing as soon as they reach the cluster boundary. The particles concentration $`u(𝐫)`$ obeys diffusion equation, which in the quasi-stationary approximation of slow growth is written as
$$^2u(𝐫)=0\mathrm{with}u(𝐫)=\{\begin{array}{cc}0,\hfill & 𝐫𝒟_{n1};\hfill \\ \frac{1}{2\pi }\mathrm{ln}|𝐫|,\hfill & |𝐫|\mathrm{}.\hfill \end{array}$$
(1)
Zero boundary condition on the cluster $`𝒟_{n1}`$ describes freezing of the $`n`$-th particle upon arrival with probability one. The points of the cluster boundary $`𝒟_{n1}`$ where subsequent additions are made are selected randomly with the probability given by the so-called harmonic measure
$$dP=|u|dl,dl𝒟_{n1},$$
(2)
where $`dl`$ is boundary element of the cluster $`𝒟_{n1}`$. As the domain changes, $`\mathrm{}𝒟_n𝒟_{n+1}\mathrm{}`$, the problem (1) has to be solved again for every new domain to determine from (2) the new particle position probability.
A considerable computational simplification of the problem can be achieved by using a sequence of analytic functions $`F_n(z)`$, $`n=0,1,2,\mathrm{}`$, to represent the domains $`𝒟_n`$. The functions $`F_n`$ are chosen so that each of them defines a conformal one-to-one mapping of the unit disk $`|z|1`$ on the domain $`𝒟_n`$, including boundary. Adding a new object to the cluster at the $`n`$-th growth step is described by changing the mapping $`F_n`$ as follows:
$$F_n(z)=F_{n1}(f_{\lambda _n,\theta _n}(z)),F_0(z)=z.$$
(3)
Here the function $`f_{\lambda _n,\theta _n}(z)`$ maps the unit circle $`|z|=1`$ onto a unit circle with a bump centered around the point $`z=e^{i\theta _n}`$ of the circle. The bump size is determined by the parameter $`\lambda _n`$ as discussed below. The angle $`\theta _n`$ is chosen randomly at each growth step.
The key simplification that arises in the conformal mapping representation (3) is due to the fact that the harmonic measure (2) is translated into a uniform probability distribution for $`\theta _n`$, so that $`dP(\theta )=d\theta /2\pi `$. Also, there is no statistical correlation between subsequent $`\theta `$’s.
The form of the function $`f_{\lambda ,\theta }(z)`$ growing bumps can be chosen according to the computational needs . In this article we use
$$f_{\lambda ,\theta }(z)=e^{i\theta }g^1\left(\stackrel{~}{f}_\lambda \left(g(e^{i\theta }z)\right)\right),$$
(4)
where the function $`g(z)=(z1)/(z+1)`$ maps the unit disk $`|z|1`$ onto the left half-plane $`\mathrm{Re}z0`$, and the function
$$\stackrel{~}{f}_\lambda (z)=h_\lambda (z)/h_\lambda (1),h_\lambda (z)=z+\sqrt{z^2+\lambda ^2}$$
(5)
grows a semicircle of radius $`\lambda `$, as shown in Fig. 1. The function $`\stackrel{~}{f}_\lambda (z)`$ is defined in (5) so that $`\stackrel{~}{f}_\lambda (1)=1`$. This ensures that the mapping (4) maps $`z=\mathrm{}`$ onto itself.
Ideally, the values $`\lambda _n`$ defining particle size should be chosen so that all particle areas are equal. In the conformal mapping model this is approximately realized via predicting the bump size to be obtained at the $`n`$-th step using the Jacobian of the already-grown cluster mapping $`F_{n1}`$.
The argument is as follows. The area of the semicircular bump grown using $`\stackrel{~}{f}_{\lambda _n}(z)`$ is $`\pi \lambda _n^2(1+O(\lambda _n^2))/8`$. The area of the corresponding bump produced by $`F_n(z)`$, at small $`\lambda _n`$, is approximately $`|J_{n1}|^2\pi \lambda _n^2/2`$, where $`J_{n1}`$ is the Jacobian of $`F_{n1}`$ evaluated at the position of the $`n`$-th bump:
$$J_{n1}=F_{n1}^{}(z=e^{i\theta _n}).$$
(6)
Hence, to compensate for stretching due to the Jacobian $`J_{n1}`$, one has to choose the values of $`\lambda _n`$ as follows:
$$\lambda _n=|J_{n1}|^1\lambda _0,$$
(7)
where the parameter $`\lambda _0`$ defines particle size. For the growth involving particles of very small size the rule (7) would have been sufficient to ensure identical areas of all bumps. For our problem, in which the bump sizes are small but finite, the areas are only approximately equal. However, one can demonstrate that, after certain improvements discussed in Section III, the rule (7) produces bumps with sufficiently close areas.
The form (4), (5) of the mapping $`\stackrel{~}{f}`$ has several nice features. First, since fractional linear function $`g(z)`$ maps a circle onto a circle, the mapping $`f_{\lambda ,\theta }`$ produces a crescent-shaped particle with circular boundary (see Fig. 1). Second, a simple calculation shows that the particle curvature radius equals $`\lambda `$. The latter has the following consequence. Consider growth starting from a circular cluster of radius $`r`$, described by the mapping $`F(z)=rz`$. The mapping $`F(f_{\lambda ,\theta }(z))`$ then produces a particle of curvature radius equal to $`\lambda r`$. After the value of $`\lambda `$ is chosen according to the rule (7), $`\lambda =\lambda _0/F^{}=\lambda _0/r`$, the particle radius becomes equal to $`\lambda _0`$, independently of the cluster radius $`r`$. The area of this particle is readily evaluated:
$`a(\lambda _0,r)`$ $`=`$ $`a_{}+\lambda _0r(r^2\lambda _0^2)\mathrm{tan}^1(\lambda _0/r),`$ (9)
$`a_{}=\pi \lambda _0^2/2.`$
The area $`a(\lambda _0,r)`$ varies between $`a_{}`$ for $`\lambda _0r`$ and $`2a_{}`$ for $`\lambda _0r`$. For a generic non-circular cluster the particle area cannot be found analytically. Statistics of the areas will be discussed in Section III. We will see that typical area of a particle is of the order of $`a_{}`$.
The overall size of the cluster $`𝒟_n`$ grown according to (3), (7) is well characterized by the mapping $`F_n`$ stretching factor at large scales:
$$R_n=F_n^{}(z\mathrm{})=\underset{k=1}{\overset{n}{}}f_{\lambda _k,\theta _k}^{}(z\mathrm{}).$$
(10)
The cluster radius $`R_n`$ can be conveniently evaluated using (10) together with the property $`g(\mathrm{})=1`$ as follows:
$$R_n=\underset{k=1}{\overset{n}{}}\left(\stackrel{~}{f}_{\lambda _k,\theta _k}^{}(z=1)\right)^1=\underset{k=1}{\overset{n}{}}\left(1+\lambda _k^2\right)^{1/2}.$$
(11)
The reason for $`R_n`$ to be an accurate measure of the cluster $`𝒟_n`$ radius lies in the properties of so-called univalent functions .
At large $`n`$ the cluster radius $`R_n`$ is expected to grow as $`n^\alpha `$, where $`\alpha `$ is a numerical constant. This is consistent with (11) provided that $`\lambda _n^2n2\alpha `$ at large $`n`$ . The growth problem (1), (2) is believed to give rise to fractal objects with fractal dimension $`d<2`$. There are several conventional definitions of fractal dimension of a growing cluster . In this article we employ scaling of the cluster size with its area. Also, one can use box counting, or the relation between average growth velocity in a strip and the strip width.
Taking $`R_n`$ defined in (10) as a cluster radius provides a numerically efficient method for calculating fractal dimension. For that one looks for a scaling relation of the form
$$R_nA_n^{1/d},$$
(12)
where $`A_n`$ is total area of the cluster $`𝒟_n`$. The dimension $`d`$ is related with the parameter $`\alpha `$ describing scaling of $`\lambda _n`$ as follows: $`d=\alpha ^1`$, which is true provided $`A_nn`$. In our simulation we make sure that individual particle areas have a sufficiently narrow distribution (see Section III), and thus total cluster area is indeed proportional to the particle number.
Scaling properties of the growth problem described above have been explored by several groups . It was concluded that the properties of the growth resulting from the conformal mapping model match those of the lattice DLA models. Below we revisit the relation between the problem (1), (2) and the conformal mapping model (3), (4), (7) and discuss several interesting extensions of this model.
## III Area distribution
In the original work it was assumed that the rule (7) is sufficient to produce particles with nearly equal areas. This assumption was apparently consistent with the cluster images in which each particle is represented by one or few points. To investigate this issue more closely, in this work we have chosen a different method of representing particles, in which the exact boundary of each particle is shown. An example of a cluster with the boundaries of all individual particles displayed (see Fig. 2), demonstrates that the areas of almost all of the particles are indeed quite close. However, there is also a number of exceptional particles of large areas.
Large particles tend to appear within fjords and seal the space between well developed branches. Typically, this happens when particle growth is attempted on the periphery of an actively growing region. Insufficiency of the rule (7) for keeping particle areas small is caused by fluctuations of the Jacobian $`F_n^{}(z)`$ over the unit circle $`|z|=1`$. These fluctuations can be large in magnitude and also very abrupt, happening on the scale of the order of $`\lambda _n`$ within the circle arc mapped onto the particle boundary. In the case when a newly grown particle overlaps with such a fluctuation, it can be “artificially stretched” under the mapping.
The appearance of large particles has been reported in Ref. and a method for eliminating them was proposed, based on choosing an optimal shape of particles produced by the mapping $`f_{\lambda ,\theta }(z)`$. It was argued that the best value of the parameter $`0<a<1`$ in the mapping defined in Ref. is given by $`a=2/3`$. This value provides a compromise between abundance of large particles at $`a0`$ and needle-like particles shape at $`a1`$. Since in this article the particle shape will be used to tune noise, we employ a different method for eliminating large particles, as described below.
To study the role of exceptional particles, we calculate the particle areas generated by the growth model (3), (7). The method employed to evaluate particle areas is the following. For the particle grown at a step $`n`$, we start with few points on the unit circle $`|z|=1`$ which are mapped by $`F_n`$ on the particle boundary. Subsequently, we add new points on the circle $`|z|=1`$ in between the old points, and compute distances between images of neighboring points under all mappings
$$f_n,f_{n1}f_n,\mathrm{},F_n=f_1f_2\mathrm{}f_n,$$
(13)
where $`f_k`$ is a shorthand notation for $`f_{\lambda _k,\theta _k}`$ and $``$ stands for mapping composition. We keep adding new points until the distances between the images of all neighbors will not exceed $`\gamma \lambda _0`$, where $`\gamma 1`$ is a numerical factor. We used the above procedure with $`\gamma =0.05`$, which produces about two to four hundred points per particle.
This method enables one to have an accurate graphical representation of each particle, as demonstrated in Fig. 2, and also to evaluate particle areas with the accuracy on the level of $`0.1\%`$. This is illustrated in Fig. 3(a) showing how the cluster total area $`A_n`$ is changing during the growth of the cluster displayed in Fig. 2. The area $`A_n`$ grows as a function of $`n`$ in small steps of order $`a_{}=\pi \lambda _0^2/2`$, alternated with occasional jumps of a much larger magnitude. The jumps correspond to the appearance of large particles which seal inner cluster regions. The decomposition of the growth of $`A_n`$ into the smooth and singular parts is revealed more clearly in Fig. 3(b) showing the dependence of $`A_n`$ versus $`n`$ for the same growth as in Fig. 3(a), with a linear function $`2.1a_{}n`$ subtracted.
A histogram of the individual particles areas $`a_n=A_nA_{n1}`$ is plotted in Fig. 4(a). The area distribution was obtained by averaging over $`10`$ realizations of the first $`1000`$ growth steps with the parameter $`\lambda _0=0.2`$ (same as in Figs. 23). The area distribution $`𝒫(a_n)`$ is sharply peaked about $`2.1a_{}`$ and has a number of other interesting properties that will be discussed in Appendix A.
The feature of the distribution $`𝒫(a_n)`$ displayed in Fig. 4(a) which corresponds to the exceptionally large particles is the tail stretching far to the right from the main peak. Note that only the beginning of the tail is displayed in Fig. 4(a) because the weight of the tail in the probability distribution is insignificant, and so the values of $`𝒫(a_n)`$ far in the tail are too small to be visible on the scale of the peak.
To display the tail we replot the distribution $`𝒫(a)`$ on a log–log scale, as shown in Fig 5. The right tail of $`𝒫(a)`$ is power-like, $`𝒫(a)a^\mu `$ with $`\mu 2.5`$. Since $`\mu >2`$, the first moment $`a=a𝒫(a)𝑑a`$ is finite (and thus $`aa_{}`$). However, since $`\mu <3`$, the second moment of $`𝒫(a)`$ is divergent. The existence of the mean particle area $`a`$ means that $`A_nn`$. However, the absence of variance implies that the fluctuations of $`A_n`$ about the mean value are non-Gaussian and larger than required by the central limit theorem. Both features are clear in the sample $`A_n`$ dependence in Fig. 3.
Let us remark that the tail in Fig. 5 in its far end is apparently somewhat steeper than $`a^\mu `$. We believe that this deviation from the $`a^\mu `$ behavior is due to finite number $`N=2000`$ of time steps in the growth samples used to obtain $`𝒫(a)`$. In a finite cluster there is an upper cutoff on possible particle areas. This makes the far tail of $`𝒫(a)`$ non-stationary, shifting the cutoff to larger areas as $`N`$ increases.
Clearly, one would like to inhibit the appearance of large particles with areas in the tail of $`𝒫(a)`$. This is desirable because, even though the tail is quite thin and large particles are rare, occasionally appearing extremely large particles may affect macroscopic characteristics of the growth. In particular, the relatively slow power law decrease in the tail, $`𝒫(a)a^\mu `$, may affect scaling of the cluster size $`R_n`$ and/or the numerical accuracy of the scaling exponent. To eliminate the growth of large particles, we choose an acceptance threshold $`a_{\mathrm{max}}=3a_{}`$ to truncate the tail of the distribution in Fig. 4(a). Then, for each growth step, we calculate the new particle area $`a_n`$. The particle is accepted only if $`a_na_{\mathrm{max}}`$, otherwise the particle is discarded and a new attempt of particle growth is made. An example of the cluster grown according to these rules is displayed in Fig. 6.
One can notice immediately that the overall structure of the branches in Fig. 6 is much more regular than that in Fig. 2. The distribution of areas for such a growth is shown in Fig. 4(b). Within the acceptance window $`[0,a_{\mathrm{max}}]`$ the distribution $`𝒫(a)`$ repeats in all details the distribution shown in Fig. 4(a) for the growth with all particle areas accepted.
For the growth with large particles eliminated, the area $`A_n`$ increases as a linear function of the step number. Average increment of $`A_n`$ is given by the mean value of $`a_n`$ taken from the distribution shown in Fig. 4(b). To verify this, we plot $`A_n`$ versus $`n`$ in Fig. 3 for the same growth parameters as those used in Fig. 2, where only the areas up to $`3a_{}`$ are accepted. Note a small difference between the slope of the dependence at $`n50`$ and at larger $`n`$ that appears because of relatively smaller size of the primary particles growing directly on the unit circle.
The growth model augmented with the area acceptance criterion has a new parameter $`a_{\mathrm{max}}/a_{}`$. In principle, choosing different values of $`a_{\mathrm{max}}`$ gives rise to different growth patterns. However, as long as the window $`[0,a_{\mathrm{max}}]`$ contains much of the $`𝒫(a)`$ peak area, we do not observe any qualitative change in the growth.
Scaling properties of the growth can be studied in several ways. Previous studies of scaling are based on the relation $`R_nn^\alpha `$, where $`R_n`$ was obtained for the growth with unrestricted areas. However, it would be more in the spirit of the notion of a fractal to use the relation (12) between cluster size and its area. This clearly would not work well with large particles being present, because statistical fluctuations of the cluster area $`A_n`$ are quite large in this case. On the other hand, for the growth with restricted areas used in this work, the fluctuations of $`A_n`$ are reduced to the level consistent with the central limit theorem, and thus one can employ the relation (12) to study scaling.
Of course, it is not clear a priori whether the growth with area cutoff is equivalent macroscopically to the growth with unrestricted areas. On general grounds, one may expect the growth to be significantly altered by eliminating large particles. Whether this is true can be indirectly tested by comparing the $`R_n`$ vs. $`n`$ dependence for the growth with unrestricted areas with the fractal dimension obtained from the relation (12) for the growth with area cutoff. We note in that regard that the scaling exponent $`d=1.7`$ found below (see Section VI) matches exactly the value found in Ref. . However, although the presence or absence of large particles seems to be irrelevant for the cluster size scaling, other growth characteristics, such as the structure of branches and fjords, are likely to be more sensitive to the method of treating large particles.
We postpone the discussion of various details and features of the area distribution $`𝒫(a)`$ to Appendix A. In the remaining part of the article we use the conformal mapping model augmented with the area acceptance criterion to study several interesting Laplacian growth problems.
## IV Noise-reduced Laplacian growth
Roughness of the growing cluster is mainly due to two factors: shot noise and the Mullins-Sekerka instability . The shot noise results from the randomness of the aggregating particles positions, and so it contributes to the fluctuations equally on all spatial scales down to the particle size. The Mullins-Sekerka instability is due to aggregation rate enhancement near the tips, which leads to incremental growth of perturbations of a smooth front. The wavenumber dependence of the growth rate for a harmonic modulation of an interface moving with average velocity $`v`$ is given by $`\gamma _k=v|k|`$. The linear $`k`$-dependence of $`\gamma _k`$ implies that the instability develops first on the smallest scale, in our problem given by the particle size.
To study the ultraviolet cutoff role, i.e., the effect of short distances on the noise and the instability, it is of interest to introduce a parameter in the problem which allows to shift the value of the cutoff scale to values larger than the particle size. One expects that upon doing so both the noise and the instability growth rate will be reduced.
In the mapping model, the noise level can be controlled by altering the shape of aggregating particles. Below we show how by changing the function $`\stackrel{~}{f}_\lambda (z)`$, defined by (5), one can create “flat” particles which are wide along the interface and thin in the growth direction. The reason that noise is suppressed due to using flat particles is the following. In this growth, a particular displacement of the growing cluster boundary amounts to a larger number of layers than in the case of rounded particles used in Ref. . Then, due to statistical averaging over many particle layers the boundary displacement becomes less erratic, and so the noise is reduced. Quantitatively, the noise suppression factor can be estimated as a square root of the particle aspect ratio.
Flat particles can be produced by modifying $`\stackrel{~}{f}_\lambda (z)`$ as follows:
$`\stackrel{~}{f}_{\lambda ,p}(z)`$ $`=`$ $`Wh_{\lambda _p}^1\left(\frac{1}{p}h_{\lambda _p}\left(h_\lambda (z)\right)\right),\lambda _p={\displaystyle \frac{2\lambda }{p+1/p}},`$ (15)
$`W=\left(h_{\lambda _p}^1\left(\frac{1}{p}h_{\lambda _p}(h_\lambda (1))\right)\right)^1,p1.`$
The function $`h_\lambda (z)`$ is defined in (5), and its inverse has the form $`h_\lambda ^1(z)=\frac{1}{2}(z\lambda ^2/z)`$. The factor $`W`$ is introduced in order to have $`\stackrel{~}{f}_{\lambda ,p}(1)=1`$, like for the function $`\stackrel{~}{f}_\lambda (z)`$ defined by (5) above. The resulting function (4) satisfies $`f_{\lambda ,\theta }(\mathrm{})=\mathrm{}`$, which ensures the property $`F_n(\mathrm{})=\mathrm{}`$ for all $`n`$.
The mapping produced by the function (15) is illustrated in Fig. 7. Note that because of $`h_{\lambda _p}(i\lambda )/p=i\lambda _p`$, the square root singularities in $`\stackrel{~}{f}_{\lambda ,p}`$ at $`z=\pm i\lambda `$ are absent for all $`p>1`$. Instead, the mapping composition (15) produces weaker singularities of the form $`(z\pm i\lambda )^{3/2}`$. This smoothens the corners of the particles, as shown in Fig. 7.
Qualitatively, under variation of $`p`$ the particle shape evolves as follows. At $`p=1`$ the mapping $`\stackrel{~}{f}_{\lambda ,p}`$ form (15) coincides with (5). Increasing $`p`$ produces particles with growing aspect ratio, as can be seen from comparing the zoom parts of Figs. 1011 and 8.
To illustrate the effect of $`p`$ on the particle shape, consider the mapping function (15) in the limit $`p1`$. First, one can rewrite (15) as
$$\stackrel{~}{f}_{\lambda ,p}(z)=W^{}\left(h_\lambda (z)\frac{\lambda _p^2(p^21)}{2h_{\lambda _p}(h_\lambda (z))}\right),$$
(16)
where $`W^{}`$ is a prefactor chosen so that $`\stackrel{~}{f}_{\lambda ,p}(1)=1`$. Expanding (16) to lowest order in $`1/p`$, one obtains
$$\stackrel{~}{f}_{\lambda ,p}(z)=z+\frac{\lambda ^2}{p^2}\left(\frac{z+2\sqrt{z^2+\lambda ^2}}{(z+\sqrt{z^2+\lambda ^2})^2}Cz\right),$$
(17)
where $`C=(2h_\lambda (1)1)/h_\lambda ^2(1)`$ and $`h_\lambda `$ is defined in (5). The boundary of the particle produced by $`\stackrel{~}{f}_{\lambda ,p}(z)`$ of the form (17), to lowest order in $`1/p`$, is
$$x=\frac{2}{\lambda ^2p^2}\left(\lambda ^2y^2\right)^{3/2},$$
(18)
where $`x+iy=z`$. The area of this particle is $`3\pi \lambda ^2/4p^2`$. Mapped by $`g^1`$, according to (4), the area is multiplied by a factor equal to $`4`$ at $`\lambda 1`$.
One can use the growth mapping model (3), (7) with the new function (15) to grow clusters in a pretty much the same way as it was done for the model with $`p=1`$ in Section III. The first step is to study the particle areas distribution for the growth with unrestricted areas. The distribution looks similar to that in Fig. 4, containing central peak and the tails corresponding to very large and very small particles. In this case the peak is somewhat wider than for the $`p=1`$ case. However, much of its weight in the distribution $`𝒫(a)`$ is still contained in the window $`[0,3a_{}]`$. Here the “standard area” $`a_{}`$ is defined, by analogy with (9), as the area of a particle grown over a perfectly flat interface. (For $`p1`$ there is no closed form expression for the particle area, like (9), and so one has to calculate $`a_{}`$ numerically.)
As before, at each growth step we choose $`\theta _n`$ randomly, $`0\theta _n<2\pi `$, and calculate the parameter $`\lambda _n`$ using (7), i.e., based on the particle area predicted from the Jacobian $`J_{n1}`$. Then we evaluate the actual area $`a_n`$ of the particle. To inhibit the appearance of large particles, we use the acceptance window $`[0,3a_{}]`$. If $`a_n>3a_{}`$, the particle is not accepted and a new growth attempt is made.
An example of growth with $`p=3`$ and $`\lambda _0=0.2`$ is displayed in Fig. 8. In the inset we zoom on the details of one finger. Note that individual particles are indeed quite flat and are evenly spread over the cluster boundary, indicating reduced noise. The growing interface is overall very smooth, without sharp tips or corners. Also, the fingers are much thicker than for the $`p=1`$ growth (see. Fig. 6).
The cluster size $`R_n`$ is defined by (10). As in Section III, the terms in the product (10) can be evaluated using the relation $`f_{\lambda _p,\theta }^{}(\mathrm{})=1/\stackrel{~}{f}_{\lambda ,p}^{}(1)`$, where
$$\stackrel{~}{f}_{\lambda ,p}^{}(1)=\frac{h_\lambda (1)/\sqrt{1+\lambda ^2}}{\sqrt{h_\lambda ^2(1)+\lambda _p^2}}\frac{h_{\lambda _p}^2(h_\lambda (1))+p^2\lambda _p^2}{h_{\lambda _p}^2(h_\lambda (1))p^2\lambda _p^2}.$$
(19)
In the following Section VI we use (19) along with (10) to evaluate the cluster radius $`R_n`$ and study its scaling.
The appearance of the cluster in Fig. 8 shows that using flat particles indeed helps to reduce statistical fluctuations. In this model, effective averaging of harmonic measure is due to the presence of a tangential-to-boundary length scale set by particle larger dimension. This length scale is controlled by the parameter $`p`$ and becomes large at $`p1`$, if measured in the units of particle size $`\sqrt{a_{}}`$. Noise reduction takes place due to the absence of fluctuations with wavelength smaller than the particle larger dimension, resulting in the shift of the shot noise spectrum cutoff wavenumber from $`2\pi /\sqrt{a_{}}`$ to lower values as the parameter $`p`$ is increased. Because of reduced noise, as compared to the $`p=1`$ case, more agregation events of flat particles are needed to reach a given radius of the cluster.
Averaging over a tangential length scale is somewhat similar to that used in the on-lattice DLA models to simulate surface tension . In these works freezing of random walkers upon each ecounter with the cluster was described by a finite probability $`t<1`$ which could be a function of occupancy of the sites around freezing point. Since freezing of each particle typically takes place after about $`t^1`$ encounters with the cluster, at $`t1`$ these models are characterized by a large length scale over which Laplacian measure is probabilistically averaged. Similarly, the flat particles used in our model can be thought of as a result of averaging over possible particle positions within a finite length scale taken over harmonic measure. Moreover, there is a slight dependence of particle size on their growth position: the particles appearing near the tips are somewhat smaller than those appearing in the concave regions (see Fig. 8). This correlation is consistent with the surface tension interpretation.
The crucial difference, however, is that particle positions in our model are chosen according to the unaltered harmonic measure, whereas in the surface tension models particle freezing depends on local boundary geometry. From that point of view our model is more similar to the multiple hits models in which statistical averaging of the harmonic measure over particle growth attempts is used to control noise. In these models noise reduction is achieved by averaging over independent random walkers with a threshold on the minimal number of visits of each site required before freezing at this site. Since independent walkers arrive at very distant points of cluster boundary, this averaging is not characterized by an additional large length scale and thus bears no resemblance to surface tension.
The models using finite freezing probability $`t<1`$ have been shown to give rise to clusters with thick branches. The Laplacian character of the dynamics and the analogy of the averaging length scale with capillary radius was pointed out and a relation with the Saffman–Taylor problem with surface tension has been conjectured . Because of the large length scale appearing in our averaging scheme, here a similar relation to the problems with surface tension can be conjectured. Indeed, the growth displayed in Fig. 8 looks like a typical fingering pattern observed in the Saffman-Taylor problem with surface tension. As a word of caution, one should realize that all available evidence for the equivalence between the problem with surface tension and our large $`p`$ growth, however similar they appear to be, is rather indirect. The issue of whether or not this growth model is indeed characterized by an effective surface tension and how the latter compares to the noise will be discussed elsewhere.
## V Anisotropic growth model
The iterated mapping model (3), (7) can be generalized to describe spatial anisotropy of local growth rate. Such an anisotropy is characteristic of crystal growth, in which all particles arriving at the crystal-liquid interface have to accomodate to the anisotropic crystal structure .
Anisotropic growth oftenly gives rise to anisotropic irregular fingering patterns called dendrites . The dynamics of dendrites growth obeys scaling laws similar to that of the isotropic Laplacian growth . One of the outstanding theoretical questions is how the scaling exponents depend on the anisotropy.
In this problem, the cluster grows due to spatially isotropic diffusion and aggregation of particles. Thus the quasistationary probability distribution still obeys Eq. (1). The new element, compared to the isotropic model, is that the cluster change due to particle freezing at the boundary depends on the local growth direction $`𝐯`$, $`|𝐯|=1`$. (The unit vector $`𝐯`$ is normal to the boundary.) Accordingly, the probability of joining the cluster becomes a function of $`𝐯`$, and Eq. (2) is replaced by
$$dP=\mathrm{\Omega }(𝐯)|u|dl,dl𝒟_{n1}.$$
(20)
where the function $`\mathrm{\Omega }(𝐯)`$ describes anisotropy.
In order to include anisotropy in the mapping model (3), (7), at the $`n`$-th growth step one has to be able to predict local growth direction $`𝐯_n`$ from particle positions described by randomly chosen angles $`\theta _k`$, $`k=1,2,\mathrm{},n1`$. This is possible because complex-valued Jacobian of a conformal mapping keeps track of the angle change under the mapping. Specifically, consider $`\mathrm{\Theta }_n=\theta _n+\mathrm{arg}J_{n1}`$, where $`J_{n1}`$ is given by (6). Then $`\mathrm{\Theta }_n`$ defines a normal to the cluster boundary, $`𝐯_n=\mathrm{cos}\mathrm{\Theta }_n\widehat{x}+\mathrm{sin}\mathrm{\Theta }_n\widehat{y}`$, at the growth point $`F_{n1}(e^{i\theta _n})`$.
Now, there are several possible ways to account for the growth anisotropy. For instance, one can introduce the anisotropy by making $`\lambda _n`$ a function of $`\mathrm{\Theta }_n`$, e.g., $`\lambda _n\mathrm{\Omega }^{1/2}(𝐯_n)`$. Another way is to introduce acceptance probability for the particles which depends on $`\mathrm{\Theta }_n`$ in some way. In the simulations reported below we use an acceptance window for $`M\mathrm{\Theta }_n`$ with $`M=3,4,\mathrm{}`$, corresponding to the growth with an $`M`$-fold symmetry. Only particles with $`\mathrm{\Theta }_n`$ such that
$$\theta _{\mathrm{max}}M\mathrm{\Theta }_n\theta _{\mathrm{max}}$$
(21)
are accepted. Here $`\theta _{\mathrm{max}}`$ is a parameter in the interval $`[0,\pi ]`$ controlling the degree of anisotropy. Small values $`\theta _{\mathrm{max}}\pi `$ correspond to highly anisotropic growth, whereas fully isotropic growth is recovered in the limit $`\theta _{\mathrm{max}}\pi `$.
Other aspects of the simulation are the same as in Section IV. We employed the elementary mapping $`\stackrel{~}{f}_{\lambda ,p}(z)`$ of the form (15) with the noise level controlled by the parameter $`p1`$. Particles with large areas were eliminated using the acceptance window $`[0,3a_{}]`$ defined in Section IV. An example of growth with the three-fold symmetry ($`M=3`$) is shown in Fig. 9. In this case, we used $`\lambda _0=0.8`$, $`p=1.5`$, and $`\theta _{\mathrm{max}}=\mathrm{cos}^1(0.9)0.451`$. The cluster is characterized by overall symmetric main branches covered with numerous side branches.
In our model one has a separate control over the degree of anisotropy and over noise, via the parameters $`\theta _{\mathrm{max}}`$ and $`p`$. This is convenient for studying the effects of noise on the ordering of branches in dendrites. To illustrate that, we compare two growths with four-fold symmetry, displayed in Figs. 1011. The cluster in Fig. 10 is obtained using the $`p=1`$ model without noise reduction, as described in Sections IIIII. The cluster in Fig. 11 is grown using the noise reduced model of Section IV with the parameter $`p=2`$. In both cases, we use the same anisotropy parameter: $`\theta _{\mathrm{max}}=\mathrm{cos}^1(0.95)0.318`$. One notes high anisotropy of the growth present at small scales in both cases, which is significantly suppressed at larger scales for the noisy growth with $`p=1`$ (see Fig. 10). However, the $`p=2`$ growth with low noise remains very anisotropic at all scales (see Fig. 11).
It is known from studies of on-lattice DLA models that noise, no matter how strong, gives way to anisotropy at sufficiently large scales . We thus expect that a similar effect may take place in the noisy growth with $`p=1`$, making the growth shown in Fig. 10 at larger scales looking like that in Fig. 11.
In agreement with the studies of the off-lattice DLA models , we observed that the dendritic growth with the symmetry of order $`M=3,\mathrm{\hspace{0.17em}4}`$ is much more stable with respect to noise than that with $`M=5,\mathrm{\hspace{0.17em}6}`$ or higher. Scaling properties of the anisotropic growth will be studied in Section VI.
## VI Scaling properties
Scaling of $`R_n`$ for all growth models introduced above is studied here using the following procedure. The cluster radius $`R_n`$ obtained from (10), (11), (19), is plotted against the cluster area $`A_n`$, evaluated as a sum of individual particle areas $`a_n`$. Asymptotically, at large $`n`$, one has $`R_nA_n^{1/d}`$. To determine $`d`$ more accurately we optimize initial conditions of the growth, represented in our model by nondimensionalized particle size $`\lambda _0`$, as described below.
In the log–log plot of $`R_n`$ versus $`A_n`$ one can clearly distinguish two regimes, initial growth and developed, or regular growth, characterized by somewhat different slopes of the corresponding parts of the $`\mathrm{ln}R`$ vs. $`\mathrm{ln}A`$ curves. Geometrical meaning of these regimes is as follows. For isotropic growth, with or without noise suppression, the cluster initially consists of branches growing practically independently. Later, at the regular growth stage, the number of main branches is reduced to four or five, all interacting and competing with each other. For anisotropic growth with the $`M`$-fold symmetry the number of main branches is $`M`$ at all stages of growth. Regular growth in this case is distinguished by many fingers appearing on the sides of $`M`$ main branches.
The initial stage is more pronounced when the particle size, determined by the value of $`\lambda _0`$, is much smaller than the unit circle from which the growth starts. Since we are interested in the regular growth scaling, in each case studied we tried to optimize the value of $`\lambda _0`$ to shorten the initial growth stage, carefully checking that the variation of $`\lambda _0`$ has no detectable effect on the asymptotic slope of the $`\mathrm{ln}R`$ vs. $`\mathrm{ln}A`$ curve. The benefit of shortening the initial growth stage is that, at constant number of particles, it leads to longer regular growth and thus allows to extract the scaling exponent with higher precision. The resulting curves are presented in Fig. 12, as described in the figure caption and below.
The optimal value of $`\lambda _0`$ determined for the isotropic growth with $`p=1`$ is close to $`\lambda _0=0.8`$. For the scaling analysis we used the growth displayed in Fig. 6, in which $`N=17545`$, $`\lambda _0=0.8`$, $`p=1`$. In Fig. 12, it corresponds to the lowest of the curves marked by $`\mathrm{a}_1`$. To eliminate the effect of fluctuations at the initial stage of the growth, we also generated the curves $`\mathrm{a}_2`$, $`\mathrm{a}_3`$, and $`\mathrm{a}_4`$, by averaging $`\mathrm{ln}R`$ over 5, 10, and 50 growth realizations with $`N=5000`$, $`1000`$, and $`200`$ time steps, respectively.
For isotropic growth with reduced noise we analyzed two growths with $`p=3`$:
curve $`\mathrm{b}_1`$ with $`N=11611`$, $`\lambda _0=0.1`$;
curve $`\mathrm{b}_2`$ with $`N=15043`$, $`\lambda _0=0.2`$.
The curve $`\mathrm{b}_2`$ corresponds to the growth displayed in Fig. 8. At low noise, the fluctuations of $`R_n`$ are quite small, which makes additional averaging over realizations unnecessary.
Note that for the isotropic growth models the strategy of optimizing $`\lambda _0`$ works quite well, allowing to almost entirely eliminate the initial growth region. The scaling dimension found from the slope of the best straight line fits is close to $`1.7`$. To study the deviation from $`1.7`$, we subtract from all curves the linear function $`\mathrm{ln}R=\mathrm{ln}A/1.7`$ and plot the result in the lower part of Fig. 12. Note that upon this subtraction the curves for isotropic growth, with or without noise suppression, become nearly perfectly horizontal. Estimate of the deviation from the best horizontal line fit shows that the value $`1.7`$ is accurate within $`1\%`$.
For anisotropic models, we consider three different growths:
curve $`\mathrm{c}_1`$ with $`N=10146`$, $`\lambda _0=0.8`$, $`p=1.5`$, $`M=3`$;
curve $`\mathrm{c}_2`$ with $`N=7635`$, $`\lambda _0=0.3`$, $`p=1`$, $`M=4`$;
curve $`\mathrm{c}_3`$ with $`N=6782`$, $`\lambda _0=0.8`$, $`p=2`$, $`M=4`$.
These curves correspond to the growths displayed in Figs. 9, 10, and 11, respectively. As above, we subtract the linear function $`\mathrm{ln}R=\mathrm{ln}A/1.7`$. However, after this subtraction, the curves $`\mathrm{c}_1`$ and $`\mathrm{c}_3`$ retain some residual slope. Estimating it, we conclude that the best value for the fractal dimension is $`d1.5`$ for $`\mathrm{c}_1`$, and $`d1.62`$ for $`\mathrm{c}_3`$. The latter value agrees with the values $`d1.58`$ and $`d1.63`$ of the growth with $`M=4`$ reported in Refs. .
For the curve $`\mathrm{c}_2`$ corresponding to anisotropic growth with noise, after subtracting $`\mathrm{ln}R=\mathrm{ln}A/1.7`$ we do not find any significant residual slope. It is possible, however, that the dimension $`1.7`$ corresponds to the crossover regime and changes to a lower value at larger $`N`$. Similar behavior is known to take place in the on-lattice DLA growth , where the dimension $`1.7`$ observed for not very large clusters crosses over to $`1.63`$ at $`N410^6`$.
To understand possible sources of errors in determining the fractal dimension from $`\mathrm{ln}R`$ vs. $`\mathrm{ln}A`$ curves, here we consider how $`R_n`$ and $`A_n`$ fluctuate with $`n`$. The fluctuations of $`\mathrm{ln}R_n`$ gradually decrease with increasing $`n`$, as can be clearly seen in the lower panel of Fig. 12. A convenient way to analyze fluctuations is to plot pairs $`(\mathrm{ln}R_N,\mathrm{ln}A_N)`$ for particular $`N`$, repeating growth many times. In Fig. 13 we present results for $`10^3`$ growth samples and several values of $`N`$. The resulting clouds become more compact as $`N`$ increases, indicating that the fluctuations of $`\mathrm{ln}R_N`$ and $`\mathrm{ln}A_N`$ are decreasing.
Let us first discuss fluctuations of $`\mathrm{ln}A_n`$. The total area $`A_n`$ is the sum of individual particles areas $`a_k`$, $`k=1,2,\mathrm{},n`$. Assuming that the areas $`a_k`$ are independent or, more precisely, have only short correlations, one obtains a Gaussian distribution with the variance $`n`$. (As we argue below, there exist long negative correlation of particle areas, which may further reduce fluctuations of $`A_n`$.) The fluctuations of $`\mathrm{ln}A_n`$ are simply given by relative fluctuations $`\delta A_n/A_n`$, which means that for large $`n`$ the distribution of $`\mathrm{ln}A_n`$ is also Gaussian, with the variance $`n^1`$.
On the other hand, the radius $`R_n`$ is a product (10) of stretching factors $`J_k^{(\mathrm{})}=f_k^{}(z\mathrm{})`$. Since $`J_k^{(\mathrm{})}>1`$ for all $`k`$, the quantity $`R_n`$ grows monotonically, so that $`R_nA_n^{1/d}`$ at large $`n`$. Thus the noise due to fluctuations of $`J_k^{(\mathrm{})}`$ is of a multiplicative nature. One can write
$$\mathrm{ln}R_n=\underset{k=1}{\overset{n}{}}\mathrm{ln}f_k^{}(z\mathrm{}),$$
(22)
which suggests that the distribution of $`\mathrm{ln}R_n`$ is Gaussian, i.e., the distribution of $`R_n`$ is log-normal. Indeed, the log-normal fit perfectly describes the statistics of $`R_n`$, as demonstrated in Fig. 14. However, attempting a Gaussian fit produces an asymmetric distribution deviating from the observed distribution of $`\mathrm{ln}R_n`$. Thus, even though the relative fluctuations of $`R_n`$ are small, the statistics is best described as log-normal.
Naively, Eq. (22) implies growth of the variance of $`\mathrm{ln}R_n`$ with $`n`$. However, Figs. 1314 demonstrate that, on the contrary, the width of $`\mathrm{ln}R_n`$ distribution is decreasing with $`n`$. This nontrivial behavior was first mentioned, without explanation, in Ref. .
To rationalize the observed sharpening of the distribution of $`R_n`$, one can argue as follows. We note that the dynamics of $`R_n`$ is characterized by a negative feedback. Consider growth of a cluster which at the $`n`$-th step has a radius smaller than average. Then the Jacobian of $`F_n`$ is typically smaller than its mean value at this number of particles. In this case, according to (7), subsequently growing particles will have larger $`\lambda _k`$’s, and thus larger areas, until the cluster radius will approach the average value. The evolution of a cluster which at a certain step has a radius larger than average can be considered in a similar way. This long-time anticorrelation of $`\lambda _k`$’s suppresses the fluctuations of $`R_n`$. Also, it produces long negative correlation of particle areas.
## VII Summary
To conclude, growth models using conformal mappings have large flexibility allowing for independent control over noise and growth anisotropy. We generalized the model by using flat particles to suppress noise. It is essential that these models lead to an intrinsically isotropic growth with reduced noise, in contrast with other previously studied models. Also, we demostrated that favoring growth in certain directions can be used to simulate anisotropy of the growth rate.
Having separate control on the noise and anisotropy, we have been able to analyze their effect on scaling properties. We found that the fractal dimension $`d=1.7`$, universally for any isotropic growth, regardless of the noise level. However, the fractal dimension is somewhat reduced in the presence of anisotropy.
It was assumed that particle size fluctuations, present in the conformal mapping model, are insignificant. We observed that the growth rules used in Ref. lead to occasional appearance of exceptionally large particles. We have shown that by augmenting the model with an area acceptance criterion this problem is fixed.
Clearly, more work has to be done to establish a relation of the introduced models with real physical processes, like viscous fingering or dendritic crystal growth. Another interesting open question is how to introduce an effective surface tension.
## ACKNOWLEDGMENTS
We thank G.E. Falkovich, M.B. Hastings, V.A. Kazakov, V.V. Lebedev, B.Z. Spivak, and P.B. Wiegmann for useful discussions. Hospitality of Weizmann Institute made our collaboration possible. This work was partially supported by the Minerva Foundation (M.S.) and by the NSF Award 67436000IRG (L.L.).
## A Insides of the particle area distribution
Here we discuss in more detail the distribution of particle areas. The main feature manifest in the area histogram plotted in Fig. 4 is a sharp asymmetric peak at $`2.1a_{}`$. This peak has its origin in the dependence of particle size on the growth point.
The argument is as follows. First we note that the growth is taking place predominantly at the tips of the branches. Because of that, for several particles growing on each other, there is a tendency to preserve growth direction. This leads to formation of relatively long chains of particles growing in a particular direction, clearly seen in the inset of Fig. 6. The chains are mostly formed at the tips of outer branches.
Now, consider a particle growing near one of the tips. The area of this particle has some dependence on the position of the growth point relative to the tip. The peak in the histogram in Fig. 4 is explained if one assumes that the particle area has a local minimum in the forward growth direction. The minimum in the area leads to a caustic in the histogram. Ideally, this would have produced an asymmetric square root singularity with probability equal zero on the left side. Because of particle size variation among branches, the singularity is smeared into a peak.
To verify the above assumption, we consider areas for the first few particles grown on the $`|z|=1`$ circle with the parameter $`\lambda _0=0.2`$. The area of the very first particle is close to $`1.2a_{}`$ and, according to (9), is independent of its position. The area of the second particle $`a_2`$ depends on its position $`\theta `$ relative to the first particle, as shown by a solid line in Fig. 15(a). Note that the area is the same as that of the first particle when the particles are far apart, $`\theta \lambda _0`$, and is overall substantially larger when the particle overlap, $`\theta \lambda _0`$. Partially, this is explained by the dependence (9) of particle area on the circle radius. (Assuming that $`a_2(\theta \lambda _0)`$ can be crudely estimated by (9) with $`r=\lambda _0`$.) Another effect that contributes to the area $`a_2`$ increase for $`\theta \lambda _0`$ is the variation of the Jacobian as a function of $`\theta `$, leading to additional stretching of the second particle.
The feature in Fig. 15(a) which is of interest in connection to the peak in the area distribution $`𝒫(a)`$ is the minimum of $`a_2(\theta )`$ at $`\theta =0`$. Translated to the histogram of areas, it leads to a caustic described by a square root singularity. However, as a possible explanation of the peak in Fig. 4 this is only partially satisfying, since one has to understand why similar caustics due to the two maxima of $`a_2(\theta )`$ are not observed in Fig. 4.
The reason for the difference between the effects of maxima and minima can be seen from a comparison with the case of three and four particles. Consider the situation when the second particle is centered exactly on the first particle, and the third particle is grown at an angular position $`\theta `$ relative to the first two particles. The area of the third particle $`a_3(\theta )`$ is plotted in Fig. 15(a) in a dashed line. Note that, since the curvature at the minimum of $`a_3(\theta )`$ at $`\theta =0`$ is much smaller than for $`a_2(\theta )`$, the corresponding caustic in $`𝒫(a)`$ will be much stronger. On the other hand, the curvature at the maxima of $`a_3(\theta )`$ is about the same as that for $`a_2(\theta )`$. Both observations remain correct for any number of particles. To illustrate this we plot the area $`a_4(\theta )`$ of the fourth particle in the presence of three particles grown exactly on the top of each other — see dotted line in Fig. 15(a).
Another notable feature in the plots of $`a_{2,3,4}(\theta )`$ is that the area becomes much smaller than $`a_{}`$, approaching zero near certain values of $`\theta `$. This behavior is related to the growth near particle corners, which are the points of divergence of the Jacobian. According to (7), larger Jacobian translates into smaller particle area. The particles growing near corners form the tail of the area distribution $`𝒫(a)`$ at small areas $`aa_{}`$. The behavior of $`𝒫(a)`$ in this tail, $`𝒫(a)a^{1/2}`$, follows from the square root divergence of the Jacobian at particle corners. The slope $`1/2`$ is clearly seen in the $`\mathrm{ln}𝒫`$ vs. $`\mathrm{ln}a`$ plot in Fig. 5.
The features in $`a_{2,3,4}(\theta )`$ discussed above evolve in an interesting way for the models with lower noise corresponding to $`p>1`$ — see Fig. 15(b). The plots of $`a_{2,3,4}(\theta )`$ in this figure are produced for the model with $`\lambda _0=0.2`$ and $`p=3`$ in the same way as above for $`p=1`$. Note that relative changes of the area as a function of $`\theta `$ are smaller than for $`p=1`$. One reason for this is in weaker curvature variation for flat particles, which makes particle area less sensitive to the growth point position. Another reason is that at $`p>1`$ the particles corners have no cusps, and thus particles with small areas do not appear.
## B Discussion of the numerical method
Here we comment on the optimal choice of the numerical procedure. First, since the areas of new particles are evaluated before the particles are accepted, one could, instead of eliminating large particles, change the growth algorithm so that all particles areas become equal. This can be achieved by adjusting the parameter $`\lambda _n`$ for each particle until its area converges to a given value. Although this would somewhat slow down the speed of computation, an obvious gain would be in a more direct relation with the standard DLA growth.
Also, one could attempt at increasing the speed and efficiency of the growth algorithm by introducing in it various improvements:
(i) Coarsening of the mappings which correspond to particles sufficiently deep in the stagnation regions. It was demonstrated in Ref. that an accurate envelope of the cluster can be obtained by using truncated Laurent series of $`F_n(z)`$. One can implement this observation as follows. At the growth step $`n`$ choose some $`1<m<n`$ in such a way that all particles with the numbers $`m`$ are located sufficiently deep inside the stagnation part of the cluster. Then one can replace the mapping $`F_n=f_1\mathrm{}f_n`$ by
$$F_n^{(\mathrm{approx})}=\left[f_1\mathrm{}f_m\right]_{\mathrm{truncated}}f_{m+1}\mathrm{}f_n,$$
(B1)
where the mapping in parentheses which is replaced by truncated series is nothing but $`F_m`$. One can choose $`m`$ so that the finite series representation of the mapping $`F_m(z)`$ is accurate for $`z`$ in the active growth region. By this trick, instead of computing a composition of $`n`$ functions, one has to deal with only $`nm`$ functions at each growth step. Since at large $`n`$ most of the particles are in the stagnation regions, one may have $`nmn`$.
(ii) Evaluating the particle area with lower precision. We used several hundred points on each particle’s boundary, which produces areas accurate within $`0.1\%`$. Practically, such a high precision may not be necessary. Instead, one can predict particle areas by estimating the Jacobian at several points chosen within the bump according to some rule or randomly.
(iii) Using an area acceptance window to discriminate against very small particles with areas $`a_{}`$. These particles practically do not change the structure of the cluster branches, except near the corners between adjacent particles. However, due to the presence of small particles additional mappings appear in the composition sequence $`f_n\mathrm{}f_1`$, which slows down the computation.
We have not used these procedures in the simulations described above and neither systematically studied their efficiency. We felt that, at the initial stage, keeping growth algorithm as precise and simple as possible, even at the price of somewhat slowing it down, makes the results more solid.
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# Cosmological two-fluid thermodynamics
## 1 Introduction
The thermodynamics of two fluids with different temperatures represents a framework which is sufficiently general to apply to entirely different epochs of the cosmological evolution. This unifying feature may be used to establish surprising similarities between otherwise quite independent phenomena in the expanding universe. In this paper we focus on aspects of the temperature evolution during periods with decay and production of particles to demonstrate the universal power of the thermodynamic description. In particular, we show that the same simple law for the cooling rate of a fluid in the expanding universe governs a wide range of phenomena implying the cosmological electron-positron annihilation after neutrino decoupling at about 1 MeV, the evaporation of primordial black holes (PBHs) with a narrow mass range, and the “deflationary” transition from an initial de Sitter stage to a subsequent FLRW period, equivalent to a phenomenological vacuum decay model. All these processes are characterized by a strong back reaction of decay and production processes on the thermal evolution of the universe. It is the possibility of taking into account this back reaction in a rather straightforward but general way, which admits an application to such a variety of different phenomena.
More specifically, we shall first reproduce the factor $`\left(11/4\right)^{1/3}`$ by which the temperature of the neutrino background differs from that of the photon background as a consequence of electron-positron annihilation. Secondly, we show that the black hole temperature behavior $`T_{_{\left(BH\right)}}m_{_{\left(BH\right)}}^1`$, where $`m_{_{\left(BH\right)}}`$ is the black hole mass, is consistent with the general fluid temperature law for a PBH “fluid”, a configuration in which all its members are assumed to have the same mass $`m_{_{\left(BH\right)}}`$. On this basis we discuss thermodynamical aspects of PBH evaporation. The third example is the evolution of the radiation temperature in a “deflationary” scenario of the early universe which implies an initial increase to a maximum value as a result of the production of relativistic particles out of a decaying vacuum, followed by a decrease which finally approaches the familiar FLRW behavior.
None of these results is really new. The first case is cosmological textbook physics (see, e.g. ), the second one was investigated in , the scenario characterizing the third case is based on (see also ). What is new however, is the unifying view which allows the discussion of so different cosmological effects starting from the same set of basic equations. To highlight the underlying common thermodynamical features of the mentioned phenomena is the main purpose of this paper. The material is organized as follows. In section 2 we recall the basic relations of two-fluid thermodynamics in an expanding universe. These relations are used in section 3 to discuss the cosmological electron-positron annihilation. In section 4 the general formalism is applied to a mixture of radiation and a component of PBHs which are assumed to have the same mass. It may be shown that under this condition they share essential properties with a pressureless gas. Thermodynamic aspects of a smooth transition from a de Sitter stage to a radiation dominated FLRW phase including an intermediate temperature maximum are investigated on the same general basis in section 5, while the final section 6 is devoted to concluding remarks. Units have been chosen so that $`c=k_B=\mathrm{}=1`$.
## 2 Basic thermodynamic relations
We assume the cosmic medium to consist of two components which share the same 4-velocity $`u^i`$. Each of the components has a perfect fluid structure with the energy-momentum tensor $`T_{_{\left(A\right)}}^{ik}`$, where $`A=1,2`$, and a corresponding particle flow vector $`N_{_{\left(A\right)}}^i`$,
$$T_{_{\left(A\right)}}^{ik}=\rho _{_{\left(A\right)}}u^iu^k+p_{_{\left(A\right)}}h^{ik}\text{ ,}N_{_{\left(A\right)}}^i=n_{_{\left(A\right)}}u^i\text{ , }\text{ }\text{ }(A=1,2).$$
(1)
Here, $`\rho _{_{\left(A\right)}}`$ is the energy density of component $`A`$, measured by a comoving observer, $`p_{_{\left(A\right)}}`$ is the corresponding equilibrium pressure, $`h^{ik}=g^{ik}+u^iu^k`$ is the spatial projection tensor, and $`n_{_{\left(A\right)}}`$ is the number density of species-$`A`$ particles. Neither $`T_{_{\left(A\right)}}^{ik}`$ nor $`N_{_{\left(A\right)}}^i`$ are required to be conserved, i.e., interactions and interparticle conversions are admitted:
$$T_{{}_{\left(A\right)}{}^{};k}^{ik}=t_{_{\left(A\right)}}^i\text{ , }\text{ }\text{ }\text{ }N_{{}_{\left(A\right)}{}^{};i}^i=\dot{n}_{_{\left(A\right)}}+3Hn_{_{\left(A\right)}}=n_{_{\left(A\right)}}\mathrm{\Gamma }_{_{\left(A\right)}}.$$
(2)
The quantity $`H`$ is the Hubble parameter $`H=\dot{a}/a`$ with the scale factor $`a`$ of the Robertson-Walker metric. $`\mathrm{\Gamma }_{_{\left(A\right)}}\dot{N}_{_{\left(A\right)}}/N_{_{\left(A\right)}}`$ denotes the rate of change of the number $`N_{_{\left(A\right)}}n_{_{\left(A\right)}}a^3`$ of particles in a comoving volume $`a^3`$. The $`T_{_{\left(A\right)}}^{ik}`$ and $`N_{_{\left(A\right)}}^i`$ add up to the corresponding quantities for the medium as a whole:
$$T^{ik}=T_{_{\left(1\right)}}^{ik}+T_{_{\left(2\right)}}^{ik}\text{ , }\text{ }\text{ }N^i=N_{_{\left(1\right)}}+N_{_{\left(2\right)}}.$$
(3)
It is well known that in general the energy-momentum tensor $`T^{ik}`$ does not take the form of a perfect fluid, but will contain a non-equilibrium pressure $`\mathrm{\Pi }`$ . Different from the $`T_{_{\left(A\right)}}^{ik}`$, the overall energy-momentum tensor has to be conserved, which establishes a relation between $`t_{_{\left(1\right)}}^i`$ and $`t_{_{\left(2\right)}}^i`$:
$$T^{ik}=\rho u^iu^k+\left(p+\mathrm{\Pi }\right)h^{ik},\text{ }\text{ }T_{;k}^{ik}=0t_{_{\left(1\right)}}^i=t_{_{\left(2\right)}}^i.$$
(4)
We do not require, however, conservation of the total particle number , i.e.,
$$N_{;a}^a=\dot{n}+3Hn=n\mathrm{\Gamma }n_{_{\left(1\right)}}\mathrm{\Gamma }_{_{\left(1\right)}}+n_{_{\left(2\right)}}\mathrm{\Gamma }_{_{\left(2\right)}},$$
(5)
where $`nn_{_{\left(1\right)}}+n_{_{\left(2\right)}}`$ is the overall particle number density. Each component is governed by its own Gibbs equation which provides us with an expression for the time behaviour of the entropy per particle $`s_{_{\left(A\right)}}`$,
$$T_{_{\left(A\right)}}\text{d}s_{_{\left(A\right)}}=\text{d}\frac{\rho _{_{\left(A\right)}}}{n_{_{\left(A\right)}}}+p_{_{\left(A\right)}}\text{d}\frac{1}{n_{_{\left(A\right)}}}\text{ , }n_{_{\left(A\right)}}T_{_{\left(A\right)}}\dot{s}_{_{\left(A\right)}}=u_at_{_{\left(A\right)}}^a\left(\rho _{_{\left(A\right)}}+p_{_{\left(A\right)}}\right)\mathrm{\Gamma }_{_{\left(A\right)}}\text{ .}$$
(6)
In general, the temperatures $`T_{_{\left(A\right)}}`$ of both components are different. With the help of the equations of state
$$p_{_{\left(A\right)}}=p_{_{\left(A\right)}}(n_{_{\left(A\right)}},T_{_{\left(A\right)}}),\text{ }\rho _{_{\left(A\right)}}=\rho _{_{\left(A\right)}}(n_{_{\left(A\right)}},T_{_{\left(A\right)}})\text{ , }$$
(7)
one obtains the evolution law for the temperatures $`T_{_{\left(A\right)}}`$. Namely, differentiating $`\rho _{_{\left(A\right)}}`$ in (7) along the fluid flow lines and applying the balances (2) we find
$$\frac{\dot{T}_{_{\left(A\right)}}}{T_{_{\left(A\right)}}}=3H\left(1\frac{\mathrm{\Gamma }_{_{\left(A\right)}}}{3H}\right)\frac{p_{_{\left(A\right)}}}{\rho _{_{\left(A\right)}}}+\frac{n_{_{\left(A\right)}}\dot{s}_{_{\left(A\right)}}}{\rho _{_{\left(A\right)}}/T_{_{\left(A\right)}}}\text{ ,}$$
(8)
where
$$\frac{p_{_{\left(A\right)}}}{\rho _{_{\left(A\right)}}}\frac{\left(p_{_{\left(A\right)}}/T_{_{\left(A\right)}}\right)_{n_{_{\left(A\right)}}}}{\left(\rho _{_{\left(A\right)}}/T_{_{\left(A\right)}}\right)_{n_{_{\left(A\right)}}}},\text{ }\frac{\rho _{_{\left(A\right)}}}{T_{_{\left(A\right)}}}\left(\frac{\rho _{_{\left(A\right)}}}{T_{_{\left(A\right)}}}\right)_{n_{_{\left(A\right)}}}.$$
The general temperature law (8) provides the unifying basis for the discussions of the following sections. It will play a central role in our investigations of both the electron-positron annihilation and the PBH evaporation and a specific inflationary scenario.
An important special case which we frequently will refer to is characterized by the condition $`\dot{s}_{_{\left(A\right)}}=0`$, which means constant entropy per particle . This condition simplifies the temperature law,
$$\dot{s}_{_{\left(A\right)}}=0\frac{\dot{T}_{_{\left(A\right)}}}{T_{_{\left(A\right)}}}=3H\left(1\frac{\mathrm{\Gamma }_{_{\left(A\right)}}}{3H}\right)\frac{p_{_{\left(A\right)}}}{\rho _{_{\left(A\right)}}}\text{ .}$$
(9)
Moreover, according to (6) it establishes the link $`u_at_{_{\left(A\right)}}^a=\left(\rho _{_{\left(A\right)}}+p_{_{\left(A\right)}}\right)\mathrm{\Gamma }_{_{\left(A\right)}}`$ between the source terms $`\mathrm{\Gamma }_{_{\left(A\right)}}`$ and $`t_{_{\left(A\right)}}^i`$ which together with the last relation of (4) provides us with a relation between the rates $`\mathrm{\Gamma }_{_{\left(1\right)}}`$ and $`\mathrm{\Gamma }_{_{\left(2\right)}}`$:
$$n_{_{\left(2\right)}}\mathrm{\Gamma }__2=\frac{h_{_{\left(1\right)}}}{h_{_{\left(2\right)}}}n_{_{\left(1\right)}}\mathrm{\Gamma }_{_{\left(1\right)}}\dot{N}_{_{\left(2\right)}}=\frac{h_{_{\left(1\right)}}}{h_{_{\left(2\right)}}}\dot{N}_{_{\left(1\right)}}\text{ }\text{ }\text{ }\text{ }\text{ }\text{ }\left(\dot{s}_{_{\left(A\right)}}=0\right),$$
(10)
where $`h_{_{\left(A\right)}}\left(\rho _{_{\left(A\right)}}+p_{_{\left(A\right)}}\right)/n_{_{\left(A\right)}}`$ are the enthalpies per particle.
## 3 Electron-positron annihilation
Let us consider the cosmological period of electron-positron annihilation and the corresponding creation of photons, shortly after neutrino decoupling (see, e.g. , §3.1.2). Annihilation becomes predominant as soon as the radiation temperature drops below the mass of the electron. Before the electron-positron annihilation there are two bosonic degrees of freedom (photons), four fermionic degrees of freedom due to the electrons and positrons, and 12 fermionic degrees of freedom due to the different neutrino species. The fermionic energy density is
$$\rho _F=\frac{7}{6}\frac{\pi ^4}{30\zeta \left(3\right)}n_FT_F\text{ , }\text{ }n_F=\frac{3}{4}\frac{\zeta \left(3\right)}{\pi ^2}g_FT_F^3\text{ , }$$
(11)
where $`\zeta \left(x\right)`$ is Riemann’s Zeta-function, and the bosonic one
$$\rho _B=\frac{\pi ^4}{30\zeta \left(3\right)}n_BT_B\text{ , }\text{ }n_B=\frac{\zeta \left(3\right)}{\pi ^2}g_BT_B^3\text{ .}$$
(12)
The factors $`g_F`$ and $`g_B`$ are the numbers of fermionic and bosonic degrees of freedom, respectively.
After the electron-positron annihilation we are left with the two bosonic degrees of freedom (photons) and the 12 neutrino degrees of freedom, i.e., four fermionic degrees of freedom have disappeared. The neutrino degrees of freedom are not affected at all by this process. The neutrino temperature behaves as $`T_{_{\left(\nu \right)}}=T_{_{\left(\nu \right)}}\left(t_0\right)a\left(t_0\right)/a\left(t\right)`$, where the initial time $`t_0`$ is assumed to be a time before the beginning of the annihilation process, i.e., electrons, positrons and photons are still at equilibrium at $`t_0`$ with $`T_{_{\left(\nu \right)}}\left(t_0\right)=T_{_{\left(e^\pm \right)}}\left(t_0\right)=T_{_{\left(\gamma \right)}}\left(t_0\right)=T_{_{\left(0\right)}}`$. Let us consider the subsystem of four fermionic and two bosonic degrees of freedom . The four fermionic degrees of freedom due to the electrons and positrons are identified with fluid 1 of our general analysis, i.e., $`\left(1\right)\left(e^\pm \right)`$ while the two bosonic degrees of freedom due to the photons are fluid 2, i.e., $`\left(2\right)\left(\gamma \right)`$. With these specifications the number densities become
$$n_{_{\left(e^\pm \right)}}=4\frac{3}{4}\frac{\zeta \left(3\right)}{\pi ^2}T_{_{\left(e^\pm \right)}}^3\text{ , }\text{ }n_{_{\left(\gamma \right)}}=2\frac{\zeta \left(3\right)}{\pi ^2}T_{_{\left(\gamma \right)}}^3\text{ .}$$
(13)
The corresponding enthalpies per particle are
$$h_{_{\left(e^\pm \right)}}=\frac{14}{9}\frac{\pi ^4}{30\zeta \left(3\right)}T_{_{\left(e^\pm \right)}}\text{ , }\text{ }h_{_{\left(\gamma \right)}}=\frac{4}{3}\frac{\pi ^4}{30\zeta \left(3\right)}T_{_{\left(\gamma \right)}}\text{ .}$$
(14)
Assuming $`h_{_{\left(e^\pm \right)}}/h_{_{\left(\gamma \right)}}`$ to be given by their (constant) initial ratio, Eq. (10) integrates to
$$N_{_{\left(\gamma \right)}}\left(t\right)=N_{_{\left(\gamma \right)}}\left(t_0\right)+\frac{h_{_{\left(e^\pm \right)}}}{h_{_{\left(\gamma \right)}}}\left[N_{_{\left(e^\pm \right)}}\left(t_0\right)N_{_{\left(e^\pm \right)}}\left(t\right)\right].$$
(15)
The final value $`N_{_{\left(\gamma \right)}}\left(t_f\right)`$ corresponds to the case where the electrons and positrons have been annihilated, i.e., $`N_{_{\left(e^\pm \right)}}\left(t_f\right)=0`$:
$$\frac{N_{_{\left(\gamma \right)}}\left(t_f\right)}{N_{_{\left(\gamma \right)}}\left(t_0\right)}=1+\frac{h_{_{\left(e^\pm \right)}}}{h_{_{\left(\gamma \right)}}}\frac{n_{_{\left(e^\pm \right)}}\left(t_0\right)}{n_{_{\left(\gamma \right)}}\left(t_0\right)}.$$
(16)
This result has to be coupled to the temperature law (9), which for photons becomes
$$\frac{\dot{T}_{_{\left(\gamma \right)}}}{T_{_{\left(\gamma \right)}}}=\frac{\dot{a}}{a}+\frac{1}{3}\frac{\dot{N}_{_{\left(\gamma \right)}}}{N_{_{\left(\gamma \right)}}}T_{_{\left(\gamma \right)}}\left(t\right)=T_{_{\left(0\right)}}\frac{a\left(t_0\right)}{a\left(t\right)}\left(\frac{N_{_{\left(\gamma \right)}}\left(t\right)}{N_{_{\left(\gamma \right)}}\left(t_0\right)}\right)^{1/3}.$$
(17)
For $`tt_f`$ the ratio $`N_{_{\left(\gamma \right)}}\left(t\right)/N_{_{\left(\gamma \right)}}\left(t_0\right)`$ is fixed by the value $`N_{_{\left(\gamma \right)}}\left(t_f\right)/N_{_{\left(\gamma \right)}}\left(t_0\right)`$. Since $`h_{_{\left(e^\pm \right)}}/h_{_{\left(\gamma \right)}}=7/6`$ \[cf. (14)\] and $`n_{_{\left(e^\pm \right)}}\left(t_0\right)/n_{_{\left(\gamma \right)}}\left(t_0\right)=3/2`$ \[cf. (13)\] we obtain $`N_{_{\left(\gamma \right)}}\left(t_f\right)/N_{_{\left(\gamma \right)}}\left(t_0\right)=11/4`$. Consequently, the temperature evolution law for $`tt_f`$ is
$$T_{_{\left(\gamma \right)}}\left(t\right)=T_{_{\left(0\right)}}\frac{a\left(t_0\right)}{a\left(t\right)}\left(\frac{11}{4}\right)^{1/3}\frac{T_{_{\left(\gamma \right)}}\left(t\right)}{T_{_{\left(\nu \right)}}\left(t\right)}=\left(\frac{11}{4}\right)^{1/3}\text{ }\text{ }\text{ }\left(tt_f\right).$$
(18)
Thus we have reproduced the well-known difference between photon and neutrino background temperatures on the basis of the temperature law (9). Usually, this result is obtained by calculating the entropy transfer from the $`e^\pm `$ pairs to the photons under the condition of entropy conservation .
## 4 Evaporation of primordial black holes
In a variety of scenarios with copious production of primordial black holes the latter may substantially contribute to the energy density of the universe (see, e.g., ). Some of these models are characterized by a narrow mass spectrum . Under such circumstances it is a good approximation to ascribe the same mass $`m_{_{\left(BH\right)}}`$ to all members of the population. On the other hand, a black hole mass is known to be characterized by a temperature $`T_{_{\left(BH\right)}}m_{_{\left(BH\right)}}^1`$. Consequently, with a single mass population of PBHs one may associate a single temperature $`T_{_{\left(BH\right)}}`$ as well. Furthermore, one may show that a PBH population in the expanding universe may be regarded as an ensemble of non-interacting particles . These properties suggest a description of the PBH component as a pressureless “fluid”, in which $`T_{_{\left(BH\right)}}`$ in some respect plays the role of a fluid temperature. Since PBHs are known to evaporate, it is tempting to establish a two-fluid model along the lines of Sec. 2 with one component being the PBH “fluid”, the second one radiation. The equations of state (7) for the PBH component are
$$p_{_{\left(BH\right)}}=0,\text{ }\text{ }\text{ }\rho _{_{\left(BH\right)}}=n_{_{\left(BH\right)}}m_{_{\left(BH\right)}}.$$
(19)
The black hole temperature is related to its mass by the well-known formula
$$T_{_{\left(BH\right)}}=\frac{1}{8\pi m_{_{\left(BH\right)}}}.$$
(20)
This temperature is attributed to each PBH individually, i.e., primarily it is not a conventional fluid temperature. The number $`N_{_{\left(BH\right)}}`$ of PBHs in a comoving volume $`a^3`$, $`N_{_{\left(BH\right)}}=n_{_{\left(BH\right)}}a^3`$, is not preserved and, according to Eq. (2), we may write down a balance equation for the corresponding PBH number flow vector $`N_{_{\left(BH\right)}}^i=n_{_{\left(BH\right)}}u^i`$,
$$N_{{}_{\left(BH\right)}{}^{};i}^i=\dot{n}_{_{\left(BH\right)}}+3Hn_{_{\left(BH\right)}}=n_{_{\left(BH\right)}}\mathrm{\Gamma }_{_{\left(BH\right)}}\text{ . }$$
(21)
The black hole energy balance becomes \[cf. Eq. (2) with (1)\]
$$\dot{\rho }_{_{\left(BH\right)}}+3H\rho _{_{\left(BH\right)}}=u_at_{_{\left(BH\right)}}^a=\rho _{_{\left(BH\right)}}\left[\mathrm{\Gamma }_{_{\left(BH\right)}}+\frac{\dot{m}_{_{\left(BH\right)}}}{m_{_{\left(BH\right)}}}\right].$$
(22)
Using $`p_{_{\left(BH\right)}}=0`$ as well as Eq. (6) in the fluid temperature law (8) we find
$$\dot{T}_{_{\left(BH\right)}}=\frac{u_at_{_{\left(BH\right)}}^a\mathrm{\Gamma }_{_{\left(BH\right)}}\rho _{_{\left(BH\right)}}}{\rho _{_{\left(BH\right)}}/T_{_{\left(BH\right)}}}=\frac{\rho _{_{\left(BH\right)}}}{\rho _{_{\left(BH\right)}}/T_{_{\left(BH\right)}}}\frac{\dot{m}_{_{\left(BH\right)}}}{m_{_{\left(BH\right)}}}\text{ .}$$
(23)
We emphasize that we have used here the same symbol, $`T_{_{\left(BH\right)}}`$, for the the PBH “fluid” temperature and for the temperature (20), which is ascribed to the individual black holes. The consistency of this identification becomes obvious if we combine the equations of state (19) with (20) and introduce the result for $`\rho _{_{\left(BH\right)}}/T_{_{\left(BH\right)}}`$ into (23):
$$\frac{\rho _{_{\left(BH\right)}}}{T_{_{\left(BH\right)}}}=\frac{\rho _{_{\left(BH\right)}}}{T_{_{\left(BH\right)}}}\dot{T}_{_{\left(BH\right)}}=T_{_{\left(BH\right)}}\frac{\dot{m}_{_{\left(BH\right)}}}{m_{_{\left(BH\right)}}}\text{ .}$$
(24)
It is the crucial point of our analysis that Hawking’s temperature law (20) for individual black holes fits together with the general fluid temperature law (8) for the equations of state (19) with (20). This circumstance provides the basis for our thermodynamical discussion of the PBH evaporation process. To this purpose we identify component 1 of the general analysis in section 2 with the PBH “fluid” and component 2 with ulrarelativistic matter (radiation, subscript r), i.e., $`\left(1\right)\left(BH\right)`$ and $`\left(2\right)\left(r\right)`$. For the latter we require constant entropy per particle, i.e., \[cf. Eq. (6)\]
$$\dot{s}_{_{\left(r\right)}}=0u_at_{_{\left(r\right)}}^a=\frac{4}{3}\rho _{_{\left(r\right)}}\mathrm{\Gamma }_{_{\left(r\right)}}\frac{\dot{T}_{_{\left(r\right)}}}{T_{_{\left(r\right)}}}=H\left(1\frac{\mathrm{\Gamma }_{_{\left(r\right)}}}{3H}\right).$$
(25)
Combination with $`t_{_{\left(BH\right)}}^a=t_{_{\left(r\right)}}^a`$ from (4) yields
$$\mathrm{\Gamma }_{_{\left(r\right)}}=\frac{4}{3}\frac{\rho _{_{\left(BH\right)}}}{\rho _{_{\left(r\right)}}}\left[\mathrm{\Gamma }_{_{\left(BH\right)}}+\frac{\dot{m}_{_{\left(BH\right)}}}{m_{_{\left(BH\right)}}}\right].$$
(26)
The total entropy flow $`S^a`$ is the sum of the contributions $`S_{_{\left(BH\right)}}^a=n_{_{\left(BH\right)}}s_{_{\left(BH\right)}}u^a`$ and $`S_{_{\left(r\right)}}^a=n_{_{\left(r\right)}}s_{_{\left(r\right)}}u^a`$. With $`s_{_{\left(BH\right)}}=4\pi m_{_{\left(BH\right)}}^2`$ we obtain the following expression for the overall entropy production density :
$$S_{;a}^a=\rho _{_{\left(BH\right)}}\mathrm{\Gamma }_{_{\left(BH\right)}}\left[\frac{1}{2T_{_{\left(BH\right)}}}\frac{1}{T_{_{\left(r\right)}}}\right]+\rho _{_{\left(BH\right)}}\frac{\dot{m}_{_{\left(BH\right)}}}{m_{_{\left(BH\right)}}}\left[\frac{1}{T_{_{\left(BH\right)}}}\frac{1}{T_{_{\left(r\right)}}}\right].$$
(27)
It is obvious from (26) that negative values of $`\dot{m}_{_{\left(BH\right)}}`$ (and $`\mathrm{\Gamma }_{_{\left(BH\right)}}`$) correspond to a positive quantity $`\mathrm{\Gamma }_{_{\left(r\right)}}`$. This case is equivalent to the creation of radiative particles at the expense of the PBH mass (and its number), i.e., to PBH evaporation. The inverse process, namely “accretion” with $`\dot{m}_{_{\left(BH\right)}}>0`$ and $`\mathrm{\Gamma }_{_{\left(r\right)}}<0`$ is described by the general formula (27) as well. Which of the two processes is thermodynamically preferred depends on the ratio of the temperatures. Given a specific initial ratio, the further evolution is entirely governed by the temperature laws in (24) and (25). Let’s assume an initial configuration with $`T_{_{\left(BH\right)}}\left(t_0\right)=T_{_{\left(r\right)}}\left(t_0\right)`$. A non-negative entropy production density then requires $`\mathrm{\Gamma }_{_{\left(BH\right)}}0`$. Since one expects $`\dot{m}_{_{\left(BH\right)}}`$ and $`\mathrm{\Gamma }_{_{\left(BH\right)}}`$ to have the same sign, this implies a positive value of $`\mathrm{\Gamma }_{_{\left(r\right)}}`$, i.e., radiation particles are produced which makes the PBH masses shrink. The further evolution depends on a subtle interplay between the rates $`\mathrm{\Gamma }_{_{\left(BH\right)}}`$ and $`\mathrm{\Gamma }_{_{\left(r\right)}}`$ and their respective back reactions on the temperature laws (24) and (25). A positive $`\mathrm{\Gamma }_{_{\left(r\right)}}`$ may either be larger or smaller than the expansion rate $`3H`$. For $`\mathrm{\Gamma }_{_{\left(r\right)}}<3H`$ the fluid temperature decreases according to Eq. (25), while the BH temperature increases according to Eq. (24). It follows that $`T_{_{\left(BH\right)}}>T_{_{\left(r\right)}}`$ at $`t>t_0`$. The evaporation process will continue since $`T_{_{\left(r\right)}}<T_{_{\left(BH\right)}}`$ requires $`\mathrm{\Gamma }_{_{\left(BH\right)}}<0`$ and $`\dot{m}_{_{\left(BH\right)}}/m_{_{\left(BH\right)}}<0`$ to guarantee $`S_{;a}^a>0`$ in Eq. (27). For $`\mathrm{\Gamma }_{_{\left(r\right)}}>3H`$, however, hypothetically realized e.g. by a large initial ratio $`\rho _{_{\left(BH\right)}}/\rho _{_{\left(r\right)}}`$, the fluid temperature increases. If this increase is smaller than the increase in $`T_{_{\left(BH\right)}}`$ we have again $`T_{_{\left(r\right)}}<T_{_{\left(BH\right)}}`$ and the PBH evaporation goes on since it remains thermodynamically favored ($`S_{;a}^a>0`$). But an increase in $`T_{_{\left(r\right)}}`$ stronger than that in $`T_{_{\left(BH\right)}}`$ results in a fluid temperature which is higher than $`T_{_{\left(BH\right)}}`$. For $`T_{_{\left(r\right)}}>2T_{_{\left(BH\right)}}`$ a positive entropy production (27) requires $`\mathrm{\Gamma }_{_{\left(BH\right)}}>0`$ and $`\dot{m}_{_{\left(BH\right)}}/m_{_{\left(BH\right)}}>0`$, implying a quick transition to a negative $`\mathrm{\Gamma }_{_{\left(r\right)}}`$, i.e., the process can no longer continue. A strong “reheating” of the fluid will stop the evaporation and reverse the process. Now, the second law requires PBHs to be formed out of the radiation and to accrete mass. A negative $`\mathrm{\Gamma }_{_{\left(r\right)}}`$, on the other hand, will make $`T_{_{\left(r\right)}}`$ subsequently decrease \[cf. Eq. (25)\]. If $`T_{_{\left(r\right)}}`$ has fallen below $`T_{_{\left(BH\right)}}`$, the evaporation process may set in again. In particular, this self-confining property implies that a catastrophic growth of the PBHs is thermodynamically forbidden. The point is that a PBH growth, i.e. mass accretion with $`\dot{m}_{_{\left(BH\right)}}/m_{_{\left(BH\right)}}>0`$, back reacts on the temperature of the ambient radiation. For a fixed PBH number, i.e. $`\mathrm{\Gamma }_{_{\left(BH\right)}}=0`$, the corresponding radiation temperature changes as
$$\frac{\dot{T}_{_{\left(r\right)}}}{T_{_{\left(r\right)}}}=H\left(1+\frac{1}{4}\frac{\rho _{_{\left(BH\right)}}}{\rho _{_{\left(r\right)}}}\frac{\dot{m}_{_{\left(BH\right)}}}{m_{_{\left(BH\right)}}}H^1\right).$$
It is obvious that for $`\dot{m}_{_{\left(BH\right)}}>0`$ from some time on the temperature $`T_{_{\left(r\right)}}`$ will cool off faster than $`T_{_{\left(BH\right)}}`$ \[cf. Eq. (24)\]. Consequently, $`T_{_{\left(r\right)}}`$ will approach $`T_{_{\left(BH\right)}}`$. As soon as $`T_{_{\left(r\right)}}`$ has fallen below $`T_{_{\left(BH\right)}}`$, mass accretion stops since for $`T_{_{\left(r\right)}}<T_{_{\left(BH\right)}}`$ the rate $`\dot{m}_{_{\left(BH\right)}}/m_{_{\left(BH\right)}}`$ has to be negative in order to guarantee a positive entropy production, i.e., the process now proceeds in the reverse direction and the PBHs can no longer grow but start to evaporate again. This completes our thermodynamic discussion of PBH evaporation based on the temperature law (8) (and its special case (9)).
## 5 “Deflationary” universe
In the two previous examples we did not consider the impact of the decay and production processes on the expansion behavior of the universe. As was shown in and , the general tendency of this influence is to increase the cosmic expansion rate. Namely, processes of the type discussed in sections 3 and 4 give rise to an effective viscous pressure of the cosmic medium as a whole \[cf. Eq. (4)\]. Since this contribution to the overall pressure is negative, its net effect is to accelerate the expansion. While this effect is small for the cases dealt with in sections 3 and 4, it is essential in the “deflationary” universe model of the present section. This model relies on Einstein’s field equations with the energy-momentum tensor (4) of a bulk viscous fluid. In a homogeneous and isotropic universe one has
$$\kappa \rho =3H^2,\text{ }\dot{H}=\frac{\kappa }{2}\left(\rho +p+\mathrm{\Pi }\right)\kappa \mathrm{\Pi }=3\gamma H^22\dot{H},$$
(28)
where $`\kappa `$ is Einstein’s gravitational constant and $`\gamma 1+p/\rho `$. In case $`\mathrm{\Pi }`$ is not a “conventional” viscous pressure but represents a quantity describing cosmological particle production on a phenomenological level , it may be related to the production rate $`\mathrm{\Gamma }`$ introduced in (5). For “adiabatic” particle production this relation is
$$\mathrm{\Pi }=\left(\rho +p\right)\frac{\mathrm{\Gamma }}{3H}.$$
(29)
Combination with the field equations (28) then yields
$$\frac{\mathrm{\Gamma }}{3H}=1+\frac{2}{3\gamma }\frac{\dot{H}}{H^2}\frac{H^{}}{H\left[\frac{\mathrm{\Gamma }}{3H}1\right]}=\frac{3}{2}\frac{\gamma }{a},$$
(30)
where $`H^{}\text{d}H/\text{d}a`$. Strictly speaking, the rate $`\mathrm{\Gamma }`$ has to be calculated on the quantum level (see, e.g., ). In a phenomenological setting an ansatz for $`\mathrm{\Gamma }/H`$ is required. For a dependence $`\mathrm{\Gamma }\rho H^2`$ and $`\gamma =4/3`$ we obtain
$$H=2\frac{a_e^2}{a^2+a_e^2}H_e\frac{\mathrm{\Gamma }}{3H}=\frac{a_e^2}{a^2+a_e^2},$$
(31)
where we have chosen the constants such that $`\dot{H}_e=H_e^2`$, i.e, $`\ddot{a}>0`$ for $`a<a_e`$ and $`\ddot{a}<0`$ for $`a>a_e`$. $`H`$ starts with a constant value $`H_0=2H_e`$ at $`aa_e`$ and then “deflates” towards the typical $`Ha^2`$ behaviour of a radiation dominated universe for $`aa_e`$. This Hubble rate has originally been obtained in the context of phenomenological approaches to cosmological vacuum decay . Again, this is a two-component model with one component playing the role of the cosmological “vacuum”. Our point here is to demonstrate that such kind of model fits into the general structure of section 2 and admits a similar thermodynamic discussion as the cases of electron-positron annihilation and PBH evaporation. We will identify the first component of the general formalism in Sec. 2 with the “vacuum” (subscript $`v`$), i.e., $`\left(1\right)\left(v\right)`$, the second one again with radiation, i.e., $`\left(2\right)\left(r\right)`$. The sketched scenario may then be obtained on the basis of an interacting two-fluid model with $`\rho =\rho _{_{\left(v\right)}}+\rho _{_{\left(r\right)}}`$ where
$$\rho _{_{\left(v\right)}}=\frac{3H_e^2}{2\pi }m_P^2\left[\frac{a_e^2}{a^2+a_e^2}\right]^3,\text{ }\text{ }\text{ }\rho _{_{\left(r\right)}}=\frac{3H_e^2}{2\pi }m_P^2\left(\frac{a}{a_e}\right)^2\left[\frac{a_e^2}{a^2+a_e^2}\right]^3,$$
(32)
and $`m_P^2=8\pi /\kappa `$ is the square of the Planck mass. The part $`\rho _{_{\left(v\right)}}`$ is finite for $`a0`$ and decays as $`a^6`$ for $`aa_e`$, while the part $`\rho _{_{\left(r\right)}}`$ describes relativistic matter with $`\rho _{_{\left(r\right)}}0`$ for $`a0`$ and $`\rho _{_{\left(r\right)}}a^4`$ for $`aa_e`$. The energy balances are ($`A=v,r`$)
$$\dot{\rho }_{_{\left(A\right)}}+3H\left[\rho _{_{\left(A\right)}}+p_{_{\left(A\right)}}\right]=\mathrm{\Gamma }_{_{\left(A\right)}}\left[\rho _{_{\left(A\right)}}+p_{_{\left(A\right)}}\right]$$
(33)
with
$$\frac{\mathrm{\Gamma }_{_{\left(v\right)}}}{3H}=\left(1\frac{1}{2}\frac{a^2}{a_e^2}\right)\frac{a_e^2}{a^2+a_e^2},\text{ }\text{ }\text{ }\frac{\mathrm{\Gamma }_{_{\left(r\right)}}}{3H}=\frac{3}{2}\frac{a_e^2}{a^2+a_e^2}.$$
(34)
The equation for $`\rho _{_{\left(v\right)}}`$ may be written as
$$\dot{\rho }_{_{\left(v\right)}}+3H\left(\rho _{_{\left(v\right)}}+p_{_{\left(v\right)}}+\mathrm{\Pi }_{_{\left(v\right)}}\right)=0,$$
(35)
where
$$\mathrm{\Pi }_{_{\left(v\right)}}\frac{\mathrm{\Gamma }_{_{\left(v\right)}}}{3H}\left(\rho _{_{\left(v\right)}}+p_{_{\left(v\right)}}\right)=\left(1\frac{1}{2}\frac{a^2}{a_e^2}\right)\frac{a_e^2}{a^2+a_e^2}\left(\rho _{_{\left(v\right)}}+p_{_{\left(v\right)}}\right).$$
(36)
This corresponds to an effective equation of state
$$P_{_{\left(v\right)}}p_{_{\left(v\right)}}+\mathrm{\Pi }_{_{\left(v\right)}}=\frac{a^2a_e^2}{a^2+a_e^2}\rho _{_{\left(v\right)}}.$$
(37)
Although we have always $`p_{_{\left(v\right)}}=\rho _{_{\left(v\right)}}/3`$, the effective equation of state for $`a0`$ approaches $`P_{_{\left(v\right)}}=\rho _{_{\left(v\right)}}`$. Effectively, this component behaves as a vacuum contribution. For $`aa_e`$ it represents stiff matter with $`P_{_{\left(v\right)}}=\rho _{_{\left(v\right)}}`$. The radiation component may be regarded as emerging from the decay of the initial vacuum according to
$$\dot{\rho }_{_{\left(r\right)}}+4H\rho _{_{\left(r\right)}}=\dot{\rho }_{_{\left(v\right)}}.$$
(38)
The radiation temperature is obtained from the general law (9), which in the present case specifies to
$$\frac{\dot{T}_{_{\left(r\right)}}}{T_{_{\left(r\right)}}}=H\left(1\frac{\mathrm{\Gamma }_{_{\left(r\right)}}}{3H}\right)=H\left[1\frac{3}{2}\frac{a_e^2}{a^2+a_e^2}\right].$$
(39)
Integration yields
$$T_{_{\left(r\right)}}=2^{3/4}T_{_{(e,r)}}\frac{a_e}{a}\left[\frac{a^2}{a^2+a_e^2}\right]^{3/4}T_{_{\left(r\right)}}\rho _{_{\left(r\right)}}^{1/4}.$$
(40)
$`T_{_{(e,r)}}`$ is the value of the radiation temperature at $`a=a_e`$. This temperature starts at $`T_{_{\left(r\right)}}=0`$ for $`a=0`$, then increases to a maximum value, given by $`\mathrm{\Gamma }_{_{\left(r\right)}}=3H`$, equivalent to $`a^2=\frac{1}{2}a_e^2`$,
$$T_{_{\left(r\right)}}^{max}=\left(\frac{32}{27}\right)^{1/4}T_{_{(e,r)}},$$
(41)
and finally decreases as $`a^1`$ for large values of $`a`$. Our formalism allows us to ascribe a temperature $`T_{_{\left(v\right)}}`$ to the “vacuum” as well, which is determined analogously by
$$\frac{\dot{T}_{_{\left(v\right)}}}{T_{_{\left(v\right)}}}=H\left(1\frac{\mathrm{\Gamma }_{_{\left(v\right)}}}{3H}\right)=\frac{3}{2}H\frac{a^2}{a^2+a_e^2}.$$
(42)
The “vacuum” temperature behaves as
$$T_{_{\left(v\right)}}=T__0\left[\frac{a_e^2}{a^2+a_e^2}\right]^{3/4}.$$
(43)
It starts from a maximum value at $`a=0`$ and decreases as $`a^{3/2}`$ for large $`a`$. The “vacuum” is radiative in the sense that $`\rho _{_{\left(v\right)}}T_{_{\left(v\right)}}^4`$ is valid. As a final remark we mention that it is also possible to introduce a temperature $`T`$ of the cosmic medium as a whole with a behavior
$$\frac{\dot{T}}{T}=H\left(1\frac{\mathrm{\Gamma }}{3H}\right)=H\frac{a^2}{a^2+a_e^2}T=T__0\left[\frac{a_e^2}{a^2+a_e^2}\right]^{1/2},$$
(44)
which “interpolates” between (43) for small $`a`$ and (40) for $`aa_e`$. These considerations clarify the central role played by the general temperature law (9) also under conditions where the relevant back reaction substantially affects the entire cosmological dynamics.
## 6 Conclusions
Cosmological thermodynamics allows us to establish a unifying view on a broad range of different phenomena and to uncover joint underlying structures. In this paper we have explored similar thermodynamic features of matter creation in the early universe, primordial black hole evaporation, and electron-positron annihilation after neutrino decoupling. All these processes are governed by the same basic temperature law for a fluid with variable particle number which takes into account the back reaction of the relevant interactions on the thermal history of the universe. A particular aspect of our considerations is the consistency of this law with Hawking’s black hole temperature formula $`T_{_{\left(BH\right)}}m_{_{\left(BH\right)}}^1`$. This circumstance provides the basis for a two-fluid model for the evaporation of a single-mass PBH component into radiation.
Acknowledgment
This paper was supported by the Deutsche Forschungsgemeinschaft, the Spanish Ministry of Education (grant PB94-0718) and NATO (grant CRG 940598).
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# 1 Introduction
## 1 Introduction
Discussions have recently arisen about the possibility that expectations from OPE for some types of semi-leptonic rates may be violated by terms of order $`1/m_Q`$. The argument of Nathan Isgur is founded on general considerations ; namely the duality is obtained in the infinite mass limit through cancellation between the falloff of the ground state contribution and the rise of the excitations (the Bjorken sum rule indeed relates the derivatives of these contributions with respect to $`w`$, near $`w=1`$), but at finite mass there is some mismatch near zero recoil, which could be of order $`1/m_Q`$. Indeed, in terms of $`t`$, the quadri-dimensional transfer<sup>3</sup><sup>3</sup>3$`t`$ is we use the old standard notation $`t`$, to avoid confusion with the tridimensional $`|\stackrel{}{q}|^2`$, which will be used extensively in this non relativistic (NR) context. :
$`t=(q^0)^2\stackrel{}{q}{}_{}{}^{2},`$ (1)
the respective $`t_{max}`$ do not coincide anymore. The argument is then given by the author further likeliness by some calculations within a very simple “toy” model : the non relativistic harmonic oscillator (HO) potential model.
In the present letter, we will not discuss directly the issue about QCD (see our article ). We simply stick to the very model used in , and show that within this model, calculating the total integrated rate $`\mathrm{\Gamma }_{inclusive}`$ by summation on the relevant final (bound) states, duality with free quark decay rate is in fact satisfied, in the SV (Shifman-Voloshin) limit<sup>4</sup><sup>4</sup>4By SV limit, we do not mean simply that the recoil velocity is small, but also, as in the original paper , that $`\frac{\delta m}{m}`$ is small ; in addition, $`\delta m`$ is taken large with respect to light quark parameters ; in non relativistic quantum mechanics, we assume : $`\mathrm{\Delta }=\frac{1}{m_dR^2}\delta mandm_d,\delta mm_b,m_c`$, $`\mathrm{\Delta }`$ being the level spacing.; this means that the difference $`\mathrm{\Gamma }_{inclusive}\mathrm{\Gamma }_{freequark}`$ comes out as expected, which implies in particular (as discussed below) cancellation of terms of relative order $`(\delta m)^2/m_b^2`$ and $`\delta m/m_b^2`$ (by relative, we mean with respect to the free quark decay rate ; note that such terms correspond to $`(1/m_Q)^0,(1/m_Q)^1`$ in the usual $`1/m_Q`$ expansion). Our argument is for integrated decay rates, so we do not claim anything on possible effects in differential or partially integrated rates. Also, of course, we cannot exclude by such argument that something odd may happen in QCD.
One very interesting point raised in the discussion of is about the very specific cancellations which are necessary for duality to hold, and about the contributions of the various regions of phase space. We try to analyze through our demonstration how such cancellations occur in subleading order for total widths. An interesting consequence of the analysis is that to find the required cancellations, one needs not to consider only the sum rule of Bjorken ; one has to take into account in addition the Voloshin sum rule (the fact that one needs the sum rules has been suggested by the Minnesota group in their discussion with N. Isgur , but is made here quite explicit ; for related discussions in QCD by the same group, see ). In fact, the Voloshin sum rule is exactly what is needed for cancellation of terms of relative order $`\delta m/m_b^2`$ in the difference $`\mathrm{\Gamma }_{inclusive}\mathrm{\Gamma }_{freequark}`$. The sum rules are trivially satisfied in the HO model, but it is not so trivial in general. Our conclusion is not in contradiction with the mismatch occuring near zero recoil, considered in , because the latter is very small parametrically with respect to the terms we consider in the total width.
Note that the use of SV limit is not essential to demonstrate duality in this way, and neither is the use of an HO potential. Their choice is pedagogical. Indeed we have also done the demonstration for an arbitrary potential () and also for fixed $`m_c/m_b`$ ratio. Nevertheless, the particular case considered here is of pedagogical interest, because on the one hand the discussion in the SV limit is much simpler, and the similar discussion in QCD can hardly be made beyond the SV limit, and because on the other hand, within HO model, we can give explicit expressions. Moreover, we are able to give a complete proof that in the HO model $`1/m_Q`$ terms are absent in the ratio $`\mathrm{\Gamma }_{inclusive}/\mathrm{\Gamma }_{freequark}`$ beyond the SV limit (article to appear ). Note also that the demonstration is independent of the leptonic tensor, as we have also shown elsewhere, but we choose here one specific for illustration. On the other hand, the coefficient of the terms of order $`\frac{1}{R^2m_b^2}`$, which we also evaluate, is model-dependent (in particular it depends on the choice of the leptonic tensor; we choose here one for illustration).
## 2 Model
$``$ The model for hadrons is the non relativistic harmonic oscillator quark model (the motion of quarks both internal and due to overall hadron are both treated non relativistically), describing the initial (quarks $`b`$ and $`d`$) and final ($`c`$ and $`d`$) hadrons. The potential is assumed to be flavor independent, which is crucial for the demonstration. The great advantage of the harmonic oscillator, which appears in the summation on final states, is that very few states contributes to the transition rates in the limited expansion in $`1overm_b`$ which we perform (see next section). Energy levels, for a state labelled by $`(n_x,n_y,n_z)`$, $`n=n_x+n_y+n_z`$, write :
$`E_n=m_{b,c}+m_d+\left({\displaystyle \frac{3}{2}}+n\right){\displaystyle \frac{1}{\mu _{b,c}R_{b,c}^2}}`$ (2)
where $`\mu _{b,c}`$ are the reduced masses $`\frac{m_{b,c}m_d}{m_{b,c}+m_d}`$ and the radii $`R_{b,c}^2`$ can be written as :
$`R_{b,c}^2=\sqrt{{\displaystyle \frac{m_d}{\mu _{b,c}}}}R_{\mathrm{}}^2`$ (3)
$`R_{\mathrm{}}`$ being the radius in the infinite mass limit. We will often denote the first level excitation energy in the infinite mass limit as :
$`\mathrm{\Delta }={\displaystyle \frac{1}{m_dR_{\mathrm{}}^2}}.`$ (4)
For simplicity, from now on, we denote :
$`R_{\mathrm{}}=R`$ (5)
$``$ Quarks are then coupled to lepton pairs : $`bc\mathrm{}\nu `$, through a quark vector current $`j^0=1`$, $`\stackrel{}{j}=0`$ <sup>5</sup><sup>5</sup>5Note that we do not claim to make a systematic non relativistic expansion of a relativistic theory, but only to consider a non relativistic Hamiltonian for the bound states; we can choose freely the weak interaction current. The essential point is then to treat consistently the matrix elements according to the chosen interactions, in the specified SV expansion. (or equivalently we can speak of spinless quarks), and a leptonic tensor, which will be described by functions denoted generically through letter $`L`$ and some arguments and indices. $`\stackrel{}{P}`$ and $`\stackrel{}{P^{}}`$ are the initial and final hadron momenta ; the total energies of hadrons are $`P^0=E+\stackrel{}{P}^2/2(m_b+m_d),P^0=E^{}+\stackrel{}{P^{}}^2/2(m_c+m_d)`$, with $`E,E^{}`$ the energies at rest ; but, in practice, we will always work in the initial hadron rest frame : $`\stackrel{}{P}=0`$ ; $`\stackrel{}{P}`$ and $`\stackrel{}{P^{}}`$ are the initial and final quark momenta. We denote $`\stackrel{}{q}=\stackrel{}{P}\stackrel{}{P^{}}=\stackrel{}{P^{}}`$. The basic equation is then, in the initial state rest frame :
$`\mathrm{\Gamma }_{inclusive}=K{\displaystyle \underset{n}{}}{\displaystyle _0^{|\stackrel{}{q}|_{max,n}}}d|\stackrel{}{q}||\stackrel{}{q}|^2L_n(|\stackrel{}{q}|){\displaystyle \underset{n=n_x+n_y+n_z}{}}|j_{0(n_x,n_y,n_z)}|^2.`$ (6)
The constant $`K`$ depends only on the decay interaction strength. The constant $`K`$ will be ommitted in the rest of the letter. $`_{n=n_x+n_y+n_z}|j_{0(n_x,n_y,n_z)}|^2`$ only depends on $`|\stackrel{}{q}|`$. The angular integration has been performed. The notations $`L_n(|\stackrel{}{q}|)`$ and $`|\stackrel{}{q}|_{max,n}`$ are explained now. A priori, after angular integration, the leptonic tensor appears through a function of energy loss $`q^0`$ and $`\stackrel{}{q}^2`$, $`L(q^0,|\stackrel{}{q}|^2)`$. However, for the decay from the ground state to a h.o. state labelled by $`(n_x,n_y,n_z)`$, by energy conservation, $`q^0=P^0P^0`$ is just a function of $`|\stackrel{}{q}|`$ and $`(n_x,n_y,n_z)`$. Moreover, the energy loss $`q^0`$ will depend only on $`n=n_x+n_y+n_z`$. and we then denote as $`L_n(|\stackrel{}{q}|)`$ the result of $`L(q^0,|\stackrel{}{q}|^2)`$, when the energy loss $`q^0`$ is assumed to be calculated for the corresponding $`n`$, as a function of $`|\stackrel{}{q}|`$. Indeed, for a state with degree of excitation $`n`$ :
$`q^0(n,|\stackrel{}{q}|)=m_bm_c+{\displaystyle \frac{3}{2}}({\displaystyle \frac{1}{\mu _bR_b^2}}{\displaystyle \frac{1}{\mu _cR_c^2}})n{\displaystyle \frac{1}{\mu _cR_c^2}}{\displaystyle \frac{|\stackrel{}{q}|^2}{2(m_c+m_d)}}.`$ (7)
Now $`q_{max}`$ is determined by the equation $`t=(q^0)^2|\stackrel{}{q}|^2=0`$, $`q^0(|\stackrel{}{q}|)=|\stackrel{}{q}|`$ :
$`|\stackrel{}{q}|_{max,n}={\displaystyle \frac{2(m_c+m_d)(\delta E)_n}{2(m_c+m_d)+\sqrt{(m_c+m_d)^2+2(m_c+m_d)(\delta E)_n}}}`$ (8)
where
$`(\delta E)_n=q^0(n,\stackrel{}{q}=0)=m_bm_c+{\displaystyle \frac{3}{2}}({\displaystyle \frac{1}{\mu _bR_b^2}}{\displaystyle \frac{1}{\mu _cR_c^2}})n{\displaystyle \frac{1}{\mu _cR_c^2}}.`$ (9)
$`|\stackrel{}{q}|_{max}`$ just depends on $`n`$. $`L(q^0,|\stackrel{}{q}|^2)`$ can be taken as an arbitrary function without spoiling any of the general statements made below, but for definiteness we will henceforth choose :
$`L(q^0,\stackrel{}{q}{}_{}{}^{2})=3(q^0)^2|\stackrel{}{q}|^2,`$ (10)
inspired by a static quark approximation of the V-A current. The corresponding free quark decay rate is :
$`\mathrm{\Gamma }_{free}=K{\displaystyle _0^{|\stackrel{}{q}|_{max,free}}}d|\stackrel{}{q}||\stackrel{}{q}|^2L(q^0,|\stackrel{}{q}|^2)`$ (11)
with :
$`q^0(free,|\stackrel{}{q}|)=m_bm_c{\displaystyle \frac{|\stackrel{}{q}|^2}{2m_c}}`$ (12)
$`|\stackrel{}{q}|_{max,free}={\displaystyle \frac{2m_c\delta m}{(m_c+\sqrt{m_c^2+2m_c\delta m})}}`$ (13)
with $`\delta m=m_bm_c`$.
## 3 SV expansion and demonstration of duality
$``$ We have then to consider the expansion of
$`ϵ={\displaystyle \frac{\mathrm{\Gamma }_{inclusive}\mathrm{\Gamma }_{free}}{\mathrm{\Gamma }_{free}}}`$ (14)
in powers of $`\frac{1}{m_b}`$, and the aim is in principle to show that it begins with order $`\frac{1}{m_b^2}`$ only, as expected from a formal OPE (the NR version of OPE will be explained in the more developped article). More precisely this holds in the limit $`m_b\mathrm{}`$ with $`r=\frac{m_c}{m_b}`$ fixed, for which we reserve for clarity the term ” usual $`1/m_Q`$ expansion”. However, we will work in the SV (Shifman-Voloshin) limit, which corresponds to making in addition an expansion in $`1r`$. Namely, with :
$`\delta m=m_bm_c,`$ (15)
we write $`m_c=m_b\delta m`$ and we expand in powers of $`\frac{1}{m_b}`$, keeping $`\delta m`$ fixed, as well as the light quark parameters, $`m_d`$, $`1/R`$ ; then, we make a second limited expansion, taking $`\mathrm{\Delta }=\frac{1}{m_dR^2}`$ small with respect to $`\delta m`$. The terms have the form $`\frac{(\delta m)^k^{}}{(m_b)^k}`$ times light quark factors. But then the aim must be more than just showing the absence of powers <sup>6</sup><sup>6</sup>6Note that, in the present letter, we term generically as power $`\frac{1}{(m_b)^k}`$ all the terms which contain the factor $`\frac{1}{(m_b)^k}`$, whatever the powers of $`\delta m`$ and light quantities. $`\frac{1}{(m_b)^k}`$, $`k<2`$ in $`ϵ`$.
Indeed, if it is true, this would not in principle preclude terms of the type $`\frac{(\delta m)^k^{}}{m_b^2}`$ ($`k^{}>0`$) in $`ϵ`$. Such terms would be large in practice, since $`\delta m`$ is not so small. And in fact, they would correspond, in terms of the usual $`1/m_Q`$ expansion, to terms of order $`(1/m_Q)^0,(1/m_Q)`$, since $`\delta m`$ would be then $`m_Q`$. Such terms are not expected from OPE. We must therefore show that such terms do not exist in the final result, and we show it in fact. More precisely, we show that potentially large terms of the type $`\frac{(\delta m)^2}{m_b^2}`$, $`\frac{m_d\delta m}{m_b^2}`$, which appear in particular contributions, do finally cancel out, leaving us with terms of the type $`\frac{1}{R^2m_b^2}`$ (terms with $`k^{}>2`$ simply do not appear in the way we calculate $`ϵ`$, neither do terms with power $`\frac{1}{(m_b)^0}`$ or $`\frac{1}{m_b}`$ \- in fact, the delicate part consists in showing the cancellation of $`\frac{m_d\delta m}{m_b^2}`$ terms). This is all that is required by duality with free quarks, as concerns the terms with power $`(\frac{1}{m_b})^k`$, $`k2`$. We will calculate the terms of type $`\frac{1}{R^2m_b^2}`$, which do not vanish in general. Note that such terms are small with respect to $`\frac{m_d\delta m}{m_b^2}`$ by a factor $`\frac{\mathrm{\Delta }}{\delta m}`$. In the usual $`1/m_Q`$ expansion they correspond to order $`1/m_Q^2`$. On the other hand, we will not calculate in the expansion of $`ϵ`$ smaller terms having also the power $`\frac{1}{m_b^2}`$, but which contain still additional powers of $`\frac{\mathrm{\Delta }}{\delta m}`$ with respect to $`\frac{1}{R^2m_b^2}`$, corresponding in $`\mathrm{\Gamma }_{inclusive}\mathrm{\Gamma }_{free}`$ to terms like $`\frac{(\delta m)^4\times \mathrm{\Delta }}{m_b^2}`$, $`\frac{(\delta m)^3\times \mathrm{\Delta }^2}{m_b^2}`$, etc… and retain only the terms proportional to $`\mathrm{\Gamma }_{free}(\delta m)^5`$. The neglected terms correspond to terms of relative order $`1/m_Q^3`$ or beyond in the $`1/m_Q`$ expansion. For sake of simplicity, we will neither examine further checks of duality in terms of the type $`\frac{(\delta m)^k^{}}{(m_b)^k}`$, with $`k>2`$.
In any case, we see that we do have to calculate terms with a power $`\frac{1}{m_b^2}`$ and not $`\frac{1}{m_b}`$ only, since the terms with a power $`\frac{1}{m_b^2}`$ may correspond to terms of the order $`(1/m_Q)^0,(1/m_Q)^1`$ in the usual expansion. The method precisely consists in writing the difference $`\mathrm{\Gamma }_{inclusive}\mathrm{\Gamma }_{free}`$ as a sum of terms which contain already a power $`\frac{1}{m_b^2}`$, and then to demonstrate the above additional cancellations.
$``$ The advantage of harmonic oscillator (HO) model is that the level $`n=1`$ (which corresponds to $`L=1`$ states) appears only with a power $`\frac{1}{m_b^2}`$, and that higher levels come only with a power $`\frac{1}{m_b^3}`$ at least. Since we keep terms with a power $`\frac{1}{m_b^i}`$,$`i2`$ , we only need consider $`n=0,1`$ states. For sake of simplicity, we denote their respective contributions $`\mathrm{\Gamma }_{0,1}`$.
We have at this order, by expanding the matrix elements :
$`\mathrm{\Gamma }_0{\displaystyle _0^{|\stackrel{}{q}|_{max,0}}}d|\stackrel{}{q}||\stackrel{}{q}|^2L_{n=0}(|\stackrel{}{q}|)(1\rho ^2{\displaystyle \frac{|\stackrel{}{q}|^2}{m_b^2}}),`$ (16)
where $`\rho ^2=\frac{1}{2}m_d^2R^2`$ is the standard slope of the ground state form factor with respect to $`w`$ ($`w1\frac{1}{2}\frac{|\stackrel{}{q}|^2}{m_b^2}`$) ; note that effect of non complete overlapping between hadrons with $`b`$ and $`c`$ quarks is completely negligible here, since it contributes at order $`1/R^2\frac{(\delta m)^2}{m_b^4}`$. For $`L=1`$ states :
$`\mathrm{\Gamma }_1{\displaystyle _0^{|\stackrel{}{q}|_{max,1}}}d|\stackrel{}{q}||\stackrel{}{q}|^2L_{n=1}(|\stackrel{}{q}|)\tau ^2{\displaystyle \frac{|\stackrel{}{q}|^2}{m_b^2}},`$ (17)
with $`\tau =\frac{m_dR}{\sqrt{2}}`$ corresponding to the $`\tau _{1/2,3/2}(w=1)\times \sqrt{3}`$. The other excitations do not contribute at this order, because the matrix element $`<n|\stackrel{}{r}|0>`$ is non zero only if $`n=1`$. From the explicit expressions, we have the relations :
$`\rho ^2\tau ^2=0,`$ (18)
$`\mathrm{\Delta }\tau ^2={\displaystyle \frac{m_d}{2}},`$ (19)
($`\mathrm{\Delta }`$ being the level spacing, 4) as non relativistic analogues of the Bjorken and Voloshin sum rules. These sum rules could then be used to generalise the present analysis.
In fact, we will try as much as possible not to specify separately $`\mathrm{\Delta },\rho ,\tau `$, but to use only the above sum rules and expressions for $`\mathrm{\Gamma }_{0,1}`$.
$``$ The strategy is to note that the difference between $`\mathrm{\Gamma }_0+\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_{free}`$ can be reexpressed by successive steps :
1) Decomposition into the same difference with $`L_{0,1}(|\stackrel{}{q}|)`$ substituted by their free counterpart $`L_{free}(|\stackrel{}{q}|)`$ (contribution (I)) plus a $`\frac{1}{m_b^2}`$ term (contribution (II)). 2) Then the first contribution (I) is rewritten trivially as a difference between two contributions having a power $`\frac{1}{m_b^2}`$ relative to the free quark decay integrand, further shown to be of relative order $`\frac{1}{R^2m_b^2}`$ or smaller. 3) It is also shown that in contribution (II), which contains manifestly a power $`\frac{1}{m_b^2}`$, there are only terms of the type $`\frac{1}{R^2m_b^2}`$ or smaller.
$``$ In a first step, using the respective expressions given above for the $`q^0(|\stackrel{}{q}|)`$’s, and expanding it to the required order, we find :
$`q^0(free,|\stackrel{}{q}|)\delta m{\displaystyle \frac{|\stackrel{}{q}|^2}{2m_c}}`$ (20)
$`q^0(n=0,|\stackrel{}{q}|)\delta m(1{\displaystyle \frac{3}{4m_b^2R^2}}){\displaystyle \frac{|\stackrel{}{q}|^2}{2(m_c+m_d)}}`$ (21)
$`q^0(n=1,|\stackrel{}{q}|)\delta m\mathrm{\Delta }`$ (22)
Note that in our expansion, the main term in the three quantities is $`\delta m`$. The main term of $`q_{max}`$ is then also $`\delta m`$ ($`q_{max}=q^0`$ at $`t=0`$). We use these expansions to make, in the integrals for $`\mathrm{\Gamma }_{0,1}`$ :
$`L_0(|\stackrel{}{q}|)L_{free}(|\stackrel{}{q}|)+6\delta m({\displaystyle \frac{3\delta m}{4R^2m_b^2}}+{\displaystyle \frac{m_d|\stackrel{}{q}|^2}{2m_b^2}})`$ (23)
$`L_1(|\stackrel{}{q}|)L_{free}(|\stackrel{}{q}|)+6\delta m(\mathrm{\Delta })+3\mathrm{\Delta }^2`$ (24)
The second terms in the r.h.s. come from the difference between the respective $`q^0`$, as a function of $`|\stackrel{}{q}|`$. Note that in the expansion of $`L_1(|\stackrel{}{q}|)`$, one can neglect terms in $`\frac{1}{m_b^2}`$ because $`\mathrm{\Gamma }_1`$ has already a power $`\frac{1}{m_b^2}`$. We get $`\mathrm{\Gamma }_{inclusive}\mathrm{\Gamma }_{free}\delta \mathrm{\Gamma }_I+\delta \mathrm{\Gamma }_{II}`$ with :
$`\delta \mathrm{\Gamma }_I=`$ $`{\displaystyle _0^{|\stackrel{}{q}|_{max,0}}}`$ $`d|\stackrel{}{q}||\stackrel{}{q}|^2L_{free}(|\stackrel{}{q}|)(1\rho ^2{\displaystyle \frac{|\stackrel{}{q}|^2}{m_b^2}})+`$ (25)
$`{\displaystyle _0^{|\stackrel{}{q}|_{max,1}}}`$ $`d|\stackrel{}{q}||\stackrel{}{q}|^2L_{free}(|\stackrel{}{q}|)\tau ^2{\displaystyle \frac{|\stackrel{}{q}|^2}{m_b^2}}`$ (26)
$`{\displaystyle _0^{|\stackrel{}{q}|_{max,free}}}`$ $`d|\stackrel{}{q}||\stackrel{}{q}|^2L_{free}(|\stackrel{}{q}|)`$ (27)
and
$`\delta \mathrm{\Gamma }_{II}=`$ $`{\displaystyle _0^{|\stackrel{}{q}|_{max,0}}}`$ $`d|\stackrel{}{q}||\stackrel{}{q}|^26\delta m({\displaystyle \frac{3\delta m}{4R^2m_b^2}}+{\displaystyle \frac{m_d|\stackrel{}{q}|^2}{2m_b^2}})+`$ (28)
$`{\displaystyle _0^{|\stackrel{}{q}|_{max,1}}}`$ $`d|\stackrel{}{q}||\stackrel{}{q}|^2(6\delta m(\mathrm{\Delta })+3\mathrm{\Delta }^2)(\tau ^2{\displaystyle \frac{|\stackrel{}{q}|^2}{m_b^2}}).`$ (29)
$``$ Contribution I. One can write it as the difference of two integrals which have already manifestly a factor $`\frac{1}{m_b^2}`$, i.e. the terms with power $`\frac{1}{(m_b)^0}`$ or $`\frac{1}{m_b}`$ are already cancelled (this amounts to using $`\rho ^2\tau ^2=0`$, which is the particular form of the Bjorken sum rule in the model):
$`\delta \mathrm{\Gamma }_I=`$ $`{\displaystyle _{|\stackrel{}{q}|_{max,free}}^{|\stackrel{}{q}|_{max,0}}}`$ $`d|\stackrel{}{q}||\stackrel{}{q}|^2L_{free}(|\stackrel{}{q}|)`$ (30)
$`{\displaystyle _{q_{max,1}}^{|\stackrel{}{q}|_{max,0}}}`$ $`d|\stackrel{}{q}||\stackrel{}{q}|^2L_{free}(|\stackrel{}{q}|)\tau ^2{\displaystyle \frac{|\stackrel{}{q}|^2}{m_b^2}}`$ (31)
We first expand each integral. \- One has :
$`|\stackrel{}{q}|_{max,0}|\stackrel{}{q}|_{max,free}\delta m({\displaystyle \frac{1}{2}}{\displaystyle \frac{m_d\delta m}{m_b^2}}{\displaystyle \frac{3}{4}}{\displaystyle \frac{1}{R^2m_b^2}}),`$ (32)
whence
$`{\displaystyle _{|\stackrel{}{q}|_{max,free}}^{|\stackrel{}{q}|_{max,0}}}d|\stackrel{}{q}||\stackrel{}{q}|^2L_{free}(|\stackrel{}{q}|)\delta m({\displaystyle \frac{1}{2}}{\displaystyle \frac{m_d\delta m}{m_b^2}}{\displaystyle \frac{3}{4}}{\displaystyle \frac{1}{R^2m_b^2}})(\delta m)^2L_{free}(|\stackrel{}{q}|_{max}=\delta m).`$ (33)
One can make $`|\stackrel{}{q}|\delta m`$ in the integrand, because the integration interval contains already a power $`1/m_b^2`$, and the difference between $`|\stackrel{}{q}|`$ and $`\delta m`$ contains a further $`1/m_b`$ factor.
\- The second integral is more delicate, because the integration interval has not a factor $`1/m_b^2`$; it is just $`\mathrm{\Delta }`$ ; the variation of $`|\stackrel{}{q}|`$ is not negligible. One must do a limited expansion of the integrand in powers of $`\frac{\mathrm{\Delta }}{\delta }`$, so as to retain at least terms of the type $`\frac{1}{R^2m_b^2}`$. It is there that the second expansion,in powers of $`\frac{\mathrm{\Delta }}{\delta }`$, enters the game :
$`{\displaystyle _{|\stackrel{}{q}|_{max,1}}^{|\stackrel{}{q}|_{max,0}}}`$ $`d|\stackrel{}{q}||\stackrel{}{q}|^2L_{free}(|\stackrel{}{q}|)\tau ^2{\displaystyle \frac{|\stackrel{}{q}|^2}{m_b^2}}`$ (34)
$`{\displaystyle _{\delta m\mathrm{\Delta }}^{\delta m}}`$ $`d|\stackrel{}{q}||\stackrel{}{q}|^2L_{free}(|\stackrel{}{q}|)\tau ^2{\displaystyle \frac{|\stackrel{}{q}|^2}{m_b^2}}`$ (35)
$`\mathrm{\Delta }(\delta m)^2L_{free}(|\stackrel{}{q}|=\delta m)\tau ^2{\displaystyle \frac{(\delta m)^2}{m_b^2}}`$
$`{\displaystyle \frac{\mathrm{\Delta }^2}{2}}{\displaystyle \frac{\tau ^2}{m_b^2}}{\displaystyle \frac{d}{d|\stackrel{}{q}|}}(|\stackrel{}{q}|^4L_{free}(|\stackrel{}{q}|))(|\stackrel{}{q}|=\delta m).`$
Let us note that to estimate the relative order of the different terms, one has to divide by a reference rate, which will be taken to be the free quark decay rate ; now $`(\delta m)^3L_{free}(|\stackrel{}{q}|=\delta m)`$, as well as $`\frac{d}{d|\stackrel{}{q}|}(|\stackrel{}{q}|^4L_{free}(|\stackrel{}{q}|))(|\stackrel{}{q}|=\delta m)`$, are of the order of the free quark decay rate (with our choice $`L_{free}(|\stackrel{}{q}|=\delta m)(\delta m)^2`$).
Then, one can first observe that in fact not only all the terms written in eq. (33), (35) have a relative power $`1/m_Q^2`$, but that they are more precisely of relative order $`m_d\delta m/m_b^2`$ at most ; terms of relative order $`(\delta m)^2/m_b^2`$ are already cancelled. This will be obtained more generally thanks to Bjorken sum rule. Now, the term of relative order $`m_d\delta m/m_b^2`$ encountered in the r.h.s. of the first integral (33) is cancelled by the first term in the r.h.s. of the second integral (35), just using Voloshin sum rule (19), i.e. $`\mathrm{\Delta }\tau ^2=m_d/2`$. All the remaining contributions are of the type $`(\delta m)^5\frac{1}{R^2m_b^2}`$. We can evaluate them readily and find them to cancel too for the particular choice made for $`L(q^0,|\stackrel{}{q}|)`$. Finally :
$`\delta \mathrm{\Gamma }_I=`$ $`{\displaystyle _{|\stackrel{}{q}|_{max,free}}^{|\stackrel{}{q}|_{max,0}}}`$ $`d|\stackrel{}{q}||\stackrel{}{q}|^2L_{free}(|\stackrel{}{q}|)`$ (36)
$`{\displaystyle _{|\stackrel{}{q}|_{max,1}}^{|\stackrel{}{q}|_{max,0}}}`$ $`d|\stackrel{}{q}||\stackrel{}{q}|^2L_{free}(|\stackrel{}{q}|)\rho ^2{\displaystyle \frac{|\stackrel{}{q}|^2}{m_b^2}}0`$ (37)
It must be emphasized that the cancellation can occur because the difference between $`|\stackrel{}{q}|_{max,n}`$ and $`|\stackrel{}{q}|_{max,free}`$ is changing sign between the ground state and the excitations. With our assumption $`\mathrm{\Delta }\delta m`$, one has $`|\stackrel{}{q}|_{max,1}<|\stackrel{}{q}|_{max,free}<|\stackrel{}{q}|_{max,0}`$
$``$ Contribution II. It is also obvious that it contains already a power $`\frac{1}{m_b^2}`$. On factorising $`(\delta m)^5`$, one sees that $`\frac{m_d\delta m}{m_b^2}`$ terms are present in the first integral (second term of the bracket in the integrand) : $`_0^{|\stackrel{}{q}|_{max,0}}dqq^2(6\delta m\frac{m_dq^2}{2m_b^2}`$)) and in the second one (first term of the bracket in the integrand) : $`_0^{|\stackrel{}{q}|_{max,0}}dqq^2(6\delta m(\mathrm{\Delta })(\tau ^2\frac{|\stackrel{}{q}|^2}{m_b^2})`$)), the rest being smaller. It is easily seen that these $`\frac{m_d\delta m}{m_b^2}`$ terms cancel at this order, just using Voloshin sum rule $`\mathrm{\Delta }\tau ^2=m_d/2`$, to leave a smaller contribution, which is only of order $`\frac{1}{R^2m_b^2}`$ ; the latter is found by performing a limited expansion of the integrand as above eq. (33) (the interval is once more $`𝒪(\mathrm{\Delta })`$), in powers of $`\frac{\mathrm{\Delta }}{\delta m}`$ :
$`3{\displaystyle \frac{m_d\delta m}{m_b^2}}{\displaystyle _{|\stackrel{}{q}|_{max,1}}^{|\stackrel{}{q}|_{max,0}}}d|\stackrel{}{q}||\stackrel{}{q}|^43(\delta m)^5{\displaystyle \frac{1}{R^2m_b^2}}`$ (38)
The other terms in the integrals are already manifestly of this order, and one ends with :
$`\delta \mathrm{\Gamma }_{II}={\displaystyle \frac{9}{5}}(\delta m)^5{\displaystyle \frac{1}{R^2m_b^2}}.`$ (39)
This result has been checked by a systematic expansion using Mathematica.
Finally, with $`\mathrm{\Gamma }_{free}\frac{4}{5}(\delta m)^5`$:
$`ϵ={\displaystyle \frac{\mathrm{\Gamma }_0+\mathrm{\Gamma }_1\mathrm{\Gamma }_{free}}{\mathrm{\Gamma }_{free}}}{\displaystyle \frac{\frac{9}{5}(\delta m)^5\frac{1}{R^2m_b^2}}{\frac{4}{5}(\delta m)^5}}={\displaystyle \frac{9}{4}}{\displaystyle \frac{1}{R^2m_b^2}}`$ (40)
Let us reinsist that it is of the order expected from OPE, unlike terms of the type $`\frac{(\delta m)^2}{m_b^2}`$ or $`\frac{m_d\delta m}{m_b^2}`$, which duely cancel, as has been shown.
## 4 Relative magnitude of Isgur contribution
$``$ Let us now return briefly to the very discussion raised by ref. . One could be worried why it is found there some duality violating effect, while we do not. The contradiction is only apparent. The answer seems to be that in totally integrated widths, the effect considered in is finally relatively small parametrically with respect to the ones we have considered. Let us show that. The mismatch near zero recoil considered in is the integral of the ground state contribution over $`w_0(t)=\frac{m_B^2+m_D^2t}{2m_Bm_D}`$ between $`w_0(t=(m_Bm_D)^2)=1`$ and the threshold for the excited state production $`w_0(t=(m_Bm_D^{})^2)`$ ( the variable $`w`$ for the ground state contribution is considered as a function of $`t`$, $`w_0(t)`$). Let us pass through the variable $`\stackrel{}{q}`$, which is more adapted to the NR problem, and denote as $`|\stackrel{}{q}|_n(t)`$ the value of $`|\stackrel{}{q}|`$ which corresponds to some $`t`$ for a state $`n`$; the total ground state contribution can be decomposed into two parts :
$`\mathrm{\Gamma }_0{\displaystyle _{|\stackrel{}{q}|_0(t=(m_Bm_D)^2)}^{|\stackrel{}{q}|_0(t=0)}}d|\stackrel{}{q}||\stackrel{}{q}|^2L_{n=0}(|\stackrel{}{q}|)(1\rho ^2{\displaystyle \frac{|\stackrel{}{q}|^2}{m_b^2}})=`$
$`{\displaystyle _{|\stackrel{}{q}|_0(t=(m_Bm_D)^2)}^{|\stackrel{}{q}|_0(t=0)}}d|\stackrel{}{q}||\stackrel{}{q}|^2L_{n=0}(|\stackrel{}{q}|){\displaystyle _{|\stackrel{}{q}|_0(t=(m_Bm_D)^2)}^{|\stackrel{}{q}|_0(t=(0)}}d|\stackrel{}{q}||\stackrel{}{q}|^2L_{n=0}(|\stackrel{}{q}|)(\rho ^2{\displaystyle \frac{|\stackrel{}{q}|^2}{m_b^2}})`$
. (41)
In the infinite mass limit, $`m_Bm_D^{}m_Bm_Dm_bm_c`$ and the functions $`|\stackrel{}{q}|_0(t)`$ and $`|\stackrel{}{q}|_1(t)`$, as well as $`|\stackrel{}{q}|_{free}(t)`$, become identical, and the functions $`L(q^0,\stackrel{}{q}{}_{}{}^{2})`$ become also identical for all states. Then, the first contribution equates the free quark decay rate, while the second one :
$`\delta \mathrm{\Gamma }_0{\displaystyle _{|\stackrel{}{q}|_0(t=(m_Bm_D)^2)}^{|\stackrel{}{q}|_0(t=0)}}d|\stackrel{}{q}||\stackrel{}{q}|^2L_{n=0}(|\stackrel{}{q}|)(\rho ^2{\displaystyle \frac{|\stackrel{}{q}|^2}{m_b^2}}),`$ (42)
is exactly cancelled by the excited state contribution :
$`\mathrm{\Gamma }_1{\displaystyle _{|\stackrel{}{q}|_1(t=(m_Bm_D^{})^2)}^{|\stackrel{}{q}|_1(t=0)}}d|\stackrel{}{q}||\stackrel{}{q}|^2L_{n=1}(|\stackrel{}{q}|)\tau ^2{\displaystyle \frac{|\stackrel{}{q}|^2}{m_b^2}},`$ (43)
due to Bjorken sum rule. Whence duality. However, when quark masses are finite, there is a small part of the integral (42) which is uncancelled, in spite of the Bjorken sum rule, by the corresponding excited state contribution, in particular because $`t=(m_Bm_D^{})^2`$ now differs from $`t=(m_Bm_D)^2`$. We estimate the mismatch as :
$`\delta \mathrm{\Gamma }{\displaystyle _{|\stackrel{}{q}|_0(t=(m_Bm_D)^2)}^{|\stackrel{}{q}|_0(t=(m_Bm_D^{})^2)}}d|\stackrel{}{q}||\stackrel{}{q}|^2L_{n=0}(|\stackrel{}{q}|)(\rho ^2{\displaystyle \frac{|\stackrel{}{q}|^2}{m_b^2}}).`$ (44)
In this calculation, following , we disregard all other sources of difference, in particular the fact that the leptonic tensor functions are no more equal, and neither are the functions $`|\stackrel{}{q}|_n(t)`$ for $`n=0`$ and $`n=1`$ respectively, and that also the first contribution in 41 no longer equates the free quark decay rate. Then, our point is that this mismatch of total widths is very small with respect to the terms we have retained. Indeed, the integral runs over a small part of the phase space, but in addition the integrand is much smaller near zero recoil, where the mismatch takes place, first because of the leptonic factor $`L(q^0,\stackrel{}{q}{}_{}{}^{2})=3(q^0)^2|\stackrel{}{q}|^2`$, second because of the factor $`(\rho ^2\frac{|\stackrel{}{q}|^2}{m_b^2})`$. Since $`L(q^0,\stackrel{}{q}{}_{}{}^{2})=3(q^0)^2|\stackrel{}{q}|^2`$, using $`|\stackrel{}{q}|_0(t=(m_Bm_D^{})^2)\sqrt{\frac{2\mathrm{\Delta }}{\delta m}}|\stackrel{}{q}|_{0,max}`$ :
$`\delta \mathrm{\Gamma }{\displaystyle \frac{\rho ^2}{m_b^2}}3(\delta m)^2{\displaystyle \frac{(\frac{2\mathrm{\Delta }}{\delta m})^{\frac{5}{2}}}{5}}|\stackrel{}{q}|_{0,max}^5,`$ (45)
and, relative to the free quark decay rate (i.e. contribution to $`ϵ`$) :
$`{\displaystyle \frac{\delta \mathrm{\Gamma }}{\mathrm{\Gamma }_{free}}}{\displaystyle \frac{\rho ^2}{m_b^2}}{\displaystyle \frac{3}{4}}(\delta m)^2({\displaystyle \frac{2\mathrm{\Delta }}{\delta m}})^{\frac{5}{2}},`$ (46)
which is parametrically small, because of the factor $`(\frac{2\mathrm{\Delta }}{\delta m})^{\frac{5}{2}}`$ (since $`\mathrm{\Delta }\delta m`$ in the SV limit). In fact, in our calculation we have not retained such terms.
Numerically too, we find it very small, with real physical masses. It is true, as noticed in ref. , that numerically the region of Dalitz plot which is concerned is physically not very small, because one is far from the SV limit; with our approximative formula, we find around $`20\%`$ of the free decay rate in this region of phase space, not far from the $`30\%`$ estimated in ref. ; but the factors considered above nevertheless combine to yield a very small effect for $`\frac{\delta \mathrm{\Gamma }}{\mathrm{\Gamma }_{free}}`$, around $`10^3\rho ^2`$. This is due to the fact that the factor $`(\rho ^2\frac{|\stackrel{}{q}|^2}{m_b^2})`$ is very small in this region of phase space.
## 5 Conclusion
Stimulated by the worries raised by N. Isgur, we have noticed mismatches between the sum of exclusive decays and the free quark total decay rate, which, considered separately, could convey the impression that quark hadron duality between total widths is violated at order $`\delta m/m_b^2`$, because all these mismatches are of this order. Let us recapitulate them :
1) The upper limit in terms of $`|\stackrel{}{q}|`$ (corresponding to $`t=0`$) of the integrals for the ground state and the excited states contributions do not coincide. Therefore, the contributions from the falloff of ground state and rise of excited states do not cancel near $`|\stackrel{}{q}|_{max}`$ ($`t=0`$).
2) The upper limit in $`|\stackrel{}{q}|`$ of the integrals for the ground state contribution and the free quark decay do not coincide for similar reasons.
3) The leptonic tensors of the various contributions are different, because the function $`q^0(|\stackrel{}{q}|)`$ depends on the transition considered.
At order $`𝒪(\frac{\delta m}{m_b^2})`$, 1) and 2) cancel between each other, while 3) has a zero net effect, by internal cancellation of the differences of leptonic tensors, when integrated (taking into account the difference in upper limits of integration in $`|\stackrel{}{q}|`$, near maximum recoil, is once more necessary).
It must be emphasized that even in this simple model and in the SV limit, it is by no means trivial to check duality, because the check requires to take into account detailed effects, such as the dependence of ground state binding energy on the heavy quark masses through their different radii, which itself reflects the flavor independence of the quark potential, etc…
In both cases, the cancellation occurs because of Voloshin sum rule. The consideration of it is absolutely necessary, in addition to Bjorken one, to demonstrate duality of total widths through summation of exclusive states at subleading order. Note that, if we have an independent mean to demonstrate duality, for example by a rigourous demonstration of OPE to the required order, we can use the result on the sum of exclusive states, on the reverse, to demonstrate these sum rules.
Acknowledgements
A.L.Y., L.O., O.P. and J.-C. R. acknowledge partial support from the EEC-TMR Program, contract N. CT 98-0169.
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# Generalized Entropy approach to far-from-equilibrium statistical mechanics.
## Abstract
We present a new approach to far–from–equilibrium statistical mechanics, based on the concept of generalized entropy, which is a microscopically-defined generalization of Onsager-Machlup functional. In the case when a set of slow (adiabatic) variables can be chosen, our formalism yields a general form of the macroscopic evolution law (Generalized Langevin Equation) and extends Fluctuation Dissipative Theorem. It also provides for a simple understanding of recently–discovered Fluctuation Theorem.
PACS numbers: 05.70.Ln, 05.20.-y
Foundations of non-equilibrium statistical physics remain in focus of intensive research for nearly a century. Over past several decades, there has been a dramatic progress in application of stochastic equations to wide variety of complex systems, as well as in understanding some of their generic features. However, the methods of microscopic derivation of such coarse-grained descriptions (e.g. starting with a Hamiltonian), remain essentially limited to traditional kinetic theory and linear response scheme . In this work, we emphasize the reductionistic mission of non-equilibrium statistical mechanics, by building a constructive formalism whose conceptual framework resembles that of equilibrium theory. Our approach is not limited to the vicinity of thermal equilibrium, and becomes equivalent to the classical linear response theory in this limit.
The central concept in our discussion is Generalized Entropy (GE), which goes back to works by Onsager and Machlup . Further development of this paradigm includes its generalization for far-from-equilibrium case by Graham and action-principle approach to Marcovian stochastic dynamics by Eyink , whose technique and conclusions have many common points with ours. The distinct feature of the present approach is an explicit microscopic interpretation of GE, which enables us to derive the macroscopic evolution equations starting from a microscopic level. Normally, entropy is assigned to a macroscopic state of a system as the logarithm of its statistical weight. Similarly, we introduce GE, as a logarithm of the statistical weight of a given macroscopic evolution (trajectory). We limit ourselves to the class of microscopic dynamics (parameterized with microvariables $`q_j`$, $`j=1\mathrm{}N`$) which are deterministic and phase-volume-preserving (as e.g. in Hamiltonian systems). In this case, there is a natural measure in the space $`q_j`$, i.e. the statistical weight of any subset of the space is its phase volume. Let $`𝐀=\widehat{𝐀}(q)`$ be a macroscopic (possibly, multicomponent) variable. The formal definition of the GE associated with a given evolution $`𝐀(t)`$ between two moments of time, $`t_i`$ and $`t_f`$, is the following:
$$S\left[𝐀(t)\right]_{t_i}^{t_f}\mathrm{log}\left[D[q_j]\underset{t=t_i}{\overset{t_f}{}}\delta \left(\widehat{𝐀}(q(t))𝐀(t)\right)\right].$$
(1)
Here, integration is performed over all actual microscopic trajectories $`q_j(t)`$, i.e. those satisfying the microscopic equations of motion. Since we assume the deterministic and phase–volume–preserving evolution on that level, all the actual trajectories are equally weighted.
Our next step is to distinguish between two contributions to GE: $`S\left[𝐀(t)\right]_{t_i}^{t_f}=S_{\mathrm{\Delta }t}^{(0)}(\overline{𝐀})+S^{(kin)}\left[\dot{𝐀}(t)\right]_{t_i}^{t_f}`$. The first term, $`S^{(0)}`$ is the logarithm of the total number of trajectories of length $`\mathrm{\Delta }tt_ft_i`$ with a fixed ”midpoint” $`\overline{𝐀}\left(𝐀_{t_i}+𝐀_{t_f}\right)/2`$. The other contribution, $`S^{(kin)}`$, to which we will refer as kinetic entropy, is the logarithm of the probability of a given evolution $`\dot{𝐀}(t)`$ for the fixed $`t_i`$, $`t_f`$ and $`\overline{𝐀}`$:
$`S^{(kin)}\left[\dot{𝐀}(t)\right]_{t_i}^{t_f}\mathrm{log}{\displaystyle \underset{t=t_i}{\overset{t_f}{}}}\delta ({\displaystyle \frac{d}{dt}}\widehat{𝐀}(q(t))\dot{𝐀}(t))_{\overline{𝐀}}=`$ (2)
$`\mathrm{log}{\displaystyle D[𝐗(t)]\mathrm{exp}\left(\mathrm{\Sigma }[𝐗(t)]_{t_i}^{t_f}𝐗(t)\dot{𝐀}(t)𝑑t\right)}.`$ (3)
Here we have transformed the expression by using the exponential representation of $`\delta `$-function, which resulted in introduction of a new variable $`𝐗`$ conjugated to $`\dot{𝐀}`$. $`\mathrm{\Sigma }`$ is the generating functional for the variable $`\dot{𝐀}`$, which in particularly means that its variation of any order in $`𝐗`$ coincides with the corresponding irreducible correlator of the conjugated field (in our case, $`\dot{𝐀}`$) :
$`\mathrm{\Sigma }[𝐗(t)]\mathrm{log}\mathrm{exp}({\displaystyle _{t_i}^{t_f}}𝐗(t){\displaystyle \frac{d}{dt}}\widehat{𝐀}(q_j(t))dt)_{\overline{𝐀}}=`$ (4)
$`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}{\displaystyle 𝑑t_1\mathrm{}𝑑t_n\dot{𝐀}(t_1)\mathrm{}\dot{𝐀}(t_n)_{\overline{𝐀}}𝐗(t_1)\mathrm{}𝐗(t_n)}.`$ (5)
Here $``$ is the irreducible part of the correlator, i.e. the n-point average with subtracted contribution reducible to the lower–order correlation functions. Here and in the future $`\dot{𝐀}(t_1)\dot{𝐀}(t_2)`$ denotes direct tensor product, which is to be distinguished from scalar one, e.g. $`𝐗(𝐭)\dot{𝐀}(t)`$. The averaging is performed over all trajectories with the given initial and final times ($`t_i`$, $`t_f`$) and fixed midpoint, $`\overline{𝐀}`$.
Now, after we have related kinetic entropy $`S^{kin}`$, to the statistics of $`\dot{𝐀}`$, it becomes possible to express another contribution to GE, $`S^{(0)}(\overline{𝐀})`$, in terms of regular thermodynamic entropy $`S(𝐀)`$. For doing so we notice that integration of the statistical weights of all trajectories originating from a given point $`𝐀`$ is, by definition, the weight of the initial state, $`\mathrm{exp}S(𝐀)`$. After making a simple calculation based on this observation, one gets $`S(𝐀)=S^{(0)}(𝐀)+\mathrm{\Sigma }\left[\frac{1}{2}\delta S^{(0)}/\delta 𝐀\right]`$. The time interval $`\mathrm{\Delta }t`$, and the associated change of $`𝐀`$ is assumed to be sufficiently small so that the linear expansion of $`S^{(0)}`$ in $`\mathrm{\Delta }𝐀`$ be valid.
In a general case, $`\mathrm{\Sigma }`$-functional is an awkward mathematical object, because of its non–local structure. However, until this point we have not restricted the choice of the macroparameters, $`𝐀`$. From the practical point of view, it is clear that the coarse–grained description of a system may be reasonable if one can choose a set of relatively slow variables as the macroparameters. Below we specify this choice in a more formal way.
By definition, the form of $`\mathrm{\Sigma }`$–functional (and the correlators of $`\dot{𝐀}`$) depends on the midpoint position, $`\overline{𝐀}`$. However, there naturally exists an interval $`\delta \overline{𝐀}`$ within which such dependence can be neglected. This allows one to introduce a concept of drift time, $`\tau _{drift}`$, over which most of the trajectories remain within this interval of constant statistics of $`\dot{𝐀}`$. Suppose there exists a shorter time scale, $`\tau _0\tau _{drift}`$ such that any correlator $`\dot{𝐀}(t_1)\mathrm{}\dot{𝐀}(t_n)`$ becomes negligible when $`\left|t_1t_n\right|>\tau _0`$. In this case, one can choose the initial and final times such that $`\tau _0\mathrm{\Delta }tt_ft_i\tau _{drift}`$. This considerably simplifies the expression for $`\mathrm{\Sigma }`$–functional: if we are only interested in the behavior of the system on times larger than $`\tau _0`$, $`\mathrm{\Sigma }`$ becomes local in $`𝐗`$:
$$\mathrm{\Sigma }[𝐗(t)]=_{t_i}^{t_f}\mathrm{\Xi }\left(𝐗(t)\right)𝑑t,$$
(6)
here
$$\mathrm{\Xi }\left(𝐗\right)=\underset{n=1}{\overset{\mathrm{}}{}}\frac{𝐗^n}{n!}𝑑t_1\mathrm{}𝑑t_{n1}\dot{𝐀}(0)\mathrm{}\dot{𝐀}(t_{n1}).$$
(7)
We will refer to time scale $`\tau _0`$ as ergodicity time. It can be shown that from the methodological point of view, the assumption of the existence of such time scale does play a role similar to that of the ergodic principle in equilibrium theory.
Collecting the above results gives the following expression for GE:
$$S[𝐀(t)]_{t_i}^{t_f}=S\left(𝐀(t_i)\right)+\mathrm{log}D[𝐗(t)]\mathrm{exp}S^{}[𝐀(t),𝐗(t)],$$
(8)
The first contribution here is the conventional entropy of the initial state, and the other one is the logarithm of the probability of a given evolution starting at that point. The latter is expressed in terms of the functional $`S^{}`$, which has a meaning of GE in $`(𝐀,𝐗)`$ space:
$$𝒮^{}[𝐀(t),𝐗(t)]=_{t_i}^{t_f}𝑑t\left[\dot{𝐀}\left(\frac{1}{2}\frac{\delta S}{\delta 𝐀}𝐗\right)+\mathrm{\Xi }\left(𝐗\right)\mathrm{\Xi }\left(\frac{1}{2}\frac{\delta S}{\delta 𝐀}\right)\right].$$
(9)
This quantity becomes additive on the time scales exceeding $`\tau _0`$, which means that the dynamics becomes Markovian. This fact allows us to extend the applicability of the above expression to the case when $`t_ft_i>\tau _{drift}`$, i.e. when $`\mathrm{\Xi }(𝐗)`$ is no longer independent of $`𝐀`$.
In a particular case of reversible microscopic dynamics, $`\mathrm{\Xi }`$ does not change if the sign of $`𝐗`$ reversed. This implies that the ratio of probabilities of direct and reversed evolutions along the same path $`𝐀(t)`$ is independent of the form of $`\mathrm{\Sigma }`$-functional and is given by $`\mathrm{exp}(𝐀(t_f)𝐀(t_i)`$ (see (8)-(9)). This property is known as Fluctuation Theorem , which has been recently established for a wide variety of non-equilibrium systems.
One can eliminate the ”fictitious” variable $`𝐗`$ from Eqs. (8)-(9):
$`S[𝐀(t)]_{t_i}^{t_f}={\displaystyle \frac{S\left(𝐀(t_i)\right)+S\left(𝐀(t_f)\right)}{2}}+`$ (10)
$`{\displaystyle _{t_i}^{t_f}}𝑑t\left[\mathrm{\Lambda }\left(\dot{𝐀}(t)\right)\mathrm{\Xi }\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{\delta S}{\delta 𝐀}}\right)\right],`$ (11)
here the kinetic entropy, which is now a local functional, has been expressed as a time integral of pseudo-Lagrangian $`\mathrm{\Lambda }`$:
$`S^{(kin)}\left[\dot{𝐀}(t)\right]={\displaystyle _{t_i}^{t_f}}\mathrm{\Lambda }(\dot{𝐀}(t))𝑑t=`$ (12)
$`\mathrm{log}{\displaystyle D[𝐗(t)]\mathrm{exp}_{t_i}^{t_f}𝑑t\left[\mathrm{\Xi }(𝐗(t))𝐗(t)\dot{𝐀}(t)\right]}.`$ (13)
Because of the obvious analogy with classical mechanics, we will refer the above $`(𝐀,𝐗)`$ and $`(𝐀,\dot{𝐀})`$ forms for GE as pseudo–Hamiltonian and pseudo–Lagrangian ones, respectively. Though they are completely equivalent, the pseudo-Hamiltonian formalism requires introduction of additional variables $`𝐗`$ (which plays the role of momentum conjugated to measurable $`𝐀`$), while the pseudo–Lagrangian form obscures the relationship between the conjugated functions $`\mathrm{\Lambda }`$ and $`\mathrm{\Xi }`$, given by Eq.(12).
In the case of a distributed system, when both $`𝐀`$ and $`𝐗`$ are fields, $`\mathrm{\Lambda }`$ and $`\mathrm{\Xi }`$ would typically become local functionals. If there is a global conservation law for one or several components of $`𝐀`$, the $`\mathrm{\Sigma }`$-functional, Eq. (4) is invariant with respect to the global transformation $`𝐗(\stackrel{}{r})𝐗(\stackrel{}{r})+_\alpha \delta ^{(\alpha )}𝐧^{(\alpha )}`$ (here index $`\alpha `$ counts all the conserved components of field $`𝐀`$, $`𝐧^{(\alpha )}𝐀`$). Existence of this transformation means that $`\mathrm{\Sigma }`$ and $`\mathrm{\Xi }`$ should depend only on gradients (and higher spatial derivatives) of the corresponding components of the field $`𝐗`$: $`\mathrm{\Sigma }=\xi (X^{(\alpha )})𝑑t𝑑\stackrel{}{r}`$. If the microscopic fluxes $`\stackrel{}{j}^{(\alpha )}`$ of the conserved parameters can be introduced, the expansion of $`\xi `$ in powers of $`X^{(\alpha )}`$ is given by the form similar to Eq. (7), with all the correlators of $`\dot{𝐀}`$ replaced with those of microfluxes $`\stackrel{}{j}^\alpha `$.
We now discuss the dynamics of the system in the deterministic limit, which correseponds to the settle point of the GE functional. One has to emphasize that although the structure of the functional is similar to regular action, the result of the variation procedure is dramatically different from that in mechanics. In our case, only initial point $`𝐀(t_i)`$ should be kept fixed, while the initial rate $`\dot{𝐀}(t_i)`$ (or, equivalently, the final point $`𝐀(t_f)`$), is subject to optimization. As a result, the current value of $`𝐀`$, rather than the pair $`(𝐀,\dot{𝐀})`$ determine the future dynamics of the system, and the equation of motion is of the first, rather than the second order. Namely, maximization of the pseudo-Lagrangian form of GE, Eq. (10), with respect to $`𝐀(t_f)`$, yields the following dynamic law:
$$\frac{\delta \mathrm{\Lambda }(\dot{𝐀})}{\delta \dot{𝐀}}=\frac{1}{2}\frac{\delta S}{\delta 𝐀}.$$
(14)
This equation can be interpreted as a balance between thermodynamic driving force (the right hand side) and the dissipative force (the apparent physical meaning of the left hand side). An alternative description can be obtained by variation of the functional in Eq.(9) with respect to $`\dot{𝐀}`$ and $`𝐗`$:
$`\dot{𝐀}={\displaystyle \frac{\delta \mathrm{\Xi }}{\delta 𝐗}}|_{𝐗=\frac{1}{2}\delta S/\delta 𝐀}.`$ (15)
The equation shows how the system moves in response to the thermodynamic driving force, which in the deterministic limit appears to be identical to variable $`𝐗`$.
In the case of conserved components of $`𝐀`$-field, Eq. (15), will be replaced with regular continuity equation, $`\dot{A}^{(\alpha )}=\stackrel{}{J}^{(\alpha )}`$, in which the macroscopic flux is given by constitutive equation, $`\stackrel{}{J}^{(\alpha )}=\delta \xi /\delta (X^\alpha )`$. It is interesting to note that the macroscopic flux can be introduced even if there is no well-defined fluxes on microscopic scale: this concept follows from the spatial locality of $`\mathrm{\Sigma }`$-functional and conservation of quantity $`A^{(\alpha )}𝑑\stackrel{}{r}`$.
“Kinetic potentials” $`\mathrm{\Lambda }`$ and $`\mathrm{\Xi }`$ in the deterministic limit are related through Legander transform, i.e. $`\mathrm{\Xi }(𝐗)=\mathrm{\Lambda }(\dot{𝐀})+𝐗\dot{𝐀}`$, where $`𝐗=\delta \mathrm{\Lambda }(\dot{𝐀})/\delta \dot{𝐀}`$. In the vicinity of equilibrium the driving forces are small, and therefore only leading terms in expansion (7) remain relevant: $`\mathrm{\Xi }=\dot{𝐀}𝐗+\frac{1}{2}\mathrm{\Gamma }^{(2)}𝐗^2`$. After substituting this expression for $`\mathrm{\Xi }`$ into equation of motion (15), we recover the classical linear response result: $`\dot{𝐀}=\dot{𝐀}+\frac{1}{2}\widehat{\mathrm{\Gamma }}^{(2)}\delta S/\delta 𝐀`$, i.e. the dissipative contribution to $`\dot{𝐀}`$ is proportional to the thermodynamic driving force, and our result for the corresponding kinetic coefficient (see (7)) coincides with the one given by Fluctuation–Dissipative Theorem (FDT). It is a general practice to assume the same linear rate-force relationship even in the regime in which entropy (free energy) is no longer a harmonic function of deviations from equilibrium. Langevin equation is one of the most noticeable examples of such an approach. Furthermore, Öttinger et al have recently proposed an elegant unified way of representing most of the known stochastic models in a single generic form, which again assumes a linear rate-force relationship for dissipative dynamics. Our Eq. (15) (or (14)) in a general case would result in a nonlinear relationship between them and, in this sense, can be referred to as Generalized Langevin Equation (GLE) (an additional noise term will be discussed below). The fact that a particular form of the evolution equation depends on the correlators of macroparameters, suggests a possibility for a synergy between our scheme and earlier field–theoretical approaches to non-equilibrium statistical mechanics .
In order to demonstrate how GLE works outside the linear regime, we consider a trivial kinetic problem: an ensemble of independent two–state systems, each of which has the same transition rate $`\kappa `$ in either direction. The relaxation dynamics for the population difference between the two states, $`N_{}N_1N_2`$ is given by equation $`\dot{N}_{}=2\kappa N_{}`$. Although linear, it is not a linear response result. The thermodynamic driving force conjugated to $`N_{}`$ is the chemical potential difference between the two states, i.e. in our notations $`2X_{}=S/N_{}=\mathrm{log}N_2/N_1`$. This expression can be linearized in $`N_{}`$ (for constant $`N_+N_1+N_2`$) only sufficiently close to equilibrium, i.e. when $`N_{}N_+`$. This means that the simple linear kinetic equation is a result of a non-linear dependence of the response $`\dot{N}_{}`$ on the driving force $`X_{}`$. In the considered case, one can calculate $`\mathrm{\Xi }(X_{})`$ exactly. In the limit of large $`N_+`$, the number of switches happening over small time $`\delta t`$ is $`N_+\kappa \delta t`$. Since all the switches are completely uncorrelated, the original formula for $`\mathrm{\Sigma }`$–potential, Eq. (4), results in the following expression for $`\mathrm{\Xi }(X)`$: $`\mathrm{\Xi }=N_+\kappa \mathrm{log}[\mathrm{cosh}(2X_{})]`$. Here we have taken into account the fact that the change of population difference, $`\delta N_{}`$ is either $`2`$ or $`2`$ for any individual switch and the both possibilities are equally weighted. Indeed, any microscopic trajectory (sequence of individual switches) can be reversed, and this does not change the position of its midpoint, $`\overline{N_{}}(N_{}(t_i)+N_{}(t_f))/2`$. After using Eq.(15), we recover the expected linear equation for $`N_{}`$.
We now proceed with the discussion of fluctuations around the deterministic dynamics. Let $`𝐀^{(\mathrm{𝟎})}(t)`$ be a solution to equation of motion (14), and $`𝐚(t)𝐀(t)𝐀^{(\mathrm{𝟎})}(t)`$ is the deviation of an actual trajectory from it ($`𝐚(t_i)=0`$). The corresponding deviation of the generalized entropy from its local maximum is given by the following quadratic expression:
$`\delta S\left[𝐚(t)\right]={\displaystyle _{t_i}^{t_f}}{\displaystyle \frac{\widehat{\lambda }^1}{2}}\left(\dot{𝐚}+\widehat{\lambda }\widehat{\kappa }𝐚\right)^2𝑑t=`$ (16)
$`{\displaystyle \frac{1}{2}}{\displaystyle 𝐚_\omega \left(\omega ^2\widehat{\lambda }^1+\widehat{\kappa }\widehat{\lambda }\widehat{\kappa }\right)𝐚_\omega \frac{d\omega }{2\pi }},`$ (17)
here $`\widehat{\lambda }\delta ^2\mathrm{\Xi }/\delta 𝐗^2`$ and $`\widehat{\kappa }\frac{1}{2}\delta ^2S/\delta 𝐀^2`$. Note that the deterministic trajectory is stable only if the response matrix $`\widehat{\lambda }\widehat{\kappa }`$ is positive-definite. Otherwise, any two trajectories originating from the same point diverge exponentially fast. By definition, $`\mathrm{exp}\left(\delta S\left[𝐚(t)\right]\right)`$ is the statistical weight of a given trajectory. Therefore, the above quadratic functional, Eq. (16), corresponds to Gaussian statistics of $`𝐚`$, with $`𝐚_\omega 𝐚_\omega =\left(\omega ^2\widehat{\lambda }^1+\widehat{\kappa }\widehat{\lambda }\widehat{\kappa }\right)^1`$. Equivalently, this result can be represented by introduction of a random Gaussian noise $`\eta (t)`$ to the second equation of (15), with $`\eta (t)\eta (t^{})=\widehat{\lambda }\delta (tt^{})`$.
The fact that the same matrix $`\delta ^2\mathrm{\Xi }/\delta 𝐗^2`$ controls both the strength of the fluctuations and the response to a small variation in driving force (see (16) or (15)), is a signature of FDT , which is conventionally applied only in linear-response regime. In order to extend FDT to our far-from-equilibrium case, we probe the system with time-dependent perturbation introduced as an addition to entropy: $`S(𝐀,t)=S_0(𝐀)+𝐡(t)𝐀`$. In a particular case of a Hamiltonian system with fast coupling to a thermal bath, $`h`$ is proportional to field $`𝐡^{}`$ canonically-conjugated to $`𝐀`$: $`𝐡=\beta 𝐡^{}`$. The perturbation results in adding a coupling term $`_\omega i\omega 𝐡_\omega 𝐚_\omega /2`$ to functional $`\delta S`$, Eq.(16). The response now can be determined by variation of $`\delta S`$ with respect to $`𝐚_\omega `$. An important aspect of this procedure is that the causality principle should be taken into account: in $`t`$-representation, $`𝐚(t)`$ should be varied with the specified past and unknown future (as we did while deriving equations of motion, (14) and (15)). In $`\omega `$ representation, the response is given by
$$\widehat{\mathrm{\Gamma }}(\omega )\frac{\delta 𝐚_\omega }{\delta 𝐡_\omega }=\frac{i\omega }{2}𝐚_\omega 𝐚_\omega _{}.$$
(18)
Here index $``$ stands for the retarded part of the correlator, i.e. the one with all poles at $`\mathrm{}(\omega )>0`$ half-plane. The above relationship extends the classical FDT towards strongly non-equilibrium regime. It has to be stressed that $`\widehat{\mathrm{\Gamma }}(\omega )`$ does determine the response to small perturbations but, in contrast to linear response regime, it does not relate the total driving force $`\delta S/\delta 𝐗`$ to the evolution rate $`\dot{𝐀}`$ (the relationships is given by non-linear Eq. (14) or (15)).
Surprisingly enough, the applicability of FDT might be extended even further, to time scales comparable or shorter than ergodicity time $`\tau _0`$. Namely, if the frequencies of interest are considerably higher than the characteristic relaxation rates (eigenvalues of relaxation matrix $`\widehat{\lambda }\widehat{\kappa }`$), we may neglect the dependency of the the driving force on $`𝐚`$. In this case, the generalized entropy in $`(𝐚,𝐱)`$–representation can be written in the following form:
$$𝒮^{}=\sigma [𝐱(t)]\sigma [𝐡(t)/2]+_{t_i}^{t_f}\dot{𝐚}\left(\frac{𝐡(t)}{2}𝐱(t)\right)𝑑t.$$
(19)
Here $`𝐱(t)`$ is the deviation of $`𝐗`$ from its deterministic value $`𝐗_0`$, and $`\sigma [𝐱(t)]\mathrm{\Sigma }[𝐗_0+𝐱(t)]\mathrm{\Sigma }[𝐗_0]`$ is the corresponding deviation of $`\mathrm{\Sigma }`$-functional, which is no longer assumed to be local, i.e. $`\omega \tau _0`$ is not small. In this regime, to which one may refer as sub-ergodic, the statistics of $`𝐚`$ need not to be Gaussian.
It is easy to show that $`\sigma [𝐱(t)]`$ is the generating functional for $`\dot{𝐚}`$, i.e. its variations coincide with the corresponding correlators of $`\dot{𝐚}(t)`$. On the other hand, in accordance with (19), these variations of $`\sigma `$ determine the response of the system in any order of $`𝐡`$, $`\widehat{\mathrm{\Gamma }}^{(n)}(\omega _1\mathrm{}\omega _n)\delta ^n𝐚_{\mathrm{\Sigma }\omega _k}/\delta 𝐡_{\omega _1}\mathrm{}\delta 𝐡_{\omega _n}`$. This results in the following extension of FDT to sub-ergodic time scales:
$$\widehat{\mathrm{\Gamma }}^{(n)}(\omega _1\mathrm{}\omega _n)=𝐚_{\mathrm{\Sigma }\omega _k}𝐚_{\omega _1}\mathrm{}𝐚_{\omega _n}_{}\underset{k=1}{\overset{n}{}}\frac{i\omega _k}{2}.$$
(20)
This version of FDT is remarkable: (i) it may be applicable to systems with considerable memory effects, e.g. glasses; (ii) it establishes the relationship between the non-linear response of a system and the deviation of its fluctuation statistics from Gaussian. The above relationship is very similar to recent results for Markovian stochastic processes, . It should be noted that the direct experimental or numeric check of Eq. (20) may be difficult to perform without correct interpretation of field $`𝐡`$ (its physical meaning is straightforward only in the case of ideal (fast) thermal bath coupled to the system). In particular, such interpretation may involve frequency-dependent temperature.
In conclusion, we have proposed a framework for construction of non-equilibrium macroscopic theory of a complex system, starting with its fundamental non-dissipative dynamics. This approach assumes the possibility of choosing a set of relatively slow (adiabatic) variables. Our major results include the general form of equation of motion of the system under a given thermodynamic driving force (GLE) and extended FDT. Among the immediate possible applications of our scheme is the development of non-equilibrium statistical theories of various complex systems starting with their model Hamiltonians, such as Heisenberg and $`XY`$ models , or Gross-Pitaevsky model of Bose condensate. Another intriguing direction of the development of the generalized entropy approach is its use for Landau-type description of bifurcations. It also provides us with an apparatus to study the problem of kinetic tunneling between various steady states (attractors) of a non-equilibrium system. A natural extension of our theory would be its quantum generalization.
Acknowledgement The author is grateful to T. Witten, P. Cvitanovic, A. Sengupta, B. Shraiman, Y. Rabin, C. Varma, E. Balkovski for useful discussions.
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# Multi-plateau magnetization curves of one-dimensional Heisenberg ferrimagnets
## I Introduction
Ground-state magnetization curves of low-dimensional quantum spin systems have been attracting much recent interest due to their nontrivial appearance contrasting with classical behaviors. A few years ago there appeared an epochal argument in the field. Generalizing the Lieb-Schultz-Mattis theorem , Oshikawa, Yamanaka, and Affleck proposed that magnetization plateaux of quantum spin chains should be quantized as
$$S_{\mathrm{unit}}m=\text{integer},$$
(1)
where $`S_{\mathrm{unit}}`$ is the sum of spins over all sites in the unit period and $`m`$ is the magnetization $`M`$ divided by the number of the unit cells. Their argument caused renewed interest in the pioneering calculations of a bond-trimerized spin-$`\frac{1}{2}`$ chain and further stimulated extensive investigations into quantum magnetization process. Quantized magnetization plateaux were reasonably detected for spin-$`\frac{1}{2}`$ , spin-$`1`$ , and spin-$`\frac{3}{2}`$ chains with modulated and/or anisotropic interactions. Totsuka , Cabra and Grynberg , and Honecker developed calculations of general polymerized spin chains. Numerous authors have been making further theoretical explorations into extended systems including spin ladders and layered magnets . Experimental observations have also followed.
Mixed-spin chains, which have vigorously been studied in recent years , also stimulate us in this context. Theoretical investigations into them are all the more interesting and important considering an accumulated chemical knowledge on ferrimagnetic materials. Kahn et al. succeeded in synthesizing a series of bimetallic chain compounds such as MM(pba)(H<sub>2</sub>O)<sub>3</sub>$``$$`2`$H<sub>2</sub>O (pba $`=`$ $`1,3`$-propylenebis(oxamato) $`=`$ C<sub>7</sub>H<sub>6</sub>N<sub>2</sub>O<sub>6</sub>) and MM(pbaOH)(H<sub>2</sub>O)<sub>3</sub> (pbaOH $`=`$ $`2`$-hydroxy-$`1,3`$-propylenebis(oxamato) $`=`$ C<sub>7</sub>H<sub>6</sub>N<sub>2</sub>O<sub>7</sub>), where the alternating magnetic ions M and M are flexible variables and therefore the low-dimensional ferrimagnetic behavior could systematically be observed. Caneschi et al. synthesized another series of mixed-spin chain compounds of general formula M(hfac)<sub>2</sub>NITR, where metal ion complexes M(hfac)<sub>2</sub> with hfac $`=`$ hexafluoroacetylacetonate are bridged by nitronyl nitroxide radicals NITR. There also exist purely organic molecule-based ferrimagnets , where sufficiently small magnetic anisotropy, whether of exchange-coupling type or of single-ion type, is advantageous for observations of essential quantum mixed-spin phenomena.
Magnetization curves of Heisenberg ferrimagnetic chains were revealed by Kuramoto . His argument covered an effect of next-nearest-neighbor interactions but the constituent spins were restricted to $`1`$ and $`\frac{1}{2}`$. Although an alignment of alternating spins $`S`$ and $`s`$ ($`S>s`$) in a field, as described by the Hamiltonian
$`={\displaystyle \underset{j=1}{\overset{N}{}}}`$ $`[`$ $`(1+\delta )(𝑺_j𝒔_j)_\alpha +(1\delta )(𝒔_j𝑺_{j+1})_\alpha `$ (2)
$``$ $`H(S_j^z+s_j^z)],`$ (3)
with $`(𝑺𝒔)_\alpha =S^xs^x+S^ys^y+\alpha S^zs^z`$, is so interesting as to possibly exhibit a series of quantized magnetization plateaux at $`m=\frac{1}{2}(1),\frac{3}{2}(2),\mathrm{},S+s1`$, its magnetization curves have not systematically been studied so far. In spite of the vigorous argument, there are few reports on multi-plateau magnetization curves. It is true that a double-plateau structure lies in NH<sub>4</sub>CuCl<sub>3</sub> , but it is owing to the variety of exchange interactions. We here demonstrate that the ferrimagnetic chain (3) generally exhibits a $`2s`$-plateau magnetization curve without any anisotropy and any bond polymerization, namely, even at $`\alpha =1`$ and $`\delta =0`$. We believe that the present calculations will accelerate physical measurements on vast ferrimagnetic chain compounds lying unexploited in the field of both inorganic and organic chemistry.
The ground state of the isotropic Hamiltonian (3) without the Zeeman term, which is a multiplet of spin $`(Ss)N`$, exhibits elementary excitations of two distinct types . The excitations of ferromagnetic aspect, reducing the ground-state magnetization, form a gapless dispersion relation, whereas those of antiferromagnetic aspect, enhancing the ground-state magnetization, are gapped from the ground state. Therefore we can readily understand the initial step at $`m=Ss`$ in the magnetization curve. In the Ising limit $`\alpha \mathrm{}`$, the initial plateau is nothing but the gapped excitation from the Néel-ordered state. The classical gap-generation mechanism is unique. Thus any magnetization curve in the Ising limit only has a single plateau. The scenario is qualitatively unchanged for an arbitrary $`\alpha `$ as long as we consider the classical vectors $`𝑺_j`$ and $`𝒔_j`$ of magnitude $`S`$ and $`s`$ instead of quantum spins. Therefore, the second and higher plateaux, if any, should generally be based on a quantum mechanism.
## II Numerical Procedure
We perform a scaling analysis on the numerically calculated energy spectra of finite clusters up to $`N=12`$. With $`E(N,M)`$ being the lowest energy in the subspace with a fixed magnetization $`M`$ for the Hamiltonian (3) without the Zeeman term, the upper and lower bounds of the field which induces the ground-state magnetization $`M`$ are expressed as
$$H_\pm (N,M)=\pm E(N,M\pm 1)E(N,M).$$
(4)
The length of the plateau with the unit-cell magnetization $`mM/N`$ is obtained as
$$\mathrm{\Delta }_N(m)=H_+(N,M)H_{}(N,M).$$
(5)
Therefore, calculating $`E(N,M)`$ at each sector of $`M`$ and extrapolating the resultant $`H_\pm (N,M)`$ with respect to $`N`$, we can obtain the thermodynamic-limit magnetization curves. Since the correlation length of the present system is considerably small , this scaling analysis works very well. The precision of the obtained magnetization curves is generally down to three decimal places. There is at most slight uncertainty in the second decimal place.
## III Results and Discussion
We show in Fig. 1 the thus-obtained magnetization curves. Making use of the Schwinger boson representation:
$$\begin{array}{cc}S_j^+=A_j^{}B_j,\hfill & S_j^z=\frac{1}{2}(A_j^{}A_jB_j^{}B_j),\hfill \\ s_j^+=a_j^{}b_j,\hfill & s_j^z=\frac{1}{2}(a_j^{}a_jb_j^{}b_j),\hfill \end{array}$$
(6)
the $`M=Ss`$ ground state of the decoupled dimers ($`\delta =1`$) are described as $`_j(A_j^{})^{Ss}(A_j^{}b_j^{}B_j^{}a_j^{})^{2s}|0`$, whose schematic representation is given in Fig. 2(a). Therefore, any plateau is enhanced by the bond alternation and the magnetization curve ends up with $`2s`$ steps, which are attributable to the crackion excitations , in the decoupled-dimer limit. Hence we here concentrate on the uniform-bond case ($`\delta =0`$). Surprisingly, in a certain region of $`\alpha `$ including the Heisenberg point, the spin-$`(S,s)`$ chain generally possesses a $`2s`$-plateau magnetization curve. To the best of our knowledge, this is the first report on the multi-plateau structure depending on neither anisotropy nor bond polymerization. $`2s`$ plateaux appear, but still, that does not mean the plateaux are dominated only by the smaller spin. Ferrimagnets have both ferromagnetic and antiferromagnetic features . The mixed aspect is explicitly exhibited, for instance, in their thermodynamics, where the specific heat and the magnetic susceptibility times temperature behave like $`T^{1/2}`$ and $`T^1`$ at low temperatures, respectively, whereas they exhibit a Schottky-like peak and a round minimum at intermediate temperatures. Figure 2(b), as well as Fig. 2(a), shows that the spin amplitude $`Ss`$ plays the ferromagnetic role, while $`2s`$ plays the antiferromagnetic one . Considering that any magnetization plateau originates from an antiferromagnetic interaction, it is convincing that $`2s`$ of the total spin amplitude $`S+s`$ contributes to the plateaux appearing.
Let us turn back to Fig. 1 and observe the plateaux more carefully, especially as functions of $`\alpha `$. In the cases of $`s=\frac{1}{2}`$, the plateaux are quite tough against the $`XY`$-like anisotropy. They are stable all over the antiferromagnetic-coupling region. These observations are in contrast with the classical behavior. The classical spin-$`(S,s)`$ Heisenberg Hamiltonian also exhibits the magnetization plateau at $`m=Ss`$, but it survives only a small amount of $`XY`$-like anisotropy. For instance, the critical value for $`(S,s)=(1,\frac{1}{2})`$ is estimated as $`\alpha _\mathrm{c}=0.943(1)`$ . The contrast between quantum spins and classical vectors suggests that the single plateaux in the quantum-spin magnetization curves may be attributed to the valence-bond excitation gap (valence-bond gap) rather than the Néel-state excitation gap (Néel gap). From this point of view, the $`\alpha `$-dependences of the two coexistent plateaux in the cases of $`s=1`$ are interesting. The tiny second plateaux are much more stable against the $`XY`$-like anisotropy than the steady-looking initial steps. Since the Néel state reaches the saturation via a single-step excitation, the second and higher plateaux should originate from the valence-bond gap. The magnitude of the gap exponentially decreases with the increase of $`m`$, but the quantum gap-generation mechanism itself is rather tough against the $`XY`$-like anisotropy. The initial plateaux in the multi-step magnetization curves, which are relatively unstable against the $`XY`$-like anisotropy, may be attributed to the Néel gap. It seems that the lowest-lying magnetization plateaux are of quantum appearance for $`s=\frac{1}{2}`$ but are of classical aspect for $`s1`$.
In order to characterize the plateaux, we introduce a variational wave function for the ground state of the model (3) as
$`|\mathrm{g}=c_\mathrm{N}{\displaystyle \underset{j=1}{\overset{N}{}}}(A_j^{})^{2S}(b_j^{})^{2s}|0`$ (7)
$`+{\displaystyle \underset{l=0}{\overset{2s}{}}}c_{\mathrm{VB}}^{(l)}{\displaystyle \underset{j=1}{\overset{N}{}}}(A_j^{})^{2Sl}(a_j^{})^{2sl}(A_j^{}b_j^{}B_j^{}a_j^{})^l|0,`$ (8)
where $`c_\mathrm{N}`$ and $`c_{\mathrm{VB}}^{(l)}`$ are the mixing coefficients. Since all the elemental states are asymptotically orthogonal to each other, the thermodynamic-limit variational ground states are considerably simple, as shown in Fig. 3, where we consider the Heisenberg point. The phase diagrams are exact on the line of $`\delta =1`$, where the spin-$`(\frac{3}{2},1)`$ chain, for example, reaches the saturation (S) via the double-bond dimer (DBD) and single-bond dimer (SBD) states. The point is that the $`\delta =0`$ ground states are better approximated by the decoupled-dimer states than by the Néel (N) states in the cases of $`s=\frac{1}{2}`$, while vice versa in all other cases. For $`s=\frac{1}{2}`$, the variational wave function (8) ends up with $`c_\mathrm{N}=0`$ all over the $`\delta `$-$`H`$ plain. The Néel-dimer crossover point $`\delta _\mathrm{c}`$ is given by
$$\delta _\mathrm{c}=\frac{2SsA+B+C}{AB+C},$$
(9)
with
$`A={\displaystyle \underset{l=1}{\overset{2s}{}}}\{{\displaystyle \frac{(2S2s+l)!(2S2s+l1)!}{[(2S2s)!]^2l!(l1)!}}`$ (10)
$`\times [S(S+1)(S2s+l1)(S2s+l)]`$ (11)
$`\times [s(s+1)(sl)(sl+1)]\}^{1/2}`$ (12)
$`/{\displaystyle \underset{l=0}{\overset{2s}{}}}{\displaystyle \frac{(2S2s+l)!}{(2S2s)!l!}},`$ (13)
$`B={\displaystyle \underset{l=0}{\overset{2s}{}}}{\displaystyle \frac{(2S2s+l)!}{(2S2s)!l!}}(S2s+l)(sl)`$ (14)
$`/{\displaystyle \underset{l=0}{\overset{2s}{}}}{\displaystyle \frac{(2S2s+l)!}{(2S2s)!l!}},`$ (15)
$`C={\displaystyle \underset{l=0}{\overset{2s}{}}}{\displaystyle \frac{(2S2s+l)!}{(2S2s)!l!}}(S2s+l)`$ (16)
$`\times {\displaystyle \underset{l=0}{\overset{2s}{}}}{\displaystyle \frac{(2S2s+l)!}{(2S2s)!l!}}(sl)`$ (17)
$`/{\displaystyle \underset{l=0}{\overset{2s}{}}}{\displaystyle \frac{(2S2s+l)!}{(2S2s)!l!}}.`$ (18)
$`\delta _\mathrm{c}`$ does not exist for $`s=\frac{1}{2}`$, whereas it is an increasing function of $`s`$ for $`s1`$, as shown in Fig. 4. Thus, the Néel-gap-like character of the initial plateaux in the multi-step magnetization curves become more and more settled as $`s`$ increases, which leads to the instability of the plateaux against the $`XY`$-like anisotropy. The exceptional cases of $`s=\frac{1}{2}`$ may be recognized as the quantum limit. Another simple calculation also supports this scenario. Let us consider a spin-wave description and a perturbation treatment of the antiferromagnetic excitation gap from the ground state. The spin-wave excitations are based on the Néel-order background, whereas the perturbation from the decoupled-dimer limit assumes the crackion-like excitations to appear in the valence-bond background. We compare in Table I both estimates with the exact values, namely, the upper critical fields for the initial plateaux. In the cases of $`s=\frac{1}{2}`$, the perturbation calculations are better than the spin-wave estimates, while vice versa in the cases of $`s1`$. We are again convinced that the single plateaux for $`s=\frac{1}{2}`$ are relatively of quantum aspect, while the lowest-magnetization plateaux in the multi-step process for $`s1`$ are relatively of classical aspect. All other plateaux, the second and higher steps, should essentially be based on the quantum mechanism. Now here is a question: As $`\alpha `$ increases, at which point do the quantum plateaux for $`m>Ss`$ disappear? Our numerical investigations estimate that they survive the whole region of $`\alpha >1`$ and disappear in the Ising limit.
Finally we draw the ground-state phase diagrams on the $`\alpha \delta `$-plane. If the system is massive at the sector of magnetization $`M`$, $`H_\pm (N,M)`$ are extrapolated to different thermodynamic-limit values $`H_\pm (m)`$ with exponential size corrections. On the other hand, in the critical phase, $`H_\pm (N,M)`$ converge to the same value as
$$H_\pm (N,M)H(m)\pm \frac{\pi v_\mathrm{s}\eta }{N},$$
(19)
where $`v_\mathrm{s}`$ is the spin-wave velocity and $`\eta `$ is the critical index for the relevant spin-correlation function. Therefore, we can visualize the phase transition by plotting the scaled gap $`N\mathrm{\Delta }_N(m)`$ as a function of $`\alpha `$ as shown in Fig. 5. The phase boundaries could in principle be extracted
from the phenomenological renormalization-group equation taking $`\mathrm{\Delta }_N(m)`$ as the order parameter. However, the anisotropy-induced breakdown of the plateau is a transition of the Kosterlitz-Thouless type and the fixed point could only be determined with great uncertainty . Thus we here rely upon the critical exponent $`\eta `$ which should cross over the value $`\frac{1}{4}`$ on the phase boundary. Provided $`v_\mathrm{s}`$ is given, we can estimate $`\eta `$ using the scaling law (19). We obtain $`v_\mathrm{s}`$ directly from the dispersion relation. Using the scaling relation
$$\frac{E(N,M)}{N}\epsilon (m)\frac{\pi cv_\mathrm{s}}{N^2},$$
(20)
we further verify the central charge being unity in the critical region, though it is not so useful in determining the phase boundary. The thus-obtained phase boundaries are shown by solid lines in Fig. 6. The single plateaux, with a quantum base, are stable over the whole antiferromagnetic-coupling region (a,b), while the initial plateaux in the multi-step process, taking on a classical character, less survive the $`XY`$-like anisotropy than the quantum higher plateaux (c,d). The existence of the second plateaux without any bond polymerization, which is the main issue in the present article, should be verified very carefully. So then we further employ an idea of the level spectroscopy , thus called, in analyzing the second plateau. Comparing the relevant excitation energies whose scaling dimensions are $`2`$ at the critical point, we recognize the level crossing of them as the phase boundary. The thus-detected transitions are also plotted by broken lines in Fig. 6. The slight difference between the two estimates inevitably arises from the logarithmic corrections to the scaling law (19), where the level spectroscopy is more reliable than the naivest scaling analysis. Anyway, we may now fully be convinced of the existence of the novel multi-plateau magnetization curves.
## IV Summary and Future Aspect
The one-dimensional Heisenberg ferrimagnet with alternating spins $`S`$ and $`s`$ exhibits a $`2s`$-plateau magnetization curve even at the most symmetric point. It is interesting to compare the present spin-$`(S,s)`$ ferrimagnetic chain with the spin-$`\frac{1}{2}`$ bond-polymerized chain of period $`2(S+s)`$, that is, the $`(2S1)`$-times-ferromagnetic-antiferromagnetic-$`(2s1)`$-times-ferromagnetic-antiferromagnetic chain. In the strong ferromagnetic-coupling limit, the latter may be regarded as equivalent to the former. In the same meaning, the spin-$`S`$ antiferromagnetic Heisenberg chain can be viewed as a spin-$`\frac{1}{2}`$ bond-polymerized chain of period $`2S`$. Such replica chains generally exhibit magnetization plateaux in certain regions of the ratio of the ferromagnetic coupling $`J_\mathrm{F}`$ to the antiferromagnetic one $`J_\mathrm{A}`$, $`\gamma J_\mathrm{F}/J_\mathrm{A}`$. However, in the case of the spin-$`\frac{1}{2}`$ ferromagnetic-ferromagnetic-antiferromagnetic chain which is the replica model of the spin-$`\frac{3}{2}`$ antiferromagnetic Heisenberg chain, Okamoto and Hida reported that the plateau vanishes at $`\gamma =45`$ and therefore the pure Heisenberg chain exhibits no plateau. Thus the plateaux observed here in the most symmetric Heisenberg ferrimagnets still interest us to a great extent.
As we come up the steps, the plateau length exponentially decreases. It is hard to numerically observe the higher-lying plateaux, still harder experimentally. Only the first and second plateaux may lie within the limits of measurement. In this context, we are fortunate to have a series of bimetallic quasi-one-dimensional complexes MM(EDTA)$``$6H<sub>2</sub>O ($`\text{M},\text{M}^{}=\text{Mn},\text{Co},\text{Ni},\text{Cu}`$) . Their exchange coupling constants are all about $`10k_\mathrm{B}`$\[K\] and thus the complete magnetization curves could technically be observed. Magnetization measurements on them, especially with $`\text{M}=\text{Mn}(S=\frac{5}{2}),\text{Co}(S=\frac{3}{2})`$ and $`\text{M}^{}=\text{Ni}(s=1)`$, are encouraged. The plateaux would more or less be obscured in any actual measurement, but they, however small, should necessarily be detected by some anomaly in the magnetic susceptibility. The chemical modification of the bond alternation $`\delta `$ and/or the exchange anisotropy $`\alpha `$ must help us to directly observe the second-step plateaux, though we take main interest in the Heisenberg point.
###### Acknowledgements.
The authors thank Dr. K. Okamoto for useful discussion. This work is supported by the Japanese Ministry of Education, Science, and Culture through Grant-in-Aid No. 11740206 and by the Sanyo-Broadcasting Foundation for Science and Culture. The computation was done in part using the facility of the Supercomputer Center, Institute for Solid State Physics, University of Tokyo.
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# Gauge Field Theory Coherent States (GCS) : IV. Infinite Tensor Product and Thermodynamical Limit
## 1 Introduction
Quantum General Relativity (QGR) has matured over the past decade to a mathematically well-defined theory of quantum gravity. In contrast to string theory, by definition QGR is a manifestly background independent, diffeomorphism invariant and non-perturbative theory. The obvious advantage is that one will never have to postulate the existence of a non-perturbative extension of the theory, which in string theory has been called the still unknown M(ystery)-Theory.
The disadvantage of a non-perturbative and background independent formulation is, of course, that one is faced with new and interesting mathematical problems so that one cannot just go ahead and “start calculating scattering amplitudes”: As there is no background around which one could perturb, rather the full metric is fluctuating, one is not doing quantum field theory on a spacetime but only on a differential manifold. Once there is no (Minkowski) metric at our disposal, one loses familiar notions such as causality structure, locality, Poincaré group and so forth, in other words, the theory is not a theory to which the Wightman axioms apply. Therefore, one must build an entirely new mathematical apparatus to treat the resulting quantum field theory which is drastically different from the Fock space picture to which particle physicists are used to.
As a consequence, the mathematical formulation of the theory was the main focus of research in the field over the past decade. The main achievements to date are the following (more or less in chronological order) :
* Kinematical Framework
The starting point was the introduction of new field variables for the gravitational field which are better suited to a background independent formulation of the quantum theory than the ones employed until that time. In its original version these variables were complex valued, however, currently their real valued version, considered first in for classical Euclidean gravity and later in for classical Lorentzian gravity, is preferred because to date it seems that it is only with these variables that one can rigorously define the kinematics andf dynamics of Euclidean or Lorentzian quantum gravity .
These variables are coordinates for the infinite dimensional phase space of an $`SU\left(2\right)`$ gauge theory subject to further constraints besides the Gauss law, that is, a connection and a canonically conjugate electric field. As such, it is very natural to introduce smeared functions of these variables, specifically Wilson loop and electric flux functions. (Notice that one does not need a metric to define these functions, that is, they are background independent). This had been done for ordinary gauge fields already before in and was then reconsidered for gravity (see e.g. ).
The next step was the choice of a representation of the canonical commutation relations between the electric and magnetic degrees of freedom. This involves the choice of a suitable space of distributional connections and a faithful measure thereon which, as one can show , is $`\sigma `$-additive. The corresponding $`L_2`$ Hilbert space and its generalization will be henceforth called the Ashtekar-Isham-Lewandowski-Baez-Sawin (AILBS) Hilbert space. The proof that the AILBS Hilbert space indeed solves the adjointness relations induced by the reality structure of the classical theory as well as the canonical commutation relations induced by the symplectic structure of the classical theory can be found in . Independently, a second representation of the canonical commutation relations, called the loop representation, had been advocated (see e.g. and especially and references therein) but both representations were shown to be unitarily equivalent in (see also for a different method of proof).
This is then the first major achievement : The theory is based on a rigorously defined kinematical framework.
* Geometrical Operators
The second major achievement concerns the spectra of positive semi-definite, self-adjoint geometrical operators measuring lengths , areas and volumes of curves, surfaces and regions in spacetime. These spectra are pure point (discete) and imply a discrete Planck scale structure. It should be pointed out that the discreteness is, in contrast to other approaches to quantum gravity, not put in by hand but it is a prediction !
* Regularization- and Renormalization Techniques
The third major achievement is that there is a new regularization and renormalization technique for diffeomorphism covariant, density-one-valued operators at our disposal which was successfully tested in model theories . This technique can be applied, in particular, to the standard model coupled to gravity and to the Poincaré generators at spatial infinity . In particular, it works for Lorentzian gravity while all earlier proposals could at best work in the Euclidean context only (see, e.g. and references therein). The algebra of important operators of the resulting quantum field theories was shown to be consistent . Most surprisingly, these operators are UV and IR finite ! Notice that, at least as far as these operators are concerned, this result is stronger than the believed but unproved finiteness of scattering amplitudes order by order in perturbation theory of the five critical string theories, in a sense we claim that the perturbation series converges. The absence of the divergences that usually plague interacting quantum fields propagating on a Minkowski background can be understood intuitively from the diffeomorphism invariance of the theory : “short and long distances are gauge equivalent”. We will elaborate more on this point in future publications.
* Spin Foam Models
After the construction of the densely defined Hamiltonian constraint operator of , a formal, Euclidean functional integral was constructed in and gave rise to the so-called spin foam models (a spin foam is a history of a graph with faces as the history of edges) . Spin foam models are in close connection with causal spin-network evolutions , state sum models and topological quantum field theory, in particular BF theory . To date most results are at a formal level and for the Euclidean version of the theory only but the programme is exciting since it may restore manifest four-dimensional diffeomorphism invariance which in the Hamiltonian formulation is somewhat hidden.
* Finally, the fifth major achievement is the existence of a rigorous and satisfactory framework for the quantum statistical description of black holes which reproduces the Bekenstein-Hawking Entropy-Area relation and applies, in particular, to physical Schwarzschild black holes while stringy black holes so far are under control only for extremal charged black holes.
Summarizing, the work of the past decade has now culminated in a promising starting point for a quantum theory of the gravitational field plus matter and the stage is set to pose and answer physical questions.
The most basic and most important question that one should ask is : Does the theory have classical general relativity as its classical limit ? Notice that even if the answer is negative, the existence of a consistent, interacting, diffeomorphism invariant quantum field theory in four dimensions is already a quite non-trivial result. However, we can claim to have a satisfactory quantum theory of Einstein’s theory only if the answer is positive.
In order to address this question with a mathematically well-defined procedure we have developed in a theory of coherent states for the matter content of the standard model (with possible supersymmetric extensions) coupled to gravity. These states are labelled by classical solutions to the field equations and have the property that a) the expectation values of densely defined field operators with respect to these states take the value prescribed by the classical solution and b) they saturate the Heisenberg uncertainty bound without quenching.
The way this has been achieved so far is the following : The degrees of freedom of, say, the gravitational field, are labelled by piecewise analytic (smooth) graphs (webs) composed of a finite number of edges (paths) only. For each such graph one finds a subspace of the AILBS Hilbert space which is the finite tensor product of mutually isomorphic Hilbert spaces, one for each edge (path) of the graph (web). The closure of finite linear combinations of vectors from these subspaces labelled by graphs (webs), which turn out to be mutually orthogonal, is forms the AILBS Hilbert space. What has been done in is to develop a theory of coherent states for each of these Hilbert spaces labelled by a finite graph $`\gamma `$. More precisely, one constructs coherent states $`\psi _e^s`$ for each of the Hilbert spaces labelled by a single edge $`e`$ of $`\gamma `$ and the classicality parameter $`s`$ ($`s0`$ is the classical limit) and then the coherent state for the whole graph $`\gamma `$ is simply the tensor product of those for each of its edges.
This framework is sufficient if the initial data hypersurface $`\mathrm{\Sigma }`$ is compact since one can describe the quantum metric as precisely as one wishes in terms of finite graphs by taking the graph to be finer and finer, filling $`\mathrm{\Sigma }`$ more and more densely. However, if $`\mathrm{\Sigma }`$ is non-compact, say of the topology of $`\text{ }\mathrm{R}^3`$ as required for Minkwoski space or the Kruskal extension of the Schwarzschild spacetime which in turn are the most important spacetimes if we want to make contact with the low energy physics of the standard model, scattering theory, Hawking radiation and thus the semiclassical approximation of quantum gravity by the theory of ordinary Quantum Field Theory on (curved) backgrounds, then the above framework is insufficient. What one needs in this case is an infinite graph no matter how coarse the graph is, that is, no matter whether the lattice spacing is 1mm or of the order of the Planck length, in order to fill $`\mathrm{\Sigma }`$ everywhere we need an infinite graph, no region of $`\mathrm{\Sigma }`$ of infinite volume must be empty if we wish to approximate a non-degenerate metric as all the classical metrics are.
One may think that one can get away by taking an infinite superposition of states labelled by mutually different finite graphs. However, such states have infinite norm with respect to the AILBS scalar product as the following simple example shows : Namely, let $`\gamma _{\mathrm{}}`$ be a cubic lattice, an infinte graph filling all of $`\mathrm{\Sigma }:=\text{ }\mathrm{R}^3`$ as densely as we wish and construct the state $`\psi ^s:=_ez_eT_e`$ where the sum runs over all edges of $`\gamma _{\mathrm{}}`$, $`z_e`$ are complex coefficients and $`T_e`$ is some linear combination of spin-network states over $`e`$. Obviously, this state is an infinite linear combination of states over finite graphs. Then, because of homogenity, this state produces the correct classical limit, corresponding to, say, Minkowski space, for each of the holonomy operators $`\widehat{h}_e`$ at most if $`z_e=z`$ is independent of $`e`$ and $`T_e\left(h_e\right)=T\left(h_e\right)`$ is the same linear combination of spin-network states for each $`e`$. But then the norm of the state is formally $`\left|z\right|^2T^2_e1=\mathrm{}`$ and badly diverges.
On the other hand, we will show that one can give meaning to states of the form $`\psi _\gamma _{\mathrm{}}^s:=_e\psi _e^s`$ where $`\gamma _{\mathrm{}}`$ is an infinite graph and if one defines the inner product to be the product of the inner products of the tensor product factors, then $`\psi _\gamma _{\mathrm{}}^s=1`$ while the semiclassical behaviour with respect to every possible operator over $`\gamma _{\mathrm{}}`$ is preserved and identical to the one for finite tensor products. Notice that in our case the dimension of the Hilbert space over each edge is countably infinite. If $`\gamma _{\mathrm{}}`$ is countably infinite, the case to which we restrict in this paper, then the direct sum Hilbert space is still separable. But even if the Hilbert space over each edge would be only two-dimensional then the countably infinite tensor product Hilbert space is non-separable !
This article is organized as follows :
In section two we recall the basic kinematical structure of canonical Quantum General Relativity.
In section three we list the essential properties of our family of coherent states for finite tensor products as needed for the purpose of the present paper.
In section four we give an account of von Neumann’s theory of the Infinite Tensor Product for the general case, in particular the occurance of von Neumann algebras of different factor types induced by the operator algebras on each tensor product factor.
Section five contains the new results of this paper. We apply the general ITP theory to our situation focussing on general and abstract properties only. We extend the quantum kinematical framework of Ashtekar, Isham and Lewandowski to piecewise analytical, infinite graphs, connect it with the (semi)classical analysis for canonical quantum field theories over non-compact initial data hypersurfaces and finally discuss the transfer of dynamical results as obtained earlier for finite graphs.
In particular, for any possible solution of the Einstein field equations we are able to identify an element of the ITP Hilbert space, a so-called $`C_0`$-vector $`\mathrm{\Omega }`$ in von Neumann’s terminology, which in the theory of quantum fields propagating on curved background spacetimes, plays the role of the vacuum or ground state and which can be constructed purely in terms of our coherent states. Perturbations of this vacuum, which in von Neumann’s terminology lie in the subspace of the ITP Hilbert space generated by the strong equivalence class of the $`C_0`$-vector $`\mathrm{\Omega }`$, can naturally be identified with the usual Fock states of QFT on the curved background that $`\mathrm{\Omega }`$ approximates. This opens the possibility to make contact with the usual perturbation theory defined in terms of Fock states.
In fact, in we show that it is possible to map a precisely defined subspace of the ITP Hilbert space for Einstein-Maxwell theory, to the Fock space defined in terms of, say, $`n`$Photon states propagating on Minkowski spacetime up to corrections due to pure quantum gravity effects caused by the fluctuating nature of the quantum metric and which one hopes to measure in experiment. More precisely, this subspace is generated by the operator algebra of the Maxwell field acting on a $`C_0`$-vector $`\mathrm{\Omega }`$ of the Einstein-Maxwell ITP Hilbert space and which is cyclic for that subspace. This vector $`\mathrm{\Omega }`$ is a minimal uncertainty vector for Einstein-Maxwell theory approximating the Minkowski metric and vanishing electromagnetic field respectively. It should be noted, however, that all the states so constructed are states of the fully interacting Einstein-Maxwell theory and not only of the free Maxwell theory propagating on Minkowski space (an example of a free quantum field theory on a fixed curved background). The two sets of states so constructed are in a one-to one and onto correspondence, leading to expectation values for physical operators which coincide to lowest order in the Planck length. As far as quantum gravity corrections are concerned, however, these states are physically very different, the states of the interacting theory give rise to the so-called $`\gamma `$ray-burst effect which is just one way to measure the Poincaré non-invariance of the present state of our universe at the fundamental level. In we will explicitly compute the size of this effect from first principles by a down-to-the-ground-computation, thereby significantly improving the results of .
## 2 Kinematical Structure of Diffeomorphism Invariant Quantum Gauge Theories
In this section we will recall the main ingredients of the mathematical formulation of (Lorentzian) diffeomorphism invariant classical and quantum field theories of connections with local degrees of freedom in any dimension and for any compact gauge group. See and references therein for more details.
### 2.1 Classical Theory
Let $`G`$ be a compact gauge group, $`\mathrm{\Sigma }`$ a $`D`$dimensional manifold and consider a principal $`G`$bundle with connection over $`\mathrm{\Sigma }`$. Let us denote the pull-back (by local sections) to $`\mathrm{\Sigma }`$ of the connection by $`A_a^i`$ where $`a,b,c,..=1,..,D`$ denote tensorial indices and $`i,j,k,..=1,..,dim\left(G\right)`$ denote indices for the Lie algebra of $`G`$. Likewise, consider a vector bundle of electric fields, whose projection to $`\mathrm{\Sigma }`$ is a Lie algebra valued vector density of weight one. We will denote the set of generators of the rank $`N1`$ Lie algebra of $`G`$ by $`\tau _i`$ which are normalized according to $`\text{tr}\left(\tau _i\tau _j\right)=N\delta _{ij}`$ and $`[\tau _i,\tau _j]=2f_{ij}^k\tau _k`$ defines the structure constants of $`Lie\left(G\right)`$.
Let $`F_i^a`$ be a Lie algebra valued vector density test field of weight one and let $`f_a^i`$ be a Lie algebra valued covector test field. We consider the smeared quantities
$$F\left(A\right):=_\mathrm{\Sigma }d^DxF_i^aA_a^i\text{ and }E\left(f\right):=_\mathrm{\Sigma }d^DxE_i^af_a^i$$
(2.1)
While both objects are diffeomorphism covariant, only the latter is gauge covariant, one reason to introduce the singular smearing discussed below. The choice of the space of pairs of test fields $`(F,f)\stackrel{~}{𝒮}`$ depends on the boundary conditions on the space of connections and electric fields which in turn depends on the topology of $`\mathrm{\Sigma }`$ and will not be specified in what follows.
Consider the set $`M`$ of all pairs of smooth functions $`(A,E)`$ on $`\mathrm{\Sigma }`$ such that (2.1) is well defined for any $`(F,f)𝒮`$. We define a topology on $`M`$ through the globally defined metric :
$`d_{\rho ,\sigma }[(A,E),(A^{},E^{})]`$
$`:=`$ $`\sqrt{{\displaystyle \frac{1}{N}}{\displaystyle _\mathrm{\Sigma }}d^Dx\left[\sqrt{det\left(\rho \right)}\rho ^{ab}\text{tr}\left(\left[A_aA_a^{}\right]\left[A_bA_b^{}\right]\right)+{\displaystyle \frac{[\sigma _{ab}\text{tr}\left([E^aE^a][E^bE^b]\right)}{\sqrt{det\left(\sigma \right)}}}\right]}`$
where $`\rho _{ab},\sigma _{ab}`$ are fiducial metrics on $`\mathrm{\Sigma }`$ of everywhere Euclidean signature. Their fall-off behaviour has to be suited to the boundary conditions of the fields $`A,E`$ at spatial infinity (if $`\mathrm{\Sigma }`$ is spatially non-compact). Notice that the metric (2.1) on $`M`$ is gauge invariant. It can be used in the usual way to equip $`M`$ with the structure of a smooth, infinite dimensional differential manifold modelled on a Banach (in fact Hilbert) space $``$ where $`𝒮\times 𝒮`$. (It is the weighted Sobolev space $`H_{0,\rho }^2\times H_{0,\sigma ^1}^2`$ in the notation of ).
Finally, we equip $`M`$ with the structure of an infinite dimensional symplectic manifold through the following strong (in the sense of ) symplectic structure
$$\mathrm{\Omega }((f,F),(f^{},F^{}))_m:=_\mathrm{\Sigma }d^Dx\left[F_i^af_a^iF_i^af_a^i\right]\left(x\right)$$
(2.3)
for any $`(f,F),(f^{},F^{})`$. We have abused the notation by identifying the tangent space to $`M`$ at $`m`$ with $``$. To see that $`\mathrm{\Omega }`$ is a strong symplectic structure one uses standard Banach space techniques. Computing the Hamiltonian vector fields (with respect to $`\mathrm{\Omega }`$) of the functions $`E\left(f\right),F\left(A\right)`$ we obtain the following elementary Poisson brackets
$$\{E\left(f\right),E\left(f^{}\right)\}=\{F\left(A\right),F^{}\left(A\right)\}=0,\{E\left(f\right),A\left(F\right)\}=F\left(f\right)$$
(2.4)
As a first step towards quantization of the symplectic manifold $`(M,\mathrm{\Omega })`$ one must choose a polarization. As usual in gauge theories, we will use connections as the configuration variables and electric fields as canonically conjugate momenta. As a second step one must decide on a complete set of coordinates of $`M`$ which are to become the elementary quantum operators. The analysis just outlined suggests to use the coordinates $`E\left(f\right),F\left(A\right)`$. However, the well-known immediate problem is that these coordinates are not gauge covariant. Thus, we proceed as follows :
The idea is to construct the theory from smaller building blocks, labelled by graphs embedded into $`\mathrm{\Sigma }`$. In the literature, two sets of graphs, labelling the so-called cylindrical functions, have been proposed : the set of finite piecewise analytical graphs $`\mathrm{\Gamma }_0^\omega `$ in and in the restriction $`\mathrm{\Gamma }_0^{\mathrm{}}`$ to so-called “webs” of the set of all piecewise smooth graphs $`\mathrm{\Gamma }^{\mathrm{}}`$. (We do not discuss here a third alternative, the set of finite piecewise linear graphs ). Here we call a graph $`\gamma `$ finite if its sets of oriented edges $`e`$ and vertices $`v`$ respectively, denoted by $`E\left(\gamma \right)`$ and $`V\left(\gamma \right)`$ respectively, have finite cardinality. A web is a special kind of a piecewise smooth graph which may not be finite but which can be obtained as the union of a finite number of smooth curves with finite range (the diffeomorphic image in $`\mathrm{\Sigma }`$ of a closed interval in $`\mathrm{R}`$) and such that its vertex set has a finite number of accumulation points. (In fact, this is the essential difference between $`\mathrm{\Gamma }_0^\omega `$ and $`\mathrm{\Gamma }_0^{\mathrm{}}`$ since a graph generated by a finite number of analytical curves is a piecewise analytical, finite graph which cannot have any accumulation points). There are some additional restrictions on the common intersections of the curves in a web which we do not need to explain here, see for all details. It is not difficult to prove that both $`\mathrm{\Gamma }_0^\omega ,\mathrm{\Gamma }_0^{\mathrm{}}`$ are closed under forming finite numbers of intersections and unions.
In this paper we are going to extend the framework to truly infinite graphs. That is, a priori, we do not impose any finiteness restriction neither on the number of edges or vertices of a graph nor on the range of its edges. Various extensions are possible. A simple possibility is the set $`\mathrm{\Gamma }^\omega `$ of piecewise analytic graphs with possibly a countably infinite number of edges. Such graphs can still have accumulation points of edges and vertices (e.g. the graph which looks like a ladder in a two-plane whose spokes are mutually parallel and come arbitrarily close to each other). An even simpler choice is the set $`\mathrm{\Gamma }_\sigma ^\omega `$ of piecewise analytic, $`\sigma `$-finite graphs which can be considered in locally compact manifolds $`\mathrm{\Sigma }`$ (every point has a compact neighbourhood) which, of course, is satisfied for any finite-dimensional manifold that we have in mind here. They are characterized by the fact that $`\gamma U\mathrm{\Gamma }_0^\omega `$, i.e. the restriction of $`\gamma `$ to any compact set is a piecewise analytic finite graph whose number of edges is uniformly bounded. More precisely :
###### Definition 2.1
Let $`\mathrm{\Sigma }`$ be a locally compact manifold. A graph $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$ is said to be a piecewise analytic, $`\sigma `$-finite graph, if for each compact subset $`U\mathrm{\Sigma }`$ the restriction of the graph is a finite graph, $`\gamma U\mathrm{\Gamma }_0^\omega `$. Moreover, for any compact cover $`𝒰`$ of $`\mathrm{\Sigma }`$ the set $`\{|E(\gamma U)|;U𝒰\}`$ is bounded.
Clearly, truly infinite piecewise analytic graphs exist only if $`\mathrm{\Sigma }`$ is not compact and in this case $`\mathrm{\Gamma }_0^\omega `$ is a proper subset of $`\mathrm{\Gamma }^\omega `$. In order to obtain maximally nice graphs we will make the further restriction that $`\mathrm{\Sigma }`$ is paracompact, see section 5.1.
The next simple choice is the set $`\mathrm{\Gamma }^{\mathrm{}}`$ of all piecewise smooth graphs with possibly a countable number of edges and possibly a countable number of accumulation points. More properly, we should call them the set of infinite webs, that is, the web $`\gamma `$ is allowed to be generated by a countably infinite number of smooth curves such that for each accumulation point $`p_i,i=1,..,N\mathrm{}`$ there exists a neighbourhood $`U_i`$ such that the $`U_i`$ are mutually disjoint and such that $`\gamma `$ restricted to $`U_i`$ is an element of $`\mathrm{\Gamma }_0^{\mathrm{}}`$. It is a non-trivial task to decide whether any of the three sets $`\mathrm{\Gamma }^\omega ,\mathrm{\Gamma }_\sigma ^\omega ,\mathrm{\Gamma }^{\mathrm{}}`$ are closed under taking finite unions and we will do this in this paper only for $`\mathrm{\Gamma }_\sigma ^\omega `$, leaving the remaining cases for future publications.
Finally, we could consider $`\mathrm{\Gamma }`$, the set of all piecewise smooth, oriented graphs $`\gamma `$ embedded into $`\mathrm{\Sigma }`$. That is, we do not impose any restriction on the cardinality of the sets $`E\left(\gamma \right),V\left(\gamma \right)`$, or on the nature of the accumulation points. This set is trivially closed under arbitrary unions but it is beyond present analytical control, furthermore, it is not clear whether $`\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }^{\mathrm{}}`$ are really different and to analyze these questions is beyond the scope of the present paper, too.
Suffice it to say that for the purposes that we have in mind, to take the classical limit, it is sufficient to work with the set $`\mathrm{\Gamma }_\sigma ^\omega `$ that is technically much easier to handle. Thus, from now on we will assume that $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$, the typical graph that we will need in our applications and that is good to have in mind as an example is a regular cubic lattice in $`\text{ }\mathrm{R}^3`$.
Let $`\gamma `$ be a graph and $`e`$ an edge of $`\gamma `$. We denote by $`h_e\left(A\right)`$ the holonomy of $`A`$ along $`e`$ and say that a function $`f`$ on $`𝒜`$ is cylindrical with respect to $`\gamma `$ if there exists a function $`f_\gamma `$ on $`G^{\left|E\left(\gamma \right)\right|}`$ such that $`f=p_\gamma ^{}f_\gamma =f_\gamma p_\gamma `$ where $`p_\gamma \left(A\right)=\left\{h_e\left(A\right)\right\}_{eE\left(\gamma \right)}`$. The set of functions cylindrical over $`\gamma `$ is denoted by Cyl<sub>γ</sub>. Holonomies are invariant under reparameterizations of the edge and in this article we assume that the edges are always analyticity preserving diffeomorphic images from $`[0,1]`$ to a one-dimensional submanifold of $`\mathrm{\Sigma }`$ if it has compact range and from $`[0,1),(0,1],(0,1)`$ if it has semi-finite or infinite range. Gauge transformations are functions $`g:\mathrm{\Sigma }G;xg\left(x\right)`$ and they act on holonomies as $`h_eg\left(e\left(0\right)\right)h_eg\left(e\left(1\right)\right)^1`$ where in the (semi)finite case $`e\left(0\right)`$ or $`e\left(1\right)`$ or both are not points in $`\mathrm{\Sigma }`$ and we simply set $`g\left(e\left(0\right)\right)=1`$ or $`g\left(e\left(1\right)\right)=1`$, which is justified by the the boundary conditions, restricting gauge transformations to be trivial at spatial infinity.
Next, given a graph $`\gamma `$ we choose a polyhedronal decomposition $`P_\gamma `$ of $`\mathrm{\Sigma }`$ dual to $`\gamma `$. The precise definition of a dual polyhedronal decomposition can be found in but for the purposes of the present paper it is sufficient to know that $`P_\gamma `$ assigns to each edge $`e`$ of $`\gamma `$ an open “face” $`S_e`$ (a polyhedron of codimension one embedded into $`\mathrm{\Sigma }`$) with the following properties :
(1) the surfaces $`S_e`$ are mutually non-intersecting,
(2) only the edge $`e`$ intersects $`S_e`$, the intersection is transversal and consists only of one point which is an interiour point of both $`e`$ and $`S_e`$,
(3) $`S_e`$ carries the normal orientation which agrees with the orientation of $`e`$.
Furthermore, we choose a system $`\mathrm{\Pi }_\gamma `$ of paths $`\rho _e\left(x\right)S_e,xS_e,eE\left(\gamma \right)`$ connecting the intersection point $`p_e=eS_e`$ with $`x`$. The paths vary smoothly with $`x`$ and the triples $`(\gamma ,P_\gamma ,\mathrm{\Pi }_\gamma )`$ have the property that if $`\gamma ,\gamma ^{}`$ are diffeomorphic, so are $`P_\gamma ,P_\gamma ^{}`$ and $`\mathrm{\Pi }_\gamma ,\mathrm{\Pi }_\gamma ^{}`$.
With these structures we define the following function on $`(M,\mathrm{\Omega })`$
$$P_i^e(A,E):=\frac{1}{N}\text{tr}\left(\tau _ih_e(0,1/2)\left[_{S_e}h_{\rho _e\left(x\right)}E\left(x\right)h_{\rho _e\left(x\right)}^1\right]h_e(0,1/2)^1\right)$$
(2.5)
where $`h_e(s,t)`$ denotes the holonomy of $`A`$ along $`e`$ between the parameter values $`s<t`$, $``$ denotes the Hodge dual, that is, $`E`$ is a $`\left(D1\right)`$form on $`\mathrm{\Sigma }`$, $`E^a:=E_i^a\tau _i`$ and we have chosen a parameterization of $`e`$ such that $`p_e=e\left(1/2\right)`$.
Notice that in contrast to similar variables used earlier in the literature the function $`P_i^e`$ is gauge covariant. Namely, under gauge transformations it transforms as $`P^eg\left(e\left(0\right)\right)P^eg\left(e\left(0\right)\right)^1`$, the price to pay being that $`P^e`$ depends on both $`A`$ and $`E`$ and not only on $`E`$. The idea is therefore to use the variables $`h_e,P_i^e`$ for all possible graphs $`\gamma `$ as the coordinates of $`M`$.
The problem with the functions $`h_e\left(A\right)`$ and $`P_i^e(A,E)`$ on $`M`$ is that they are not differentiable on $`M`$, that is, $`Dh_e,DP_i^e`$ are nowhere bounded operators on $``$ as one can easily see. The reason for this is, of course, that these are functions on $`M`$ which are not properly smeared with functions from $`\stackrel{~}{𝒮}`$, rather they are smeared with distributional test functions with support on $`e`$ or $`S_e`$ respectively. Nevertheless, one would like to base the quantization of the theory on these functions as basic variables because of their gauge and diffeomorphism covariance. Indeed, under diffeomorphisms $`h_eh_{\phi ^1\left(e\right)},P_i^eP^{\phi ^1\left(e\right)}`$ where we abuse notation since $`P^e`$ depends also on $`S_e,\rho _e`$, see for more details. We proceed as follows.
###### Definition 2.2
By $`\overline{M}_\gamma `$ we denote the direct product $`[G\times Lie(G)]^{|E(\gamma )|}`$. The subset of $`\overline{M}_\gamma `$ of pairs $`(h_e(A),P_i^e(A,E))_{eE(\gamma )}`$ as $`(A,E)`$ varies over $`M`$ will be denoted by $`(\overline{M}_\gamma )_{|M}`$. We have a corresponding map $`p_\gamma :M\overline{M}_\gamma `$ which maps $`M`$ onto $`(\overline{M}_\gamma )_{|M}`$.
Notice that the set $`\left(\overline{M}_\gamma \right)_{|M}`$ is, in general, a proper subset of $`\overline{M}_\gamma `$, depending on the boundary conditions on $`(A,E)`$, the topology of $`\mathrm{\Sigma }`$ and the “size” of $`e,S_e`$. For instance, in the limit of $`e,S_eeS_e`$ but holding the number of edges fixed, $`\left(\overline{M}_\gamma \right)_{|M}`$ will consist of only one point in $`\overline{M}_\gamma `$. This follows from the smoothness of the $`(A,E)`$.
We equip a subset $`M_\gamma `$ of $`\overline{M}_\gamma `$ with the structure of a differentiable manifold modelled on the Banach space $`_\gamma =\text{ }\mathrm{R}^{2dim\left(G\right)\left|E\left(\gamma \right)\right|}`$ by using the natural direct product manifold structure of $`\left[G\times Lie\left(G\right)\right]^{\left|E\left(\gamma \right)\right|}`$. While $`\overline{M}_\gamma `$ is a kind of distributional phase space, $`M_\gamma `$ has suitable regularity properties similar to (2.1).
In order to proceed and to give $`M_\gamma `$ a symplectic structure derived from $`(M,\mathrm{\Omega })`$ one must regularize the elementary functions $`h_e,P_i^e`$ by writing them as limits (in which the regulator vanishes) of functions which can be expressed in terms of the $`F\left(A\right),E\left(f\right)`$. Then one can compute their Poisson brackets with respect to the symplectic structure $`\mathrm{\Omega }`$ at finite regulator and then take the limit pointwise on $`M`$. The result is the following well-defined strong symplectic structure $`\mathrm{\Omega }_\gamma `$ on $`M_\gamma `$.
$`\{h_e,h_e^{}\}_\gamma `$ $`=`$ $`0`$
$`\{P_i^e,h_e^{}\}_\gamma `$ $`=`$ $`\delta _e^{}^e{\displaystyle \frac{\tau _i}{2}}h_e`$
$`\{P_i^e,P_j^e^{}\}_\gamma `$ $`=`$ $`\delta ^{ee^{}}f_{ij}^kP_k^e`$ (2.6)
Since $`\mathrm{\Omega }_\gamma `$ is obviously block diagonal, each block standing for one copy of $`G\times Lie\left(G\right)`$, to check that $`\mathrm{\Omega }_\gamma `$ is non-degenerate and closed reduces to doing it for each factor together with an appeal to well-known Hilbert space techniques to establish that $`\mathrm{\Omega }_\gamma `$ is a surjection of $`_\gamma `$. This is done in where it is shown that each copy is isomorphic with the cotangent bundle $`T^{}G`$ equipped with the symplectic structure (2.1) (choose $`e=e^{}`$ and delete the label $`e`$).
Now that we have managed to assign to each graph $`\gamma `$ a symplectic manifold $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ we can quantize it by using geometric quantization. This can be done in a well-defined way because the relations (2.1) show that the corresponding operators are non-distributional. This is therefore a clean starting point for the regularization of any operator of quantum gauge field theory which can always be written in terms of the $`\widehat{h}_e,\widehat{P}^e,eE\left(\gamma \right)`$ if we apply this operator to a function which depends only on the $`h_e,eE\left(\gamma \right)`$.
The question is what $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$ has to do with $`(M,\mathrm{\Omega })`$. In it is shown that there exists a partial order $``$ on the set of triples $`(\gamma ,P_\gamma ,\mathrm{\Pi }_\gamma )`$ and one can form a generalized projective limit $`M_{\mathrm{}}`$ of the $`M_\gamma `$ (in particular, $`\gamma \gamma ^{}`$ means $`\gamma \gamma ^{}`$). Moreover, the family of symplectic structures $`\mathrm{\Omega }_\gamma `$ is self-consistent in the sense that if $`(\gamma ,P_\gamma ,\mathrm{\Pi }_\gamma )(\gamma ^{},P_\gamma ^{},\mathrm{\Pi }_\gamma ^{})`$ then $`p_{\gamma ^{}\gamma }^{}\{f,g\}_\gamma =\{p_{\gamma ^{}\gamma }^{}f,p_{\gamma ^{}\gamma }^{}g\}_\gamma ^{}`$ for any $`f,gC^{\mathrm{}}\left(M_\gamma \right)`$ and $`p_{\gamma ^{}\gamma }:M_\gamma ^{}M_\gamma `$ is a natural projection.
Now, via the maps $`p_\gamma `$ of definition 2.2 we can identify $`M`$ with a subset of $`M_{\mathrm{}}`$. Moreover, in it is shown that there is a generalized projective sequence $`(\gamma _n,P_{\gamma _n},\mathrm{\Pi }_{\gamma _n})`$ such that $`lim_n\mathrm{}p_{\gamma _n}^{}\mathrm{\Omega }_{\gamma _n}=\mathrm{\Omega }`$ pointwise in $`M`$. This displays $`(M,\mathrm{\Omega })`$ as embedded into a generlized projective limit of the $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$, intuitively speaking, as $`\gamma `$ fills all of $`\mathrm{\Sigma }`$, we recover $`(M,\mathrm{\Omega })`$ from the $`(M_\gamma ,\mathrm{\Omega }_\gamma )`$. On non-compact manifolds $`\mathrm{\Sigma }`$ this is possible only if the label set $`\mathrm{\Gamma }_\sigma ^\omega `$ contains infinite graphs.
It follows that quantization of $`(M,\mathrm{\Omega })`$, and conversely taking the classical limit, can be studied purely in terms of $`M_\gamma ,\mathrm{\Omega }_\gamma `$ for all $`\gamma `$. The quantum kinematical framework for this will be given in the next subsection.
### 2.2 Quantum Theory
At this point there is a clash with the previous subsection because the quantum kinematical structure has so far been defined only for the finite category of graphs $`\mathrm{\Gamma }_0^\omega `$. We thus have to extend this framework which we will do in section 5.1. However, as the structure from $`\mathrm{\Gamma }_0^\omega `$ can be nicely embedded into the more general context, we will repeat here the relevant notions for finite, piecewise analytical graphs $`\gamma `$.
Let us denote the set of all smooth connections by $`𝒜`$. This is our classical configuration space and we will choose for its coordinates the holonomies $`h_e\left(A\right),e\gamma ,\gamma \mathrm{\Gamma }_0^\omega `$. $`𝒜`$ is naturally equipped with a metric topology induced by (2.1).
Recall the notion of a function cylindrical over a graph from the previous subsection. A particularly useful set of cylindrical functions are the so-called spin-netwok functions which so far have been introduced only for $`\mathrm{\Gamma }_0^\omega `$, in fact, it is unclear whether one can define spin-network functions for all elements of $`\mathrm{\Gamma }_0^{\mathrm{}}`$, see for a discussion. As we will see in section 5, the spin-network bases proves to be of modest practical use in the context of $`\mathrm{\Gamma }_\sigma ^\omega `$ only, to be replaced by what we will call von Neumann bases based on $`C_0`$-vectors. To see what the problem is, we anyway have to introduce them here.
A spin-network function is labelled by a graph $`\gamma \mathrm{\Gamma }_0^\omega `$, a set of non-trivial irreducible representations $`\stackrel{}{\pi }=\left\{\pi _e\right\}_{eE\left(\gamma \right)}`$ (choose from each equivalence class of equivalent representations once and for all a fixed representant), one for each edge of $`\gamma `$, and a set $`\stackrel{}{c}=\left\{c_v\right\}_{vV\left(\gamma \right)}`$ of contraction matrices, one for each vertex of $`\gamma `$, which contract the indices of the tensor product $`_{eE\left(\gamma \right)}\pi _e\left(h_e\right)`$ in such a way that the resulting function is gauge invariant. We denote spin-network functions as $`T_I`$ where $`I=\{\gamma ,\stackrel{}{\pi },\stackrel{}{c}\}`$ is a compound label. One can show that these functions are linearly independent. From now on we denote by $`\stackrel{~}{\mathrm{\Phi }}_\gamma `$ finite linear combinations of spin-network functions over $`\gamma `$, by $`\mathrm{\Phi }_\gamma `$ the finite linear combinations of elements from any possible $`\stackrel{~}{\mathrm{\Phi }}_\gamma ^{},\gamma ^{}\gamma `$ a subgraph of $`\gamma `$ and by $`\mathrm{\Phi }`$ the finite linear combinations of spin-network functions from an arbitrary finite collection of graphs. Clearly $`\stackrel{~}{\mathrm{\Phi }}_\gamma `$ is a subspace of $`\mathrm{\Phi }_\gamma `$ which by itself is a proper subspace of the set Cyl$`{}_{\gamma }{}^{}{}_{}{}^{\mathrm{}}`$ of smooth cylindrical functions over $`\gamma `$. To express this distinction we will say that functions in $`\stackrel{~}{\mathrm{\Phi }}_\gamma `$ are labelled by “coloured graphs” $`\gamma `$ while functions in $`\mathrm{\Phi }_\gamma `$ are labelled simply by graphs $`\gamma `$, abusing the notation by using the same symbol $`\gamma `$.
The set $`\mathrm{\Phi }`$ of finite linear combinations of spin-network functions forms an Abelian algebra of functions on $`𝒜`$. By completing it with respect to the sup-norm topology it becomes an Abelian C algebra $``$ (here the compactness of $`G`$ is crucial). The spectrum $`\overline{𝒜}`$ of this algebra, that is, the set of all algebraic homomorphisms $`\text{ }\mathrm{C}`$ is called the quantum configuration space. This space is equipped with the Gel’fand topology, that is, the space of continuous functions $`C^0\left(\overline{𝒜}\right)`$ on $`\overline{𝒜}`$ is given by the Gel’fand transforms of elements of $``$. Recall that the Gel’fand transform is given by $`\widehat{f}\left(\overline{A}\right):=\overline{A}\left(f\right)\overline{A}\overline{𝒜}`$. It is a general result that $`\overline{𝒜}`$ with this topology is a compact Hausdorff space. Obviously, the elements of $`𝒜`$ are contained in $`\overline{𝒜}`$ and one can show that $`𝒜`$ is even dense . Generic elements of $`\overline{𝒜}`$ are, however, distributional.
The idea is now to construct a Hilbert space consisting of square integrable functions on $`\overline{𝒜}`$ with respect to some measure $`\mu `$. Recall that one can define a measure on a locally compact Hausdorff space by prescribing a positive linear functional $`\chi _\mu `$ on the space of continuous functions thereon. The particular measure we choose is given by $`\chi _{\mu _0}\left(\widehat{T}_I\right)=1`$ if $`I=\{\left\{p\right\},\stackrel{}{0},\stackrel{}{1}\}`$ and $`\chi _{\mu _0}\left(\widehat{T}_I\right)=0`$ otherwise. Here $`p`$ is any point in $`\mathrm{\Sigma }`$, $`0`$ denotes the trivial representation and $`1`$ the trivial contraction matrix. In other words, (Gel’fand transforms of) spin-network functions play the same role for $`\mu _0`$ as Wick-polynomials do for Gaussian measures and like those they form an orthonormal basis in the Hilbert space $`:=L_2(\overline{𝒜},d\mu _0)`$ obtained by completing their finite linear span $`\mathrm{\Phi }`$.
An equivalent definition of $`\overline{𝒜},\mu _0`$ is as follows :
$`\overline{𝒜}`$ is in one to one correspondence, via the surjective map $`H`$ defined below, with the set $`\overline{𝒜}^{}:=\text{Hom}(𝒳,G)`$ of homomorphisms from the groupoid $`𝒳`$ of composable, holonomically independent, analytical paths into the gauge group. The correspondence is explicitly given by $`\overline{𝒜}\overline{A}H_{\overline{A}}\text{Hom}(𝒳,G)`$ where $`𝒳eH_{\overline{A}}\left(e\right):=\overline{A}\left(h_e\right)=\stackrel{~}{h}_e\left(\overline{A}\right)G`$ and $`\stackrel{~}{h}_e`$ is the Gel’fand transform of the function $`𝒜Ah_e\left(A\right)G`$. Consider now the restriction of $`𝒳`$ to $`𝒳_\gamma `$, the groupoid of composable edges of the graph $`\gamma `$. One can then show that the projective limit of the corresponding cylindrical sets $`\overline{𝒜}_\gamma ^{}:=\text{Hom}(𝒳_\gamma ,G)`$ coincides with $`\overline{𝒜}^{}`$. Moreover, we have $`\left\{\left\{H\left(e\right)\right\}_{eE\left(\gamma \right)};H\overline{𝒜}_\gamma ^{}\right\}=\left\{\left\{H_{\overline{A}}\left(e\right)\right\}_{eE\left(\gamma \right)};\overline{A}\overline{𝒜}\right\}=G^{\left|E\left(\gamma \right)\right|}`$. Let now $`f`$ be a function cylindrical over $`\gamma `$ then
$$\chi _{\mu _0}\left(\stackrel{~}{f}\right)=_{\overline{𝒜}}𝑑\mu _0\left(\overline{A}\right)\stackrel{~}{f}\left(\overline{A}\right)=_{G^{\left|E\left(\gamma \right)\right|}}_{eE\left(\gamma \right)}d\mu _H\left(h_e\right)f_\gamma \left(\left\{h_e\right\}_{eE\left(\gamma \right)}\right)$$
where $`\mu _H`$ is the Haar measure on $`G`$. As usual, $`𝒜`$ turns out to be contained in a measurable subset of $`\overline{𝒜}`$ which has measure zero with respect to $`\mu _0`$. It turns out that it is this definition of the measure which can be extended to the category of infinite graphs.
Let, as before, $`\mathrm{\Phi }_\gamma `$ be the finite linear span of spin-network functions over $`\gamma `$ or any of its subgraphs and $`_\gamma `$ its completion with respect to $`\mu _0`$. Clearly, $``$ itself is the completion of the finite linear span $`\mathrm{\Phi }`$ of vectors from the mutually orthogonal $`\stackrel{~}{\mathrm{\Phi }}_\gamma `$. Our basic coordinates of $`M_\gamma `$ are promoted to operators on $``$ with dense domain $`\mathrm{\Phi }`$. As $`h_e`$ is group-valued and $`P^e`$ is real-valued we must check that the adjointness relations coming from these reality conditions as well as the Poisson brackets (2.1) are implemented on our $``$. This turns out to be precisely the case if we choose $`\widehat{h}_e`$ to be a multiplication operator and $`\widehat{P}_j^e=i\mathrm{}\kappa X_j^e/2`$ where $`\kappa `$ is the gravitational constant, $`X_j^e=X_j\left(h_e\right)`$ and $`X^j\left(h\right),hG`$ is the vector field on $`G`$ generating left translations into the $`jth`$ coordinate direction of $`Lie\left(G\right)T_h\left(G\right)`$ (the tangent space of $`G`$ at $`h`$ can be identified with the Lie algebra of $`G`$) and $`\kappa `$ is the coupling constant of the theory. For details see .
The question is now whether all of this structure can be extended to the infinite analytic category. In particular, in what sense does a spin-network function converge, what is the sup-norm for a function which is a finite linear combination of infinite products of holonomy functions etc. Obviously, at this point one must invoke the theory of the Infinite Tensor Product. We therefore have to postpone the answer to these questions to section 5.
## 3 Gauge Field Theory Coherent States
For a rather general idea of how to obtain coherent states for arbitrary canonically quantized quantum (field) theories and quantum gauge field theories in particular, see which is based on . Here we will stick with the heat kernel family initialized by the mathematician Brian Hall who proved that the associated Segal-Bargmann space is unitarily equivalent with the usual $`L_2`$ space. These results were extended to the Hilbert spaces underlying cylindrical functions of section 2.2 in . However, the semiclassical properties of these states were only later analyzed in .
### 3.1 Compact Group Coherent States
We will begin with only one copy of $`G`$ and consider the space of square integrable functions over $`G`$ with respect to the Haar measure $`d\mu _H`$, that is, the Hilbert space $`_G=L_2(G,d\mu _H)`$. Let $`s`$ be a positive real number, $`\pi `$ a (once and for all fixed, arbitrary representant from its equivalence class) irreducible representation, $`\chi _\pi `$ its character and $`d_\pi `$ its dimension. Let $`\mathrm{\Delta }`$ be the Laplacian on $`G`$, then it is well known that the $`dim_\pi ^2`$ functions $`\pi _{AB}`$ are eigenfunctions of $`\mathrm{\Delta }`$ with eigenvalue $`\lambda _\pi 0`$ which vanishes if and only if $`\pi `$ is the trivial representation.
Let $`hG`$ denote an element of $`G`$ and $`gG_{\text{ }\mathrm{C}}`$ an element of its complexification (for instance, if $`G=SU\left(2\right)`$ then $`G_{\text{ }\mathrm{C}}=SL(2,\text{ }\mathrm{C})`$). Then the (non-normalized) coherent state $`\psi _g^s`$ at classicality parameter $`s`$ and phase space point $`g`$ (the reason for this notation will be explained shortly) is defined by
$$\psi _g^s\left(h\right):=\underset{\pi }{}d_\pi e^{s\lambda _\pi /2}\chi _\pi \left(gh^1\right)=\left(e^{s\mathrm{\Delta }/2}\delta _h^{}\right)\left(h\right)_{|h^{}g}$$
(3.1)
that is, it is given by heat kernel evolution with time parameter $`s`$ of the $`\delta `$-distribution on $`G`$ followed by analytic continuation.
On $`_G`$ we introduce multiplication and derivative operators on the dense domain $`𝒟:=C^{\mathrm{}}\left(G\right)`$ by
$$\left(\widehat{h}_{AB}f\right)\left(h\right):=h_{AB}f\left(h\right)\text{ and }\left(\widehat{p}_jf\right)\left(h\right)=\frac{is}{2}\left(X_jf\right)\left(h\right)$$
(3.2)
where $`h_{AB}`$ denote the matrix elements of the defining representation of $`G`$ and $`i,j,k,..=1,2,..,dim\left(G\right)`$ and $`X_j\left(h\right)=\text{tr}\left(\left[\tau _jh\right]^T/h\right)`$ denotes the generator of right translations on $`G`$ into the $`j`$’th coordinate direction of $`Lie\left(G\right)`$, the Lie algebra of $`G`$. We choose a basis $`\tau _j`$ in $`Lie\left(G\right)`$ with respect to which tr$`\left(\tau _j\tau _k\right)=N\delta _{jk},[\tau _j,\tau _k]=2f_{jk}^l\tau _l`$ where $`N1`$ is the rank of $`G`$. The operators (3.2) enjoy the canonical commutation relations
$$[\widehat{h}_{AB},\widehat{h}_{CD}]=0,[\widehat{p}_j,\widehat{h}_{AB}]=\frac{is}{2}\left(\tau _j\widehat{h}\right)_{AB},[\widehat{p}_j,\widehat{p}_k]=isf_{jk}^l\widehat{p}_l$$
(3.3)
mirroring the classical Poisson brackets
$$\{h_{AB},h_{CD}\}=0,\{p_j,h_{AB}\}=\frac{1}{2}\left(\tau _jh\right)_{AB},\{p_j,p_k\}=f_{jk}^lp_l$$
(3.4)
on the phase space $`T^{}G`$, the cotangent bundle over $`G`$, where $`s`$ plays the role of Planck’s constant. It is easy to check that the CCR of (3.3) and the adjointness relations coming from $`\overline{p}_j=p_j,\overline{h_{AB}}=f_{AB}\left(h\right)`$ are faithfully implemented on $`_G`$. Here, $`f_{AB}`$ depends on the group, e.g. $`f_{AB}\left(h\right)=\left(h^1\right)_{BA}`$ for $`G=SU\left(N\right)`$, and $`\widehat{p}_j`$ is essentially self-adjoint with core $`𝒟`$.
We now consider $`G`$ as a subgroup of some unitary group so that the $`\tau _j`$ are antihermitean. We then identify $`G_{\text{ }\mathrm{C}}`$ with $`T^{}G`$ by the diffeomorphism
$$\varphi :T^{}GG^{\text{ }\mathrm{C}};(h,p)g:=e^{ip^j\tau _j/2}h=:Hh$$
(3.5)
where the inverse is simply given by the unique polar decomposition of $`gG_{\text{ }\mathrm{C}}`$. One can show that the symplectic structure (3.4) is compatible with the complex structure of $`G_{\text{ }\mathrm{C}}`$, displaying the complex manifold $`G_{\text{ }\mathrm{C}}`$ as a Kähler manifold.
Next we define on $`𝒟`$ the annihilation and creation operators
$$\widehat{g}_{AB}:=e^{s\mathrm{\Delta }/2}\widehat{h}_{AB}e^{s\mathrm{\Delta }/2}\text{ and }\left(\widehat{g}_{AB}\right)^{}:=e^{s\mathrm{\Delta }/2}f\left(\widehat{h}\right)_{AB}e^{s\mathrm{\Delta }/2}$$
(3.6)
Then, as one can show, $`\widehat{g}=e^{Ns/4}e^{i\widehat{p}_j\tau _j}\widehat{h}`$ so that the operator $`\widehat{g}`$ qualifies as a quantization of the polar decomposition of $`g`$.
To call these operators annihilation and creation operators is justified by the following list of properties with respect to the coherent states (3.1).
* Eigenstate Property
The states (3.1) are simultaneous eigenstates of the operators $`\widehat{g}_{AB}`$ with eigenvalue $`g_{AB}`$.
$$\widehat{g}_{AB}\psi _g^s=g_{AB}\psi _g^s$$
(3.7)
* Expectation Value Property
From this it follows immediately that the expectation values of the operators (3.6) with respect to the states (3.1) exactly equal their classical ones as prescribed by the phase space point $`g`$.
$$\frac{<\psi _g^s,\widehat{g}_{AB}\psi _g^s>}{\psi _g^s^2}=g_{AB},\frac{<\psi _g^s,\left(\widehat{g}_{AB}\right)^{}\psi _g^s>}{\psi _g^s^2}=\overline{g_{AB}}$$
(3.8)
* Saturation of the Heisenberg Uncertainty Bound
Consider the self-adjoint operators $`\widehat{x}_{AB}=\left(\widehat{g}_{AB}+\left[\widehat{g}_{AB}\right]^{}\right)/2,\widehat{y}_{AB}=\left(\widehat{g}_{AB}\left[\widehat{g}_{AB}\right]^{}\right)/\left(2i\right)`$ an their classical counterparts analogously built from $`g_{AB}`$. Then these operators saturate the Heisenberg uncertainty obstruction bound, moreover, the coherent states are unquenched for $`\widehat{x},\widehat{y}`$.
$$<\left[\widehat{x}_{AB}x_{AB}\right]^2>_g^s=<\left[\widehat{y}_{AB}y_{AB}\right]^2>_g^s=\frac{\left|<[\widehat{x}_{AB},\widehat{y}_{AB}]>_g^s\right|}{2}$$
(3.9)
where $`<.>_g^s`$ denotes the expectation value with respect to $`\psi _g^s`$. Thus they are minimal uncertainty states.
* Completeness and Segal-Bargmann Hilbert Space
There exists a measure $`\nu _s`$ on $`G_{\text{ }\mathrm{C}}`$ and a unitary map
$$\widehat{U}_s:_G_{G^{\text{ }\mathrm{C}}}:=\text{Hol}\left(G_{\text{ }\mathrm{C}}\right)L_2(G_{\text{ }\mathrm{C}},d\nu _s);f\left(h\right)\left(\widehat{U}_sf\right)\left(g\right):=\left(e^{s\mathrm{\Delta }/2}f\right)\left(h\right)_{hg}$$
(3.10)
between $`_G`$ and the space of $`\nu _s`$-square integrable, holomorphic functions, the Segal-Bargmann space. Moreover, the coherent states satisfy the overcompleteness relation
$$1__G=_{G_{\text{ }\mathrm{C}}}𝑑\nu _s\left(g\right)\widehat{P}_{\psi _g^s}$$
(3.11)
where $`\widehat{P}_f`$ denotes the projection onto the one-dimensional subspace of $`_G`$ spanned by the element $`f`$.
* Peakedness Properties
As usual, semiclassical behaviour of the system is most conveniently studied in terms of $`_{G_{\text{ }\mathrm{C}}}`$ because wave functions depend on phase space rather than on configuration space only. For instance, we have the peakedness property of the overlap function
$$\frac{|<\psi _g^s,\psi _g^{}^s>|^2}{\psi _g^s^2\psi _g^{}^s^2}=e^{\frac{F(p,p^{})+G(\theta ,\theta ^{})}{s}}\left(1K_s(g,g^{})\right)$$
(3.12)
where $`g=e^{ip_j\tau _j/2}e^{\theta _j\tau _j}`$ (and similar for $`g^{}`$) is the polar decomposition of $`g`$. $`K_s`$ is a positive function, uniformly bounded by a constant $`K_s^{}`$ independent of $`g,g^{}`$ that approaches zero exponentially fast as $`s0`$. $`F,G`$ are positive definite functions which take the value zero if and only if $`p_j=p_j^{}`$ and $`\theta _j=\theta _j^{}`$, moreover for small $`p_j^{}p_j,\theta _j^{}\theta _j`$ we have $`F(p,p^{})\left(p_j^{}p_j\right)^2,G(\theta ,\theta ^{})\left(\theta _j^{}\theta _j\right)^2`$ which shows that these states generalize the familiar $`T^{}\text{ }\mathrm{R}`$ coherent states to the non-linear setting of $`T^{}G`$. It can be shown that (3.12) is the probability density, with respect to the Liouville measure on $`T^{}G`$, to find the system at the phase space point $`g^{}`$ if it is in the state $`\psi _g^s`$ and that density equals unity at $`g=g^{}`$ and is otherwise strongly, Gaussian suppressed as $`s0`$ with width $`\sqrt{s}`$. Similar peakedness properties can be established in the configuration or momentum representation.
* Ehrenfest Theorems
The expectation value property holds for the operators $`\widehat{g}_{AB}`$ and $`\widehat{g}_{AB}^{}`$ at any value of $`s`$. For the remaining operators one can show
$$\underset{s0}{lim}<\widehat{h}_{AB}>_g^s=h_{AB}\text{ and }\underset{s0}{lim}<\widehat{p}_j>_g^s=p_j$$
(3.13)
where $`g=e^{ip_j\tau _j/2}h`$ and the convergence is exponentially fast. The result (3.13) extends to arbitrary polynomials of $`\widehat{h}_{AB},\widehat{p}_j`$ and even to non-polynomial, non-analytic functions of $`\widehat{p}_j`$ of the type that occur in quantum gravity, most importantly the volume operator mentioned in the introduction.
These beautiful properties of the states introduced by Hall will be extended to the gauge field theory case in the next subsection.
### 3.2 Graph Coherent States
Let $`\gamma \mathrm{\Gamma }_0^\omega `$ be a piecewise analytic, finite graph, that is, with a finite number of edges $`eE\left(\gamma \right)`$. For each edge $`e`$ we introduce the functions $`h_e,P_j^e`$ on $`(M,\omega )`$ as in subsection (2.1). Furthermore, we introduce the dimensionless quantities
$$p_j^e:=\frac{P_j^e}{a^{n_D}}\text{ and }s:=\frac{\mathrm{}\kappa }{a^{n_D}}$$
(3.14)
Here $`n_D=n_D^{}`$ if $`n_D^{}0`$ and $`n_D=1`$ otherwise where $`n_D^{}=D3`$ for Yang-Mills theory and $`n_D^{}=D1`$ for general relativity. Furthermore, if $`n_D^{}0`$ then $`a`$ is some fixed, arbitrary parameter of the dimension of a length (e.g. $`a=`$1cm), if $`n_D^{}=0`$ then $`a`$ is dimensionfree and $`\mathrm{}\kappa `$ is the Feinstruktur constant. Then the Poisson bracket relations of (2.1) become
$`[\widehat{h}_e,\widehat{h}_e^{}]_\gamma `$ $`=`$ $`0`$
$`[\widehat{p}_j^e,\widehat{h}_e^{}]_\gamma `$ $`=`$ $`is\delta _e^{}^e{\displaystyle \frac{\tau _j}{2}}\widehat{h}_e`$
$`[\widehat{p}_i^e,\widehat{p}_j^e^{}]_\gamma `$ $`=`$ $`is\delta ^{ee^{}}f_{ij}^k\widehat{p}_k^e`$ (3.15)
where the notation $`[.,.]_\gamma `$ indicates that all operators are restricted to the subspace $`_\gamma `$ of $``$. It is trivial to see that these relations classically carry over from the category $`\mathrm{\Gamma }_0^\omega `$ to the category $`\mathrm{\Gamma }_\sigma ^\omega `$.
We can now introduce the graph coherent states
$$\psi _{\gamma ,\stackrel{}{g}}^s\left(\stackrel{}{h}\right):=\underset{eE\left(\gamma \right)}{}\psi _{g_e}^s\left(h_e\right)$$
(3.16)
which are obviously neither gauge invariant nor diffeomorphism invariant. In it was indicated how to obtain diffeomorphism invariant coherent states from those in (3.16) and in the same was done in order to obtain gauge invariant ones, employing the group averaging technique . Since at the moment we are interested in issues related to the classical limit of the theory, in particular, whether we obtain in the classical limit the classical Einstein equations in an appropriate sense, we will not use those invariant states for two reasons :
1) In order to check the correctness of the classical limit we must verify, in particular, whether the quantum constraint algebra of the the quantum theory becomes the Dirac algebra in the classical limit. However, one cannot check an algebra on its kernel, see for a discussion.
2) As far as the gauge – and diffeomorphism constraint are concerned, it is perfectly fine to work with non-invariant coherent states because the corresponding gauge groups are represented as unitarily on the Hilbert space. This implies that expectation values of gauge – and diffeomorphism invariant operators are automatically also gauge – and diffeomorphism invariant and so qualify as expectation values of the reduced theory. Famously, the time reparameterizations associated with the Hamiltonian constraint of the theory cannot be unitarily represented and so the argument just given does not carry over to operators commuting with the Hamiltonian constraint. Presumably, the Hamiltonian constraint cannot be exponentiated at all and one will then have to work with its infinitesimal version. To pass then to the reduced theory one would need to work with coherent states that are annihilated by the Hamiltonian constraint (trivial representation of the “would be time reparameterization group”).
The coherent states (3.16) then form a valid starting point for adressing semiclassical questions in the case that $`\mathrm{\Sigma }`$ is compact, say, in some cosmological situations. To cover the case that $`\mathrm{\Sigma }`$ is asymptotically flat we must blow up the framework and pass to the Infinite Tensor Product.
## 4 The Abstract Infinite Tensor Product
> Bei Systemen mit N Teilchen ist der Hilbertraum das Tensorprodukt von den N Hilberträumen der einzelnen Teilchen.
> Das unendliche Tensorprodukt öffnet die Tür zu den mathematischen Finessen der Feldtheorie.
> (Walter Thirring)
The Infinite Tensor Product (ITP) of Hilbert spaces is a standard construction in statistical physics (through the thermodynamic (or infinite volume) limit) as well as in Operator Theory (von Neumann Algebras). In fact, the first examples of von Neumann algebras which are not of factor type $`I_n,I_{\mathrm{}}`$ (isomorphic to an algebra of bounded operators on a (separable) Hilbert space) have been constructed by using the ITP.
On the other hand, since the concept of separable (Fock) Hilbert spaces plays such a dominant role in high energy physics, presumably many theoretical physicists belonging to that community have never come across the concept of the Infinite Tensor Product (ITP) of Hilbert spaces which produces a non-separable Hilbert space in general. In fact, let us quote from Streater&Wightman, p. 86, 87 in that respect :
“…It is sometimes argued that in quantum field theory one is dealing with a system of an infinite number of degrees of freedom and so must use a non-separable Hilbert space….. Our next task is to explain why this is wrong, or at best is grossly misleading…..All these arguments make it clear that that there is no evidence that separable Hilbert spaces are not the natural state spaces for quantum field theory….”
Because of this, we have decided to include here a rough account of the most important concepts associated with the abstract Infinite Tensor Product. As it will become clear shortly, the ITP decomposes into an uncountable direct sum of Hilbert spaces which in most applications are separable. Each of these tiny subspaces of the complete ITP are isomorphic with the usual Fock spaces of quantum field theory on Minkowski space (or some other background). Presumably, the fact that one can do with separable Hilbert spaces in ordinary QFT is directly related to the fact that one fixes the background since this fixes the vacuum. The necessity to deal with the full ITP in quantum gravity could therefore be based on the fact that, in a sense, one has to consider all possible backgrounds at once ! More precisely, the metric cannot be fixed to equal a given background but becomes itself a fluctuating quantum operator.
We follow the beautiful and comprehensive exposition by von Neumann who invented the Infinite Tensor Product (ITP) more than sixty years ago already. The reader is recommended to consult this work for more details.
### 4.1 Definition of the Infinite Tensor Product of Hilbert Spaces
Let $``$ be some set of indices $`\alpha `$. We will not restrict the cardinality $`\left|\right|`$, rather for the sake of maximal generality we will allow $`\left|\right|`$ to take any possible value in the set of Cantor’s Alephs . The cardinality of the countably infinite sets is given by the non-standard number $`\mathrm{}`$. Then the cardinality of any other infinite set can be written as a function of $`\mathrm{}`$ (usually exponentials (of exponentials of..) $`\mathrm{}`$), e.g. the set $`\mathrm{R}`$ has the cardinality $`2^{\mathrm{}}`$. The mathematical justification for this amount of generality is because, following von Neumann ,
“…while the theory of enumerably infinite direct products $`_{n=1}^{\mathrm{}}_n`$ presents essentially new features, when compared with that of the finite $`_{n=1}^N_n`$, the passage from $`_{n=1}^{\mathrm{}}_n`$ to the general $`_\alpha _\alpha `$ presents no further difficulties…., the generalizations of the direct product lead to higher set-theoretical powers (G. Cantor’s “Alephs”), and to no measure problems at all.”
###### Definition 4.1
Let $`\{z_\alpha \}_\alpha `$ be a collection of complex numbers. The infinite product
$$\underset{\alpha }{}z_\alpha $$
(4.1)
is said to converge to the number $`z\text{ }\mathrm{C}`$
$`\delta >0_0(\delta ),|_0(\delta )|<\mathrm{}|z_{\alpha 𝒥}z_a|<\delta _0(\delta )𝒥,|𝒥|<\mathrm{}`$.
From the definition it is also straightforward to prove that if $`_\alpha z_\alpha ,_\alpha z_\alpha ^{}`$ converge to $`z,z^{}`$ respectively then $`_\alpha z_\alpha z_\alpha ^{}`$ converges to $`zz^{}`$.
Recall that a series $`_\alpha z_\alpha `$ converges if and only if it converges absolutely which in turn is the case if and only if $`z_\alpha =0`$ for all but countably infinitely many $`\alpha `$. The following theorem gives a useful convergence criterion for infinite products.
###### Theorem 4.1
1)
Let $`\rho _\alpha 0`$.
i) If $`\alpha _0\rho _{\alpha _0}=0`$ then $`_\alpha \rho _\alpha =0`$.
ii) If $`\rho _\alpha >0\alpha `$ then $`_\alpha \rho _\alpha `$ converges if and only if $`_\alpha \text{max}(\rho _\alpha 1,0)`$ converges.
iii) If $`\rho _\alpha >0\alpha `$ then $`_\alpha \rho _\alpha `$ converges to $`\rho >0`$ if and only if $`_\alpha |\rho _\alpha 1|`$ converges.
2)
Let $`z_\alpha =\rho _\alpha e^{i\phi _\alpha }\text{ }\mathrm{C}`$ where $`\rho _\alpha =|z_\alpha |,\phi _\alpha [\pi ,\pi ]`$. Then $`_\alpha z_\alpha `$ converges if and only if
i) either $`_\alpha \rho _\alpha `$ converges to zero in which case $`_\alpha z_\alpha =0`$,
ii) or $`_\alpha \rho _\alpha `$ converges to $`\rho >0`$ and $`_\alpha |\phi _\alpha |`$ converges in which case $`_\alpha z_\alpha =\rho e^{i_\alpha \phi _\alpha }`$.
In contrast to the case of an infinite series, absolute convergence of an infinite product does not imply convergence, the phases of the factors could fluctuate too wildly. This motivates the following definition.
###### Definition 4.2
Let $`z_\alpha \text{ }\mathrm{C}`$. We say that $`_\alpha z_\alpha `$ is quasi-convergent if $`_\alpha |z_\alpha |`$ converges. In this case we define the value of $`_\alpha z_\alpha `$ to equal $`_\alpha z_\alpha `$ if $`_\alpha z_\alpha `$ is even convergent and to equal zero otherwise.
This definition assigns a value to the infinite product of numbers which converge absolutely but not necessarily non-absolutely. As a corollary of theorem 4.1 we have
###### Corollary 4.1
Quasi-convergence of $`_\alpha z_\alpha `$ to a non-vanishing value is equivalent with convergence to the same value. A necessary and sufficient criterion is that $`z_\alpha 0\alpha `$ and that $`_\alpha |z_\alpha 1|`$ converges.
After having defined convergence for infinite products of complex numbers we are ready to turn to the ITP of Hilbert spaces.
###### Definition 4.3
Let $`_\alpha ,\alpha `$ be an arbitrary collection of Hilbert spaces. For a sequence $`f:=\{f_\alpha \}_\alpha ,f_\alpha _\alpha `$ the object
$$_f:=_\alpha f_\alpha $$
(4.2)
is called a $`C`$-vector provided that $`_\alpha f_\alpha _\alpha `$ converges, where $`||.||_\alpha `$ denotes the Hilbert norm of $`_\alpha `$. The set of $`C`$-vectors will be called $`V_C`$.
The following property holds for $`C`$-vectors, enabling us to compute their inner products.
###### Lemma 4.1
For two $`C`$-vectors $`_f=_\alpha f_\alpha ,_g=_\alpha g_\alpha `$ the inner product
$$<_f,_g>:=\underset{\alpha }{}<f_\alpha ,g_\alpha >_\alpha $$
(4.3)
is a quasi-convergent product of the individual inner products $`<f_\alpha ,g_\alpha >_\alpha `$ on $`_\alpha `$.
There are $`C`$-vectors $`_f`$ such that $`_\alpha f_\alpha _\alpha =0`$ although $`f_\alpha _\alpha >0\alpha `$. Thus, it is conceivable that it happens that $`<\mathrm{\Phi }_f,\mathrm{\Phi }_g>0`$ for some $`C`$-vector $`\mathrm{\Phi }_g`$. If that would be the case, the Schwarz inequality would be violated for the inner product (4.3) on $`C`$-vectors. That this is not the case is the content of the following lemma.
###### Lemma 4.2
Let $`_f`$ be a $`C`$-sequence with $`_\alpha f_\alpha _\alpha =0`$. Then $`<_f,_g>=0`$ for any $`C`$-vector $`_g`$.
To distinguish trivial $`C`$-vectors from non-trivial ones we define
###### Definition 4.4
A sequence $`(f_\alpha )`$ defines a $`C_0`$-vector $`_f=_\alpha f_\alpha `$ iff
$$\underset{\alpha }{}\left|f_\alpha _\alpha 1\right|$$
(4.4)
converges. The set of $`C_0`$-vectors will be denoted by $`V_0`$.
It is easy to prove by means of theorem 4.1 that every $`C_0`$-vector is a $`C`$-vector but only those $`C`$-vectors are $`C_0`$-vectors for which $`<_f,.>`$, considered as a linear functional on $`C`$-vectors, does not equal zero which by lemma 4.2 implies, in particular, that $`_\alpha f_\alpha _\alpha 0`$. It follows that the norm of a $`C_0`$-vector does not vanish, as the following lemma reveals.
###### Lemma 4.3
For any complex numbers, the convergence of one of $`_\alpha ||z_\alpha |1|,_\alpha ||z_\alpha |^21|`$ implies the convergence of the other.
Thus, since by definition of a $`C_0`$-vector $`_f`$ and theorem 4.11)iii) $`z_\alpha =f_\alpha _\alpha `$ satisfies the assumption of lemma 4.3, by that lemma and again theorem 4.11iii) in the opposite direction we find that $`\left|\right|_f\left|\right|>0`$.
Obviously we will construct the ITP Hilbert space from the linear span of $`C_0`$-vectors (we can ignore the $`C`$vectors which are not $`C_0`$-vectors by lemma 4.2). For this it will be useful to know how the Hilbert space decomposes into orthogonal subspaces. The following definition serves this purpose.
###### Definition 4.5
If $`_f`$ is a $`C_0`$-vector, we will call the sequence $`f=\{f_\alpha \}`$ a $`C_0`$-sequence. We will call two $`C_0`$-sequences $`f,g`$ strongly equivalent, denoted $`fg`$, provided that
$$\underset{\alpha }{}|<f_\alpha ,g_\alpha >_\alpha 1|$$
(4.5)
converges.
###### Lemma 4.4
Strong equivalence of $`C_0`$-sequences is an equivalence relation (reflexive, symmetric, transitive).
This lemma motivates the following definition.
###### Definition 4.6
The strong eqivalence class of a $`C_0`$ sequence $`f`$ will be denoted by $`[f]`$. The set of strong equivalence classes of $`C_0`$-sequences will be called $`𝒮`$.
The subsequent theorem justifies the notion of strong equivalence.
###### Theorem 4.2
i) If $`f^0[f][g]g^0`$ then $`<_{f^0},_{g^0}>=0`$.
ii) If $`f^0,g^0[f]`$ then $`<_{f^0},_{g^0}>=0`$ if and only if there exists $`\alpha `$ such that $`<f_\alpha ,g_\alpha >_\alpha =0`$.
So, $`C_0`$-vectors from different strong equivalence classes are always orthogonal and those from the same strong equivalence class are orthogonal if and only if they are orthogonal in at least one tensor product factor.
The following theorem gives two useful criteria for strong equivalence.
###### Theorem 4.3
i) $`[f]=[g]`$ if and only if $`_\alpha f_\alpha ^0g_\alpha ^0_\alpha ^2`$ and $`_\alpha |\mathrm{}(<f_\alpha ^0,g_\alpha ^0>_\alpha )|`$ converge for some $`f^0[f],g^0[g]`$.
ii) If $`f_\alpha =g_\alpha `$ for all but finitely many $`\alpha `$ then $`fg`$.
Obviously, it will be convenient to choose a representant $`f^0\left[f\right]`$ which is normalized in each tensor product factor. This is always possible.
###### Lemma 4.5
For each $`[f]𝒮`$ there exists $`f^0f`$ such that $`f_\alpha ^0_\alpha =1`$ for all $`\alpha `$.
The next lemma reveals that caution is due when trying to extend multilinearity from the finite to the infinite tensor product.
###### Lemma 4.6
Let $`_\alpha z_\alpha `$ be quasi-convergent. Then
i) If $`f`$ is a $`C`$-sequence, so is $`zf`$ with $`(zf)_\alpha :=z_\alpha f_\alpha `$.
ii) If moreover $`_\alpha ||z_\alpha |1|`$ converges and $`f`$ is a $`C_0`$sequence, so is $`zf`$.
iii) The product formula
$$_{zf}=\left[\underset{\alpha }{}z_\alpha \right]_f$$
(4.6)
fails to hold only if 1) $`_\alpha z_\alpha `$ is not convergent and 2) $`<_f,.>0`$ considered as a linear functional on $`C`$-vectors. In that case, $`\{z_\alpha \},f`$ satisfy the assumptions of ii), moreover $`z_\alpha 0\alpha `$.
iv) If $`\{z_\alpha \},f`$ satisfy the assumptions of ii) then $`[zf]=[f]`$ iff $`_\alpha |z_\alpha 1|`$ converges. If even $`z_\alpha 0\alpha `$, the latter condition implies convergence of $`_\alpha z_\alpha `$.
An important conclusion that we can draw from this lemma is the following. If (4.6) fails then, by iii), $`f,zf`$ are both $`C_0`$-sequences while $`_\alpha z_\alpha `$ is only quasi-convergent. Thus, both $`_f,_{zf}0`$ while $`_\alpha z_\alpha =0`$ by definition 4.2. Thus, $`\left[_\alpha z_\alpha \right]_f=0_{zf}`$.
Next, since, also by iii), $`z_\alpha 0\alpha `$ we have from collorary 4.1 that $`_\alpha \left|z_\alpha 1\right|`$ cannot be convergent as otherwise $`_\alpha z_\alpha `$ would be convergent which cannot be the case as $`_\alpha z_\alpha `$ is only quasi-convergent. Thus, by iv) $`f,zf`$ lie in different strong equivalence classes and therefore by theorem 4.2 $`<_f,_{zf}>=0`$.
###### Definition 4.7
By $`_C`$ we denote the completion of the complex vector space of finite linear combinations of elements from $`V_C`$ equipped with the sesquilinear form $`<.,.>`$ obtained by extending (4.3) from $`V_C`$ to $`_C`$ by sesquilinearity.
Notice that for $`C`$-vectors which are not $`C_0`$-vectors we have $`_f=0`$ as an element of $`_C`$.
###### Lemma 4.7
$`<\xi ,\xi >0\xi _C`$ and we define $`||\xi ||^2=<\xi ,\xi >`$. In particular, $`<.,.>`$ satisfies the Schwarz inequality and $`\xi =0`$ if and only if $`\xi =0`$.
We can now give the definition of the ITP.
###### Definition 4.8
We will denote by
$$^{}:=_\alpha _\alpha $$
(4.7)
the Cauchy-completion of the pre-Hilbert space $`_C`$. It is called the complete ITP of the $`_\alpha `$.
To analyze the structure $`^{}`$ in more detail, the strong equivalence classes provide the basic tool.
###### Definition 4.9
For a strong equivalence class $`[f]𝒮`$ we define the closed subspace of $`^{}`$
$$_{\left[f\right]}:=\overline{\{\underset{k=1}{\overset{N}{}}z_k_{f^k};z_k\text{ }\mathrm{C},f^k\left[f\right],N<\mathrm{}\}}$$
(4.8)
by the closure of the finite linear combinations of $`_f^{}`$’s with $`f^{}[f]`$. It is called the $`[f]`$-adic incomplete ITP of the $`_\alpha `$’s.
Notice that we could absorb the $`z_k`$ in (4.8) into one of the $`f_\alpha ^k`$. Now we have the fundamental theorem which splits $`^{}`$ into simpler pieces.
###### Theorem 4.4
The complete ITP decomposes as the direct sum over strong equivalence classes $`[f]`$ of $`[f]`$-adic ITP’s,
$$^{}=\overline{_{\left[f\right]𝒮}_{\left[f\right]}}$$
(4.9)
Also each $`\left[f\right]`$-adic ITP can be given a simple description.
###### Lemma 4.8
For a given $`[f]𝒮`$, fix any $`f^0[f]`$. By lemma 4.5 we can choose an $`f^0`$ with $`f_\alpha ^0_\alpha =1`$. Then $`_{[f]}`$ is the closure of the vector space of finite linear combinations of $`_f^{}`$’s where $`f^{}[f]`$ and $`f_\alpha ^{}=f_\alpha ^0`$ for all but finitely many $`\alpha `$.
It is easy to provide a complete orthonormal basis for an $`\left[f\right]`$-adic ITP if we know one in each $`_\alpha `$.
###### Lemma 4.9
Let $`f^0[f]𝒮,f_\alpha ^0_\alpha =1\alpha `$. Let $`d_\alpha =dim(_\alpha )`$ (takes the value of some higher Cantor aleph if $`_\alpha `$ is not separable). Let $`𝒥_\alpha ,\mathrm{\hspace{0.33em}0}𝒥_\alpha \alpha `$ be a set of indices of cardinality $`d_\alpha `$ and choose a complete orthonormal basis $`e_\alpha ^\beta ,\beta 𝒥_\alpha `$ such that $`e_\alpha ^0=f_\alpha ^0`$.
Consider the set $``$ of functions
$$\beta :\times _\alpha 𝒥_\alpha ;\alpha \left\{\beta \left(\alpha \right)\right\}_\alpha $$
(4.10)
such that 1) $`\beta (\alpha )𝒥_\alpha `$ and 2) $`\beta (\alpha )0`$ for finitely many $`\alpha `$ only. Let
$$_{e^\beta }:=_\alpha e_\alpha ^{\beta \left(\alpha \right)}$$
(4.11)
Then $`e^\beta [f]`$ and the set of $`C_0`$-vectors $`\{_{e^\beta };\beta \}`$ forms a complete orthonormal basis of $`_{[f]}`$, called a von Neumann basis.
The following corollary establishes that the $`\left[f\right]`$-adic ITP’s are mutually isomorphic.
###### Corollary 4.2
Each $`[f]`$-adic ITP is unitarily equivalent to the Hilbert space $`_{}=L_2(,d\nu _0)`$ of square summable functions on $``$, $`\widehat{\xi }:\text{ }\mathrm{C};\beta \widehat{\xi }(\beta )`$, where $`\nu _0`$ is the counting measure. The unitary map is given by
$$\widehat{U}_{\left[f\right]}:_{}_{\left[f\right]};\widehat{\xi }\underset{\beta }{}\widehat{\xi }\left(\beta \right)_{e^\beta }$$
(4.12)
The inverse map is given by
$$\widehat{U}_{\left[f\right]}^1:_{\left[f\right]}_{};\xi \widehat{\xi }\left(\beta \right)=<_{e^\beta },\xi >$$
(4.13)
In particular, since each $`\left[f\right]`$-adic ITP has a complete orthonormal basis labelled by $``$ and since the ITP is the direct sum of (the mutually isomorphic) $`\left[f\right]`$-adic ITP’s we have $`dim\left(^{}\right)=\left|\right|\left|𝒮\right|`$ where the appearing cardinalities will be aleph-valued in general (already in the simplest non-trivial case $`dim\left(_\alpha \right)=2,=\text{ }\mathrm{N}`$).
Notice that the index set $``$ is not required to have any ordering structure, thus we have identities of the form $`_\alpha f_\alpha =f_{\alpha _0}\left[_{\alpha \alpha _0}f_\alpha \right]`$, these are just different notations for the same object. However, is important to realize that the associative law generically does not hold for the ITP. By this we mean the following :
Let us decompose $``$ into mutually disjoint index sets $`_l`$ with $`l`$ then we can form the following Hilbert spaces : $`^{}=_\alpha _\alpha `$ and $`^{}:=_l_l`$ where $`_l:=_{\alpha _l}_\alpha `$. The $`C_0`$-vectors of $`^{}`$ are given by $`_f=_\alpha f_\alpha `$ while the $`C_0`$-vectors of $`^{}`$ are given by $`_f^{}=_lf_l^{}`$ where $`f_l^{}_l`$ is a (Cauchy limit of a) finite linear combination of vectors of the form $`_f^l=_{\alpha _l}f_\alpha `$. Inner products between $`C_0`$-vectors are computed as $`<_f,_g>=_\alpha <f_\alpha ,g_\alpha >_\alpha `$ and $`<_f^{},_g^{}>=_\alpha <f_l,g_l>_l`$ respectively where $`<_f^l,_g^l>_l=_{\alpha _l}<f_\alpha ,g_\alpha >_\alpha `$.
It is easy to see that if $`f=\left\{f_\alpha \right\}_\alpha `$ is a $`C_0`$-sequence for $`^{}`$ then $`f^{}=\left\{f_l^{}:=_f^l\right\}_l`$ is a $`C_0`$-squence for $`^{}`$. However, the obvious map between $`C_0`$-sequences given by
$$C:ff^{}$$
(4.14)
in general does not preserve the decomposition into strong equivalence classes of $`^{}`$ and $`^{}`$ respectively. We will give a few examples to illustrate this point.
* Let $`==\text{ }\mathrm{N}`$ and $`_l=\{2l1,2l\}`$ so that $`=_{l=1}^{\mathrm{}}_l`$. Consider the following two $`C_0`$-sequences : $`f_\alpha ,\alpha \text{ }\mathrm{N}`$ is just some normal vector in $`_\alpha `$, that is, $`f_\alpha _\alpha =1`$ and $`g_\alpha =f_\alpha `$. Then certainly their strong equivalence classes with respect to $`^{}`$ are different, $`\left[f\right]\left[g\right]`$ since $`|<f_\alpha ,g_\alpha >1|=2`$ so that (4.5) blows up. On the other hand we have $`f_l^{}=_f^l=f_{2l1}f_{2l}=[f_{2l1}][f_{2l}]=_g^l=g_l^{}`$. Thus, trivially $`\left[f^{}\right]^{}=\left[g^{}\right]^{}`$ where the prime at the bracket indicates that the class is with respect to $`^{}`$.
* Even multiplication by complex numbers is problematic : Take the same index sets as in i) and consider the complex numbers $`z_\alpha =1`$. Then $`_\alpha z_\alpha `$ is quasi-convergent but not convergent and therefore by definition $`_\alpha z_\alpha =0`$. Our map (4.14) now sends $`zf`$ to $`z^{}f^{}`$ with $`z_l^{}=z_{2l1}z_{2l}`$. Now $`z_l^{}=1`$ and thus $`_lz_l^{}`$ is convergent to $`1`$. It follows that $`_{zf}\left[_\alpha z_\alpha \right]_f=0`$ but $`_{z^{}f^{}}^{}=\left[_lz_l^{}\right]_f^{}^{}=_f^{}^{}`$, in particular, $`\left[f\right]\left[zf\right]`$ but $`\left[f^{}\right]^{}=\left[z^{}f^{}\right]^{}`$.
* Our map is certainly not invertible : Consider, for the same index sets as in i), the vector
$$f_l^{}:=\frac{1}{\sqrt{2}}\left[e_{2l1}^1e_{2l}^1+e_{2l1}^2e_{2l}^2\right]$$
(4.15)
where we assume that $`_\alpha `$ is at least two-dimensional and we choose two orthonormal vectors $`e_\alpha ^j,j=1,2`$ for each $`\alpha `$. Then $`f_l^{}_l=1`$ and $`f^{}`$ is a $`C_0`$-sequence for $`^{}`$. However, we cannot write $`f^{}`$ as a finite linear combination of $`C_0`$-vectors of $`^{}`$ : Any attempt to use the distributive law and to write it as a linear combination of $`C_0`$-vectors for $`^{}`$ of the form $`_l\left[e_{2l1}^{j_l}e_{2l}^{j_l}\right]`$ with $`j_l\{1,2\}`$ fails because all of these vectors are orthogonal (with respect to $`^{}`$) to $`_f^{}^{}`$ :
$$<_l[e_{2l1}^{j_l}e_{2l}^{j_l}],_lf_l^{}>=\underset{l}{}\frac{1}{\sqrt{2}}=0$$
(4.16)
It is plausible and one can indeed show that these complications do not arise if $`\left|\right|<\mathrm{}`$.
### 4.2 Von Neumann Algebras on the Infinite Tensor Product
The set of von Neumann algebras that one can define on the Infinite Tensor Product Hilbert space is of a surprisingly rich structure. In fact, every possible type of von Neumann’s factors (I, II<sub>1</sub>, II, III<sub>0</sub>, III<sub>1</sub>, III$`{}_{\lambda }{}^{};\lambda (0,1)`$) can be realized on the ITP. Physically, one will start from the local operators that “come from the various $`_\alpha `$”. However, there are many more operators which are not local and which are well-defined on the ITP. All the algebras that we consider are assumed to be unital.
###### Definition 4.10
We denote by $`(_\alpha )`$ the set of bounded operators on $`_\alpha `$ and by $`^{}:=(^{})`$ the set of bounded operators on the ITP $`^{}`$.
The restriction to bounded operators is not a severe one since every unbounded operator can be written (up to domain questions) as a linear combination of self-adjoint ones and those are known if we know their spectral projections which are bounded operators.
An operator on one of the tensor product factors is not a priori defined on the ITP. The following lemma embeds $`\left(_\alpha \right)`$ into $`^{}`$.
###### Lemma 4.10
Let $`\alpha _0`$ and $`A_{\alpha _0}(_{\alpha _0})`$. Then there exists a unique operator $`\widehat{A}_{\alpha _0}^{}`$ such that for any $`C`$-sequence $`f`$
$$\widehat{A}_{\alpha _0}_f=_f^{}\text{ where }f_\alpha ^{}=\{\begin{array}{cc}f_\alpha & :\alpha \alpha _0\\ A_{\alpha _0}f_{\alpha _0}& :\alpha =\alpha _0\end{array}$$
(4.17)
We will use the notation
$$\widehat{A}_{\alpha _0}_f=\left[A_{\alpha _0}f_{\alpha _0}\right]\left[_{\alpha \alpha _0}f_\alpha \right]$$
(4.18)
This lemma gives rise to the following definition.
###### Definition 4.11
We denote by $`_\alpha `$ the extension of $`(_\alpha )`$ to the ITP, that is,
$$_\alpha =\left\{\widehat{A}_\alpha ;A_\alpha \left(_\alpha \right)\right\}$$
(4.19)
Obviously $`_\alpha ^{}`$.
It is not difficult to prove that $`A_\alpha \widehat{A}_\alpha `$ is in fact a algebra isomorphism. The algebras $`^{},\left(_\alpha \right)`$ are C algebras by definition. Recall that, on the other hand, a von Neumann algebra over a Hilbert space is a weakly (equivalently strongly) closed sub- algebra of the algebra of bounded operators on that Hilbert space.
###### Lemma 4.11
For all $`\alpha `$, the algebra $`_\alpha `$ is a von Neumann algebra (v.N.a.) over $`^{}`$.
The idea of proof is quite simple : One writes $`^{}=\left(_\alpha _{\overline{\alpha }}\right)`$ where $`\overline{\alpha }=\alpha `$. Next, it is almost obvious that $`_\alpha `$ coincides with $`_{\overline{\alpha }}^{}=\left\{\widehat{B}^{};[\widehat{A},\widehat{B}]=0\widehat{A}_{\overline{\alpha }}\right\}`$, the commutant of $`_{\overline{\alpha }}`$. Then an appeal to the bicommutant (or von Neumann density) theorem finishes the proof.
Actually, the correspondence of lemma 4.10 extends to von Neumann algebras $`\left(_\alpha \right)\left(_\alpha \right)`$ as we state in the subsequent theorem.
###### Theorem 4.5
The one to one correspondence $`(_\alpha )A_\alpha \widehat{A}_\alpha _\alpha `$ extends to a isomorphism between von Neumann algebras $`(_\alpha )(_\alpha )_\alpha =\{\widehat{A}_\alpha ;A_\alpha (_\alpha )\}`$
The largest von Neumann algebra on $`^{}`$ that we can construct from the algebras $`\left(_\alpha \right)`$ is the following one.
###### Definition 4.12
By $`^{}`$ we denote the smallest v.N.a. that contains all the $`_\alpha `$, that is, the weak closure of the set
$$_\alpha _\alpha $$
(4.20)
It turns out that not surprisingly $`^{}`$ is a proper subalgebra of $`^{}`$ unless $`\left|\right|<\mathrm{}`$. Physically, the indices $`\alpha `$ label local degrees of freedom and therefore the elements of $`_\alpha `$ correspond to local operators of a quantum field theory. Thus the algebra $`^{}`$ is the algebra of local observables represented on the ITP $`^{}`$. The remainder $`^{}^{}`$ can therefore be identified with a set of non-local operators. Thus, while the algebra $`^{}`$ is rather important from the point of view of local (or algebraic) quantum field theory it is the remainder which offers challenging possibilities in the sense that it could be the universal home for operators that map a given physical system to a drastically different one. Examples for this could be the change of energy by an infinite amount or topology change of the underlying spacetime manifold. We will come back to this point in section 5. These issues should be particularly important for quantum general relativity since there all the (Dirac) observables are supposed to be non-local.
In any case we should investigate the subalgebra $`^{}`$ in more detail. To that end, recall from lemma 4.6 that the equation $`_{zf}=\left[_\alpha z_\alpha \right]_f`$ is false only if both $`f,zf`$ are $`C_0`$-vectors, $`z_\alpha 0\alpha `$ but $`_\alpha z_\alpha `$ is only quasi-convergent. This fact gives rise to the next definition.
###### Definition 4.13
Two $`C_0`$-sequences $`f,g`$ are said to be weakly equivalent, denoted by $`fg`$, provided that there are complex numbers $`z_\alpha `$ such that $`zf`$ and $`g`$ are strongly equivalent, that is, $`zfg`$.
Important facts about weak equivalence are contained in the following lemma which also contains a necessary and suffient criterion.
###### Lemma 4.12
i) Definition 4.13 remains unchanged if we restrict to complex numbers with $`|z_\alpha |=1`$.
ii) Weak equivalence is an equivalence relation (reflexive, symmetric, transitive).
iii) $`fg`$ if and only if
$$\underset{\alpha }{}\left|\right|<f_\alpha ,g_\alpha >_\alpha |1|$$
(4.21)
converges.
Comparing with definition 4.5 we see that the “only” difference between strong and weak equivalence is the additional modulus for $`<f_\alpha ,g_\alpha >_\alpha `$ in (4.21).
###### Definition 4.14
i) For a $`C_0`$-sequence $`f`$ its weak equivalence class is denoted by $`(f)`$. The set of weak equivalence classes is denoted by $`𝒲`$.
ii) For given $`(f)𝒲`$ we denote by $`_{(f)}`$ the closure of the set of finite linear combinations of $`_f^{}`$’s where $`f^{}(f)`$.
Obviously, weak equivalence is weaker than strong equivalence. Thus, each $`\left(f\right)𝒲`$ decomposes into mutually disjoint $`\left[f^{}\right]𝒮,f^{}\left(f^{}\right)`$. It follows from this and the mutual orthogonality of the $`_{\left[f^{}\right]}`$’s (theorem 4.2) that we may write
$$_{\left(f\right)}=\overline{_{\left[f^{}\right]𝒮\left(f\right)}_{\left[f^{}\right]}}$$
(4.22)
###### Lemma 4.13
i) For every sequence of complex numbers $`\{z_\alpha \}_\alpha `$ such that $`|z_\alpha |=1\alpha `$ there exists a unique, unitary operator $`\widehat{U}_z`$, densely defined on (finite linear combinations of) $`C_0`$-vectors $`f`$ such that $`\widehat{U}_z_f=_{zf}`$.
ii) Given $`s𝒮,w𝒲`$ respectively, denote by $`\widehat{P}_s,\widehat{P}_w`$ respectively the projection operators from $`^{}`$ onto the closed subspaces $`_s,_w`$ respectively. Then :
a) $`[\widehat{U}_z,\widehat{P}_w]=0`$,
b) $`[\widehat{U}_z,\widehat{P}_s]=0`$ if and only if $`_\alpha z_\alpha `$ converges to $`z,|z|=1`$ in which case $`\widehat{U}_z=z\text{1}_{^{}}`$ and
c) if $`[\widehat{U}_z,\widehat{P}_s]0`$ then $`\widehat{U}_z_s=_s^{}`$ where $`ss^{}𝒮`$, that is, $`\widehat{U}_z`$ maps different $`s`$-adic ITP subspaces onto each other which are thus unitarily equivalent.
The following theorem describes much of the structure of $`^{}`$.
###### Theorem 4.6
i) An operator $`\widehat{A}^{}`$ belongs actually to $`^{}`$ if and only if it commutes with all the $`\widehat{U}_z,\widehat{P}_s`$ of lemma 4.13. In particular, the elements of $`^{}`$ leave all the $`_s,s𝒮`$ invariant.
ii) For each $`w𝒲`$, fix once and for all an element $`s_w𝒮w`$. Suppose that we are given a family of bounded operators $`\widehat{A}_w`$ on $`_{s_w}`$ for each $`w𝒲`$. Then there exists an operator $`\widehat{A}^{}`$ such that its restriction $`\widehat{A}_{s_w}`$ to $`_{s_w}`$ coincides with $`\widehat{A}_w`$, provided that the set of non-negative numbers $`\{\widehat{A}_w;w𝒲\}`$ is bounded. In that case, $`\widehat{A}_w`$ is actually unique.
iii) The norm of the operator $`\widehat{A}`$ of ii) is given by
$$\widehat{A}=sup\left\{\widehat{A}_s;s𝒮\right\}=sup\left\{\widehat{A}_w;w𝒲\right\}$$
(4.23)
This theorem tells us the following about $`^{}`$:
1)
As $`\widehat{P}_w=_{s𝒮w}\widehat{P}_s`$, item i) reveals that each $`_w,w𝒲`$ is an invariant subspace for any element $`\widehat{A}^{}`$, it is “block diagonal” with respect to $`^{}`$ where the blocks correspond to the $`_w,w𝒲`$. Within each of these blocks, $`\widehat{A}`$ is further reduced by each $`_s,s𝒮w`$. Moreover, since $`\widehat{U}_z`$ commutes with $`\widehat{A}`$ and we obtain any $`_s,s𝒮w`$ by mapping $`_{s_w}`$ of theorem 4.6 with $`\widehat{U}_z`$, knowlegde of $`\widehat{A}`$ on $`_{s_w}`$ is sufficient to determine it all over $`_w`$. 2)
Item ii) tells us that certainly not every element of $`^{}`$ lies in $`^{}`$, actually it is easy to construct bounded operators, e.g. the $`\widehat{U}_z^{}`$, which do not lie in $`^{}`$.
Finally we determine the cardinality of the set $`𝒮w`$.
###### Lemma 4.14
i) If $`||<\mathrm{}=\mathrm{}`$ then $`𝒮=𝒲,|𝒮|=1`$ and $`^{}=_w=_s`$.
ii) If $`||\mathrm{}`$ then $`|𝒮w|=2^{||}`$.
iii) If the number of $`\alpha `$’s such that $`dim(_\alpha )2`$ is finite, then $`|𝒲|=1`$. Otherwise, $`|𝒲|2^{\mathrm{}}`$.
To investigate the structure of $`^{}`$ further we need to recall some of the notions from the theory of von Neumann algebras, e.g. .
###### Definition 4.15
i)
Let $`()`$ be a v.N.a. over the Hilbert space $``$. The commutant of $``$, denoted by $`^{}`$ is the set of operators in $`()`$ that commute with all elements of $``$. For a v.N.a. we have $`^{\prime \prime }=`$. $`𝒵()=^{}`$ is called the center of $``$. The v.N.a. is called a factor if $`𝒵()=\{\lambda \text{1}_{};\lambda \text{ }\mathrm{C}\}`$, that is, the center consists only of the scalars.
ii)
Let $`\widehat{P},\widehat{Q}`$ be projections. We say that
a) $`\widehat{Q}`$ is a subprojection of $`\widehat{P}`$, denoted $`\widehat{Q}\widehat{P}`$, iff $`\widehat{Q}\widehat{P}`$.
b) $`\widehat{Q},\widehat{P}`$ are equivalent, denoted $`\widehat{Q}\widehat{P}`$, iff there exists a partial isometry with initial subspace $`\widehat{P}`$ and final subspace $`\widehat{Q}`$.
c) $`\widehat{P}0`$ is a minimal projection if there is no proper subprojection $`\widehat{Q}0`$ of $`\widehat{P}`$.
d) $`\widehat{P}0`$ is an infinite projection if there is a proper subprojection $`\widehat{Q}0`$ of $`\widehat{P}`$ to which it is equivalent.
iii)
Let $``$ be a factor. Then we call $``$ of type
I : if $``$ contains a minimal projection. If 1<sub>H</sub> is an infinite projection, then the type is $`I_{\mathrm{}}`$ otherwise it is $`I_n`$ where $`n=dim()`$.
III : every non-zero projection of $``$ is infinite. A further systematic classification of type III factors is due to Connes, see e.g. and references therein. One distinguishes between type III<sub>0</sub> (the Krieger factor ), type III$`{}_{\lambda }{}^{},\lambda (0,1)`$ (the Powers factor ) and III<sub>1</sub> (the factor of Araki and Woods ).
II : if R is neither of type I nor of type III. If 1<sub>H</sub> is an infinite projection, then $``$ is called type II otherwise type II<sub>1</sub>.
One can show that factors of type I are isomorphic to algebras of bounded operators on some Hilbert space. Factors of type II are generated by operators of the form $`A_1\text{1}__2,\text{1}__1A_2`$ acting on the Hilbert space $`_1_2`$ where $`A_1`$ belongs to a factor of type I over $`_1`$ and $`A_2`$ to one of type II<sub>1</sub> over $`_2`$. For factors of type I and II it is possible to introduce a dimension function for projections, that is, a positive definite function $`dim\left(\widehat{P}\right)0`$ vanishing only if $`\widehat{P}=0`$, uniquely determined by the two properties that
1) $`dim\left(\widehat{P}+\widehat{Q}\right)=dim\left(\widehat{P}\right)+dim\left(\widehat{Q}\right)`$ if $`\widehat{P}\widehat{Q}`$ and
2) $`dim\left(\widehat{P}\right)=dim\left(\widehat{Q}\right)`$ if $`\widehat{P}\widehat{Q}`$.
The range of that function is $`0,1,2,..,n`$ for type I<sub>n</sub>, $`0,1,2,..,\mathrm{}`$ for type I, $`[0,1]`$ for type II<sub>1</sub> and $`[0,\mathrm{}]`$ for type II. For type III a dimension function can be introduced but it takes only the values $`0,\mathrm{}`$ and therefore cannot be used to obtain the finer subdivision of type III factors outlined above for which the use of modular (or Tomita-Takesaki) theory and the Connes invariant is necessary (a self-contained exposition aimed at mathematical phyicists can be found in ).
We close this section by mentioning that the more unfamiliar factors of type II and III are not only of academic interest. In fact, they appear already in systems as simple as the infinite spin chain (see the second reference of ). If one represents the abstract CCR $`C^{}`$algebra of spin $`1/2`$ operators $`\widehat{\sigma }_l^j;j=1,2,3;l=1,2,..`$, via the GNS theorem , for a state $`\omega _s`$ ($`s[0,1]`$) for which we get GNS data $`(\mathrm{\Omega }_s^{\mathrm{}},_s^{\mathrm{}},\pi _s^{\mathrm{}})`$ where the cyclic vector is
$$\mathrm{\Omega }_s^{\mathrm{}}=_{l=1}^{\mathrm{}}\mathrm{\Omega }_s,\mathrm{\Omega }_s=[\sqrt{\frac{1+s}{2}}e_1e_1+\sqrt{\frac{1s}{2}}e_2e_2],$$
the Hilbert space is the ITP $`_s^{\mathrm{}}=_{l=1}^{\mathrm{}}[\text{ }\mathrm{C}^2\text{ }\mathrm{C}^2]`$ corresponding to the index set $``$ of pairs $`\alpha =(l,\tau ),\tau =1,2`$, and the representation is $`\pi _s^{\mathrm{}}\left(\widehat{\sigma }_l^j\right)`$ acting only on the Hilbert $`\text{ }\mathrm{C}^2`$ (with standard orthonormal basis $`e_1,e_2`$) corresponding to $`\alpha =(l,1)`$, then upon weak closure a factor of type I or II<sub>1</sub> or III<sub>s</sub> results for $`s=1`$ or $`s=0`$ or $`s(0,1)`$. The physical interpretation of the parameter $`s`$ is that $`\mathrm{\Omega }_s`$ is the GNS datum for the mixed state
$$\omega _s\left(A\right)=\frac{\text{tr}\left(Ae^{\beta \sigma ^3}\right)}{\text{tr}\left(e^{\beta \sigma ^3}\right)}$$
with $`s=\text{th}\left(\beta \right)`$ on $`_s=\text{ }\mathrm{C}^4=\text{ }\mathrm{C}^2\text{ }\mathrm{C}^2`$, thus type I, II<sub>1</sub> and III<sub>s</sub> respectively means zero, infinite or finite temperture respectively.
Finally, the type of local algebras $`\left(𝒪\right)`$ appearing in algebraic quantum field theory is the unique hyperfinite factor of type III<sub>1</sub> for diamond regions $`𝒪`$ (intersections of past and future light cones in the obvious way; this result can be extended to arbitrary $`𝒪`$ in case that the theory has a scaling limit (short distance conformal invariance) ). Here a v.N.a. is said to be hyperfinite if it is the inductive limit of finite dimensional algebras. This brings to the next topic.
### 4.3 Inductive Limits of Hilbert Spaces and von Neuman Algebras
For the applications that we have in mind, specifically quantum gravity and quantum gauge theory coupled to gravity, the framework of the Infinite Tensor Product is not general enough for the following reason. Recall from section 2 that the degrees of freedom of these field theories are labelled by graphs. Moreover, given a graph $`\gamma `$ the degrees of freedom associated with it are labelled by the edges of that graph. Thus, it seems that we are in position to apply the theory outlined in sections 4.1 and 4.2 by choosing $`=E\left(\gamma \right)`$. While that is indeed true for the given graph $`\gamma `$, rather than working with a fixed, infinite graph $`\gamma `$ we are working with all of them because we do not have a lattice gauge field theory but a continuum one. So we actually get an uncountably infinite family of ITP’s. That would not pose any problems if we could treat each of them independently, however, this is not the case, e.g. not if a graph is contained in a bigger one. The inductive limit construction is well suited to handle this problem.
###### Definition 4.16
0)
Let $``$ be a partial order (that is, a reflexive, antisymmetric and transitive relation) on the index set $``$. The index set is said to be directed if for any $`l,l^{}`$ there exists $`l^{\prime \prime }`$ such that $`ll^{\prime \prime }`$ and $`l^{}l^{\prime \prime }`$.
1i)
Let $`\{_l\}_l`$ be a family of $`C^{}`$ algebras (v.N.a.’s) labelled by a directed index set $``$. Suppose that for all $`l,l^{}`$ with $`ll^{}`$ there is a monomorphism (injective homomorphism) $`F_{ll^{}}:_l_l^{}`$ satisfying
1) $`F_{ll^{}}(\text{1}__l)=\text{1}__l^{}`$ and
2) $`F_{ll^{}}F_{l^{}l^{\prime \prime }}=F_{ll^{\prime \prime }}`$ for any $`ll^{}l^{\prime \prime }`$.
Then the pair of families $`\{_l,F_{ll^{}}\}`$ is called a directed system of $`C^{}`$ algebras (v.N.a.’s).
1ii)
Let $`\{_l\}_l`$ be a family of $`C^{}`$ algebras (v.N.a.’s) labelled by a directed index set $``$. A $`C^{}`$ algebra (v.N.a.) $``$ is said to be the $`C^{}`$ ($`W^{}`$) inductive limit of the $`_l`$ provided there exist monomorphisms $`F_l:_l`$ such that
1) $`F_l(\text{1}__l)=\text{1}_{}`$ and
2) $`_lF_l(_l)`$ is uniformly (weakly) dense in $``$.
2i)
Let $`\{_l\}_l`$ be a family of Hilbert spaces labelled by a directed index set $``$. Suppose that for all $`l,l^{}`$ with $`ll^{}`$ there is an isometric monomorphism $`\widehat{U}_{ll^{}}:_l_l^{}`$ such that $`\widehat{U}_{ll^{}}\widehat{U}_{l^{}l^{\prime \prime }}=\widehat{U}_{ll^{\prime \prime }}`$ for any $`ll^{}l^{\prime \prime }`$.
Then the pair of families $`\{_l,\widehat{U}_{ll^{}}\}`$ is called a directed system of Hilbert spaces.
2ii)
Let $`\{_l\}_l`$ be a family of Hilbert spaces labelled by a directed index set $``$. A Hilbert space $``$ is said to be the inductive limit of the $`_l`$ provided there exist isometric monomorphisms $`\widehat{U}_l:_l`$ such that $`_l\widehat{U}_l_l`$ is dense in $``$.
3i)
Given a directed system of Hilbert spaces $`_l`$, suppose that we are given a family of operators $`\widehat{A}_l_l(_l)`$ such that
1) $`sup\{\widehat{A}_l_l;l\}<\mathrm{}`$ and
2) there exists $`l_0`$ so that $`\widehat{U}_{ll^{}}\widehat{A}_l=\widehat{A}_l^{}\widehat{U}_{ll^{}}`$ for any $`l_0ll^{}`$.
Then the family is called a directed system of operators.
3ii)
Given an inductive limit $``$ of Hilbert spaces $`_l`$ together with a family of operators $`\widehat{A}_l_l(_l)`$, an operator $`\widehat{A}()`$ is called the inductive limit of the $`\widehat{A}_l`$ provided that there exists $`l_0`$ so that $`\widehat{U}_l\widehat{A}_l=\widehat{A}\widehat{U}_l`$ for any $`l_0l`$.
The connection between ni) and nii), $`n=1,2,3`$ is made through the following theorem.
###### Theorem 4.7
1)
Given a directed system of $`C^{}`$ algebras (v.N.a.’s) $`\{_l,F_{ll^{}}\}`$ there exists a unique (up to isomorphisms) $`C^{}`$ ($`W^{}`$) inductive limit $``$ of the $`_l`$ where the corresponding monomorphisms $`F_l`$ satisfy the compatibility condition $`F_l^{}F_{ll^{}}=F_l`$.
2)
Given a directed system of Hilbert spaces $`\{_l,\widehat{U}_{ll^{}}\}`$ there exists a unique (up to unitarity) inductive limit $``$ of the $`_l`$ where the corresponding isometric monomorphisms $`\widehat{U}_l`$ satisfy the compatibility condition $`\widehat{U}_l^{}\widehat{U}_{ll^{}}=\widehat{U}_l`$.
3)
Given a directed system of operators $`\{\widehat{A}_l\}`$ on a directed system of Hilbert spaces $`_l`$, there exists a unique (up to unitarity) inductive limit operator $`\widehat{A}`$ on the inductive limit Hilbert space $``$.
The proof of this theorem can be found in the second volume of the first reference of . Notice that inductive and projective limits (as used, e.g. in ) are essentially identical, just that the projective limit employs “projections downwards” a chain in the directed system while the inductive limit employs “embeddings upwards” the chain.
The importance of the inductive limit for our purposes lies in the following construction. Suppose we are given an index set $``$ and consider the set $``$ of all possible subsets of $``$ (notice that we allow the cardinality of $`l`$ to be infinite). Then $``$ is a directed set where the partial order $``$ is given by the inclusion relation $``$. For each $`l`$ we can form the Infinite Tensor Product $`_l^{}`$ of the $`_\alpha ,\alpha l`$ and the corresponding von Neumann algebra $`_l^{}`$. Moreover, we have for $`ll^{}`$ the obvious monomorphism $`F_{ll^{}}`$ assigning to $`\widehat{A}_l_l`$ the operator $`\widehat{A}_l\left[_{\alpha l^{}l}\text{1}__\alpha \right]_l^{}`$. Finally, choose for each $`\alpha `$ a fixed standard unit vector $`\mathrm{\Omega }_\alpha _\alpha `$, then for $`ll^{}`$ we have isometric monomorphisms $`\widehat{U}_{ll^{}}`$ mapping $`\xi _l_l`$ to $`\xi _l\left[_{\alpha l^{}l}\mathrm{\Omega }_\alpha \right]`$. It is easy to see that $`F_{ll^{}},\widehat{U}_{ll^{}}`$ satisfy the requirements of definition 4.14 and so we can form the inductive limit von Neumann algebra $`_{\mathrm{}}^{}`$ and inductive limit Hilbert space $`_{\mathrm{}}^{}`$ respectively which are the universal objects from which our various “lattice” algebras $`_l`$ and Hilbert spaces $`_l`$ respectively can be obtained by theorem 4.7.
## 5 Infinite Tensor Products and Continuum Quantum Gauge Theories
We will now apply the machinery of section 4 to quantum gauge field theories on globally hyperbolic, spatially non-compact manifolds along the lines suggested by the exposition of section 2 and make contact with the semi-classical analysis machinery in connection with the coherent states as outlined in section 3. We proceed in several steps.
### 5.1 Kinematical Framework
In this subsection we carefully carry over the Ashtekar-Isham-Lewandowski kinematical framework developed for the finite analytical category to the infinite analytical one.
#### 5.1.1 Properties of Infinite Graphs
Notice that in order that $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$ has an infinite number of edges, $`\mathrm{\Sigma }`$ must not be compact by the very definition of compactness.
Next, while $`\mathrm{\Gamma }_\sigma ^\omega `$ contains graphs with an infinite number of edges, the number of these edges is at most countably infinite if $`\mathrm{\Sigma }`$ is paracompact as we assume here as otherwise integration theory cannot be employed. To see this, notice that since a finite dimensional manifold $`\mathrm{\Sigma }`$ is locally compact we can apply the theorem in chapter I, paragraph 9 which says that a (connected) locally compact space is paracompact if and only if it is the countable union of compact sets. Assume now that $`\gamma `$ has an uncountably infinite number of edges and let $`U_n,n=1,2,..`$ be a countable compact cover of $`\mathrm{\Sigma }`$. We conclude that at least one of the $`U_n`$ must contain an uncountably infinite number of edges of $`\gamma `$ because $`\gamma `$ has an uncountable number of edges and the countable union of countable sets is countable. But this cannot happen if $`\gamma `$ is piecewise analytic and $`\sigma `$-finite by definition.
We conclude that each element $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$ is of a rather controllable form with at most a countable number of edges and vertices and no accumulation points as it would happen for webs. They thus resemble maximally the lattices that one is used to from lattice gauge theory and this is the form of graphs which are clearly most suitable for semiclassical analysis and the continuum limit. (The typical element of $`\mathrm{\Gamma }_0^{\mathrm{}}`$ has at least one accumulation point of vertices and on such graphs one will certainly not approximate actions, Hamiltonians and the like).
Moreover, we have the following basic lemma and this is where analyticity comes in.
###### Lemma 5.1
The set $`\mathrm{\Gamma }_\sigma ^\omega `$ is a directed set by inclusion.
Proof of Lemma 5.1 :
Notice that if $`\gamma ,\gamma ^{}`$ are two piecewise analytical, $`\sigma `$-finite graphs then $`\gamma ^{\prime \prime }:=\gamma \gamma ^{}`$ is also piecewise analytic. We claim that it is also $`\sigma `$-finite. Suppose this was not the case. Then, either a) there exists a compact subset $`U\mathrm{\Sigma }`$ such that $`\gamma ^{\prime \prime }U`$ is an infinite graph or b) there exists a compact cover $`𝒰`$ such that the set $`\left\{\left|E\left(\gamma ^{\prime \prime }U\right)\right|;U𝒰\right\}`$ is unbounded.
As for case a), we know that $`\gamma U,\gamma ^{}U`$ are both finite graphs with finite number of edges $`e,e^{}`$ respectively. Since $`\gamma ^{\prime \prime }U=\left[\gamma U\right]\left[\gamma ^{}U\right]`$ the only way that $`\gamma ^{\prime \prime }U`$ can possibly be infinite is that there is at least one edge $`e`$ of a $`\gamma `$ and one edge $`e^{}`$ of $`\gamma ^{}`$ such that $`ee^{}`$ intersect each other in an infinite number of isolated points. (The possibility that they overlap in a finite segment is excluded by analyticity as they would be analytic extensions of each other in this case and thus would make up a single analytical curve). However, two analytical curves that coincide in an infinite number of points are analytical extensions of each other. Thus, case a) cannot occur.
As for case b), we find compact sets $`U_n`$ labelled by natural numbers $`n`$ such that $`\gamma ^{\prime \prime }U_n`$ has at least $`n`$ edges. However, we know that there are natural numbers $`N,N^{}`$ such that $`\left|E\left(\gamma U_n\right)\right|<N,\left|E\left(\gamma ^{}U_n\right)\right|<N^{}`$ for all $`n`$. It follows that $`U_{\mathrm{}}`$ has the property of case a) which we excluded already. Thus, case b) can also not occur.
$`\mathrm{}`$
That $`\mathrm{\Gamma }_\sigma ^\omega `$ is a directed set is of paramount importance for inductive limit constructions.
#### 5.1.2 Quantum Configuration Space
Recall that in section 2.2 the quantum configuration space $`\overline{𝒜}`$ arose as the Gel’fand spectrum of the Abelian $`C^{}`$ algebra generated by finite linear combinations of functions of smooth connections, restricted to finite graphs (cylindrical functions) and completed in the supremum norm. It is natural to ask whether we can extend this construction to functions of smooth connections restricted to infinite graphs and to see if the size of the quantum configuration space is changed. The following simple example reveals that a naive transcription of this method is problematic :
Take $`\mathrm{\Sigma }=\text{ }\mathrm{R}^3,G=SU\left(2\right)`$ and let $`\gamma `$ be the $`x`$-axis split into the countably infinite number of intervals $`e`$ of equal unit length. Thus, $`\gamma `$ is a piecewise analytic, $`\sigma `$-finite graph. Let us consider the following function of smooth connections
$$Af\left(A\right):=\underset{e}{}\left[k\chi _j\left(h_e\left(A\right)\right)\right]$$
(5.1)
where $`k`$ is a constant, $`\chi _j\left(h\right)=\text{tr}\left(\pi _j\left(h\right)\right)`$ is the character of the spin $`j`$ representation and the convergence of (5.1) is meant in the sense of definition 4.1. By definition, the sup-norm of that function is $`f=sup_{A𝒜}_e\left|k\chi _j\left(h_e\left(A\right)\right)\right|`$. Now the zero connection is certainly an element of $`𝒜`$, so $`fsup_{A𝒜}_e\left|k\left(2j+1\right)\right|`$ and this infinite product converges to $`0`$ if $`\left|k\right|<1/\left(2j+1\right)`$, to $`1`$ if $`\left|k\right|=1/\left(2j+1\right)`$ and diverges otherwise. Now it is easy to see that for any $`hSU\left(2\right)`$ we have in fact $`\left|\chi _j\left(h\right)\right|2j+1`$ and equality is reached for $`h=1`$ so that indeed $`f=sup_{A𝒜}_e\left|k\left(2j+1\right)\right|`$. It follows that in the only case that the norm is finite and non-vanishing, we have that $`f\left(A\right)`$ is non-vanishing iff $`_e\left|k\chi _j\left(h_e\left(A\right)\right)1\right|`$ converges which means that for each $`ϵ>0`$ the set of $`e^{}s`$ such that $`\left|k\chi _j\left(h_e\left(A\right)\right)1\right|ϵ`$ is finite. In other words, $`f\left(A\right)`$ is almost given by $`_e\delta _{h_e\left(A\right),1}`$, an infinite product over Kronecker $`\delta `$’s rather than $`\delta `$-distributions and it is almost granted that its support is of measure zero for every reasonable measure even if we extend $`𝒜`$ to $`\overline{𝒜}`$. We will prove shortly that this is indeed the case with respect to the Ashtekar-Lewandowski measure which turns out to be extendible to our context. Thus, we face the problem that the (Gel’fand transforms of) functions of finite sup-norm (and thus all the elements of the Abelian $`C^{}`$ algebra) seem to be supported on measure zero subsets of interesting measures on the resulting spectrum.
On the other hand, it is physically plausible that the quantum configuration space as obtained from $`\mathrm{\Gamma }_0^\omega `$ should not change when we extend to $`\mathrm{\Gamma }_\sigma ^\omega `$. The reason is that, by the very definition of $`\sigma `$-finiteness, if we consider a function depending on the infinite number of degrees of freedom labelled by the edges of $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$ but restrict its dependence to a finite number of degrees of freedom by “freezing” all degrees of freedom labelled by $`\gamma \left[\gamma U\right]`$ for any compact set $`U\mathrm{\Sigma }`$ then we get a function cylindrical over $`\gamma U\mathrm{\Gamma }_0^\omega `$ whose behaviour is certainly not different from the ones considered in section 2.2. In other words, functions over $`\gamma `$ satisfy a locality property.
Thus, rather than deriving the spectrum $`\overline{𝒜}`$ as the Gel’fand spectrum arising from an Abelian C algebra of cylindrical functions over truly infinite graphs it is the characterization of the Ashtekar-Isham spectrum derived in for finite graphs which we simply extend to the infinite category ! This works as follows :
We need the set $`W_0^\omega `$ that one obtains as the union of a finite number of, not necessarily compactly supported anlytical paths. Since analytical paths of non-compact range can intersect each other in an infinite number of isolated points and since generic elements of $`\mathrm{\Gamma }_\sigma ^\omega `$ cannot be otained as the union of a finite number of paths we see that we have the proper inclusions $`\mathrm{\Gamma }_0^\omega W_0^\omega \mathrm{\Gamma }_\sigma ^\omega `$. The set $`W_0^\omega `$ is trivially directed by inclusion, in a sense it is very similar to the set $`\mathrm{\Gamma }_0^{\mathrm{}}`$. In fact, if one would blow up the neighbourhood of the source of a tassel by an infinite amount then one gets, apart from the difference between smooth and analytic paths, precisely the kind of objects that lie in $`W_0^\omega `$. For this reason, we will call them analytical webs. Notice that in contrast to smooth webs the paths that determine an analytical web are obviously holonomically independent because 1) they cannot overlap each other in a finite segment due to analyticity, they can only intersect each other in a possibly infinite number of isolated points and 2) because they always have a (non-overlapped) segment in the bulk of $`\mathrm{\Sigma }`$ where no fall-off conditions on $`A`$ restrict the range of the holonomy along that segment.
Let now $`𝒜`$ be the classical configuration space of section 2.1 where appropriate fall-off conditions at spatial infiniy are obeyed. Then the holonomy of $`A𝒜`$ along an analytic path of infinite range is in fact well-defined precisely due to the fall-off conditions on $`A`$ at spatial infinity. As in the context of $`\mathrm{\Gamma }_0^\omega `$ we can now consider the algebra of cylindrical functions of $`A𝒜`$ which are simply finite linear combinations of functions of the form $`f\left(A\right)=f_w\left(\left\{h_e\left(A\right)\right\}_{ew}\right)`$ where $`wW_0^\omega `$ denotes the analytical web and $`f_w`$ is a complex valued function on $`G^{\left|w\right|},\left|w\right|`$ the number of paths that determine $`w`$. Now the complications that we observed above in the context of cylindrical functions over $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$ are out of the way because the cylindrical functions for webs depend on a finite number of arguments only. We can therefore complete the algebra in the sup-norm just as in section 2.2, obtain a $`C^{}`$ algebra and can follow exactly the same steps reviewed there for $`\mathrm{\Gamma }_0^\omega `$ to arrive at the Ashtekar-Isham spectrum $`\overline{𝒜}`$ as the Gel-fand spectrum of that algebra. Finally, by following exactly the same proofs as in we find $`\overline{𝒜}`$ to be in one to one correspondence with the set of all homomorphisms from the groupoid $`𝒳`$ of (composable) analytic paths in $`\mathrm{\Sigma }`$ into the gauge group $`G`$. The isomorphism is the same as the one from , that is,
$$\overline{𝒜}\overline{A}H_{\overline{A}}\text{Hom}(𝒳,G);\left(H_{\overline{A}}\left(e\right)\right)_{mn}:=\overline{A}\left(\left(h_e\right)_{mn}\right)=\left(\widehat{h}_e\right)_{mn}\left(\overline{A}\right)$$
(5.2)
Here $`m,n`$ are the indices of the matrix elements of the defining represenation of $`G`$ and $``$ denotes the Gel’fand transform. Notice that in contrast to it was not necessary to consider the one point compactification $`\widehat{\mathrm{\Sigma }}`$ of $`\mathrm{\Sigma }`$. In fact, we refrained from doing that because we now can consider the paths $`e𝒳`$ that determine an analytical web $`w`$ also as possible edges of a truly infinite graph $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$. Clearly, considering the one point compactification $`\widehat{\mathrm{\Sigma }}`$ with an embedded generic element of $`\mathrm{\Gamma }_\sigma ^\omega `$ results in a highly singular object and therefore we do not have the luxury to do this.
In summary, essentially we do not change the spectrum $`\overline{𝒜}`$ as compared to except that the correspondence (5.2) is now extended to paths with non-compact range and therefore all the properties of $`\overline{𝒜}`$ derived in the literature are preserved.
One could ask whether there is a more fundamental reason for this choice, trying to define, as in the finite category, an Abelian $`C^{}`$ algebra of cylindrical functions depending on an infinite graph. This meets mathematical difficulties which are once more related to the fact that the associative law does not hold in general for the ITP and boils down to saying that one cannot really build an algebra of cylindrical functions over infinite graphs, only a vector space. We thus just adopt the above point of view with respect to definition of $`\overline{𝒜}`$. However, an outline of these difficulties will be given in the subsequent digression since it is instructive and gives rise to some natural definitions.
A natural way to proceed with infinite graphs comes from the observation that the set $`\mathrm{\Gamma }_0^\omega `$ is a subset of $`\mathrm{\Gamma }_\sigma ^\omega `$ which consists of compactly supported graphs. This observation motivates the following definition.
###### Definition 5.1
Let $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$.
i)
A function $`f`$ on $`𝒜`$ is said to be a $`C`$ function (not to be confused with the $`C`$ vectors of section 4.1) over $`\gamma `$ with values in $`\text{ }\mathrm{C}\{\mathrm{}\}`$ provided that for each $`eE(\gamma )`$ there exist functions $`f_e`$ on $`𝒜`$ of the form $`f_e(A)=F_e(h_e(A))`$, where $`F_e`$ is a complex valued function on $`G`$, such that
$$f\left(A\right)=\underset{eE\left(\gamma \right)}{}f_e\left(A\right)$$
(5.3)
and convergence is defined as in definition 4.1 where we set $`f(A)=\mathrm{}`$ if $`_e|f_e(A)|=\mathrm{}`$ irrespective of the phases of the $`f_e(A)`$.
ii)
A function $`f`$ on $`𝒜`$ is said to be cylindrical over $`\gamma `$ if it is a finite linear combination of $`C`$ functions over $`\gamma `$. The set of cylindrical functions over $`\gamma `$ is denoted by Cyl<sub>γ</sub>.
iii)
A function $`f`$ on $`𝒜`$ is said to be cylindrical if it is a finite linear combination of cylindrical functions over some graphs $`\gamma `$. The set of cylindrical functions is denoted Cyl.
iv)
An element $`0f=_{n=1}^Nz_n_{eE(\gamma )}f_e^{(n)}\text{Cyl}_\gamma ,z_n\text{ }\mathrm{C}`$ is said to be $`\sigma `$-bounded if and only if
$$f_\gamma :=\underset{U\mathrm{\Sigma }}{sup}\underset{A𝒜}{sup}\left|\underset{n=1}{\overset{N}{}}z_n\underset{eU\mathrm{}}{}f_e^{\left(n\right)}\left(A\right)\right|$$
(5.4)
is finite where $`U`$ runs over all compact subsets of $`\mathrm{\Sigma }`$. For $`f=0`$ we set $`f=0`$. Notice that the argument of the modulus in (5.4) is a cylindrical function in the sense of section 2.2. We will denote the set of $`\sigma `$-bounded, cylindrical functions by Cyl$`{}_{\gamma }{}^{}{}_{}{}^{b}`$. Notice that Cyl$`{}_{\gamma }{}^{}{}_{}{}^{b}`$ is not empty precisely due to the usual boundary conditions on smooth connections $`𝒜`$ for non-compact $`\mathrm{\Sigma }`$.
The norm (5.4) assigns a finite value to functions $`f`$ even if there is $`A𝒜`$ such that $`f\left(A\right)=\mathrm{}`$ which corresponds to our motivation to take over the structure from $`\mathrm{\Gamma }_0^\omega `$.
###### Lemma 5.2
The space of $`\sigma `$-bounded cylindrical functions over $`\gamma `$ forms a algebra with the $`C^{}`$ property.
Proof of Lemma 5.1 :
That Cyl$`{}_{\gamma }{}^{}{}_{}{}^{b}`$ is closed under linear combination, multiplication by scalars and factor-wise complex conjugation is obvious. Suppose now that $`f=_{m=1}^Mu_m_{eE\left(\gamma \right)}f_e^{\left(m\right)}`$, $`g=_{n=1}^Nv_n_{eE\left(\gamma \right)}g_e^{\left(n\right)}`$ are given and we define
$$fg:=\underset{m,n}{}u_mv_n\underset{eE\left(\gamma \right)}{}f_e^{\left(m\right)}g_e^{\left(n\right)}$$
(5.5)
simply by factorwise multiplication. Then
$`fg_\gamma `$ $`=`$ $`\underset{A,U}{sup}\left|{\displaystyle \underset{m,n}{}}u_mv_n{\displaystyle \underset{eU\mathrm{}}{}}f_e^{\left(m\right)}\left(A\right)g_e^{\left(n\right)}\left(A\right)\right|`$ (5.6)
$`=`$ $`\underset{A,U}{sup}\left|\left[{\displaystyle \underset{m=1}{\overset{M}{}}}u_m{\displaystyle \underset{eU\mathrm{}}{}}f_e^{\left(m\right)}\left(A\right)\right]\left[{\displaystyle \underset{n=1}{\overset{N}{}}}v_n{\displaystyle \underset{eU\mathrm{}}{}}g_e^{\left(n\right)}\left(A\right)\right]\right|`$
$``$ $`\left[\underset{A,U}{sup}\left|{\displaystyle \underset{m=1}{\overset{M}{}}}u_m{\displaystyle \underset{eU\mathrm{}}{}}f_e^{\left(m\right)}\left(A\right)\right|\right]\left[\underset{A,U}{sup}\left|{\displaystyle \underset{n=1}{\overset{N}{}}}v_n{\displaystyle \underset{eU\mathrm{}}{}}g_e^{\left(n\right)}\left(A\right)\right|\right]`$
$`=`$ $`f_\gamma g_\gamma `$
is also bounded. The $`C^{}`$ property follows from $`\left|\overline{f_U\left(A\right)}\right|=\left|f_U\left(A\right)\right|`$ and $`sup_{U,A}\left|f_U\left(A\right)\right|^2=\left(sup_{U,A}\left|f_U\left(A\right)\right|\right)^2`$.
$`\mathrm{}`$
So far we have considered only one cylindrical algebra Cyl$`{}_{\gamma }{}^{}{}_{}{}^{b}`$. Can we consider the algebra Cyl<sup>b</sup> of finite linear combinations of elements of Cyl$`{}_{\gamma }{}^{}{}_{}{}^{b}`$ for some $`\gamma `$’s ? As we have shown in lemma 5.1, $`\mathrm{\Gamma }_\sigma ^\omega `$ is a directed set so that for any finite collection $`\gamma _1,..,\gamma _n`$ there exists a $`\gamma `$ containing each of them. However, it may no longer be true that a given $`f_kCyl_{\gamma _k}^b,k=1,..,n`$ can be written as a finite linear combination of $`C`$ functions over $`\gamma `$, in fact, this will almost never be the case. Thus, while linear combinations pose no problem, products do as we then can no longer multiply factor-wise without having to consider infinite linear combinations of $`C`$ functions. In other words, as soon as we allow linear combinations of functions cylindrical over different infinite graphs, we end up having no algebra any more, products are ill-defined. The only exception is that for each of $`\gamma _1,..,\gamma _n`$ only a finite number of edges have to be decomposed into a finite number of segments each of which is an edge of $`\gamma `$. In that case, each of $`f_k`$ can be considered already as a function in Cyl$`{}_{\gamma }{}^{}{}_{}{}^{b}`$ so that nothing new is gained. Thus, the only way to proceed along the lines of is therefore to consider all the Cyl$`{}_{\gamma }{}^{}{}_{}{}^{b}`$ separately.
Once this is agreed on, the remainder is now standard. We complete the algebra Cyl$`{}_{\gamma }{}^{}{}_{}{}^{b}`$ in the norm (5.4) and obtain an Abelian $`C^{}`$ algebra $`_\gamma `$ which now depends on $`\gamma `$, in contrast to section 2.2. By the Gel’fand theorem we obtain the spectrum $`\overline{𝒜}_\gamma `$ of this algebra and $`_\gamma `$ is, via the Gel’fand transform $`f\widehat{f}`$, isometrically isomorphic to the algebra of continuous functions $`C^0\left(\overline{𝒜}_\gamma \right)`$ over the compact Hausdorff space $`\overline{𝒜}_\gamma `$. But now we meet the next difficulty and this finishes our trial to proceed this way : Namely, the set $`𝒜`$ is now no longer a subset of $`\overline{𝒜}_\gamma `$. Namely, let $`A_0𝒜\overline{𝒜}_\gamma `$ then we have from isometricity
$$f=\widehat{f}=\underset{\overline{A}\overline{𝒜}_\gamma }{sup}\left|\widehat{f}\left(\overline{A}\right)\right|=\underset{\overline{A}\overline{𝒜}_\gamma }{sup}\left|\overline{A}\left(f\right)\right|\left|A_0\left(f\right)\right|=\left|f\left(A_0\right)\right|$$
(5.7)
which from the definition (5.4) can be true only if $`A_0`$ has compact support. However, we are precisely interested in (distributional) connections which are supported everywhere in $`\mathrm{\Sigma }`$ as this corresponds to the intended physical application in connection with the clasical limit for non-compact $`\mathrm{\Sigma }`$. There is no claim that one could not introduce a different $`C^{}`$ norm on cylindrical functions which would lead to the desired distributional extension of $`𝒜`$ but there seems to be no obvious, natural candidate as the above discussion reveals. We leave the question on the existence of such a norm for future research. This terminates our digression.
We thus will not use the norm (5.4) any further but simply consider the vector space Cyl of arbitrary cylindrical functions of $`𝒜`$ without any convergence requirements, to begin with. As we will see, a subset of this space, extended to distributional connections, is dense in the Hilbert space which we are going to construct and although it is not an algebra, inner products can be computed even if we have linear combinations of functions over different infinite graphs.
This extension works as follows : Since every cylindrical function is a finite linear combination of $`C`$ functions over some $`\gamma `$ we can also extend any $`f`$Cyl to a function $`\widehat{f}`$ on $`\overline{𝒜}`$ simply by the pull-back of the Gel’fand transform on $`C`$ functions
$$\widehat{f}:=\underset{eE\left(\gamma \right)}{}\widehat{f}_e\text{ where }\widehat{f}_e\left(\overline{A}\right)=F_e\left(\overline{A}\left(h_e\right)\right)=F_e\left(\widehat{h}_e\left(\overline{A}\right)\right)=\left(^{}f_e\right)\left(\overline{A}\right)$$
(5.8)
extended by linearity. The notation means that $`\widehat{f}`$ is the Gel’fand transform of $`f=_ef_e,f_e=F_eh_e`$ extended from finite to infinite graphs. We will continue to call the extensions $`\widehat{f}`$ cylindrical functions.
Although a general $`\widehat{f}`$Cyl will take an infinite value on almost every point $`\overline{A}\overline{𝒜}`$ it is still possible to equip Cyl with a topology which is weaker than the Hilbert space topology that we are going to construct, moreover, the Hilbert space measure is such that these infinite values are integrable. This is important in order to have a framework for solving quantum constraints via (analogs of) Gel’fand triples . we postpone the definition of this topology to subsection 5.2.1.
#### 5.1.3 Measure and Hilbert Space
Consider for a moment the set $`𝒞^\omega `$ of all possible oriented, analytic curves in $`\mathrm{\Sigma }`$. Clearly, at most countable collections of elements $`e𝒞_\sigma ^\omega `$ constitute an element $`\gamma 𝒢_\sigma ^\omega `$ through their union if that union is $`\sigma `$-finite. The idea is now to construct the Infinite Tensor Product Hilbert spaces $`_\gamma ^{}`$ associated with the Hilbert spaces $`_e,eE\left(\gamma \right)`$ of section 3.2, that is,
$$_\gamma ^{}:=_{eE\left(\gamma \right)}_e$$
(5.9)
Using the notation of section 4.3 we would have index sets $`=𝒞^\omega `$ and the set of arbitrary index subsets $``$ (or power set) of $``$ of which $`\mathrm{\Gamma }_\sigma ^\omega `$ is a proper subset.
The reader may now wonder why we do not use the full power of the Infinite Tensor Product of being able to deal with index sets of arbitrary cardinality and rather stick with $`\mathrm{\Gamma }_\sigma ^\omega `$. Indeed, an interesting observation is now the following : Consider instead of $``$ the slightly smaller set $`𝒫`$ of arbitrary subsets $`C`$ of $`𝒞^\omega `$ (not necessarily elements of $`\mathrm{\Gamma }_\sigma ^\omega `$) such that no element $`eC`$ can be written as a composition of elements of $`C\left\{e\right\}`$ and their inverses. Then we say $`CC^{}`$ if every element $`eC`$ can be written as a composition of elements $`e^{}C^{}`$ and their inverses which gives also $`𝒫`$ a partial order. For $`CC^{}`$ we define $`CC^{}=C^{}`$. Recall that a subset $`P𝒫`$ is called a chain if all elements $`CP`$ are in relation $``$. Given a chain $`P`$, consider the element $`C_P:=_{CP}C`$ which is an element of $`𝒫`$ (not necessarily of $`P`$), moreover, $`CC_PCP`$. In other words, every chain in $`𝒫`$ has an upper bound in $`𝒫`$ and by the lemma of Zorn we obtain that $`𝒫`$ has a maximal element $`C_{\mathrm{}}`$, that is, $`CC_{\mathrm{}}`$ for all $`C𝒫`$. Certainly, there are infinitely many such maximal elements each of which we will call a “supergraph”. By construction, every element $`eC_{\mathrm{}}`$ is not composition of elements of $`C_{\mathrm{}}\left\{e\right\}`$ and thus they are holonomically independent. (This construction can obviously be repeated for the smooth category of curves as well). Notice that the existence of $`C_{\mathrm{}}`$, while of theoretical interest since it allows us to construct the universal ITP $`_{\mathrm{}}^{}:=_{eC_{\mathrm{}}}_e`$, universal in the sense that every possible piecewise analytic graph $`\gamma `$ can be written as composition of elements of $`C_{\mathrm{}}`$, it is practically so far of modest interest only because 1) no one knows how to describe $`C_{\mathrm{}}`$ explicitly and 2) even if one knew $`C_{\mathrm{}}`$ explicitly, given $`\gamma \mathrm{\Gamma }^\omega `$, every edge $`e`$ of $`\gamma `$ would generically decompose into an infinite number of segments each of which is an element of $`E\left(C_{\mathrm{}}\right)`$. Thus, even a very simple function from the point of view of $`\gamma `$ would look very complicated from the point of view of $`C_{\mathrm{}}`$. In particular, as we have seen already in section 4.1, the associative law fails for the ITP and it will in general happen that a function on an incomplete ITP associated with some $`\gamma `$ cannot be written as an element of the universal ITP. We are therefore forced to work with all the $`_\gamma ^{}`$ simultaneously rather than with the single universal object $`_{\mathrm{}}^{}`$ only.
However, the supergraph $`\gamma _{\mathrm{}}`$ allows us to give a simple proof of the existence of a $`\sigma `$-additive, faithful, Borel measure on $`\overline{𝒜}`$ with respect to which we can compute arbitrary inner products of cylindrical functions. This is a simple corollary of the Kolmogorov theorem for the case of an uncountably infinite tensor product of probability measures and works as follows in the present context :
The supergraph $`C_{\mathrm{}}𝒫`$ is a generating system of holonomically independent analytic curves for every element $`P𝒫`$, in particular, for every element $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$. Each element $`\overline{A}`$ of the Ashtekar-Isham space $`\overline{𝒜}`$ of generalized connections assigns to each curve $`eC_{\mathrm{}}`$ an element $`\overline{A}\left(h_e\right)=\widehat{h}_e\left(\overline{A}\right)`$ of $`G`$ and as $`\overline{A}`$ varies, this map is onto (except if $`e`$ is just a point in which case $`\overline{A}\left(h_e\right)=1_G`$). Given $`P𝒫`$ we consider the $`\sigma `$-algebra $`_P`$ generated by preimages of Borel subsets of $`G^{\left|P\right|}`$ under the map $`p_P:\overline{𝒜}G^{\left|P\right|};\overline{A}\left\{\overline{A}\left(h_e\right)\right\}_{eP}`$ where $`\left|P\right|`$ denotes the cardinality of the set $`P`$. Consider the $`\sigma `$-algebra $``$ generated by all the $`_P`$ displaying $`(\overline{𝒜},)`$ as a measurable space. We say that a function $`f`$ is measurable if it is of the form $`f=Fp_P`$ for some $`P𝒫`$ and some function $`F`$ on $`G^{\left|P\right|}`$. A measure on $`\overline{𝒜}`$ can now be defined on measurable functions by
$$\mu _0\left(f\right):=_{G^{\left|P\right|}}_{eP}d\mu _H\left(h_e\right)F\left(\left\{h_e\right\}_{eP}\right)$$
(5.10)
where $`\mu _H`$ is the Haar measure on $`G`$. The normalization, right – and left invariance and the invariance under inversion display this measure as a consistently defined measure on measurable functions, the proof is completely analogous to the one displayed in so that we can omit it here. Notice, however, that in contrast to we allow $`\left|P\right|=\mathrm{}`$. Being consistently defined, the measure qualifies as one to have a $`\sigma `$-additive extension to $``$. Rather than using the existence theorem of we simply write it down :
$$\mu _0(.):=_{G^{\left|C_{\mathrm{}}\right|}}_{eC_{\mathrm{}}}d\mu _H\left(h_e\right)(.)$$
(5.11)
We will call it the extended Ashtekar-Lewandowski measure. Again, the remarkable properties of the Haar measure on $`G`$ reveal that (5.11) reduces to (5.10) for $`f=Fp_P`$ and theorem 12.1 in guarantees that (5.11) has the required properties. Thus, although we do not know the object $`C_{\mathrm{}}`$, its mere existence can be used to define $`\mu _0`$. Henceforth we will denote $`:=L_2(\overline{𝒜},d\mu _0)`$.
A more explicit construction of that Hilbert space is as follows and this provides a simple way to embed the kinematical framework of section 2.2 for finite piecewise analytical graphs into the context of $`\mathrm{\Gamma }_\sigma ^\omega `$ of infinite piecewise analytical $`\sigma `$-finite graphs.
Given $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$ we can use the inductive limit construction of section 4.3 to obtain $`_\gamma ^{}`$ for infinite $`\gamma `$ from the Hilbert spaces $`_\gamma `$ constructed in for finite $`\gamma `$. Notice that we get this way a genuine extension of the so-called Ashtekar-Lewandowski Hilbert space
$$_{AL}:=\overline{_{\gamma \mathrm{\Gamma }_0^\omega }_\gamma }$$
(5.12)
as we will see in a moment. In fact, the Hilbert space that we will construct is defined by
$$^{}:=\overline{_{\gamma \mathrm{\Gamma }_\sigma ^\omega }_\gamma ^{}}$$
(5.13)
Let us then proceed to the explicit construction. Recall that $`_\gamma `$ is the completion, with respect to the Ashtekar-Lewandowski measure $`\mu _0`$ of section 2.2, of the space of cylindrical functions $`\text{Cyl}_\gamma `$ over $`\gamma \mathrm{\Gamma }_0^\omega `$. Since $`\text{Cyl}_\gamma `$ can be replaced by the finite linear combinations of (non-coloured) spin-network functions over $`\gamma `$ we see that $`_\gamma `$ can be equivalently described as the closure of the finite linear combinations of $`C_0`$-vectors of the finite tensor product
$$_\gamma ^{}=_{eE\left(\gamma \right)}_e$$
(5.14)
where each $`_e`$ is isometric isomorphic with $`L_2(G,d\mu _H)`$ where $`\mu _H`$ is the Haar measure. Thus, $`_\gamma =_\gamma ^{}`$ for $`\gamma \mathrm{\Gamma }_0^\omega `$.
Indeed, as it is immediately obvious from the cylindrical consistency of the measure $`\mu _0`$, it reduces on $`_\gamma `$ precisely to the tensor product Haar measure, corresponding to one copy of $`G`$ for each $`e\gamma `$ and this is precisely the original definition of the Ashtekar-Lewandowski measure in terms of its cylindrical projections given in .
Let now $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$ be given, then we find a sequence of elements $`\gamma _n\mathrm{\Gamma }_0^\omega `$ such that $`\gamma _n\gamma `$ and $`\gamma _n\gamma _{n+1}`$ for each $`n=1,2..`$, moreover $`_{n=1}^{\mathrm{}}\gamma _n=\gamma `$. By means of the isometric monomorphisms defined on $`C_0`$-vectors for $`mn`$
$$\widehat{U}_{\gamma _m\gamma _n}:_{\gamma _m}^{}_{\gamma _n}^{};_{eE\left(\gamma _m\right)}f_e\left[_{eE\left(\gamma _m\right)}f_e\right]\left[_{eE\left(\gamma _n\gamma _m\right)}1\right]$$
(5.15)
and extended by linearity, where $`1\left(A\right)=1`$ is the unit function, we display the system of Hilbert spaces $`_\gamma _\gamma ^{}`$ as a directed system. By theorem 4.7 the unique inductive limit of the $`_{\gamma _n}`$ exists and can be identified with the Infinite Tensor Product for each $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$
$$_\gamma ^{}=_{eE\left(\gamma \right)}_e$$
(5.16)
and indeed the required isometric isomorphisms are given by
$$\widehat{U}_{\gamma _n}:_{\gamma _n}^{}_\gamma ^{};_{eE\left(\gamma _n\right)}f_e\left[_{eE\left(\gamma _n\right)}f_e\right]\left[_{eE\left(\gamma \gamma _n\right)}1\right]$$
(5.17)
Thus, for a truly infinite graph $`\gamma `$ the Hilbert space $`_\gamma ^{}`$ is hyperfinite, that is, it is the inductive limit of the finite dimensional Hilbert spaces $`_{\gamma _n}^{}`$.
Several remarks are in order :
* From the point of view of $`_\gamma ^{}`$ the vectors of $`_{\gamma _n}_{\gamma _n}^{}`$ lie in the strong equivalence class of the $`C_0`$-sequence $`f^0=\left\{f_e^0\right\}_{eE\left(\gamma \right)}`$ where $`f_e^0=1`$ for each $`e`$. This is an immediate consequence of lemma 4.8. It follows that the Hilbert space (5.12) is just a tiny subspace of the Hilbert space (5.13) since every vector over $`\gamma `$ which is not in the strong equivalence class of $`f^0`$ is orthogonal to all of the $`_{\gamma _n}`$ and there are uncountably infinitely many different strong eqivalence classes even for fixed $`\gamma `$ as follows from lemma 4.14 since $`\left|E\left(\gamma \right)\right|=\mathrm{}`$. To see this in more detail, notice that if a generic element $`\xi _\gamma ^{}`$ would be a Cauchy sequence of elements $`\xi _n_{AL}`$ then for any $`ϵ>0`$ we would find $`n_0\left(ϵ\right)`$ such that $`\xi \xi _n<ϵn>n_0\left(ϵ\right)`$. Now each $`\xi _n`$ can be chosen to be in some $`_{\gamma _n}`$ with $`\mathrm{\Gamma }_0^\omega \gamma _n\gamma `$ since any vector in the completion of $`_{AL}`$ can be approximated by vectors of that form and since any vector depending on a coloured graph which is not contained in $`\gamma `$ is automatically orthogonal to $`\xi `$. However, if we choose, e.g., $`\xi `$ to be a linear combination of $`C_0`$-vectors each of which lies in a different strong equivalence class Hilbert space than the vector $`f^0`$ above then we get the contradiction $`\xi ^2<\xi _n^2+\xi ^2<ϵn>n_0\left(ϵ\right)`$.
* While the Ashtekar-Lewandowski Hilbert space is just a tiny subspace of $`^{}`$ in (5.5), the Ashtekar-Lewandowski measure is still the appropriate measure to use in our extended context.
Indeed, it has been identified already as the $`\sigma `$-additive extension of the cylindrically defined measure of to the projective (or inductive) limit of arbitrarily large and complicated, but finite piecewise analytic graphs in . Therefore, it could be used to date only in order to integrate special functions depending on an infinite number degrees of freedom (i.e. depending on infinite graphs) : Namely those which can be written as infinite sums of functions each of which depends only on a finite graph, an exception being where some sort of infinite volume limit has been taken. One contribution of the present paper is to show that the measure can be used to integrate more general functions depending on an infinite number of degrees of freedom : namely those which are infinite products of functions each of which depends only on a finite graph.
* In the context of finite graphs we can (even for non-gauge-invariant states) write down a complete orthonormal (with respect to the Ashtekar-Lewandowski measure $`\mu _0`$) basis, the so-called spin-network basis of section 2.2. It is frequently stressed that the Ashtekar-Lewandowski measure can then be dispensed with by just requiring these functions to be orthogonal and to check that a positive definite sesquilinear form results in this way . Adopting this point of view, given arbitrary functions in $`_{AL}`$ one can explicitly compute their inner products by writing them in terms of spin-network functions and using sesquilinearity. This is no longer possible in the context of the Infinite Tensor Product (spatially non-compact $`\mathrm{\Sigma }`$), here the Ashtekar-Lewandowski measure is the only way to calculate inner products ! To see this we just need to display one simple example :
Consider an infinite graph $`\gamma `$ with a countable number of analytic edges $`e`$ (say a cubic lattice in $`\mathrm{\Sigma }=\text{ }\mathrm{R}^3`$). Consider the $`C_0`$-sequence $`f:=\left\{f_e\right\}`$ where (from now on we drop the bar in $`\overline{A}`$ for a distributonal connection and we write $`h_e`$ instead of the Gel’fand transform $`\widehat{h}_e`$)
$$f_e\left(A\right):=f^0\left(h_e\left(A\right)\right):=\frac{1+\chi _j\left(h_e\left(A\right)\right)}{\sqrt{2}}$$
(5.18)
where $`\chi _j`$ is again the character in the spin $`j>0`$ representation of $`SU\left(2\right)`$. Using the extended Ashtekar Lewandowski measure (5.11), which on this infinite graph just reduces to $`d\mu _{0\gamma }=_{eE\left(\gamma \right)}d\mu _H\left(h_e\right)`$ we immediately verify that the norm of the $`C_0`$-vector $`_f`$ equals unity. Suppose now we wanted to use only the knowledge that the set of functions
$$\underset{k=1}{\overset{n}{}}\chi _j\left(h_{e_k}\right)$$
(5.19)
for finite $`n`$ and mutually distinct $`e_1,..,e_nE\left(\gamma \right)`$ are mutually orthogonal spin-network functions. Then, in order to compute the norm of $`_f`$ we would need to decompose this vector into the latter set of functions which at least formally can be done using the distributive law over and over again. However, it is easy to see that each of these infinite number of terms comes with the coefficient $`\left(1/\sqrt{2}\right)^{\mathrm{}}`$ and so our attempt to compute the norm would result in the ill-defined expression $`0\mathrm{}`$. More precisely, this ill-defined result is due to the fact that the inner product between the vectors (5.18) and (5.19) for $`n\mathrm{}`$ equals zero.
Concluding, in the ITP or $`\mathrm{\Gamma }_\sigma ^\omega `$ category the spin-network functions no longer provide a basis (a related observation has been made independently already in in the context of webs or $`\mathrm{\Gamma }_0^{\mathrm{}}`$), simply because, even for a single $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$, the orthonormal set of functions given by
$$A\underset{e}{}\left[\sqrt{2j_e+1}\pi _{j_e}\left(h_e\left(A\right)\right)_{m_en_e}\right]$$
(5.20)
where $`eE\left(\gamma \right);\mathrm{\hspace{0.33em}2}j_e=0,1,2,..;m_e,n_e=j_e,j_e+1,..,j_e`$ and which from experience with spin-network functions one might think to provide a basis, is not complete ! For instance, the unit $`C_0`$-vector
$$A\underset{e}{}\chi _j\left(h_e\left(A\right)\right)$$
(5.21)
is orthogonal to all of them for any $`j>0`$, even if we choose $`j_e=j`$ for all $`e`$ since $`|<\sqrt{2j+1}\pi _{jmn},\chi _j>_{L_2(SU\left(2\right),d\mu _H)}|1/\sqrt{2j+1}<1`$. The ITP Hilbert space has many more orthogonal directions than one is used to due to its non-separability. A complete orthonormal bases on a single $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$ is not given by spin-network functions but rather by a von Neumann basis defined in lemma 4.9 and corollary 4.2, one for each $`\left[f\right]`$-adic Infinite Tensor Product subspace of $`_\gamma `$. The only $`\left[f\right]`$-adic ITP that has indeed a spin-network basis is the one given by the trivial strong equivalence class $`\left[f^0\right]`$ where for any given $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$ we have $`f_e^0=1`$ for each $`eE\left(\gamma \right)`$.
Our treatment is still incomplete because, while we can compute inner products between finite linear combinations of $`C_0`$ vectors over a single $`\gamma `$, nothing has so far been said about inner products between finite linear combinations of $`C_0`$ vectors over different $`\gamma `$’s and this is what we need if we wish to glue together the $`_\gamma ^{}`$ as displayed in (5.13). The idea is, of course, to use the inductive limit construction once again, however, as far as inner products are concerned we have to go somewhat beyond von Neumann’s theory which tells us only how to compute inner products between finite linear combinations of $`C_0`$ vectors over the same $`\gamma `$.
A concrete and natural definition can be given employing the extended Ashtekar-Lewandowski measure. Let us derive it, proceeding formally to begin with :
Let $`\gamma ,\gamma ^{}\mathrm{\Gamma }_\sigma ^\omega `$ and let $`f_\gamma ,g_\gamma ^{}`$ respectively be $`C_0`$-sequences over $`\gamma ,\gamma ^{}`$ respectively. Consider the graph $`\gamma ^{\prime \prime }:=\gamma \gamma ^{}`$. The idea is to define the inner product between the corresponding $`C_0`$-vectors $`_{f_\gamma },_{g_\gamma ^{}}`$ by
$`<_{f_\gamma },_{g_\gamma ^{}}>`$ $`:=`$ $`{\displaystyle _{G^{\left|C_{\mathrm{}}\right|}}}\left[_{e^{\prime \prime }C_{\mathrm{}}}d\mu _H\left(h_{e^{\prime \prime }}\right)\right]\overline{\left[_{f_\gamma }\right]}\left[_{g_\gamma ^{}}\right]`$ (5.22)
$`=`$ $`{\displaystyle _{G^{\left|E\left(\gamma ^{\prime \prime }\right)\right|}}}\left[_{e^{\prime \prime }E\left(\gamma ^{\prime \prime }\right)}d\mu _H\left(h_{e^{\prime \prime }}\right)\right]\overline{\left[_{f_\gamma }\right]}\left[_{g_\gamma ^{}}\right]`$
where the second equality follows from cylindrical consistency. The problem, that by now we are already used to with these infinite tensor products, is that the associative law does not hold. In other words, the ITP
$$_{\gamma ^{\prime \prime }}^{}:=_{e^{\prime \prime }E\left(\gamma ^{\prime \prime }\right)}_{e^{\prime \prime }}$$
(5.23)
is in general quite different from the subdivisions
$`\left(_{eE\left(\gamma \right)}\left[_{e^{\prime \prime }E\left(\gamma ^{\prime \prime }\right)e}_{e^{\prime \prime }}\right]\right)\left(_{e^{\prime \prime }E\left(\gamma ^{\prime \prime }\right)\gamma }_{e^{\prime \prime }}\right)\text{ and}`$
$`\left(_{e^{}E\left(\gamma ^{}\right)}\left[_{e^{\prime \prime }E\left(\gamma ^{\prime \prime }\right)e^{}}_{e^{\prime \prime }}\right]\right)\left(_{e^{\prime \prime }E\left(\gamma ^{\prime \prime }\right)\gamma ^{}}_{e^{\prime \prime }}\right)`$ (5.24)
to which $`f_\gamma ,g_\gamma ^{}`$ belong respectively. This is precisely the problem outlined at the end of section 4.1 : The correspondence with the notation there is that $`=E\left(\gamma ^{\prime \prime }\right),=E\left(\gamma \right)\left\{\gamma ^{\prime \prime }\gamma \right\},_l=E\left(\gamma ^{\prime \prime }\right)e`$ for $`l=e`$ and $`_l=E\left(\gamma ^{\prime \prime }\right)\gamma `$ for $`l=\gamma ^{\prime \prime }\gamma `$ and similar for $`\gamma ^{}`$. Here we have identified $`_e`$ with $`_{e^{\prime \prime }E\left(\gamma ^{\prime \prime }\right)e}_{e^{\prime \prime }}`$ and similar for $`e^{}`$. Notice that in general $`\left|E\left(\gamma ^{\prime \prime }\right)e\right|,\left|E\left(\gamma ^{\prime \prime }\right)e^{}\right|=\mathrm{}`$, an example being given by two graphs consisisting of a single edge only, $`\gamma =e,\gamma ^{}=e^{}`$, which however both have non-compact range and intersect each other an infinite number of times in isolated points. This is not excluded by piecewise analyticity since there is no accumulation point of intersection points (take, e.g. $`e=(x,0)`$ and $`e^{}=(x,\mathrm{sin}\left(x\right))`$ in $`\text{ }\mathrm{R}^2`$).
Step I)
In order to proceed, we subdivide $`\gamma ^{\prime \prime }`$ into the mutually disjoint sets $`\gamma ^{}:=\gamma \gamma ^{},\overline{\gamma }=\gamma ^{\prime \prime }\gamma ,\overline{\gamma }^{}=\gamma ^{\prime \prime }\gamma ^{}`$. Then we formally embed $`_{f_\gamma }`$ into $`_{\gamma ^{\prime \prime }}`$ by identifying it with $`\left(_{f_\gamma }\right)\left(_{e^{\prime \prime }E\left(\gamma ^{\prime \prime }\right)\overline{\gamma }}1\right)`$ and similarly we identify $`_{g_\gamma ^{}}`$ with $`\left(_{g_\gamma ^{}}\right)\left(_{e^{\prime \prime }E\left(\gamma ^{\prime \prime }\right)\overline{\gamma }^{}}1\right)`$. Clearly, we will now perform first the easy integrals corresponding to the tensor products factors of the unit function. In order to do this, for given $`eE\left(\gamma \right)`$ we recall that we can write $`f_e\left(h_e\right)=_\pi f_{e\pi }^{mn}\pi _{mn}\left(h_e\right)`$ by the Peter& Weyl theorem where the sum is over a complete set of equivalence classes of irreducible representations of $`G`$, $`\pi _{mn}`$ denotes the matrix elements of a group element in the representation $`\pi `$ and $`f_{e\pi }^{mn}`$ are the Fourier coefficients of $`f_e`$. Now suppose that $`e^{\prime \prime }e`$ and that $`e^{\prime \prime }\gamma ^{}`$, that is, $`e^{\prime \prime }E\left(\gamma ^{\prime \prime }\right)\left(e\overline{\gamma }^{}\right)`$. Then we can write $`f_e\left(h_e\right)=_\pi f_{e\pi }^{mn}\pi _{mn}\left(h_e^{\left(1\right)}h_eh_e^{\left(2\right)}\right)`$ where $`h\left(1\right)_e,h_e^{\left(2\right)}`$ depend on $`ee^{\prime \prime }`$. In other words, we can consider it as a function $`F`$ of $`h_{e^{\prime \prime }}`$ only and as far as the integral over $`h_{e^{\prime \prime }}`$ is concerned it reduces to evaluating $`<F,1>_{L_2(G,d\mu _H)}=\overline{f_{e\pi _0}}=<f_e,1>__e`$. It follows that for any $`eE\left(\gamma \right)`$ which is not fully overlapped by edges of $`E\left(\gamma ^{}\right)`$ we can replace $`f_e`$ by $`<f_e,1>`$ (we drop the index at the inner product). Likewise, for any $`e^{}E\left(\gamma \right)`$ which is not fully overlapped by edges of $`E\left(\gamma \right)`$ we can replace $`g_e^{}`$ by $`<1,g_e^{}>`$. This is the result of performing the integral over all $`h_{e^{\prime \prime }}`$ with $`e^{\prime \prime }E\left(\gamma ^{\prime \prime }\right)\left[\overline{\gamma }\overline{\gamma }^{}\right]=E\left(\gamma ^{\prime \prime }\right)\left[\gamma ^{\prime \prime }\gamma ^{}\right]`$. It remains to perform the integral over the edges of $`E\left(\gamma ^{\prime \prime }\right)\gamma ^{}`$.
Step II)
Now notice that from $`_{eE\left(\gamma \right)}f_e`$ only those factors are left corresponding to edges $`e`$ fully overlapped by edges of $`E\left(\gamma ^{}\right)`$ and from $`_{e^{}E\left(\gamma ^{}\right)}g_e^{}`$ only those factors are left corresponding to edges $`e^{}`$ fully overlapped by edges of $`E\left(\gamma \right)`$. Let us denote the corresponding subsets by $`E\left(\gamma \right)_{|\gamma ^{}}E\left(\gamma \right),E\left(\gamma ^{}\right)_{|\gamma ^{}}E\left(\gamma ^{}\right)`$. The union of both sets of edges is contained in $`\gamma ^{}`$. Suppose now that $`eE\left(\gamma \right)_{|\gamma ^{}}`$ is overlapped by a collection of edges $`e^{}`$ of $`E\left(\gamma ^{}\right)`$, that is, there is a countable number of edges $`e_{10}^{},..,e_{11}^{}`$ of $`E\left(\gamma ^{}\right)`$ so that the endpoint of one is the starting point of the next, such that $`e`$ is contained in their union and such that it is not any more contained if we remove $`e_{10}^{}`$ or $`e_{11}^{}`$ from the collection. It follows that $`e_{10}^{},..,e_{11}^{}`$ are analytical continuations of each other.
Step III)
Let us first focus on $`e_{10}^{}`$. Now either, A) $`e_{10}^{}`$ is also contained in $`e`$ or, B) it is not. In case B), if there are no other edges of $`E\left(\gamma \right)`$ overlapping the remaining segment of $`e_{10}^{}`$ not contained in $`e`$ then the edge $`e_{10}^{}`$ does not appear any more in $`E\left(\gamma ^{}\right)_{|\gamma ^{}}`$. By the same argument as in Step I), if we now perform the integral over any $`h_{e^{\prime \prime }}`$ with $`e^{\prime \prime }`$ contained in $`ee_{10}^{}`$ then we can replace $`\overline{f_e}`$ by $`<f_e,1>`$ and that factor also drops out of the integral. Thus, we can focus on the case that $`e_{10}^{}E\left(\gamma ^{}\right)_{|\gamma ^{}}`$ , that is, there are such other edges $`e_0,..,e_1`$ of $`E\left(\gamma \right)`$ where $`e_0`$ is adjacent to $`e`$, an endpoint of one is the starting point of the next and if $`e_1`$ is removed, the collection $`e,e_0,..,e_1`$ no longer overlaps $`e_{10}^{}`$. We see that $`e,e_0,..,e_1`$ are analytical continuations of each other. Now either, A) $`e_1`$ is also contained in $`e_{10}^{}`$ or B), it is not. In case B), if there are no other edges of $`E\left(\gamma ^{}\right)`$ overlapping the remaining segment of $`e_1`$ not contained in $`e_{10}^{}`$ then $`e_1`$ does not belong to $`E\left(\gamma \right)_{|\gamma ^{}}`$ and so as in Step I) we can replace $`g_{e_{10}^{}}`$ by $`<1,g_{e_{10}^{}}>`$ and so that factor drops out of the integral. However, then as just explained also $`\overline{f_e}`$ drops out of the integral. Thus, we may assume that $`e_1E\left(\gamma \right)_{|\gamma ^{}}`$ and there are new edges $`e_{20}^{},..e_{21}^{}`$ with $`e_{21}^{}`$ adjacent to $`e_{01}^{}`$, the endpoint of one is the starting point of the next and such that $`e_1`$ is no longer overlapped if we remove $`e_{20}^{}`$.
Let us now rename $`ee_0..e_1`$ by $`e`$ and the collection $`e_{20}^{},..,e_{21}^{},e_{10}^{},..,e_{11}^{}`$ by $`e_{10}^{},..,e_{11}^{}`$. Then we are in the same situation as in the beginning of Step III). Iterating, we conclude that either we end up with case A) or with case B) but that $`e_{10}^{}`$ is no longer overlapped. In case B) the whole chain collapses like a cardhouse and we can replace $`\overline{f_e}`$ by $`<f_e,1>`$. In case A) we see that we found a maximal analytical continuation of the original $`e`$, into the direction of its starting point, by other edges of $`E\left(\gamma \right)`$ all of which are overlapped by edges of $`E\left(\gamma ^{}\right)`$ and those edges are also contained in that maximal analytical continuation contained in $`\gamma `$.
Step IV)
Now we focus on $`e_{11}^{}`$ and proceed completely analogously. The end result is that $`\overline{f_e}`$ can be replaced by $`<f_e,1>`$ unless there exists a maximal bothsided maximal analytic continuation $`\stackrel{~}{e}`$ of $`e`$ by edges of $`E\left(\gamma \right)`$ completely overlapped by edges of $`E\left(\gamma ^{}\right)`$ and those edges of $`E\left(\gamma ^{}\right)`$ are also completely overlapped by $`\stackrel{~}{e}`$.
Step V)
As the argument is completely symmetric with respect to $`\gamma ,\gamma ^{}`$ we conclude that the remaining integral depends only on the graph $`\stackrel{~}{\gamma }`$ consisting of analytical edges $`\stackrel{~}{e}`$ which can be written simultaneously as compositions of edges of $`E\left(\gamma \right)`$ alone and edges of $`E\left(\gamma ^{}\right)`$ alone. For all other edges $`eE\left(\gamma \right)\stackrel{~}{\gamma }`$ we can replace $`\overline{f_e}`$ by $`<f_e,1>`$ and for all other edges $`e^{}E\left(\gamma ^{}\right)\stackrel{~}{\gamma }`$ we can replace $`g_e^{}`$ by $`<1,g_e^{}>`$. We are thus left with
$`<_{f_\gamma },_{g_\gamma ^{}}>`$
$`=`$ $`[{\displaystyle \underset{eE\left(\gamma \right)\stackrel{~}{\gamma }}{}}<f_e,1>][{\displaystyle \underset{e^{}E\left(\gamma ^{}\right)\stackrel{~}{\gamma }}{}}<1,g_e^{}>][{\displaystyle \underset{\stackrel{~}{e}E\left(\stackrel{~}{\gamma }\right)}{}}<\left[_{eE\left(\gamma \right)\stackrel{~}{e}}f_e\right],\left[_{e^{}E\left(\gamma ^{}\right)\stackrel{~}{e}}g_e^{}\right]>]`$
Step VI)
It remains to compute the inner product labelled by $`\stackrel{~}{e}`$ in (5.1.3). For each $`\stackrel{~}{e}\stackrel{~}{\gamma }`$ consider its unique breakup into segments $`e^{\prime \prime }E\left(\gamma ^{\prime \prime }\right)`$ defined by the breakpoints given by the union of the endpoints of the $`E\left(\gamma \right)e\stackrel{~}{e}`$ and $`E\left(\gamma ^{}\right)e^{}\stackrel{~}{e}`$ respectively. Then the last inner product in (5.1.3) is defined by
$`<_{eE\left(\gamma \right)\stackrel{~}{e}}f_e,_{e^{}E\left(\gamma ^{}\right)\stackrel{~}{e}}g_e^{}>`$
$`:=`$ $`{\displaystyle _{G^{\left|\stackrel{~}{e}\right|}}}\left[_{e^{\prime \prime }\stackrel{~}{e}}d\mu _H\left(h_{e^{\prime \prime }}\right)\right]\overline{\left[{\displaystyle \underset{eE\left(\gamma \right)\stackrel{~}{e}}{}}f_e\left({\displaystyle \underset{e^{\prime \prime }e}{}}h_{e^{\prime \prime }}\right)\right]}\left[{\displaystyle \underset{e^{}E\left(\gamma ^{}\right)\stackrel{~}{e}}{}}g_e^{}\left({\displaystyle \underset{e^{\prime \prime }e^{}}{}}h_{e^{\prime \prime }}\right)\right]`$
where we have symbolically written the holonomies along the edges $`e,e^{}`$ respectively as products of holonomies along the $`e^{\prime \prime }`$. The integral (5.1.3) is already well-defined if the number of $`e^{\prime \prime }\stackrel{~}{e}`$ is finite, if not, then we proceed as follows :
Since $`\gamma ,\gamma ^{}`$ are both $`\sigma `$-finite graphs, $`\stackrel{~}{e}`$ must be an infinite curve in $`\mathrm{\Sigma }`$ with either A) one or B) both ends at infinity, otherwise there would be an accumulation point. If only one endpoint is at infinity, choose the other point as the starting point of $`\stackrel{~}{e}`$. If both endpoints are at infinity, choose an arbitrary breakpoint $`p`$ on $`\stackrel{~}{e}`$ and choose it as the startpoint of the the resulting semi-infinite curves, that is, $`\stackrel{~}{e}=\left(\left[\stackrel{~}{e}^{\left(1\right)}\right]^1\left[\stackrel{~}{e}^{\left(2\right)}\right]\right)^{\pm 1}`$ is a choice of orientation of $`\stackrel{~}{e}`$. Since $`\gamma ,\gamma ^{}`$ are both $`\sigma `$-finite graphs, the number of $`e^{\prime \prime }\stackrel{~}{e}`$ is at most countable and we can label them by integers which are increasing into the direction of the orientation of $`\stackrel{~}{e}`$ in case A) and of $`\stackrel{~}{e}^{(1,2)}`$ respectively in case B), that is, $`\stackrel{~}{e}=e_1^{\prime \prime }e_2^{\prime \prime }\mathrm{}`$ and $`\stackrel{~}{e}^{(1,2)}=e_1^{(1,2)\prime \prime }e_2^{(1,2)\prime \prime }\mathrm{}`$ respectively. The integral is then defined by performing the integrals over the $`h_{e_n^{\prime \prime }}`$ in case A) and over the pairs $`h_{e_n^{1\prime \prime }},h_{e_n^{2\prime \prime }}`$ in case B) in both cases in the order of increasing $`n`$. It is easy to see that the prescription in case B) is independent of the choice of breakpoint $`p`$ because the two integrals differ by a change of the order of a finite number of integrations which is irrelevant by properties of the measure $`\mu _0`$ and the compactness of $`G`$. Namely, all appearing functions are certainly absolutly integrable in any order and the assertion follows from Fubini’s theorem.
Steps I)-VI) provide the motivation for the following definition.
###### Definition 5.2
Let $`\gamma ,\gamma ^{}\mathrm{\Gamma }_\sigma ^\omega `$ and let $`f_\gamma ,g_\gamma ^{}`$ respectively be $`C_0`$-sequences over $`\gamma ,\gamma ^{}`$ respectively. Let $`\gamma ^{\prime \prime }=\gamma \gamma ^{}`$ and $`\stackrel{~}{\gamma }\gamma \gamma ^{}`$ be the piecewise analytic, $`\sigma `$-finite graph consisting of analytic edges $`\stackrel{~}{e}`$ which can be written simultaneously as the (countable) composition of edges of $`E(\gamma )`$ alone and of edges $`E(\gamma ^{})`$ alone. For $`\stackrel{~}{e}\stackrel{~}{\gamma }`$ we define
$`<_{f_\gamma },_{g_\gamma ^{}}>_{\stackrel{~}{e}}`$
$`:=`$ $`{\displaystyle _{G^{\left|\stackrel{~}{e}\right|}}}\left[_{e^{\prime \prime }E\left(\gamma ^{\prime \prime }\right)\stackrel{~}{e}}d\mu _H\left(h_{e^{\prime \prime }}\right)\right]\overline{\left[{\displaystyle \underset{eE\left(\gamma \right)\stackrel{~}{e}}{}}f_e\left({\displaystyle \underset{e^{\prime \prime }e}{}}h_{e^{\prime \prime }}\right)\right]}\left[{\displaystyle \underset{e^{}E\left(\gamma ^{}\right)\stackrel{~}{e}}{}}g_e^{}\left({\displaystyle \underset{e^{\prime \prime }e^{}}{}}h_{e^{\prime \prime }}\right)\right]`$
where $`e,e^{}`$ have been written as their decompositions over $`\gamma ^{\prime \prime }`$ and the order of integrations is defined in step VI) above. Then the scalar product between the $`C_0`$ vectors over $`\gamma ,\gamma ^{}`$ is defined by
$`<_{f_\gamma },_{g_\gamma ^{}}>`$
$`:=`$ $`[{\displaystyle \underset{eE\left(\gamma \right)\stackrel{~}{\gamma }}{}}<f_e,1>][{\displaystyle \underset{e^{}E\left(\gamma ^{}\right)\stackrel{~}{\gamma }}{}}<1,g_e^{}>][{\displaystyle \underset{\stackrel{~}{e}E\left(\stackrel{~}{\gamma }\right)}{}}<_{f_\gamma },_{g_\gamma ^{}}>_{\stackrel{~}{e}}]`$
where the separate convergence of the infinite products in the square brackets is in the sense of definition 4.1.
In order to define a scalar product on finite linear combinations of $`C_0`$ vectors over different $`\gamma `$’s we extend definition 5.2 by sesquilinearity. Notice that the definition reduces to the scalar product on $`_\gamma `$ if both vectors are finite linear combinations of $`C_0`$ vectors over $`\gamma `$.
Of course, in order to serve as a scalar product we must check that the scalar product is positive definite. However, this is obvious from the explicit measure theoretic expression (5.22) and can be verified by direct means as well.
###### Definition 5.3
The pre-Hilbert space of finite linear combinations of $`C_0`$ vectors over graphs $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$ completed in the scalar product (5.2) defines the Hilbert space $`^{}`$ of (5.13).
Why do we choose the Hilbert space $``$ of definition 5.3 as our quantum mechanical starting point ? The reason is the same as in the case of the original Hilbert space $`_{AL}`$ in which finite linear combinations of cylindrical functions over finite graphs were dense : the basic operators of the theory are still the same local operators as in section 2.2. They can be realized as operators on the infinite tensor product following the operator extension procedure of lemma 4.10 in section 4.2. Therefore, canonical commutation relations and adjointness relations are completely unchanged compared to the finite category.
### 5.2 Inductive Limit Structure
In the previous subsection we showed that any $`_\gamma ^{}`$ for $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$ can be obtained as the inductive limit of a sequence of Hilbert spaces $`_{\gamma _n}^{}`$ where $`\gamma _n\mathrm{\Gamma }_0^\omega `$. It is therefore natural to ask whether not all of $`^{}`$ arises in turn as the inductive limit of the $`_\gamma ^{}`$ for $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$.
The answer turns out to be negative, however, there is an inductive substructure which we now describe.
Notice that if either 1) any of $`\gamma \stackrel{~}{\gamma },\gamma ^{}\stackrel{~}{\gamma }`$ is an infinite graph or 2) any $`\stackrel{~}{e}\stackrel{~}{\gamma }`$ is a composition of an infinite number of edges of $`\gamma `$ or $`\gamma ^{}`$, or 3) the number of those $`\stackrel{~}{e}`$, which are compositions of more than one edge of $`\gamma `$ or $`\gamma ^{}`$ respectively, is infinite then almost always the expression (5.2) will vanish, simply again because the associative law does not hold on the ITP. It follows that if $`\gamma \gamma ^{}`$ but one of the three cases 1), 2), 3) just listed applies, a generic function $`f_\gamma _\gamma ^{}`$ cannot be written as a linear combination of functions $`f_\gamma ^{}_\gamma ^{}^{}`$.
This implies that, although $`\mathrm{\Gamma }_\sigma ^\omega `$ is a set directed by inclusion, we cannot simply define $`^{}`$ as the inductive limit of the $`_\gamma ^{}`$. To see this, notice that given $`\gamma ,\gamma ^{}\mathrm{\Gamma }_\sigma ^\omega `$ with $`\gamma \gamma ^{}`$ there is only one natural candidate for a unitary map $`\widehat{U}_{\gamma \gamma ^{}}:_\gamma _\gamma ^{}`$ :
For any $`eE\left(\gamma \right)`$, find its breakup $`e=e_1^{n_1}..e_N^{n_N},N\mathrm{}`$ into edges of $`\gamma ^{}`$ where $`n_k=\pm 1`$. We then consider the functon $`p_{\gamma \gamma ^{}}:\overline{𝒜}_\gamma \overline{𝒜}_\gamma ^{};h_eh_{e_1^{n_1}}..h_{e_N^{n_N}}`$ and then define
$$\widehat{U}_{\gamma \gamma ^{}}f_\gamma :=\left[p_{\gamma \gamma ^{}}^{}f_\gamma \right]\left[_{e^{}E\left(\gamma ^{}\right)\gamma }1\right]$$
(5.29)
This map is unitary when considered as a map from $`_\gamma ^{}`$ into $`^{}`$, with $`^{}`$ as defined in the previous section, since the extended Ashtekar Lewandowski measure is consistently defined. However, the right hand side of (5.29) will for a generic element $`f_\gamma _\gamma ^{}`$ simply not define an element of $`_\gamma ^{}^{}`$ for the reason already explained.
This state of affairs is in sharp contrast with the situation for the category $`\mathrm{\Gamma }_0^\omega `$ where the Hilbert space could indeed be written as the inductive limit of the various $`_\gamma _\gamma ^{}`$. For the category $`\mathrm{\Gamma }_\sigma ^\omega `$ the only way to define the Hilbert space structure is through (5.2).
The inductive limit still has a limited application in the following sense :
First, we define a new partial order $``$ on $`\mathrm{\Gamma }_\sigma ^\omega `$, motivated by the conditions 1), 2) and 3) at the beginning of this subsection.
###### Definition 5.4
For $`\gamma ,\gamma ^{}\mathrm{\Gamma }_\sigma ^\omega `$ we define $`\gamma \gamma ^{}`$ if and only if 1) $`\gamma \gamma ^{}`$ and 2) there exist disjoint (up to common vertices) unions $`\gamma =\stackrel{~}{\gamma }\gamma _1,\gamma ^{}=\stackrel{~}{\gamma }\gamma _1^{}`$ where $`\stackrel{~}{\gamma }\mathrm{\Gamma }_\sigma ^\omega `$ and $`\gamma _1,\gamma _1^{}\mathrm{\Gamma }_0^\omega `$.
It is important to notice that condition and 2) is not equivalent with 2’) that $`\gamma ^{}\gamma \mathrm{\Gamma }_0^\omega `$. This is because $`\gamma \gamma ^{}`$ only means that every edge $`eE\left(\gamma \right)`$ is a countable composition of edges $`e^{}E\left(\gamma \right)`$ and 2’) then does not exclude the existence of either a) at least one edge of $`\gamma `$ which is a countably infinite composition of edges of $`\gamma ^{}`$ or b) an infinite number of edges which are composed of at least two edges of $`\gamma ^{}`$. Both possiblities a) and b) are excluded by condition 2) which, in addition, implies 2’).
###### Lemma 5.3
The relation $``$ of definition 5.4 is a partial order.
Proof of Lemma 5.3 :
Only transitivity is nontrivial to prove. If $`\gamma \gamma ^{}\gamma ^{\prime \prime }`$ then first of all $`\gamma \gamma ^{}\gamma ^{\prime \prime }`$ so that $`\gamma \gamma ^{\prime \prime }`$. Secondly, if $`\gamma =\stackrel{~}{\gamma }\gamma _1,\gamma ^{}=\stackrel{~}{\gamma }\gamma _1^{}=\stackrel{~}{\gamma }^{}\gamma _2^{},\gamma ^{\prime \prime }=\stackrel{~}{\gamma }^{}\gamma _1^{\prime \prime }`$ are the corresponding disjoint unions with $`\stackrel{~}{\gamma },\stackrel{~}{\gamma }^{}\mathrm{\Gamma }_\sigma ^\omega `$ and $`\gamma _1,\gamma _1^{},\gamma _2^{},\gamma _1^{\prime \prime }\mathrm{\Gamma }_0^\omega `$ then we may define $`\stackrel{~}{\gamma }^{\prime \prime }:=\gamma ^{}\left(\gamma _1^{}\gamma _2^{}\right)=\stackrel{~}{\gamma }\gamma _2^{}=\stackrel{~}{\gamma }^{}\gamma _1^{}`$ which is obviously a subgraph of both $`\stackrel{~}{\gamma },\stackrel{~}{\gamma }^{}`$ and an element of $`\mathrm{\Gamma }_\sigma ^\omega `$ since $`\gamma _1^{}\gamma _2^{}\mathrm{\Gamma }_0^\omega `$ (recall that $`\mathrm{\Gamma }_0^\omega `$ is closed under unions due to piecewise analyticity and compact support of all its edges). It follows that there exist $`\gamma _2,\gamma _2^{\prime \prime }\mathrm{\Gamma }_0^\omega `$ such that $`\gamma =\stackrel{~}{\gamma }^{\prime \prime }\gamma _2`$ and $`\gamma ^{\prime \prime }=\stackrel{~}{\gamma }^{\prime \prime }\gamma _2^{\prime \prime }`$ are disjoint unions.
$`\mathrm{}`$
It is easy to see that $`\mathrm{\Gamma }_\sigma ^\omega `$ equipped with this partial order is not a directed set. This motivates to construct directed subsets.
###### Definition 5.5
Two graphs $`\gamma ,\gamma ^{}\mathrm{\Gamma }_\sigma ^\omega `$ are said to be finitely related, $`\gamma \gamma ^{}`$, provided that $`\gamma ,\gamma ^{}\gamma \gamma ^{}`$.
###### Lemma 5.4
Finite relatedness is an eqivalence relation.
Proof of Lemma 5.4 :
Reflexivity and symmetry are trivial to check. To see transitivity notice that $`\gamma ,\gamma ^{}\gamma \gamma ^{}`$ implies the existence of $`\stackrel{~}{\gamma },\stackrel{~}{\gamma }^{}\mathrm{\Gamma }_\sigma ^\omega `$ and of $`\gamma _1,\gamma _1^{\prime \prime },\gamma _1^{},\gamma _2^{\prime \prime }\mathrm{\Gamma }_\sigma ^\omega `$ such that we obtain disjoint unions $`\gamma =\stackrel{~}{\gamma }\gamma _1,\gamma ^{}=\stackrel{~}{\gamma }^{}\gamma _1^{},\gamma \gamma ^{}=\stackrel{~}{\gamma }\gamma _1^{\prime \prime }=\stackrel{~}{\gamma }^{}\gamma _2^{\prime \prime }`$. The last equality demonstrates that we may write the disjoint union $`\gamma \gamma ^{}=\stackrel{~}{\gamma }^{\prime \prime }\left(\gamma _1^{\prime \prime }\gamma _2^{\prime \prime }\right)`$ with $`\mathrm{\Gamma }_\sigma ^\omega \stackrel{~}{\gamma }^{\prime \prime }=\gamma \gamma ^{}(\gamma _1^{\prime \prime }\gamma _2^{\prime \prime })=\stackrel{~}{\gamma }\gamma _2^{\prime \prime }=\stackrel{~}{\gamma }^{}\gamma _1^{\prime \prime }`$. This in turn implies that we may actually write also disjoint unions $`\gamma =\stackrel{~}{\gamma }^{\prime \prime }\gamma _2,\gamma ^{}=\stackrel{~}{\gamma }^{\prime \prime }\gamma _2^{}`$ for some $`\gamma _2,\gamma _2^{}`$. In other words, $`\gamma \gamma ^{}`$ guarantees property 2) of definition 5.4. Transitivity now follows from the transitivity part of the proof of lemma 5.3.
$`\mathrm{}`$
We conclude that $`\mathrm{\Gamma }_\omega ^\sigma `$ decomposes into equivalence classes $`\left(\gamma _0\right)`$, called clusters and labelled by representants $`\gamma _0`$, called sources. Now, by construction, each cluster is directed by $``$. Moreover, since also by construction for any $`\gamma ,\gamma ^{}\left(\gamma _0\right)`$ the three conditions 1), 2) and 3) observed at the beginning of this subsection are not met, we find, in particular, that the operator (5.29) for $`\gamma \gamma ^{}`$ is now indeed a unitary operator which obviously satisfies the consistency condition $`\widehat{U}_{\gamma ^{}\gamma ^{\prime \prime }}\widehat{U}_{\gamma \gamma ^{}}=\widehat{U}_{\gamma \gamma ^{\prime \prime }}`$ for $`\gamma \gamma ^{}\gamma ^{\prime \prime }`$. The general results of section 4.3 now reveal the existence of the inductive limit Hilbert space $`_{\left(\gamma _0\right)}^{}`$ for any cluster $`\left(\gamma _0\right)`$ and corresponding unitarities $`\widehat{U}_\gamma :_\gamma ^{}_{\left(\gamma _0\right)}^{}`$ for any $`\gamma \left(\gamma _0\right)`$ such that $`\widehat{U}_\gamma =\widehat{U}_\gamma ^{}\widehat{U}_{\gamma \gamma ^{}}`$.
It would be a very pretty result if one could establish that the Hilbert spaces $`_{\left(\gamma _0\right)}`$ corresponding to different clusters are mutually orthogonal with respect to (5.2). But this is certainly not the case, just take any $`\gamma \left(\gamma _0\right)\left(\gamma _0^{}\right)\gamma ^{}`$ and consider the $`C_0`$-vectors $`f_\gamma :=_{eE\left(\gamma \right)}1`$ and $`f_\gamma ^{}:=_{e^{}E\left(\gamma ^{}\right)}1`$. Then trivially $`<_{f_\gamma },_{f_\gamma ^{}}>=1`$.
Denote by $`𝒞_\sigma ^\omega `$ the set of clusters in $`\mathrm{\Gamma }_\sigma ^\omega `$. Then we have the following equivalent definition of the full Hilbert space
$$^{}=\overline{_{\left(\gamma _0\right)𝒞_\sigma ^\omega }_{\left(\gamma _0\right)}^{}}$$
(5.30)
which displays it as a kind of cluster decomposition. The decomposition is, however, not a direct sum decomposition. Obviously, the cluster Hilbert spaces $`_{\left(\gamma _0\right)}^{}`$ are mutually isomorphic and therefore in particular isomorphic with the original Ashtekar Lewandowski Hilbert space $`_{AL}_{\left(\mathrm{}\right)}^{}`$ based on $`\mathrm{\Gamma }_0^\omega `$ obtained by choosing as the source the empty graph. The fact that the cluster Hilbert spaces can be written as inductive limits is then not any more surprising because they are isomorphic with $`_{AL}`$ of which we knew already that it is an inductive limit.
#### 5.2.1 Rigging Triple Structures
Finally we can equip the space Cyl with a topology in analogy with the one defined in , the difference coming from the fact that we do not know an explicit orthonormal basis for $``$.
###### Definition 5.6
Choose for any $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$ once and for all a von Neumann basis $`T_{\gamma s\beta }`$ over $`\gamma `$ where $`s𝒮_\gamma `$ runs through the strong equivalence classes in $`_\gamma `$ and $`\beta `$ through the set of functions $`_\gamma `$ defined in (4.10). Let $`f`$Cyl be a cylindrical function.
i)
The family of Fourier semi-norms of $`f`$ is defined by
$$\left|\right|\left|f\right|||_\gamma :=\underset{s\beta }{}|<T_{s\beta },f>|$$
(5.31)
where the inner product in (5.31) is defined by (5.2). Notice that indeed $`|f+g|_\gamma |f|_\gamma +|g|_\gamma ,|zf|_\gamma =|z||f|_\gamma `$ for all $`\gamma \mathrm{\Gamma }_\sigma ^\omega ,z\text{ }\mathrm{C},f,g`$Cyl. Obviously the family separates the points of Cyl since $`|f|_\gamma <\mathrm{}`$ for all $`\gamma `$ implies $`f`$ and $`|f|_\gamma =0`$ for all $`\gamma `$ implies $`f=0`$, that is, $`f`$ is the zero $`C_0`$ sequence and so $`f=0`$, the zero $`C`$ function in Cyl.
ii)
Consider the subspace $`\mathrm{\Phi }`$ of Cyl consisting of elements which are finite linear combinations of $`C`$ functions $`f`$ with the property that
$$\left|f\right|:=\underset{\gamma }{sup}\left|f\right|_\gamma <\mathrm{}$$
(5.32)
iii)
Item i) displays $`\mathrm{\Phi }`$ as a locally convex vector space. Upon equipping it with the natural topology (the weakest topology such that all the $`|||.|||_\gamma `$ and addition are continuous) it becomes a topological vector space $`\mathrm{\Phi }`$.
An alternative choice for a topology for $`\mathrm{\Phi }`$, upon which it would become a normed (but not necessarily complete, Banach) topological vector space, is given by the norm (5.32). The natural topology defined in iii) is not generated by a countable set of seminorms, therefore it is not metrizable and (upon completion) cannot be a Fréchet space. On the other hand, the norm $`\left|\right||.|\left|\right|`$ certainly defines a metric. The two topologies are therefore not equivalent, clearly the natural topology is weaker than the norm topology.
###### Lemma 5.5
We have $`\mathrm{\Phi }`$ and $`\mathrm{\Phi }`$ is dense in $``$.
Proof of Lemma 5.5 :
To see this, consider an arbitrary element $`f=_{n=1}^Nf_{\gamma _n}`$ where $`f_{\gamma _n}`$ is a $`C`$ function over $`\gamma _n`$in $`\mathrm{\Phi }`$ and $`N<\mathrm{}`$. Therefore, $`\left|f_{\gamma _n}\right|_\gamma <\mathrm{}`$ for any $`\gamma ,n=1,..,N`$. In particular, $`\left|f_{\gamma _n}\right|_{\gamma _n}<\mathrm{}`$ for any $`\gamma `$. Thus,
$$\left|\right|\left|f_{\gamma _n}\right|||_{\gamma _n}=\underset{s,\beta }{}|<T_{\gamma s\beta },f_{\gamma _n}>|<\mathrm{}$$
and therefore
$$\left|\right|f_{\gamma _n}||^2=\underset{s,\beta }{}|<T_{\gamma s\beta },f_{\gamma _n}>|^2<\left(\right|\left|\right|f_{\gamma _n}\left|\right||_{\gamma _n})^2<\mathrm{}$$
from which we see that $`f_{\gamma _n}`$ is a $`C_0`$ vector. It follows that $`f`$ is a finite linear combinations of $`C_0`$ vectors and thus an element of $``$. As the finite linear combinations of $`C_0`$ vectors form a dense subset of $``$ which we just showed to be contained in $`\mathrm{\Phi }`$, we conclude that $`\mathrm{\Phi }`$ is dense in $``$.
$`\mathrm{}`$
With the natural topology on $`\mathrm{\Phi }`$ we are equipped with the rigging triple $`\mathrm{\Phi }\mathrm{\Phi }^{}`$ and can take over the framework of to solve constraints also in the context of the ITP.
### 5.3 Contact with Semiclassical Analysis
In this subsection we will make contact with the semi-classical states of section 3.
Let $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$ be an infinite graph, filling all of $`\mathrm{\Sigma }`$ arbitrarily densely (in the absence of a background metric by this we mean simply that for an arbitrary choice of neighbourhoods of each point of $`\mathrm{\Sigma }`$, $`\gamma `$ can be chosen to intersect all of them). Suppose we are given a solution of the classical field equations (say the Einstein equations in the absence of matter or the Einstein-Yang-Mills equations in the presence of matter, in the latter case $`G`$ is the direct product of the gravitational $`SU\left(2\right)`$ with the $`SU\left(3\right)\times SU\left(2\right)\times U\left(1\right)`$ of the standard model), that is, for each time slice $`\mathrm{\Sigma }_t,t\text{ }\mathrm{R}`$ we have an initial data set $`(A_t^0\left(x\right),E_t^0\left(x\right)),x\mathrm{\Sigma }`$ satisfying the field equations and in particular the constraint equations. Moreover, we will have to choose a certain gauge to write down the solution explicitly. Then, with the techniques of section 3, for each edge $`eE\left(\gamma \right)`$ and given classicality parameter $`s`$ we obtain a normalized coherent state
$$\xi _{g_e^t(A^0,E^0)}^s:=\frac{\psi _{g_e^t(A^0,E^0)}^s}{\psi _{g_e^t(A^0,E^0)}^s__e}$$
(5.33)
where $`g_e^t\left((A^0,E^0)\right)=\mathrm{exp}\left(i\tau _jP_j^e(E_t^0,A_t^0)/\left(2a^2\right)\right)h_e\left(A_t^0\right)`$ for pure general relativity. Finally, we consider the $`C_0`$-vector over $`\gamma `$ given by
$$\xi _{\gamma ,(A_t^0,E_t^0)}^s:=_{eE\left(\gamma \right)}\xi _{g_e^t(A^0,E^0)}^s$$
(5.34)
These states comprise a preferred set of coherent states over the infinite graph $`\gamma `$ and provide the basic tool with which to address the following list of fascinating physical problems :
* Given one and the same graph $`\gamma `$ and classicality parameter $`s`$, when are the strong and weak equivalence classes of the states (5.34) equal to each other ? What is the physical significance of strong and weak equivalence anyway ? From experience with model systems one expects that two different weak equivalence classes correspond to drastically different physical situations such as an infinite difference in ground state energies or topologically different situations while the general analysis of section 4.1 (lemma 4.14) states that two incomplete ITP’s corresponding to different strong equivalence classes within the same weak one are unitarily equivalent.
Of course, different topological situations can be described within the same complete ITP only if we get rid off the embedding spacetime that one classically started with. One way to do this would be roughly as follows :
Consider $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$ not as an embedded graph but merely as a countable, combinatorical one. A countable combinatorical graph is simply a countable collection of edges (which are analytic curves when embedded into any given $`\mathrm{\Sigma }`$), and vertices together with its connectivity relations, that is, information telling us at which vertices a given edge ends. Now recall that the spectrum $`\text{spec}_\mathrm{\Sigma }\left(\widehat{O}\right)`$ of important operators $`\widehat{O}`$ of the theory such as the area operator (see ), as obtained on the Hilbert spaces corresponding to graphs embedded in a concrete $`\mathrm{\Sigma }`$, depends on the topology of $`\mathrm{\Sigma }`$ and it should be true that a complete set of operators encodes full information about the topology of $`\mathrm{\Sigma }`$ via the range of their respective spectra. We can now define the universal operator $`\widehat{O}`$ acting on Hilbert spaces over combinatorical graphs by a new kind of summming over topologies, namely, one allows the spectrum of $`\widehat{O}`$ to take all possible values, that is, $`\text{spec}\left(\widehat{O}\right)=_\mathrm{\Sigma }\text{spec}_\mathrm{\Sigma }\left(\widehat{O}\right)`$. One would then say that a given closed subspace of the Hilbert space, carrying a representation of the operator algebra describes a concrete topology of $`\mathrm{\Sigma }`$ provided the spectra of the operators $`\widehat{O}`$ resricted to that subspace are compatible with th spectra $`\text{spec}_\mathrm{\Sigma }\left(\widehat{O}\right)`$.
Now the Infinite Tensor Product Hilbert space as obtained from combinatorical graphs space comes in as follows. It is expected that closed $`\left[f\right]`$-adic subspaces of that ITP corresponding to different topologies of $`\mathrm{\Sigma }`$ in the way just described also correspond to strong equivalence classes within different weak equivalence classes. Now while elementary operators of the theory will leave these subspaces invariant, there are densely defined operators on the complete ITP which mediate between the two. Thus, the ITP might be used to describe topology change in Quantum General Relativity and would then wipe away one of the main criticisms directed towards the whole programme.
A related interesting question is, whether classical states ($`C_0`$-vectors) corresponding to Minkowski and Kruskal spacetime respectively are orthogonal, more generally, whether one can superimpose classical states corresponding to globally different spacetimes within the same strong equivalence class. Interestingly, all this can be analyzed by performing relatively straightforward calculations of the type outlined in .
* Given one and the same graph $`\gamma `$ and solution $`(A_t^0,E_t^0)`$, are the $`C_0`$ vectors (5.34) for different values of $`s`$ in the same weak equivalence class for non-compact $`\mathrm{\Sigma }`$ ? Since the parameter $`s`$ plays a role similar to a mass parameter in free scalar field theory on Minkowski space one might expect this not to be the case as the Fock representations over Minkowski space with different mass are not unitarily equivalent. Indeed, it is easy to see that for Minkowski space in the gauge $`A_a^j=0,E_j^a=\delta _j^a`$ we have $`<\xi _e^s,\xi _e^s^{}>=|<\xi _e^s,\xi _e^s^{}>|=q<1`$ for $`ss^{}`$ so that we obtain different weak eqivalence classes.
* Given one and the same $`(A_t^0,E_t^0)`$ and $`s`$, what happens under refinements of the graph $`\gamma `$ ? Again, since under refinements of an infinite graph we perform an infinite change on the graph, from expression (5.2) we expect the corresponding $`C_0`$ vectors to be orthogonal. This turns out to be correct.
* Given a $`C_0`$ vector $`\xi `$ of the type (5.34) we know from lemma 4.8 that any other vector in the $`\left[\xi \right]`$-adic incomplete ITP can be obtained as a (Cauchy sequence of) linear combinations of $`C_0`$ vectors each of which differs from $`\xi `$ in only a finite number of entries $`eE\left(\gamma \right)`$. Now consider a deformation of a classical solution $`(A_t^0,E_t^0)`$, say the Minkowski metric plus a graviton or photon. A plane wave graviton is everywhere excited over $`\mathrm{\Sigma }`$, that is, differs everywhere significantly from the Minkowski background, and therefore will not lie in the closure of the states just described. In the case of the electromagnetic field this is expected because plane wave solutions have infinite energy. But the anyway more physical graviton wave packets, although also everywhere excited, are Gaussian damped and thus have a chance to lie in that $`\left[\xi \right]`$-adic strong equivalence class of Minkowski space (or any other background). We expect that there is a unitary map between the usual Fock space description of gravitons and the coherent state Gaussian wave packet gravitons of the present framework. If that turns out to be correct, we can also describe Einstein-Maxwell-Theory this way and consider photons propagating on quantum spacetimes. These issues will be examined in .
The same analysis can, of course, be performed on any background and this is the way we will try to describe the Hawking effect in this approach .
* Related to this is the question if we can recover the spectacular results of Quantum Field Theory on Curved Backgrounds anyway . The philosophy of that approach is that if the backreaction of matter to geometry can be safely neglected then treating the metric as a given, classical background should be a good approximation to the physics of the system. Recently, there has been a quantum leap in this field of research due to a precise formulation of the microlocal spectrum condition on arbitrary, globally hyperbolic but not necessarily stationary backgrounds.
It is to be expected, and an important consistency check, that to zeroth order in the Planck length the full quantum gravity calculation should agree with the predictions of Quantum Field Theory on curved backgrounds, in other words, QFT on curved backgrounds is a semi-classical limit of quantum gravity. The way to check this expectation is of course the following : The total Hilbert space of the system matter plus geometry is the tensor product of the Hilbert spaces for the matter sector and the gravity sector respectively. Given a classical background metric, we will choose states $`\psi _{total}`$ which are tensor products of arbitrary states $`\psi _{matter}`$ from the matter Hilbert space with one fixed state $`\psi _{grav}^0`$ from the gravity Hilbert space, namely the coherent $`C_0`$ vector for the metric to be approximated, symbolically $`\psi _{total}=\psi _{matter}\psi _{grav}^0`$. The matter Hamiltonian operator of, say, bosonic matter coupled to quantum gravity is roughly a linear combination of operators of the form $`\widehat{H}_{total}=\widehat{A}_{matter}\widehat{A}_{grav}`$. Thus we find for the matrix elements of that operator
$$<\psi _{total},\widehat{H}_{total}\psi _{total}^{}>_{total}=<\psi _{matter},\widehat{A}_{matter}\psi _{matter}^{}>_{matter}<\psi _{grav}^0,\widehat{A}_{grav}\psi _{grav}^0>_{grav}$$
which shows that we obtain an effective matter Hamiltonian given by $`\widehat{H}_{total}^{eff}=\widehat{A}_{matter}<\psi _{grav}^0,\widehat{A}_{grav}\psi _{grav}^0>_{grav}`$ which by the properties of the operator of is finite ! The corrections to the classical background metric are of course contained in the difference $`<\psi _{grav}^0,\widehat{A}_{grav}\psi _{grav}^0>_{grav}A_{grav}^{class}`$ where $`A_{grav}^{class}`$ is the classical limit of the gravity operator evaluated for the given classical background. This quantity is certainly at least of order $`\mathrm{}_p`$. However, this is not the only correction to Quantum Field Theory on curved backgrounds. A second correction comes from the fact that our theory is non-perturbative which means, in particular, that the matter Hilbert space is not the usual perturbative Fock space. Therefore it is not at all obvious that the spectrum of the operator $`\widehat{A}_{matter}`$ on the non-perturbative Hilbert space coincides with the spectrum of the usual matter Hamiltonian on the usual perturbative Fock space. The spectra better agree, at least modulo corrections of order at least $`\mathrm{}_p`$, in order that we can claim to have quantized a theory which has general relativity plus standard matter as the classical limit.
* A first application of this procedure to discover new physical effects due to quantum gravity is the exact treatment of the so-called $`\gamma `$-ray-burst effect which we will do in .
To date the exact astrophysical explanation or source for high energetic $`\gamma `$-photons (up to TeV !) is unclear but what is important for us is that these photons were created billions of years ago, they can come from distances comparable to the Hubble radius. The idea is that these photons on their way to us constantly are influenced by the vacuum fluctuations of the gravitational field and although the influence is very, very, very small, it can accumulate due to the long travelling time of the photons. Now the higher energetic the photon, the more it should probe the small scale discreteness of quantum geometry and we thus expect an energy-dependent dispersion law. The dispersion law being energy and therefore (Minkowski) frame dependent, it violates Poincaré invariance. The effect therefore cannot come from any perturbative theory (interacting QFT on Minkowski space, perturbative quantum (super)gravity, perturbative string theory) all of whose $`S`$-matrix elements or $`n`$-point functions are by definition (or (Wightman) axiom) Poincaré invariant. For observational purposes it is convenient that the intensity peak of the burst can have a time width as small as of the order of 1ms. The idea is then to calibrate the detector to detect events at energies $`E_2>E_1`$ at times $`t_2>t_1`$ due to the energy dependence of the speed of light which according to is speculated to be of the form $`c\left(E\right)/c\left(0\right)=1k\left(E/E_{eff}\right)^\alpha `$. Here $`k`$ is a coefficient of the order of $`1`$, $`E_{eff}`$ is the effective quantum gravity scale and $`\alpha `$ is a power which is hopefully of the order $`1`$ for the effect to be detectable. For $`\alpha =1`$ one finds $`t_2t_1=k\left[\left(E_2E_1\right)/E_{eff}\right]\left[L/c\left(0\right)\right]`$ where $`L`$ is the distance of the source inside a galaxy and so can be determined from its redshift. Inserting the numbers for a burst which is a billion lightyears away and $`E_2E_1=1`$TeV we get $`t_2t_1`$ of the order of a second (!) if we set $`E_{eff}=m_p`$ which is large enough compared to the width of the signal. Thus, for $`\alpha =1`$ the effect could be indeed observable, say by a Čerenkov observatory (but not for $`\alpha =2`$), at least in principle, however, experimentally it is a highly non-trivial task to take into account all possible errors (dark matter, gravitational lensing, dust, atmosphere, …) and to make sure that the measured intensities really came from the same burst.
Our aim in will be to compute $`k`$ and $`\alpha `$, or more generally, the precise dispersion law, exactly along the lines outlined in item v). It is important to realize, that the effect is an inevitable theoretical prediction of quantum general relativity in the present formulation due to the Heisenberg Uncertainty Obstruction. Namely, the quantum metric depends on magnetic (connection $`A`$) and electric (conjugate momentum $`E`$) degrees of freedom which upon quantization become noncommutative operators as we have seen in section 2.2 and therefore cannot be simultaneously diagonalized. Thus, the best we can do is to write down a best approximation eigenstate of the metric operator, that is, a coherent state which saturates the Heisenberg uncertainty bound. As we showed in , our states of section 3 have precisely these semiclassical properties. However, while a best approximation state, it is not an exact eigenstate and thus cannot be Poincaré invariant.
It is also important to see that our analysis is more ambitious than the pioneering work for three reasons : First of all, instead of coherent states only weaves were used, however, these approximate only half of the degrees of freedom and are more similar to momentum eigenstates than semiclassical states. Secondly, the matter field was treated classically and one was computing only the dispersion law coming from the changed d’Alembert operator, an option which we also have. Thirdly, in contrast to our coherent states, the weave with the assumed semi-classical behaviour was not proved to exist as a normalizable state of the Hilbert space.
* The results of are a small indication that quantum gravity plus quantum matter combine to a finite quantum field theory. An elementary particle physicist who computes Feynman diagrammes and has to renormalize divergent quantities all the time will rightfully ask what happened to the ultraviolet divergencies of his everyday life. A short answer seems to be, that in a diffeomorphism invariant, background independent theory there is no room for UV divergencies since there is no difference between “large” and “small” distances, the renormalization group gets “absorbed” into the diffeomorphism group. While plausible, to the best of our knowledge nobody has so far investigated these speculations in detail. This will be the topic of .
* The result of shows that the geometry and matter Hamiltonian (constraints) are densely defined operators on the unextended Ashtekar-Lewandowski Hilbert space. How does the situation change with the huge extension of the Hilbert space performed in section 5.1 ? The answer, investigated in detail in is that, not unexpectedly, these operators are not densely defined on all of $``$, however, they are on physical $`\left[f\right]`$-adic Hilbert spaces for $`f`$ a coherent $`C_0`$-vector. Here we call a coherent $`C_0`$-vector physical if it is labelled by classical field configuration that obeys the fall-off conditions at spatial infinity. To see how this roughly comes about suppose we have a coherent $`C_0`$ vector $`_f=_ef_e`$ over some infinite graph $`\gamma `$. The Hamiltonian constraint operators applied to $`_f`$ are of the form $`\widehat{H}_f=\left[_v\widehat{H}_v\right]_f`$ where the sum is over all vertices of $`\gamma `$ and $`\widehat{H}_v`$ influences only those $`f_e`$ for which its edge $`e`$ is incident at $`v`$, that is, it is a local operator. Our first observation is that therefore $`\widehat{H}_v_f_{\left[f\right]}`$ as guaranteed again by lemma 4.8. So $`\widehat{H}_f`$ is a countably infinite sum of vectors in $`_{\left[f\right]}`$ and the question is whether it is convergent, that is, whether $`\left|\right|\widehat{H}_f||^2<\mathrm{}`$. Suppose now that $`f_e=\xi _{g_e^t(A^0,E^0)}^s`$ and $`_f`$ is a coherent $`C_0`$ vector peaked at $`A_t^0,E_t^0`$. Then we have by the Ehrenfest theorem of
$$\left|\right|\widehat{H}_f||^2=<_f,\widehat{H}^{}\widehat{H}_f>=|\underset{v}{}H_v(A_t^0,E_t^0)|^2[1+O\left(s\right)]$$
where $`H_v(A_t^0,E_t^0)`$ is by construction a dicretization of an integral over a small region in $`\mathrm{\Sigma }_t`$ of the classical Hamiltonian density $`H(A_t^0,E_t^0)`$ and the sum over vertices is a Riemann sum for the integral $`_{\mathrm{\Sigma }_t}d^3xH(A_t^0,E_t^0)`$. It follows that the norm exists if and only if the field configuration $`A_t^0,E_t^0`$ satisfies the fall-off conditions, that is, if it is a point in the classical phase space (no constraints or field equations yet being imposed). Next, again from lemma 4.8 we see that for such coherent $`C_0`$ sequences $`f`$ the Hamiltonian constraint is densely defined on the $`\left[f\right]`$-adic Hilbert space because a dense domain is given by the vector space of finite linear combinations of $`C_0`$-vectors differing from $`_f`$ in at most a finite number of entries $`f_e`$ and the convergence proof just outlined certainly goes through since the finite number of changes made affect a finite number of vertices only and these are of $`d^3x`$ measure zero in the limit of infinitely fine graphs.
This result is rather pretty because it tells us that the classical theory still has some effect on the quantum theory, not every $`[.]`$-adic incomplete ITP carries a representation of the operator algebra. The set of $`C_0`$ vectors whose $`[.]`$-adic ITP’s do carry a representation includes the physical coherent $`C_0`$ vectors but excludes the non-physical ones. Since coherent states form an overcomplete basis on the complete ITP, this statement seems to give a rather complete classification of $`[.]`$-adic ITP’s carrying a representation.
However, there are also other $`[.]`$-adic ITP’s which are not of that form : An example is provided by the Ashtekar-Lewandowski state $`\omega _{AL}`$ which on the complete ITP $`^{}`$ is given by the GNS cyclic vector $`\mathrm{\Omega }_{AL}=_{eC_{\mathrm{}}}1`$ where as before $`C_{\mathrm{}}`$ denotes a supergraph. Now it is easy to check, using the overcompleteness of Hall’s coherent states, that $`1=1_e=_{G^{\text{ }\mathrm{C}}}𝑑\nu _s\left(g_e\right)\psi _{g_e}^s`$ where $`\nu _t`$ is Hall’s measure , that is, the Ashtekar-Lewandowski $`C_0`$-vector is an infinite-fold superposition of all coherent states formally given by the (kind of functional integral)
$$\mathrm{\Omega }_{AL}=\left[_ed\nu _s\left(g_e\right)\right]_e\psi _{g_e}^s$$
and so includes non-physical coherent $`C_0`$-vectors which, however, come with the appropriate weight enabling it to carry a representation of the observable algebra. This way to write $`\mathrm{\Omega }_{AL}`$ makes it obvious that it is not peaked on a particular metic at all although it is annihilated by all momentum operators which might lead one to assume that $`\mathrm{\Omega }_{AL}`$ approximates the zero metric ! This is clearly not the case.
Another way to write $`\mathrm{\Omega }_{AL}`$ is $`\mathrm{\Omega }_{AL}=lim_s\mathrm{}_{eC_{\mathrm{}}}\psi _{g_e}^s`$ which displays it as a $`C_0`$-vector, but only in the anticlassical limit $`s\mathrm{}`$ ! Both ways to write $`\mathrm{\Omega }_{AL}`$ reaffirm one more time the impression that in the spatially non-compact case the $`\left[1\right]`$-adic incomplete ITP, or, in other words, the original, unextended Ashtekar-Lewandowski Hilbert space is a pure quantum representation of the observable algebra which seems to have no obvious semi-classical correspondence. All the solutions found in are (diffeomorphism invariant versions of) states in that Hilbert space and thus have presumably no (semi)classical relevance, as speculated by many, it is indeed correct that not only infinite linear combinations but indeed infinite products are necessary to catch the physically relevant sector of the complete ITP. This could also explain the discrepancy with respect to the number of degrees of freedom in 2+1 gravity pointed out in : There the $`\left[1\right]`$-adic incomplete ITP was used and gave rise to an infinite number of physical states each of which describes zero volume almost everywhere. As in that case $`\mathrm{\Sigma }`$ was assumed to be compact, the $`\left[1\right]`$-adic incomplete ITP coincides in fact with $`^{}`$, however, the true physical states are only obtained as infinite linear combinations of zero volume states which presumably builds up an $`\left[f\right]`$-adic incomplete ITP describing an everywhere non-degenerate metric. We will come back to this in a future publication .
* A heuristic method to use the coherent states in order to derive a Hamiltonian constraint operator with the correct classical limit is as follows : We choose a point $`(A_{t=0}^0\left(x\right),E_{t=0}^0\left(x\right)),x\mathrm{\Sigma }`$ on the constraint surface of the phase space of, say, general relativity in some gauge (that is, a field configuration on $`\mathrm{\Sigma }_0`$) and obtain its trajectory under the Einstein evolution to first order in the time parameter $`t`$ as $`A_t^0\left(x\right)=A_0^0\left(x\right)+t\{H_0\left(N\right),A_0^0\left(x\right)\}`$ (that is, a field configuration on $`\mathrm{\Sigma }_t`$) and likewise for $`E_t\left(x\right)`$. Here, $`H_0\left(N\right)`$ is the Hamiltonian constraint on $`\mathrm{\Sigma }_0`$. We thus obtain coherent states as in (5.34) over any $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$ and now define a family of operators $`\widehat{H}_\gamma \left(N\right)`$, each of which is densely defined on the $`[.]`$-adic incomplete ITP corresponding to the strong equivalence class of the $`C_0`$-sequence associated with the coherent $`C_0`$-vector (5.34) by the definition
$$\widehat{H}_\gamma \left(N\right)\xi _{\gamma ,(A_0^0,E_0^0)}^s:=\frac{1}{i}\left(\frac{d}{dt}\right)_{t=0}\xi _{\gamma ,(A_t^0,E_t^0)}^s$$
(5.35)
Notice that time evolution preserves the kinematical constraints of the phase space, in particular, the fall-off conditions and so the strong equivalence classes of the $`C_0`$ sequences defined by the $`C_0`$-vectors $`\xi _{\gamma ,(A_0^0,E_0^0)}^s,\xi _{\gamma ,(A_t^0,E_t^0)}^s`$ are equal to each other. Since the map $`t(A_t^0,E_t^0)`$ is smooth on the continuum phase space $`M`$, it induces a smooth map on the subset $`\overline{M}_{\gamma |M}`$ of the graph phase space of section 2.1. Therefore, $`t\xi _{\gamma ,(A_t^0,E_t^0)}^s`$ is strongly continuous since $`\xi _{\gamma ,(A_t^0,E_t^0)}^s\xi _{\gamma ,(A_0^0,E_0^0)}^s`$ depends smoothly on $`\overline{M}_{\gamma |M}`$. If we could verify also that $`\xi _{\gamma ,(A_t^0,E_t^0)}^s=\xi _{\gamma ,(A_0^0,E_0^0)}^s`$ up to terms of order $`t^2`$ then time evolution would be given to first order in $`t`$ by a one-parameter group of unitary operators and the existence of $`\widehat{H}_\gamma \left(N\right)`$ would follow from Stone’s theorem . However, due to the complicated classical constraint algebra $`\{H\left(N\right),H\left(M\right)\}=d^3x\left(MN_{,a}M_{,a}N\right)q^{ab}V_b0`$, where $`V_b`$ is the infinitesimal generator of the vector constraint, the quantum evolution is better not unitary in order that it has the correct classical limit. Therefore, Stone’s theorem will not apply and if we interprete (5.35) as a strong limit then it might be ill-defined.
* A great surprise of quantum theory was the resolution of a classical paradoxon, the explanation for the stability of atoms. According to classical electrodynamics the electrons orbiting the nucleus should emit radiation and fall into it after a finite time. The discreteness of the bound state energy spectrum bounded from below prevents this from happening and displays a mechanism for the avoidance of a classical singularity. It is an interesting speculation that something like this could also happen in quantum gravity, that the classical singularities predicted by the singularity theorems of Hawking and Penrose are actually absent in quantum gravity, providing a resolution of the information paradox. This question can also be naturally adressed within the framework of coherent states : given a classical black hole spacetime with its singularity, say the Kruskal spacetime, we could compute the expectation value, with respect to a coherent state for that black hole spacetime, of an operator whose classical counterpart becomes singular there. If the singularity is quantum mechanically absent, then the operator should be bounded from above and the expectation value should be finite. From the point of view of the Bargmann-Segal Hilbert space, the coherent state is peaked at the singular spacetime but there is a non-zero probability to be away from it just in the right way to be square integrable. This is in analogy with the eigenstates of the electron energy operator of the hydrogenium atom whose probability density at the origin is finite and which are also square integrable (notice that coherent states are approximate eigenstates of any operator). More generally, one would like to treat quantum black holes with the new semi-classical input provided by coherent states which come out of the quantum theory and are not a purely classical input such as classically encoding the presence of an isolated horizon into the topology of $`\mathrm{\Sigma }`$, thus inducing corresponding boundary conditions, in quantum general relativity or the identification of classical supergravity black hole solutions with D-brane configurations protected against quantum corrections due to the BPS nature of the corresponding states in string theory (see, e.g. and references therein). These and related questions will be examined in .
* The coherent state framework of so far is worked out in full detail only for the compact groups of rank one and direct products of those. As argued there, the extension to groups of higher rank should be straightforward given the strategy for $`U\left(1\right),SU\left(2\right)`$ but it is yet a lot of work. It would be important to give the full details at least for the physically important case of $`SU\left(3\right)`$. The analysis will be given in .
* The exposition in section 4.2 underlines the relevance of von Neumann algebras and operator theory for the Infinite Tensor Product. This provides a pretty interface with the methods of Algebraic Quantum Field Theory. In particular, it would be interesting to work out the Tomita Takesaki Theory for the appearing operator algebras as it was done for scalar field theory on Minkowski space for diamond regions (the Bisognano Wichmann theorem) where the challenge in our context is that we have only a spacetime background topology and differentiable structure but not a spacetime background metric. Modular theory is the basic tool to determine the precise type of the type III factors which from experience with scalar field theory seem to be the most relevant types of factors in quantum field theory and will be studied in .
At this point the careful reader will wonder how we can apply the theory outlined in section (4.2) which was geared only at bounded operators. However, our basic operators are $`\widehat{h}_e,\widehat{p}_i^e,eE\left(\gamma \right)`$ and the latter is unbounded although essentially self-adjoint on $`_e`$, a property which trivially extends to the ITP. This mismatch can be cured by considering instead of $`\widehat{p}_i^e`$ the Weyl kind of operator
$$\widehat{H}_e:=e^{\tau _j\widehat{p}_j^e/2}$$
(5.36)
which takes values in the set of group valued bounded operators and transforms as $`\widehat{H}_e\text{Ad}_{g_e\left(0\right)}\widehat{H}_e`$ under gauge transformations. Its boundedness, for instance for $`G=SU\left(2\right)`$, is evident from the formula
$$\widehat{H}_e=\frac{e^{is\mathrm{\Delta }_e/2}\widehat{h}_ee^{is\mathrm{\Delta }_e/2}\widehat{\left(h_e\right)^1}}{e^{is}}$$
(5.37)
which can be proved from a similar formula in the first reference of by analytical continuation of the classicality parameter $`s`$. Here, $`\mathrm{\Delta }_e`$ is the Laplacian on the copy of $`G`$ corresponding to $`h_e`$. Formula (5.37) displays $`\widehat{H}_e`$ as a product of four bounded operators. Taking the operator adjoint of (5.37) one finds that the operators $`\widehat{H}_e,\widehat{h}_e`$ obey the kind of Weyl algebra
$$\widehat{H}_e\widehat{h}_e=e^{2is}ϵ\left(\widehat{h}_{e^1}\widehat{H}_{e^1}\right)^T$$
(5.38)
where $`(.)^T`$ denotes transposition, $`ϵ`$ is the skew symmetric spinor of second rank of unit determinant and $`\widehat{H}_{e^1}=\left(\widehat{H}_e\right)^1=\left(\left(\widehat{H}_e\right)^{}\right)^T`$ (matrix and operator inverse but only operator adjoint) clarifies the adjointness relations. Notice that one could also consider the objects $`\widehat{H}_e^j:=e^{i\widehat{p}_j^e}`$ which satisfy the simpler Weyl algebra
$$\widehat{H}_e^j\widehat{h}_e=e^{s\tau _j/2}\widehat{h}_e\widehat{H}_e^jϵ^1$$
but the $`\widehat{H}_e^j`$ do not transform covariantly under gauge transformations.
Although it is slightly inconvenient to work with $`\widehat{H}_e`$ in place of $`\widehat{p}_e^j`$ since physical operators are more easily expressed in terms of the latter, it can be done and will be useful to prove abstract theorems. In that respect it is worthwhile mentioning that one could also try to work directly with the unbounded operators but the general theory for this does not yet exist due to complications associated with the fact that domains do not interact with the linear structure, efforts have, to the best of our knowledge, so far been restricted to the discussion of essential self-adjointness of (infinite) sums of operators restricted to one $`\left[f\right]`$-adic Hilbert space, where $`\left[f\right]`$ is a strong eqivalence class .
* The present framework could also be employed to make contact with the so-called spin foam models . These are a class of state sum models including the one that one obtains by studying the transition amplitudes associated with the Hamiltonian constraint constructed in and which since then have attracted a large amount of researchers. The procedure will be to exploit the fact that coherent states are the most convenient (over)complete bases to construct a formal Feynman path integral (see e.g. and references therein). Certainly, a lot of work will be necessary to make that formal path integral rigorous but coherent states provide a definite starting point. The coherent state path integral should then be equivalent to a spin foam model which can be considered as a path integral using “momentum eigenstate bases”. These issues will be worked out in .
* Finally, the coherent states that we have constructed are coherent over a fixed graph only, they are pure states. However, our techniques readily combine with the random lattice approach developed in to produce mixed coherent states, that is, trace class operators on $``$ (in physical terms : density matrices). We can outline some of the ideas already here : Given a density parameter $`\lambda `$, an infinite volume cut-off parameter $`r`$ and a spatial metric $`q_{ab}`$, to be approximated by a mixed coherent state $`\widehat{\rho }_{A^0,E^0}^{s\lambda r}`$, we choose a number of $`1N^r<\mathrm{}`$ points at random in $`\mathrm{\Sigma }^r`$ where $`\mathrm{\Sigma }^r`$ is a compact subset of $`\mathrm{\Sigma }`$ which tends to $`\mathrm{\Sigma }`$ as $`r\mathrm{}`$. This is done in such a way that the density of points as measured by $`q_{ab}`$ is roughly constant and equal to $`\lambda `$. More precisely, a region $`R\mathrm{\Sigma }^r`$ is macroscopic if its volume satisfies $`V_R\left(q\right)=_Rd^3x\sqrt{det\left(q\left(x\right)\right)}\mathrm{}_p^3`$. Then we find approximately $`N_R\left(q\right)=\lambda V_R\left(q\right)`$ points inside this region (it will be convenient to choose $`\lambda =1/a^3`$ where $`a`$ is the length parameter of equation (3.14)). We will also set $`V^r:=V_{\mathrm{\Sigma }^r}\left(q\right)`$. It follows that
$$d\mu _q^r\left(x\right):=\frac{\sqrt{det\left(q\right)\left(x\right)}}{V^r}d^3x$$
(5.39)
is a probability measure on $`\mathrm{\Sigma }_r`$ (the necessity for the cut-off $`r`$ is evident). The probability to find the $`N^r=\lambda V^r`$ points $`p_k`$ in the infinitesimal volumes $`\sqrt{det\left(q\right)\left(p_k\right)}d^3p_k`$ is given by $`_{k=1}^{N^r}d\mu _q^r\left(p_k\right)`$. For each such random distribution of points we can construct a four-valent lattice $`\gamma _{p_1,..,p_{N^r}}^q`$ by the generalized Dirichlet-Voronoi construction which depends on $`q`$. Automatically, also a dual lattice is generated which we can use for the polyhedronal decomposition of $`\mathrm{\Sigma }^r`$ dual to $`\gamma _{p_1,..,p_{N^r}}^q`$ and which goes into the definition of the momenta $`p_j^e,eE\left(\gamma _{p_1,..,p_{N^r}}^q\right)`$. For this lattice, let $`\widehat{P}_{\gamma _{p_1,..,p_{N^r}}^q}^{s\lambda r}(A_0,E_0)`$ be the one-dimensional projector on the coherent $`C_0`$-vector $`\xi _{\gamma _{p_1,..,p_{N^r}}^q,(A_0,E_0)}^s`$ as in (5.34). Then, the task is to show that for the following operator, (which is trace class at $`r<\mathrm{}`$),
$$\widehat{\rho }_{A^0,E^0}^{s\lambda r}:=_{\left(\mathrm{\Sigma }^r\right)^{N^r}}\underset{k=1}{\overset{N^r}{}}d\mu _r^q\left(x_k\right)\widehat{P}_{\gamma _{p_1,..,p_{N^r}}^q}^{s\lambda r}(A_0,E_0)$$
(5.40)
the limit $`r\mathrm{}`$ exists as a trace class operator on $``$ which can presumably be proved by invoking inductive limit methods. That (5.40) is trace class at finite $`r`$ follows from
$$\text{tr}\left[\widehat{P}_{\gamma _{p_1,..,p_{N^r}}^q}^{s\lambda r}(A_0,E_0)\right]=\xi _{\gamma _{p_1,..,p_{N^r}}^q,(A_0,E_0)}^s^2=1$$
by construction so that actually $`\text{tr}\left(\widehat{\rho }_{A^0,E^0}^{s\lambda r}\right)=1`$ for all $`r`$ as it should be for a mixed state. This is a strong indication that the limit within the trace class ideal of the set of bounded operators exists. Practically, it might even be unnecessary to actually perform the limit as long as one measures only local operators : if the surfaces and paths with respect to which the operator is smeared lie within $`\mathrm{\Sigma }_{r_0}`$ then the measurement should be the same for all $`r>r_0`$.
This state is an average over a huge class of graphs and should have an improved semi-classical behaviour as compared to the pure ones. Notice that it is here that the possibility to compute inner products between $`C_0`$ vectors over different graphs becomes important. The details of this construction will appear in .
As the above list reveals there exists a plethora of fascinating and challenging open question and a huge programme is to be performed. In particular, the formalism is expectedly rather complicated as far as computations are concerned. The development of approximation techniques and error estimates as outlined in will become important. The coherent states together with the Infinite Tensor Product beautifully combine three main research streams in general relativity :
A) Quantum Gravity, since these are states of a quantum theory of general relativity,
B) Classical Mathematical General Relativity, since the states are labelled by classical solutions of Einstein’s equations and
C) Numerical Relativity, since the computations will need the help of supercomputers, the stage is prepared for numerical canonical quantum general relativity. In fact, since the graphs that we are using are not too different from the grids employed in numerical general relativity and lattice gauge theory, some codes in classical numerical relativity or lattice gauge theory might be easily adaptable to our purposes, although many new codes have to be written as well, for instance a fast diagonalization code for the volume operator.
Remark :
In the authors observe that the quantum fluctuations for the holonomy operator of a macroscopic loop, being the product of a large number of holonomies along “plaquettes” or elementary loops, are always large and it seems that there is no state that can approximate such holonomy operators. First of all, this “problem” is not tied to, say, lattice gauge theories but applies to any theory in which operators that are products of a large (or infinite) number of elementary operators play a role. Next, while the observation is certainly correct, given a large loop $`\alpha `$ on a graph $`\gamma `$ we can trade it for a single plaquette loop $`\beta `$ while keeping the number of holonomically independent loops constant. With this relabelling of our degrees of freedom over $`\gamma `$ the loop $`\alpha `$ is now elementary and we can write down a coherent state which approximates it arbitrarily well. From the point of view of the ITP, while the coherent states with either $`\alpha `$ or $`\beta `$ considered as elementary are defined over the same $`\gamma `$, they correspond to different regroupings of Hilbert spaces labelled by edges in the infinite tensor product. Thus, we see once more that the observation of is directly related to the fact that the associative law is generally wrong for the infinite tensor product of Hilbert spaces.
### 5.4 Dynamical Framework
So far our discussion has not touched the question whether it is possible to construct coherent states which are at the same time physical, that is, annihilated by the constraint operators in an appropriate (generalized) sense. At least with respect to the gauge – and diffeomorphism constraint one might think that the answer should be given by the group averaging proposal applied in to finite graphs. This section is intended to point out in which sense this can be carried over to infinite graphs. We first consider the averaging of general functions and after that averaging of coherent states.
#### 5.4.1 Gauge Group Averaging
The following trivial example demonstrates that the group averaging proposal requires due modification in the ITP context already at the level of the Gauss constraint :
Recall that the group $`𝒢`$ of local (generalized, i.e. without continuity requirements) gauge transformations $`g:\mathrm{\Sigma }G;xg\left(x\right)`$ is unitarily represented on $``$ by extending its action on $`C_0`$-vectors over $`\gamma `$
$$\widehat{U}\left(g\right)_f=_{eE\left(\gamma \right)}f_e\left(g\left(e\left(0\right)\right)h_eg\left(e\left(1\right)\right)^1\right)$$
(5.41)
to a dense subset of $``$ by linearity and to all of $``$ by continuity for any $`g`$.
Consider once more for $`\gamma `$ the $`x`$-axis in $`\text{ }\mathrm{R}^3`$ subdivided into unit intervals $`e_n=[n1,n],n\text{ }\mathrm{Z}`$. On this graph we can introduce the non-gauge invariant $`C_0`$-vector
$$\chi _\pi :=\underset{n}{}\chi _\pi \left(h_{e_n}\right)$$
(5.42)
where each edge carries the same irreducible representation $`\pi `$. Group averaging this vector with respect to the Gauss constraint means to compute the infinite dimensional integral
$$\eta _G\chi _j:=\underset{n}{}_G𝑑\mu _H\left(g\left(n\right)\right)\delta (g\left(\mathrm{}\right),1)\delta (g\left(\mathrm{}\right),1)\underset{m}{}\chi _\pi \left(g\left(m1\right)h_{e_m}g\left(m\right)^1\right)$$
(5.43)
where the $`\delta `$ distributions are due to the boundary condition that $`g\left(\pm \mathrm{}\right)=1`$. We consider (5.43) as the limit as $`N\mathrm{}`$ of
$$\eta _G^N\chi _j:=\underset{n=N}{\overset{N}{}}_G𝑑\mu _H\left(g\left(n\right)\right)\delta (g\left(N\right),1)\delta (g\left(N\right),1)\underset{m=N+1}{\overset{N}{}}\chi _\pi \left(g\left(m1\right)h_{e_m}g\left(m\right)^1\right)$$
(5.44)
which can be readily computed and gives
$$\eta _G^N\chi _j=\frac{\chi _j\left(h_{e_{N+1}..e_N}\right)}{d_\pi ^{2N1}}$$
(5.45)
Thus the norm of this vector is $`\eta _G^N\chi _j=1/d_\pi ^{2N1}`$ and so (5.43) vanishes unless $`d_\pi =1`$.
In order to cure this we must obviously factor out the power of $`d_\pi `$. This can be achieved by requiring that group averaging should produce a norm one vector from a norm one vector, that is, we propose (see )
$$\chi _Gf=\frac{_{vV\left(\gamma \right)}_G𝑑\mu _H\left(g_v\right)_{eE\left(\gamma \right)}f_e\left(g\left(e\left(0\right)\right)h_eg\left(e\left(1\right)\right)^1\right)}{_{vV\left(\gamma \right)}_G𝑑\mu _H\left(g_v\right)_{eE\left(\gamma \right)}f_e\left(g\left(e\left(0\right)\right)h_eg\left(e\left(1\right)\right)^1\right)}$$
(5.46)
where one makes sense of the formally zero numerator and denominator through a limiting procedure as outlined above. (One does not need to check that the result is independent of the way one performs the limit, if one gets different gauge invariant answers one simply gets different gauge invariant vectors which is all that we want from the group averaging machinery anyway for the case of the gauge group since, due to its finiteness, we can still use the extended Ashtekar-Lewandowski measure on group averaged states). This makes group averaging a non-(anti)linear procedure. It means, in particular, that we produce completely new Infinite Tensor Product Hilbert Spaces. Namely, in the case of the example the procedure (5.46) gives us the gauge invariant vector $`\mathrm{\Xi }_j=\chi _j\left(h_e\right)`$ where $`e`$ is the $`x`$-axis, the prototype of a tangle . So in this case the original graph $`\gamma `$ with its countable number of edges has collapsed to a graph with a single edge, a finite tensor product Hilbet space. Following definition 5.2 to compute ITP inner products for $`C_0`$ vectors over different graphs we see that the scalar product between $`\chi _j`$ and $`\mathrm{\Xi }_j`$ vanishes, the averaged and unaveraged vectors are orthogonal to each other.
That this happens is not an accident but generic : Consider a graph $`\gamma `$ which is the union of an infinite number of mutually disjoint, finite graphs $`\gamma _n,n=1,2,..`$. Then a $`C_0`$ vector over $`\gamma `$ is of the form
$$f=_n\left[_{eE\left(\gamma _n\right)}f_e\right],\left|\right|f_e\left|\right|=1$$
and defines an element of $`_\gamma =_{eE\left(\gamma \right)}_e`$. Group averaging evidently turns this $`C_0`$-vector into a new $`C_0`$-vector of the form
$$\eta _Gf=_nf_{\gamma _n},\left|\right|f_{\gamma _n}\left|\right|=1$$
which now is an element of $`_\gamma ^{}=_n_{\gamma _n}`$. This once more demonstrates the source of the trouble : the associative law does not hold on the ITP and the latter vector simply cannot be written, in general, as a linear combinations of vectors of the former Hilbert space.
#### 5.4.2 Diffeomorphism Group Averaging
Next we turn to group averaging with respect to the diffeomorphism constraint. Recall that this is done by relying explicitly on the spin-network basis. This is necessary because only if a function cylindrical over a graph $`\gamma `$ depends on each of its edges through non-trivial irreducible representations does group averaging over the diffeomorphism group produce a well-defined distribution, the complication being due to the infinite volume of the diffeomorphism group with respect to the “counting measure” which produces a singularity each time we sum over diffeomorphisms which do not modify the graph on which a function depends. Another complication associated with so-called graph symmetries can be satisfactorily dealt with, see for details.
However, as we have seen in section 5.1.3 spin-network functions do not provide a basis in the ITP context. It follows that not all functions of the ITP Hilbert space can be group averaged with respect to the diffeomorphism group.
More precisely, recall that the group Diff$`\left(\mathrm{\Sigma }\right)`$ of analyticity preserving diffeomorphisms of $`\mathrm{\Sigma }`$ is unitarily represented on $``$ by extending its action on $`C_0`$-vectors over $`\gamma `$
$$\widehat{U}\left(\phi \right)_f=_{eE\left(\gamma \right)}f_e\left(h_{\phi \left(e\right)}\right)$$
(5.47)
to a dense subset of $``$ by linearity and to all of $``$ by continuity for any $`\phi \text{Diff}\left(\mathrm{\Sigma }\right)`$. Given a $`C_0`$ sequence $`f`$ we define its orbit $`\left\{f\right\}`$ to be the set of $`C_0`$ sequences given by
$$\left\{f\right\}=\{f^{};\phi \text{Diff}\left(\mathrm{\Sigma }\right)_f^{}=\widehat{U}\left(\phi \right)_f\}$$
(5.48)
The following $`C_0`$ sequences lie in the range of the group average map.
###### Definition 5.7
A $`C_0`$ sequence $`f`$ is called a spin-network $`C_0`$ sequence over $`\gamma `$ if and only if $`<1,f_e>=0`$ for all $`eE(\gamma )`$. A spin-network $`C_0`$ sequence is called finite if its graph symmetry group is finite. For finite spin-network $`C_0`$ sequences we define the group average of its associated $`C_0`$-vector with respect to the diffeomorphism group by
$$\eta _{Diff}_f:=\left[_f\right]:=\underset{f^{}\left\{f\right\}}{}_f^{}$$
(5.49)
where we have assumed that the graph symmetry group of $`\gamma `$ is trivial, otherwise we modify the procedure as in or . The object (5.49) lies in $`\mathrm{\Phi }^{}`$, the topological dual of $`\mathrm{\Phi }`$ as follows from results of .
Graphs with infinite graph symmetry group are excluded from the domain of the average map, similar as in . Notice that for the typical graphs that we have in mind (e.g. cubic lattices) the graph symmetry group is in fact infinite due to the infinite number of translations which leave the graph invariant, but in order to cure this it is enough to replace a single edge by a kink. With this problem out of the way, this defines $`\eta _{Diff}`$ on finite spin-network $`C_0`$ vectors over typical lattices and can be extended by linearity to finite linear combinations of those. That this indeed defines a linear operation is granted due to our treatment of graph symmetries.
#### 5.4.3 Averaging of Coherent States
The interesting question is, of course, whether the coherent states that we defined are in the domain of the average map.
A) Gauge group averaging.
Returning to the example graph already discussed in equation (5.42) above, consider the (non-normalized) coherent state over the graph with $`2N`$ adjacent unit intervals as edges symmetrically around the origin along the $`x`$-axis, that is,
$$\psi _{g_N}^s\left(A\right):=_{n=N+1}^N\psi _{g_{e_n}}^s\left(h_{e_n}\left(A\right)\right)$$
(5.50)
where $`\psi _g^s`$ was defined in (3.1). Under a gauge transformation, represented by the unitary operator $`\widehat{U}\left(g\right)`$, the tensor product factor with label $`n`$ is transformed into
$$\psi _{g_{e_n}}^s\left(g\left(e_n\left(0\right)\right)h_{e_n}g\left(e_n\left(1\right)\right)^1\right)=\psi _{g\left(e\left(0\right)\right)^1g_{e_n}g\left(e_n\left(1\right)\right)}^s\left(h_{e_n}\right)$$
and integrating over $`g\left(e_n\left(1\right)\right),n=N+1,..,N1`$ with the Haar measure produces the state
$$\mathrm{\Psi }_{g_N}^s\left(A\right):=\psi _{g_{e_N}}^{2Ns}\left(h_{e_N}\left(A\right)\right)$$
(5.51)
where $`e_N=e_{N+1}..e_N`$ and $`g_{e_N}=g_{e_{N+1}}..g_{e_N},h_{e_N}=h_{e_{N+1}}..h_{e_N}`$. In other words, the finite number of integrations produce a coherent state with the correct dependence on $`h_e,g_e`$, however, the classicality parameter gets augmented from $`s`$ to $`2Ns`$ which in the limit $`N\mathrm{}`$, of course, does not show any classical behaviour any longer. Thus, in order to produce a gauge invariant coherent state form a non-gauge invariant one on the ITP by group averaging not only do we have to go through a limiting procedure as $`N\mathrm{}`$ as already discussed above with an associated “renormalization” of the norms of the vectors before and after averaging, but also one has to rescale the classicality parameter $`s`$ appropriately.
Thus, gauge group averaging becomes very difficult to perform if the graph $`\gamma `$ has at least one infinite connected component. At the opposite extreme are the infinite cluster graphs which are infinite graphs obtained by the infinite disjoint union of finite graphs, called clusters. Obviously, each of these finite graphs can be gauge group averaged (and renormalized) individually. In particular, if all clusters are diffeomorphic then group averaging reduces to the infinite repetition of one averaging for functions over one finite graph. An example to keep in mind is a cubic lattice in which we remove some edges to obtain disjoint finite cubic sub-lattices.
B) Diffeomorphism group averaging.
Recall (see the first reference in )
$$\sqrt{\frac{\mathrm{sinh}\left(p^e\right)}{2\sqrt{\pi }p^e}s^{3/2}\left(1K_s\right)}<1,\xi _{g_e}^s>=1/\left|\right|\psi _{g_e}^s\left|\right|\sqrt{\frac{\mathrm{sinh}\left(p^e\right)}{2\sqrt{\pi }p^e}s^{3/2}\left(1+K_s^{}\right)}$$
where the constants $`K_s,K_s^{}`$ tend to zero exponentially fast with $`s0`$. Since $`p^e`$ is bounded, tending to zero for ever and ever finer lattice at least for classical configurations we see that for sufficiently fine lattices at given (small) $`s`$ we have not only
$`|<1,\xi _{g_e}^s>|<1`$ as granted by the Schwarz inequality but moreover that there exist numbers $`0<q,q^{}<1`$ with $`q<c_e:=|<1,\xi _{g_e}^s>|q^{}`$ for all $`e`$ for sufficiently fine lattices which is precisely the application that we are aiming at.
Splitting $`\xi _{g_e}^s=\delta \xi _{g_e}^s+c_e1`$ we may want to write for given $`\gamma \mathrm{\Gamma }_\sigma ^\omega `$ the state
$$\xi _{\gamma g_\gamma }^s:=_{eE\left(\gamma \right)}\xi _{g_e}^s$$
(5.52)
as
$$\xi _{\gamma g_\gamma }^s=\underset{N=0}{\overset{\mathrm{}}{}}\underset{\{e_1,..,e_N\}E\left(\gamma \right)}{}\left[\underset{k=1}{\overset{N}{}}c_{e_k}\right]\left[_{eE\left(\gamma \right)\{e_1,..,e_N\}}\delta \xi _{g_e}^s\right]$$
(5.53)
or as
$$\xi _{\gamma g_\gamma }^s=\underset{N=0}{\overset{\mathrm{}}{}}\underset{\{e_1,..,e_N\}E\left(\gamma \right)}{}\left[\underset{eE\left(\gamma \right)\{e_1,..,e_N\}}{}c_e\right]\left[_{k=1}^N\delta \xi _{g_{e_k}}^s\right]$$
(5.54)
However, both attempts fail since in (5.53) all appearing vectors have zero norm (in fact $`\delta \xi _{g_e}^s^2=1c_e^2<1q^2<1`$) and in (5.54) all coefficients vanish identically. Thus, the vector (5.52) does not lie in the domain of the average map.
A substitute for averaging and to deal with the diffeomorphism group is to work with representatives, i.e. from each diffeomorphism class $`\left\{f\right\}`$ we choose an element $`f_{\left\{f\right\}}^0`$. Thus $`f^0:\left\{f\right\}f_{\left\{f\right\}}^0`$ is a choice function, its existence being granted by the lemma of choice. We specify this choice function further by choosing from each graph diffeomorphism class $`\left\{\gamma \right\}`$ a representative $`\gamma _{\left\{\gamma \right\}}^0`$. Given a function $`f`$, let $`\gamma _f`$ be the minimal graph on which it depends non-trivially. Then $`f_{\left\{f\right\}}^0`$ can be chosen to depend on $`\gamma _{\left\{\gamma _f\right\}}^0`$. If $`\gamma _f`$ has graph symmetries then this prescription does not yet fix $`f_{\left\{f\right\}}^0`$ uniquely and we must further choose from one of the $`\widehat{U}\left(\phi _n\right)f_{\left\{f\right\}}^0`$ where $`\phi _n`$ is a symmetry of $`\gamma _{\left\{\gamma _f\right\}}^0`$. A kind of group averaging map is now defined by $`\eta _{Diff}f:=f_{\left\{f\right\}}^0`$ which obviously satisfies $`\eta _{Diff}\widehat{U}\left(\phi \right)=\eta _{Diff}`$ and the inner product on these “solutions to the diffeomorphism constraint” is just the usual inner product between representatives. This makes the whole proposal unfortunately very choice dependent and thus less attractive. Notice, however, that diffeomorphism invariant operators which are defined on the kinematical Hilbert space obviously keep their adjointness properties.
A different way to deal with diffeomorphism invariance is by gauge fixing (alternatively, one has to construct gauge and diffeomorphism invariant coherent states from scratch) :
Given a collection $`g_\gamma =\left\{g_e\right\}_{eE\left(\gamma \right)}`$, a local gauge transformation $`g𝒢`$ and a diffeomorphism $`\phi \text{Diff}\left(\mathrm{\Sigma }\right)`$ we define $`g_\gamma ^g:=\left\{g_e^g\right\}_{eE\left(\gamma \right)}`$ with $`g_e^g:=g\left(e\left(0\right)\right)^1g_eg\left(e\left(1\right)\right)`$ and $`g_\gamma ^\phi :=\left\{g_e^\phi \right\}_{eE\left(\gamma \right)}`$ with $`g_e^\phi :=g_{\phi ^1\left(e\right)}`$. It is then easy to see, using unitarity (invariance of norms) that
$$\widehat{U}\left(g\right)\psi _{\gamma g_\gamma }^s=\psi _{\gamma g_\gamma ^g}^s\text{ and }\widehat{U}\left(\phi \right)\psi _{\gamma g_\gamma }^s=\psi _{\phi \left(\gamma \right)g_{\phi \left(\gamma \right)}^\phi }^s$$
(5.55)
Given classical initial value data $`(A^0,E^0)`$ in a certain gauge the $`g_\gamma =g_\gamma \left((A^0,E^0)\right)`$ are fixed and we require that $`g_\gamma ^g=g_\gamma =g_\gamma ^\phi `$ for all $`\gamma `$ which (generically) trivializes the residual gauge freedom to $`g=1,\phi =`$id.
Thus, as far as the gauge and diffeomorphism constraints are concerned, we can fix a gauge to take care of gauge and diffeomorphism invariance. The issue lies much harder with respect to the Hamiltonian constraint because its action is much more complicated than the action of the kinematical constraints, and almost no Hamiltonian invariant observables are known with respect to which one would need to construct the invariant coherent states. Fortunately, there are certain “simple” solutions to the Hamiltonian constraint corresponding to states whose underlying graph is out of the range of graphs that the Hamiltonian constraint produces. If we build (non-distributional) coherent states on such graphs, then they lie in the kernel of the Hamiltonian constraint in the sense of generalized eigenvectors with eigenvalue zero. Thus, at least for these simple solutions, together with fixing of gauge and diffeomorphism freedom, we can incorporate the quantum dynamics of general relativity.
All these observations reveal that group averaging non-gauge and/or non-diffeomorphism invariant coherent states over the gauge or diffeomorphism group is a non-trivial task, at least not if $`\mathrm{\Sigma }`$ is non-compact and applied naively without some sort of renormalization leads to meaningless results. More work is needed in order to construct rigorous solutions to all constraints which at the same time behave semi-classically. However, at the moment we are not so much interested in obtaining semi-classical solutions to all constraints. Rather, besides the applications already mentioned in section 5.3, it is of paramount importance to test the consistency of a quantum representation of the classical constraint algebra and the verification of its correct classical limit . In order to do this one obviously must not have semi-classical states which solve the constraints.
Acknowledgements
O. W. thanks the “Studienstiftung des Deutschen Volkes” for financial support.
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# Astronomers and the Science Citation Index, 1981–1997
## 1 INTRODUCTION
For the past several decades, we in astronomy have relied on the Science Citation Index, as compiled and published by the Institute for Science Information (ISI), as our source for citation statistics of our papers. At the dawn of the 21st Century, we in astronomy are becoming more and more reliant on electronic databases for the papers we read, and for the citation statistics on those papers. The two main internet sites we access for listings of our published papers (as opposed to preprints) are the Astrophysics Data System (ADS) (adsabs.harvard.edu) and the Institute for Science Information’s (ISI) Science Citation Index (www.webofscience.com).
Indeed, the Science Library at our University no longer subscribes to the ISI’s printed Science Citation Index; rather there is now complete reliance on the ISI’s web-accessible Web of Science. Yet, as this author has discovered during the research conducted for this paper, many of us do not clearly understand the contraints and limitations of either the hard-copy Science Citation Index and the new Web of Science.
However, what prompted the research done in this paper was the a new kind of citation analysis produced by the ISI research group over the past few years. At least two new citation lists were generated by the ISI research group before this author contacted them: A list of the chemists cited 500 times or more and a similar list for physicists/astronomers. Each list differs from what we can access either from the Web of Science or from the hard-copy Science Citation Index, in that time intevals for both citing and cited papers are specified.
A French chemist, Dr. Armel Le Bail (Laboratoire des Flouresces, CNRS, ESA) purchased the ISI’s most-cited chemist list (for $1000), and posted it on the Web (pcb4122.univ-lemans.fr/cgi-bin/physiciens.pl). Along with the most-cited chemist list, the ISI sent Dr. Le Bail the first 1120 names of physicists/astronomers on that most-cited list, which Dr. Le Bail scanned into his computer and posted on the web. One of the reasons Dr. Le Bail posted these lists on the web was given by Dr. David Pendlebury of the ISI, who points out that, of the 50 most-cited chemists, 7 have been awarded the Nobel Prize (cf. Garfield & Welljam-Dorof 1992).
The interest of the present author in this subject was piqued when one of his colleagues pointed out that his name was in the most-cited physicists/astronomers list on the web page of Dr. Le Bail. Thus began the journey of this author down the rabbit hole of web-accessible and web-generated paper and citation information. The present paper, with its lessons learned and data gathered, is the net result of that journey.
Aside from the not-inconsiderable curiosity factor (e.g., most of us would like to know where we stand relative to others in the number of times our papers have been cited) and job-related factors (my own promotion to professor was aided by such a list compiled by a colleague for one year of citations for our department), why would one do such a study? This author can think of several questions one would like to answer. First, at the very least, if citation information is to be used in connection with job-related decisions, should not the available data be of highest possible quality? Such data should be treated like any other data, and investigated as to random and systematic errors. Should it not also be clear what assumptions go into the data being used? Given that a relationship exists between being most-cited and winning a Nobel prize among chemists, does the same relationship exist for honors received by astronomers? As a guide to those scientists just entering our field of study, what do these data tell us about how the way we put our names on our papers, and where we publish our papers, influence how we are honored by our peers?
Previous papers which tried to assess citation information for astronomers (Abt 1981a,b, 1982, 1983, 1984a,b, 1985, 1987a,b, 1988a,b 1989, 1990a,b, 1992a,b, 1996, 1998a,b; Abt & Zhou 1996; Trimble 1985, 1986a,b, 1988, 1991, 1993a,b, 1996; White 1992; Girard & Davoust 1997; Davoust & Schmadel 1987, 1992) were limited by time and data access to asking statistical questions that are more restricted than those that now can be addressed electronically. The methodology employed by this paper are detailed in Section 2, where we address, in detail, what one can do, and what cannot do, with the present databases made available to this author by the ISI as well as those databases generally available on the web. The data we have generated from this study are discussed in Section 3. The statistical studies of citations for astronomers are discussed in Section 4, both among themselves and in comparison to the life-time honors bestowed to individuals. The main results of this paper are summarized in Section 5, where a “modest proposal” is made towards solving the ever-pervasive name confusion problem.
## 2 METHODOLOGY
### 2.1 The ISI Databases: Definition and Restrictions
The first list of most-cited physicists/astronomers generated by Dr. David Pendlebury and his collaborators and given to Dr. Le Bail, has a cutoff at 500 citations/name. When the “unique-two” issue was discovered (cf. Secs. 2.3, 3.1), this writer requested a new list be generated by the ISI that placed its cutoff at 100 citations/name. This new list, provided by Dr. Pendlebury to this writer, gives 62,813 astronomer and physicist “names” cited in the “usual” refereed journals at least 100 times or more during the time period January, 1981 (1981.0) through June, 1997 (1997.5), for papers published in the same journals during the same time period (hereafter referred to as the “P&A-100” list). Note that what the ISI provides in its databases is literally a last name and first initial(s), without uniquely identifying that name, per se, with an actual individual.
Separately, Dr. Pendelbury’s group has generated a list of the top 200 papers cited in astronomy each year from 1981-1996 (“High Impact Papers in Astronomy, 1981–1996”; hereafter referred to as the “AST-top-papers” list) from a set of astronomer-used journals (see Table 1), for citations made up to 1998.0. Among the 3,200 papers in this database are 5,035 astronomer ISI “names.” For the purposes of this paper alone, the ISI has given this author access to these two databases, and it is from these datasets that this author has compiled the citation data in this paper for astronomers. In return, this author has communicated to the ISI the various issues one finds when one attempts to do a complete survey of citations for scientists in a relatively well-defined field of study, such as our own.
Abbreviations for the journals used for the ISI P&A-100 list are given in Table 1, divided into four classes (ISI code in parentheses): space science (SP), condensed matter and applied physics (APP), optical and acoustic (O/A) and general physics (PHS). Essentially all of the standard journals in which astronomers publish their papers are in the SP category. The journal abbreviations come from both the Astrophysical Data System (ADS) (generally mixed upper and lower case letters) and from the ISI (an 11-character code, all capital letters). Those SP journals whose abbrevations are given in bold letters in Table 1 are used for the AST-top-papers database. As is evident, the AST-top-papers list employs the “usual” astronomy journals, so the “names” in that database can be all considered as those of astronomers for this analysis. In contrast, the P&A-100 list includes all of the journals in Table 1, and thus includes far more physicists than astronomers (as is evident from the marked contrast in numbers of names in the two lists).
Both the ISI databases have limitations placed upon them in terms of how the ISI does its business in handling certain well-known “sticky” citation analysis issues:
First, no meeting papers or books are used for these compilations, either for the cited papers or the citing papers. Indeed, if you currently use the Web of Science’s “Cited Search” option, you have to specify at least the first author of any article for that search to bear fruit (more discussion on the issues of using this webpage in Secs. 2.4-2.7). The ISI is not alone in making apparently draconian decisions about meeting papers. The ADS also tends to list only a few authors on multi-authored meeting papers, as well as for many journal papers published pre-electronic submission (see Secs. 2.4-2.7). It is apparently difficult for electronic databases to account for authorship of meeting papers/catalogs/books or multi-authored papers, in their citation statistics for scientists.
Second, while the authorship of a paper that is cited comes from the cited paper itself (minimizing spelling errors of authors’ names), the citation for that paper comes from the citing paper. Human error being as it is, citations in the citing papers can sometimes be in error. This leads to what is, in effect, a random error for the citations that is proportional to the number of citations for a given paper (i.e., more hands in the pot, more likely an error). An estimate of this random error will be made in Sec. 2.6.
Third, only the first 16 authors of a given journal paper are used in the two ISI databases used here. This decision was made by the ISI in the interests of economy, and from the fact that relatively few physics/astronomy papers have many authors. Fortunately, one can use the AST-top-papers list to make corrections to the P&A-100 list to account for this exclusion (Sec. 2.2) so that the resulting error in citation statistics is neglgible.
Fourth, the long-standing decision of the ISI is that every author of a paper be given full credit for that paper. Whether or not one agrees with this decision, this is what you get when you access the Web of Science or the ADS for citation purposes. Whether or not this author (or a reader of this paper) agrees with this decision by ISI or not, this is what the ISI does to generate the Web of Science and the data lists used in the present analysis. It is also in keeping with the philosophy they used for the hard-copy Science Citation Index, but could not completely employ in that presentation mode. This important issue is discussed in detail in Sec. 2.4.
Fifth, the lists are not lists of individual astronomers and physicists, but, as stated above, the “names” as given on their papers. These “names” uniformly have a last name and first initial(s) only. In cross-correlating these two datasets with lists of actual astronomers, name confusion is inevitable if all we have for identification are the ISI names. Many people share the same last names and first initial(s).
### 2.2 A Step-by-Step Name Search Process
The two ISI databases provide citation data for astronomers in a complementary way. Comparison of the names of astronomers in the AST-top-papers list with those in the P&A-100 list shows that not all astronomers who have 100 or more citations are included in the P&A-100 list, owing to the multi-author problem (“sticky” issue 3, above). Comparison of the names of astronomers in the AST-top-papers list to those in the AAS 1998 Membership Directory finds many astronomers not in the AAS Directory. Hence, by using the P&A-100 list as the primary data source, and the AST-top-papers list as the secondary data source, we can assemble a more complete citation database for astronomers than by using either ISI database alone. The assembly of the citation database for astronomers was done through five time-consuming steps, using all of the information, both written and electronic, available to this author.
Step 1. The names of astronomers in the 1998 AAS Membership Directory were compared with the last names, first initial(s) in the P&A-100 list. Those in common were noted. The AAS Directory was used as this is the only membership directory for astronomers generally available. Subsequent to placing a copy of this paper on the Web, this author did the same comparison for those astronomers in the 2000 directory of the Astronomical Society of India.
Step 2. For each of the found “names,” a web search was made using the ADS to see whether more than one person, physicist or astronomer, could have that “name.” The ADS was used for this search, as it gives, in many cases, first names for individuals. Depending on what was found on the ADS, the listed “name” was either uniquely assigned to an astronomer, or that “name” was marked as confused with other astronomers/physicists. At this point in the process, a preliminary Astonomy Citation Database (ACD) was assembled ($`3700`$ names).
Step 3. To further help to overcome the obvious AAS-driven, North American bias in the preliminary ACD, the “names” of astronomers in the AST-top-paper list were compared to the names in both the P&A-100 list and in the preliminary ACD. Those AST-top-papers “names” found in the P&A-100 list but not yet in the ACD were added to the ACD via the same Web search process as in Step 2 above ($`2500`$ names).
Step 4. To overcome the 16 author/paper limit of the ISI databases, a search was made of all papers (105) in the AST-top-papers list that had more than 16 authors. These included papers with as few as 15 citations and as many as 908 citations. The remaining authors on these papers were then found either by looking up the paper on the Web of Science or from the actual hardcopy of the journal. As before, the additional data for astronomers was added into the ACD via the process outlined in Step 2. The citation data for new astronomers ($`170`$) found by this process, and the citation data that was added to astronomers already on the ACD, are noted as “A16” in the ACD.
Step 5. We generated a list of astronomers (278) honored by the AAS since 1949, honored by the Nobel Committee and/or who are current members of the U.S. National Academy of Sciences. Any astronomer in this list not already in the existing ACD ($`50`$) was added to that list if the astronomer’s name was also in the P&A-100 list.
The citation data from the two ISI databases were modified in just one way to handle the special cases of three well-used data catalogs that were published during the cited time period (1981.0 to 1997.5): the Third Reference Catalog of Bright Galaxies (de Vaucouleurs et al. 1991), the Revised Shapley Ames Catalog (Sandage & Tammann 1981, 1987) and the Bright Star Catalog (Hoffleit 1982). The printed Science Citation Index and the Web of Science were both used to estimate ($`\pm 100`$) the number of citations to add for these catalogs for the individuals involved: 1000 citations for the Third Reference Catalog, 1250 citations for the Revised Shapley–Ames Catalogs, and 1395 citations for the Bright Star Catalog. Citation data for these catalogs follow the same rules as for the citation data in the ISI lists — only the those citations as given in papers in refereed journals during the stated citing time period.
### 2.3 Name Search Leads to Name Confusion
The ISI procedure of taking the names of the authors from the papers cited, rather than from the citing papers eliminates most, but not quite all, spelling errors. Of the few specific instances of misspelling caught by this author in the two ISI databases, most have been fixed by the ISI. There still remain three individuals whose names that are likely mispelled in the journal proceedings (or mis-scanned by the ISI): RP Kirshner/Kirschner; A Dyachkov/Dyatchkov; H Luhr/Luehr. The case of RP Kirshner is curious, as in 1998, the ADS recognized that both spellings are associated with the same person. As such, while citations are added together for Kirshner/Kirschner, the numbers of papers cited is kept at the value listed for Kirshner. In the cases of A Dyachkov and H Luhr, the mispellings are quite evident, so the papers and citations for both spellings are added together.
The same person can also have different last names, owing to marraige. Four persons having two names in the databases were so identified (with kind aid from V. Trimble): J Bland–Hawthorn/Bland, MJ Rieke/Lebofsky, S Viegas-Aldrovandi/Aldrovandi and S Collin-Souffrin/Collin. The citation data for the first three of these astronomers were added together, as the different names are on different papers. The citation data for S Collin-Souffrin cannot be added together with those for S Collin, as the S Collin name is confused with that of another astronomer as well as with names of physicists.
In an attempt to disentangle as much of the name confusion as possible, this writer took advantage of an aspect of the ADS website that is apparently not yet available on the Web of Science website: The ADS website will search all names using the string of letters it is given, and will do so for an exact name match even if you do not give the whole name. Moreover, the ADS will often give the first names of the individuals found. As such, it was possible to search both the astronomy and the physics ADS websites for each ISI “name” suspected to be name-confused. This search was eventually done for almost all “names” thought to be those of an astronomer. When an unambiguous first name is found for a last name having only one first initial, that first name is noted in the ACD. In addition, all last names that are in fairly common usage were searched for name confusion for all sets of first initials. Some name confusion was solved by comparing citation data for astronomers as obtained from the P&A-100 list with those obtained for the same astronomers in the AST-top-papers list. In the former list, astronmers and physicists can have their names confused; such is not the case in the latter list.
It was quickly found in this analysis that some astronomers are listed under two or more sets of initials in the ISI databases (e.g. JP Huchra also commonly has J Huchra on his papers). In the 357 cases currently in the ACD for which this writer could identify the author’s last name with two or more sets of first initials, and that set of initials/last name are all uniquely identified with an individual, the citations are summed for all sets of initials with the same last name. The number of papers is summed as well, as the ISI would find these as separate papers from those with another initial for the given last name.
Name confusion can occur for two reasons. First, many last names are common in each culture (e.g., Smith, Jones, Wang, Suzuki, Singh), so that those individuals with only one, or even two, first initials can have their names confused with those of other astronomers and, especially physicists. In certain cases, astronomers related to each other may have the same last name and first initials. Second, given that we find many astronomers using two sets of initials for their papers, we have to allow that other astronomers may do this, but that one or both sets of names they use can be confused with those of others (who also may or may not use two sets of initials). This then leads to all possible combinations of confused/unconfused names: one set of initials for a given last name uniquely identified with an individual, the other set confused; both sets of initials confused; a last name with one first initial having possible confusion with 2 or more (up to 6) names with that first initial plus different middle initials. In other words, far more than 357 astronomers publish using two or more sets of initials for their papers.
If this all sounds confusing, well, it is!
### 2.4 Methodology Summary
Investigating the citation data for astronomers is, indeed, a trip down the rabbit hole. To help make sense of this trip, here we answer seven questions one can ask about this kind of citation survey:
1. Which journals/meeting–proceedings/books does one choose from which to gather citations? The ISI has defined the answer. Only refereed journals as defined by the ISI (Table 1) are used here, both for papers cited and for papers doing the citations. Thus, papers published in meeting proceedings, and authorship of books and catalogs, are not used for citation statistics (Sec. 2.2).
2. The papers that do the citations are published during what time periods? The ISI defines these papers to be published between 1981.0 and 1997.5 in all of the journals in Table 1, including both physics and astronomy journals. The AST-top-papers list covers the citing papers during the period 1981.0 through 1998.0 (i.e., 1/2 year longer than the other list), using the astronomy journals whose abbreviations are in bold face in Table 1.
3. Papers during what time period are cited? The ISI defines this time period for each list slightly differently. For the P&A-100 list, it is from 1981.0 to 1997.5, same as the paper-citing time interval. For the AST-top-papers list, it is for papers published from 1981.0 to 1997.0, ending one year before the paper-citing period. In both cases, the cited papers are from the same journals as those that do the citing for each database. Again, it is stressed that meeting proceeding papers are neither cited nor are citations taken from them by the ISI. Hence, no meeting proceeding data are included in ACD.
4. Credit only the first author, or credit all authors on a given paper? The ISI standard practice for its Web of Science, and for its generation of citation lists is that all authors on a journal paper are credited for each paper. As stated earlier, this is not a change in ISI policy, but rather full expression the ISI policies given the freedom of the web.
5. Fractional, or unitary credit of citations for authors on multi-authored papers? The ISI gives each author of a journal paper full credit for that paper (i.e., unitary credit), up to 16 authors per paper.
6. How does one handle name confusion when only first initials and last names are available? This is handled in a complicated manner, using all available means available. The sobering fact is that even having full first names available does not completely remove the confusion issue, as many individuals also have the same first and last names.
7. What is the accuracy of the citation estimates? As dicussed below (Sec. 2.6), the random error in the ACD citations is proportional to the number of citations received, at about the 4% level. Other small, systematic effects exist as well, stemming from several different issues.
Of all of the choices that the ISI makes for its databases, the assignment of one paper credit to each author of a paper is the most-controversial. While this writer agrees with this choice, others with whom the author has discussed this paper do not. The plain fact is that any present or future study that uses the ISI citation database is using citation data which assigns each author of a paper all of the citations for that paper.
Given the lively controversy this decision by ISI has engendered among this authors colleagues, this writer feels it necessary to state why he agrees with ISI’s choice in this matter. First, and foremost, this writer does not know of a universally-accepted fair way to give fractional credit. On the small sampling of astronomer colleagues, whether or not one votes for fractional credit seems to depend on whether or not one has been involved in a multi-authored paper. Those of us who have been involved in multi-astronomer projects tend to vote for giving full credit to each author of a paper. Those who have not been involved much with such projects tend to vote to give fractional credit.
Yet, if ISI, or any reader of this paper, decided to give fractional credit for the citations of a paper, how would such fractional credit be calculated? Would it be fairer to divide the number of citations strictly by numbers of authors, and give each author fractional credit? Or do we try to credit some authors (say, first author) more than others? As shown in Sec. 2.6, if we take the most straightforward way to calculate citations for individuals in a fractional sense — even division by numbers of authors — there would be substantial revision of names among the top-cited astronomer list.
### 2.5 Two ways we lessen the impact of our papers
In the compilation of the astronomer citation list, this author has found that getting credit for the papers you publish can be lessened in two ways.
Both the ISI databases and the ADS (Sec. 2.1) have problems in registering citations for papers published in meeting proceedings. While books and catalogs can be found in the Web of Science via a “cited reference” search of the correct first author name and year of publication, these citations come only from those made in journal papers. So, while meeting papers are cited there, these data are not incorporated into the main ISI databases. Meeting papers do not figure at all into the citations statistics for the Web of Science, nor for the two ISI databases given this author to develop the ACD.
The other way credit for citations of your paper can be harmed is if your name is either confused, or if you use more than one first initial on your papers. For those of us with common last names, especially our Asian colleagues, solving name confusion will not easy (cf. Sec. 5). If you permit your papers to be published with two or more sets of first initials, the Web of Science will put you in two or more different places, fuzzing the credit you will get for your citations. This issue was one thing when we used the Science Citation Index hardcopy, as we could easily see the two different entries. Such is not the case when you access the same data electronically. Indeed, if using electronic databases, finding citations data for an individual who uses different first initials requires knowledge aforethought that this problem exists. While someone working in the field of study can sort this problem out with a lot of effort, it is impossible for people not working in that field of study (such as the ISI personnel) to do the correct sorting.
### 2.6 Citation Errors: “Random” and Systematic
A “random error” infects the ISI databases owing to the manner in which the ISI gets its citations per paper: the journal volume, page number of the cited papers are taken from the citing papers. Human error being what it is, a certain percentage of those cites give the wrong journal number, wrong page number (sometime switching one for the other), confuse ApJ with ApJL or ApJS, etc. One does not discover that this can lead to errors in the ISI databases until one accesses the Web of Science, pushes on the general search button, then opt for a “cited ref search.”
As opposed to the “general search” mode, “cited reference search” is best used by entering both the first author and the correct year of publication. What one then finds is a list of the citations for the papers of that author. Those that are underlined in blue you will also find in the general search. Those not highlighted either are misentered in a citing journal paper, or are referenced in citing journal papers to non-journal papers, private communications, books, catalogs, etc. In cited reference search mode you will see all of the misentered entries for a given cited paper.
This writer has accessed the cited reference search for 36 papers in Table 3 to estimate the misentry incidence as a function of journal. Note that statistics cannot be done for the citation data in the ACD per se, but rather for the current (2000.0) citation data for the papers involved. Citation data for individual papers is affected as much as 126 misentries for 989 1999.9 citations (for 1992 ApJ Letter paper of Smoot et al.) to a low of 2 misentries for 532 1999.9 citations (the A&A paper of Renzini & Voli 1981). If we express these errors in terms of percentages, they range from a high of 19% to a low of 0.004%, with a mean of 4%. This error is strongly journal–dependent, being the most for the ApJL (with a relatively high percentage of papers not citing the journal as ApJL), and lowest for Nature and PASP.
Given that none of us publish our papers solely in one journal, a reasonable estimate for a one-sigma random error of citations in the ACD is 4%. While this random error can be corrected in principle, in practice it would be highly labor-intensive, requiring a scientist from our field to work directly with the ISI to make the corrections, paper-by-paper. As such, a one-sigma random error of 4% in the number of our citations can be viewed as the dues we pay for using the digital computers to do the citation calculations.
Several systematic errors also affect the citations in the ACD:
Omission of astronomers: Non-inclusion of astronomers at the low end of the ACD comes about in a somewhat convoluted manner. 173 astronomer names from the “A16” analysis were excluded from the ACD because the ISI names of those astronomers had fewer than 100 citations accredited to them in the A16-corrected AST-top-papers list. The AST-top-papers list references only a subset of the journals used for the P&A-100 list, of which only a subset of those papers are used. Hence, it is likely that the names of at least some of these astronomers would be in the ACD if all papers on which they are authors had been included in the P&A-100 accounting (i.e., their other papers could total up to 99 citations). As such, these 173 astronomer names are given in a separate “honorable mention” list that will be supplied electronically.
Separately, inclusion of people who became deceased during the sampled period is handled both from the memory of this writer, and through the comparison of the astronomer names in the AST-top-papers list to those in the P&A-100 list. It is not expected that many such astronomer have been omitted from the ACD. As an additional check, a scan of the full P&A-100 list by this writer of those ISI names with 3000 or more citations uncovered no other astronomer names than those given in the ACD.
Wrong papers: One can compare the authors on the papers in Table 3 to the astronomer names in Table 2, to see that papers that publish wrong results do not substantially contribute to the citations for these astronomers. Moreover, one can verify that very few of the 3000+ astronomers in the top half of the ACD have reputations built on the publication of wrong papers. The conclusion of this paper is that citation of wrong papers negligibly influences the citation statistics.
Self-citation: The present available websites make it difficult to quantiatively assess the effect of self-citations on the citations for a given astronomer. What is important here is the variance around the percentage of our citations that are our own papers, as this is something we all do. On can most directly assess this for individual papers, which this writer has done for about a dozen, many-cited papers. The result of this non-statistical sampling is these papers self-cite in the range 3–15%, with a mean about 8%. Since it is variance about this mean that affects the citations of one astronomer vs. another, we can expect a variance of 5–7% in citations that can be ascribed to self-citation. Whether or not to assign a “self-citation” variance of $`6\%`$ to the ACD data (to then be added in quadrature to the random error estimate of 4%) is a matter of choice. As with giving/not giving full citation credit for each author of a paper, honest people can differ on whether or not to account for self-citation variance in the ACD.
Fractional or Full Credit: While a scientifically correct test of this issue cannot be made from either the Web of Science (no restriction on citing years) or from the two ISI databases. As such, compiling a list that is complete as the ACD that gives only fractional author credit is impossible without the full cooperation of the ISI. Nonetheless, an illustrative differential test can be made using the data in the AST-top-papers list. This writer took the top 13 individuals cited in the AST-top-papers list (which are not necessarily the same as the top-cited authors in the ACD), and counted the number of authors of each paper for each individual, as well as counted the number of papers that individual was first author, and the number of papers the name of that individual was placed in alphabetical order (for papers with 3 or more authors). To this list this writer added the same data for three other astronomers known to this writer to mainly publish significant, single-authored papers.
For a first test, fractional credit was given by dividing the number of citations for each paper by the number of authors on that paper, and then summing the fractional credit for each individual. The ratio of (fractional citations/full citations) so obtained ranged from 0.05 (for an individual publishing mainly with large groups) to 1.00 (for an individual with 2 sole-authored papers), with a median of 0.30. If these 15 individuals were the only ones in the database, a substantial reordering would be done from the full citation credit list. The net result would be to keep 7 of the 13 top-cited individuals in the top positions, but move 6 individuals from lower on the full credit list to near the top of the fractional credit list.
For a second test, we find the ratio of papers for which the author is first or sole, to the number of individual’s papers included in the AST-top-papers list. This ratio varies from a high of 1.00 to a low of 0.067, with a median of 0.27. If we then would calculate fractional credit by giving more credit to the first author of a paper, the order of listing of individuals would be different from that of either full credit or strict fractional credit lists.
The bottom line here is that there is simply no way one can calculate the number of citations for individuals that everyone can agree with. Depending on which way one chooses, the result will be different. Rather, given the now wide-spread use of the Web of Science, it seems most logical to accept the practices of the ISI in this regard in giving full citation credit to each author on a given paper.
Interestingly, how the ISI has presented its database has been a product of techonology. From the start, for the the printed Science Citation Index, ISI obviously had to make hard choices of how to present data, both in number of letters for names and number of authors per paper used. One choice that was made for the printed version was to list papers under only first authors in the Author part of the Index. All authors of a paper (up to an author limit given in the Index explanations of that year/summary list) would then be referenced to the each paper in the Source part of the Index. As such, previous investigations of citation-related issues would, indeed, have just found first-authored papers if only the Author part of the Index was used.
In contrast, the Web of Science now permits the ISI to show all of its data for each author. This means that when one accesses a given author name, one gets the full citations for each paper that author is on, not just the ones for which that author is first. The ISI notes such other papers by preferencing the Web-cited reference with several dots.
In summary, errors in data lists are introduced due to human error or human choice, whether using a computer or not. When data lists are of numbers, errors or choices are one thing; when data lists are of the accomplishments of real people, errors or choices are quite another thing. It is likely that human errors of omission and comission in ACD exist. Unforunately, there is no easy way for this writer to find all of them. The most serious of the likely systematic errors in the ACD concern exclusion of astronomers from the list, while the random errors of citations in the list are at about the 4% level (modulo accounting/non-accounting for 6% self-citation variance). The most intractable issue is to whether or not give full citation credit to each author. Whether or not one views this as an error in the data base depends on ones opinion. Interestingly, if one had only used the Author section of the hard-copy Science Citation Index in past years, one would have only found those papers for which the author was first.
While the current version of the ACD has been constructed to be as complete as possible, it is the ongoing aim of this writer to make this list as complete and accurate as is reasonably possible. Given that the time periods of citing and cited papers are fixed, this is doable. Towards this end, the reader is encouraged to contact this writer if the reader feels a name should be in the ACD that is not, of if one feels the data entered for a person is badly in error. All such cases will be addressed individually by this author and corrections to the ACD made, if warranted.
### 2.7 Are the ISI Lists the Best We Can Do?
The ISI lists do something important that this writer has not yet found among the available websites or in written form: Specify both a time interval for when the papers are published and a time interval for when the papers citing those papers are published, and sum all of the found citations per author name. This permits “snapshots” in time of citations to be assembled, of which the current ACD is just the first.
The Web of Science gives the number of citations per paper, but not per author for that year. This latter number is something one would have to manually extract from the information given. The Web of Science also does not currently permit the citing years to be restricted. As such, when checking for the random error problem (Sec. 2.4), one could only assess the errant citations against ISI-registered citations at the current time. However, a specific interval for the years the papers of an astronmer were published can be specified. Hence, what one gets is a running number of all citations of a given paper up until the date you access the citation information. In absence of a world-wide directory for astronomers, It is essentially impossible to predefine all astronomer names to do the kind of analysis with the Web of Science that one can do with the P&A-100 and AST-top-papers databases.
The ADS website has a button one can push which will give the total citations per astronomer from their database (such a button is not available for the physicist webpage, though). By pushing this button you get a summary total of citations for that astronomer, and a list of the citing papers, but not the number of citations per paper, nor the number of cited papers, nor a list of cited journals/meetings/etc. A perusal of citing papers for several individuals indicates that papers in meeting proceedings are only included in the ADS list if those proceedings could be accessed electronically. Combine this with the problems the ADS has in listing authors for multi-authored, pre-electronic submission papers, the citation information one gathers from the ADS website has more systematic problems within it than that we can get from the ISI website. As with the current ISI website, the current ADS does not permit defining a specific range of years for citing papers.
Comparing numbers of citations for astronomers between the P&A-100 and AST-top-papers lists shows that the P&A-100 list always has more citations, with comparable citations for the most-cited individuals. This indicates that the half-year sampling advantage of the AST-top-papers list is more than compensated by the greater completeness of the P&A-100 list in terms of journals covered. Suprisingly, it is also found that the ADS citations are generally less than those of the P&A-100 list for the same astronomer (at least in late 1999). This difference could be a result of referencing only a restricted number of journals (the ADS currently accesses only astronomy journals for astronomers, as does the AST-top-papers list).
## 3 THE DATA PRODUCTS
### 3.1 The Astronomy Citation Database (ACD)
The Astronomy Citation Database (ACD) is comprised of the names of actual astronomers that are associated with the “names” in the P&A-100 and AST-top-papers. For each associated astronomer name we give the “name” as given in the ISI databases: last name used (ISI format without hyphens), first initial(s); the total number of citations received, the total number of papers cited, and the average citation/paper (cite/pap). Each author of a paper is given full credit for the citations of that paper. First names are given for those astronomers for which first names are known/could be found. The citation data are adjusted for the “A16” issue, the three specified data catalogs (Sec. 2.2), the changed/mispelled names, and for those individuals who use two or more first initials. All additions to the citation data from the P&A-100 list are noted for the P&A-100 “name” commonly used for that astronomer. Note that when the P&A-100 citation data has been modified for a given “name,” the added data are specified as well. If one wishes to see the original P&A-100 data for a given person, subtract the number of citations and number of papers that were added.
For statistical purposes we divide the ACD into three subsets, owing to the degree to which each ISI “name” can be associated with a unique astronomer:
“Unique-One:” These are the astronmers whose last names and first initial(s) are unambiguosly associated with a single, identifiable person. These are most often authors whose names have two or more first initials and those with uncommon last names. The term “unique” here refers to the fact that a unique person is identified with a unique last name and first set of initial(s). There are 4617 “unique-one” astronomers in the ACD.
“Unique-Two:” These are the astronomers whose last names are unambiguously associated with two or more sets of first initials. Again, the use of the word “unique” here denotes association with a unique, identifiable person. While there are 357 “unique-two” astronomers the ACD, the number of astronomers using two or more first initials for their papers is likely much more.
“Confused-Named:” These are the 1484 cases for which unique ownership of a given last name and first initial(s) is intrinsically confused. Here the term “name” literally means a “name,” not a unique person. As stated earlier, name confusion occurs both because a particular last name is confused for the first initial(s) used, or because one set of initials for an individual is confused, but another set is not (the other side of the “unique-two” issue). Of the names that are confused among astronomers and/or physicists, 398 combine one version of the name confused, and the other version uniquely-identified. Hence the same person might be found in two places in the ACD. In addition, there are 94 others pairs of names that have the same last name, two sets of first initials, and both sets of initials/names are confused among those of astronomers and/or physicists. Finally, there are 40 sets in which one name with one first initial can combine with 2 or more names with the same first initial but different middle initials, and all sets of this name are confused among astronomers and/or physicists. Hence, nearly half of the confused-named cases in the ACD could involve confusion among “unique-two”–related issues. (Again, sorting through confused names is confusing!)
To the main ACD we add an “honorable mention” list comprised of the available AST-top-papers citation data for 173 astronomer names from the “A16” analysis whose names were not found in the P&A-100 list (cf. Sec. 2.6).
The full ACD is provided in electronic form from this writer, and will be made available through the Astrophysics Data Center. For the purposes of further statistics, in Table 2 we publish those unique-named and confused-name astronomers cited 3000 times or more. It is possible to separate out those names least-confused for top-cited astronomers in two ways. First, compare the number of citations between the P&A-100 list and the AST-top-papers list. Second, in the case of physicist-related name confusion, use the ADS and the P&A-100 list to check for citations of co-authors on physics papers. These tests have only been applied for those “names” with 3000 or more citations. Astronomers with confused names who have a large ratio of AST-top-papers citations to P&A-100 citations, and are not found to have much physics confusion, are designated as “Cg.” The Cg designation indicates that while the astronomer’s name is confused, a substantial number of citations, if not the vast majority, are for the identified astronomer.
In addition, there are also several possible persons that could be among those with 3000 or more citations, but are not owing to name confusion. These include B/BA Brown, J/JH Lee, J/JW Lee, J/JM Stone, J/JC Wang and R/RW Wilson. The multi-faceted properties of name confusion also affects the number of citations for individuals/confused names already in Table 2, with one confused version of the name qualifying for the list, while the other does not: CL Bennett has C Bennett with 108 citations, M Cohen has MH Cohen with 1862 citations, RF Green has R Green with 740 citations, RH Koch has R Koch with 1494 citations, and DW Murphy has D Murphy with 110 citations.
The data in Table 2 are divided by ACD subset, in descending order of number of citations/name: Column (1) – the last name of the astronomer, no spaces/hyphens, in the ISI standard format for last names; (2) – the first initial(s) of that name (most commonly cited in the case of names in the “unique-two” subset); (3) – the citation/paper ratio (cite/pap); (4) – the number of papers cited, including all papers that can be attributed to the individual or name; (5) – the total number of citations for these papers; (6) – Typ is the sample into which author is placed (1 = “unique-one,” 2 = “unique-two,” C = “confused-named”, and Cg as explained above); (7) – life-time honor awards: Nobel prize, National Academy membership, Russell prize, Heineman Prize, with the last two digits of the year awarded given with the code; (8) – country of citizenship, other than U.S.A., coded as given in Section 4; (9) – comments, including: for astronomers in the “unique-two” subset, the other initial(s) used by the author (including mispellings), and for two names (N/NZ Scoville and R/RS Ellis, the fact they were put into the Unique-2 category is based on comparison of P&A-100 to AST-top-papers citations); for astronomers in the “unique-one” subset, first names and middle initials (if any) for those astronomers using just one first initial, papers added due to mispelling of last name (+sp:); added catalog citations (+ca:); and corrections for the “A16” flaw in the original ISI databases. The “confused-named” subsets give the first names of the astronomers having this/these initial(s) (if the names easily fit into the table), as well as an indication (“+ others”) if there are physicists with the same last name and set of first initial(s).
### 3.2 Top Ten Astronomy Papers Cited, 1981–1996
While the full list of 3200 papers in the AST-top-papers list is an item the ISI has made available for purchase only ($1995), the ISI has given this author permission to print the top 10 papers cited for each year. Table 3 does so for the years 1981–1996 in standard reference format: Column (1) given the number of citations for that paper; (2) the absolute ranking of the paper in the full list of 3200 top–200 cited papers of each year sampled (if number of citations/paper is 90 or over, this is an absolute ranking; if less than this number, no absolute ranking is given, as the AST-top-papers list is not complete for papers with 89 citations or less); (3) the ranking within a given year; (4) the standard astronomical reference for the paper, using the “et al.” for other authors when the number of authors exceeds five (NCC is NUOV CIM C = Nuovo Cimento della Soc. Italiana di Fisica C, Geophysics); (5) the year of publication; (6) the journal in which the paper was published; (7) the journal volume; and (8) the starting page number.
### 3.3 The Honored Few
We in the U.S.A. astronomical community have developed a large number of ways to honor those among us who we feel have done outstanding work in their scientific fields. These awards are separate from that considered the ultimate of our profession — the Nobel Prize. As such, we can go further in our comparison of citations to honors than can likely be done in other fields (cf. the Nobel prize comparison for chemist made by Garfield & Welljams-Dorof 1992). While honors are not awarded solely on the basis of numbers of citations, it is certainly the case that the more our papers are cited, the more likely our work is known to others. As such, one expects a good correlation between the names of those astronomers who have been honored and those most cited. The ACD was partly assembled to test this hypothesis.
All nine astrophysicists to whom the Physics Nobel Prize has been awarded since 1964 are included among the 6458 names in the ACD: CH Townes (1964); HA Bethe (1967): A Hewish (1974); AA Penzias, RW Wilson (1978); S Chandrasekhar, WA Fowler (1983); JH Taylor, RA Hulse (1993). That their Nobel prizes do not correspond to their positions in the ACD owes much more to the fact that for all but the 1993 Nobel prize, the work for which the prize was given well pre-dates the time interval sampled here. Additionally, name confusion affects the citations for JH Taylor and RW Wilson, both of whose names are found in the “confused-named” subset of the ACD, illustrating well that problem.
A total of 278 astrophysicists have been honored either by the Nobel Prize committee (since 1964), and/or the American Astronomical Society and its scientific divisions (since 1949), and/or by current membership in the National Academy of Sciences (in either the Astronomy, Physics or Geophysics sections). While only the lifetime awards are given in Table 2 for individuals so-honored by them, the electronically-available ACD lists all of the awards given by the AAS and its divisions.
Of the 278 individuals honored by the AAS, NAS memberships or Nobel Prizes, 41 are not in the ACD. Of these 41 individuals, 29 received their last honor before 1980, hence most of their papers are likely published before 1981. Three other honored astronomers not in the ACD received only one honor that is more related to public service than to science, per se. Of the remaining nine honored astronomers not in the ACD, it is likely that most of their distinguished work was published before 1981.
## 4 The Statistics of Citations
### 4.1 Unique-One vs. Unique-Two vs. Confused-Named
Three data products are produced from this analysis: the number of citations per astronomer or astronomer “name;” the number of papers cited for these citations; and the ratio of the number of citations to number of papers (cite/pap). Figures 1a,b,c show the histogram distribution for the unique-one, unique-two and confused-name lists for each of these parameters. The data that go into these figures are given in Table 4.
The adopted procedure of handling the number of papers for the unique-two astronomers is verified by the fact that the unique-two authors have a satistically higher cite/pap ratio than do the unique-one authors. This confirms the choice of adding together both papers and citations for unique-two astronomers. In contrast, it is to be expected that the confused-named subset have a signficantly lower cite/pap ratio than the subset of unique-one astronomers.
As shown in Figure 1b, the distribution of the number of citations is nearly identical for the astronomers in the unique-one and confused-named subsets, but skews to signficantly higher values for those in the unique-two subsett. In Figure 1c we see that the numbers of papers cited have the lowest distribution for unique-one astronomers, but are of comparable higher values for the confused-named and unique-two astronomers (note the relatively large number of confused names with 300 or more papers).
Confused-named astronomers have a systematically lower cite/pap ratio than either of the unique-listed astronomers, while having a similar distribution of numbers of citations as those in the unique-one list. This difference was noted early in the analysis and was used to help seek out confused names by investigating the ownership of those names with low cite/pap ratios. The most cited names in the ACD are those astronomers with asian names whose ISI names are multiply-confused, even among astronomers.
The big surprise of this analysis is the finding that the statistics for the unique-two astronomers exceed those for the unique-one astronomers in all three catagories. This difference is also seen when median values for these parameters are compared: The median values for number of citations are 381, 425 and 863 for unique-one, confused-named, and unique-two astronomers, respectively. Similarly, the median values for papers cited are 24, 40 and 38, and for cite/pap are 18.3, 11.5 and 24.2.
Part of the difference in the citation pattern for unique-one authors vs. unique-two authors is consistent with the idea that the median value of citations is likely be twice for people cited in two places in the original list, compared to that for people cited only once. This is true for the numbers of citations, as the median value for unique-two authors (863) is a bit over twice the median value for unique-one authors (381). This, however, does not explain the difference in cite/pap ratio between unique-one and unique-two astronomers. Moreover, examination of the astronomer names in the unique-two subset shows that many are well-known astronomers whose papers are generally highly regarded.
We are forced to conclude that real differences exist between the citation patterns for unique-two astronomers versus those for unique-one astronomers. Why this is so is open to speculation. Is it possible the confidence a person has in giving two sets of initials for her/his papers is related to how useful the papers are for astronomy? Or is this more a product that a person collaborating with other individuals puts his/her on the paper without verifying how her/his name is given? Explaining this result is beyond the scope of this paper, and likely would involve sociologists more than astronomers.
### 4.2 Honors versus Citations
As enumerated earlier (Section 3.3), 237 of the 278 astrophysicists in the honored list are also included in the citation list. Following what was done for the chemist list by Dr. Pendlebury, we concentrate on how the most-cited people in each subset of the ACD (unique-one, unique-two, confused-named) have been honored; i.e., those astronomers listed in Table 2. This list includes 54 astronomers in the unique-one category (top 1.2%); 22 astronomers in the unique-two category (top 6.2%) and 54 astronomers in the confused-name category (top 3.6%), of which 12 can be placed in the Cg category.
The honors received by each person are coded after the citation information by letter code with the last two digits of the date awarded by the code (many of which are seen in Table 2, all of which are used in the full table of honored individuals that will be made electronically available): International: NOB = Nobel Prize; National, NAS = member of the National Academy of Sciences; R = Henry Norris Russell Prize, H = Heinemann Prize.
We will view the relationship between honors and citations from the direction of first selecting on citations. This is because name-confusion fuzzes the numbers to which we might compare for unique-two astronomers versus the other two categories. Until such name confusion is settled, the only valid statistics go from citations to honors, not the other way around. A glance at Table 2 pretty much tells the whole story. While only 1 of 22 Unique-Two astronomers (4.6%), and only 2 of 42 full-confused-named astronomers (4.9%) have been given lifetime honors, 17 of 54 unique-one astronomers (31.5%) and four of 12 Cg astronomers (33.3%) have been so-honored. Similarly, of 59 non-U.S. astronomers, 3 (5.1%) have been so-honored, compared to 21 of 71 U.S. astronomers (29.6%) so-honored. Finally, while 24 of the 130 astronomers/names cited 3000 times or more are so-honored (18.5%), only 78 of the 6328 astronomers/names (1.2%) are so-honored.
The data in Table 2 and in the ACD tell the following story: To first order, the number of citations your papers receive are, indeed, reasonably correlated with whether you are honored by your peers. To second order, whether or not an individual with high numbers of citations has been honored to date is a function of several variables: into which citation subset (unique-one, unique-two or confused) her/his name falls; how many citations that person has; and from what country that person mostly does his/her work. This last effect is understandable, as most of the honors listed in Table 2 are given by the Amercian Astronomical Society and the U.S. National Academy of Sciences.
## 5 Summary and A Modest Proposal
An Astronomy Citation Database (ACD) has been assembled, which gives the number of citations, papers cited and the cite/pap ratio for 6331+ astronomers for a 16.5 year shapshot of time. These data correspond to citation information assembled by the Institute for Science Information (ISI) for astronomers and physicists for papers published in the years 1981-1997.0, as cited in papers published in 1981-1997.5. The data for the astronomy list was assembled from two databases given to this writer by the ISI. One list (P&A-100) contains citation data for 62,813 physics and astronomy ISI “names” cited 100 or more times during the specific intervals. The other list (AST-top-papers) contains both citation data for the 200 most-cited papers published each year in astronomy during 1981-1996, as well as for the astronomy-related citations for 5,035 astronomer names from those papers.
The databases given to this author by the ISI give full citation credit to each author for each of her/his papers, whether first author or not, whether multi-authored or not. This is also the methodology used for ISI’s Web of Science. While the methodology of the ISI in this regard has not changed over the years, how it presents its data has. In particular, users familiar with the hard-copy Science Citation Index will note that only papers on which you are a first author are listed in the Author section, while all papers on which you were an author (first or not) are given the Source section. As such, readers should be aware that the ISI Web of Science now gives full credit for the citations for all papers on which each of us is an author, but only for refereed journals listed from their list of sampled journals (cf. Table 1).
In order to make sense of these citation data, this writer had to engage in a series of laborious, time-consuming tasks. These tasks also involved discovering and correcting for a number of biases in the original data lists, some of which are inherent to any electronically-assembled database.
The main file for ACD is divided into three subsets: 4617 “unique-one” astronomers (those whose names are uniquely identified with individuals); 357 “unique-two” astronomers (those whose last names are cited with two or more sets of first initials); and 1484 “confused-named” astronomers (those with names and initials that are confused with those of other astronomers and/or physicists). The use of the word “unique” to name two of the subsets refers to the fact that these are names singly assigned to an individual astronomer. Such is not the case for the “confused” ISI names, which are associated with two or more individuals in astronomy and/or physics. Due to name confusion, it is likely than many astronomers are listed twice in the ACD, either two times in the confused-named list, or once in the unique-one list and once in the confused-named list.
Two other files are provided with the ACD. One is an “honorable mention” list of 173 names for those astronomers whose names are in the AST-top-papers list but not in the P&A-100 list, and which have less than 100 citations from the AST-top-papers list. The other is the “honors” list, which correlates the citation data information from the main database with the honors 278 individuals have received for their astronomical work.
A list of the 10 most-cited papers per year, from 1981 to 1996 is provided in this paper (but not in the electronic database). Comparison of the authors’ names on the 10 most-cited papers to those names most cited overall shows a good correspondence. Moreover, the papers in the 10 most-cited-per-year list are of a wide range of paper type (e.g., review, data, theory, observation) and are all known to be of high quality. Hence, the old shibboleth that one can get many citations from publising papers with wrong results is shown to be the myth that it is.
The errors for citations for individuals in the ACD are both random and systematic. Random errors exist in proportion (estimated to be 4%) to the number of citations, owing to the way in which the ISI compiles citations. Systematic errors include errors of omission and comission. Errors in citations owing to variance in self-citations among astronomers can also affect the statistics at the $`6\%`$ level, if one chooses to apply such a criterion to these data. The likely most-egregious error in this database in the eyes of some of the readers of this paper is the use of full citations for each author of a paper. Tests using the papers of top-cited authors in the AST-top-papers database shows that no way we can think of to calculate citations for authors will give the same results. The plain fact is, however, whether one agrees with using full citations or not for astronomers, this is what the ISI gives us in the Web of Science, as well as in the P&A-100 and AST-top-papers lists.
We show that, due to the manner in which citations are assembled by the Institute for Science Information, which maintains the science citations for our field of study, we hurt the impact of our papers if we publish many papers in meeting proceedings, or if our name is confused with those of others.
Two more of our findings are of specific sociological interest. The first is that astronomers who divide their papers among two or more first sets of initials on average publish more papers, have more citations, and have a significantly higher citations/paper ratio than astronomers who publish under just one set of initials.
The second is that there is a very good, but not perfect, correspondence of a person being near the top of the citation list and the chance that person has been honored by her/his peers in U.S. astronomy for life-time achievement. Closer examination of this good correlation reveals that assigning honors to our peers is as human an enterprise as any that we do. The evidence asuggests that anything we do, intended or not, to blur the focus of our papers to our peers has a measurable effect on how we are honored by them. Such blurring effects can come in several forms: authoring papers using two or more sets of initials; having a name confused with those of others; publishing your papers predominantly in non–U.S.–based journals or meeting proceedings. Other means of blurring the focus of our papers (having high citations, but low cite/pap ratio; mostly publishing in large groups or with more well-known authors; publishing in several different scientific fields) also likely exist.
In a paper devoted to analyzing the citation data for astronomers, it is relevent to note that citation statistics give us but one view of the impact of individuals on our science. This is best evidenced by the most recent Nobel Prize awarded to astronomers, that to JH Taylor and RA Hulse in 1993. Where these two individuals stand in the citation list has little relationship to their scientific impact on our field.
The two ISI databases used in this paper were generated with ISI software and given to this author expressly for the analysis done in this paper. These databases are different from those we can access via the Web in three ways: ability to specify specific ranges of years for cited and citing papers; kinds of papers cited, and number of citations attributed to each author. The idiosyncracies discovered in the course of this survey of the use of current web-based database for the bibliography of astronomers are detailed.
The full ACD (all three files) will be made available via anonymous ftp from samuri.la.asu.edu, as well as through the Astrophysics Data Center. The top ten papers cited for each year between 1981-1996 are given in this paper. The full top-200 cited papers database produced by the ISI is proprietary (termed “High-Impact Papers in Astronomy”) and is available for a fee separately from the ISI.
Since the ACD is a database about the accomplishments of people, any error in the ACD is a serious error. As such, readers are encouraged to contact this writer if errors of omission or comission are found in the ACD. Those appropriate modifications to the ACD that should be made, will be made. It is hoped that with the help of readers, a full unconfused, ACD can be eventually made of cited papers during 1981–1997.5. Such a list can then act as a baseline against which future investigations of this kind may be made. It also hoped that the lessons learned in this paper about the idiosyncracies of the various electronic databases will aid others in their own searches.
Towards this end, I end this paper by unashamedly borrowing from Jonathan Swift in putting forward a “modest proposal” for eliminating name confusion in our field. We are used to having social security numbers, university ID numbers, shopper ID numbers; each of us is now various numbers in various databases. If this is so, why do we not assign a “publishing ID number” (PID for short) for each person who publishes a paper in our journals? I suggest that this be done, and we start by assigning PIDs to all astronomers who have published papers in the journals in the past. The ACD could be a start, but only a start, as to solve name confusion the process must work iteratively among the ISI, the individuals involved and the ACD. The PID number would then be carried by each journal in the author list (but not necessarily printed out for each paper).
If we, and other scientists in other fields of study, are interested in having a true, honest assessment of citation data for astronomers as a function of time, then the problem of name confusion should be, and can be, solved.
This work could not have been done without the active participation and cooperation of Dr. David Pendlebury of the ISI, who supplied the two ISI datasets used in this study, and who also did a very careful reading of the first draft of this paper. This author owes Dr. Pendlebury much thanks. Conversations with Helmut Abt, Sandra Faber and Anne Cowley were also helpful in writing this paper. My thanks to Arnab Choudhuri for pointing me in the direction of the online Astronomy Society of India directory. Insightful comments from the editor, Bob Milkey, also greatly helped the presentation of this paper.
Figure Caption
Figure 1: The histograms for the average citations/paper (cite/pap) for the unique-one (4617 names), unique-two (357 names) and confused-named (1484 names) subsets in the Astronomy Citation Database. a) The cite/pap ratio is binned by 5, and the histogram values are in terms of the fraction of the astronomers in each subset, falling in each cite/pap bin. Note the higher median value of cite/pap, and longer tail towards higher values of cite/pap for the unique-two astronomers as opposed to the unique-one astronomers. b) Numbers of citations are binned by 100. Histogram values as in Figure 1a. Note, as with Figure 1a, the higher median value of citations, and longer tail towards higher numbers of citations for the unique-two astronomers as opposed to the unique-one astronomers. c) Numbers of papers cited for the unique-one, unique-two, and confused names in our astronomer sample. The number of papers cited is binned by 10, and the histogram values are in terms of the fraction of the astronomers in each list, falling in each paper-cited bin. Note the higher median value of numbers of papers cited for the unique-two astronomers relative to those for the unique-one astronomers, and the longer tails towards high values of papers cited for both unique-two and confused-named astronomers.
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# Low Temperature Anomaly in Mesoscopic Kondo Wires
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## Abstract
We report the observation of an anomalous magnetoresistance in extremely dilute quasi-one-dimensional AuFe wires at low temperatures, along with a hysteretic background at low fields. The Kondo resistivity does not show the unitarity limit down to the lowest temperature, implying uncompensated spin states. We suggest that the anomalous magnetoresistance may be understood as the interference correction from the accumulation of geometric phase in the conduction electron wave function around the localized impurity spin.
\]
A localized magnetic moment interacts with the conduction electrons in a metal resulting in a logarithmic increase of the resistivity as the temperature is lowered. This is known as the Kondo effect . Below the Kondo temperature, $`T_K`$, an electron cloud begins to screen the impurity until its spin is completely compensated, forming a singlet state at low temperature. The nature of this state and the extent of the screening cloud has been studied for decades. Recently this effect has been explored in mesoscopic systems in an attempt to understand whether the screening is affected by the finite sample size , including high temperature large concentration experiments on layered Kondo systems , and 2D films . Interference effects in mesoscopic Kondo systems containing impurity concentrations $`c>50\text{ ppm}`$ do not generally contribute significantly to the measured magnetoresistance or resistivity because of the strong suppression of long range phase coherence due to spin-flip scattering. In spite of its relevance to mesoscopic systems, a complete study of the low temperature magnetoresistance in very dilute alloys ($`c<10\text{ ppm}`$), where the Kondo screening length is comparable to the phase coherence length $`L_\varphi `$, has not been done. In this regime, an interference experiment which can reveal new information on the development of the Kondo screening cloud is possible. The three-dimensional character of the local dipolar magnetic field from the impurity spin coupled with an externally applied field should provide an additional interference contribution to the electron wavefunction. This is analogus to the Berry phase effect predicted for coherent electrons in a ring traversing in an externally applied 3D magnetic field texture concentric with the ring.
In this paper, we report the magnetoresistance and the temperature dependence of the resistivity down to 38 mK for five quasi-1D AuFe wires in the concentration range of $`3<c<10`$ ppm. We determine both the spin-flip scattering rate and the phase decoherence rate by fitting the low field magnetoresistance to standard weak localization theory . We find that the unitarity limit corresponding to the formation of the singlet state is not yet reached at our lowest temperature in spite of the fact that AuFe Kondo systems are known to have a Kondo temperature of 1 K. At intermediate fields we observe a negative magnetoresistance that is characteristic in temperature dependence and shape of an interference correction, and different from the expected standard Kondo magnetoresistance. At low temperatures this magnetoresistance shows hysteresis which vanishes if the magnetic field is swept to a larger value or if the temperature is increased. We argue that our data is not consistent with a spin glass model but rather with a new interference correction similar to a Berry phase effect .
Our studies are done on pure (99.9995$`\%`$) samples of gold (Au) before, and after, the ion implantation of 3 to 10 ppm of iron (Fe) impurities. This provides a clear advantage over earlier works on layered or flash-evaporated samples in that the contribution to the magnetoresistance at various field scales coming solely from the magnetic impurities could be easily identified. Sample dimensions, diffusion constant $`D`$, and $`L_\varphi `$ measured after implantation are given in Table 1. These samples are quasi-1D, since $`w,tL_T,L_\varphi `$, where $`L_T=\sqrt{\mathrm{}D/k_BT}`$ is the thermal diffusion length. The Kondo contribution to the resistivity $`\mathrm{\Delta }\rho (T)`$ is found to have the expected logarithmic increase: $`\mathrm{\Delta }\rho (T)=AB\mathrm{ln}(T)`$ (See Fig. 1), after the subtraction of the electron-electron interaction(EEI) contribution measured before the ion implantation, which has the expected theoretical value, $`\mathrm{\Delta }\rho _{ee}(2e^2R^2wt/hL^2)L_T`$.
Total scattering rate $`1/\tau `$ relevant for resistance is $`1/\tau =1/\tau _n+1/\tau _s`$; $`1/\tau _n`$ is the nonmagnetic scattering rate. Phase-breaking rate $`1/\tau _\varphi `$ in the presence of magnetic scattering is given by $`\tau _\varphi ^1=2\tau _s^1+\tau _{\varphi (nonmag)}^1`$. Fig. 2 displays the temperature dependence of the magnetic scattering rate $`1/\tau _s`$ obtained from WL measurements for samples AuFe1 and AuFe2. $`1/\tau _s`$ is obtained from WL after subtracting the inelastic rate $`1/\tau _{\varphi (nonmag)}`$ due to nonmagnetic sources, measured in the same Au wires before ion implantation. The $`1/\tau _\varphi `$ correction term does not produce the observed behavior seen in Fig. 2 because $`1/\tau _s`$ is much larger than $`1/\tau _\varphi `$ in the corresponding clean system. The maxima near 0.2 K-0.4 K represent the previously observed resonant spin-flip scattering processes .
As shown in Fig. 1, the unitarity limit is not reached down to 40 mK, even in the presence of disorder and a finite magnetic field required to quench WL, both of which should help form the singlet state. This is consistent with earlier observations . The impurity spin is thus not completely screened. However, at a larger magnetic field, a resistivity plateau is observed corresponding perhaps to the unitarity limit (See Fig. 3(a)). The plateau shifted to higher temperatures with increasing magnetic field. Additionally, we observed a maximum around $`T_K`$ (See Fig. 3(b)). This observation is consistent with earlier experiments on (LaCe)Al<sub>2</sub> and (LaCe)B<sub>6</sub> , consequently explained by a wave description of the spin density . This implies that there is a substantial spin polarization around the impurities with a potential $`V(r)=V_0\mathrm{cos}(2k_Fr)/r^3`$. The local magnetic field
of the spin polarization can be on the order of a Tesla within a couple of nanometers from the impurity, though it is negligible on the scale of the typical inter-impurity distance of $``$ 10 nm. Strength of this potential $`V_0`$ is experimentally known to be very large for AuFe, decreasing exponentially with increasing concentration $`c`$ . Thus, there are strong local magnetic fields for purer samples with longer $`L_\varphi `$. NMR measurements of the conduction-electron spin density around Fe atoms in a Cu matrix also find a nonvanishing radial component above and below $`T_K`$ .
That there exists a distribution of local magnetic fields from the impurity spins is further confirmed by the observation of hysteresis in the low-field MR. As shown in Fig. 4, the background of the WL curve is asymmetric with a positive or negative slope depending on the field history. Hysteresis disappears at high temperatures, typically between 0.4 K and 1.5 K depending on the sample. In contrast to what is observed in a spin glass, we find this hysteresis to be stronger for systems with longer $`L_\varphi `$ (hence for lower concentration samples). Hysteresis is expected for a spin glass system below $`T_g`$; so if it were a spin glass, we would have observed stronger hysteresis for higher concentration samples, contrary to our data. Our experiment suggests that hysteresis arises because of different realizations of the three-dimensional local field distribution. As the sample gets cold, impurity spins freeze out in random orientations, providing a particular configuration for the local-field distribution. This distribution is modified by a magnetic field due to spin alignment. Magnetic field cycling between $`\pm `$ 1 Tesla removes the hysteresis and flattens the background of the low-field MR, while cycling between $`\pm `$ 0.05 Tesla does not.
All our samples are in the single-impurity regime and the logarithmic increase of resistivity scales with concentration. It is unlikely that these systems behave like a spin glass for a number of reasons: (a) In AuFe, spin glass behavior is not observed for $`c100\text{ ppm}`$, as is well known; (b) Second, spin glass temperature $`T_g`$
for a system with 3-10 ppm Fe in Au would be 1 mK or lower; (c) The resistivity maximum expected for a spin glass is also not seen in Fig. 1; and, (d) The possibility of inhomogeneous pockets of impurities, or clustering, is ruled out by measuring different segments in a sample. The observed behavior is found to be independent of the choice of the segment, suggesting a homogeneous mechanism. For these reasons, the spin glass formation can be ruled out.
High-field magnetoresistance of a representative sample, AuFe4, is shown in Fig. 5(a). WL is observed at a field scale of $`B<0.03\text{ T}`$. At high fields, due to the cyclotron orbits of the electrons a classical MR is expected: $`\mathrm{\Delta }R_c/R(\omega \tau _e)^2`$, with $`\omega eB/m`$, and $`\tau _e`$ being the electron mean free path. This classical $`B^2`$ dependence is displayed in Fig. 5(b), which is subtracted out in Fig. 5(a). At the intermediate field scale $`(1\text{ T})`$, we observe a negative magnetoresistance in all our samples at $`T<T_K`$ that is very sensitive to temperature. Theoretically, in the standard Kondo model, one expects a negative MR due to the suppression of the spin-flip scattering by the alignment of the spins with the field: $`\mathrm{\Delta }R_2/R(gS\beta )^2(H/T)^2`$, where $`\beta `$ is the Bohr magneton. The data is not described by this contribution, as evident in the shape of the MR at various temperatures. We have observed this anomalous MR in all our samples along with the WL dip at zero field. At 40 mK, the conductance change, $`\mathrm{\Delta }G=\mathrm{\Delta }R/R^2`$, for all our samples in units of $`e^2/h`$ is : $`0.001,0.002,0.018,0.028,`$ and 0.004 for samples AuFe1 through AuFe5 respectively.
Earlier experiments on higher concentration AuFe samples revealed a behavior compatible with the standard expected form, and different from what we observe. Above $`T_K`$, the standard high-field magnetoresistance is essentially a function of the thermal average of the local moment in the field direction $`S_Z`$. As temperature is increased, the field scale increases with the height of the MR decreasing, ultimately becoming flat at a very high temperature due to thermal fluctuations of the localized spin. This behavior is observed in 2D
Kondo films of AuFe at 1.4 K and 4.2 K . However, in another experiment on AuFe wires with a much higher concentration of Fe impurities ($`50`$ ppm), temperature dependence of the MR was not studied. There are two important characteristics of our low temperature MR, different from the bulk Kondo behavior. First, the magnetoresistance as a function of temperature cannot be explained by $`S_Z`$, since the field scale is expected to grow with increasing temperature while conserving area under the curve. Second, $`S_Z`$ as a function of temperature is expected to increase with decreasing temperature, becoming flat at low temperatures, whereas the dependence shown in Fig. 5(a) displays no saturation down to 40 mK. High concentrations of impurities in the earlier experiment on AuFe wires and high temperature range in the experiments on 2D AuFe films imply a very short $`\tau _\varphi `$ in these systems, yielding non-mesoscopic bulk behavior. The local magnetic field due to polarization in these high concentration samples is expected to be extremely weak, in contrast to our samples.
It is clear from our resistivity and scattering rate measurements that the long range polarization of the conduction electrons around the localized spin is effective at low temperatures for our low concentration mesoscopic systems. From our observation of hysteresis, we believe that this polarization or the local magnetic field causes the anomalous high-field magnetoresistance. Furthermore, the shape of the magnetoresistance and its temperature dependence are very much similar to what is expected from a quasi-1D interference effect, which suggests a similarity to weak localization. These effects were seen in long $`L_\varphi `$ samples, implying an essential role played by the phase coherence of electrons. Considering all this, we propose a connection of this new interference correction to Berry phase.
It is possible for the phase coherent mesoscopic Kondo wires to show a weak-localization-like magnetoresistance driven by a geometric phase $`\mathrm{\Gamma }=_{t_i}^{t_f}𝐀_g𝑑𝐑`$, similar to the standard weak localization driven by the Aharonov-Bohm phase . $`𝐀_g`$ is the geometric gauge potential, and $`𝐑`$ is the position vector describing the tip of the spin. The spin part of the wave function of the phase coherent electron picks up a geometric phase as it aligns along the local magnetic field of the uncompensated spin. This is further helped by disorder in the sample , since the electron spends more time around the spin than it would in a ballistic sample. The corresponding geometric phase is equal to half of the solid angle subtended by the area enclosed by the tip of the electron spin vector due to its evolution in a closed loop. A complementary path, going in the opposite direction, will contribute an opposite phase shift. Interference of two such paths around the local field results in a correction to conductivity, analogous to the anticipated Berry phase correction in a ring structure. There are no oscillations as in the Aharonov-Bohm effect, but just half a period in resistivity, because the maximum Berry phase acquired is $`\pi `$, half of the maximum solid angle of $`2\pi `$. An externally applied perpendicular field aligns the electron spin. If the spin is completely aligned along the external field, the solid angle subtended by the tip of the spin is zero, resulting in the complete suppression of the Berry phase correction.
Berry phase changes sign under time-reversal. This leads to a contribution similar to the Cooperon propagator in WL. Correction to the resistance contains the disorder average of all possible loops acquiring Berry phase. As temperature is increased, $`L_\varphi `$ (which includes spin fluctuations) reduces greatly, thus reducing the magnetoresistance correction as seen in Fig. 5(a). This dependence is similar to that of WL. In the spirit of WL, a geometric length $`L_B`$ can be introduced, which is the length scale over which the net accumulated geometric phase is on the order of $`\pi `$. $`L_B`$ may be defined by $`\mathrm{\Delta }R_g/R^2(e^2/\mathrm{})L_B/L`$. For the data from the sample AuFe4( shown in Fig. 5(a)) at 40 mK, the geometric length $`L_B18\mu m`$ ($`L_\varphi 3\mu m`$ at 40 mK), implying that within $`L_\varphi `$ the acquired (disorder-averaged) geometric phase is on the order of $`\pi L_\varphi /L_B\pi /6`$ for this sample.
To summarize, we have observed an unusual temperature dependence of the magnetoresistance along with hysteresis in quasi one-dimensional disordered Kondo systems at $`T<T_K`$. We believe that this arises from the adiabatic evolution of the phase coherent electron around the impurity spin, which results in a Berry phase effect. We thank B. Altshuler, H. Fukuyama, D. Loss, P. Schwab, J. Schwarz, and A. Zawadowski for conversations. This work is supported by the NSF (DMR9510416) and the ARO (DAAG559710330).
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# Prym varieties and the canonical map of surfaces of general type
## 1 Introduction
Let $`X`$ be a smooth surface of general type and let $`\varphi :X\mathrm{\Sigma }𝐏^{p_g(X)1}`$ be the canonical map of $`X`$, where $`\mathrm{\Sigma }`$ is the image of $`\varphi `$. Suppose that $`\mathrm{\Sigma }`$ is a surface and that $`\varphi `$ has degree $`\delta 2`$. Let $`ϵ:S\mathrm{\Sigma }`$ be a desingularization of $`\mathrm{\Sigma }`$. A classical result, which goes back to Babbage \[Bab\], and has been more recently proved by Beauville, \[B2\] (see also \[Cat1\]), says that either $`p_g(S)=0`$ or $`S`$ is of general type and $`ϵ:S\mathrm{\Sigma }`$ is the canonical map of $`S`$. In the latter case we have a dominant rational map $`\psi :XS`$ of degree $`\delta `$, which we call a good canonical cover of degree $`\delta `$ (see definition 2.3 for a slightly more general definition).
While there is no problem at all in exhibiting as many examples as one likes of the former type, i.e. where $`p_g(S)=0`$ (see \[B2\]), not so many good canonical covers are available in the current literature. In few sporadic examples of such covers the surface $`X`$ is regular (see \[VdGZ\], \[B2\] proposition 3.6, \[Cat1\] theorem 3.5, \[C1\], \[Pa2\]). On the other hand, there is an interesting construction, due to Beauville (see \[Cat2\], 2.9 and \[MP\]), which produces an infinite series of such covers of degree $`2`$ where $`X`$ is irregular, precisely of irregularity $`2`$. Beauville’s construction is recalled in §4 and in example 3.1. The resulting canonical covers have been extensively studied in \[MP\], where they have been classified in terms of their birational invariants.
In our attempts to find more examples of canonical covers, we have been lead to understand Beauville’s construction better. In particular we extracted from it its main features, and this lead us to give a definition, the one of a good generating pair (see 2.4 for a more general definition), which, roughly speaking, is the following. A good generating pair $`(h:VW,L)`$ is the datum of a finite morphism $`h:VW`$ of degree $`2`$ between surfaces, $`V`$ smooth and irreducible, $`W`$ with isolated double points of type $`A_1`$, and $`L`$ a nef and big line bundle on $`W`$. Furthermore one requires that $`|L|`$ has at least dimension $`1`$ and contains a smooth, irreducible, non–hyperelliptic curve $`C`$, that $`h^{}K_W=K_V`$ (this means that $`h`$ has only isolated ramification points, corresponding to the double points of $`W`$) and that the pull-back of the adjoint linear system $`|K_W+L|`$ is the complete linear system $`|K_V+h^{}L|`$.
However cumbersome and un–motivated this definition may appear at a first glance, it turns out to be rather useful for constructing canonical covers. Indeed one finds many of these in the following way (see §3 for details). Consider the map $`\stackrel{~}{h}=h\times Id:V\times 𝐏^1W\times 𝐏^1`$ and the projections $`p_i`$, $`i=1,2`$, of $`W\times 𝐏^1`$ onto the two factors. A general surface $`\mathrm{\Sigma }|p_1^{}Lp_2^{}𝒪_{𝐏^1}(n)|`$, $`n3`$, has only points of type $`A_1`$ as singularities. We set $`X=\stackrel{~}{h}^{}(\mathrm{\Sigma })`$, $`\varphi =\stackrel{~}{h}|_X:X\mathrm{\Sigma }`$, $`ϵ:S\mathrm{\Sigma }`$ the minimal desingularization, $`\psi =ϵ^1\varphi `$. Then, using adjunction both on $`V\times 𝐏^1`$ and $`W\times 𝐏^1`$, one sees that $`\psi :XS`$ is a good canonical cover of degree $`2`$.
General properties of generating pairs $`(h:VW,L)`$ are studied in §5 (see also §8, where some information about higher degree generating pairs has been collected). In particular, we see that $`V`$ and $`W`$ have the same Kodaira dimension (see proposition 5.5), and, while $`W`$ is always regular, $`V`$,instead, is irregular, and its irregularity can be expressed in terms of the genus $`g`$ of the general curve $`C|L|`$ and of the degree of $`h`$ (see proposition 5.4). We also give formulas for the invariants of the canonical covers arising from a given generating pair (see proposition 2.7).
It is interesting to notice that Beauville’s example is essentially characterized by the fact that $`V`$ and $`W`$ have Kodaira dimension $`\kappa =0`$ (see proposition 8.2 for a more precise statement). The case of Kodaira dimension $`1`$ is also rather restricted, as proposition 8.3 shows.
Beauville’s example corresponds to the case in which $`V`$ is a principally polarized abelian surface, $`W`$ is its Kummer surface, and $`L`$ is the polarization on $`W`$ which lifts to a symmetric principal polarization on $`V`$. Unfortunately, more generating pairs do not easily show up. The only ones which we know about are listed in section §3. These give rise to more infinite series of good canonical covers which wait for a deeper understanding, like, as we said, in \[MP\] has been done for Beauville’s examples.
The difficulty in finding generating pairs is not casual. This is explained in §6, and this is where Prym varieties come into the picture. If $`(h:VW,L)`$ is a generating pair, and $`C`$ is a general curve in $`|L|`$, of genus $`g`$, then $`h^{}C=C^{}C`$ is an unramified double cover, with a related Prym variety $`Prym(C^{},C)`$. In theorem 6.1 we prove that $`Prym(C^{},C)`$ is naturally isomorphic to the Albanese variety of $`V`$. As a consequence we find that, if the generating pair is good, the Albanese image of $`V`$ is a surface and therefore the Kodaira dimension of $`V`$ and $`W`$ is non–negative (see corollary 6.2). Moreover, some general facts about irregular surfaces and isotrivial systems of curves on them, which have been collected in §4, imply that the Prym map has an infinite fibre at the cover $`h:C^{}C`$ (see proposition 6.6). This, together with results about the fibre of the Prym map due to several authors (see §6 for references), enable us to prove that, if $`(h:VW,L)`$ is a good pair, then one has the bounds $`g12`$ for the genus $`g`$ of $`C`$ and is $`q11`$ for the irregularity $`q`$ of $`V`$ (see theorem 6.9 and proposition 6.11). We suspect that, along the same lines, it should be possible to improve this bound for $`g`$ and $`q`$, but this would preliminarly require a deepening of our understanding of the fibres of the Prym map. For instance we would like to know answers to questions like: when may these fibres contain rational curves? Problems, of course, of independent interest.
Finally, using Reider’s method, we obtain the bound $`L^24`$ (see proposition 7.3), so that one really sees why there are not so many possibilities for a good generating pair. We give a complete classification of good pairs with $`L`$ ample and $`h^0(L,W)>2`$.
These satisfy $`h^0(W,L)4`$ and $`L^2=3,4`$. The only example with $`h^0(W,L)=4`$ is Beauville’s one (see corollary 7.4). The cases $`h^0(W,L)=3`$ and $`L^2=3`$ or $`L^2=4`$ but $`|L|`$ with a base point are studied in §7. (see theorem 7.10); we find that the former case corresponds either to example 3.3 or to a suitable modification of Beauville’s example, while the latter does not occur. These cases share the feature that the general curve $`C`$ is trigonal, and we take advantage, in the proof of our classification theorem 7.10, of a globalization to $`V`$ of the well known trigonal construction (\[Ca2\]), which is the inverse of the equally famous Recillas’ construction. A different proof of the same result is sketched in remark 7.9.
Finally we prove that $`L^2=4`$, $`h^0(W,L)=3`$ does not occur (see corollary 7.7). In the pencil case $`h^0(W,L)=2`$ (in which there are examples, like 3.2, but a classification is still lacking) we show that the possibility $`L^2=4`$ is severely restricted (see corollary 7.7).
Using similar ideas, we are able to construct an infinite family of good canonical covers with $`X`$ regular. We will be back on this in a forthcoming paper.
Notation and conventions: all varieties are defined over the field of complex numbers. A map between varieties is a rational map, while a morphism is a rational map that is regular at every point. We do not distinguish between Cartier divisors and line bundles and use the additive and multiplicative notation interchangeably. The Kodaira dimension of a variety $`X`$ is denoted by $`\kappa (X)`$. We denote by $`_{num}`$ the numerical equivalence between divisors on a smooth surface.
## 2 Canonical covers and generating pairs
###### Notation 2.1
Let $`S`$ be a surface with canonical singularities, i.e. either smooth or with rational double points, so that in particular $`S`$ is Gorenstein. We denote by $`K_S`$ the canonical divisor of $`S`$, and we let $`p_g(S)=h^0(S,K_S)=h^2(S,𝒪_S)`$ be the geometric genus and $`q(S)=h^1(S,𝒪_S)`$ the irregularity. If $`p_g(S)2`$, the canonical map of $`S`$ is the rational map $`\varphi :S𝐏^{p_g(S)1}`$ defined by the moving part of the canonical system $`|K_S|`$ of $`S`$. If $`S_0`$ is the open set of smooth points of $`S`$ and $`ϵ:S^{}S`$ is any desingularization, then $`p_g(S)=p_g(S^{})`$ and $`q(S)=q(S^{})=h^0(S^{},\mathrm{\Omega }_S^{}^1)=h^0(S_0,\mathrm{\Omega }_{S_0}^1)`$. The Albanese map of $`S^{}`$ factors through $`ϵ`$, since the exceptional locus of $`ϵ`$ is a union of rational curves, and so we can speak of the Albanese map of $`S`$.
Let $`X`$ be a smooth surface of general type and let $`\varphi :X\mathrm{\Sigma }𝐏^{p_g(X)1}`$ be the canonical map of $`X`$, where $`\mathrm{\Sigma }`$ is the image of $`\varphi `$. We assume that $`\mathrm{\Sigma }`$ is a surface and that $`\varphi `$ has degree $`d2`$, and we denote by $`ϵ:S\mathrm{\Sigma }`$ a desingularization of $`\mathrm{\Sigma }`$. We recall the following theorem due to Beauville, \[B2\], Thm. 3.4.
###### Theorem 2.2
Under the above assumptions, either:
(i) $`p_g(S)=0`$ or;
(ii) $`S`$ is of general type and $`ϵ:S\mathrm{\Sigma }`$ is the canonical map of $`S`$.
We introduce some terminology for surfaces verifying condition (ii) of Theorem 2.2:
###### Definition 2.3
Let $`X,S`$ be smooth surfaces of general type. Let $`\psi :XS`$ be a dominant rational map of degree $`d2`$. Assume that:
* $`p_g(X)=p_g(S)`$;
* the canonical image of $`S`$ is a surface $`\mathrm{\Sigma }`$.
In this case the canonical map $`\varphi :X\mathrm{\Sigma }`$ of $`X`$ is the composition of $`\psi `$ and the canonical map $`ϵ:S\mathrm{\Sigma }`$ of $`S`$, and we say that $`\psi :XS`$ is a canonical cover of degree $`d`$. If $`ϵ:S\mathrm{\Sigma }`$ is birational, then we say that the canonical cover is good.
A few sporadic examples of canonical covers are available in the literature (\[VdGZ\],\[B2\] prop. 3.6, \[Cat1\] thm. 3.5, \[C1\], \[Pa2\]). However, so far, there is only one construction, due to Beauville (see \[Cat2\], 2.9 and \[MP\]), which produces an infinite series of such covers. We recall it next.
Let $`V`$ be a principally polarized abelian surface such that the principal polarization $`D`$ is irreducible, and let $`h:VW`$ be the quotient map onto the Kummer surface $`W=V/<1>`$. The surface $`W`$ can be embedded into $`𝐏^3`$ as a quartic surface via a complete linear system $`|L|`$ such that $`h^{}|L|=|2D|`$. Consider the map $`\stackrel{~}{h}=h\times Id:V\times 𝐏^1W\times 𝐏^1`$ and the projections $`p_i`$, $`i=1,2`$, of $`W\times 𝐏^1`$ onto the two factors. A general surface $`\mathrm{\Sigma }|p_1^{}Lp_2^{}𝒪_{𝐏^1}(n)|`$, $`n3`$, has only points of type $`A_1`$ as singularities. We set $`X=\stackrel{~}{h}^{}(\mathrm{\Sigma })`$, $`\varphi =\stackrel{~}{h}|_X:X\mathrm{\Sigma }`$, $`ϵ:S\mathrm{\Sigma }`$ the minimal desingularization, $`\psi =ϵ^1\varphi :XS`$. Then it is easy to check, using adjunction both on $`V\times 𝐏^1`$ and $`W\times 𝐏^1`$, that $`\psi :XS`$ is a good canonical cover of degree $`2`$.
We wish to study to what extent this construction can be generalized. We introduce a class of pairs $`(h:VW,L)`$, where $`h:VW`$ is a finite morphism of surfaces and $`L`$ is a line bundle on $`W`$, in such a way that by applying the above construction to $`(h:VW,L)`$ one gets an infinite series of canonical covers.
###### Definition 2.4
Consider a pair $`(h:VW,L)`$, where $`h`$ is a finite morphism of degree $`d2`$ between irreducible surfaces, $`V`$ smooth, $`W`$ with at most canonical singularities and $`L`$ is a line bundle on $`W`$, such that:
* $`K_V=h^{}K_W`$;
* $`h^0(W,L)2`$ and $`L`$ is big, i.e. $`L^2>0`$;
* the general curve $`C`$ of $`|L|`$ is smooth of genus $`g2`$ and the curve $`C^{}:=h^{}C`$ is not hyperelliptic;
* $`p_g(V)=p_g(W)`$, $`h^0(V,K_V+h^{}L)=h^0(W,K_W+L)>0`$.
We call $`(h:VW,L)`$ a degree $`d`$ and genus $`g`$ generating pair of canonical covers, and we denote by $`L^{}`$ the line bundle $`h^{}L`$ on $`V`$. The pair is said to be minimal if both $`V`$ and $`W`$ are minimal.
The generating pair is called good if the general $`C`$ of $`|L|`$ is not hyperelliptic (hence $`g3`$ in this case).
Notice that condition (GP1) is equivalent to the fact that $`h`$ is ramified only over the singular points of $`W`$. Condition (GP3) and Bertini’s theorem imply that the general curve $`C`$ in $`|L|`$ is smooth and irreducible, hence $`L`$ is also nef, i.e. $`LD0`$ for every effective divisor $`D`$ on $`V`$. The assumption that $`C^{}`$ is not hyperelliptic is a technical condition whose meaning will be clearer later (cf. for instance theorem 6.1). Finally, the base points of $`|L|`$, if any, are smooth points of $`W`$.
In the rest of this section we show that by applying the original construction of Beauville, to a (good) generating pair one obtains an infinite series of (good) canonical covers, and we compute the invariants of such canonical covers. In order to do this, we need the following result, that will be proven later (cf. proposition 5.4):
###### Proposition 2.5
If $`(h:VW,L)`$ is a generating pair, then $`q(W)=0`$.
We introduce now some more notation:
###### Notation 2.6
Given a generating pair $`(h:VW,L)`$ of degree $`d`$ and genus $`g`$, we denote by $`p_i,i=1,2`$, the projections of $`W\times 𝐏^1`$ onto the two factors and we write $`\stackrel{~}{h}=h\times Id:V\times 𝐏^1W\times 𝐏^1`$. We denote by $`(n)`$ the line bundle $`p_1^{}Lp_2^{}𝒪_{𝐏^1}(n)`$, where $`n`$ is a positive integer.
In addition, we let $`\mathrm{\Sigma }|(n)|`$ be a general surface, $`Y=\stackrel{~}{h}^{}(\mathrm{\Sigma })`$. We denote by $`ϵ:S\mathrm{\Sigma }`$ and $`ϵ^{}:XY`$ the minimal desingularizations, by $`f`$ the map $`\stackrel{~}{h}|_X:X\mathrm{\Sigma }`$, and by $`\psi `$ the map $`ϵ^1fϵ^{}:XS`$.
###### Proposition 2.7
We use notation 2.6.
Let $`(h:VW,L)`$ be a generating pair of degree $`d`$ and genus $`g`$. If $`n3`$, then one has:
i) let $`\mathrm{\Sigma }|(n)|`$ be general and let $`Y=h^{}\mathrm{\Sigma }`$; $`\mathrm{\Sigma }`$ and $`Y`$ are surfaces of general type with at most canonical singularities. If in addition $`L`$ is ample, then $`S`$ and $`X`$ are both minimal;
ii) $`\psi :XS`$ is a canonical cover of degree $`d`$, that is said to be $`n`$–related to the generating pair $`(h:VW,L)`$. If the generating pair is good, then $`\psi :XS`$ is a good canonical cover, while if the generating pair is not good then the canonical map of $`S`$ is $`2`$-to-$`1`$ onto a rational surface;
iii) the invariants of $`S`$ are: $`p_g(S)=np_g(W)+(n1)g`$, $`q(S)=0`$, $`K_S^2=n(K_W^2L^2)+8(n1)(g1)`$;
iv) the invariants of $`X`$ are: $`p_g(X)=p_g(S)`$ and $`q(X)=(d1)(g1)`$, $`K_X^2=dK_S^2`$.
Proof: Recall that by condition (GP2) of definition 2.4 the general curve of $`|L|`$ is smooth, and thus, in particular, $`|L|`$ has no fixed components. Thus also the linear system $`|(n)|`$ has no fixed components and is not composed with a pencil. Therefore its general member $`\mathrm{\Sigma }`$ is irreducible. Moreover the set of base points of $`|(n)|`$ is the inverse image via $`p_1`$ of the set of base points of $`|L|`$ and thus it is a finite union of fibres of $`p_1`$. Using Bertini’s theorem and the fact that the general curve of $`|L|`$ is smooth, one proves that the singularities of the general $`\mathrm{\Sigma }|(n)|`$ at points of the fixed locus of $`|(n)|`$ are finitely many rational double points of type $`A_r`$. Now, the projection $`p_1`$ restricts to a generically finite map $`p:\mathrm{\Sigma }W`$ of degree $`n`$ which, by Bertini’s theorem again, is unramified over the singular points of $`W`$. So the general $`\mathrm{\Sigma }`$ has, over each singular point $`x`$ of $`W`$, $`n`$ singularities which are analytically equivalent to the one $`W`$ has in $`x`$ (i.e. $`n`$ canonical singularities) and it is smooth at points that are smooth for $`W\times 𝐏^1`$ and are not base points of $`|(n)|`$. To describe the singularities of $`Y=\stackrel{~}{h}^{}(\mathrm{\Sigma })`$, we notice that the restriction $`Y\mathrm{\Sigma }`$ of $`\stackrel{~}{h}`$ is ramified precisely over the singularities of $`\mathrm{\Sigma }`$ that occur at singular points of $`W\times 𝐏^1`$; so $`Y`$ has $`d`$ singularities analytically isomorphic to those of $`\mathrm{\Sigma }`$ over each of those singular points of $`\mathrm{\Sigma }`$ that occur at base points of $`|(n)|`$ and it is smooth elsewhere, since it is general in $`\stackrel{~}{h}^{}|(n)|`$. In conclusion the singularities of $`Y`$ and $`\mathrm{\Sigma }`$ are canonical, and their invariants, which we now compute, are equal to those of $`X`$, $`S`$, respectively.
By the adjunction formula and condition (GP1) in definition 2.4, one has $`K_\mathrm{\Sigma }=(K_{W\times 𝐏^1}+\mathrm{\Sigma })|_\mathrm{\Sigma }=(p_1^{}K_W+(n2))|_\mathrm{\Sigma }`$ and $`K_X=(K_{V\times 𝐏^1}+X)|_X=\stackrel{~}{h}^{}(K_{W\times 𝐏^1}+\mathrm{\Sigma })|_X=\psi ^{}(K_\mathrm{\Sigma })`$, and thus $`K_S^2=K_\mathrm{\Sigma }^2=n(K_W^2L^2)+8(n1)(g1)`$, $`K_X^2=dK_S^2`$ . To compute the remaining invariants of $`S`$, $`\mathrm{\Sigma }`$ and $`X`$, one considers the long cohomology sequences associated to the restriction sequences
$$0K_{W\times 𝐏^1}K_{W\times 𝐏^1}+(n)K_\mathrm{\Sigma }0$$
and
$$0K_{V\times 𝐏^1}K_{V\times 𝐏^1}+\stackrel{~}{h}^{}(n)K_X0.$$
By Kawamata-Viehweg’s vanishing theorem, we have $`h^i(W\times 𝐏^1,K_{W\times 𝐏^1}+(n))=h^i(V\times 𝐏^1,K_{V\times 𝐏^1}+\stackrel{~}{h}^{}(n))=0`$ for $`i>0`$. Hence:
$$p_g(S)=h^0(\mathrm{\Sigma },K_\mathrm{\Sigma })=$$
$$=h^0(W\times 𝐏^1,K_{W\times 𝐏^1}+(n))+h^1(W\times 𝐏^1,K_{W\times 𝐏^1})=$$
$$=h^0(W,K_W+L)(n1)+p_g(W)=np_g(W)+(n1)g$$
where the last equality follows again from Kawamata–Viehweg’s vanishing and the last equality but one follows from $`q(W)=0`$. Therefore, by the definition of a generating pair:
$$p_g(X)=h^0(V,K_V+L)(n1)+p_g(V)=p_g(S).$$
A similar computation gives $`q(S)=q(W)=0`$, $`q(X)=q(V)=(d1)(g1)`$.
The linear system $`|K_\mathrm{\Sigma }|`$ contains the restriction of the system $`|p_1^{}K_W+(n2)|`$, whose fixed locus is the inverse image via $`p_1`$ of the fixed locus of $`|K_W+L|`$. Let $`C|L|`$ be a smooth curve; since $`W`$ is regular, the linear system $`|K_W+L|`$ restricts to the complete canonical system $`|K_C|`$. Thus $`C`$ does not contain any base point of $`|K_W+L|`$. If $`L`$ is ample, this implies that $`|K_W+L|`$ has a finite number of base points, none of which is also a base point of $`|L|`$. Thus in this case the fixed locus of $`|p_1^{}K_W+(n2)|`$ intersects the general $`\mathrm{\Sigma }`$ in a finite number of points and, a fortiori, the canonical system of $`\mathrm{\Sigma }`$ has no fixed components and the surfaces $`X`$, $`S`$ are minimal.
Notice now that $`|p_1^{}K_W+(n2)|`$ separates the fibres of $`p_2|_\mathrm{\Sigma }`$, since $`n3`$. A fibre $`F`$ of $`p_2|_\mathrm{\Sigma }`$ is identified by $`p_1`$ with a curve $`C|L|`$ and the restriction of $`|p_1^{}K_W+(n2)|`$ to $`F`$ is identified with the restriction of $`|K_W+L|`$ to $`C`$, which is the complete canonical system $`|K_C|`$, since $`W`$ is regular. Thus, if the general $`C`$ is not hyperelliptic, then the canonical map of $`S`$ is birational and $`\psi :XS`$ is a good canonical cover, while if the general $`C`$ is hyperelliptic then the canonical map of $`S`$ is of degree $`2`$ onto a rational surface and $`\psi :XS`$ is a non good canonical cover. $``$
Since we aim at a classification of generating pairs, we find useful to introduce a notion of blow-up. We will show (cf. corollary 6.3) that in most cases that almost every generating pair is obtained from a minimal one by a sequence of blow-ups.
###### Definition 2.8
Let $`(h:VW,L)`$ be a generating pair of degree $`d`$ and genus $`g`$. Let $`xW`$ be a smooth point. Then we can consider the cartesian square:
$$\begin{array}{ccccc}& V^{}& & V& \\ \hfill h^{}& & & & h\hfill \\ & W^{}& \stackrel{f}{}& W& \end{array}$$
where $`f:W^{}W`$ is the blow-up of $`W`$ at $`x`$, with exceptional divisor $`E`$ and, accordingly, $`V^{}`$ is the blow-up of $`V`$ at the $`d`$ points $`x_1,\mathrm{},x_d`$ of the fibre of $`h`$ over $`x`$. Fix $`m=0`$ or $`1`$ and assume that:
(i) $`L^2>m^2`$;
(ii) $`h^0(W^{},f^{}LmE)2`$ and the general curve $`C|h^{}LmE|`$ is smooth.
Then the pair $`(h^{}:V^{}W^{},f^{}LmE)`$ is again a generating pair. We say that it is obtained from $`(h:VW,L)`$ by a simple blow-up. The blow-up is said to be essential if $`m=1`$ and inessential if $`m=0`$.
The reason why we only consider $`m1`$ in the above definition is that generating pairs satisfy the inequality $`L^24`$ (cf. prop. 7.3 and prop. 8.4).
## 3 Examples of generating pairs
In this section we describe some examples of generating pairs.
###### Example 3.1
Beauville’s example. (see \[Cat2\], 2.9, \[MP\], example 4 in section 3). This example has already been described in section 2: $`V`$ is a principally polarized abelian surface with an irreducible polarization $`D`$, $`W`$ is the Kummer surface of $`V`$, $`h:VW`$ is the projection onto the quotient, and $`L`$ is an ample line bundle on $`W`$ such that the class of $`L^{}=h^{}L`$ is equal to $`2D`$. This generating pair is good and therefore so is any related canonical cover. More precisely, by proposition 2.7, an $`n`$-related canonical cover $`\psi :XS`$ is minimal, with geometric genus $`4n3`$. The invariants of $`S`$ and $`X`$ satisfy the relations:
$$K_S^2=3p_g(S)7;K_X^2=6p_g(X)14;q(X)=2.$$
According to \[MP\], Thm. 4.1, this is the only good generating pair such that the related canonical covers satisfy $`K_X^2=6p_g(X)14`$ and $`K_X`$ is ample.
Notice that, if, in the above situation, the polarization $`D`$ on $`V`$ is not irreducible, then the same construction produces a generating pair which is no longer good (cf. also \[MP\], example 2 in section 3). We will refer to this example as to the non good Beauville’s example.
###### Example 3.2
A good generating pair of degree $`2`$ and genus $`3`$. (cf. also \[C2\], example (c), page 70). Let $`A`$ be an abelian surface with an irreducible principal polarization $`D`$, let $`p:VA`$ be the double cover branched on a symmetric divisor $`B|2D|`$ and such that $`p_{}𝒪_V=𝒪_A𝒪_A(D)`$. Since $`K_V=p^{}(D)`$, the invariants of the smooth surface $`V`$ are: $`p_g(V)=2`$, $`q(V)=2`$, $`K_V^2=4`$. By the symmetry of $`B`$, multiplication by $`1`$ on $`A`$ can be lifted to an involution $`i`$ of $`V`$ that acts as the identity on $`h^0(V,K_V)`$. We denote by $`h:VW=V/<i>`$ the projection onto the quotient. We observe that $`p_g(W)=2`$, $`q(W)=0`$, $`K_W^2=2`$ and the only singularities of the surface $`W`$ are $`20`$ ordinary double points. In addition, $`h^0(W,2K_W)=\chi (𝒪_W)+K_W^2=4=h^0(V,2K_V)`$, so that the bicanonical map of $`V`$ factors through $`h:VW`$. An alternative description of $`W`$ is as follows. One embeds, as usual, the Kummer surface $`Kum(A)`$ of $`A`$ as a quartic surface in $`𝐏^3=𝐏(H^0(A,2D)^{})`$. The surface $`W`$ is a double cover of $`Kum(A)`$ branched over the smooth plane section $`H`$ of $`Kum(A)`$ corresponding to $`B`$ and on $`6`$ nodes (corresponding to the six points of order 2 of $`A`$ lying on $`D`$). The ramification divisor $`R`$ of $`WKum(A)`$ is a canonical curve isomorphic to $`H`$, and thus it is not hyperelliptic. This completes the proof that $`(h:VW,K_W)`$ is a good generating pair. Notice that, under suitable generality assumptions, $`K_W`$, as well as $`K_V`$, is ample. An $`n`$-related canonical cover $`\psi :XS`$ has geometric genus $`5n3`$ and is, in general, minimal. The invariants of $`S`$ and $`X`$ satisfy the relations:
$$5K_S^2=16p_g(S)32;5K_X^2=32p_g(X)64;q(X)=2.$$
(1)
###### Example 3.3
A good generating pair of degree $`2`$ and genus $`4`$ (cf. \[C2\], example 3.13). Let $`\mathrm{\Gamma }`$ be a non–hyperelliptic curve of genus $`3`$ and let $`V:=Sym^2(\mathrm{\Gamma })`$. The surface $`V`$ is smooth minimal of general type with invariants: $`K_V^2=6`$, $`p_g(V)=q(V)=3`$. If we embed $`\mathrm{\Gamma }`$ into $`𝐏^2`$ via the canonical system, then the canonical map of $`V`$ sends the unordered pair $`\{p,q\}`$ of $`V`$ to the line $`<p,q>𝐏^2`$, hence it is a degree $`6`$ morphism onto the plane. There is an involution $`i`$ on $`V`$ that maps $`\{p,q\}V`$ to $`\{r,s\}`$, where $`<p,q>\mathrm{\Gamma }=p+q+r+s`$. The fixed points of $`i`$ correspond to the $`28`$ bitangents of $`\mathrm{\Gamma }`$ and the canonical map of $`V`$ clearly factors through the quotient map $`h:VW=V/<i>`$. Hence the invariants of $`W`$ are: $`p_g(W)=p_g(V)=3`$, $`K_W^2=K_V^2/2=3`$, $`\chi (W)=(\chi (V)+7)/2=4`$, and thus $`q(W)=0`$. In addition we have $`h^0(W,2K_W)=\chi (𝒪_W)+K_W^2=7=h^0(V,2K_V)`$ and thus $`|2K_V|=h^{}|2K_W|`$. In order to complete the proof that $`(h:VW,K_W)`$ is a good generating pair we remark that the general canonical curve $`C`$ of $`W`$ is not hyperelliptic, since the restriction of $`|K_W|`$ to $`C`$ is a base–point free $`g_3^1`$. Notice that $`K_W`$ and $`K_V`$ are ample. An $`n`$-related canonical cover $`\psi :XS`$ is minimal, of geometric genus $`7n4`$. The invariants of $`S`$ and $`X`$ satisfy the relations:
$$7K_S^2=24p_g(S)72;7K_X^2=48p_g(X)144;q(X)=3.$$
(2)
An interesting question, concerning this example and the previous one, is whether these are the only generating pairs such that the related canonical covers have invariants satisfying (1) and (2).
###### Example 3.4
A non good generating pair of degree $`3`$ and genus $`2`$. Let $`p:W𝐏^2`$ be the double cover of $`𝐏^2`$ ramified on an irreducible sextic $`B`$ with 9 cusps ($`B`$ is the dual of a smooth cubic). The surface $`W`$ is a K3 surface whose singularities are 9 double points of type $`A_2`$. According to \[BdF\] (cf. also \[BL2\], \[Ba\]), there exists a smooth cover $`h:VW`$ of degree 3 ramified only at the 9 double points. The surface $`V`$ is an abelian surface. Let $`L=p^{}(𝒪_{𝐏^2}(1))`$. Since $`L`$ is ample, we have $`h^0(W,L)=\chi (𝒪_W)+\frac{1}{2}L^2=3=h^0(V,h^{}L)`$, and thus $`(h:VW,L)`$ is a non good generating pair. An $`n`$-related, minimal, canonical cover $`\psi :XS`$ has geometric genus $`4n2`$. The invariants of $`S`$ and $`X`$ satisfy the relations:
$$K_S^2=2p_g(S)4;K_X^2=6p_g(X)12;q(X)=2.$$
It is perhaps worth remarking that the surfaces $`S`$ thus obtained have invariants lying on the Noether’s line $`K_S^2=2p_g(S)4`$. It would be interesting to know whether there are other canonical covers with so low geometric genus.
###### Example 3.5
A series of non good generating pairs of degree $`2`$ with unbounded invariants. For $`i=1,2`$, let $`\varphi _i:C_i𝐏^1`$ be a double cover, where $`C_i`$ is a smooth curve of genus $`g_i>0`$, and let $`\sigma _i`$ be the involution on $`C_i`$ induced by $`\varphi _i`$. We set $`V=C_1\times C_2`$, $`W=V/<\sigma _1\times \sigma _2>`$ and we denote by $`h:VW`$ the projection onto the quotient. We remark that there exists a double cover $`f:W𝐏^1\times 𝐏^1`$ such that $`\varphi _1\times \varphi _2:V𝐏^1\times 𝐏^1`$ factors as $`\varphi _1\times \varphi _2=fh`$. We denote by $`H`$ a divisor of type $`(1,1)`$ on $`𝐏^1\times 𝐏^1`$ and we set $`L=f^{}H`$. Both systems $`|K_V|`$ and $`|K_V+h^{}L|`$ are clearly pull-back via $`\varphi _1\times \varphi _2:V𝐏^1\times 𝐏^1`$. This immediately implies that $`(h:VW,L)`$ is a non good generating pair of degree $`2`$ and genus $`g_1+g_2+1`$. One has: $`p_g(W)=g_1g_2`$, $`q(V)=g_1+g_2`$. An $`n`$-related canonical cover $`\psi :XS`$ has geometric genus $`ng_1g_2+(n1)(g_1+g_2+1)`$ and moreover $`q(X)=g_1+g_2`$.
Notice that, if $`g_1=g_2=1`$, we find again the non good Beauville’s example (cf. example 3.1).
## 4 Auxiliary results on irregular surfaces
In this section we collect a few general facts on irregular surfaces that will be used in the rest of the paper. We use notation 2.1.
###### Proposition 4.1
Let $`h:VW`$ be a finite morphism of surfaces with canonical singularities such that $`K_V=h^{}K_W`$ and $`p_g(V)=p_g(W)`$. If $`q(V)>q(W)>0`$, then the Albanese image of $`V`$ is a curve.
Proof: The critical set $`\mathrm{\Delta }`$ of $`h`$ is finite by assumption. We let $`W_0=W(\mathrm{Sing}(W)\mathrm{\Delta })`$ and $`V_0=h^1W_0`$, so that the restricted map $`h:V_0W_0`$ is a finite étale map between smooth surfaces. In particular $`h`$ is flat, and there is a canonical vector bundle isomorphism $`h_{}𝒪_{V_0}𝒪_{W_0}E`$, where $`h_{}𝒪_{V_0}`$ and $`E`$ are locally free of ranks $`d=\mathrm{deg}h`$ and $`d1`$ respectively. Since $`\mathrm{\Omega }_{V_0}^i=h^{}\mathrm{\Omega }_{W_0}^i`$, $`i=1,2`$, one has $`h_{}\mathrm{\Omega }_{V_0}^i=\mathrm{\Omega }_{W_0}^ih_{}𝒪_{V_0}=\mathrm{\Omega }_{W_0}^i(\mathrm{\Omega }_{W_0}^iE)`$. Notice that this decomposition as a direct sum is canonical. We set $`M_+^i=H^0(W_0,\mathrm{\Omega }_{W_0}^i)`$ and $`M_{}^i=H^0(W_0,\mathrm{\Omega }_{W_0}^iE)`$. We deduce that $`H^0(V_0,\mathrm{\Omega }_{V_0}^i)=H^0(W_0,h_{}\mathrm{\Omega }_{V_0}^i)=M_+^iM_{}^i`$ and we denote by $`\pi _+^i`$ the projection onto the first factor of this decomposition. To ensure that the Albanese image is a curve, we show that $`\tau _1\tau _2=0`$ for every choice of $`\tau _1,\tau _2H^0(V_0,\mathrm{\Omega }_{V_0}^1)=M_+^1M_{}^1`$. Noticing that both $`M_+^1`$ and $`M_{}^1`$ are non-zero (since $`q(V)>q(W)>0`$), we only need to show that $`h^{}\sigma \tau =0`$ for every choice of $`\sigma M_+^1`$ and $`\tau M_{}^1`$. Indeed, to show that $`^2M_+^1=0`$ we fix $`(0)\tau M_{}^1`$: if $`\sigma _1,\sigma _2M_+^1`$, the vanishing $`h^{}\sigma _i\tau =0`$ ($`i=1,2`$) means that $`h^{}\sigma _i`$ is pointwise proportional to $`\tau `$ ($`i=1,2`$), so that $`h^{}\sigma _1`$ and $`h^{}\sigma _2`$ are mutually pointwise proportional. Similarly one proves that $`^2M_{}^1=0`$.
Since $`p_g(V)=p_g(W)`$, $`\pi _+^2`$ is an isomorphism. Notice also that $`\pi _+^i(h^{}\sigma )=\sigma `$ for any $`\sigma M_+^i`$, and that $`\pi _+^2(h^{}\sigma \tau )=\sigma \pi _+^1(\tau )`$ for $`\sigma M_+^1`$ and $`\tau H^0(V_0,\mathrm{\Omega }_{V_0}^1)`$. Therefore $`h^{}\sigma \tau =0`$ for any $`\sigma M_+^1`$ and $`\tau ker\pi _+^1=M_{}^1`$, as we wanted. $``$
We recall the following results:
###### Proposition 4.2
(Serrano, \[Se\], section 1) Let $`V`$ be a smooth surface, let $`C`$ be a smooth curve, and let $`p:VC`$ be an isotrivial fibration with fibre $`D`$. Then there exist a curve $`B`$, a finite group $`G`$ acting both on $`B`$ and $`D`$, an isomorphism $`f:CB/G`$, and a birational map $`r:V(D\times B)/G`$, where $`G`$ acts diagonally on $`D\times B`$, such that the following diagram commutes:
$$\begin{array}{ccccc}& V& \stackrel{r}{}& (D\times B)/G& \\ \hfill p& & & & p^{\prime \prime }\hfill \\ & C& \stackrel{f}{}& B/G& \end{array}$$
where $`p^{\prime \prime }`$ is the map induced by the projection $`D\times BB`$. The irregularity $`q(V)`$ is equal to $`g(C)+g(D/G)`$. In particular, if $`q(V)>0`$ and $`g(C)=0`$, then the Albanese image of $`V`$ is a curve isomorphic to $`D/G`$ and the Albanese pencil is given by the composition $`p^{}r`$, where $`p^{}`$ is the map induced by the projection $`D\times BD`$.
###### Proposition 4.3
(Xiao, \[Xi\], Thm.1) Let $`p:V𝐏^1`$ be a fibration with fibres of genus $`\gamma `$. If $`p`$ is not isotrivial, then $`\gamma 2q(V)1`$.
The next proposition combines the previous results.
###### Proposition 4.4
Let $`V`$ be a smooth surface with a pencil $`|D|`$ such that the general curve $`D`$ of $`|D|`$ is smooth and irreducible of genus $`\gamma >1`$; if the Albanese image of $`V`$ is a curve, then one (and only one) of the following holds:
i) there exists a birational map $`r:VD\times 𝐏^1`$ such that $`D`$ is the strict transform via $`r`$ of a fibre of the projection $`D\times 𝐏^1𝐏^1`$. In this case $`\gamma =q(V)`$;
ii) there exist an hyperelliptic curve $`B`$, a free involution $`i`$ on $`D`$, and a birational map $`r:V(D\times B)/𝐙_2`$, where $`𝐙_2`$ acts on $`B`$ as the hyperelliptic involution, on $`D`$ via $`i`$ and diagonally on $`D\times B`$, such that $`D`$ is the strict transform via $`r`$ of a fibre of the projection $`(D\times B)/𝐙_2B/𝐙_2=𝐏^1`$. In this case $`\gamma =2q(V)1`$.
iii) $`\gamma >2q(V)1`$.
In particular, if $`p`$ is not isotrivial, then iii) holds.
Proof: Since the statement is essentially birational, up to blowing up the base locus of $`|D|`$, we may assume that $`D`$ defines a morphism $`p:V𝐏^1`$. Denote by $`\alpha :VC`$ the Albanese pencil. If $`p`$ is not isotrivial, then $`\gamma 2q(V)1`$ holds by proposition 4.3. If $`\gamma =2q(V)1`$ , then by the Hurwitz formula the restriction of $`\alpha `$ to a smooth curve $`D`$ is an étale cover of $`C`$, whose degree is $`2`$. Thus $`p`$ is isotrivial, contradicting the previous assumption.
Assume now that $`p`$ is isotrivial. By proposition 4.2, there is a commutative diagram:
$$\begin{array}{ccccc}& V& \stackrel{r}{}& (D\times B)/G& \\ \hfill p& & & & p^{\prime \prime }\hfill \\ & 𝐏^1& \stackrel{f}{}& B/G& \end{array}$$
where $`B`$ is a curve, $`G`$ is a finite group acting on $`B`$ and on $`D`$ and acting diagonally on $`D\times B`$, $`r`$ is a birational map and $`f:𝐏^1B/G`$ is an isomorphism. Again by proposition 4.2, the Albanese image of $`V`$ is isomorphic to $`D/G`$. So we have either $`G=\{1\}`$, corresponding to case i), or $`2q(V)1\gamma `$, with equality if and only if $`G=𝐙_2`$ acts freely on $`D`$. The latter case corresponds to case ii). $``$
## 5 General properties of generating pairs
In this section we give some useful information on the degree, genus and Kodaira dimension of a generating pair.
###### Notation 5.1
If $`(h:VW,L)`$ is a generating pair of degree $`d`$ and genus $`g`$, we write $`C`$ for a general curve of $`|L|`$ and $`C^{}=h^{}C`$, so that $`C`$ and $`C^{}`$ are smooth curves of genera $`g`$ and $`d(g1)+1`$ respectively, and $`h`$ restricts to an unramified cover $`\pi :C^{}C`$ of degree $`d`$.
###### Lemma 5.2
Let $`(h:VW,L)`$ be a generating pair of degree $`d`$ and genus $`g`$. If the Albanese image of $`V`$ is a curve, then $`d(g1)+1>2q(V)1`$.
Proof: According to proposition 4.4, we distinguish three cases. Setting $`D=C^{}`$, $`\gamma =d(g1)+1`$ and keeping the rest of notation of proposition 4.4, we only need to exclude the occurrence of the first two cases:
i) $`V`$ is ruled and $`C^{}`$ is a section: in this case the adjoint system $`|K_V+L^{}|`$ is empty, contradicting assumption (GP3) of definition 2.4;
ii) there are two subcases:
* Assume $`B=𝐏^1`$. Then $`V`$ is ruled over $`C^{}/𝐙_2`$ and $`C^{}`$ is a bisection of $`V`$ meeting each fibre of the map $`p:VC^{}/𝐙_2`$ in two distinct points interchanged by the free $`𝐙_2`$ action. By repeatedly blowing down $`1`$ curves $`E`$ such that $`EL^{}1`$, one obtains a map $`f:VV^{}`$ such that $`V^{}`$ is minimal and the map $`p`$ factors as $`p^{}f`$, where $`p^{}:V^{}C^{}/𝐙_2`$. The curve $`C^{\prime \prime }=f(C^{})`$ is smooth and the induced map $`f:C^{}C^{\prime \prime }`$ is an isomorphism. Moreover, the map $`p^{}:V^{}C^{}/𝐙_2`$ is a projective bundle, i.e. there exists a rank $`2`$ vector bundle $`M`$ on $`C^{}/𝐙_2`$ such that $`V^{}=\mathrm{Proj}_{C^{}/𝐙_2}(M)`$, and $`C^{\prime \prime }`$ meets each fibre of $`p^{}`$ in two distinct points interchanged by the free $`𝐙_2`$ action.
If we denote by $`H`$ the tautological section of $`V^{}`$ and by $`L^{\prime \prime }`$ the line bundle determined by $`C^{\prime \prime }`$ on $`V^{}`$, then the condition that the projection map $`C^{}C^{}/𝐙_2`$ is unramified of degree $`2`$ is equivalent to $`L^{\prime \prime }`$ being numerically equivalent to $`2H\mathrm{deg}(M)F`$, and thus we have $`L^{\prime \prime 2}=0`$. This would imply $`L^20`$, contradicting the fact that $`L^{}`$ is big.
* Assume that $`B`$ is not rational. Notice that $`(C^{}\times B)/𝐙_2`$ is the quotient of $`C^{}\times B`$ by a free $`𝐙_2`$ action. Hence it is smooth. In addition it is minimal, since it is a free quotient of the minimal surface $`C^{}\times B`$. This implies that the birational map $`r:V(C^{}\times B)/𝐙_2`$ is a morphism. Let $`C^{\prime \prime }`$ be a fibre of the morphism $`(C^{}\times B)/𝐙_2B/𝐙_2=𝐏^1`$. Since, by proposition 4.4, $`C^{}`$ is the strict transform of $`C^{\prime \prime }`$ via $`r`$ and since $`C^{\prime \prime 2}=0`$, we have again that $`L^20`$, which is impossible since $`L^{}`$ is big. $``$
###### Lemma 5.3
Let $`(h:VW,L)`$ be a generating pair of degree $`d`$ and genus $`g`$. Then $`d(g1)+12q(V)1`$, and if equality holds then the Albanese image of $`V`$ is a surface.
Proof: Consider a pencil $`𝒫|L|`$ such that the general curve is smooth and irreducible. Up to blowing up, we may assume that the pull-back of $`𝒫`$ on $`V`$ via $`h`$ is a base point free pencil. If the corresponding fibration is not isotrivial, then the claim holds by proposition 4.3. If the fibration is isotrivial, then the Albanese image of $`V`$ is a curve according to proposition 4.2, and by lemma 5.2 we have $`d(g1)+1>2q(V)1`$. $``$
###### Proposition 5.4
If $`(h:VW,L)`$ is a generating pair of degree $`d`$, then $`q(W)=0`$, and the list of possibilities is as follows:
* $`d=2`$, $`q(V)=g1`$,
* $`d=3`$, $`g3`$, $`q(V)=2(g1)`$
* $`d=4`$, $`g=2`$, $`q(V)=3`$.
If the pair is good, then case i) holds; in case ii), $`g=3`$, and in case iii) the Albanese image of $`V`$ is a surface.
Proof: By Kawamata-Viehweg’s vanishing theorem one has $`h^0(W,K_W+L)=\chi (W)+g1`$. Analogously, one has $`h^0(V,K_V+L^{})=\chi (V)+d(g1)`$, and thus $`q(V)q(W)=(d1)(g1)>0`$ by condition (GP3) of definition 2.4. Assume that $`q(W)>0`$. By proposition 4.1, the Albanese image of $`V`$ is a curve and lemma 5.2 implies that $`d(g1)+1>2q(V)1=2(d1)(g1)+2q(W)1`$, but this is impossible, since $`d,g2`$. So, $`q(W)=0`$ and, according to lemma 5.3, one has $`d(g1)+12q(V)1=2(d1)(g1)1`$. The statement includes all possible solutions. In cases ii) with $`g=3`$ and iii) we also apply lemma 5.3.
Assume now that the pair is good. By the above discussion, we have $`d3`$. By (\[B2\], Prop. 4.1 and Rem. 4.2), if $`\psi :XS`$ is a good canonical cover of degree $`3`$, then $`q(X)3`$. On the other hand, by proposition 2.7, canonical covers arising from a good generating pair of degree $`3`$ and genus $`g`$ satisfy $`q(X)=2(g1)4`$. $``$
###### Proposition 5.5
Let $`(h:VW,L)`$ be a generating pair: then $`\kappa (V)=\kappa (W)`$.
Proof: Remark first of all that $`\kappa (V)\kappa (W)`$. Hence we may assume $`\kappa (W)1`$. Consider the following commutative diagram:
$$\begin{array}{ccccc}& V^{}& \stackrel{b}{}& V& \\ \hfill h^{}& & & & h\hfill \\ & W^{}& \stackrel{f}{}& W& \end{array}$$
where $`f:W^{}W`$ is a minimal desingularization and $`h^{}:V^{}W^{}`$ is obtained by taking base change, normalizing and finally solving the singularities of the surface thus obtained. We notice the following facts:
(i) since $`V`$ and $`V^{}`$ are smooth surfaces, $`b`$ is a sequence of blow-ups and thus $`K_V^{}=b^{}K_V+E`$, where $`E`$ is an effective divisor supported on the $`b`$-exceptional locus. In addition, for every $`m1`$ we have $`|mK_V^{}|=b^{}|mK_V|+mE`$. (ii) Since $`W`$ has only canonical singularities, one has $`K_W^{}=f^{}K_W`$. Therefore we have $`b^{}K_V=b^{}(h^{}K_W)=h^{}K_W^{}`$.
Suppose that $`\kappa (W)=\mathrm{}`$, i.e. $`W`$ is rational by proposition 5.4. Hence also $`W^{}`$ is rational, and therefore there is an effective irreducible big divisor $`D`$ on $`W^{}`$ such that $`DK_W^{}<0`$. By remark (ii) above, there is an effective big divisor $`D^{}`$ on $`V^{}`$ such that $`D^{}(b^{}K_V)<0`$. This, together with remark (i), shows that $`\kappa (V^{})=\kappa (V)=\mathrm{}`$.
Assume now that $`\kappa (W)=0`$; then there exists a nef and big line bundle $`H`$ on $`W^{}`$ such that $`HK_W^{}=0`$. Thus $`(h^{}H)(b^{}K_V)=(h^{}H)(h^{}K_W^{})=0`$, and thus $`h^{}H`$ is a nef and big divisor that has zero intersection with the moving part of any pluricanonical system. Thus it follows that $`\kappa (V^{})0`$. If $`\kappa (W^{})=1`$, then there exists a fibration $`f:W^{}D`$, where $`D`$ is a smooth curve, such that the general fibre $`E`$ of $`f`$ is an elliptic curve. So $`(h^{}E)(b^{}K_V)=0`$, and thus the maps given by the pluricanonical systems are all composed with the fibration $`f^{}=fh`$. This shows that $`\kappa (V)1`$. $``$
According to the previous proposition, we may, and will, speak of the Kodaira dimension of a generating pair. Generating pairs of degree $`2`$ and Kodaira dimension $`0`$ are completely described in proposition 8.2.
## 6 Pairs of degree $`2`$ and Prym varieties
We consider here the case of generating pairs of degree $`2`$. The relevance of this case is underlined in proposition 5.4, where it is shown that all good generating pairs have degree $`2`$ and that all generating pairs of degree $`>2`$ have genus $`3`$.
If $`C|L|`$ is a general curve and $`C^{}=h^{}C`$, then the map $`h`$ induces an étale double cover $`\pi :C^{}C`$. If one denotes by $`J`$ (resp. $`J^{}`$) the Jacobian of $`C`$ (resp. $`C^{}`$), then the connected component of the kernel of the norm map $`\pi _{}:J^{}J`$ is a $`(g1)`$–dimensional containing the origin is an abelian variety, on which the principal polarization of $`J^{}`$ induces the double of a principal polarization. This principally polarized abelian variety is called the Prym variety of $`C^{}C`$ and it is denoted by $`Prym(C^{},C)`$. The connection between generating pairs and Prym varieties is explained in the following theorem.
###### Theorem 6.1
Let $`(h:VW,L)`$ be a generating pair of degree 2. Let $`C|L|`$ be a general curve. Then there is a natural isomorphism $`\phi :Prym(C^{},C)A`$, where $`A=Alb(V)`$ is the Albanese variety of $`V`$. In particular $`Prym(C^{},C)`$ does not depend on $`C|L|`$.
Proof: Under the present assumption, the singular points of $`W`$ form a set of $`t`$ ordinary double points, where $`t`$ satisfies the relation $`\chi (V,𝒪_V)=2\chi (W,𝒪_W)t/4`$. Evaluating the Euler characteristic of $`V`$ and $`W`$ as in proposition 5.4, one deduces that $`t=4(g+p_g(W))>0`$. So, one can choose a ramification point $`x_0`$ for $`h`$ in $`V`$. Since $`W`$ is regular by proposition 5.4, the Albanese map of $`W`$ with base point $`h(x_0)`$ is the zero map. The Albanese map of $`V`$ with base point $`x_0`$, denoted by $`\alpha :VA`$, is equivariant with respect to the involution induced by $`h`$ and the multiplication by $`1`$ on $`A`$. In particular, the restriction $`\alpha |_C^{}:C^{}A`$ is also equivariant.
Now we use the universal property of Prym varieties (cf. \[BL1\], page 382). Let $`\beta :C^{}Prym(C^{},C)`$ be the Abel–Prym map with respect to a point $`c^{}C^{}`$ and let $`\tau :AA`$ be the translation by $`\alpha (c^{})`$. Then there is a unique homomorphism $`\phi :Prym(C^{},C)A`$, independent of $`c^{}C^{}`$, such that $`\alpha |_C^{}=\tau \phi \beta `$.
Denote by $`J^{}`$ the Jacobian of $`C^{}`$. Let $`j:C^{}J^{}`$ be the Abel map with base point $`c^{}`$ and $`\gamma :J^{}Prym(C^{},C)`$ the map such that $`\beta =\gamma j`$. Let $`i_{}:J^{}A`$ be the homomorphism induced by the inclusion $`i:C^{}V`$ and the choice of $`c^{}C^{}`$. Notice that, up to a translation, we have $`\alpha _{|C^{}}=i_{}j`$. Then it is clear that $`i_{}`$ factors, up to a translation, as $`\phi \gamma `$. The differential of $`i_{}`$ at the origin of $`J^{}`$ is dual to the map $`H^1(V,𝒪_V)H^1(C^{},𝒪_C^{})`$, which is injective since $`H^1(V,𝒪_V(L^{}))=0`$ because $`L^{}`$ is big and nef. So $`i_{}`$ is surjective and $`\phi `$ is an isogeny since $`A`$ and $`Prym(C^{},C)`$ both have dimension $`g1`$ by proposition 5.4. To show that $`\phi `$ is an isomorphism, it is enough to prove that $`i_{}`$ has connected fibres. In turn, this follows if we show that the map $`H_1(C^{},𝐙)H_1(V,𝐙)`$ induced by the inclusion $`i:C^{}V`$ is surjective. The system $`|L^{}|`$ has no fixed part by assumption, so by theorem $`6.2`$ of \[Za\] there exists an integer $`k`$ such that $`|kL|`$ gives a morphism $`g:V𝐏^N`$; the image of $`g`$ is a surface, since $`L^{}`$ is big. So there exists an hyperplane $`H`$ in $`𝐏^N`$ such that $`g^1H=C^{}`$ as sets. By Theorem $`1.1`$, page $`150`$, of \[GM\], the map $`\pi _1(C)\pi _1(V)`$ is surjective, and thus $`H_1(C^{},𝐙)H_1(V,𝐙)`$ is surjective too. $``$
###### Corollary 6.2
Let $`(h:VW,L)`$ be a generating pair of degree $`2`$ and genus $`g`$; then the Albanese image of $`V`$ is a surface. In particular, the Kodaira dimension of the pair is non–negative.
Proof: Assume that the Albanese image of $`V`$ is a curve $`\mathrm{\Gamma }`$. Then $`\mathrm{\Gamma }`$ has genus $`g1`$. On the other hand, by theorem 6.1, the Albanese image of $`V`$ contains the Abel–Prym image of $`C^{}`$, which is isomorphic to $`C^{}`$ (cf. \[BL1\], prop. 12.5.2), since $`C^{}`$ is not hyperelliptic. This is a contradiction and thus the claim is proven. $``$
###### Corollary 6.3
Let $`(h:VW,L)`$ be a generating pair of degree $`2`$; then $`(h:VW,L)`$ is obtained from a minimal pair by a sequence of simple blow-ups of weight $`0`$ or $`1`$.
Proof: Denote by $`i:VV`$ the involution induced by $`h`$ and let $`E`$ be a $`1`$ curve of $`V`$. We claim that either $`L^{}E=0`$ or $`L^{}E=1`$. Let $`ϵ:VV_0`$ be the blow–down of $`E`$, let $`C^{}|L^{}|`$ be smooth and let $`C_0=ϵ(C^{})`$; notice that $`C_0`$ is singular if and only if $`L^{}E>1`$. Let $`\alpha :VA`$ be the Albanese map of $`V`$; $`A`$ is also the Albanese variety of $`V_0`$ and, if we denote by $`\alpha _0:V_0A`$ the Albanese map of $`V_0`$, one has $`\alpha =\alpha _0ϵ`$. Thus $`\alpha (C^{})=\alpha _0(C_0)`$; by theorem 6.1, $`\alpha (C^{})`$ is isomorphic to $`C^{}`$, since $`C^{}`$ is not hyperelliptic, and thus $`C_0`$ is smooth and $`L^{}E1`$. Let $`E^{}`$ be the image of $`E`$ via $`i`$; $`E^{}`$ is also a $`1`$–curve and thus, since $`\kappa (V)0`$ by corollary 6.2, either $`E=E^{}`$ or $`E`$ and $`E^{}`$ are disjoint. If $`E=E^{}`$, then $`E`$ contains precisely $`2`$ fixed points of $`i`$, but this contradicts the fact that $`E^2`$ is odd. So $`EE^{}`$ and $`F=h(E)=h(E^{})`$ is a $`1`$ curve contained in the smooth part of $`W`$. Let $`V^{}`$ be the surface obtained by blowing down $`E`$ and $`E^{}`$, let $`W^{}`$ be the surface obtained by blowing down $`F`$ and let $`h^{}:V^{}W^{}`$ be the double cover induced by $`h`$; if one denotes by $`M`$ the direct image of $`L`$, then it is easy to check that $`(h^{}:V^{}W^{},M)`$ is also a generating pair. By iterating this process finitely many times, one eventually obtains a generating pair with $`V`$ minimal. Thus $`K_V=h^{}K_W`$ is nef, and it follows that $`K_W`$ is also nef and $`W`$ is minimal, too. $``$
###### Corollary 6.4
Let $`(h:VW,L)`$ be a generating pair of genus $`g`$ and degree $`2`$. Then:
(i) $`p_g(V)=p_g(W)g2>0`$;
(ii) if the Kodaira dimension of the pair is $`2`$, then $`p_g(V)=p_g(W)max\{g1,2g6\}`$; if $`p_g(V)=2g6`$ then $`V`$ is birational to the product of a curve of genus $`2`$ and a curve of genus $`g3`$.
Proof: By corollary 6.2, $`g1=q(V)>1`$. Thus we have $`\chi (V)0`$, $`p_g(W)=p_g(V)q(V)1=g2>0`$. The case of Kodaira dimension $`2`$ follows from the theorem at pg. 345 of \[B4\]. $``$
###### Corollary 6.5
Let $`(h:VW,L)`$ be a generating pair of genus $`g`$ and degree $`2`$. If $`|C^{}||L^{}|`$ is a pencil containing a smooth curve, then $`|C^{}|`$ is not isotrivial.
Proof: Follows from corollary 6.2 and proposition 4.2. $``$
If we denote by $`_g`$ the moduli space of étale double covers of curves of genus $`g`$ and by $`𝒜_{g1}`$ the moduli space of principally polarized abelian varieties of dimension $`g1`$, then the Prym map $`𝒫_g:_g𝒜_{g1}`$ associates to every isomorphism class of étale double covers the corresponding Prym variety. The geometry of Prym varieties has been extensively studied by many authors. We are going to use some of these results in order to give a bound on the genus of good generating pairs.
###### Proposition 6.6
Let $`(h:VW,L)`$ be a generating pair of genus $`g`$ and degree $`2`$. Let $`C|L|`$ be general and let $`C^{}=h^{}C`$. Then the fibre of the Prym map $`𝒫_g:_g𝒜_{g1}`$ at the point of $`_g`$ corresponding to the double cover $`C^{}C`$ has positive dimension.
Proof: Follows from theorem 6.1 and corollary 6.5. $``$
It is known that the Prym map is generically finite for $`g6`$ (cf. the survey \[B5\] and the references quoted therein). However there exist positive dimensional fibres of $`𝒫_g`$ for any value of $`g`$. In order to state Naranjo’s theorem 6.7 that characterizes the positive dimensional fibres of $`𝒫_g`$ for high values of $`g`$, we recall that a curve $`C`$ is called bi-elliptic if and only if it admits a double cover $`CE`$ onto an elliptic curve $`E`$.
###### Theorem 6.7
(Naranjo, see \[Na2\], page 224 and \[Na1\], theorem (10.10)) Let $`C^{}C`$ be an unramified double cover of a genus $`g`$ curve $`C`$.
(i) If $`g13`$, then the fibre of $`𝒫_g`$ at the point of $`_g`$ corresponding to $`C^{}C`$ is positive dimensional if and only if $`C`$ is either hyperelliptic or is bi-elliptic. In addition, in the latter case, if $`CE`$ is a double cover of an elliptic curve, then the Galois group of the composition $`C^{}CE`$ is $`G=𝐙_2\times 𝐙_2`$ and each quotient of $`C^{}`$ under an element of $`G`$ has genus strictly greater than $`1`$.
(ii) If $`g10`$, the fibre of $`𝒫_g`$ at the point of $`_g`$ corresponding to $`C^{}C`$ is positive dimensional and $`C`$ is bi-elliptic, then the Galois group of the composition $`C^{}CE`$ is $`G=𝐙_2\times 𝐙_2`$, and each quotient of $`C^{}`$ under an element of $`G`$ has genus strictly greater than $`1`$.
From the point of view of generating pairs, the hyperelliptic case in Theorem 6.7 corresponds to the case of generating pairs of degree $`2`$ which are not good, and example 3.5 shows that these exist for arbitrary values of $`g`$. On the other hand, the bielliptic case can be excluded for good generating pairs with $`g`$ large, as theorem 6.9 below shows.
We recall some general and elementary properties of bi-elliptic curves and bi-double covers, i.e. finite covers with Galois group $`𝐙_2\times 𝐙_2`$ (cf. \[Na1\], page 50 and ff.; \[Pa1\]). If $`C`$ is bi-elliptic, then the double cover $`CE`$ with $`E`$ elliptic is unique up to automorphisms of $`E`$ if $`g6`$. Analogously, a bi-elliptic curve $`C`$ is not hyperelliptic if $`g4`$ and it is not trigonal if $`g6`$.
If $`C^{}C`$ is an étale double cover of a bi-elliptic curve $`CE`$, then the composition $`C^{}CE`$ is a degree $`4`$ cover of $`E`$ whose Galois group $`G`$ contains $`𝐙_2`$. Assume that $`G=𝐙_2\times 𝐙_2`$, denote by $`\sigma `$ the element of $`G`$ such that $`C^{}/<\sigma >=C`$ and by $`\sigma _i`$ ($`i=1,2`$) the remaining non trivial elements. For $`i=1,2`$, set $`p_i:C^{}C_i=C^{}/<\sigma _i>`$ the corresponding projection and notice that $`C_i`$ is a smooth curve of genus $`g_i`$, where $`g_1+g_2=g+1`$. Then there exists a cartesian diagram:
$$\begin{array}{ccccc}& C^{}& \stackrel{\pi _2}{}& C_2& \\ \hfill \pi _1& & & & \varphi _2\hfill \\ & C_1& \stackrel{\varphi _1}{}& E& \end{array}$$
(3)
where, for $`i=1,2`$, $`\varphi _i:C_iE`$ is a double cover, and the branch loci $`\mathrm{\Delta }_i`$ of $`\varphi _i`$, $`i=1,2`$, are disjoint. Moreover, $`\sigma =\sigma _1\sigma _2`$. The group $`G`$ also acts on $`Prym(C^{},C)`$. We denote by $`P_i`$ the connected component containing the origin of the fixed locus of the action of $`\sigma _i`$ on $`Prym(C^{},C)`$ ($`i=1,2`$), and we observe that $`(P_1,P_2)`$ is a pair of complementary abelian subvarieties of $`Prym(C^{},C)`$ of dimensions $`g_11`$ and $`g_21`$, respectively.
###### Lemma 6.8
Let $`(h:VW,L)`$ be a good generating pair of genus $`g10`$. Then the general curve $`C|L|`$ is not bielliptic.
Proof: By corollary 6.3 we may assume that the pair is minimal. Suppose, by contradiction, that the general curve $`C|L|`$ admits an elliptic involution $`CE`$, which, as we saw, is unique up to automorphisms of $`E`$. Moreover, by part (ii) of theorem 6.7 and by proposition 6.6, the Galois group $`G`$ of the composition $`C^{}CE`$ can be identified with $`𝐙_2\times 𝐙_2`$. Theorem 6.7 also ensures that there exists a cartesian diagram as in (3), with $`C_i`$ of genus $`g_i>1`$. We wish to extend this construction to $`V`$.
In order to do this, we prove first that we may choose the involutions $`\{\sigma _1,\sigma _2\}`$ consistently on the curves $`C^{}=h^{}C`$ as $`C`$ varies in $`|C|`$. In other words, there is a double cover $`\mathrm{\Psi }\mathrm{\Phi }`$ of the open subset $`\mathrm{\Phi }`$ of $`|C|`$ parametrizing smooth curves, such that its fibre at a general point $`C\mathrm{\Phi }`$ is the pair of involutions $`\{\sigma _1,\sigma _2\}`$ acting on $`C^{}=h^{}C`$. We want to prove that $`\mathrm{\Psi }`$ is the union of two irreducible components both mapping birationally to $`\mathrm{\Phi }`$. In order to do this, we have to prove that there are two sections of $`\mathrm{\Phi }\mathrm{\Psi }`$ mapping the general point $`C\mathrm{\Phi }`$ to $`\sigma _1`$, resp. $`\sigma _2`$, namely that we can rationally distinguish $`\sigma _1`$ from $`\sigma _2`$.
Recall that, by theorem 6.1, $`Prym(C^{},C)`$ is isomorphic, in a canonical way, to the Albanese variety $`A`$ of $`V`$. In this isomorphism, the connected component $`P_i`$ of the origin of the fixed locus of the action of $`\sigma _i`$ on $`Prym(C^{},C)`$ maps to an abelian subvariety $`B_i`$ of $`A`$ ($`i=1,2`$). The pair $`(B_1,B_2)`$ of complementary subvarieties can vary only in a discrete set, and therefore it is constant, independent of $`C`$. This proves our claim about the reducibility of $`\mathrm{\Psi }`$.
Next we claim that there are involutions $`\tau _i`$ on $`V`$ inducing $`\sigma _i`$ on the general $`C^{}`$, for $`i=1,2`$. Indeed, let $``$ be a general pencil inside $`|C|`$. If $`xV`$ is a general point, define $`\tau _i(x)`$ as $`\sigma _i(x)`$, where $`\sigma _i`$ is the involution defined on the unique curve $`C^{}`$ in $`h^{}()`$ passing through $`x`$. Since $`V`$ is minimal, $`\tau _i`$ extends to an automorphism of $`V`$. Notice that $`\tau _i`$ is independent of $``$, otherwise, as $``$ varies in a general rational $`1`$–parameter family of pencils, the point $`\tau _i(x)`$, $`xV`$ general, would describe a rational curve, hence $`\kappa (V)`$ would be negative, against proposition 6.4.
We denote by $`S_i`$ the quotient surface $`V/<\tau _i>`$, by $`h_i:VS_i`$ the projection onto the quotient and by $`C_i`$ the image in $`S_i`$ of a general $`C^{}`$. The singularities of $`S_1`$ and $`S_2`$, if any, are $`A_1`$ points and $`q(S_i)=g_i1`$. By proposition 4.4, if the curves $`C_i`$ vary in moduli, then $`g_i3`$, thus $`g=g_1+g_2+17`$, a contradiction. If the curves $`C_i`$ do not vary in moduli, then the Albanese image of $`S_i`$ is a curve by proposition 4.2 and the inequality $`g_i3`$ ($`i=1,2`$) holds by proposition 4.3, since $`q(S_i)g_i`$. $``$
Now we are ready to prove the following basic result:
###### Theorem 6.9
Let $`(h:VW,L)`$ be a good generating pair of genus $`g`$. Then $`g12`$, $`q(V)11`$.
Proof: Suppose, by contradiction, that $`g13`$. According to proposition 6.6 and to part (i) of theorem 6.7, we can assume that the general $`C|L|`$ is bi-elliptic. This, on the other hand, contradicts lemma 6.8. $``$
###### Corollary 6.10
Let $`(h:VW,L)`$ be a good generating pair of genus $`g`$. Assume that, in addition, $`V`$ is of general type. Then $`K_W^2529`$.
Proof: Follows by applying the index theorem to $`L`$ and $`K_W`$ on $`W`$. $``$
A more precise statement is the following:
###### Proposition 6.11
Let $`(h:VW,L)`$ be a good generating pair of genus $`g`$. Then $`g12`$ and:
* if the general curve $`C`$ in $`|L|`$ is bi-elliptic or trigonal, then $`g9`$ and $`q(V)8`$;
* if $`10g12`$ then either the general curve $`C`$ in $`|L|`$ is a smooth plane sextic (and $`g=10`$) or it is not bi-elliptic and has a base point free $`g_4^1`$.
Proof: The proof follows from theorem 6.9, theorem 6.7, and from the following results:
* (Green-Lazarsfeld \[GL\]) Assume $`g10`$. If the fibre of the Prym map $`𝒫_g`$ is positive dimensional at the point of $`_g`$ corresponding to a double cover $`C^{}C`$, then either $`C`$ has a $`g_4^1`$ or it is a smooth plane curve of degree six (and genus 10).
* (Naranjo \[Na2\]) Assume $`g10`$. Then the fibre of $`𝒫_g`$ over the point of $`_g`$ corresponding to a double cover of a trigonal curve $`C`$ is finite. $``$
## 7 Good generating pairs with $`h^0(W,L)3`$
This section is devoted to the proof of the following:
###### Theorem 7.1
Let $`(h:VW,L)`$ be a good generating pair such that $`L`$ is ample and $`h^0(W,L)3`$; then $`(h:VW,L)`$ is one of the following:
(i) example 3.1, and in this case $`h^0(W,L)=4`$, $`L^2=4`$, $`g=3`$;
(ii) a blow–up of weight $`1`$ of case (i) above, and in this case $`h^0(W,L)=3`$, $`L^2=3`$, $`g=3`$;
(iii) example 3.3, and in this case $`h^0(W,L)=3`$, $`L^2=3`$, $`g=4`$.
Theorem 7.1 will follow from a series of auxiliary results (proposition 7.5 and theorem 7.10), containing also some additional information on generating pairs. We also make use of the following result, which is proven in section 8 (it follows from proposition 8.2 and 8.3).
###### Proposition 7.2
If $`(h:VW)`$ is a good generating pair with $`\kappa (V)1`$ and $`h^0(W,L)3`$, then it is obtained from
example 3.1 by a sequence of blow-ups, at most one of which is essential, of weight $`1`$.
We start by using Reider’s method to give an upper bound for $`L^2`$ for most generating pairs.
###### Proposition 7.3
If $`(h:VW,L)`$ is a generating pair of degree $`2`$ and non–negative Kodaira dimension, then $`L^24`$.
Proof: Since $`L^2=2L^2`$, it suffices to show that $`L^29`$. We assume that $`L^210`$ and we observe that, by the hypothesis, the linear system $`|K_V+L^{}|`$ is not birational on $`V`$. Indeed, if $`xW`$ is a general point and $`h^1(x)=\{x_1,x_2\}`$, then $`x_1`$, $`x_2`$ are identified by $`|K_V+L^{}|`$. According to Reider’s Theorem (\[Re\], Thm. 1 and Cor. 2), there exists an effective divisor $`B=B_x`$ passing through $`x_1`$ and $`x_2`$, such that $`L^{}B=1`$ or $`2`$ and $`B^20`$. Since $`x`$ is general we must have $`B^2=0`$ and, by standard arguments, we may assume that $`B`$ moves in a base point free pencil and $`L^{}D>0`$ for each component of a general $`B`$ . Since the general curve $`C^{}`$ of $`|L^{}|`$ is irreducible and meets the general curve $`B_x`$ at the points $`x_1`$ and $`x_2`$, it follows that $`L^{}B=2`$. If the general $`B`$ is reducible, then $`B=B_1+B_2`$, where $`B_1`$, $`B_2`$ are numerically equivalent irreducible curves. Then $`L^{}B_i=1`$ and $`V`$ is covered by rational curves, contradicting the assumption $`\kappa (V)0`$. So, $`L^{}B=2`$ and $`B`$ is irreducible. Furthermore, the general fibre of $`h`$ is contained in some curve of the pencil described by $`B_x`$, as $`x`$ varies in $`W`$. This immediately implies that each curve of this pencil is invariant under the involution $`\iota `$ determined by $`h`$. On the other hand $`|L^{}|`$ cuts on a general curve $`B`$ a $`g_2^1`$, which of course induces on $`B`$ the restriction of $`\iota `$. This means that the image of $`B`$ via $`h`$ is a rational curve on $`W`$, which therefore has a pencil of rational curves. But this is a contradiction to $`\kappa (W)0`$. $``$
###### Lemma 7.4
Let $`(h:VW,L)`$ be a generating pair such that $`h^0(W,L)3`$. Then there are the following possibilities:
(i) $`h^0(W,L)=3`$ and $`2L^24`$;
(ii) $`h^0(W,L)=4`$, $`L^2=4`$ and $`|L|`$ is base point free.
Proof: If $`h^0(W,L)=r`$, then the restriction of the system $`|L|`$ to a general $`C`$ is a linear system $`|D|`$ of dimension $`r2>0`$ and degree $`L^24`$, according to proposition 7.3. We denote by $`|M|`$ the moving part of $`D`$. If $`r=3`$, then $`L^2\mathrm{deg}M2`$, since $`C`$ is not rational, and (i) is proven. If $`r>3`$, then $`4L^2\mathrm{deg}M2dim|M|=2(r2)4`$. Thus $`L^2=\mathrm{deg}M=r=4`$ and $`|L|`$ is base point free. $``$
###### Proposition 7.5
Let $`(h:VW,L)`$ be a good generating pair such that $`h^0(W,L)3`$. The possible cases are:
* $`L^2=4`$, $`h^0(W,L)=3`$ and $`|L|`$ is base point free. In particular, the general $`C|L|`$ is tetragonal.
* $`h^0(W,L)=3`$ and either $`L^2=3`$ or $`L^2=4`$ and $`|L|`$ has a simple base point. In particular, the general $`C|L|`$ is trigonal;
* $`L^2=4`$, $`h^0(W,L)=4`$, and the pair is obtained from Beauville’s example 3.1 via unessential blow-ups.
Proof: We denote by $`|D|`$ the restriction of $`|L|`$ to a general $`C`$ of $`|L|`$. Assume that we are in case (i) of lemma 7.4: then (i) and (ii) follow by remarking that the moving part of $`|D|`$ has degree $`>2`$, since $`C`$ is not hyperelliptic.
Assume that we are in case (ii) of lemma 7.4. Then Clifford’s theorem implies that $`g=3`$, $`|D|`$ is the canonical system and $`C`$ is embedded by $`|D|`$ as a smooth plane quartic. So the linear system $`|L|`$ maps the surface $`W`$ birationally onto a quartic $`Q𝐏^3`$. The Kodaira dimension of $`W`$ is non–negative by corollary 6.4 and thus it is zero. Claim (iii) now follows by proposition 7.2. $``$
###### Proposition 7.6
If $`(h:VW,L)`$ is a good generating pair of Kodaira dimension 2 such that $`L^2=4`$ and $`h^0(W,L)3`$, then $`10g12`$.
Proof: Notice first of all that the inequality $`g12`$ follows from theorem 6.9.
By corollary 6.3 it follows that the pair is obtained from a minimal pair by unessential blow–ups. Thus we may assume that the pair is minimal. Write $`h^0(W,L)=2+l`$, so that either $`l=0`$ or $`l=1`$.
Since $`W`$ is of general type, one has $`0<K_WL=2g2L^2=2g6`$, hence $`g4`$. For a general $`C|L|`$, consider the exact sequence:
$$0H^0(W,K_WL)H^0(W,K_W)H^0(C,K_CL_{|C})$$
and notice that $`h^0(C,K_CL_{|C})=g4+l`$, by the regularity of $`W`$ and by Riemann–Roch applied to $`C`$.
Assume first $`g6`$. Then by corollary 6.4, we have $`h^0(W,K_WL)g1(g4+l)=3l2`$. Notice that $`h^0(W,K_W2L)=0`$, since $`(K_W2L)L=2g14<0`$ and $`L`$ is nef. Thus $`|K_WL|`$ cuts out on $`C`$ a linear series of dimension $`1`$ and of degree $`2g102`$, contradicting the assumption that $`C`$ is not hyperelliptic.
Therefore we have $`g7`$. Corollary 6.4 and the above exact sequence yield: $`h^0(W,K_WL)2g6(g4+l)=g2l`$. Let $`h^0(W,K_W2L)=r`$. Then $`|K_WL|`$ cuts out on $`C`$ a special linear series of dimension $`g3rl`$ and degree $`2g10`$. By Clifford’s theorem and by the fact that $`C`$ is not hyperelliptic, we have $`g3rl<g5`$, namely $`r+l3`$. Thus we either have $`l=0`$, $`r3`$ or $`l=1`$, $`r2`$. If $`g8`$, then $`L(K_W3L)=2g18<0`$, and thus $`h^0(W,K_W3L)=0`$ and $`|K_W2L|`$ cuts out on $`C`$ a special linear series of dimension at least $`r11`$ and of degree $`2g142`$, contradicting again the assumption that $`C`$ is not hyperelliptic.
If $`g=9`$, then $`L(K3L)=0`$ and thus $`h^0(W,K_W3L)1`$, since $`L`$ is nef and big. Assume that $`h^0(W,K_W3L)=1`$; then we have $`K_W^236=K_W(K_W3L)0`$, since $`W`$ is of general type. On the other hand, the index theorem gives $`K_W^236`$. It follows that $`K_W^2=36`$ and $`K_{num}3L`$. Therefore $`K_W=3L`$, since $`K_W3L`$ is effective. So $`r=h^0(W,K_W2L)=h^0(W,L)=2+l`$, and for $`l=0`$ this contradicts the above inequality $`r+l3`$. If $`l=1`$, consider the exact sequence:
$$0(k1)LkLkL_{|C}0.$$
(4)
By Clifford’s theorem we have $`h^0(C,2L_{|C})4`$. So, for $`k=2`$, the sequence (4) implies $`h^0(W,2L)7`$. Using this and sequence (4) for $`k=3`$, one gets $`p_g(W)=h^0(W,3L)13`$, and thus $`\chi (V)=2+p_g(W)g6`$. On the other hand, Miyaoka–Yau’s inequality would give $`72=K_V^29\chi (V)=54`$, a contradiction.
So we are left with the case $`h^0(W,K_W3L)=0`$. If $`l=0`$, then $`r3`$ and the restriction of $`|K_W2L|`$ to $`C`$ is a $`g_4^2`$, contradicting again the fact that $`C`$ is non-hyperelliptic of genus $`9`$. Thus the case $`g=9`$ and $`l=0`$ does not occur.
If $`l=1`$, then we have $`r=2`$, since for $`r>2`$ we can argue as above and show that $`|K_W2L|`$ restricts to a $`g_4^2`$ on $`C`$. So we have $`h^0(W,K_WL)2+h^0(C,(K_WL)_{|C})2+4=6`$, where the last inequality follows again by Clifford’s theorem. In turn, $`p_g(W)6+h^0(C,K_{W}^{}{}_{|C}{}^{})=12`$. On the other hand, by proposition 6.4, one has $`p_g(W)12`$, with equality holding iff $`V=C_1\times C_2`$, with $`C_1`$ a curve of genus $`2`$ and $`C_2`$ a curve of genus $`6`$. Thus $`p_g(W)=12`$ and the restriction map $`H^0(W,K_W)H^0(C,K_{W}^{}{}_{|C}{}^{})`$ is surjective. Since the canonical map of $`V`$ factors through the map $`h:VW`$ and $`q(W)=0`$, the curve $`C_2`$ is hyperelliptic, and the canonical map of $`V`$ has degree $`4`$. The canonical image $`\mathrm{\Sigma }`$ of $`V`$ (and of $`W`$) is $`𝐏^1\times 𝐏^1`$ embedded via the system $`|𝒪_{𝐏^1\times 𝐏^1}(1,5)|`$. The curves of $`|L|`$ are mapped $`2`$–to–$`1`$ onto curves $`D`$ of $`\mathrm{\Sigma }`$ with $`D^2=2`$. So the curves $`D`$ are of type $`(1,1)`$, and thus rational. It follows that the general curve of $`|L|`$ is hyperelliptic, contradicting again the assumption that the pair be good. $``$
###### Corollary 7.7
If $`(h:VW,L)`$ is a good generating pair such that $`L^2=4`$ and $`h^0(W,L)=3`$, then $`10g12`$.
Proof: Proposition 7.2 implies that the Kodaira dimension of the pair is $`2`$. The thesis follows from proposition 7.6. $``$
The next result shows that case (i) of proposition 7.5 does not occur for $`L`$ ample.
###### Proposition 7.8
Let $`(h:VW,L)`$ be a good generating pair with $`h^0(W,L)=3`$, $`L^2=4`$ and $`L`$ ample; then $`|L|`$ has one simple base point.
Proof: First of all we remark that the assumption that $`L`$ is ample, corollary 6.3 and proposition 7.3 imply that the pair is minimal. By lemma 7.5 we only have to exclude that $`|L|`$ is base point free. So we assume that $`|L|`$ has no base points and we show that this leads to a contradiction. As usual we denote by $`C`$ a general curve of $`|L|`$ and by $`C^{}`$ the inverse image of $`C`$ via $`h`$; we denote by $`\varphi :W𝐏^2`$ the finite degree $`4`$ morphism given by $`|L|`$. Notice that $`\varphi `$ is flat, since it is a projective morphism with finite fibres from a normal surface to a smooth one. Our proof requires various steps.
Step 1: the polarized Abelian variety $`Prym(C^{},C)`$ is not a Jacobian or a product of Jacobians. By theorem (4.10) of \[B1\] and the assumption that $`C`$ is not hyperelliptic, if $`Prym(C^{},C)`$ is a Jacobian then one of the following holds: (i) $`C`$ is trigonal, (ii) $`C`$ is bielliptic, (iii) $`g6`$. Case (iii) is impossible by corollary 7.6. Since $`C`$ has a free $`g_4^1`$, case (i) implies $`g6`$, and therefore it is also excluded. Finally case (ii) is excluded by lemma 6.8.
Step 2: The curves of $`|L^{}|`$ are $`2`$-connected. Notice first of all that the curves of $`|L^{}|`$ are $`1`$-connected since $`|L^{}|`$ is ample. Assume that $`D|L^{}|`$ is not $`2`$-connected, namely that $`D=A+B`$ with $`A`$, $`B`$ effective and $`AB=1`$. Then $`1L^{}A=A^2+AB=A^2+1`$ and $`1L^{}B=B^2+1`$, hence $`A^20`$, $`B^20`$. If, say, $`A^2=0`$, then $`L^{}A=1`$ and $`A`$ is irreducible, since $`L^{}`$ is ample, and rational, since $`|L|`$ is base point free. This contradicts $`\kappa (V)=2`$, hence $`A^2`$ and $`B^2`$ are both positive. Since $`8=L^2=A^2+B^2+2`$, we have $`A^2+B^2=6`$ whereas the index theorem gives $`A^2B^2(AB)^2=1`$.
Step 3: The branch divisor $`Z`$ of $`\varphi `$ in $`𝐏^2`$ is not a union of lines. Here we need to consider the intersection number of Weil divisors on $`W`$. We recall that, since the singularities of $`W`$ are $`A_1`$ points, given Weil divisors $`A`$, $`B`$ on $`W`$, the intersection number $`AB`$ is an element of $`\frac{1}{2}𝐙`$, and it is an integer whenever $`A`$ or $`B`$ is Cartier. Assume that $`Z`$ is a union of lines and let $`R`$ be a line contained in $`Z`$. Then $`C_0=\varphi ^{}(R)|L|`$ is of the form $`C_0=mA+B`$, with $`2m4`$ and with $`A`$, $`B`$ effective, non-zero Weil divisors such that $`A`$ is irreducible and not contained in $`B`$. We set $`C_0^{}=h^{}C_0`$, $`A^{}=h^{}A`$, $`B^{}=h^{}B`$, so that $`C_0^{}=mA^{}+B^{}`$.
Notice that $`4=L^2=mLA+LB2LA`$ yields $`1LA2`$. Assume first $`LA=2`$. Then one has $`B=0`$, $`m=2`$, and thus $`C_0^{}=2A^{}`$. Recall that by proposition 6.1 the abelian variety $`Prym(C^{},C)`$ is naturally isomorphic to the Albanese variety of $`V`$ and denote by $`\alpha :VPrym(C^{},C)`$ the Albanese map. If $`\mathrm{\Xi }`$ is the principal polarization of $`Prym(C^{},C)`$, then by Welters criterion $`\alpha _{}C^{}`$ is homologically equivalent to $`\frac{2}{(g2)!}^{g2}\mathrm{\Xi }`$ for every $`C^{}|L|`$. Thus $`\alpha _{}A^{}`$ is homologically equivalent to $`\frac{1}{(g2)!}^{g2}\mathrm{\Xi }`$, and it follows that $`A^{}`$ is smooth and $`Prym(C^{},C)`$ is isomorphic to the Jacobian of $`A^{}`$. This is impossible by step 1, and therefore $`LA=1`$. The condition $`LA=1`$ implies that $`m3`$, $`B`$ is nonempty and $`A`$ is smooth and irreducible. Assume that $`m=3`$, and let $`R_1Z`$ be another line; write $`\varphi ^{}R_1=m_1A_1+B_1`$, with $`2m_13`$, $`A_1`$ irreducible and not contained in $`B_1`$. The equality $`1=A_1L=3A_1A+A_1B`$ gives $`AA_1=0`$, $`A_1B=1`$, and thus $`1=BL=mBA_1+BB_1m2`$, a contradiction. Thus, for every line $`RZ`$, we have $`\varphi ^{}R=2A+B`$, with $`A`$ irreducible and not contained in $`B`$. In particular, $`Z`$ is reduced. Notice also that $`AB>0`$, since the curves of $`|L^{}|`$ are $`2`$–connected by step $`2`$, and thus $`A`$ and $`B`$ have nonempty intersection. Let $`x_0AB`$, let $`y_0=\varphi (x_0)`$, and let $`C`$ be the pull–back of a general line through $`y_0`$; then $`C(2A+B)=L(2A+B)=4`$ and $`x_0`$ accounts at least for $`3`$ intersections. Thus $`\varphi ^1(y_0)`$ either consists of $`x_0`$ only, or contains also a point $`x_0^{}`$ that is not a branch point of $`\varphi `$; in either case $`x_0`$ is not a simple ramification point of $`\varphi `$ and therefore $`Z`$ is not smooth at $`x_0`$. Thus there is another line $`R_1Z`$ that contains $`x_0`$. Write $`\varphi ^{}R_1=C_1=2A_1+B_1`$. From $`1=AL=2AA_1+AB_1`$ we see that either $`AA_1=\frac{1}{2}`$ and $`AB_1=0`$ or $`AA_1=0`$ and $`AB_1=1`$. On the other hand, $`A_1`$ contains $`x_0`$ and thus we have $`A_1B>0`$ and $`A_1A>0`$. Thus we have a contradiction, and $`Z`$ is not a union of lines.
We can now consider an irreducible component $`Z^{}`$ of $`Z`$ that is not a line and a general tangent line $`R`$ to $`Z^{}`$. The curve $`C_0=\varphi ^{}(R)`$ is reduced, but singular at some point $`x`$. It moves in a base point free continuous system on $`W`$. Set $`h^1(x)=\{x_1,x_2\}`$ and let $`C_0^{}=h^{}C_0`$. Notice that the map $`h:C_0^{}C_0`$ is étale. Moreover $`C_0^{}`$ is singular at $`x_1`$ and $`x_2`$, and we can apply theorem (3.2) from \[ML\]. Then we have $`C_0^{}=A^{}+B^{}`$ with $`A^{}`$, $`B^{}`$ reduced and with no common component, since $`C_0^{}`$ is reduced as well as $`C_0`$, and $`A^{}B^{}=2`$. Actually $`A^{}B^{}=\{x_1,x_2\}`$, which proves that $`A^{}`$ and $`B^{}`$ are smooth at $`x_1`$ and $`x_2`$.
Step 4: One has $`A^2=B^2=2`$ hence $`2A^{}`$ and $`2B^{}`$ are numerically equivalent to $`L^{}`$. Since $`A^{}`$ and $`B^{}`$ move without fixed components on $`V`$, we have $`A^20`$ and $`B^20`$. Furthermore we have $`A^2+B^2=4`$ and $`L^{}A^{}=A^2+2`$ and $`L^{}B^{}=B^2+2`$. Suppose $`A^2=0`$, hence $`L^{}A^{}=A^{}B^{}=2`$. We claim that in this case $`A^{}`$ is irreducible: in fact, if $`A^{}=A_1+A_2`$, then $`A_1^2`$, $`A_2^20`$, since $`A_1`$ and $`A_2`$ move, and thus $`A_1^2=A_2^2=0`$ and $`A_1_{num}A_2`$, $`A_1B^{}=A_2B^{}=1`$, contradicting the fact that the curves of $`|L^{}|`$ are $`2`$–connected. Thus the general curve $`A^{}`$ is irreducible and moves in an irrational pencil $`𝒜^{}`$ on $`V`$. The involution $`\iota `$ determined by $`h:VW`$ fixes $`C_0^{}`$, hence it maps $`A^{}`$ to an irreducible curve $`A^{\prime \prime }`$. If $`A^{}=A^{\prime \prime }`$, then there exists $`AC_0`$ on $`W`$ such that $`A^{}=h^{}A`$, $`A^2=0`$, $`LA=1`$. Thus $`A`$ is smooth rational and, since $`A^{}`$ is general, $`W`$ is covered by rational curves, contradicting $`\kappa (W)0`$. So $`A^{\prime \prime }`$ is contained in $`B^{}`$. The curve $`A^{\prime \prime }`$ also moves in an irrational pencil $`𝒜^{\prime \prime }`$, and $`A^{}A^{\prime \prime }2`$, since $`A^{}`$ and $`A^{\prime \prime }`$ both contain $`x_1`$ and $`x_2`$. Write $`B=A^{\prime \prime }+D`$, $`C_0^{}=A^{}+A^{\prime \prime }+D`$; since $`C_0^{}A^{}=L^{}A^{}=2`$, we get $`A^{}A^{\prime \prime }=2`$ and $`DA^{}=0`$. Since $`D`$ also moves on $`V`$ without fixed components, it consists of curves of $`𝒜^{}`$, hence $`D^2=0`$. Since $`L^{}`$ is fixed by $`\iota `$, we have $`L^{}A^{\prime \prime }=L^{}A^{}=2`$, $`A^{\prime \prime 2}=A^2=0`$, and therefore: $`2=A^{\prime \prime }L^{}=A^{\prime \prime }(A^{}+A^{\prime \prime }+D)=2+A^{\prime \prime }D`$, $`A^{\prime \prime }D=0`$. Thus $`D`$ and $`A^{\prime \prime }`$ are numerically equivalent, but this contradicts $`A^{}A^{\prime \prime }=2`$.
Suppose that $`A^2=1`$. By proposition (0.18) of \[CCML\] we deduce that $`A^{}`$ is smooth irreducible and $`V`$ is isomorphic to the symmetric product of $`A^{}`$. The canonical maps of symmetric products are well known. Thus, the fact that the canonical map of $`V`$ is not birational, since it factors through $`h`$, tells us that either $`3p_a(A^{})=q(V)=g1`$ or $`A^{}`$ is hyperelliptic of genus $`p_a(A^{})4`$. The former case is impossible by corollary 7.7. The latter case is also impossible because $`|L^{}|`$ cuts out on $`A^{}`$ a base point free $`g_3^1`$. Hence we are left with the only possibility $`A^22`$ and, similarly, $`B^22`$, which implies the assertion.
Step 5: the divisors $`A^{}`$ and $`B^{}`$ are exchanged by $`\iota `$. The divisor $`\iota (A^{})=A^{\prime \prime }`$ is contained in $`C_0^{}`$ and is numerically equivalent to $`A^{}`$, since $`2A^{}`$ and $`2A^{\prime \prime }`$ are both numerically equivalent to $`L^{}`$. If $`A^{}=A^{\prime \prime }`$, then there exists $`A`$ on $`W`$ such that $`h^{}A=A^{}`$, $`A^2=1`$. We apply proposition (0.18) of \[CCML\] to the pull-back of $`A`$ to the minimal desingularization $`\stackrel{~}{W}`$ of $`W`$ and deduce that $`\stackrel{~}{W}`$ is birational to the symmetric product of $`A`$, contradicting $`q(W)=0`$. If $`A^{}`$ is irreducible, this is enough to prove that $`A^{\prime \prime }=B^{}`$. So assume that $`A^{}`$ is reducible and write $`A^{}=N+M`$, with $`N,M`$ effective nonzero. Then $`2=A^2=A^{}N+A^{}M`$, hence $`A^{}N=A^{}M=1`$ since $`A`$, as well as $`L^{}`$, is ample. This proves that $`N,M`$ are both irreducible. Since they move on $`V`$, we have $`N^20`$ and $`M^20`$ and the index theorem yields $`N^2=M^2=0`$, $`NM=1`$ and $`N`$ and $`M`$ both describe base point free pencils on $`V`$. Since $`A^{}A^{\prime \prime }`$, $`B^{}`$ and $`A^{\prime \prime }`$ have at least a common component. Thus we may write $`B^{}=M^{}+N^{}`$, where $`M^{}`$ is equal to, say, $`\iota (M)`$. We have $`M^{}B^{}=MA=1`$, $`M_{}^{}{}_{}{}^{2}=M^2=0`$, hence $`B^{}N^{}=1`$ and $`N^{}`$ is irreducible by the ampleness of $`B^{}`$. If $`\iota (N)=N^{}`$, then the claim is proven. So assume $`\iota (N)N^{}`$. Then we have $`\iota (N)=N`$ and there exists $`N_0W`$ such that $`h^{}N_0=N`$. It follows that $`LN_0=1`$, and thus $`N_0`$ is a rational curve. This is simpossible, since otherwise $`W`$ would be covered by rational curves. Thus $`\iota (N)=N^{}`$ and $`\iota `$ exchanges $`A^{}`$ and $`B^{}`$.
Step 6: conclusion of the proof. We use the notation introduced in step $`1`$. By step 5, if the base point of the Albanese map $`\alpha :VPrym(C^{},C)`$ is invariant for $`i`$, then $`\alpha _{}B^{}=(1)_{}\alpha _{}A^{}`$, since $`\alpha `$ is equivariant with respect to $`\iota `$ on $`V`$ and multiplication by $`1`$ on $`Prym(C^{},C)`$. Thus $`\alpha _{}(C_0^{})=\alpha _{}(A^{})+\alpha _{}(B^{})`$ and $`2\alpha _{}A^{}`$ represent the same cohomology class. By Welters criterion, this implies that $`\alpha _{}A^{}`$ is equivalent in cohomology to $`\frac{1}{(g2)!}^{g2}\mathrm{\Xi }`$. By the criterion of Matsusaka–Ran, $`Prym(C^{}C)`$ is isomorphic as a principally polarized abelian variety either to a Jacobian or to a product of Jacobians. This contradicts step $`1`$, and the proof is complete. $``$
###### Remark 7.9
The same ideas we exploited in the proof of the previous proposition would also yield the following result: in case (ii) of proposition 7.5, the generating pair is either the pair of the example 3.3 or it is obtained from Beauville’s example 3.1, with a blow-up procedure. We will next prove the same theorem with a different technique, which also seems illuminating to us. Hence we give here only an idea of its proof with the present methods.
Assume for simplicity $`L^2=3`$. Then one considers the finite, degree $`3`$ map $`\varphi :W𝐏^2`$ determined by $`|L|`$. First one shows that no line is in the branch divisor of $`\varphi `$. Then one proves the existence of a $`1`$-dimensional family of reduced curves $`C_0^{}|L|`$ which split as $`C_0=A+B`$, with $`AB=2`$. This implies that $`A^2=B^2=1`$. At this point one uses proposition (0.18) from \[CCML\] and proves that $`V`$ is birational to the symmetric product of $`A=B`$. The fact that the canonical map of $`V`$ is not birational tells us that either $`g3`$, which leads to the two cases which actually occur, or $`A`$ is hyperelliptic of genus $`g4`$. But this not possible because, via $`h`$, $`A`$ is birational to the image $`C_0`$ of $`C_0^{}`$, and $`C_0`$ has a base point free $`g_3^1`$.
The rest of this section is devoted to the analysis of case (ii) of proposition 7.5 under the hypothesis that $`L`$ be ample. We prove the following result:
###### Theorem 7.10
Let $`(h:VW,L)`$ be a good generating pair of genus $`g`$ such that $`L`$ is ample and $`h^0(W,L)=3`$. Then $`L^2=3`$ and:
(i) either there exists a smooth plane quartic $`\mathrm{\Gamma }`$ such that $`(h:VW,L)`$ is constructed from $`\mathrm{\Gamma }`$ as explained in example 3.3;
(ii) or $`(h:VW,L)`$ is obtained from Beauville’s example 3.1 via a simple blow-up of weight $`1`$.
By propositions 7.5 and 7.8, a pair satisfying the assumption of theorem 7.10 either has $`L^2=3`$ and $`|L|`$ is base point free or has $`L^2=4`$ and $`|L|`$ has a simple base point. So, up to a simple blow up of the pair, we may assume that $`L^2=3`$ and $`|L|`$ is base point free. Thus for the rest of the section we make the following assumption:
###### Assumption 7.11
$`(h:VW,L)`$ is a good generating pair of genus $`g`$ such that $`L`$ is ample, $`h^0(W,L)=3`$, $`L^2=3`$ and $`|L|`$ is base point free.
If assumption 7.11 holds, then $`|L|`$ defines a finite morphism $`f:W𝐏^2`$ of degree $`3`$. The restriction of $`f`$ to the general curve $`C|L|`$ exhibits $`C`$ as a triple cover of $`𝐏^1`$ showing that $`C`$ is trigonal. Given a curve $`C`$ of genus $`g`$, a degree $`3`$ map $`f:C𝐏^1`$, and an unramified double cover $`\pi :C^{}C`$, the trigonal construction (\[Rec\], cf. \[B3\]) yields a degree $`4`$ map $`\varphi :D𝐏^1`$, where $`D`$ is a smooth curve of genus $`g1`$ and $`\varphi `$ has no double fibre, such that the Jacobian of $`D`$ and $`Prym(C^{},C)`$ are isomorphic as principally polarized abelian varieties. We briefly recall the trigonal construction. One considers the induced morphism $`\pi ^{(3)}:C^{(3)}C^{(3)}`$ between the symmetric products of $`C^{}`$ and $`C`$. The curve $`\stackrel{~}{D}=\pi ^{(3)^1}(g_3^1)`$ has a natural morphism $`\stackrel{~}{D}𝐏^1`$; it turns out that $`\stackrel{~}{D}𝐏^1`$ splits as the disjoint union of two isomorphic smooth connected degree $`4`$ covers $`\varphi _i:D_i𝐏^1`$, $`i=1,2`$, and one can set $`D=D_1`$, $`\varphi =\varphi _1`$.
The trigonal construction is a one–to–one correspondence, whose inverse is the Recillas’ construction (\[Rec\], cf. \[BL1\] page 391). Given a smooth genus $`g1`$ curve $`D`$ with a degree $`4`$ morphism $`\phi :D𝐏^1`$ without double fibres, one defines a curve $`C^{}D^{(2)}`$ by setting:
$$C^{}=\{p_1+p_2D^{(2)}|p_1+p_2+p_3+p_4\text{is a fibre of }\phi \text{ for some}p_1,p_3D\}.$$
The curve $`C^{}`$ is smooth and connected of genus $`2g1`$, and has a natural free involution $`\sigma `$, which maps an element $`p_1+p_2`$ (in a fibre of $`\phi `$) to the complementary element $`p_3+p_4`$. If $`\pi :C^{}C=C^{}/<\sigma >`$ denotes the natural projection, it is easy to check that $`C`$ is trigonal.
Recillas’ correspondence has been generalized in \[Ca2\], where the author introduces the discriminant of a degree $`4`$ Gorenstein cover $`\phi :ZY`$, which is a degree $`3`$ morphism $`f:\mathrm{\Delta }(Z)Y`$. We recall that a cover $`\phi :ZY`$ is said to be a Gorenstein cover if the scheme theoretic fibre $`\phi ^1(y)`$ is Gorenstein over $`k(y)`$ for every $`yY`$ (cf. \[Ca1\]). If $`Y=𝐏^2`$, then the discriminant construction gives a one–to–one correspondence between the following objects:
(A) normal Gorenstein covers $`f:W𝐏^2`$ of degree $`3`$ such that the singularities of $`W`$ are at most RDP’s and such that there exists a double cover $`h:VW`$ branched exactly over the singularities of $`W`$;
(B) degree $`4`$ Gorenstein covers $`\phi :Z𝐏^2`$ with $`Z`$ smooth such that:
* for every $`y𝐏^2`$ the Zariski tangent space to the fibre $`\phi ^1(y)`$ has dimension $`1`$ at each point.
* the set $`R_0𝐏^2`$ of points $`y`$ such that the fibre $`\phi ^1(y)`$ is isomorphic either to $`\text{spec}𝐂[t]/(t^4)`$ or to $`\text{spec}𝐂[t,s]/(t^2+1,s^2)`$ is finite.
The properties of this correspondence ensure that the branch loci of the associated covers $`\phi `$ and $`f=\mathrm{\Delta }(\phi )`$ coincide as divisor of $`𝐏^2`$. Moreover, the singularities of $`W`$ occur precisely over the points $`yR_0`$. Notice that, in the case we are interested in, the singular locus of $`W`$ is not empty (see the proof of theorem 6.1), and therefore $`R_0`$ is not empty. Finally, fibrewise, $`Z_y`$ is the base locus of a pencil of conics whose discriminant is $`W_y`$.
Assumption 7.11 allows us to apply the trigonal contruction to the present case. Thus, given a good generating pair $`(h:VW,L)`$ as in 7.11, there exists a unique degree $`4`$ Gorenstein cover $`\phi :Z𝐏^2`$ as in (B) such that the morphism $`f:W𝐏^2`$ associated to the system $`|L|`$ is obtained from $`\phi `$ via the trigonal construction. We denote by $`|M|`$ the pull-back to $`Z`$ of the linear system of lines in $`𝐏^2`$.
###### Lemma 7.12
The smooth elements of $`|M|`$ are isomorphic curves of genus $`g12`$.
Proof: Let $`H`$ be a general line in $`𝐏^2`$, let $`D=\phi ^1H`$, let $`C=f^1H`$ and let $`C^{}C`$ be the unramified cover determined by $`h`$. By theorem 6.1, the Prym variety $`P=Prym(C^{},C)`$ is independent of $`H`$. On the other hand, $`C^{}C`$ is obtained from $`D`$ via the trigonal construction, and thus $`P`$ and the Jacobian of $`D`$ are isomorphic as p.p.a.v.’s. In particular, since the genus of $`|L|`$ is at least $`3`$, the genus of $`|M|`$ is at least $`2`$. By the global Torelli theorem for curves, the isomorphism class of $`D`$ is also independent of $`H`$. This implies that the natural map from the open set of smooth curves of $`|M|`$ to the moduli space of curves is constant. $``$
###### Lemma 7.13
Let $`y𝐏^2`$ and let $`|M_y||M|`$ be the pull-back on $`Z`$ of the pencil of lines through $`y`$. Then the general curve of $`|M_y|`$ is smooth.
Proof: The base scheme of $`|M_y|`$ is $`\phi ^1(y)`$. The statement follows by Bertini’s theorem since $`\phi :Z𝐏^2`$ satisfies condition (i) of (B). $``$
###### Lemma 7.14
$`Z`$ is a minimal geometrically ruled surface, and the smooth elements of $`|M|`$ are sections of the ruling.
Proof: Denote by $`RZ`$ and by $`B𝐏^2`$ the ramification divisor and the branch divisor of $`\phi `$. By condition (B), the ramification order of $`\phi `$ along each component of $`R`$ is $`3`$, each component of $`R`$ is mapped birationally onto its image and different components of $`R`$ are mapped to different components of $`B`$. Let $`(M_t)_{t𝐏^1}`$ be a general pencil contained in $`|M|`$ and assume that $`M_0`$ is singular. By applying stable reduction, one can replace $`M_0`$ by a stable curve $`M_0^{}`$. Lemma 7.12 implies that $`M_0^{}`$ is isomorphic to $`M_t`$, for $`t`$ general.
Assume that there exists a component $`\mathrm{\Theta }`$ of $`R`$ such that $`\mathrm{\Delta }=\phi (\mathrm{\Theta })`$ is not a line. If $`\phi `$ is ramified of order $`3`$ along $`\mathrm{\Theta }`$, then the inverse image $`M_0`$ of a generic line tangent to $`\mathrm{\Delta }`$ has an ordinary cusp over the tangency point and it is smooth elsewhere. It follows that $`M_0`$ is irreducible. Since $`p_a(M)>1`$, $`M_0`$ is not rational; the special fibre of the stable reduction of a general pencil containing $`M_0`$ is the union of the normalization $`M_0^{\prime \prime }`$ of $`M_0`$ and of a smooth elliptic curve meeting $`M_0^{}`$ at one point, but this is impossible by the remark above. If $`\phi `$ is simply ramified along $`\mathrm{\Theta }`$, take $`(M_t)`$ to be the pull-back of a pencil of lines such that $`M_0`$ is the pull-back of a line simply tangent to $`\mathrm{\Delta }`$ at a point $`y_0`$ and meeting $`B`$ transversely elsewhere. Then $`M_0`$ has an ordinary node at a point $`x_0`$ such that $`\phi (x_0)=y_0`$ and no other singularities. By the remarks above, the curve $`M_0`$ is not semistable; therefore we have $`M_0=M_0^{}+F`$, where $`F`$ is a smooth rational curve, $`M_0^{}`$ is isomorphic to $`M_t`$ for $`t`$ general, and $`M_0^{}F=1`$. We have: $`4=M^2=M_t(M_0^{}+F)`$, $`M_tM_0^{}3`$ (since $`M_0^{}`$ is not hyperelliptic), and thus $`M_tM_0^{}=3`$ and $`F^2=0`$. Noticing that $`y_0`$ is a general point of $`\mathrm{\Delta }`$, it follows that $`Z`$ is ruled. Since the system $`|M|`$ is ample, $`MF=1`$ and, by lemma 7.12, the curves of $`|M|`$ are not rational, $`Z`$ is geometrically ruled and minimal.
So we have proven that either $`Z`$ is as claimed or all the components of $`B`$ are lines. Let $`\mathrm{\Delta }B`$ be a line; by condition (B), it is not possible that $`f^{}\mathrm{\Delta }=2A`$. Thus $`\phi ^{}\mathrm{\Delta }=mA+B`$, with $`m3`$, $`A`$ irreducible, $`B`$ nonempty and $`A`$ not contained in $`B`$. Then one can argue as in step $`3`$ of the proof of proposition 7.6 and prove that $`B`$ is not a union of lines. $``$
###### Proposition 7.15
Let $`B`$ be the base curve of the ruled surface $`Z`$ of lemma 7.14 and let $`p:ZB`$ be the projection. Then there exists a birational morphism $`s:B\mathrm{\Gamma }𝐏^2`$ such that:
i) $`\mathrm{\Gamma }`$ is either a smooth quartic or a quartic with a double point;
ii) $`Z=𝐏(s^{}T_{𝐏^2}(1))`$, and $`p_{}𝒪_Z(M)=s^{}T_{𝐏^2}(1)`$.
Proof: According to lemma 7.14 there exists a rank $`2`$ bundle $`E`$ on $`B`$ such that $`Z=𝐏(E)`$ and $`p_{}𝒪_Z(M)=E`$ (in particular, $`\mathrm{deg}(detE)=4`$). Let $`D`$ be a smooth curve in $`|M|`$, which we may identify with $`B`$ via the map $`p|_D`$. Then $`M|_D`$ is identified with $`detE`$. By condition (B), (ii), if $`D`$ is general, then $`\phi |_D`$ has no multiple fibre, while if $`\phi (D)`$ contains a point of $`R_0`$ (which, as we know, is not empty) then $`\phi |_D`$ has at least one multiple fibre. So the restriction of $`|M|`$ to $`D`$ is not a complete system, i.e. $`h^0(B,detE)=3`$. Let $`s:B𝐏^2`$ be the morphism given by the linear system $`|detE|`$ and let $`\mathrm{\Gamma }=s(B)`$. If $`\mathrm{\Gamma }`$ were a conic, then the map $`\phi |_D`$ would have two multiple fibres for every smooth $`D`$ of $`|M|`$, contradicting condition (B). So $`\mathrm{\Gamma }`$ is a quartic and $`s`$ is birational. Since $`B`$ has genus $`g12`$, it follows that $`\mathrm{\Gamma }`$ is either smooth or it has one double point and $`detE=s^{}𝒪_\mathrm{\Gamma }(K_\mathrm{\Gamma })`$.
Let $`UH^0(Z,M)`$ be the subspace such that $`𝐏(U)=\phi ^{}|𝒪_{𝐏^2}(1)|`$. If we identify $`U`$ with a subspace of $`H^0(B,E)`$, then the natural sheaf map $`U𝒪_ZE`$ is surjective ($`|M|`$ is base-point free). Moreover, the map $`^2UH^0(B,s^{}K_\mathrm{\Gamma })`$ is an isomorphism (this follows from the discussion above, since we have shown that $`|M|`$ does not restrict to the same $`g_4^1`$ on all the curves of $`|M|`$).
If we choose a basis for $`U`$, then we have a short exact sequence:
$$0𝒪_B(s^{}K_\mathrm{\Gamma })𝒪_B^3E0.$$
(5)
Let the inclusion $`𝒪_B(s^{}K_\mathrm{\Gamma })𝒪_B^3`$ be given by $`(s_0,s_1,s_2)`$, where $`s_iH^0(\mathrm{\Gamma },K_\mathrm{\Gamma })`$, $`i=0,1,2`$, and let $`S`$ be the subspace of $`H^0(\mathrm{\Gamma },K_\mathrm{\Gamma })`$ spanned by $`s_0,s_1,s_2`$. Notice that $`dimS2`$, since $`E`$ is torsion free. If $`dimS=2`$, then it is clear that $`E=𝒪_B𝒪_B(s^{}K_\mathrm{\Gamma })`$ and condition ii) above is not satisfied. Thus $`s_0,s_1,s_2`$ are independent and sequence (5) is the pull-back via the map $`s`$ of the twisted Euler sequence:
$$0𝒪_{𝐏^2}(1)𝒪_{𝐏^2}^3T_{𝐏^2}(1)0.$$
Now we are ready to finish the proof of theorem 7.10:
###### Proposition 7.16
Notation as in proposition 7.15. The surface $`Z`$ is the normalization of the incidence surface $`Y=\{(p,l)\mathrm{\Gamma }\times (𝐏^2)^{}|pl\}`$, and the maps $`p:ZB`$ and $`\phi :Z𝐏^2`$ are induced by the projections of $`Y`$ onto $`\mathrm{\Gamma }`$ and $`(𝐏^2)^{}`$ respectively.
Let $`f:W𝐏^2`$ be the triple cover obtained from $`\phi :Z𝐏^2`$ via the discriminant contruction, $`h:VW`$ the corresponding double cover and $`L=f^{}𝒪_{𝐏^2}(1)`$. Then:
* if $`\mathrm{\Gamma }`$ is smooth, then $`V=Sym^2(\mathrm{\Gamma })`$, and $`(h:VW,L)`$ is as in example 3.3;
* if $`\mathrm{\Gamma }`$ has a double point, write $`p+q`$ for the only effective divisor linearly equivalent to $`s^{}K_\mathrm{\Gamma }K_B^{}`$. Then:
+ $`V`$ is the blow-up of $`Sym^2(B)`$ at $`p+q`$, namely it is the blow-up of the Jacobian $`J=J(B)=Pic^2(B)`$ of $`B`$ at the points corresponding to $`K_B`$ and $`p+q`$;
+ $`W`$ is obtained as the quotient of $`V`$ by the involution which is induced on $`V`$ by the birational involution on $`Sym^2(B)`$ which associates to the general divisor $`x+y`$ the divisor $`|s^{}K_\mathrm{\Gamma }xy|`$. Notice that $`W`$ is the blow-up of the Kummer surface $`Kum(J)`$ at a smooth point;
+ the generating pair $`(h:VW,L)`$ is obtained from example 3.1 by a simple blow-up of weight $`1`$.
Proof: We keep the notation of the proof of proposition 7.15. Assume first $`\mathrm{\Gamma }`$ is smooth. Then the first assertion immediately follows by the well known fact that $`𝐏(T_{𝐏^2}(1))`$ is the incidence correspondence inside $`𝐏^2\times (𝐏^2)^{}`$. Having in mind Recillas’ construction described at the beginning of this section, also part i) immediately follows. The case $`\mathrm{\Gamma }`$ singular is completely analogous and can be dealt with in the same way. We leave the details to the reader. $``$
## 8 The other cases
In this section we collect some information on pairs that are not good or not of Kodaira dimension $`2`$. We start by classifying non good degree $`2`$ pairs with $`L^2=4`$. (We recall that by propositions 5.5 and 7.3 such a pair always has $`L^24`$.)
###### Proposition 8.1
Let $`(h:VW,L)`$ be a non good generating pair of degree $`2`$ with $`L^2=4`$; then there exist smooth curves $`C_i`$, $`i=1,2`$, of genus $`g_i>0`$ and double covers $`\varphi _i:C_i𝐏^1`$ such that $`(h:VW,L)`$ is obtained by a sequence of unessential blow–ups from a generating pair constructed from $`\varphi _i:C_i𝐏^1`$ as in example 3.5.
Proof: By proposition 6.3 we can assume that the pair is minimal.
Let $`C|L|`$ be general and let $`C^{}=h^{}C`$. By \[Mu\], p. 346, we see that the Galois group $`G`$ of the composition of $`C^{}C`$ with the hyperelliptic involution on $`C`$ can be identified with $`𝐙_2\times 𝐙_2`$. As in the proof of lemma 6.8, denote by $`\sigma `$ the element of $`G`$ such that $`C^{}/<\sigma >=C`$ and by $`\sigma _i`$ ($`i=1,2`$) the remaining non trivial elements. For $`i=1,2`$, set $`p_i:C^{}C_i=C^{}/<\sigma _i>`$ the corresponding projection and notice that $`C_i`$ is a smooth curve of genus $`g_i`$, where $`g_1+g_2=g1`$ and there exists a cartesian diagram:
$$\begin{array}{ccccc}& C^{}& \stackrel{\pi _2}{}& C_2& \\ \hfill \pi _1& & & & \varphi _2\hfill \\ & C_1& \stackrel{\varphi _1}{}& 𝐏^1& \end{array}$$
(6)
where, for $`i=1,2`$, $`\varphi _i:C_i𝐏^1`$ is a double cover. In the present case there is an isomorphism $`Prym(C^{},C)=J(C_1)\times J(C_2)`$ as principally polarized abelian varieties, and $`A=Alb(V)`$ is also isomorphic to $`Prym(C^{},C)`$ by theorem 6.1. We can assume that $`g_1g_2`$ and the condition that $`C^{}`$ is not hyperelliptic ensures that $`g_1>0`$. Notice the existence of commutative diagrams:
$$\begin{array}{ccccc}& C^{}& \stackrel{}{}& Prym(C^{},C)& \\ \hfill \pi _i& & & & p_i\hfill \\ & C_i& \stackrel{}{}& J(C_i)& \end{array}$$
(7)
where $`C^{}Prym(C^{},C)`$ is the Abel-Prym map, $`C_iJ(C_i)`$ is the Abel-Jacobi map and $`p_i:Prym(C^{},C)=J(C_1)\times J(C_2)J(C_i)`$ is the $`i`$-th projection, $`i=1,2`$.
As in the proof of lemma 6.8, one shows that there exist involutions $`\tau _1`$, $`\tau _2`$ on $`V`$ that act on $`C^{}`$ as $`\sigma _1`$, respectively, $`\sigma _2`$. Clearly, the involution $`\iota `$ associated to $`h`$ is equal to $`\tau _1\tau _2`$. We denote by $`S_i`$ the quotient surface $`V/<\tau _i>`$, by $`h_i:VS_i`$ the projection onto the quotient and by $`C_i`$ the image in $`S_i`$ of a general $`C^{}`$. The singularities of $`S_1`$ and $`S_2`$, if any, are $`A_1`$ points and $`q(S_i)=g_i`$. More precisely, we claim that $`J(C_i)`$ is the Albanese variety of $`S_i`$. Indeed, the map $`VJ(C_i)`$ obtained by composing the Albanese map of $`V`$ with the projection onto $`C_i`$ is equivariant with respect to $`\tau _i`$, provided that the base point of the Albanese map is invariant for $`\tau _i`$. Thus we have an induced map $`S_iJ(C_i)`$ and thus the Albanese variety $`A_i`$ of $`S_i`$ is isogenous to $`J(C_i)`$. To show that this isogeny is actually an isomorphism, it is enough to remark that the map $`H_1(V,𝐙)H_1(J(C_i),𝐙)`$ is surjective, since it is the composition of $`H_1(V,𝐙)H_1(A,𝐙)`$, that is an isomorphism up to torsion, and of $`H_1(A,𝐙)H_1(J(C_i),𝐙)`$ which is surjective. On the other hand, $`H_1(V,𝐙)H_1(J(C_i),𝐙)`$ is also the composition of $`H_1(V,𝐙)H_1(S_i,𝐙)`$ and $`H_1(S_i,𝐙)H_1(J(C_i),𝐙)`$, hence the latter map is surjective and $`A_i`$ is isomorphic to $`J(C_i)`$.
We claim that $`S_i`$ is birational to $`𝐏^1\times C_i`$. Indeed, by proposition 4.3 the curve $`C_i`$ does not vary in moduli and the Albanese image of $`S_i`$ is a curve. By proposition 4.4, this concludes the proof of the claim. In particular the Albanese image of $`S_i`$ is the curve $`C_i`$.
Composing $`h_i`$ with the Albanese map $`S_iC_i`$, we get morphisms $`f_i:VC_i`$, $`i=1,2`$. Denote by $`F_i`$ a fibre of $`f_i`$. The Index theorem applied to $`F_1+F_2`$ and $`L^{}`$ gives $`2(F_1F_2)L^2[L^{}(F_1+F_2)]^2=16`$, namely $`L^28`$, $`L^24`$. If $`L^2=4`$, then $`F_1F_2=1`$ and $`L^{}`$ is numerically equivalent to $`2F_1+2F_2`$. Thus $`f=f_1\times f_2:VC_1\times C_2`$ is birational, and therefore it is an isomorphism since $`V`$ is minimal. One has: $`\tau _1=\sigma _1\times Id`$, $`\tau _2=Id\times \sigma _2`$, $`\iota =\sigma _1\times \sigma _2`$ and the curves of $`|L^{}|`$ are invariant for $`\tau _1`$, $`\tau _2`$ and it is easy to see that $`(h:VW,L)`$ is precisely as in example 3.5. $``$
Next we classify pairs of degree $`2`$ and Kodaira dimension $`0`$.
###### Proposition 8.2
Let $`(h:VW,L)`$ be a generating pair of degree $`2`$ and genus $`g`$; if the Kodaira dimension of the pair is $`0`$, then it can be obtained from example 3.1 (Beauville’s example) by a sequence of simple blow-ups, only three of which at most essential, of weight $`1`$.
Proof: Assume that the pair has Kodaira dimension $`0`$ and is minimal. By proposition 5.4 and corollary 6.4, we see that $`g=3`$ and the irregularity of $`V`$ is $`2`$, hence $`V`$ is an abelian surface. Since $`q(W)=0`$ by proposition 5.4, $`W`$ is the Kummer surface of $`V`$. By theorem 6.1, if $`C|L|`$ is general then $`V`$ is isomorphic to $`Prym(C^{},C)`$ and thus, in particular, it is principally polarized. In addition, by Welters criterion, $`C^{}`$ is a divisor of type $`(2,2)`$ and thus we have precisely example 3.1. By corollary 6.3, this implies that if the pair is not minimal, then it is obtained from example 3.1 by a sequence of blow–ups of weight $`0`$ or $`1`$. Since $`L`$ is big by assumption and example 3.1 has $`L^2=4`$, there are at most $`3`$ blow–ups of weight $`1`$ in the sequence. $``$
The next result is an almost complete classification of pairs of degree $`2`$ and Kodaira dimension $`1`$.
###### Proposition 8.3
Let $`(h:VW,L)`$ be a generating pair of degree $`2`$ and genus $`g`$ with Kodaira dimension $`1`$. Then there exist an elliptic curve $`E`$ and an hyperelliptic curve $`B`$ of genus $`g22`$ such that $`(h:VW,L)`$ is obtained by a sequence of simple blow–ups of degree $`0`$ or $`1`$ from one of the following:
(a) the pair constructed from $`E`$ and $`B`$ as in example 3.5. In this case the pair is not good;
(b) a pair $`(h_0:V_0W_0,L_0)`$ such that $`g=4`$ (and thus $`B`$ has genus $`2`$), $`V_0=B\times E`$ and $`h_0:V_0W_0`$ is the quotient map for the $`𝐙_2`$–action given by $`(b,e)(j(b),\sigma (e))`$, where $`j`$ is the hyperelliptic involution of $`B`$ and $`\sigma `$ is an involution of $`E`$ with rational quotient. In this case $`L^2=2`$, and, if the pair is good, then $`h^0(W,L)=2`$.
Proof: By corollary 6.3 we may assume that $`(h:VW,L)`$ is minimal.
Let $`VB`$ be the elliptic fibration. By a result of Beauville (see \[B4\], pg. 345) and by corollary 6.2, $`V`$ is a product $`B\times E`$, where $`E`$ is the general fibre of $`VB`$. The involution $`\iota `$ determined by $`h`$ on $`V`$ preserves the fibration $`VE`$, and, since the quotient of $`V`$ by $`\iota `$ is regular, it acts on $`B`$ as an involution $`j`$ with rational quotient. Thus $`\iota `$ can be written as $`(b,e)(j(b),\sigma _b(e))`$, where $`\sigma _b:EE`$ is an automorphism of $`E`$. If $`\sigma _b`$ is a translation for every $`bB`$, then the pull-back on $`V`$ of the nonzero $`1`$-form of $`E`$ is invariant for $`\iota `$, but this contradicts the regularity of $`W`$. So $`\sigma _b`$ is not a translation. Since $`\iota ^2=1`$, one has $`\sigma _b\sigma _{j(b)}=1`$, and, if $`b_0`$ is a fixed point of $`j`$, then $`\sigma _{b_0}^2=1`$. So $`\sigma _{b_0}`$ acts on $`H^0(E,\omega _E)`$ as multiplication by $`1`$. Since the possible actions of an automorphism of $`E`$ on $`H^0(E,\omega _E)`$ are a finite number, it follows that $`\sigma _b`$ acts on $`H^0(E,\omega _E)`$ as multiplication by $`1`$ for every $`bB`$. Thus we have $`\sigma _b^2=1`$, namely $`\sigma _b=\sigma _{j(b)}`$. So $`b\sigma _b`$ descends to a map $`B/<j>=𝐏^1Aut(E)`$ and it is therefore constant. Notice that $`B`$ has genus $`g22`$, since $`\kappa (V)=1`$.
Denote by $`F`$ the general fibre of the pencil $`p_1:W𝐏^1=B/<j>`$; $`F`$ is isomorphic to $`E`$ and $`p_1`$ has $`2g2`$ double fibres, each containing $`4`$ nodes of $`W`$. Now, with the usual notation, we take $`C|L|`$ a general curve and $`C^{}=h^{}C`$. Note that $`CF`$ is even, since $`p_1`$ has double fibres and the general $`C`$ contains no singular point of $`W`$. So we set $`CF=2l`$. The system $`|K_W|`$ is equal to $`|(g3)F|`$, and thus the adjunction formula on $`V`$ gives:
$$4(g1)=C^2+C^{}K_V=C^2+4l(g3).$$
If $`l=1`$, then we have $`C^2=8`$, namely $`L^2=4`$, and $`p_1`$ restricts to a $`g_2^1`$ on $`C`$, so that $`C`$ is hyperelliptic and the pair is not good. Thus proposition 8.1 implies that we are in case (a). Then we have $`0<C^2=4[g(1l)+3l1]`$, which leaves us with the only possibility $`l=2`$, $`g=4`$, $`C^2=4`$, and therefore $`C^2=2`$ and this is case (b). If the pair is good, then $`h^0(W,L)=2`$ by proposition 7.5. $``$
By proposition 5.4, only non good generating pairs can have degree $`3`$ or $`4`$ and only for very restricted values of the genus $`g`$. The following theorem gives some more information on this case:
###### Proposition 8.4
Let $`(h:VW,L)`$ be a generating pair of degree $`d>2`$ and genus $`g`$. Then one of the following holds:
i) $`d=3`$, $`q(V)=g=2`$, $`\kappa (V)0`$, $`p_g(V)1`$;
ii) $`d=3`$, $`q(V)=4`$, $`g=3`$, $`p_g(V)4`$, $`V`$ is a surface of general type and its Albanese image is a surface;
iii) $`d=4`$, $`q(V)=3`$, $`g=2`$, $`p_g(V)2`$ and the Albanese image of $`V`$ is a surface.
Proof: The possible values of $`d`$ and $`g`$ and the corresponding values of the irregularity of $`V`$ are given in proposition 5.4, as well as the assertion on the dimension of the Albanese image of $`V`$ in case iii). The claim on the dimension of the Albanese image of $`V`$ in the case ii) follows by lemma 5.2. So the Albanese image of $`V`$ is a surface, and thus $`V`$ is not ruled in these cases.
Assume now that we are in case ii), so that $`\kappa (V)1`$, by the Kodaira–Enriques classification of surfaces. If $`\kappa (V)=1`$, then the minimal model $`\overline{V}`$ of $`V`$ is equal to $`E\times B`$, where $`E`$ is an elliptic curve and $`B`$ is a smooth curve of genus $`3`$, since otherwise the Albanese image of $`V`$ would be a curve (cf. \[De1\], Lemma page 345), contradicting theorem 6.1.
We recall that $`W`$ is regular by proposition 5.4 and that the surfaces $`V`$ and $`W`$ have the same Kodaira dimension by proposition 5.5. Let $`\overline{p}:W𝐏^1`$ be the elliptic fibration on $`W`$. Clearly we have a commutative diagram:
$$\begin{array}{ccccc}& V& \stackrel{p}{}& B& \\ \hfill h& & & & \overline{h}\hfill \\ & W& \stackrel{\overline{p}}{}& 𝐏^1& \end{array}$$
where $`p:VB`$ is the composition of $`V\overline{V}`$ and of the projection $`\overline{V}=B\times EB`$, and $`\overline{h}:B𝐏^1`$ has degree $`3`$. The map $`h:VW`$, being finite, is obtained from $`\overline{h}`$ by base change and normalization. Let $`y𝐏^1`$ be such that $`\overline{h}`$ is branched at $`y`$ and assume that $`\overline{h}^{}y=2x_1+x_2`$, with $`x_1x_2`$. Since $`h`$ is étale in codimension $`1`$, in particular it is not ramified along $`p^1x_1`$. It follows that the fibre of $`\overline{p}`$ over $`y`$ must be a double fibre. But then the fibre of $`p`$ over $`x_2`$ is a double fibre, since the diagram is commutative and $`\overline{h}`$ is unramified at $`x_2`$. This is impossible, since $`p`$ is obtained from the projection $`B\times EB`$ by a composition of blow–ups. Thus the ramification points of $`\overline{p}`$ all have order $`3`$, and there are $`5`$ of them by the Hurwitz formula. By the classical Riemann construction, the covering $`\overline{p}:B𝐏^1`$ is determined, up to isomorphism, by the branch points $`y_1,\mathrm{}y_5𝐏^1`$ and by permutations $`\sigma _iS_3`$ describing the local monodromy at $`y_i`$. The $`\sigma _i`$ satisfy $`\sigma _1\mathrm{}\sigma _5=1`$. Up to renumbering the $`y_i`$, the only possibility is that there is a $`3`$-cycle $`\sigma S_3`$ such that $`\sigma _i=\sigma `$ for $`i=1\mathrm{}4`$ and $`\sigma _5=\sigma ^1`$. By \[Pa1\], proposition 2.1, there exists a cyclic cover with these properties, and thus $`\overline{h}:B𝐏^1`$ is cyclic, and the same is true for $`h:VW`$. The Galois group $`𝐙_3`$ of $`h`$ acts also on the minimal model $`\overline{V}=E\times B`$ of $`V`$, and the quotient is a surface $`\overline{W}`$ with rational singularities. The minimal resolution $`\stackrel{~}{W}`$ of $`\overline{W}`$ has invariants $`p_g=3`$, $`q=0`$. Arguing as in the proof of proposition 8.3, one shows that a generator $`\gamma `$ of $`𝐙_3`$ acts on $`V`$ by $`(b,e)(\gamma b,\sigma _be)`$, where the action of $`\gamma `$ on $`B`$ is the one associated to the Galois cover $`\overline{h}:V𝐏^1`$ and $`\sigma _b`$ is an automorphism of order $`3`$ of $`E`$ that is not a translation. The action of $`\sigma _b`$ on $`H^0(E,\omega _E)`$ is independent of $`b`$. Each of the curves $`\{x_i\}\times E`$, $`i=1\mathrm{}5`$, contains $`3`$ fixed points of the $`𝐙_3`$ action on $`\overline{V}`$, and these are the only fixed points. The surface $`\stackrel{~}{W}`$ has an $`A_2`$ singularity at the image of a fixed point $`P`$ of $`V`$ if the representation of $`𝐙_3`$ on the tangent space at $`\overline{V}`$ in $`P`$ is contained in $`SL(2,𝐂)`$ and has a singulairty of type $`\frac{1}{3}(1,1)`$ otherwise. From the above description of the $`𝐙_3`$ action, it follows that $`\overline{W}`$ has either $`12`$ points of type $`A_2`$ and $`3`$ points of type $`\frac{1}{3}(1,1)`$, or it has $`3`$ points of type $`A_2`$ and $`12`$ points of type $`\frac{1}{3}(1,1)`$. The Euler characteristics of $`\overline{V}`$ and of $`\stackrel{~}{W}`$ are related by the formula:
$$\chi (\overline{V})=3\chi (\stackrel{~}{W})\frac{1}{3}\alpha \frac{2}{3}\beta $$
where $`\alpha `$ is the number of singularities of type $`\frac{1}{3}(1,1)`$ and $`\beta `$ is the number of singularities of type $`A_2`$. Thus we have either $`\chi (\stackrel{~}{W})=3`$ or $`\chi (\stackrel{~}{W})=2`$, contradicting $`\chi (\stackrel{~}{W})=4`$. So this case does not occur, and we have shown that $`\kappa (V)=2`$ if $`q(V)=4`$. In particular, one has $`p_g(V)4`$.
We turn now to the case $`q(V)=g=2`$. Suppose that $`V`$ is ruled and denote by $`p:VB`$ the Albanese pencil of $`V`$, where $`B`$ is a curve of genus $`2`$. We can consider a minimal model $`\overline{V}`$ of $`V`$ that is a geometrically ruled surface over $`B`$. Arguing exactly as in the previous case, one shows that the Galois group of $`h`$ is isomorphic to $`𝐙_3`$. Therefore $`𝐙_3`$ acts also on $`\overline{V}`$ with $`8`$ fixed points, giving rise to $`4`$ singularities of type $`A_2`$ and $`4`$ singularities of type $`\frac{1}{3}(1,1)`$ of the quotient surface $`\overline{W}`$. The minimal desingularization of $`\overline{W}`$ is a rational surface and thus we have a contradition, using again the formula above for the Euler characteristics. In particular, $`\chi (V)0`$, namely $`p_g(V)q(V)1=1`$. $``$
Authors’ adddresses:
Ciro Ciliberto and Francesca Tovena
Università di Roma Tor Vergata
Dipartimento di Matematica
Via della Ricerca Scientifica
00133 Roma, Italia
Rita Pardini
Università di Pisa
Dipartimento di Matematica ”L. Tonelli”
Via F. Buonarroti 2
56127 Pisa, Italia
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# Accretion disk models and their X-ray reflection signatures. I. Local spectra.
## 1. Introduction
X-ray illumination of accretion disks in Active Galactic Nuclei (AGN) and Galactic Black Hole Candidates (GBHC) is a phenomenon of a great observational importance with implications for theories of accretion disks (AD). Since X-rays often produce a non-trivial part of the overall bolometric luminosity of AGN and GBHCs, it is clear that X-ray heating of the accretion disk surface may change the energy and ionization balance there, causing corresponding changes across the entire electromagnetic spectrum emitted by these objects. In addition, detailed calculations predict that spectra should contain many potentially observable atomic lines, edges, recombination continua, etc. Therefore, comparison of theoretical models (which are parameterized by few parameters only) and observations presents an invaluable opportunity to solve the inverse problem of Astrophysics of ADs – that is to learn about the accretion disk structure from observed spectra.
This is why so much theoretical effort has gone into studies of the X-ray illumination problem in the last few decades. A first detailed account of the problem can be found in Basko, Sunyaev & Titarchuck (1974)<sup>1</sup><sup>1</sup>1These authors considered X-ray illumination of the surface of a normal star in an X-ray binary, but this problem is almost identical to the one of X-ray illuminated ADs.. A further important development was done by Lightman & White (1988) and White, Lightman & Zdziarski (1988) who provided fitting formulae to the results of their Monte-Carlo simulations of reflection off neutral matter, which allowed X-ray observers to use these formulae to make direct fits to data. In addition, Fabian et al. (1989) have shown that the relativistically smeared fluorescent K$`\alpha `$ line emission of ADs yields characteritic line profiles that can be used to constrain geometry of accretion disks. A number of authors expanded on these studies since then (e.g., George & Fabian 1991; Done et al. 1992; Ross & Fabian 1993; Matt, Brandt & Fabian 1996; $`\dot{\mathrm{Z}}`$ycki et al. 1994; Czerny & $`\dot{\mathrm{Z}}`$ycki 1994; Krolik, Madau & $`\dot{\mathrm{Z}}`$ycki 1994; Magdziarz & Zdziarski 1995; Ross, Fabian & Brandt 1996; Matt, Fabian & Ross 1993, 1996; Poutanen, Nagendra & Svensson 1996; Blackman 1999)
However, except for Basko et al. (1974), the authors of the publications referenced above either studied a non-ionized reflection, or assumed that the density of the illuminated gas is constant with height for the more complicated ionized reflection calculations. Raymond (1993) and Ko & Kallman (1994) considered the problem of X-ray illumination of an accretion disk in Low Mass X-ray Binaries (LMXB) relaxing the assumption of the constant gas density and instead solving for the density via hydrostatic balance. Ró$`\dot{\mathrm{z}}`$ańska & Czerny (1996) also included hydrostatic balance in their semi-analytical study of X-ray illumination in AGN. Nayakshin, Kazanas & Kallman (2000; hereafter NKK) extended results of Ró$`\dot{\mathrm{z}}`$ańska & Czerny (1996) for the inner parts of ADs in AGN and GBHCs by providing an accurate radiation transfer in the optically thick illuminated slab. All these authors found that the thermal ionization instability, previously well known in the context of the AGN emission line regions (e.g., Krolik, McKee & Tarter 1981), plays a crucial role in the establishing of the equilibrium temperature and density profiles of the X-ray illuminated gas. Through these profiles, the instability is directly involved in the formation of the reflected continuum spectrum as well as fluorescent line emission, such as that of iron K$`\alpha `$ lines.
While the fixed density models are out of hydrostatic balance (e.g., the gas pressure at the top of the illuminated layer may be up to few hundred times larger than that on the bottom of the layer), one could still hope that the main results of such calculations will apply to the more realistic hydrostatic balance models. However, NKK compared spectra obtained from these two classes of models, and found that the predictions of these two models for the behavior of the iron K$`\alpha `$ lines, edges, and all the other features of the reflected spectrum are very different.
In this paper, we use the code of NKK to make a detailed analysis of local reflected spectra in two physically distinct limits<sup>2</sup><sup>2</sup>2We specifically limit our attention to single radius unsmeared spectra to expose the physics of the problem clearer. Full disk spectra with relativistic energy shifts will be presented in a separate publication. We will show that these two classes of models lead to completely different reflection spectra unless accretion rate is very small ($`\dot{m}\stackrel{<}{}10^3`$, see below).
In one limit, the X-ray flux illuminating the disk, $`F_\mathrm{x}`$, is smaller than $`F_\mathrm{d}`$, the flux of the soft thermal emission intrinsically generated in the disk. Physically, the situation $`F_\mathrm{x}\stackrel{<}{}F_\mathrm{d}`$ occurs in the so-called “lamppost model”, where the X-rays are produced high above the black hole, so that they illuminate a large portion of the innermost disk region. Roughly speaking, the X-ray illuminating flux is $`F_\mathrm{x}L_x/4\pi R^2`$, where $`R10R_S`$, where $`R_s`$ is Schwartzchild radius.
In the other case, the X-ray flux exceeds the disk flux by a large factor. This situation occurs when the X-ray luminosity is produced within magnetic flares, such that most of the X-ray reflection happens near the flare locations. The physical distinction from the lamppost model is that the same X-ray luminosity originates much closer to the disk surface because magnetic loops are expected to be of the order of few disk height scales, which is much smaller than radius for thin disks. Hence, the X-ray illumination will be spread over the disk surface very unevenly. The covering fraction of magnetic flares, $`f_c`$, may be quite small (see estimates in Nayakshin 1998b, §2.5.5). Most of the disk will receive little X-ray illumination, whereas near the flares $`F_\mathrm{x}L_x/4\pi f_cR^2L_x/4\pi R^2`$ and thus $`F_\mathrm{x}`$ is very likely to exceed $`F_\mathrm{d}`$. It is also well known that $`F_\mathrm{x}F_\mathrm{d}`$ is in fact a necessary condition for the magnetic flare model to reproduce the continuum X-ray and UV spectra of AGN and GBHCs (see Haardt, Maraschi & Ghisellini 1994; Svensson 1996; Nayakshin 1998a).
In addition, we discuss the ionization equilibria in the geometry of a full corona overlying a cold accretion disk. We consider coronae heated from below and also coronae heated by internal viscous dissipation (as those thought to exist in the transition region of the Advection Dominated Accretion Flows – e.g., Esin, McClintok & Narayan 1997). In both of these cases we show that the ionization equilibria permit only cold solutions for the material below the corona. Therefore, reflected spectra from such a material should look “neutral”. We believe that the predictions of these three models are sufficiently different to allow these models to be be meaningfully tested against observations with existing and future data.
The structure of the paper is as follows. In §2 we describe the way in which several parameters important for the X-ray illumination caluclations can be deduced for any accretion disk theory once broad band spectrum of an accreting source is known. In §3 and §4 we present our calculations for the lamppost and magnetic flare models, respectively. In §5 we consider the X-ray illumination for the full corona case. We give an extended discussion of our results in §6 and in §7 we present our conclusions.
## 2. On predictive power of X-ray illumination calculations
Calculations of the X-ray reflected spectra is a very powerful tool with which to infer the structure of ADs around compact objects. The value of such calculations is currently under-estimated, we believe, and this is why we will now explain how these calculations should be used in order to constrain accretion disk theories and why the results of such calculations are quite robust. Let us assume that we have a well resolved broad band spectrum of an AGN with the total luminosity $`L_{\mathrm{tot}}`$, and that the integrated optical-UV luminosity is $`L_{\mathrm{uv}}`$, whereas that of the entire X-ray range is $`L_x`$ (so that $`L_{\mathrm{tot}}=L_{\mathrm{uv}}+L_x`$).
One can now assume a value for the black hole mass, $`M`$, and then investigate a particular accretion disk theory. For example, if the disk structure is given by the Shakura-Sunyaev theory, then we can (1) find the dimensionless accretion rate $`\dot{m}=L_{\mathrm{tot}}/L_{\mathrm{Edd}}`$, where $`L_{\mathrm{Edd}}`$ is the Eddington luminosity for the mass $`M`$; (2) determine the disk flux, $`F_\mathrm{d}`$, and the disk height scale, $`H`$, for every radius. The next step is to use the geometry of the X-ray source appropriate for the given model to infer how the X-ray illuminating flux is distributed over the disk surface. With this, there are no uncertainties in the resulting reflected spectra beyond geometry and value of $`M`$. This is because the gravity parameter $`A`$, defined by NKK, is none other than the ratio of the vertical component of the gravitational force, $`_g`$, at the height of one disk height scale to the radiation pressure force, $`_{\mathrm{rad}}\sigma _tn_eF_\mathrm{x}/c`$, that the X-ray flux would provide if the cross section were given by the Thomson value ($`n_e=1.2n_H`$ is the electron density assuming that H and He are completely ionized):
$$A\frac{R_s\mu _mH^2}{2R^3m_e}\frac{m_ec^3}{\sigma _tF_xH}=\frac{GM\rho H}{R^3}\frac{c}{\sigma _tn_eF_\mathrm{x}}\frac{_g}{_{\mathrm{rad}}},$$
(1)
Hence, once the accretion disk structure is prescribed, and $`F_\mathrm{x}(r)`$ is known, the reflection calculations will provide a definite outcome that can be compared with observations.
We emphasis that this has not been possible in the context of the conventional constant density models. In the context of those models, the gas density in the midplane of the disk, $`n_H`$, is the primary parameter, since it enters the definition of the ionization parameter $`\xi =4\pi F_x/n_H`$ (see, e.g., $`\dot{\mathrm{Z}}`$ycki et al. 1994). Because $`n_H\alpha ^1`$ for radiation-dominated ADs, where $`\alpha `$ is the Shakura-Sunyaev viscosity parameter (see Shakura & Sunyaev 1973), ionization parameter scales as $`\xi \alpha `$. For gas-dominated disks, recent MHD simulations give $`\alpha 0.01`$ (e.g., Miller & Stone 2000). Observations of disks in Cataclismic Variables seem to imply $`\alpha 0.1`$ (Smak 1984; 1999), but it is also not unusual to invoke values of $`\alpha =0.3`$ (e.g., Esin et al. 1997) or even closer to unity. For radiation dominated disks, the value of $`\alpha `$ is even less certain, because it is not clear whether the viscosity in such disks will scale with total or only the gas pressure (see, e.g., Stella & Rosner 1984; Nayakshin, Rappaport & Melia 2000), so in principle the value of the “efective” $`\alpha `$ can be as small as $`10^6`$. Therefore, the results of the constant density models may be uncertain by a factor of $`10^210^6`$ (!).
Our calculations avoid this problem because the value of the gas density (or total pressure) in the midplane has very little influence on the final result. The important parameter is the height, $`z_b`$, at which the bottom of the ionized skin is located because that defines the value of the gravitational force $`_g`$. The pressure at $`z_b`$ is very small compared with the disk mid-plane pressure. Since at $`zz_b`$ the pressure declines very quickly with height (exponentially – see, e.g., Shakura & Sunyaev 1973), a large change in $`\alpha `$ will lead only to a logariphmic change in $`z_b`$ (see also Nayakshin 2000 on that). In other words, the most important parameter of our calculations – the gravity parameter $`A`$ – depends on $`\alpha `$ only logariphmically<sup>3</sup><sup>3</sup>3This is similar to stellar atmospheres – clearly, it is the gravity and the star’s radius that are important, not the gas density in the center of the star..
Finally, one does not have to assume that the structure of the illuminated AD is given by the Shakura-Sunyaev theory. Any other AD theory may be used to prescribe the vertical disk structure and then calculations proceed in exactly the same way as they would for a SS disk.
## 3. Lamp Post model: the “warm skin” limit.
### 3.1. Physical setup and method of calculation
In this section we will assume that all the X-rays are produced within a point-like source located $`h_x=6R_s`$ above the black hole on the symmetry axis. Although there is no solid physical justification for the location of the X-ray source directly above the black hole, this geometry can be used as a testing ground for studying complex phenomena, such as iron line reverberation (e.g., Reynolds & Begelman 1997; Reynolds et al. 1999; Young & Reynolds 2000). We will neglect special and general relativistic effects as we concentrate on the physics of the local (i.e., at a given disk radius $`R`$) ionization balance. We restrict our attention here to $`R=6R_s`$.
We use the code of NKK to solve the X-ray illumination problem. Several input parameters are (1) the disk accretion rate, $`\dot{m}`$, measured in the Eddington units such that $`\dot{m}=1`$ corresponds to the disk luminosity equal to the Eddington luminosity; (2) the luminosity of the X-ray source, parameterized in terms of the ratio of the latter to the disk integrated luminosity, $`\eta _x`$ (see Nayakshin 2000); (3) the cosine of the X-ray incidence angle that is fixed by values of $`h_x`$ and $`R`$ ($`\mu _i=1/\sqrt{2}`$). For the most tests in this paper, we fixed $`\eta _x`$ at a value of $`0.2`$.
### 3.2. Temperature profiles
Figure 1 shows the resulting temperature profiles for the X-ray illuminated upper layer of the disk for $`\eta _x=0.2`$ and the accretion rate $`\dot{m}`$ scanning a range of values from the low value of $`10^3`$ to high of $`0.512`$. The case with $`\dot{m}=10^3`$ is the least ionized one, and it corresponds to the left most curve in Fig. 1 (only the first zone is highly ionized, and all the rest are at the temperature $`kT8`$ eV). The next curve to the right was computed for $`\dot{m}=4\times 10^3`$, and all the subsequent curves are computed with $`\dot{m}`$ increasing by a factor of 2 from a previous value. We will refer to these runs by their number, such that the least ionized case is referred to as $`w1`$ and the most ionized as $`w9`$ ($`w`$ stands for “warm”).
The temperature of the ionized skin will later be shown to be decisive in establishing the nature of the reflected spectra. It is important to note that the maximum gas temperature is almost the same in all of these cases, and that it is far below $`T_x`$, the Compton temperature corresponding to the X-ray flux only. This is why we refer to the given limit as the “warm skin” limit to distinguish it from the “hot skin” limit studied in §4 below. In the Eddington approximation, one can easily show that the value of the Compton temperature at the top of the skin is approximately given by
$$T_c=\frac{J_xT_x+T_{\mathrm{bb}}J_{\mathrm{bb}}}{J_x+J_{\mathrm{bb}}}T_x\left[1+2\mu _i\frac{F_\mathrm{x}+F_\mathrm{d}}{F_\mathrm{x}}\right]^1.$$
(2)
From this equation it is apparent that the skin is much cooler than the X-ray Compton temperature because of the presence of the large intrinsic disk flux, $`F_\mathrm{d}`$. Therefore, physically, the warm skin limit corresponds to the case $`F_\mathrm{d}F_\mathrm{x}`$, whereas the hot skin results when $`F_\mathrm{d}F_\mathrm{x}`$ (and the X-ray spectrum is hard such that $`T_x`$ is high).
Note also that in all cases considered, except for the run $`w1`$, there is a rather extended “shelf” with $`kT100`$ eV, whereas this shelf is absent or weak in the hot skin limit (see Fig. 6 below). We give an explanation for this fact in §6.1.
### 3.3. Reflected spectra
Figure 2 shows angle-averaged reflected spectra in the range from 10 eV to about 100 keV. The ionizing X-ray spectrum is drawn with the dashed-dotted curve. All the spectra are normalized in such a way that $`F_\mathrm{x}^1_0^{\mathrm{}}E𝑑EF(E)=1+F_\mathrm{d}/F_\mathrm{x}`$ for convinience. The least ionized reflected spectrum corresponds to the run $`w1`$, and the subsequently more ionized spectra are for the tests $`w3`$, $`w5`$, $`w7`$, and $`w9`$. The coldest case ($`w1`$) is basically the usual neutral reflection component that has almost no reflected flux below few keV all the way down to $`E`$ few $`\times kT_{\mathrm{eff}}`$, where the emission is dominated by a quasi-black body reprocessed flux.
Figure 3 shows the same spectra in a narrower energy range. The spectra are also shifted with respect to each other to allow for a greater visibility of the individual curves. It is immediately clear that except for the least ionized case on the bottom of the Figure, all the other spectra are strikingly different from the neutral reflection spectrum. Note that the energy of the dominant line component is around 6.7 keV and it comes from He-like iron mainly (see §3.4 below). The line Equivalen Width (EW) is about 180 eV in the coldest case and it increases to about 500 eV in the run $`w9`$, which is consistent with earlier findings of Matt, Fabian & Ross (1996). Note also a very prominent ionized absorption edge and the recombination continuum at $`9`$ keV. Other significant spectral features are labeled directly in Figure 3.
### 3.4. Ionization Structure of the gas
The key to understanding the reflection spectra is the ionization structure of the gas, shown in Figure (4) for iron for run $`w3`$. Panel (a) of the Figure shows ionic fractions of different ionization stages of iron (Fe27 $``$ completely ionized iron) as a function of the Thomson depth. These fractions are defined as the ratio $`n_i/n_{\mathrm{Fe}}`$, where $`n_i`$ is the density of the ionization stage $`i`$ and $`n_{\mathrm{Fe}}`$ is the total local density of iron. The most important fact from the Figure is that the skin is highly ionized, but it is far from being completely ionized. Even in the uppermost zone, more than 50% of iron is in the form of Fe26 and Fe25, and this fraction quickly increases to almost 100% at deeper layers (and then goes to zero as less ionized ions start to dominate). Panel (b) of Figure 4 shows the K$`\alpha `$ line emissivity as a function of Thomson depth. The line emissivity is split onto several components. In particular, the dashed line shows the line emission in the energy bin that contains the H-like iron line at $`E6.97`$ keV; the dotted curve shows the line emission in the bin that contains the He-like line at $`E6.7`$ keV plus some lines from intermediate ionization stages of ions Fe17-Fe23; and finally, the solid curve shows the line emission in the bin that includes the “neutral”-like line from ions Fe1-Fe16 with $`E6.4`$ keV.
Note that the skin emits a strong He-like line. Moreover, in the test $`w9`$, whose ionization structure is shown in Figure 5, the “6.7 keV line” bin dominates the line emissivity by far. In addition, other strong features from the soft X-ray band to the iron recombination continuum allow us to state the following general point: because the warm skin is not completely ionized, it produces strong atomic emission and absorption features that should be observable in spectra of many bright AGN (if this model is the physically correct one).
## 4. X-rays from magnetic flares: hot skin limit
### 4.1. Setup
We now will study the case in which the ionizing flux is much larger than the disk soft flux, $`F_\mathrm{x}F_\mathrm{d}`$, which is thought to be appropriate for the two-phase patchy corona model of Seyfert Galaxies (e.g., Haardt, Maraschi & Ghisellini 1994 and Svensson 1996). In that model, the X-rays are produced within magnetic flares (Galeev, Rosner & Vaiana 1979). The condition $`F_\mathrm{x}F_\mathrm{d}`$ is a necessary condition for the model to work, since it comes from the requirement that Compton cooling of the X-ray producing active regions by the soft disk photons not be too strong, because otherwise the continuum spectra are too steep to explain typical Seyfert 1 spectra and much less that of GBHCs in their hard state (e.g., Gierlinski et al. 1997; Dove et al. 1997; Nayakshin 1998a). As we discussed in NKK, the X-rays induce evaporative winds close to the flare location, so that hydrostatic balance does not apply in a direct sense. However, based on earlier work on X-ray induced winds in stars (e.g., McCray & Hatchett 1975; Basko et al. 1977; London, McCray, & Auer 1981, and references therein) one expects to see the same two-layer structure for the illuminated gas. In the geometry of thin accretion disks, the winds will decrease the value of the Thomson depth of the skin by pushing the ionized material along the disk away from the flare location. A careful multi-dimensional calculation including gas dynamics, ionization and radition transfer is needed. Such calculation is beyond the scope of this paper, but we attempt to model the effects of the wind in a simple manner.
From the point of view of ionization calculations, the gas density structure is of paramaunt importance for a self-consistent solution. The gas density where the temperature discontinuity sets in is expected to be the same or roughly the same in both static and wind situations because this density is set by the energy balance that is largely given by radiation heating/cooling and thus is independent of whether the gas is static or moving non-relativistically. The density gradient is then the quantity of the primary interest. In the hydrostatic case, this gradient is given by the gravity in the approximately isothermal skin (the gas temperature in the skin varies by a factor of order of few only, see Nayakshin 2000). In the case of the wind, the gas density gradient is controlled by the radiation pressure, gravity and gas pressure forces. A simple estimate shows that the ratio of the gravitational force to the radiation pressure force is of the order of $`_g/_{\mathrm{rad}}F_\mathrm{d}/F_\mathrm{x}`$, which is much less than unity for magnetic flares. Hence, the gas density gradient is controlled by the radiation pressure mainly and is (very roughly!) $`F_\mathrm{x}/F_\mathrm{d}`$ times the one which is given by the gravity.
In the absense of a multi-dimensional radiation hydrodynamics and radiation transfer approach, we resort to treating the complications due to the wind by artificially increasing the value of the local gravity in the atmosphere of the disk, which, in some sense, leads to the same end result – a lower value of $`\tau _\mathrm{s}`$ for a given $`F_\mathrm{x}`$ (as compared with the one that would have been obtained if there were no winds). Mathematically this is done by multiplying the local gravity gradient ($`GMz/R^3`$) by a dimensionless number $`𝒜1`$. Note that this approach is similar to the one we adopted in most of the calculations presented in NKK, although the gravity parameter $`A`$ defined there is different from the parameter $`𝒜`$ defined here.
In terms of the final results for the magnetic flare model, the poorly constrained value of $`𝒜`$ means that the Thomson depth of the ionized skin is not calculated exactly and represents a current uncertainty of the model. However, we believe that, for a given value of $`\tau _\mathrm{s}`$, the ionization structure of the illuminated gas is calculated with a reasonble precision and thus this aspect of calculations is reliable. Moreover, it is encouraging that the uncertainties of the model due to the wind do not present an unsurmountable obstacle because X-ray induced winds from accretion disks and stars have been treated by many different authors (e.g., Buff & McCray 1974; McCray & Hatchett 1975; Begelman, McKee & Shields 1983; Proga, Stone & Drew 1999). We plan to incorporate evaporative winds in our calculations in the future.
### 4.2. Temperature profiles
Figure 6 shows temperature profiles for the illuminated gas. From these profiles, one can immediately see why we refer to the limit $`F_\mathrm{x}F_\mathrm{d}`$ as the hot skin limit. The Compton temperature of the gas is much higher than that in the case $`F_\mathrm{x}\stackrel{<}{}F_\mathrm{d}`$ which can be understood from equation 2. Although the X-ray Compton temperature $`T_x`$ is the same in both cases, the local Compton temperature is determined by the X-rays and the disk soft flux, so that in the case $`F_\mathrm{d}F_\mathrm{x}`$, $`T_c`$ is a small fraction of $`T_x`$. Also note that the middle branch of the S-curve (the region with $`kT90150`$ eV) does not appear in Figure 6 until the skin becomes moderately thick. This is due to the fact that the middle branch of the S-curve is available only when the average photon energy in the overall illuminating specrum is low enough (see §6.1). The latter condition can be satisfied in the given case only by Compton downscattering of the incident hard X-rays, which becomes significant when $`\tau _\mathrm{s}1`$.
### 4.3. Reflected spectra
Figures 7 and 8 show the reflected spectra for the hot skin limit. This is the limit previously discussed by NKK. The strongest line component is at 6.4 keV, the iron absorption edge is neutral-like and is weak, and only few of the spectra exhibit some soft X-ray emission. The latter is due to the fact that the skin is more strongly ionized than it is in the warm skin limit, and thus reflection in the soft X-ray range is dominated by the Compton reflection in the skin. Below $`30`$ keV, the highest ionization spectra can hardly be distinguished from the illuminating power-law, especially if one adds the additional smearing due to relativistic effects in vicinity of the black hole. This might be the reason why GBHCs do not show broad iron lines and a somewhat small reflection covering fraction (see Gierlinski et al. 1997 and references therein, and also see Nayakshin 1998a; Ross, Fabian & Young 1999; Done & Nayakshin 2000).
### 4.4. Ionization Structure of the gas
Figures 9 and 10 show the ionic fractions and line emissivity of the illuminated gas for tests $`h3`$ and $`h7`$, respectively. These figures are to be compared with the figures 4 and 5 for the warm skin tests $`w3`$ and $`w9`$. The biggest difference is in the fact that the hot skin is almost 100% ionized except for regions near the temperature discontinuity, whereas the warm skin is dominated by H- and He-like iron. This is the reason why the line emissivity is dominated by the “cold” line coming from the cool material below the skin. The skin thus only masks the presence of the cold gas. In addition, because the skin is hot, Oxygen, Sulfur, Silicon and Neon are completely stripped of their electrons, and thus there are no broad recombination edges or line emission from these elements at least from the skin itself.
## 5. Full corona above a standard disk
We will now consider X-ray reflection in the geometry of a full corona covering the whole inner accretion disk. This geometrical arrangement is physically rather different from the two other geometries that we studied so far, because for the latter two, the X-ray producing region does not directly border the illuminated surface of the disk. In other words, there is no material above the top of the skin, and hence the gas pressure at the top of the skin, $`P_{\mathrm{gas}}(z_\mathrm{t})`$, is zero. This fact is used as an explicit boundary condition and affects all aspects of the X-ray illumination problem. In the case of the full corona, the pressure at the top of the skin is not zero, and by continuity arguments, $`P_{\mathrm{gas}}(z=z_\mathrm{t})=P_{\mathrm{cor}}(z_\mathrm{t})`$, where $`P_{\mathrm{cor}}(z_\mathrm{t})`$ is the gas pressure at the bottom of the corona. Assuming that the X-ray flux produced in the corona is radiated isotropically, there is an equal amount of the X-rays emitted down to the disk and up away from the disk.
Let us first assume that the corona is heated from below. In equilibrium, there must be $`F_ϵ=2F_\mathrm{x}`$ heating flux from the cold disk to the corona. Some agent, most likely a magnetic field with magnetic pressure $`P_{\mathrm{mag}}`$ will have to carry the energy into the corona. The physical distinction between the skin and the hot corona is then such that the magnetic fields heat the corona but not the skin. Further, $`F_ϵP_{\mathrm{mag}}v`$, where $`v`$ is the speed with which the energy flux is being carried. The maximum value of $`v`$ is roughly the greater of $`c_s`$ and $`v_A`$, where $`c_s`$ is the sound speed and $`v_A`$ is the Alfvén velocity.
Let the gas pressure at the top of the skin be related to the magnetic pressure as $`P_{\mathrm{gas}}=P_{\mathrm{mag}}/\beta `$, where $`\beta `$ is a dimensionless number. Hence, $`vc_s(1+\beta ^{1/2})`$. The gas pressure at the top of the skin is
$$P_{\mathrm{gas}}(z_\mathrm{t})=\beta ^1P_{\mathrm{mag}}\frac{2c}{v\beta }\frac{F_\mathrm{x}}{c}\frac{2\times 10^3T_1^{1/2}}{\beta (1+\sqrt{\beta })}\frac{F_\mathrm{x}}{c},$$
(3)
where $`T_1`$ is the Compton temperature in units of 1 keV. It is interesting to compare this pressure with the “critical” pressure $`P_c`$. This quantity is defined to be the gas pressure at point (c) on the S-curve (see Fig. 11), i.e., at the location where the temperature discontinuity occurs. Nayakshin (2000) found this quantity to be
$$P_c=0.032T_1^{3/2}\frac{J_{\mathrm{tot}}}{c},$$
(4)
where $`J_{\mathrm{tot}}F_\mathrm{x}`$ is the total intensity of radiation integrated over all angles at $`z=z_\mathrm{t}`$ (see also Krolik et al. 1981). Thus,
$$\frac{P_{\mathrm{gas}}}{P_c}6\times 10^4\frac{1}{\beta (1+\sqrt{\beta })}T_1^21.$$
(5)
That is, unless $`\beta P_{\mathrm{mag}}/P_{\mathrm{gas}}10^3`$ in the skin, the gas pressure exceeds the one at which the hottest branch of the solution (i.e., the completely ionized skin) exists. Note that $`\beta `$, in fact, is likely to be less than unity. We defined the skin as a region where no magnetic heating occurs. However, magnetic reconnection often occurs when the magnetic pressure starts to exceed the gas pressure, i.e., when $`\beta 1`$.
This estimate was done for a corona heated from below. A different scenario arizes when the corona is heated internally – via viscous dissipation of the accretion energy of the gas flowing through the corona itself. In the latter case, the radiation flux from the corona can be written as
$$F_\mathrm{x}=ϵ\frac{3}{8\pi }\frac{GM\dot{M}_c}{R^3}J(R),$$
(6)
where $`J(R)=(1\sqrt{3R_s/R})`$, and $`ϵ=1`$ for the complete transfer of the accretion energy into radiation (as in Shakura-Sunyaev disks), while $`ϵ`$ is less than unity for an advection dominated corona (e.g., Esin et al. 1997). Further, the accretion rate through the corona, $`\dot{M}_c=2\pi R\mathrm{\Sigma }_{\mathrm{cor}}v_{rc}`$ where $`\mathrm{\Sigma }_{\mathrm{cor}}`$, and $`v_{rc}`$ are the corona mass column density and radial inflow velocity, respectively. The gas pressure in the corona may be estimated via hydrostatic balance:
$$P_{\mathrm{cor}}=\frac{1}{2}\frac{GMH_c}{R^3}\mathrm{\Sigma }_{\mathrm{cor}},$$
(7)
where $`H_c`$ is the height scale of the corona. Using these two equations, we can now conclude that
$$\frac{P_{\mathrm{cor}}c}{F_\mathrm{x}}=\frac{2}{3ϵ}\frac{H_c}{R}\frac{c}{v_{rc}}$$
(8)
If the coronal accretion energy is radiated locally, then we can use the standard equation for the radial inflow velocity: $`v_{rc}\alpha c_s(H_c/R)^2`$, where $`c_s`$ is the sound speed in the corona. Thus, one obtains
$$\frac{P_{\mathrm{gas}}}{P_c}\stackrel{>}{}21T_1^{3/2}\alpha ^1\frac{c}{c_s}\frac{R}{H_c}1.$$
(9)
If coronal cooling is dominated by advection, then $`ϵ1`$, and so even though $`H_cR`$ and $`v_{rc}v_K`$, where $`v_K`$ is Keplerian velocity (see, e.g., Narayan & Yi 1994), $`P_{\mathrm{gas}}`$ is still large compared with $`P_c`$.
Summarizing, in all of the three cases for the full corona above the disk, we concluded that the gas pressure in the skin is very much larger than the pressure at which the thermal instability operates. Thus, the gas pressure and density below the corona are too high for the Compton-bremsstrahlung stable branch of the solution to exist, and hence no ionized skin forms below the hot corona. The gas temperature below the corona is therefore close to the effective one and the reflection and the lines will be those that are produced in a “neutral” material. This conclusion holds for arbitrarily large accretion rates. Note, however, that a thin transition layer may still form due to conductive heating of the disk (e.g., Maciolek-Niedzwiecki, Krolik & Zdziarski 1997).
## 6. Discussion
In this paper, we considered three physically distinct geometries: the lamppost geometry; the magnetic flares above the disk and the full corona above the disk. Each of these geometries was shown to produce different reflected spectra. The large differences in the spectra were caused by (i) absense or presence of a large soft disk flux that influences the Compton temperature; and (ii) the additional pressure, or weight, of the corona for the full corona geometry. While it is perhaps possible to find conditions in which these three models may yield similar spectra (one example is a very low accretion rate when the ionized skin is very Thomson thin and thus negligible), the behavior of the spectra with the X-ray luminosity as well as other parameters is clearly different which should allow one to distinguish between these models observationally.
We do not discuss observational status of the lamppost and the magnetic flare models because we will present an extended discussion of this topic in a future publication where we will also present complete disk spectra that include relativistic broadening. However, it appears to us that the full corona model is the most unpromising of all and we will not study this model in our future work. In particular, since no skin forms below the hot corona, it is not possible to explain the hard continuum spectra of GBHCs (e.g., see Gierlinski et al. 1997; Dove et al. 1997, Nayakshin 1998a). Moreover, there are AGN that have more optical/UV power than the X-ray power which argues against the full corona geometry (see Haardt et al 1994) in which the reprocessed power should be rather less than the X-rays. In addition, as we have shown in §5, the reflector below the skin is cold, which means that the small reflection covering fraction seen in many GBHCs (e.g., Gierlinski et al. 1997) is problematic. The lack of the iron line reverberation discovered in the recent observations of two Seyfet 1 Galaxies (Chiang et al. 2000; Reynolds 2000; Lee et al. 2000) presents yet another uneasy challenge to this model.
Now we will concentrate our discussion on the differences between the lamppost and the magnetic flare models.
### 6.1. Changes in the S-curve due to soft disk flux
In this paper we emphasised the large difference in the resulting temperature profiles, ionization structure and the reflected spectra between the hot and the warm skin limits. For a clearer understanding of our results, it is useful to discuss this difference in the simplest case – optically thin heating/cooling balance. Figure 11 shows the energy equilibrium curves for the spectral index $`\mathrm{\Gamma }=1.8`$ and a constant angle-integrated intensity of radiation $`4\pi J=10^{16}`$ erg cm<sup>-2</sup> s<sup>-1</sup>. The overall intensity $`J`$ in these tests is a sum of the X-ray intensity, $`J_x`$ and the black-body intensity $`J_{\mathrm{bb}}`$ (with $`kT=10eV`$). The only difference between the curves is the fraction of the black-body intensity compared with that of the hard power-law. In particular, the curves are computed for $`J_{\mathrm{bb}}/J_x=`$ 0, 1, 2, 4, and 8.
The temperature profiles shown in Figures 1 and 6 can now be discussed with the help of figure 11. The maximum temperature reached in the two opposite limits is different due to the fact that the Compton-bremsstrahlung stable branch of the S-curve shifts significantly as the ionizing spectrum evolves from “all X-ray” to “all black-body”. The X-ray Compton temperature for the given spectrum is about 7.1 keV. Stars in figure 11 indicate the value of the Compton temperature calculated from equation 2, which is clearly smaller for larger values of $`J_{\mathrm{bb}}/J_x`$. The fact that the corresponding curves pass right through the stars shows that the gas temperature is indeed very close to the corresponding Compton temperature when $`\mathrm{\Sigma }_x1`$.
Another important difference between the temperature profiles for the hot and warm skin limits is the extent of the middle stable shelf in the S-curve. In the case of $`F_\mathrm{x}\stackrel{<}{}F_\mathrm{d}`$, the middle stable shelf is always prominent in the temperature profiles, whereas in the opposite limit when $`F_\mathrm{x}F_\mathrm{d}`$, the middle stable branch only appears for skin that is at least moderately Thomson thick. In order to understand that, we recall that the thermal conduction picks solutions with the smallest temperature gradients (see NKK), and thus the transition from the Compton-bremsstrahlung branch to the middle one happens after the gas pressure exceeds the pressure $`P_c`$ at point $`c`$. Thus, in order for the middle shelf to be present in the temperature profiles, the point $`c`$ should lie to the right of the turning point at $`kT80`$ eV in Figure 11.
When there is little blackbody flux, the gas pressure at point $`c`$ is large enough that the transition happens directly from the Compton-heated to the cold branch. As the blackbody intensity increases, $`T_c`$ decreases, and then the transition happens first to the middle branch and only later to the cold solution. However, with increase in the blackbody intensity, point $`c`$ moves to the right as much as to make the whole range of temperatures from $`kT80eV`$ to the Compton one to be thermally stable.
In addition, even if $`J_{\mathrm{bb}}=0`$, the incident X-rays are downscattered in the skin and thus for a moderately Thomson thick skin point $`c`$ again moves to the right (the more the thicker the skin is). This is best understood from the fact that the location of the point $`c`$ sensitively depends on the Compton temperature, which is affected by scatterings. Analytical theory of the Compton-bremsstrahlung cooled branch (see Krolik et al 1981; Nayakshin 2000) yields that the gas temperature at point $`c`$ is a third of the Compton temperature, and pressure $`P_c`$ is given by equation 4.
The square boxes in Figure 11 indicate the location of the critical point $`c`$ found via equation 4. Note that when there is no black-body contribution to the ionizing intensity, the analytical theory of the Compton-bremstrahlung cooled upper branch works perfectly. Physically, this occurs because the gas is completely ionized and atomic heating/cooling is negligible. However, when $`J_{\mathrm{bb}}\stackrel{>}{}J_x`$, the gas becomes cool enough and heating due to photo-absorption cannot be neglected; this is the reason why the analytical solutions under-estimate the value of the gas temperature at point $`c`$.
### 6.2. Iron ionic fractions and lines
Figure 12 summarizes the differences in the ionization structure of the skin in the hot and warm limits by presenting the integrated Thomoson depth of the ions Fe13 – Fe27 for both of the limiting cases. This depth is defined as the integral $`_0^{\mathrm{}}𝑑\tau _t(n_i/n_{Fe})`$, where $`n_i`$ is the ion density for ion $`i`$. The curves shown in the Figure correspond to tests presented in Figures 3 & 8, respectively. The largest difference between the ionization structure of the two sets of calculations is seen for He- and H-like iron and for completely ionized iron. For the warm skin cases, the first two ions by far outweight the presence of Fe27, as well as any of the less ionized stages. For the hot skin, the situation is reverse, and this is why the skin is “invisible” to the observer in this case and can only be uncovered in the hard X-ray energy range where relativistic rollover in the Klein-Nishina cross section produce a corresponding rollover in the reflected spectrum.
There is also a substantial difference in the equivalent width and energy of the iron K$`\alpha `$ lines between the warm and hot skin limits. In the former case, the most abundunt ion in the skin is the He-like iron, and therefore the strongest line is also He-like with energy around 6.7 keV. H-like line and lines from Fe17-23 also contribute to the complex of the emitted lines. Because He-like iron has a large fluorescence yield compared with neutral iron, the warm skin actually yields the strongest iron line emission. In a sense, it amplifies the line.
The opposite limit of the strong X-ray flux, $`F_\mathrm{x}F_\mathrm{d}`$, is characterized by the predominance of completely ionized iron in the skin. The X-rays incident on the skin are Compton scattered but not photo-absorbed with consecutive iron line fluorescence. Only those photons that are able to reach the cold layers, whether scattered or unscattered in the skin, will produce the “neutral”-like K$`\alpha `$ line, which still needs to make it through the skin to the observer (some of the line photons may be scattered in the skin, return to the cold material and be photo-absorbed without re-emission as a 6.4 keV line photon). Note also that scattering in the skin is very effective in removing the photons from the line bin, because single scattering disperses photons by $`\mathrm{\Delta }E/E0.06T_1^{1/2}\stackrel{>}{}0.1`$, or about 700 eV at the line energy. Therefore, the hot skin only degrades the cold-like line by reflecting the incident X-rays and by scattering the line photons emitted from the cool layers.
### 6.3. On generality of our results
#### 6.3.1 Ratio of X-ray and disk fluxes
In terms of the ratio of the illuminating flux $`F_\mathrm{x}`$ to the disk flux $`F_\mathrm{d}`$, we limited our attention so far to two rather extreme values – one quite small ($`\eta _x=0.2`$) and the other very large, $`F_\mathrm{x}/F_\mathrm{d}=75`$. Thus, one wonders how our results will change when $`\eta _x`$ is inbetween these two values. To address this question, we conducted several additional series of tests. In the first series, we fixed the accretion rate through the disk at $`\dot{m}=0.25`$, and varied the ratio $`\eta _x=L_x/L_d`$ from $`2^1`$ to $`2^4`$ in steps of factor of 2. These tests are shown in Figure 13. The reflected spectra show a strong He-like line, edge and the recombination continuum for all $`\eta _x`$ below 8. In fact, even for the largest values of $`\eta _x`$ (8 and 16), there still exist a substantial column density of H- and He-like iron, but it occurs rather deep in the skin, i.e., at $`\tau _T1`$, so that it is smeared out by Compton scatterings. The skin is very Thomson thick in these two latter tests because $`L_x`$ is supper-Eddington.
The fact that there is a large column density in the last two not completely ionized stages of iron means that their presence may probably be seen for lower values of $`\dot{m}`$ when the completely ionized part of the skin is not Thomson thick. To test that, we calculated reflected spectra for same values of $`\eta _x`$, but for the accretion rate through the disk varying as $`\dot{m}=0.25/(0.5+\eta _x)`$. Results of these tests are shown in Figure 14. Finally, we also simulated reflected spectra for $`\dot{m}=0.05/(0.5+\eta _x)`$, which we show in Figure 15. Both of these series of calculations indeed show that the Hydrogen and Helium-like iron still survives deep inside the ionized skin even for $`\eta _x=16`$. This is somewhat unexpected since $`F_\mathrm{x}10F_\mathrm{d}F_\mathrm{d}`$ in the latter case, and hence one could expect the skin to be in the hot limit.
Let us try to explain this result. The gas temperature in the skin is a fraction of the local Compton temperature, which is determined by equation 2. At large optical depth, one can use Eddington approximation for the radiation transfer, and obtain that the black-body radiation field is
$$J_{\mathrm{bb}}=\left\{(1a_s)+\frac{F_{\mathrm{disk}}}{F_\mathrm{x}}\right\}(1+3\tau /2),$$
(10)
where $`a_s`$ is the integrated X-ray albedo of the skin. When $`\tau _\mathrm{s}1`$, the albedo becomes large, i.e., $`1a_s1`$, and thus in this situation one really compares $`F_\mathrm{d}/F_\mathrm{x}`$ with a small quantity $`1a_s`$ rather than with unity. Finally, when the Thomson depth of the skin is greater than unity, Compton down-scattering of the incident X-rays makes $`T_x`$ in equation 2 to be dependent on the location in the skin – it is lower on the bottom of the skin than on its top. This is an additional mechanism to lower $`T_c`$. For these reasons, it turns out to be enough for the disk flux to be a small fraction of $`F_\mathrm{x}`$ to present a large cooling source for the gas at the bottom of the skin, and this is why one needs to go to rather large values of $`L_x/L_d\stackrel{>}{}10`$ to make the lamppost spectra look similar to that from magnetic flares.
#### 6.3.2 The high energy rollover of the spectrum and spectral index
Another degree of freedom which we have not explored in this paper is the rollover energy, or equivalently, the gas temperature of the X-ray producing region or regions. Our illuminating spectrum is chosen to be reminiscent of what is typically seen in Seyfert 1 Galaxies and GBHCs (e.g., Svensson 1996; Gierlinski et al. 1997), and this is why we set the rollover energy at $`E_c=200`$ keV. However, “odd” objects such as low-luminosity AGN do not have to have the same values of $`E_c`$, and it appears to us to be potentially important for the reflected spectra. In particular, $`E_c`$ figures prominently in the determination of the X-ray Compton temperature $`T_x`$ if the X-ray spectrum is hard (i.e., $`\mathrm{\Gamma }\stackrel{<}{}2`$). If $`E_c`$ is significantly lower than 200 keV assumed in this paper, then the Compton temperature of the skin may be low enough even if the spectrum is hard and $`F_\mathrm{x}F_\mathrm{d}`$. We will present concrete calculations showing these effects elsewhere.
Finally, the spectral index $`\mathrm{\Gamma }`$ also spans a range of values for real sources, and it is well known that the reflected spectra strongly depend on the actual value of $`\mathrm{\Gamma }`$ (e.g., NKK). Notwithstanding the complications connected with the two additional parameters, an encouraging fact is that both of these parameters can in principle be extracted from observations if a broad band spectrum is available.
## 7. Summary
In this paper, we have sketched the way in which X-ray reflection spectra may be computed for different accretion disk theories, and then showed example calculations for two models (lamppost and flares) and also discussed spectra from the full corona disk model. The main three parameters that determine the outcome of the X-ray illumination problem are (1) the ratio of the X-ray illuminating flux to the disk thermal flux, $`F_\mathrm{x}/F_\mathrm{d}`$, because it defines the Compton temperature of the corona and thus the degree to which it will be ionized; (2) the accretion rate through the disk since it defines the height at which the skin is located ($`H`$ to few $`H`$); (3) Compton temperature that is a function of the spectral index of the ionizing radiation, $`\mathrm{\Gamma }`$, and the spectral cutoff energy, $`E_c`$. These parameters can be readily deduced from observations once UV and X-ray spectra, luminosities and at least an estimate of the black hole mass for the given source are available. Therefore, for any accretion disk model which clearly specifies a connection between the overall observed X-ray luminosity $`L_x`$ and the radius-dependent ionizing flux $`F_\mathrm{x}`$, accurate disk-integrated spectra can be calculated. These spectra are weakly dependent on the unknown value of the viscosity parameter $`\alpha `$ (see §2).
Our main results, valid for relatively hard X-ray spectra ($`\mathrm{\Gamma }\stackrel{<}{}2`$) and the rollover energy $`E_c200`$ keV, are as following:
* If the incident X-ray flux is smaller than or comparable with the soft thermal flux generated intrinsically in the disk, then the Compton-heated skin is “warm”, i.e., $`kT\stackrel{<}{}1`$ keV. In that case iron in the skin is not completely ionized and the majority of it is in the form of hydrogen and helium-like ions. In this limit, the iron line is dominated by the He-like line at 6.7 keV and the line equivalent width increases with the ionizing X-ray luminosity.
* In addition, in the warm skin limit, the ionization physics permits the existence of a large shelf of material occupying the middle stable branch of the S-curve (see Figure 11). Medium Z-elements such as oxygen are not completely ionized there and thus they produce strong soft X-ray features which should be visible with modern X-ray telescopes such as Chandra and XMM.
* If the illuminating X-ray flux is much higher than the soft disk flux, then the skin temperature is substantially higher. Iron is then mostly completely ionized in the skin and thus the skin does not emit or absorbs photons due to atomic processes. Therefore, the hot skin only masks the presence of the cold material. The larger the ionizing X-ray luminosity, the less atomic marks one sees in the reflected spectra. In particular, the iron line is emitted almost exclusively by the cold material below the skin and its energy is close to 6.4 keV. Further, EW of the line decreases with increasing $`L_x`$ and eventually goes to zero.
* Because ionization equilibria depend on the ratio $`F_\mathrm{x}/F_\mathrm{d}`$, the middle temperature solutions are not present in the hot skin, and thus there are less soft X-ray emission features from medium Z-elements. For high $`L/L_{\mathrm{Edd}}`$ in the magnetic flare model, the reflected spectra are featurless power-laws up to the Compton rollover at $`30`$ keV.
* Full coronae above accretion disks do not have the highly ionized skin between the corona and the disk, because the gas density and pressure are too high to allow for that. This means that X-ray reflection from full corona models, if not broadened to invisibility by scattering in the corona, is always “cold” and does not change with increase in the accretion rate through the disk or the corona. We believe this model is rulled out observationally (see §6).
We hope that the significant differences in the spectra of these three models make it possible to distinguish among them observationally based on current and future data.
The authors acknowledge many stimulating discussions with Manuel Bautista, Demos Kazanas, and Julian Krolik. SN acknowledges National Research Council Associateship which fully supported this research.
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# Breather initial profiles in chains of weakly coupled anharmonic oscillators
## 1 Introduction
Chains of coupled anharmonic oscillators have many applications in condensed matter and biophysics as simple one-dimensional models of crystals or biomolecules. Of particular interest are the so-called “breather” solutions supported by such chains, that is, oscillatory solutions which are periodic in time and exponentially localized in space . The simplest possible class of models, where all the oscillators are identical, with one degree of freedom, and nearest neighbours are coupled by identical Hooke’s law springs has been widely studied in this context. The equation of motion for the position $`q_n(t)`$ of the $`n`$-th oscillator ($`n`$) is
$$\ddot{q}_n\alpha (q_{n+1}2q_n+q_{n1})+V^{}(q_n)=0$$
(1)
where $`\alpha `$ is the spring constant and $`V`$ is the anharmonic substrate potential, which we choose to normalize so that $`V^{}(0)=0`$ and $`V^{\prime \prime }(0)=1`$. MacKay and Aubry have proved the existence of breathers in the weak coupling (small $`\alpha `$) regime of this system . They noted that in the limit $`\alpha 0`$, system (1) supports a one-site breather, call it $`𝐪_0=(q_n)_n`$, where one site ($`n=0`$ say) oscillates with period $`T`$ while all the others remain stationary at $`0`$. Using an implicit function theorem argument they proved existence of a constant period continuation $`𝐪_\alpha `$ of this breather away from $`\alpha =0`$ provided $`\alpha `$ remains sufficiently small and $`T2\pi ^+`$.
Having established the existence of these breathers, the question remains: what do they look like? This question may be addressed by means of numerical analysis in several ways. One technique is to seek period $`T`$ solutions of a spatially truncated (i.e. $`|n|N<\mathrm{}`$) version of (1) by searching for fixed points of the period $`T`$ Poincaré return map using, for example, a Newton-Raphson algorithm . This method potentially allows the construction of breathers far from $`\alpha =0`$, and can be used to determine the domain of existence of continued one-site breathers in the $`(\alpha ,T)`$ parameter space (as applied by the authors to a non-standard discrete sine-Gordon system in ). This technique has its drawbacks, however. It is rather computationally expensive and there are subtleties concerning convergence of the Newton-Raphson method close to $`\alpha =0`$ . On the other hand, if the weak coupling regime is of primary interest, then much useful information may be obtained from
$$𝐪_0^{}:=\frac{𝐪_\alpha }{\alpha }|_{\alpha =0},$$
(2)
the tangent vector at $`\alpha =0`$ to the continuation curve $`𝐪_\alpha `$. This vector may be calculated by solving a certain initial value problem for a set of 4 coupled second order ODEs, a numerically trivial task. The small $`\alpha `$ behaviour of breathers can then be determined by approximating the curve $`𝐪_\alpha `$ by its tangent line at $`\alpha =0`$, that is
$$𝐪_\alpha =𝐪_0+\alpha 𝐪_0^{}+o(\alpha ).$$
(3)
Since this method is so computationally cheap, it is easy to make a systematic study of the period dependence of breather initial profiles for a variety of substrate potentials. This paper presents such a study.
The paper is organized as follows. In section 2 we give a precise statement of the MacKay-Aubry existence theorem and reduce the calculation of $`𝐪_0^{}`$ to numerically tractable form. In section 3 we present numerically generated graphs of the components of $`𝐪_0^{}`$ against period $`T`$ for various choices of potential, and extract from them information about breather initial profiles in the small $`\alpha `$ regime. The results show a generic behaviour of alternating $`T`$-bands of two qualitatively different types of breather, which we call “in-phase breathers” and “anti-phase breathers”. In section 4 we prove that this alternating behaviour is generic in a precise sense, and examine the small $`T`$ limit analytically. Section 5 contains some concluding remarks.
## 2 The direction of continuation of one-site breathers
In the following we will assume that $`V:`$ is twice continuously differentiable and has a normalized stable equilibrium point at $`0`$ ($`V^{}(0)=0`$, $`V^{\prime \prime }(0)=1`$). Consider the equation of motion for a particle moving in such a potential,
$$\ddot{x}+V^{}(x)=0$$
(4)
with $`\dot{x}(0)=0`$. Provided $`|x(0)|`$ is small enough, $`x(t)`$ must be a periodic oscillation. All the potentials we consider will be anharmonic with classical spectrum $`(2\pi ,\mathrm{})`$. Anharmonic means that the period of oscillation varies nondegenerately with $`x(0)>0`$. The classical spectrum of $`V`$ is the set of periods of the oscillations supported by $`V`$.
Let $`x_T(t)`$ denote the solution of (4) with period $`T>2\pi `$ and $`x(0)>0`$ which has even time-reversal symmetry, $`x_T(t)x_T(t)`$. From this we may construct a 1-site breather solution of system (1) with $`\alpha =0`$, call it $`𝐪_0`$:
$$q_{n,0}(t)=\{\begin{array}{cc}x_T(t)& n=0\\ 0& n0.\end{array}$$
(5)
The MacKay-Aubry theorem establishes the existence of a continuation $`𝐪_\alpha `$ of this solution away from $`\alpha =0`$ in a suitable function space, defined as follows.
###### Definition 1
For any $`n`$, let $`C_{T,n}^+`$ denote the space of $`n`$ times continuously differentiable mappings $``$ which are $`T`$-periodic and have even time reversal symmetry. Note that $`C_{T,n}^+`$ is a Banach space when equipped with the uniform $`C^n`$ norm:
$$|q|_n:=\underset{t}{sup}\{|q(t)|,|\dot{q}(t)|,\mathrm{},|q^{(n)}(t)|\}.$$
###### Definition 2
For any $`n`$,
$$\mathrm{\Omega }_{T,n}^+:=\{𝐪:C_{T,n}^+\text{such that}𝐪_n<\mathrm{}\}$$
where
$$𝐪_n:=\underset{m}{sup}|q_m|_n.$$
Note that $`(\mathrm{\Omega }_{T,n}^+,||||_n)`$ is also a Banach space.
The required function space is $`\mathrm{\Omega }_{T,2}^+`$. By construction, $`𝐪_0\mathrm{\Omega }_{T,2}^+`$ and is exponentially spatially localized.
###### Theorem 3
(MacKay-Aubry) If $`T2\pi `$ there exists $`ϵ>0`$ such that for all $`\alpha [0,ϵ)`$ there is a unique continuous family $`𝐪_\alpha \mathrm{\Omega }_{T,2}^+`$ of solutions of system (1) at coupling $`\alpha `$ with $`𝐪_0`$ as defined in (5). These solutions are exponentially localized in space and the map $`[0,ϵ)\mathrm{\Omega }_{T,2}^+`$ given by $`\alpha 𝐪_\alpha `$ is $`C^1`$.
The idea of the proof is to define a $`C^1`$ mapping $`F:\mathrm{\Omega }_{T,2}^+\mathrm{\Omega }_{T,0}^+`$,
$$F(𝐪,\alpha )_m=\ddot{q}_m\alpha (q_{m+1}2q_m+q_{m1})+V^{}(q_m),$$
(6)
so that $`F(𝐪,\alpha )=0`$ if and only if $`𝐪`$ is an even $`T`$-periodic solution of system (1). In particular, $`F(𝐪_0,0)=0`$ by construction. Using $`T2\pi ^+`$ and anharmonicity of $`V`$, one can show that the partial derivative of $`F`$ with respect to $`𝐪`$ at $`(𝐪_0,0)`$, $`DF_{q_0}:\mathrm{\Omega }_{T,2}^+\mathrm{\Omega }_{T,0}^+`$, is invertible ($`DF_{q_0}`$ is injective with $`(DF_{q_0})^1`$ bounded). Hence the implicit function theorem applies and local existence and uniqueness of the $`C^1`$ family $`𝐪_\alpha `$ satisfying
$$F(𝐪_\alpha ,\alpha )=0$$
(7)
are assured. Persistence of exponential localization is proved as a separate step.
The object of interest in this paper is $`𝐪_0^{}\mathrm{\Omega }_{T,2}^+`$, the tangent vector to the curve $`𝐪_\alpha `$ at $`\alpha =0`$. This may be constructed by implicit differentiation of (7) with respect to $`\alpha `$ at $`\alpha =0`$:
$`DF_{q_0}𝐪_0^{}+{\displaystyle \frac{F}{\alpha }}|_{(𝐪_0,0)}`$ $`=`$ $`0`$
$`𝐪_0^{}`$ $`=`$ $`(DF_{q_0})^1{\displaystyle \frac{F}{\alpha }}|_{(𝐪_0,0)}.`$ (8)
Straightforward calculation shows that
$$[DF_{q_0}𝝌]_m=\{\begin{array}{cc}\ddot{\chi }_m+\chi _m& m0\\ \ddot{\chi }_m+V^{\prime \prime }(x_T)\chi _m& m=0\end{array}$$
(9)
for any $`𝝌\mathrm{\Omega }_{T,2}^+`$, and
$$\left(\frac{F}{\alpha }|_{(𝐪_0,0)}\right)_m=\{\begin{array}{cc}2x_T& m=0\\ x_T& |m|=1\\ 0& |m|>1.\end{array}$$
(10)
So evaluating the right hand side of (8) to find $`𝐪_0^{}=𝝌\mathrm{\Omega }_{T,2}^+`$ is equivalent to solving the following infinite decoupled set of ODEs for $`\{\chi _mC_{T,2}^+:m\}`$:
$`\ddot{\chi }_0+V^{\prime \prime }(x_T)\chi _0`$ $`=`$ $`2x_T`$ (11)
$`\ddot{\chi }_m+\chi _m`$ $`=`$ $`x_T|m|=1`$ (12)
$`\ddot{\chi }_m+\chi _m`$ $`=`$ $`0|m|>1.`$ (13)
Recalling that $`T2\pi ^+`$ (so $`\mathrm{cos}C_{T,2}^+`$) one sees from (13) that $`\chi _m0`$ for $`|m|>1`$, so one need only compute $`\chi _0`$ and $`\chi _1`$ (clearly $`\chi _1\chi _1`$).
In each case we must find the even $`T`$-periodic solution of an inhomogeneous 2nd order linear ODE. For $`m=0,1`$ let $`y_m(t)`$ be the solution of the corresponding homogeneous equation with initial data $`y_m(0)=1`$, $`\dot{y}_m(0)=0`$, so
$`\ddot{y}_0+V^{\prime \prime }(x_T)y_0`$ $`=`$ $`0`$ (14)
$`\ddot{y}_1+y_1`$ $`=`$ $`0.`$ (15)
(Of course $`y_1=\mathrm{cos}`$.) Similarly, let $`z_m(t)`$ denote the particular integral of the inhomogeneous equation with initial data $`z_m(0)=\dot{z}_m(0)=0`$, so
$`\ddot{z}_0+V^{\prime \prime }(x_T)z_0`$ $`=`$ $`2x_T`$ (16)
$`\ddot{z}_1+z_1`$ $`=`$ $`x_T.`$ (17)
Since $`\chi _m(t)`$ is even, it follows that $`\chi _m(t)=\chi _m(0)y_m(t)+z_m(t)`$, and $`\chi _m(0)`$ may be determined by applying either of the $`T`$-periodicity constraints
$`\chi _m(T)=\chi _m(0)`$ $``$ $`\chi _m(0)={\displaystyle \frac{z_m(T)}{1y_m(T)}}`$ (18)
$`\text{or}\dot{\chi }_m(T)=\dot{\chi }_m(0)`$ $``$ $`\chi _m(0)={\displaystyle \frac{\dot{z}_m(T)}{\dot{y}_m(T)}}.`$ (19)
In either case, one sees that the tangent vector $`𝐪_0^{}`$, and in particular its initial value $`𝐪_0^{}(0)=(\mathrm{},0,0,\chi _1(0),\chi _0(0),\chi _1(0),0,0,\mathrm{})`$ can be constructed once the 4 initial value problems (1417) are solved on $`[0,T]`$.
Having computed the initial value of the tangent vector, we may approximate the breather initial profile for small $`\alpha `$ using (3):
$$q_{m,\alpha }(0)=\{\begin{array}{ccccccc}x_T(0)& +& \alpha \chi _0(0)& +& o(\alpha )& & m=0\hfill \\ & & \alpha \chi _1(0)& +& o(\alpha )& & |m|=1\hfill \\ & & & & o(\alpha )& & |m|>1.\hfill \end{array}$$
(20)
So to first order in $`\alpha `$, the continuation leaves all but the central ($`m=0`$) and off-central ($`m=\pm 1`$) sites at the equilibrium position. The qualitative shape of the breather initial profile depends crucially on $`\chi _1(0)`$ (but not $`\chi _0(0)`$). If $`\chi _1(0)>0`$, the continuation displaces the off-central sites from equilibrium in the same direction as the central site. The result is a hump shaped breather in which the central and off-central sites oscillate, roughly speaking, in phase (that is they attain their maxima and minima simultaneously). We shall call such breathers “in-phase breathers (IPBs).” If $`\chi _1(0)<0`$, on the other hand, the off-central sites are displaced in the opposite direction from the central site, resulting in a sombrero shaped initial profile. In this case, the central and off-central sites oscillate, roughly speaking, in anti-phase (the central site attains its maximum when the off-central sites attain their minima, and vice-versa). We shall call such breathers “anti-phase breathers (APBs).” Sievers and Takeno have argued (without giving a rigorous mathematical proof) that breathers of the latter type are supported by certain oscillator chains with no substrate potential and anharmonic nearest neighbour coupling, but their existence in systems of type (1) does not seem to have attracted attention previously in the literature. This is surprising in light of the numerical results described in section 3. We shall find that for a generic substrate potential $`V`$, system (1) supports both IPBs and APBs, the type varying in bands as $`T`$ increases through $`(2\pi ,\mathrm{})`$.
## 3 Numerical results
Calculation of the constants $`\chi _0(0)`$ and $`\chi _1(0)`$ may be performed numerically by solving the initial value problems (14), (16) and (17) using some approximate ODE solver (recall the exact solution of (15) is known). The results presented in this section were generated using a 4th order Runge-Kutta method with fixed time step $`\delta t=0.01`$. Each ODE requires the periodic oscillation $`x_T(t)`$ as input. For most potentials this function must itself be generated by numerical solution of the oscillator equation (4), which might just as well be solved in parallel with (14,16,17), making 4 coupled 2nd order ODEs in all. Since the correct initial displacement $`x_T(0)`$ for a given period $`T`$ is not known in advance, the initial value problem is parametrized by initial displacement, $`T`$ being determined from the approximate solution of (4) by counting sign changes of $`x_T(t)`$ and linear interpolation.
The results obtained depend crucially on whether the substrate potential has reflexion symmetry about the equilibrium position.
### 3.1 Asymmetric potentials
The following asymmetric potentials were investigated:
Morse: $`V_M(x)`$ $`={\displaystyle \frac{1}{2}}(1e^x)^2`$ (21)
Lennard-Jones: $`V_{LJ}(x)`$ $`={\displaystyle \frac{1}{72}}\left({\displaystyle \frac{1}{x^{12}}}{\displaystyle \frac{2}{x^6}}\right)`$ (22)
Cubic: $`V_C(x)`$ $`={\displaystyle \frac{1}{2}}x^2{\displaystyle \frac{1}{3}}x^3.`$ (23)
Note that the equilibrium position for $`V_{LJ}(x)`$ is $`x=1`$ rather than $`x=0`$, as we have been assuming so far. Hence, even though $`V_{LJ}`$ is even, it is still asymmetric, since it is not reflexion symmetric about $`x=1`$. Graphs of $`\chi _0(0)`$ and $`\chi _1(0)`$ against $`T`$ for all these potentials are presented in figures 1 and 2 respectively.
Figure 1: Graphs of $`\chi _0(0)`$ against $`T`$ for various asymmetric substrate potentials.
Figure 1 illustrates a clear difference in the large $`T`$ behaviour of $`\chi _0(0)`$ between $`V_C`$ and the other two potentials $`V_M`$ and $`V_{LJ}`$. Namely $`\chi _0(0)`$ remains bounded as $`T\mathrm{}`$ for $`V_C`$, but grows unbounded for $`V_M`$ and $`V_{LJ}`$. This phenomenon appears to be linked to the large $`T`$ behaviour of the $`T`$ periodic oscillations $`x_T(t)`$. Due to the unstable equilibrium point of $`V_C`$ at $`x=1`$, $`x_T(t)`$ tends to a bounded homoclinic orbit as $`T\mathrm{}`$ for this potential, and $`\chi _0(0)`$ is correspondingly bounded for large $`T`$. For $`V_M`$ and $`V_{LJ}`$ no such bounded homoclinic orbit exists: $`x_T(0)`$ is an unbounded function of $`T`$ and $`\chi _0(0)`$ is correspondingly unbounded as $`T\mathrm{}`$. In all cases, we note that $`\chi _0(0)>0`$, so the continuation to $`\alpha >0`$ displaces the central site further from equilibrium.
Figure 2: Graphs of $`\chi _1(0)`$ against $`T`$ for various asymmetric substrate potentials (solid: Morse; dashed: Lennard-Jones; dotted: cubic).
From figure 2 we see that the sign of $`\chi _1(0)`$ changes as $`T`$ increases, leading to the prediction of $`T`$-bands of IPBs and APBs, as explained in section 2. The sign changes are clearly associated with the vertical asymptotes of the graphs at each $`T2\pi ^+`$ (note that these do not conflict with the contents of section 2 since the existence proof for $`𝐪_\alpha `$, and hence for $`𝐪_0^{}`$, breaks down when $`T2\pi ^+`$). The presence of such asymptotes may be understood by using Green’s function techniques to write down $`z_1(t)`$, the solution of (17):
$$z_1(t)=_0^t\mathrm{sin}(ts)x_T(s)𝑑s.$$
(24)
Combining (24) with (19), and recalling that $`y_1=\mathrm{cos}`$, yields a formula for $`\chi _1(0)`$,
$$\chi _1(0)=\frac{1}{\mathrm{sin}T}_0^T\mathrm{cos}(Tt)x_T(t)𝑑t.$$
(25)
Equation (25) shows that there is a vertical asymptote at $`T=2n\pi `$, with a sign change in $`\chi _1(0)`$, unless
$$_0^{2n\pi }\mathrm{cos}tx_{2n\pi }(t)𝑑t=0,$$
(26)
that is, unless $`x_{2n\pi }(t)`$ has vanishing $`n`$-th Fourier coefficient. It might appear from (25) that there should also be asymptotes at $`T=(2n+1)\pi `$. Of course, this cannot be true since standard results on continuity of solutions of ODEs with respect to initial data imply that $`\chi _1(0)`$ is continuous for $`T2\pi ^+`$. In fact, a simple argument using periodicity and evenness of $`x_T`$ demonstrates that
$$_0^{(2n+1)\pi }\mathrm{cos}tx_{(2n+1)\pi }(t)𝑑t0$$
(27)
for all $`n^+`$ and $`V`$. By contrast, we will prove in section 4 that (26) almost never holds (in a sense which will be made precise), so that the resonant periods $`T=2n\pi `$ generically separate IPB bands from APB bands.
### 3.2 Symmetric potentials
The following symmetric potentials were investigated:
Frenkel-Kontorova: $`V_{FK}(x)`$ $`=1\mathrm{cos}x`$ (28)
Quartic: $`V_Q(x)`$ $`={\displaystyle \frac{1}{2}}x^2{\displaystyle \frac{1}{4}}x^4`$ (29)
Gaussian: $`V_G(x)`$ $`=1e^{x^2/2}`$ (30)
The graphs in figure 3 show $`\chi _0(0)`$ against $`T`$ and closely resemble those of the previous section: $`\chi _0(0)`$ is bounded as $`T\mathrm{}`$ if $`x_T`$ tends to a bounded homoclinic orbit ($`V_{FK}`$ and $`V_Q`$), and unbounded if $`x_T(0)\mathrm{}`$ as $`T\mathrm{}`$ ($`V_G`$).
Figure 3: Graphs of $`\chi _0(0)`$ against $`T`$ for various symmetric substrate potentials.
On the other hand, the graphs of $`\chi _1(0)`$ against $`T`$ (see figure 4) are strikingly different from the asymmetric case. In each case, vertical asymptotes are present at $`T=2n\pi `$ only if $`n`$ is an even positive integer. If $`n`$ is odd, no asymptote is present. To explain this, note that evenness of $`V`$ implies that $`x_T(t)`$ is $`T/2`$ antiperiodic, that is $`x_T(tT/2)x_T(t)`$, which in turn guarantees that equation (26) holds whenever $`n`$ is odd. The limit $`lim_{T2n\pi }\chi _1(0)`$ ($`n`$ odd) exists and may, in principle, be computed using L’Hospital’s rule with
$$\frac{d}{dT}z_1(T)=x_T(T)+_0^T[\mathrm{sin}(Tt)x_T(t)+\mathrm{cos}(Tt)f^{}(T)y_0(t)]𝑑t,$$
(31)
where $`y_0`$ is the solution of (14) as before and $`f(T):=x_T(0)`$.
Figure 4: Graphs of $`\chi _1(0)`$ against $`T`$ for various symmetric substrate potentials (solid: Frenkel-Kontorova; dashed: quartic; dotted: Gaussian).
### 3.3 Almost symmetric potentials
For $`m=1`$ equation (18) becomes
$$\chi _1(0)=\frac{z_1(T)}{1\mathrm{cos}T}.$$
(32)
From this equation it is clear that the behaviour of $`\chi _1(0)`$ as a function of $`T`$ is largely determined by the zeroes of $`z_1(T)`$. In particular, the sign of $`z_1(T)`$ determines the sign of $`\chi _1(0)`$. Figure 5 shows typical zero distributions for both asymmetric and symmetric potentials. It is interesting to note that no infinitesimal perturbation of the symmetric type distribution yields the asymmetric type distribution. Hence the dependence of $`\chi _1(0)`$ on $`T`$ for an almost symmetric potential must differ qualitatively from the typical asymmetric behaviour observed in section 3.1.
Figure 5: Typical zero distributions of $`z_1(T)`$ for asymmetric and symmetric substrates.
This observation motivates us to consider $`z_1(T)`$ for the following 1-parameter family of potentials:
$$V_ϵ(x)=\frac{1}{2}x^2\frac{1}{4}x^4\frac{ϵ}{5}x^5.$$
(33)
The zeroes of $`z_1(T)`$ in a region of the $`(T,\mathrm{log}ϵ)`$ plane are plotted in figure 6. In the symmetric limit ($`\mathrm{log}ϵ\mathrm{}`$), the typical symmetric type distribution is recovered, as expected. For a fixed but very small $`ϵ>0`$ (e.g. $`\mathrm{log}ϵ=30`$), the small $`T`$ oscillations $`x_T(t)`$ are pointwise essentially unchanged from the $`ϵ=0`$ case, and the distribution for small $`T`$ is numerically indistinguishable from the symmetric case. Only for large $`T`$ does the oscillation $`x_T(t)`$ detect the asymmetry of $`V_ϵ`$, and a transition from symmetric to asymmetric behaviour occurs in the distribution. As $`ϵ`$ is increased, the period at which this transition occurs becomes smaller and smaller, so that the transition “propagates” leftwards as seen in figure 6. For sufficiently large $`ϵ`$ (e.g. $`\mathrm{log}ϵ>3`$) the transition is complete and the typical asymmetric distribution is recovered. These observations appear to be essentially independent of the asymmetric perturbation considered.
Figure 6: Zero distribution of $`z_1(T)`$ for the one parameter family of almost symmetric substrates $`V_ϵ(x)=\frac{1}{2}x^2\frac{1}{4}x^4\frac{ϵ}{5}x^5`$.
## 4 Genericity of numerical results
The aim of this section is to explain some of the generic features observed in the numerical results of section 3. In particular we will consider the $`T2\pi ^+`$ behaviour of $`\chi _1(0)`$, and the generic presence of asymptotes at $`T2\pi ^+`$.
### 4.1 The small period limit
The numerical data suggest that $`\chi _1(0)+\mathrm{}`$ as $`T2\pi ^+`$ for every potential $`V`$. We may confirm this prediction, at least in the case where $`V`$ is analytic (as is the case in all our examples), by means of a Linstedt expansion for $`x_T(t)`$ . We seek to construct a solution of the oscillator equation (4) with initial data $`x(0)=ϵ>0`$, $`\dot{x}(0)=0`$, by power series expansion in $`ϵ`$. Such a solution will be periodic in $`t`$ with unknown period $`T(ϵ)=2\pi /\omega (ϵ)`$, where $`T(ϵ)`$ and hence $`\omega (ϵ)`$ depends analytically on $`ϵ`$. Defining a rescaled time variable $`s:=\omega (ϵ)t`$ (so the solution is $`s`$-periodic with period $`2\pi `$) and dependent variable $`y(s):=x(t)/ϵ`$ (so $`y(0)=1`$, $`y^{}(0)=0`$ for all $`ϵ`$), we substitute the expansions
$`\omega (ϵ)`$ $`=`$ $`a_0+a_1ϵ+a_2ϵ^2+\mathrm{}`$ (34)
$`y(s)`$ $`=`$ $`y_0(s)+ϵy_1(s)+ϵ^2y_2(s)+\mathrm{}`$ (35)
into (4) and Taylor expand $`V(x)`$ about $`x=0`$. Grouping terms by order in $`ϵ`$, this results in an infinite set of coupled ODEs for $`\{y_n:n\}`$, to be solved in turn ($`n=0,1,2,\mathrm{}`$) with initial data $`y_0(0)=1`$, $`\dot{y}_0(0)=y_n(0)=\dot{y}_n(0)=0`$. Demanding that each $`y_n(s)`$ have period $`2\pi `$ then fixes the constants $`a_0,a_1,a_2,\mathrm{}`$. For our purposes, it is sufficient to proceed only up to order $`ϵ^2`$, whereupon one finds that
$`\omega (ϵ)`$ $`=`$ $`1Aϵ^2+O(ϵ^3)`$
$`T(ϵ)`$ $`=`$ $`2\pi [1+Aϵ^2+O(ϵ^3)],`$ (36)
where $`A=V^{(4)}(0)/16(V^{(3)}(0))^2/12`$ is positive for each of our examples (since they are all softly anharmonic, i.e. $`T>2\pi `$). Then
$$x_{T(ϵ)}(t)=ϵ\mathrm{cos}\omega (ϵ)t+ϵ^2y_1(\omega (ϵ)t)+ϵ^3y_2(\omega (ϵ)t)+\mathrm{}$$
(37)
where all we need know about $`y_1,y_2,\mathrm{}`$ is that they are even with period $`2\pi `$. We can now extract the asymptotic form of $`z_1(T(ϵ))`$ for small $`ϵ`$:
$`z_1(T(ϵ))`$ $`=`$ $`{\displaystyle _0^{T(ϵ)}}\mathrm{sin}(T(ϵ)t)x_{T(ϵ)}(t)𝑑t`$ (38)
$`=`$ $`{\displaystyle \frac{1}{\omega (ϵ)}}{\displaystyle _0^{2\pi }}\mathrm{sin}\left({\displaystyle \frac{2\pi s}{\omega (ϵ)}}\right)(ϵ\mathrm{cos}s+ϵ^2y_1(s)+\mathrm{})𝑑s`$
$`=`$ $`A\pi ^2ϵ^3+O(ϵ^4).`$
Hence from equation (32),
$$\chi _1(0)=\frac{z_1(T(ϵ))}{1\mathrm{cos}T(ϵ)}=\frac{1}{2\pi Aϵ}+O(ϵ^0).$$
(39)
As $`ϵ0^+`$, equations (36) and (39) reproduce the $`T2\pi ^+`$ asymptotic behaviour observed in all the numerical experiments.
### 4.2 Vertical asymptotes
We have seen that a sign change in $`\chi _1(0)`$, and hence a transition from IPBs to APBs, must occur at $`T=2n\pi `$ ($`n^+`$) unless $`z_1(2n\pi )=0`$, equation (26). In this section we will prove that, for a given $`n`$, the set of potentials on which (26) holds is negligibly small, so that the presence of asymptotes in the graphs of $`\chi _1(0)`$ against $`T`$ really is generic. To be precise, we will show that, for potentials in a certain Banach space $`𝒫`$, the subset of potentials on which (26) holds is locally a codimension 1 submanifold.
###### Definition 4
Let $`𝒫=\{V:\text{such that }V\text{ is }C^2\text{ and }V_2<\mathrm{}\}`$ and $`F:C_{T,2}^+𝒫C_{T,0}^+`$ be the mapping
$$F(q,V)=\ddot{q}+V^{}(q).$$
Note that $`F`$ is a $`C^1`$ mapping between Banach spaces.
###### Definition 5
A potential $`V𝒫`$ is anharmonic at $`qC_{T,2}^+`$ if $`q`$ is nonconstant, $`F(q,V)=0`$, and $`\mathrm{ker}DF_{(q,V)}=\{0\}`$, where $`DF_{(q,V)}:C_{T,2}^+C_{T,0}^+`$ denotes the partial derivative of $`F`$ at $`(q,V)`$.
Clearly $`F(q,V)=0`$ means that $`q`$ is a solution of Newton’s equation for motion in potential $`V`$. That injectivity of $`DF_{(q,V)}`$ is equivalent to the standard definition of anharmonicity in Hamiltonian mechanics is shown, for example, in the original proof of Theorem 3 . For all the potentials we considered, $`V`$ is anharmonic at $`q=x_TC_{T,2}^+`$ for all $`T>2\pi `$.
###### Definition 6
A perturbation neighbourhood of $`(q,V)C_{T,2}^+𝒫`$ is a pair $`(𝒰,f)`$ where $`𝒰𝒫`$ is an open set containing $`V`$, and $`f:𝒰C_{T,2}^+`$ is a $`C^1`$ map satisfying (a) every $`W𝒰`$ is anharmonic at $`f(W)`$ and (b) $`f(V)=q`$. Note that $`𝒰`$ is trivially a Banach manifold.
###### Lemma 7
For any $`V𝒫`$ anharmonic at $`qC_{T,2}^+`$ there exists a perturbation neighbourhood $`(𝒰,f)`$ of $`(q,V)`$. The function $`f`$ is unique on sufficiently small $`𝒰`$.
Proof: Anharmonicity implies that $`DF_{(q,V)}`$ is injective, and hence the standard solvability criterion in linear ODE theory guarantees it is also surjective (see section 3.3.2 of for details). The open mapping theorem then ensures boundedness of $`DF_{(q,V)}^1`$. Hence $`DF_{(q,V)}`$ is invertible, and the implicit function theorem applied to $`F`$ ensures existence of $`(𝒰,f)`$ and local uniqueness of $`f`$. $`\mathrm{}`$
###### Definition 8
For $`T2\pi ^+`$, we shall say that $`V𝒫`$ is degenerate at $`qC_{T,2}^+`$ if $`V`$ is anharmonic at $`q`$ and
$$_0^T\mathrm{cos}tq(t)𝑑t=0.$$
###### Theorem 9
Let $`T2\pi ^+`$, $`V_0𝒫`$ be degenerate at $`q_0C_{T,2}^+`$, and $`(𝒰,f)`$ be a perturbation neighbourhood of $`(q_0,V_0)`$. The subset of $`𝒰`$ on which degeneracy persists is a codimension 1 submanifold.
Proof: Consider the $`C^1`$ mapping $`I:𝒰`$ defined by
$$I(V)=_0^T\mathrm{cos}t[f(V)](t)𝑑t$$
($`I`$ is $`C^1`$ since $`f`$ is $`C^1`$). By the Regular Value Theorem (see appendix) the result follows if we establish that $`0`$ is a regular value of $`I`$, or in other words that $`I`$ is a submersion at every $`VI^1(0)`$. Since ker $`DI_V`$ is of finite codimension in $`T_V𝒰=𝒫`$ it splits and hence it suffices to show that $`DI_V:𝒫`$ is surjective for all $`VI^1(0)`$.
Let $`VI^1(0)`$, $`q=f(V)`$. For any $`\delta V𝒫`$, let $`\delta q_{\delta V}=[DF_{(q,V)}^1(\delta V^{}q)]C_{T,2}^+`$. In other words $`\delta q_{\delta V}`$ is the unique even, $`T`$-periodic $`C^2`$ solution of
$$\ddot{\delta q}_{\delta V}+V^{\prime \prime }(q(t))\delta q_{\delta V}=\delta V^{}(q(t)),$$
whose existence follows from anharmonicity of $`V`$ at $`q`$ (invertibility of $`DF_{(q,V)}`$ on $`C_{T,0}^+`$). With this notation,
$$DI_V(\delta V)=_0^T\mathrm{cos}t\delta q_{\delta V}(t)𝑑t$$
and surjectivity of $`DI_V`$ will follow if we exhibit $`\delta V𝒫`$ such that $`DI_V(\delta V)0`$. We construct such a $`\delta V`$ in two steps.
For some $`t_0(0,T/2)`$ and $`ϵ>0`$, let $`bC_{T,2}^+`$ be a non-negative function with $`\text{supp}b[0,T/2]=[t_0ϵ,t_0+ϵ]`$. Clearly
$$_0^Tb(t)\mathrm{cos}tdt=2_0^{T/2}b(t)\mathrm{cos}tdt0$$
(where $`T2\pi ^+`$ has been used) provided we choose $`t_0\mathrm{cos}^1(0)`$ and $`ϵ`$ sufficiently small. Since the critical points of $`q(t)`$ are isolated, we may also assume that $`q`$ is invertible on $`[t_0ϵ,t_0+ϵ]`$ with $`C^2`$ inverse $`q^1:[x_1,x_2][t_0ϵ,t_0+ϵ]`$. By its construction the function $`g`$ defined by
$$g(x)=\{\begin{array}{cc}[[DF_{(q,V)}b]q^1](x)& x[x_1,x_2]\\ 0& x[x_1,x_2]\end{array}$$
is in $`C^0()`$ (since $`DF`$ maps to $`C_{T,0}^+`$). This allows us to construct a function $`\delta \stackrel{~}{V}(x)C^1()`$ such that $`\delta q_{\delta \stackrel{~}{V}}=b`$ as follows. For $`x\mathrm{ran}(q)`$, $`\delta \stackrel{~}{V}(x)=_0^xg(z)𝑑z`$ defines a $`C^1`$ function on ran($`q`$). Since $`\delta q_{\delta \stackrel{~}{V}}`$ is clearly independent of the behaviour of $`\delta \stackrel{~}{V}`$ outside $`\mathrm{ran}(q)`$, we extend $`\delta \stackrel{~}{V}`$ to a $`C^1`$ function on $``$ with support contained in some compact set $`K\mathrm{ran}(q)`$.
We cannot immediately conclude that $`DI_V`$ is surjective, since $`\delta \stackrel{~}{V}`$ is $`C^1`$ but not necessarily $`C^2`$. However, using a density argument it is straightforward to find a $`\delta VC^2`$ close enough to $`\delta \stackrel{~}{V}`$ such that $`DI_V(\delta V)0`$ stills holds. More precisely, $`DI_V`$ has a natural continuous extension, call it $`\overline{DI_V}`$ to $`C^1`$. By construction $`\delta \stackrel{~}{V}`$ satisfies $`\overline{DI_V}(\delta \stackrel{~}{V})0`$ and belongs to $`C_K^1`$ (the $`C^1`$ functions with support contained in $`K`$). Since $`C_K^2`$ is dense in $`C_K^1`$, there exists a sequence $`\delta V_mC_K^2𝒫`$ converging (in $`C^1`$ norm) to $`\delta \stackrel{~}{V}`$, and by continuity $`\overline{DI_V}(\delta V_m)0`$ for all $`m`$ sufficiently large. $`\mathrm{}`$
In the above, we have chosen $`𝒫=(C^2(),||||_2)`$ as our space of potentials. This choice was made for the sake of clarity and notational simplicity – many other choices would work. In fact, all but two ($`V_{FK}`$ and $`V_G`$) of the example potentials considered in section 3 lie, strictly speaking, outside $`𝒫`$, since they are unbounded. However, the analysis above can easily be adapted to deal with this: one simply replaces $`𝒫`$ by the affine space $`𝒫_V:=\{W:WV_2<\mathrm{}\}`$.
## 5 Concluding remarks
In this paper, breather initial profiles in the weak coupling regime of a simple class of oscillator networks have been examined, focusing on the dependence on breather period $`T`$. The direction of continuation of one-site breathers was determined numerically using a simple and inexpensive numerical scheme. Two types of breather were identified, called IPBs and APBs. The numerical data suggest that generically these two types occur in alternating bands in the $`T`$ parameter space ($`T(2\pi ,\mathrm{})`$), and that the resonant periods $`T2\pi ^+`$ separate an IPB band from an APB band. The genericity of this behaviour was proved rigorously.
The distinction between IPBs and APBs, which does not appear to have attracted much attention in the literature, may have phenomenological implications in applications of the model (1). One reason for interest in discrete breathers (particularly continued one-site breathers) is that, because of their strong spatial localization, they typically require little energy to achieve large amplitude oscillations close to the centre. This makes them good candidates for “seeds” of mechanical breakdown of the network, the idea being that the central oscillation becomes so violent as to break the chain (a similar mechanism is postulated as a mechanism for DNA denaturation, for example ). Since APBs stress the central intersite springs more than IPBs they presumably make more effective seeds of mechanical breakdown.
## Acknowledgments
This work was partially completed during a visit by MH to the Max-Planck-Institut für Mathematik in den Naturwissenschaften, Leipzig, where JMS was a guest scientist. Both authors wish to thank Prof. Eberhard Zeidler for the generous hospitality of the institute. JMS is an EPSRC Postdoctoral Research Fellow in Mathematics.
## Appendix
We recall some basic definitions and an elementary result of infinite dimensional differential topology, the Regular Value Theorem, which gives conditions under which the level set of a smooth function is guaranteed to be a manifold.
###### Definition 10
We say that a closed subspace $`S`$ of a complete topological vector space $`B`$ splits if there exists another closed subspace $`C`$ which is complementary to $`S`$, i.e. $`C+S=B`$ and $`CS=(0)`$.
Note that if $`B`$ is a Hilbert space then any closed subspace $`S`$ splits since we can take the complement to be $`C=S^{}`$. Also if $`B`$ is finite dimensional then since any subspace is closed all subspaces split. More generally any finite dimensional subspace (necessarily closed) of a Banach space splits, as does any closed subspace of finite codimension.
###### Definition 11
A $`C^1`$ map $`f:XY`$ between Banach manifolds is a submersion at $`xX`$ if $`Df_x:T_xXT_{f(x)}Y`$ is surjective and ker $`Df_x`$ splits. A value $`yY`$ is a regular value of $`f`$ if $`f`$ is a submersion for every $`xf^1(y)`$.
Note that in the case that $`Y`$ is finite dimensional then ker $`Df_x`$ is of finite codimension in $`T_xX`$ and hence is guaranteed to split. In this case we need only verify that $`Df_x`$ is surjective for $`f`$ to be a submersion at $`x`$.
###### Theorem 12
(The Regular Value Theorem) Let $`f:XY`$ be a $`C^1`$ map between Banach manifolds and $`yY`$ be a regular value of $`f`$. Then $`f^1(y)`$ is a submanifold of $`X`$ with $`T_xf^1(y)\mathrm{ker}Df_x`$.
The proof follows from working in charts around $`x`$ and $`f(x)`$ and applying the Implicit Function Theorem. For details of a proof see .
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# 1 Introduction
## 1 Introduction
After the discovery that gravity on the brane may be localized there was renewed interest in the studies of higher-dimensional (brane-world) theories. In particular, numerous works (and refs. therein) have been devoted to the investigation of cosmology (inflation) of brane-worlds. In refs. it has been suggested the inflationary brane-world scenario realized due to quantum effects of brane matter. Such scenario is based on large $`N`$ quantum CFT living on the brane . Actually, that corresponds to implementing of RS compactification within the context of renormalization group flow in AdS/CFT set-up. Note that working within large $`N`$ approximation justifies such approach to brane-world quantum cosmology as then quantum matter loops contribution is essential.
Another important aspect of brane-world Universe with localized gravity is related with the possibility to resolve the cosmological constant problem. For example, it has been shown in ref. that in the presence of bulk scalar (dilaton) one can find static solutions of equations of motion where the bulk dilatonic potential vanishes. Such self-tuning mechanism has been further studied in ref.. Unfortunately, it is usual that such solutions which localize gravity have a naked space-time singularity. (The presence of non-trivial dilatonic potential makes the situation even more complicated ). General properties of the self-tuning domain wall solutions of 5d gravity-scalar system with various potentials and brane couplings have been discussed in ref.. It has been shown there that for some specific potential the resolution of singularities (when potential and brane coupling are fine-tuned) may be achieved. The bulk spacetime is asymptotically AdS and gravity localization may occur without having singularities! However, in the studies in this direction the discussion has been done so far mainly for solutions with flat 4d domain walls (flat branes). (The explicit, non-singular example of ref. corresponds to such flat brane configuration). The reason is that in this case the second-order equations of motion may be reduced to first-order form for an arbitrary dilatonic potential, see ref. for explicit examples. In the case when the branes are not flat this procedure does not work directly, generally speaking. <sup>4</sup><sup>4</sup>4Note that brane-world theory may be often understood as completely 4d two measure theory. Such models with exponential potentials which appear due to scale invariance have been intensively studied in refs.(inflation and role of vacuum effects).
The purpose of the present work is to investigate the role of quantum matter living on the brane in the study of brane-world cosmology in 5d AdS dilatonic gravity with non-trivial dilatonic potential (bosonic sector of the corresponding gauged supergravity). We are mainly interested in the situation when the boundary of 5d AdS space represents a 4d constant curvature space whose creation (as is shown) is possible only due to quantum effects of brane matter. Thus, the possibility of dilatonic brane-world inflation induced by quantum effects is proved. In different versions of such scenario discussed here the dynamical determination of dilaton occurs as well.
The paper is organized as follows. In the next section we investigate dilatonic brane-world inflation induced by quantum effects (using anomaly induced effective action) in the situation with constant bulk potential. The bulk space represents (singular) asymptotically AdS background with non-trivial dilaton. The brane matter quantum effects (maximally SUSY Yang-Mills theory is considered as brane CFT) help to create de Sitter or AdS space on the brane. Hence, the quantum realization of brane-world inflation is possible in the presence of the dilaton which is determined dynamically in the bulk as well as on the brane. Note that an analytical treatment is done in this section. Section 3 is devoted to extension of results of previous consideration for non-constant bulk potentials. One solvable example of bulk equations of motion for exponential potential is given. In this case de Sitter (or hyperbolic) brane with small radius occurs for SUSY Yang-Mills theory. It is also interesting that without quantum corrections ($`W`$ vanishes) the dilatonic hyperbolic brane is still possible.
The conditions to get non-singular, asymptotically AdS dilatonic spacetime (when 4d gravity is trapped) are discussed. An example of a toy dilatonic potential is presented. It is shown (in some approximation) that due to quantum effects the brane represents de Sitter space and localization of gravity occurs. Hence, the role of such (fine-tuned) dilatonic potential is to make weaker (or to avoid completely) the singularity which appears for the AdS bulk solution of section 2. Dilaton is still determined dynamically. In section 4 we show that all the above picture may be well realized in the situation when brane matter is not exactly conformal invariant matter (dilaton coupled spinors are considered). Similar qualitative results as in previous sections are obtained. Some resume and perspectives are drawn in the Discussion. In the Appendix a short discussion of some equivalence between 5d dilatonic gravity and 4d dilatonic gravity coupled with CFT is done.
## 2 Dilatonic brane-world inflation induced by quantum effects: Constant bulk potential
We start with Euclidean signature for the action $`S`$ which is the sum of the Einstein-Hilbert action $`S_{\mathrm{EH}}`$ with kinetic term for dilaton $`\varphi `$, the Gibbons-Hawking surface term $`S_{\mathrm{GH}}`$, the surface counter term $`S_1`$ and the trace anomaly induced action $`W`$<sup>5</sup><sup>5</sup>5For the introduction to anomaly induced effective action in curved space-time (with torsion), see section 5.5 in .:
$`S`$ $`=`$ $`S_{\mathrm{EH}}+S_{\mathrm{GH}}+2S_1+W,`$ (1)
$`S_{\mathrm{EH}}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi G}}{\displaystyle d^5x\sqrt{g_{(5)}}\left(R_{(5)}\frac{1}{2}_\mu \varphi ^\mu \varphi +\frac{12}{l^2}\right)},`$ (2)
$`S_{\mathrm{GH}}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi G}}{\displaystyle d^4x\sqrt{g_{(4)}}_\mu n^\mu },`$ (3)
$`S_1`$ $`=`$ $`{\displaystyle \frac{3}{8\pi Gl}}{\displaystyle d^4x\sqrt{g_{(4)}}},`$ (4)
$`W`$ $`=`$ $`b{\displaystyle d^4x\sqrt{\stackrel{~}{g}}\stackrel{~}{F}A}`$ (5)
$`+b^{}{\displaystyle }d^4x\{A[2\stackrel{~}{\mathrm{}}^2+\stackrel{~}{R}_{\mu \nu }\stackrel{~}{}_\mu \stackrel{~}{}_\nu {\displaystyle \frac{4}{3}}\stackrel{~}{R}\stackrel{~}{\mathrm{}}^2+{\displaystyle \frac{2}{3}}(\stackrel{~}{}^\mu \stackrel{~}{R})\stackrel{~}{}_\mu ]A`$
$`+(\stackrel{~}{G}{\displaystyle \frac{2}{3}}\stackrel{~}{\mathrm{}}\stackrel{~}{R})A\}`$
$`{\displaystyle \frac{1}{12}}\left\{b^{\prime \prime }+{\displaystyle \frac{2}{3}}(b+b^{})\right\}{\displaystyle d^4x\left[\stackrel{~}{R}6\stackrel{~}{\mathrm{}}A6(\stackrel{~}{}_\mu A)(\stackrel{~}{}^\mu A)\right]^2}`$
$`+C{\displaystyle d^4xA\varphi \left[\stackrel{~}{\mathrm{}}^2+2\stackrel{~}{R}_{\mu \nu }\stackrel{~}{}_\mu \stackrel{~}{}_\nu \frac{2}{3}\stackrel{~}{R}\stackrel{~}{\mathrm{}}^2+\frac{1}{3}(\stackrel{~}{}^\mu \stackrel{~}{R})\stackrel{~}{}_\mu \right]\varphi }.`$
Here the quantities in the 5 dimensional bulk spacetime are specified by the suffices <sub>(5)</sub> and those in the boundary 4 dimensional spacetime are specified by <sub>(4)</sub>. The factor $`2`$ in front of $`S_1`$ in (1) is coming from that we have two bulk regions which are connected with each other by the brane. In (3), $`n^\mu `$ is the unit vector normal to the boundary. In (5), one chooses the 4 dimensional boundary metric as
$$g_{(4)}^{}{}_{\mu \nu }{}^{}=\mathrm{e}^{2A}\stackrel{~}{g}_{\mu \nu },$$
(6)
and we specify the quantities given by $`\stackrel{~}{g}_{\mu \nu }`$ by using $`\stackrel{~}{}`$. $`G`$ ($`\stackrel{~}{G}`$) and $`F`$ ($`\stackrel{~}{F}`$) are the Gauss-Bonnet invariant and the square of the Weyl tensor, which are given as
$`G`$ $`=`$ $`R^24R_{ij}R^{ij}+R_{ijkl}R^{ijkl},`$
$`F`$ $`=`$ $`{\displaystyle \frac{1}{3}}R^22R_{ij}R^{ij}+R_{ijkl}R^{ijkl},`$ (7)
<sup>6</sup><sup>6</sup>6We use the following curvature conventions: $`R`$ $`=`$ $`g^{\mu \nu }R_{\mu \nu }`$ $`R_{\mu \nu }`$ $`=`$ $`R_{\mu \lambda \nu }^\lambda `$ $`R_{\mu \rho \nu }^\lambda `$ $`=`$ $`\mathrm{\Gamma }_{\mu \rho ,\nu }^\lambda +\mathrm{\Gamma }_{\mu \nu ,\rho }^\lambda \mathrm{\Gamma }_{\mu \rho }^\eta \mathrm{\Gamma }_{\nu \eta }^\lambda +\mathrm{\Gamma }_{\mu \nu }^\eta \mathrm{\Gamma }_{\rho \eta }^\lambda `$ $`\mathrm{\Gamma }_{\mu \lambda }^\eta `$ $`=`$ $`{\displaystyle \frac{1}{2}}g^{\eta \nu }\left(g_{\mu \nu ,\lambda }+g_{\lambda \nu ,\mu }g_{\mu \lambda ,\nu }\right).`$
In the effective action (5), we now consider the case corresponding to $`𝒩=4`$ $`SU(N)`$ Yang-Mills theory, where
$$b=b^{}=\frac{C}{4}=\frac{N^21}{4(4\pi )^2}.$$
(8)
The dilaton field $`\varphi `$ which appears from the coupling with extended conformal supergravity is in general complex but we consider the case in which that only the real part of $`\varphi `$ is non-zero. Adopting AdS/CFT correspondence one can argue that in symmetric phase the quantum brane matter appears due to maximally SUSY Yang-Mills theory as above. Note that there is a kinetic term for the dilaton in the classical bulk action but also there is dilatonic contribution to the anomaly induced effective action $`W`$. Here, it appears the difference with the correspondent construction in ref. where there was no dilaton.
In the bulk, the solution of the equations of motion is given in , as follows
$`ds^2`$ $`=`$ $`f(y)dy^2+y{\displaystyle \underset{i,j=0}{\overset{d1}{}}}\widehat{g}_{ij}(x^k)dx^idx^j`$
$`f`$ $`=`$ $`{\displaystyle \frac{d(d1)}{4y^2\lambda ^2\left(1+\frac{c^2}{2\lambda ^2y^d}+\frac{kd}{\lambda ^2y}\right)}}`$
$`\varphi `$ $`=`$ $`c{\displaystyle 𝑑y\sqrt{\frac{d(d1)}{4y^{d+2}\lambda ^2\left(1+\frac{c^2}{2\lambda ^2y^d}+\frac{kd}{\lambda ^2y}\right)}}}.`$ (9)
Here $`\lambda ^2=\frac{12}{l^2}`$ and $`\widehat{g}_{ij}`$ is the metric of the Einstein manifold, which is defined by $`r_{ij}=k\widehat{g}_{ij}`$, where $`r_{ij}`$ is the Ricci tensor constructed with $`\widehat{g}_{ij}`$ and $`k`$ is a constant. We should note that there is a curvature singularity at $`y=0`$ . The solution with non-trivial dilaton would presumbly correspond to the deformation of the vacuum (which is associated with the dimension 4 operator, say $`\mathrm{tr}F^2`$) in the dual maximally SUSY Yang-Mills theory.
If one defines a new coordinate $`z`$ by
$$z=𝑑y\sqrt{\frac{d(d1)}{4y^2\lambda ^2\left(1+\frac{c^2}{2\lambda ^2y^d}+\frac{kd}{\lambda ^2y}\right)}},$$
(10)
and solves $`y`$ with respect to $`z`$, we obtain the warp factor $`\mathrm{e}^{2\widehat{A}(z,k)}=y(z)`$. Here one assumes the metric of 5 dimensional space time as follows:
$$ds^2=dz^2+\mathrm{e}^{2A(z,\sigma )}\stackrel{~}{g}_{\mu \nu }dx^\mu dx^\nu ,\stackrel{~}{g}_{\mu \nu }dx^\mu dx^\nu l^2\left(d\sigma ^2+d\mathrm{\Omega }_3^2\right).$$
(11)
Here $`d\mathrm{\Omega }_3^2`$ corresponds to the metric of 3 dimensional unit sphere. Then for the unit sphere ($`k=3`$), we find
$$A(z,\sigma )=\widehat{A}(z,k=3)\mathrm{ln}\mathrm{cosh}\sigma ,$$
(12)
for the flat Euclidean space ($`k=0`$)
$$A(z,\sigma )=\widehat{A}(z,k=0)+\sigma ,$$
(13)
and for the unit hyperboloid ($`k=3`$)
$$A(z,\sigma )=\widehat{A}(z,k=3)\mathrm{ln}\mathrm{sinh}\sigma .$$
(14)
We now identify $`A`$ and $`\stackrel{~}{g}`$ in (11) with those in (6). Then we find $`\stackrel{~}{F}=\stackrel{~}{G}=0`$, $`\stackrel{~}{R}=\frac{6}{l^2}`$ etc.
According to the assumption in (11), the actions in (2), (3), (4), and (5) have the following forms:
$`S_{\mathrm{EH}}`$ $`=`$ $`{\displaystyle \frac{l^4V_3}{16\pi G}}{\displaystyle }dzd\sigma \{(8_z^2A20(_zA)^2)\mathrm{e}^{4A}`$ (15)
$`+\left(6_\sigma ^2A6(_\sigma A)^2+6\right)\mathrm{e}^{2A}`$
$`{\displaystyle \frac{1}{2}}\mathrm{e}^{4A}(_z\varphi )^2{\displaystyle \frac{1}{2l^2}}\mathrm{e}^{2A}(_\sigma \varphi )^2+{\displaystyle \frac{12}{l^2}}\mathrm{e}^{4A}\},`$
$`S_{\mathrm{GH}}`$ $`=`$ $`{\displaystyle \frac{3l^4V_3}{8\pi G}}{\displaystyle 𝑑\sigma \mathrm{e}^{4A}_zA},`$ (16)
$`S_1`$ $`=`$ $`{\displaystyle \frac{3l^3V_3}{8\pi G}}{\displaystyle 𝑑\sigma \mathrm{e}^{4A}},`$ (17)
$`W`$ $`=`$ $`V_3{\displaystyle }d\sigma [b^{}A(2_\sigma ^4A8_\sigma ^2A)`$ (18)
$`2(b+b^{})\left(1_\sigma ^2A(_\sigma A)^2\right)^2`$
$`+CA\varphi (_\sigma ^4\varphi 4_\sigma ^2\varphi )].`$
Here $`V_3`$ is the volume or area of the unit 3 sphere:
$$V_3=2\pi ^2.$$
(19)
On the brane at the boundary, one gets the following equations
$`0`$ $`=`$ $`{\displaystyle \frac{48l^4}{16\pi G}}\left(_zA{\displaystyle \frac{1}{l}}\right)\mathrm{e}^{4A}+b^{}\left(4_\sigma ^4A16_\sigma ^2A\right)`$ (20)
$`4(b+b^{})\left(_\sigma ^4A+2_\sigma ^2A6(_\sigma A)^2_\sigma ^2A\right)`$
$`+2C\left(_\sigma ^4\varphi 4_\sigma ^2\varphi \right),`$
from the variation over $`A`$ and
$$0=\frac{l^4}{8\pi G}\mathrm{e}^{4A}_z\varphi +C\left\{A\left(_\sigma ^4\varphi 4_\sigma ^2\varphi \right)+_\sigma ^4(A\varphi )4_\sigma ^2(A\varphi )\right\},$$
(21)
from the variation over $`\varphi `$. We should note that the contributions from $`S_{\mathrm{EH}}`$ and $`S_{\mathrm{GH}}`$ are twice from the naive values since we have two bulk regions which are connected with each other by the brane. The equations (20) and (21) do not depend on $`k`$, that is, they are correct for any of the sphere, hyperboloid, or flat Euclidean space. The $`k`$ dependence appears when the bulk solutions are substituted. Substituting the bulk solution given by (2), (10) and (12), (13) or (14) into (20) and (21), one obtains
$`0`$ $`=`$ $`{\displaystyle \frac{1}{\pi Gl}}\left(\sqrt{1+{\displaystyle \frac{kl^2}{3y_0}}+{\displaystyle \frac{l^2c^2}{24y_0^4}}}1\right)y_0^2+8b^{},`$ (22)
$`0`$ $`=`$ $`{\displaystyle \frac{c}{8\pi G}}+6C\varphi _0.`$ (23)
Here we assume the brane lies at $`y=y_0`$ and the dilaton takes a constant value there $`\varphi =\varphi _0`$:
$$\varphi _0=\frac{c}{48\pi GC}.$$
(24)
Note that eq.(22) does not depend on $`b`$ and $`C`$. Eq.(23) determines the value of $`\varphi _0`$. That might be interesting since the vacuum expectation value of the dilaton cannot be determined perturbatively in string theory. Of course, (24) contains the parameter $`c`$, which indicates the non-triviality of the dilaton. The parameter $`c`$, however, can be determined from (22). Hence, in such scenario one gets a dynamical mechanism to determine of dilaton on the boundary (in our observable world).
The effective tension of the domain wall is given by
$$\sigma _{\mathrm{eff}}=\frac{3}{4\pi G}_yA=\frac{3}{4\pi Gl}\sqrt{1+\frac{kl^2}{3y_0}+\frac{l^2c^2}{24y_0^4}}.$$
(25)
We should note that the radial ($`z`$) component of the geodesic equation for the in the metric (11) is given by $`\frac{d^2x^z}{d\tau ^2}+_zA\mathrm{e}^{2A}\left(\frac{dx^t}{d\tau }\right)^2=0`$. Here $`\tau `$ is the proper time and we can normalize $`\mathrm{e}^{2A}\left(\frac{dx^t}{d\tau }\right)^2=1`$ and obtain $`\frac{d^2x^z}{d\tau ^2}+_zA=0`$. Since the cosmological constant on the brane is given by $`\frac{3}{4\pi G}`$, $`\sigma _{\mathrm{eff}}`$ gives the effective mass density: $`\frac{3}{4\pi G}\frac{d^2x^z}{d\tau ^2}=\sigma _{\mathrm{eff}}`$.
As in , defining the radius $`R`$ of the brane in the following way
$$R^2y_0,$$
(26)
we can rewrite (22) as
$$0=\frac{1}{\pi Gl}\left(\sqrt{1+\frac{kl^2}{3R^2}+\frac{l^2c^2}{24R^8}}1\right)R^4+8b^{}.$$
(27)
Especially when the dilaton vanishes ($`c=0`$) and the brane is the unit sphere ($`k=3`$), the equation (27) reproduces the result of ref. for $`𝒩=4`$ $`SU(N)`$ super Yang-Mills theory in case of the large $`N`$ limit where $`b^{}\frac{N^2}{4(4\pi )^2}`$:
$$\frac{R^3}{l^3}\sqrt{1+\frac{R^2}{l^2}}=\frac{R^4}{l^4}+\frac{GN^2}{8\pi l^3}.$$
(28)
Let us define a function $`F(R,c)`$ as
$$F(R,c)\frac{1}{\pi Gl}\left(\sqrt{1+\frac{kl^2}{3R^2}+\frac{l^2c^2}{24R^8}}1\right)R^4,$$
(29)
It appears in the r.h.s. in (27).
First we consider the $`k>0`$ case. Since
$`{\displaystyle \frac{\left(\mathrm{ln}F(R,c)\right)}{R}}`$ $`=`$ $`{\displaystyle \frac{1}{R}}\left(\sqrt{1+{\displaystyle \frac{kl^2}{3R^2}}+{\displaystyle \frac{l^2c^2}{24R^8}}}1\right)^1\left(\sqrt{1+{\displaystyle \frac{kl^2}{3R^2}}+{\displaystyle \frac{l^2c^2}{24R^8}}}\right)^1`$ (30)
$`\times \left(4+{\displaystyle \frac{kl^2}{R^2}}+4\sqrt{1+{\displaystyle \frac{kl^2}{3R^2}}+{\displaystyle \frac{l^2c^2}{24R^8}}}\right)^1`$
$`\times \left({\displaystyle \frac{8kl^2}{3R^2}}+{\displaystyle \frac{k^2l^4}{R^4}}{\displaystyle \frac{2l^2c^2}{3R^8}}\right).`$
$`F(R,c)`$ has a minimum at $`R=R_0`$, where $`R_0`$ is defined by
$$0=\frac{8kl^2}{3R_0^2}+\frac{k^2l^4}{R_0^4}\frac{2l^2c^2}{3R_0^8}.$$
(31)
When $`k>0`$, there is only one solution for $`R_0`$. Therefore $`F(R,c)`$ in the case of $`k>0`$ (sphere case) is a monotonically increasing function of $`R`$ when $`R>R_0`$ and a decreasing function when $`R<R_0`$. Since $`F(R,c)`$ is clearly a monotonically increasing function of $`c`$, we find for $`k>0`$ and $`b^{}<0`$ case that $`R`$ decreases when $`c`$ increases if $`R>R_0`$, that is, the non-trivial dilaton makes the radius smaller. Then, since $`1/R`$ corresponds to the rate of the inflation of the universe, when we Wick-rotate the sphere into the inflationary universe, the large dilaton supports the rapid universe expansion. Hence, we showed that quantum CFT living on the domain wall leads to the creation of inflationary dilatonic 4d de Sitter-brane Universe realized within 5d AdS bulk space.<sup>7</sup><sup>7</sup>7Such brane-world quantum inflation for the case of constant dilaton has been presented in refs.. In the usual 4d world the anomaly induced inflation has been suggested in ref. (no dilaton) and in ref. when a non-constant dilaton is present. Of course, such ever expanding inflationary brane-world is understood in a sense of the analytical continuation of 4d sphere to Lorentzian signature. It would be interesting to understand the relation between such inflationary brane-world and inflation in D-branes, for example, of Hagedorn type .
Since one finds
$$F(R_0,c)=\frac{klR_0^2}{4\pi G},$$
(32)
using (29) and (31), Eq.(27) has a solution if
$$\frac{klR_0^2}{4\pi G}8b^{}.$$
(33)
That puts again some bounds to the dilaton value. When $`|c|`$ is small, using (31), one obtains
$$R_0^4\frac{2c^2}{3k^2l^2},F(R_0,c)\frac{1}{4\pi G}\frac{|c|}{\sqrt{3}}.$$
(34)
Therefore Eq.(33) is satisfied for small $`|c|`$. On the other hand, when $`c`$ is large, we get
$$R_0^6\frac{c^2}{4k},F(R_0,c)\frac{\left(k|c|\right)^{\frac{2}{3}}}{4^{\frac{4}{3}}\pi G}.$$
(35)
Therefore Eq.(33) is not always satisfied and we have no solution for $`R`$ in (27) for very large $`|c|`$. Then the existence of the inflationary Universe gives a restriction on the value of $`c`$, which characterizes the behavior of the dilaton.
We now consider the $`k<0`$ case. When $`c=0`$, there is no solution for $`R`$ in (27). Let us define another function $`G(R,c)`$ as follows:
$$G(R,c)1+\frac{l^2c^2}{24R^8}+\frac{kl^2}{3R^2}.$$
(36)
Since $`G(R,c)`$ appears in the root of $`F(R,c)`$ in (29), $`G(R,c)`$ must be positive. Then
$$\frac{G(R,c)}{R}=\frac{l^2c^2}{3R^9}\frac{2kl^2}{3R^3},$$
(37)
$`G(R,c)`$ has a minimum
$$1+\frac{kl^2}{4}\left(\frac{2k}{c^2}\right)^{\frac{1}{3}},$$
(38)
when
$$R^6=\frac{c^2}{2k}.$$
(39)
Therefore if
$$c^2\frac{k^4l^6}{32},$$
(40)
$`F(R,c)`$ is real for any positive value of $`R`$. Since
$$F(0,c)=\frac{|c|}{\pi G\sqrt{24}},$$
(41)
and when $`R\mathrm{}`$
$$F(R,c)\frac{klR^2}{6\pi G}<0,$$
(42)
there is a solution $`R`$ in (27) if
$$\frac{|c|}{\pi G\sqrt{24}}>8b^{}.$$
(43)
If we Wick-rotate the solution corresponding to hyperboloid, we obtain a 4 dimensional AdS space, whose metric is given by
$$ds_{\mathrm{AdS}_4}^2=dz^2+\mathrm{e}^{\frac{2z}{R}}\left(dt^2+dx^2+dy^2\right).$$
(44)
Eq.(43) tells that there is such kind of solution due to the quantum effect if the parameter $`c`$ characterizing the behavior of the dilaton is large enough. Thus we demonstrated that due to the dilaton presence there is the possibility of quantum creation of a 4d hyperbolic wall Universe. Again, some bounds to the dilaton appear. It is remarkable that hyperbolic brane-world occurs even for usual matter content due to the dilaton. One can compare with the case in ref. where a hyperbolic 4d wall could be realized only for higher derivative conformal scalar.
In summary, in this section for constant bulk potential, we presented the nice realization of quantum creation of 4d de Sitter or 4d hyperbolic brane Universes living in 5d AdS space. The quantum dynamical determination of dilaton value is also remarkable.
## 3 Non-constant bulk potentials
We now consider the case that the dilaton field $`\varphi `$ has a non-trivial potential:
$$\frac{12}{l^2}V(\varphi )=\frac{12}{l^2}+\mathrm{\Phi }(\varphi ).$$
(45)
The surface counter terms when the dilaton field $`\varphi `$ has a non-trivial potential are given in :
$`S^{(2)}`$ $`=`$ $`S_1^\varphi +S_2^\varphi ,`$
$`S_1^\varphi `$ $`=`$ $`{\displaystyle \frac{1}{16\pi G}}{\displaystyle d^4\sqrt{g_{(4)}}\left(\frac{6}{l}+\frac{l}{4}\mathrm{\Phi }(\varphi )\right)},`$
$`S_2^\varphi `$ $`=`$ $`{\displaystyle \frac{1}{16\pi G}}{\displaystyle }d^4\{\sqrt{g_{(4)}}({\displaystyle \frac{l}{2}}R_{(4)}{\displaystyle \frac{l}{2}}\mathrm{\Phi }(\varphi )`$ (46)
$`{\displaystyle \frac{l}{4}}_{(4)}\varphi _{(4)}\varphi ){\displaystyle \frac{l^2}{8}}n^\mu \left(\sqrt{g_{(4)}}\mathrm{\Phi }(\varphi )\right)\}.`$
Following the argument in , if one replaces $`\frac{12}{l^2}`$ in (2) and $`S_1`$ in (1) with $`V(\varphi )`$ in (45) and $`S_1^\varphi `$ in (3), we obtain the gravity on the brane induced by $`S_2^\varphi `$. Then if we assume the metric in the following form
$$ds^2=f(y)dy^2+y\underset{i,j=0}{\overset{3}{}}\widehat{g}_{ij}(x^k)dx^idx^j,$$
(47)
as in (2) and $`\varphi `$ depends only on $`y`$: $`\varphi =\varphi (y)`$, we obtain the following equations of motion in the bulk:
$`0`$ $`=`$ $`{\displaystyle \frac{3}{2y^2}}{\displaystyle \frac{2kf}{y}}{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{d\varphi }{dy}}\right)^2\left({\displaystyle \frac{6}{l^2}}+{\displaystyle \frac{1}{2}}\mathrm{\Phi }(\varphi )\right)f,`$ (48)
$`0`$ $`=`$ $`{\displaystyle \frac{d}{dy}}\left({\displaystyle \frac{y^2}{\sqrt{f}}}{\displaystyle \frac{d\varphi }{dy}}\right)+\mathrm{\Phi }^{}(\varphi )y^2\sqrt{f}.`$ (49)
On the other hand, on the brane, we obtain the following equations instead of (20) and (21):
$`0`$ $`=`$ $`{\displaystyle \frac{48l^4}{16\pi G}}\left(_zA{\displaystyle \frac{1}{l}}{\displaystyle \frac{l}{24}}\mathrm{\Phi }(\varphi )\right)\mathrm{e}^{4A}+b^{}\left(4_\sigma ^4A16_\sigma ^2A\right)`$ (50)
$`4(b+b^{})\left(_\sigma ^4A+2_\sigma ^2A6(_\sigma A)^2_\sigma ^2A\right)`$
$`+2C\left(_\sigma ^4\varphi 4_\sigma ^2\varphi \right),`$
$`0`$ $`=`$ $`{\displaystyle \frac{l^4}{8\pi G}}\mathrm{e}^{4A}_z\varphi {\displaystyle \frac{l^5}{32\pi G}}\mathrm{e}^{4A}\mathrm{\Phi }^{}(\varphi )`$ (51)
$`+C\left\{A\left(_\sigma ^4\varphi 4_\sigma ^2\varphi \right)+_\sigma ^4(A\varphi )4_\sigma ^2(A\varphi )\right\}.`$
In (50) and (51), one assumes the form of the metric as in (11) instead of (47) using the change of the coordinate: $`dz=\sqrt{f}dy`$ and equations similar to (12), (13), and (14) by choosing $`l^2\mathrm{e}^{2\widehat{A}(z,k)}=y(z)`$.
Using (48) and (49), we can delete $`f`$ from the equations and we obtain an equation that contains only the dilaton field $`\varphi `$:
$`0`$ $`=`$ $`\left\{{\displaystyle \frac{5k}{2}}{\displaystyle \frac{k}{4}}y^2\left({\displaystyle \frac{d\varphi }{dy}}\right)^2+\left({\displaystyle \frac{3}{2}}y{\displaystyle \frac{y^3}{6}}\left({\displaystyle \frac{d\varphi }{dy}}\right)^2\right)\left({\displaystyle \frac{6}{l^2}}+{\displaystyle \frac{1}{2}}\mathrm{\Phi }(\varphi )\right)\right\}{\displaystyle \frac{d\varphi }{dy}}`$ (52)
$`+{\displaystyle \frac{y^2}{2}}\left({\displaystyle \frac{2k}{y}}+{\displaystyle \frac{6}{l^2}}+{\displaystyle \frac{1}{2}}\mathrm{\Phi }(\varphi )\right){\displaystyle \frac{d^2\varphi }{dy^2}}+\left({\displaystyle \frac{3}{4}}{\displaystyle \frac{y^2}{8}}\left({\displaystyle \frac{d\varphi }{dy}}\right)^2\right)\mathrm{\Phi }^{}(\varphi ).`$
First we consider a solvable case where
$$\frac{6}{l^2}+\frac{1}{2}\mathrm{\Phi }(\varphi )=\frac{2k}{y}.$$
(53)
The explicit form, or $`\varphi `$ dependence, of $`\mathrm{\Phi }(\varphi )`$ can be determined after solving the equations of motion. Then since
$$\mathrm{\Phi }^{}(\varphi )\frac{d\varphi }{dy}=\frac{4k}{y^2},$$
(54)
from (53), Eq.(52) can be rewritten as follows:
$$0=\left\{\left(\frac{d\varphi }{dy}\right)^2\frac{6}{y^2}\right\}^2.$$
(55)
Note that the $`k`$-dependence disappears in (55). The solution of (55) is trivially given by
$$\varphi =\pm \sqrt{6}\mathrm{ln}(m^2y).$$
(56)
Here $`m^2`$ is a constant of the integration.<sup>8</sup><sup>8</sup>8It is interesting that from AdS/CFT point of view the exponent of above dilaton corresponds to running gauge coupling which has a power behavior in terms of the energy parameter $`y`$. This gauge coupling corresponds to a boundary QFT with (broken) supersymmetry. Then from (53) and (54), we can find the explicit form of $`\mathrm{\Phi }(\varphi )`$:
$$\mathrm{\Phi }(\varphi )=\frac{12}{l^2}4km^2\mathrm{e}^{\frac{\varphi }{\sqrt{6}}}.$$
(57)
Note that exponential potentials of the above type often appear as the result of spherical reduction in M-theory or string theory, see discussion in ref.. One can also find that Eq.(48) is trivially satisfied. Integrating (49), we obtain
$$f=\frac{1}{\frac{2ky}{9}+\frac{f_0}{y^2}}.$$
(58)
Here $`f_0`$ is a constant of the integration and $`f_0`$ should be positive in order that $`f`$ is positive for large $`y`$. There is a (curvature) singularity at $`y=0`$. One should also note that when $`k>0`$, the horizon appears at
$$y^3=y_0^3\frac{9f_0}{2k},$$
(59)
and we find
$$yy_0.$$
(60)
Then since
$$_zA=\frac{1}{2}_z\left(\mathrm{ln}y\right)=\frac{1}{2y}\frac{dy}{dz}=\frac{1}{2y\sqrt{f(y)}},$$
(61)
Eqs.(50) and (51) have the following forms:
$`0`$ $`=`$ $`{\displaystyle \frac{1}{\pi G}}\left({\displaystyle \frac{1}{2y_0}}\sqrt{{\displaystyle \frac{f_0}{y_0^2}}{\displaystyle \frac{2ky_0}{9}}}{\displaystyle \frac{1}{2l}}+{\displaystyle \frac{kl}{3y_0}}\right)y_0^2+8b^{}`$ (62)
$`=`$ $`{\displaystyle \frac{1}{\pi G}}\left({\displaystyle \frac{1}{2R^2}}\sqrt{{\displaystyle \frac{f_0}{R^4}}{\displaystyle \frac{2kR^2}{9}}}{\displaystyle \frac{1}{2l}}+{\displaystyle \frac{kl}{3R^2}}\right)R^4+8b^{},`$
$`0`$ $`=`$ $`{\displaystyle \frac{y_0\sqrt{6}}{8\pi G}}\sqrt{{\displaystyle \frac{f_0}{y_0^2}}{\displaystyle \frac{2ky_0}{9}}}{\displaystyle \frac{kly_0\sqrt{6}}{2\pi G}}+6C\varphi _0`$ (63)
$`=`$ $`{\displaystyle \frac{R^2\sqrt{6}}{8\pi G}}\sqrt{{\displaystyle \frac{f_0}{R^4}}{\displaystyle \frac{2kR^2}{9}}}{\displaystyle \frac{klR^2\sqrt{6}}{48\pi G}}+6C\varphi _0.`$
Eq.(63) gives a value of the dilaton on the brane:
$$\varphi _0=\frac{1}{C\sqrt{6}}\frac{R^2}{8\pi G}\left(\sqrt{\frac{f_0}{R^4}\frac{2kR^2}{9}}+\frac{kl}{6}\right).$$
(64)
When $`k>0`$, Eq.(62) does not have a solution for large $`R`$ since there is an upper bound $`R\sqrt{y_0}`$ coming from (60). Even for $`k0`$, there is no solution for large $`R`$ in case of $`𝒩=4`$ Yang-Mills theory ($`b^{}<0`$) since Eq.(62) behaves for large $`R`$
$$0\frac{1}{2l\pi G}R^4+8b^{}.$$
(65)
On the other hand, if one assumes $`R`$ is small, Eq.(62) has the following form:
$$0=\frac{1}{\pi G}\left(\frac{\sqrt{f_0}}{2R^4}+\frac{kl}{3R^2}\right)R^4+8b^{}+𝒪(R^2),$$
(66)
which can be solved with respect to $`R`$:
$$R^2=\frac{3}{kl}\left(\frac{\sqrt{f_0}}{2}+8\pi Gb^{}\right).$$
(67)
Then there is a solution for $`k<0`$ ($`k>0`$) if
$$f_0>128\pi ^2G^2b_{}^{}{}_{}{}^{2}\left(f_0<128\pi ^2G^2b_{}^{}{}_{}{}^{2}\right).$$
(68)
Hence, the results are similar to those in the previous section but in the presence of non-trivial bulk potential.
One can also consider the case of no quantum corrections, i.e. $`W`$ vanishes. Putting $`C=b^{}=0`$, we obtain from (62) and (63)
$`0`$ $`=`$ $`{\displaystyle \frac{1}{2R^2}}\sqrt{{\displaystyle \frac{f_0}{R^4}}{\displaystyle \frac{2kR^2}{9}}}{\displaystyle \frac{1}{2l}}+{\displaystyle \frac{kl}{3R^2}},`$ (69)
$`0`$ $`=`$ $`\sqrt{{\displaystyle \frac{f_0}{R^4}}{\displaystyle \frac{2kR^2}{9}}}+{\displaystyle \frac{kl}{6}}.`$ (70)
Eq.(70) tells that $`k0`$ but by combining (69) and (70), we find $`R^2=\frac{kl^2}{2}`$. Then there is not consistent solution.
Note, however, that the quantum equation (65) for $`R`$ has the solution for conformally invariant higher derivative scalar whose contribution to $`b^{}`$ is positive: $`b=8/120(4\pi )^2,b^{}=28/360(4\pi )^2`$. In a similar way one can analyze other types of dilatonic potentials (numerically or using some perturbative technique) which lead to (singular) 5d AdS space with 4d constant curvature wall(s).
Let us discuss other examples in attempt to construct non-singular brane-world with inflationary brane induced by quantum effects. As the singularity usually appears at $`y=0`$, we investigate the behavior of (52) when $`y0`$. Here we only consider the case $`k>0`$. First one assumes that there is no singularity. Then $`\varphi `$, $`\frac{d\varphi }{dy}`$, and $`\frac{d^2\varphi }{dy^2}`$ would be finite and we can assume
$$\varphi \varphi _1(\text{constant})\text{when}y0.$$
(71)
It is supposed the spacetime becomes asymptotically AdS, which is presumbly the unique choice to avoid the singularity and to localize gravity on the brane . The condition to get asymptotically AdS requires
$$\mathrm{\Phi }^{}(\varphi _1)=0,$$
(72)
and one assumes
$$\mathrm{\Phi }^{}(\varphi )\beta \varphi _2^\alpha (\alpha >0),\varphi _2\varphi \varphi _1.$$
(73)
Then from (52), one gets
$$0\frac{5k}{2}\frac{d\varphi _2}{dy}+ky\frac{d^2\varphi _2}{dy^2}+\frac{3}{4}\beta \varphi _2^\alpha .$$
(74)
If we also assume $`\varphi _2`$ behaves as
$$\varphi _2\stackrel{~}{b}y^a(a>0),$$
(75)
one obtains
$`\alpha `$ $`=`$ $`1{\displaystyle \frac{1}{a}}`$ (76)
$`\beta `$ $`=`$ $`{\displaystyle \frac{4k}{3}}\stackrel{~}{b}^{\frac{1}{a}}a\left(a+{\displaystyle \frac{3}{2}}\right).`$ (77)
Eq.(76) requires $`0<\alpha <1`$ and/or $`a>1`$ and Eq.(77) tells that $`\beta `$ cannot vanish and $`\stackrel{~}{b}`$ should be positive, which tells that $`\varphi `$ increases when $`y0`$.
If we assume $`\frac{d\varphi }{dy}=0`$ at $`y=y_1>0`$ in (52), we obtain
$$0=\frac{y_1^2}{2}\left(\frac{2k}{y_1}+\frac{6}{l^2}+\frac{1}{2}\mathrm{\Phi }(\varphi (y_1))\right)\frac{d^2\varphi }{dy^2}+\frac{3}{4}\mathrm{\Phi }^{}(\varphi (y_1)).$$
(78)
In case that $`V(\varphi )=\frac{12}{l^2}+\mathrm{\Phi }(\varphi )>0`$, $`\frac{d^2\varphi }{dy^2}>0`$ $`\left(\frac{d^2\varphi }{dy^2}<0\right)`$ if $`\mathrm{\Phi }^{}(\varphi )<0`$ $`\left(\mathrm{\Phi }^{}(\varphi )>0\right)`$. Since $`\varphi `$ increases when $`y0`$, $`\varphi `$ increases monotonically if $`V(\varphi )>0`$ and $`\mathrm{\Phi }^{}(\varphi )<0`$.
One also finds that when $`y`$ is large, Eq.(52) does not depend on $`k`$:
$`0`$ $`=`$ $`\left({\displaystyle \frac{3}{2}}y{\displaystyle \frac{y^3}{6}}\left({\displaystyle \frac{d\varphi }{dy}}\right)^2\right)\left({\displaystyle \frac{6}{l^2}}+{\displaystyle \frac{1}{2}}\mathrm{\Phi }(\varphi )\right){\displaystyle \frac{d\varphi }{dy}}`$ (79)
$`+{\displaystyle \frac{y^2}{2}}\left({\displaystyle \frac{6}{l^2}}+{\displaystyle \frac{1}{2}}\mathrm{\Phi }(\varphi )\right){\displaystyle \frac{d^2\varphi }{dy^2}}+\left({\displaystyle \frac{3}{4}}{\displaystyle \frac{y^2}{8}}\left({\displaystyle \frac{d\varphi }{dy}}\right)^2\right)\mathrm{\Phi }^{}(\varphi ).`$
Let us consider the following example as a toy model:
$$l^2\mathrm{\Phi }(\varphi )=\frac{4}{3}\varphi ^{\frac{3}{2}}+\frac{3}{4}\varphi ^4\frac{1}{8}\varphi ^8+\frac{17}{24}.$$
(80)
Since
$$l^2\mathrm{\Phi }^{}(\varphi )=2\varphi ^{\frac{1}{2}}+3\varphi ^3\varphi ^7,$$
(81)
by comparing (81) with (75), one finds
$$\varphi _1=0,\alpha =\frac{1}{2},\beta =2,a=2,\stackrel{~}{b}=\frac{1}{196}.$$
(82)
Eq.(81) also tells $`\mathrm{\Phi }^{}(\varphi )=0`$ when $`\varphi =0`$ or $`1`$ and $`\mathrm{\Phi }^{}(\varphi )<0`$ when $`0<\varphi <1`$. Then if $`\varphi 0`$ when $`y0`$, we can naively expect $`\varphi 1`$ when $`y+\mathrm{}`$. This naive expectation can be confirmed by the numerical calculation, which is given in Figs.1 and 2.
In Fig.1, the behavior of $`\varphi `$ when $`y`$ is small is given and in Fig.2, the behavior of $`\varphi `$ when $`y`$ is large is drawn. From Fig.2, one can find that $`\varphi `$ goes to unity when $`y`$ is large. Then there is not any (curvature) singularity and the gravity on the brane can be localized. If we assume
$$1\varphi \eta y^\xi ,\xi <0,\text{(}\eta \text{ and }\xi \text{ are constant).}$$
(83)
the numerical calculation in Fig.2 tells
$$\xi =0.2,\eta =1.$$
(84)
Then from (48), we find the behavior of $`f(y)`$ when $`y`$ is large:
$$f\frac{l^2}{4y^2}\left\{1\frac{\xi ^2}{6}\eta ^2\left(\frac{y}{l^2}\right)^{2\xi }+\mathrm{}\right\}.$$
(85)
When $`y_0`$ is large, Eqs.(50) and (51) have the following forms:
$`0`$ $``$ $`{\displaystyle \frac{l^3\xi ^2\eta ^2}{12\pi G}}\left({\displaystyle \frac{y_0}{l^2}}\right)^{2\xi +2}+8b^{}`$ (86)
$`=`$ $`{\displaystyle \frac{l^3\xi ^2\eta ^2}{12\pi G}}\left({\displaystyle \frac{R}{l}}\right)^{4\xi +4}+8b^{},`$
$`0`$ $``$ $`{\displaystyle \frac{l^3\eta \xi }{4\pi G}}\left({\displaystyle \frac{y_0}{l^2}}\right)^{\xi +2}+6C\varphi _0`$ (87)
$`=`$ $`{\displaystyle \frac{l^3\eta \xi }{4\pi G}}\left({\displaystyle \frac{R}{l}}\right)^{2\xi +4}+6C\varphi _0.`$
Eqs.(86) and (87) can be solved with respect to $`R`$ and $`\varphi _0`$, respectively:
$`R`$ $``$ $`l\left({\displaystyle \frac{96\pi Gb^{}}{l^3\xi ^2\eta ^2}}\right)^{\frac{1}{4+4\xi }},`$ (88)
$`\varphi _0`$ $``$ $`\left({\displaystyle \frac{4b^{}}{C}}\right)\left(\eta \xi \right)^{\frac{1}{1+\xi }}\left({\displaystyle \frac{l^3}{96\pi Gb^{}}}\right)^{\frac{\xi }{2+2\xi }}.`$ (89)
Since $`b^{}=\frac{C}{4}=\frac{N^21}{4(4\pi )^2}`$ from (8) for $`𝒩=4`$ $`SU(N)`$ Yang-Mills theory and $`\frac{G}{l^3}=\frac{\pi }{2N^2}`$, Eqs.(88) and (89) tell that
$`R`$ $``$ $`l\left({\displaystyle \frac{3}{4\eta ^2\xi ^2}}\right)^{\frac{1}{4+4\xi }}=2.5\mathrm{},`$ (90)
$`\varphi _0`$ $``$ $`\left(\eta \xi \right)^{\frac{1}{1+\xi }}\left({\displaystyle \frac{4}{3}}\right)^{\frac{\xi }{2+2\xi }}=0.13\mathrm{}.`$ (91)
Here we also used (84). Since $`R`$ is not so large, the large $`y`$ or $`R`$ approximation converges slowly. Since $`0<\varphi _0<1`$, however, there is no apparent conflict. Eq.(90) shows that the brane does not lie in the asymptotically AdS region when $`y`$ is large. Anyway it suggests that there is a solution where the brane corresponds to $`S_4`$, which gives the de Sitter space after transition to lorentzian signature.
For comparison, one can consider the classical case where $`W`$ vanishes. Then (50) and (51) have the following form:
$`0`$ $`=`$ $`{\displaystyle \frac{1}{2y\sqrt{f(y)}}}{\displaystyle \frac{1}{l}}{\displaystyle \frac{l}{24}}\mathrm{\Phi }(\varphi ),`$ (92)
$`0`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{f(y)}}}{\displaystyle \frac{d\varphi }{dy}}{\displaystyle \frac{l}{4}}\mathrm{\Phi }^{}(\varphi ).`$ (93)
Here it is used
$$_zA=\frac{1}{2y\sqrt{f(y)}},\frac{d}{dz}=\frac{1}{\sqrt{f(y)}}\frac{d}{dy}.$$
(94)
Since we have
$$f(y)=\frac{\frac{3}{2y^2}\frac{1}{4}\left(\frac{d\varphi }{dy}\right)^2}{\frac{2k}{y}+\frac{6}{l^2}+\frac{1}{2}\mathrm{\Phi }(\varphi )},$$
(95)
from (48), one can delete $`f`$ and $`\frac{d\varphi }{dy}`$ in (92), (93), (95) and obtain
$$\frac{2k}{y_0}=l^2\left(\frac{3}{8}\mathrm{\Phi }^{}(\varphi )+\frac{1}{96}\mathrm{\Phi }^2\right).$$
(96)
For the potential (80), the l.h.s. in (96) is negative when $`1>\varphi >\varphi _00.000144\mathrm{}`$, which is almost all the allowed region ($`1>\varphi >0`$) in the solution for the potential in (80). Therefore there is no classical solution for the $`k>0`$ case. Then the brane solution corresponding to 4 dimensional sphere or de Sitter space cannot exist without quantum correction coming from $`W`$.
Thus, using fine-tuned dilatonic potential in AdS dilatonic gravity we presented non-singular asymptotically AdS bulk space with de Sitter brane living on the boundary. The dilatonic de Sitter brane is induced by quantum effects of the CFT on the wall. As one can see, gravity trapping occurs. The values of brane radius and of dilaton are dynamically determined.
## 4 Not exactly conformal brane quantum matter
In this section, we consider the case that the matter on the brane is not the exact CFT like super Yang-Mills theory but some exactly non-conformal theory like QED or QCD. Of course, such a theory is classically a conformally invariant one. As an explicit example in order to be able to apply large N-expansion we suppose that dominant contribution is due to $`N`$ massless Majorana spinors coupled with the dilaton, whose action is given by
$$S=\sqrt{g_{(4)}}\mathrm{e}^{a\varphi }\underset{i=1}{\overset{N}{}}\overline{\mathrm{\Psi }}_i\gamma ^\mu D_\mu \mathrm{\Psi }_i.$$
(97)
The case of minimal spinor coupling corresponds to the choice $`a=0`$. Then the trace anomaly induced action $`W`$ corresponding to (5) has the following form :
$`W`$ $`=`$ $`b{\displaystyle d^4x\sqrt{\stackrel{~}{g}}\stackrel{~}{F}A_1}`$
$`+b^{}{\displaystyle }d^4x\{A_1[2\stackrel{~}{\mathrm{}}^2+\stackrel{~}{R}_{\mu \nu }\stackrel{~}{}_\mu \stackrel{~}{}_\nu {\displaystyle \frac{4}{3}}\stackrel{~}{R}\stackrel{~}{\mathrm{}}^2+{\displaystyle \frac{2}{3}}(\stackrel{~}{}^\mu \stackrel{~}{R})\stackrel{~}{}_\mu ]A_1`$
$`+(\stackrel{~}{G}{\displaystyle \frac{2}{3}}\stackrel{~}{\mathrm{}}\stackrel{~}{R})A_1\}`$
$`{\displaystyle \frac{1}{12}}\left\{b^{\prime \prime }+{\displaystyle \frac{2}{3}}(b+b^{})\right\}{\displaystyle d^4x\left[\stackrel{~}{R}6\stackrel{~}{\mathrm{}}A_16(\stackrel{~}{}_\mu A_1)(\stackrel{~}{}^\mu A_1)\right]^2}.`$
Here
$$A_1=A+\frac{a\varphi }{3},$$
(99)
and
$$b=\frac{3N}{60(4\pi )^2},b^{}=\frac{11N}{360(4\pi )^2}.$$
(100)
We also choose $`b^{\prime \prime }=0`$ as it may be changed by finite renormalization of classical gravitational action.
First one considers a constant potential ($`\mathrm{\Phi }(\varphi )=0`$). Then the behavior of the solution in the bulk do not change with respect to those in Section 2. On the brane, we obtain the following equations corresponding to (20) and (21):
$`0`$ $`=`$ $`{\displaystyle \frac{48l^4}{16\pi G}}\left(_zA{\displaystyle \frac{1}{l}}\right)\mathrm{e}^{4A}+b^{}\left(4_\sigma ^4A_116_\sigma ^2A_1\right)`$ (101)
$`4(b+b^{})\left(_\sigma ^4A_1+2_\sigma ^2A_16(_\sigma A_1)^2_\sigma ^2A_1\right),`$
$`0`$ $`=`$ $`{\displaystyle \frac{l^4}{8\pi G}}\mathrm{e}^{4A}_z\varphi +{\displaystyle \frac{4}{3}}ab^{}\left(4_\sigma ^4A_116_\sigma ^2A_1\right).`$ (102)
Then one gets
$`0`$ $`=`$ $`{\displaystyle \frac{1}{\pi Gl}}\left\{\sqrt{1+{\displaystyle \frac{kl^2}{3R^2}}+{\displaystyle \frac{l^2c^2}{24R^8}}}1\right\}R^4+8b^{},`$ (103)
$`0`$ $`=`$ $`{\displaystyle \frac{c}{8\pi G}}+32ab^{}.`$ (104)
Note that for minimal spinor coupling the second equation does not have a solution. Eq.(103) is identical with (27). Eq.(104) can be solved with respect to $`c`$:
$$c=32\times 8\pi Gab^{},$$
(105)
but the boundary value $`\varphi _0`$ of $`\varphi `$ becomes a free parameter. Hence, for constant bulk potential there is again the possibility of quantum creation of a 4d de Sitter or a 4d hyperbolic brane living in 5d AdS bulk space. This occurs even for not exactly conformal invariant quantum brane matter. The details of this scenario are similar to those in section 2.
When there is a non-trivial potential corresponding to (53) or (57), Eq.(62) is not changed but Eq.(63) is changed into
$$0=\frac{R^2\sqrt{6}}{8\pi G}\sqrt{\frac{f_0}{R^4}\frac{2kR^2}{9}}\frac{klR^2\sqrt{6}}{2\pi G}+32ab^{}.$$
(106)
For the potential (80), Eq.(86) is not changed again and instead of (87), we obtain
$$0=\frac{l^3\eta \xi }{4\pi Gl}\left(\frac{R}{l}\right)^{2\xi +4}+32ab^{}.$$
(107)
Eqs.(106) and (107) define the parameter $`a`$, which characterizes the dilaton coupling in (97). Since the equations for $`R`$ are identical with (62) and (86), the expression of the radius is not changed. Then for the potential (80), we have
$$Rl\left(\frac{96\pi Gb^{}}{l^3\xi ^2\eta ^2}\right)^{\frac{1}{4+4\xi }},$$
(108)
if $`R`$ is large. In this case, however, the value of $`b^{}`$ in (100)is different from that of (8) for $`𝒩=4`$ $`SU(N)`$ Yang-Mills theory and we do not know the value for $`\frac{l^3}{G}`$ (it may be considered as free parameter). Then the value of $`R`$ itself will be changed from the one in the previous section.
Hence, we have shown that in case of quantum brane matter different from super Yang-Mills theory still there arises (non)-singular brane-worlds for various dilatonic potentials in d5 AdS dilatonic gravity. As in the previous section gravity is trapped. The brane represents a constant curvature space which may be considered as an inflationary phase of our observable Universe.
## 5 Discussion
In summary, the role of brane quantum matter effects in the realization of de Sitter or AdS dilatonic branes living in d5 (asymptotically) AdS space is carefully investigated. (We are working with d5 dilatonic gravity). The explicit examples of such dilatonic brane-world inflation are presented for constant bulk dilatonic potentials as well as for non-constant bulk potentials. Dilaton gives extra contributions to the effective tension of the domain wall and it may be determined dynamically from bulk/boundary equations of motion. The main part of discussion has dealing with maximally SUSY Yang-Mills theory (exact CFT) living on the brane. However, in section 4 we demonstrated that qualitatively the same results may be obtained when not exactly conformal quantum matter (like classically conformally invariant theory of dilaton coupled spinors) lives on the brane. An explicit example of toy (fine-tuned) dilatonic potential is presented for which the following results are obtained from the bulk/boundary equations of motion: 1. Non-singular asymptotically AdS space is the bulk space. 2. The brane is described by de Sitter space (inflation) induced by brane matter quantum effects. 3. The localization of gravity on the brane occurs. The price to avoid the bulk naked singularity is the fine-tuning of dilatonic potential and dynamical determination (actually, also a kind of fine-tuning) of dilaton and radius of de Sitter brane. Note also that in the same fashion as in ref. one can show that the brane CFT strongly suppresses the metric perturbations (especially, on small scales).
One can easily generalize the results of this work in different directions. For example, taking into account that it is not easy to find new dilatonic bulk solutions like asymptotically AdS space presented in this work one can think about changes in the structure of the boundary manifold. One possibility is in the consideration of a Kantowski-Sachs brane Universe. Another important question is related with the study of cosmological perturbations around the founded backgrounds and of details of late-time inflation and exit from inflationary phase in brane-world cosmology (eventual decay of de Sitter brane to FRW brane). It would be also interesting to study more examples of dilatonic potentials within the action and brane-world structure under consideration. Clearly, this can be done numerically or using some perturbative expansion of the potential.
Acknowledgements. The work by SDO has been supported in part by CONACyT (CP, ref.990356) and in part by RFBR. This research has been also supported in part by CONACyT grant 28454E.
## Appendix
AdS<sub>5</sub>/CFT<sub>4</sub> correspondence tells us that the effective action $`W_{\mathrm{CFT}}`$ of CFT in 4 dimensions is given by the path integral of the supergravity in 5 dimensional AdS space:
$`\mathrm{e}^{W_{\mathrm{CFT}}}`$ $`=`$ $`{\displaystyle [dg][d\phi ]\mathrm{e}^{S_{\mathrm{grav}}}},`$ (109)
$`S_{\mathrm{grav}}`$ $`=`$ $`S_{\mathrm{EH}}+S_{\mathrm{GH}}+S_1+S_2+\mathrm{},`$
$`S_{\mathrm{EH}}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi G}}{\displaystyle d^5x\sqrt{g_{(5)}}\left(R_{(5)}+\frac{12}{l^2}+\mathrm{}\right)},`$
$`S_{\mathrm{GH}}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi G}}{\displaystyle d^4x\sqrt{g_{(4)}}_\mu n^\mu },`$
$`S_1`$ $`=`$ $`{\displaystyle \frac{1}{8\pi Gl}}{\displaystyle d^4x\sqrt{g_{(4)}}\left(\frac{3}{l}+\mathrm{}\right)},`$
$`S_2`$ $`=`$ $`{\displaystyle \frac{1}{16\pi Gl}}{\displaystyle d^4x\sqrt{g_{(4)}}\left(\frac{1}{2}R_{(4)}+\mathrm{}\right)},`$
$`\mathrm{}.`$
Here $`\phi `$ expresses the (matter) fields besides the graviton. $`S_{\mathrm{EH}}`$ corresponds to the Einstein-Hilbert action and $`S_{\mathrm{GH}}`$ to the Gibbons-Hawking surface counter term and $`n^\mu `$ is the unit vector normal to the boundary. $`S_1`$, $`S_2`$, $`\mathrm{}`$ correspond to the surface counter terms, which cancell the divergences when the boundary in AdS<sub>5</sub> goes to the infinity.
In , two 5 dimensional balls $`B_5^{(1,2)}`$ are glued on the boundary, which is 4 dimensional sphere $`S_4`$. Instead of $`S_{\mathrm{grav}}`$, if one considers the following action $`S`$
$$S=S_{\mathrm{EH}}+S_{\mathrm{GH}}+2S_1=S_{\mathrm{grav}}+S_1S_2\mathrm{},$$
(110)
for two balls, using (109), one gets the following boundary theory in terms of the partition function :
$`{\displaystyle _{B_5^{(1)}+B_5^{(1)}+S_4}}[dg][d\phi ]\mathrm{e}^S`$ (111)
$`=`$ $`\left({\displaystyle _{B_5}}[dg][d\phi ]\mathrm{e}^{S_{\mathrm{EH}}S_{\mathrm{GH}}S_1}\right)^2`$
$`=`$ $`\mathrm{e}^{2S_2+\mathrm{}}\left({\displaystyle _{B_5}}[dg][d\phi ]\mathrm{e}^{S_{\mathrm{grav}}}\right)^2`$
$`=`$ $`\mathrm{e}^{2W_{\mathrm{CFT}}+2S_2+\mathrm{}}.`$
Since $`S_2`$ can be regarded as the Einstein-Hilbert action on the boundary, which is $`S_4`$ in the present case, the gravity on the boundary becomes dynamical. The 4 dimensional gravity is nothing but the gravity localized on the brane in the Randall-Sundrum model .
For $`𝒩=4`$ $`SU(N)`$ Yang-Mills theory, the AdS/CFT dual is given by identifying
$$l=g_{\mathrm{YM}}^{\frac{1}{2}}N^{\frac{1}{4}}l_s,\frac{l^3}{G}=\frac{2N^2}{\pi }.$$
(112)
Here $`g_{\mathrm{YM}}`$ is the coupling of the Yang-Mills theory and $`l_s`$ is the string length. Then (111) tells that the RS model is equivalent to a CFT ($`𝒩=4`$ $`SU(N)`$ Yang-Mills theory) coupled to 4 dimensional gravity including some correction coming from the higher order counter terms with a Newton constant given by
$$G_4=G/l.$$
(113)
This is an excellent explanation to why gravity is trapped on the brane.
In case that we include the dilaton field (generally with the dilaton potential), the explicit forms of the actions are given by
$`S_{\mathrm{EH}}^\varphi `$ $`=`$ $`{\displaystyle \frac{1}{16\pi G}}{\displaystyle d^5x\sqrt{g_{(5)}}\left(R_{(5)}\frac{1}{2}_\mu \varphi ^\mu \varphi +\frac{12}{l^2}+\mathrm{\Phi }(\varphi )\right)},`$
$`S_1^\varphi `$ $`=`$ $`{\displaystyle \frac{1}{16\pi G}}{\displaystyle d^4\sqrt{g_{(4)}}\left(\frac{6}{l}+\frac{l}{4}\mathrm{\Phi }(\varphi )\right)},`$
$`S_2^\varphi `$ $`=`$ $`{\displaystyle \frac{1}{16\pi G}}{\displaystyle }d^4\{\sqrt{g_{(4)}}({\displaystyle \frac{l}{2}}R_{(4)}{\displaystyle \frac{l}{2}}\mathrm{\Phi }(\varphi )`$ (114)
$`{\displaystyle \frac{l}{4}}_{(4)}\varphi _{(4)}\varphi ){\displaystyle \frac{l^2}{8}}n^\mu _\mu \left(\sqrt{g_{(4)}}\mathrm{\Phi }(\varphi )\right)\}.`$
AdS/CFT tells the effective action $`W`$ of the boundary field theory is given by
$$\mathrm{e}^W=[dg][d\varphi ][d\stackrel{~}{\phi }]\mathrm{e}^{S_{\mathrm{EH}}^\varphi S_{\mathrm{GH}}S_1^\varphi S_2^\varphi +\mathrm{}}.$$
(115)
Here $`\stackrel{~}{\phi }`$ express the fields besides the graviton and dilaton. Then if we consider the action
$$S=S_{\mathrm{EH}}^\varphi +S_{\mathrm{GH}}+2S_1^\varphi ,$$
(116)
in the two balls, instead of (110), we obtain the following boundary theory given, instead of (111):
$`{\displaystyle _{B_5^{(1)}+B_5^{(1)}+S_4}}[dg][d\varphi ][d\stackrel{~}{\phi }]\mathrm{e}^S`$ (117)
$`=`$ $`\left({\displaystyle _{B_5}}[dg][d\phi ]\mathrm{e}^{S_{\mathrm{EH}}^\varphi S_{\mathrm{GH}}S_1^\varphi }\right)^2`$
$`=`$ $`\mathrm{e}^{2W+2S_2^\varphi +\mathrm{}}.`$
As $`S_2^\varphi `$ contains the Einstein action, there appears the dilatonic gravity localized on the brane. We can also choose the trace anomaly induced action as the effective action $`W`$.
Then again in this case, by using the identifications in (112) and coupling, (117) tells that the dilatonic RS model is equivalent to a CFT coupled to 4d dilatonic gravity. Such equivalence may be useful in various explicit calculations.
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# 1 Introduction
## 1 Introduction
A quite old result in mathematical statistics concerns the eigenvalue distribution of random matrices of the form $`A=X^TX`$ where $`X`$ is a $`n\times N`$ ($`nN`$) rectangular matrix with real entries. First it was proved by Wishart that with
$$(dX):=\underset{j=1}{\overset{n}{}}\underset{k=1}{\overset{N}{}}dx_{jk},(dA):=\underset{1j<kN}{}da_{jk}$$
denoting the product of differentials of the independent elements, the change of variables from the elements of $`X`$ to the elements of $`A`$ takes place according to the formula
$$(dX)=\left(detA\right)^{(nN1)/2}(dA).$$
(1.1)
From this a description in terms of eigenvalues and eigenvectors can be undertaken by introducing the spectral decomposition
$$A=O\mathrm{\Lambda }O^T$$
where the columns of $`O`$ consist of the normalized eigenvectors of $`A`$, and $`\mathrm{\Lambda }`$ is the diagonal matrix of eigenvalues. About a decade after the work of Wishart, it was proved by a number of authors (see e.g. ) that
$$(dA)=\underset{1j<kN}{}|\lambda _k\lambda _j|\underset{j=1}{\overset{N}{}}d\lambda _j(O^TdO).$$
(1.2)
A significant qualitative feature of (1.2) is that the eigenvalue dependence factors from that of the eigenvectors.
Suppose now the elements of $`X`$ are identical, independently distributed standard Gaussian random variables so that the corresponding probability measure is proportional to
$$\underset{j=1}{\overset{n}{}}\underset{k=1}{\overset{N}{}}e^{x_{jk}^2/2}(dX)=e^{\frac{1}{2}\mathrm{Tr}(X^TX)}(dX)=e^{\frac{1}{2}_{j=1}^N\lambda _j}(dX).$$
(1.3)
Noting that $`detA=_{j=1}^N\lambda _j`$, substituting (1.2) in (1.1) and then substituting the resulting formula for $`(dX)`$ in (1.3) gives the now standard result that the eigenvalue probability density function (p.d.f.) of the matrix $`A=X^TX`$ is given by
$$\frac{1}{C}\underset{j=1}{\overset{N}{}}\lambda _j^{(nN1)/2}e^{\lambda _j/2}\underset{1j<kN}{}|\lambda _k\lambda _j|,$$
(1.4)
where $`C`$ denotes the normalization and $`\lambda _j>0`$ $`(j=1,\mathrm{},N)`$. This is referred to as the real Wishart distribution, or alternatively as the Laguerre orthogonal ensemble (LOE<sub>N</sub>), the latter name originating from the occurence of the classical Laguerre weight function $`\lambda ^ae^\lambda `$ (up to scaling of $`\lambda `$) and the invariance of (1.4) under the mapping $`AOAO^T`$.
Another random matrix structure which leads to the p.d.f. (1.4) is the block matrix
$$\left(\begin{array}{cc}0_{n\times n}& X\\ X^T& 0_{N\times N}\end{array}\right).$$
(1.5)
It is straightforward to verify that in general (1.5) has $`nN`$ zero eigenvalues, while the remaining $`2N`$ eigenvalues come in $`\pm `$ pairs. It is similarly easy to verify that with $`X`$ distributed according to (1.3) the positive eigenvalues are distributed according to (1.4) but with $`\lambda _j\lambda _j^2`$ and an additional factor of $`_{j=1}^N|\lambda _j|`$. Hence the square roots of the positive eigenvalues of (1.5) are distributed according to (1.4).
Over the past decade the p.d.f. (1.4) has found application in a number of physical problems. One example is in the theory of quantum transport through disordered wires, in which the matrix product $`X^TX`$ for $`X`$ a $`N\times N`$ random matrix modelling the top right hand block of the transmission matrix occurs in the Landauer formula for the conductance (see e.g. ). Another is in quantum chromodynamics, where the structure (1.5) models a random Dirac operator in the chiral guage, and the number of zero eigenvalues is prescribed .
Because $`A`$ is positive definite and so has positive eigenvalues the eigenvalues near $`\lambda _j=0`$ are said to be near the hard edge. For eigenvalues in this neighbourhood, it is known that with $`nN`$ fixed, the statistical properties tend to well defined limits in the $`N\mathrm{}`$ scaled limits, where the scaling is
$$\lambda _j\lambda _j/4N$$
(1.6)
. In fact the scaled $`k`$-point distribution function is known exactly . Thus with
$`a`$ $`:=`$ $`(nN1)/2`$
$`K^\mathrm{h}(X,Y)`$ $`:=`$ $`{\displaystyle \frac{J_a(X^{1/2})Y^{1/2}J_a^{}(Y^{1/2})X^{1/2}J_a^{}(X^{1/2})J_a(Y^{1/2})}{2(XY)}}`$ (1.7)
$`D_1^\mathrm{h}(x,y)`$ $`:=`$ $`{\displaystyle \frac{}{x}}S_1(x,y)`$
$`I_1^\mathrm{h}(x,y)`$ $`:=`$ $`{\displaystyle _x^y}S_1(x,z)𝑑z{\displaystyle \frac{1}{2}}\mathrm{sgn}(xy)`$
$`f_1(x,y)`$ $`:=`$ $`\left[\begin{array}{cc}S_1^\mathrm{h}(x,y)\hfill & I_1^\mathrm{h}(x,y)\hfill \\ D_1^\mathrm{h}(x,y)\hfill & S_1^\mathrm{h}(y,x)\hfill \end{array}\right]`$ (1.10)
$`\rho _{(k)}^{\mathrm{LOEh}}(x_1,\mathrm{},x_k;a)`$ $`:=`$ $`\underset{N\mathrm{}}{lim}\left({\displaystyle \frac{1}{4N}}\right)^k\rho _{(k)}^{\mathrm{LOE}_N}({\displaystyle \frac{x_1}{4N}},\mathrm{},{\displaystyle \frac{x_k}{4N}};a)`$ (1.11)
we have
$$\rho _{(k)}^{\mathrm{LOEh}}(x_1,\mathrm{},x_k;(a1)/2)=\mathrm{qdet}[f_1(x_j,x_k)]_{j,l=1,\mathrm{},k}$$
(1.12)
where qdet is the quaternion determinant introduced into random matrix theory by Dyson , and the superscipt “h” denotes the hard edge scaling (1.6).
In this work we will compute the exact distribution of the smallest eigenvalue in the scaled LOE as specified by the $`k`$-point distribution (1.12) in terms of a certain Painlevé V transcendent. We will also compute the same distribution for the scaled Laguerre symplectic ensemble (LSE), which before scaling is specified by the eigenvalue p.d.f.
$$\frac{1}{C}\underset{j=1}{\overset{N}{}}\lambda _j^ae^{\lambda _j}\underset{1j<kN}{}|\lambda _k\lambda _j|^4.$$
This p.d.f. results from positive definite matrices $`A=X^{}X`$ when the matrix $`X`$ has the $`2\times 2`$ block structure
$$\left(\begin{array}{cc}z& w\\ \overline{w}& \overline{z}\end{array}\right)$$
of a real quaternion (each eigenvalue is then doubly degenerate as well as occuring in $`\pm `$ pairs). The explict form of the corresponding scaled $`k`$-point distribution function is given in ; in taking the scaled limit with scaling (1.6) the ensemble LSE<sub>N/2</sub> is considered rather than LSE<sub>N</sub>.
Crucial to our study is knowledge, in terms of a Painlevé V transcendent, of the distribution of the smallest eigenvalue in the scaled Laguerre unitary ensemble (LUE). Before scaling, the latter distribution is specified by the eigenvalue p.d.f.
$$\frac{1}{C}\underset{j=1}{\overset{N}{}}\lambda _j^ae^{\lambda _j}\underset{1j<kN}{}|\lambda _k\lambda _j|^2,$$
(1.13)
and results from positive definite matrices $`A=X^{}X`$ when the matrix $`X`$ has complex elements. The corresponding scaled $`k`$-point distribution function
$$\rho _{(k)}^{\mathrm{LUEh}}(x_1,\mathrm{},x_k)=\underset{N\mathrm{}}{lim}\left(\frac{1}{4N}\right)^k\rho _{(k)}^{\mathrm{LUE}_N}(\frac{x_1}{4N},\mathrm{},\frac{x_k}{4N})$$
has the explicit form
$$\rho _{(k)}^{\mathrm{LUEh}}(x_1,\mathrm{},x_k)=det[K^\mathrm{h}(x_j,x_l)]_{j,l=1,\mathrm{},k},$$
(1.14)
where $`K^\mathrm{h}`$ is given by (1.7).
In general the probability that there are no eigenvalues in an interval $`J`$, $`E(0;J)`$, can be written in terms of the corresponding $`k`$-point distribution by
$$E(0;J)=1+\underset{k=1}{\overset{\mathrm{}}{}}\frac{(1)^k}{k!}_J𝑑x_1\mathrm{}_J𝑑x_k\rho _{(k)}(x_1,\mathrm{},x_k).$$
(1.15)
For $`\rho _{(k)}`$ a $`k\times k`$ determinant with entries $`g(x_j,x_k)`$ the structure (1.15) is just the expansion of the Fredholm integral operator on $`J`$ with kernel $`g(x,y)`$. Thus it follows from (1.14) that
$$E_2^\mathrm{h}(0;(0,s);a)=det(1K^\mathrm{h})$$
where $`E_2^\mathrm{h}(0;(0,s);a)`$ denotes the probability there are no eigenvalues in the interval $`(0,s)`$ for the scaled LUE (the subscript 2 characterizes the LUE via the exponent on the product of differences in (1.13)), while $`K^\mathrm{h}`$ denotes the integral operator on $`(0,s)`$ with kernel $`K^\mathrm{h}(x,y)`$. With $`p_\beta ^{\mathrm{min}}(s;a)`$ denoting the distribution of the smallest eigenvalue in the scaled LOE ($`\beta =1`$), LUE ($`\beta =2`$) or LSE ($`\beta =4`$) we have in general
$$p_\beta ^{\mathrm{min}}(s;a)=\frac{d}{ds}E_\beta ^\mathrm{h}(0;(0,s);a)$$
so to compute $`p_\beta ^{\mathrm{min}}(s;a)`$ it suffices to compute $`E_\beta ^\mathrm{h}(0;(0,s);a)`$.
For general $`a>1`$, $`E_2^\mathrm{h}(0;(0,s);a)`$ has been computed in terms of a Painlevé transcendent by Tracy and Widom . The Painlevé transcendent is denoted by $`q_\mathrm{h}`$ (in the subscript h is not present), and specified as the solution of the nonlinear equation
$$s(q_\mathrm{h}^21)(sq_\mathrm{h}^{})^{}=q_\mathrm{h}(sq_\mathrm{h}^{})^2+\frac{1}{4}(sa^2)q_\mathrm{h}+\frac{1}{4}sq_\mathrm{h}^3(q_\mathrm{h}^22)$$
(1.16)
subject to the boundary condition
$$q_\mathrm{h}(s)\underset{s0^+}{}\frac{1}{2^a\mathrm{\Gamma }(1+a)}s^{a/2}.$$
That $`q_\mathrm{h}`$ is a Painlevé transcendent follows from the transformation
$$q_\mathrm{h}(s)=\frac{1+y(x)}{1y(x)},s=x^2$$
from which one can deduce that $`y(x)`$ satisfies the Painlevé V equation
$$y^{\prime \prime }=\left(\frac{1}{2y}+\frac{1}{1y}\right)(y^{})^2\frac{1}{x}y^{}+\frac{(y1)^2}{x^2}\left(\alpha y+\frac{\beta }{y}\right)+\frac{\gamma y}{x}+\frac{\delta y(y+1)}{y1}$$
with $`\alpha =\beta =a^2/8`$, $`\gamma =0`$ and $`\delta =2`$. The result of is that the Painlevé transcendent $`q_\mathrm{h}`$ specifies $`E_2^\mathrm{h}`$ via the formula
$$E_2^\mathrm{h}(0;(0,s);a)=\mathrm{exp}\left(\frac{1}{4}_0^s(\mathrm{log}s/t)(q_\mathrm{h}(t))^2𝑑t\right).$$
(1.17)
In this work we will show that $`E_1^\mathrm{h}`$ can also be evaluated in terms of $`q_\mathrm{h}(t)`$. Specifically, we obtain the formula
$$\left(E_1^\mathrm{h}(0;(0,s);(a1)/2)\right)^2=E_2^\mathrm{h}(0;(0,s);a)\mathrm{exp}\left(\frac{1}{2}_0^s\frac{q_\mathrm{h}(t)}{\sqrt{t}}𝑑t\right).$$
(1.18)
With $`E_2^\mathrm{h}`$ and $`E_1^\mathrm{h}`$ known in terms of $`q_\mathrm{h}(t)`$ the probability $`E_4^\mathrm{h}(0;(0,s);a)`$ can also be expressed in terms of $`q_\mathrm{h}(t)`$ by using the formula
$$E_4^\mathrm{h}(0;(0,s);a+1)=\frac{1}{2}\left(E_1^\mathrm{h}(0;(0,s);(a1)/2)+\frac{E_2^\mathrm{h}(0;(0,s);a)}{E_1^\mathrm{h}(0;(0,s);(a1)/2)}\right).$$
(1.19)
Thus
$$\left(E_4^\mathrm{h}(0;(0,s);a+1)\right)^2=E_2^\mathrm{h}(0;(0,s);a)\mathrm{cosh}^2\left(\frac{1}{4}_0^s\frac{q_\mathrm{h}(t)}{\sqrt{t}}𝑑t\right).$$
(1.20)
Here $`E_4^\mathrm{h}`$ is computed by scaling the ensemble LSE<sub>N/2</sub> according to (1.6). Also, with $`E_1^\mathrm{h}(1;(0,s);a)`$ denoting the probability that there is exactly one eigenvalue in the interval $`(0,s)`$ for the scaled LOE, we have the inter-relationship
$$E_1^\mathrm{h}(1;(0,s);(a1)/2)=E_4^\mathrm{h}(0;(0,s);a+1)E_1^\mathrm{h}(0;(0,s);(a1)/2).$$
Substituting (1.18) and (1.20) shows
$$\left(E_1^\mathrm{h}(1;(0,s);(a1)/2)\right)^2=E_2^\mathrm{h}(0;(0,s);a)\mathrm{sinh}^2\left(\frac{1}{4}_0^s\frac{q_\mathrm{h}(t)}{\sqrt{t}}𝑑t\right).$$
(1.21)
Crucial to our derivation of (1.18) is a reworking of the derivation of Tracy and Widom giving the probability $`E_1^\mathrm{s}(0;(s,\mathrm{}))`$ in terms of a Painlevé transcendent. Here $`E_1^\mathrm{s}(0;(s,\mathrm{}))`$ denotes the probability that there are no eigenvalues in the interval $`(s,\mathrm{})`$ for the scaled Gaussian orthogonal ensemble (GOE). As the density falls off rapidly as $`s`$ increases, the region $`(s,\mathrm{})`$ is said to be a soft edge, thus explaining the use of the superscript “s” in $`E_1^\mathrm{s}`$.
The ensemble GOE<sub>N</sub> refers to the eigenvalue p.d.f.
$$\frac{1}{C}\underset{j=1}{\overset{N}{}}e^{x_j^2/2}\underset{1j<kN}{}|x_kx_j|.$$
(1.22)
This is realized by $`N\times N`$ real symmetric matrices, with diagonal elements having the Gaussian distribution $`N[0,1]`$, and independent off diagonal elements having the distribution $`N[0,1/\sqrt{2}]`$. Tracy and Widom begin with the quaternion determinant expression for the $`k`$-point distribution function in the finite GOE (we also draw attention to the recent work in which Painlevé type recurrence equations are obtained for the analogue of the probabilities $`E_1^\mathrm{s}`$ and $`E_1^\mathrm{h}`$ in the finite system). Instead, inspired by the observation of Baik and Rains that the square of the distribution of the largest eigenvalue in GOE is equal to the distribution of the largest eigenvalue in two independent, appropriately superimposed GOE’s, technically the ensemble
$$\mathrm{even}(\mathrm{GOE}_N\mathrm{GOE}_N),$$
(1.23)
we take as our starting point the $`k`$-point distribution of (1.23), scaled at the spectrum edge.
Now the $`k`$-point distribution of (1.23) is an ordinary determinant, whereas the $`k`$-point distribution of the GOE is a quaternion determinant. Furthermore the elements of the determinant contains terms familiar from the analysis of $`E_2^\mathrm{s}`$ given in ; this is also true of the quaternion determinant but the former involves only a subset of the latter. These facts together provide a simplified evaluation of $`E_1^\mathrm{s}`$. The power of this derivation is demonstrated by its application to the evaluation of $`E_1^\mathrm{h}`$. We find that each step used in the derivation of $`E_1^\mathrm{s}`$ has an analogous step in the case of $`E_1^\mathrm{h}`$ and this leads to (1.18).
We begin in Section 2 by providing the evaluation of the $`k`$-point distribution for the ensemble (1.23) in the scaled limit at the right hand soft edge, as well as that for the ensemble
$$\mathrm{odd}(\mathrm{LOE}_N\mathrm{LOE}_N)$$
(1.24)
in the scaled limit at the hard edge. In Section 3 we begin by using the evaluation of $`\rho _{(k)}`$ obtained in Section 2 as the starting point for the evaluation of $`E_1^\mathrm{s}(0;(s,\mathrm{}))`$, and then proceed to mimick this calculation to evaluate $`E_1^\mathrm{h}(0;(0,s);a)`$. In Section 4 our evaluations (1.18) and (1.20) are related to previously known results.
## 2 The ensembles even/odd$`(\mathrm{OE}_N(f)\mathrm{OE}_N(f))`$ for $`f`$ classical
Let OE$`{}_{N}{}^{}(f)`$ denote the matrix ensemble with eigenvalue p.d.f.
$$\underset{j=1}{\overset{N}{}}f(x_j)\underset{1j<kN}{}|x_kx_j|.$$
(2.1)
We see from (1.4) and (1.22) that the LOE is of this form with
$$f(x)=x^ae^{x/2},(x>0,a:=(nN1)/2)$$
(2.2)
while the GOE is of this form with
$$f(x)=e^{x^2/2}.$$
In fact the four special choices of $`f`$
$$f(x)=\{\begin{array}{cc}e^{x^2/2},\hfill & \mathrm{Gaussian}\hfill \\ x^{(a1)/2}e^{x/2}(x>0),\hfill & \mathrm{Laguerre}\hfill \\ (1x)^{(a1)/2}(1+x)^{(b1)/2}(1<x<1),\hfill & \mathrm{Jacobi}\hfill \\ (1+x^2)^{(\alpha +1)/2},\hfill & \mathrm{Cauchy}\hfill \end{array}$$
(2.3)
(note that the exponent of $`x`$ in the Laguerre case has been renormalized relative to (2.2)) have been shown in to possess special properties in regards to the superimposed ensembles
$$\mathrm{even}(\mathrm{OE}_N(f)\mathrm{OE}_N(f))\mathrm{and}\mathrm{odd}(\mathrm{OE}_N(f)\mathrm{OE}_N(f))$$
(2.4)
(amongst other superimposed ensembles). In particular the $`k`$-point distribution is given by a determinant formula with the same general structure in each case.
To present the formula for the $`k`$-point distributions, some additional theory from must be recalled. In particular, it is found that each of the weight functions (2.3) is one member of a pair $`(f,g)`$ which naturally occur in the study of the superimposed ensembles. Explicitly the weight functions $`g`$ are
$$g=\{\begin{array}{cc}e^{x^2},\hfill & \mathrm{Gaussian}\hfill \\ x^ae^x(x>0),\hfill & \mathrm{Laguerre}\hfill \\ (1x)^a(1+x)^b(1<x<1),\hfill & \mathrm{Jacobi}\hfill \\ (1+x^2)^\alpha ,\hfill & \mathrm{Cauchy},\hfill \end{array}$$
(2.5)
Now, let $`\{p_n(x)\}_{n=0,1,\mathrm{}}`$ denote the set of monic orthogonal polynomials associated with a particular weight function $`g`$, and let $`(p_n,p_n)_2`$ denote the corresponding normalization. Then it is shown in that for the ensemble even$`(\mathrm{OE}_N(f)\mathrm{OE}_N(f))`$
$$\rho _{(k)}(x_1,\mathrm{},x_k)=\underset{i=1}{\overset{k}{}}g(x_i)det\left[\underset{l=0}{\overset{N2}{}}\frac{p_l(x_i)p_l(x_j)}{(p_l,p_l)_2}+\frac{p_{N1}(x_i)F_{N1}(x_j)}{(p_{N1},F_{N1})_2}\right]_{i,j=1,\mathrm{},n},$$
(2.6)
where
$$F_{N1}(x)=\underset{l=N1}{\overset{\mathrm{}}{}}\frac{(p_l,I_{})_2}{(p_l,p_l)_2}p_l(x),I_{}(x):=\frac{f(x)}{g(x)}_{\mathrm{}}^xf(t)𝑑t.$$
(2.7)
Similarly for the ensemble odd$`(\mathrm{OE}_N(f)\mathrm{OE}_N(f))`$ the $`k`$-point distribution is again given by the formula (2.6) but with $`I_{}`$ in (2.7) replaced by
$$I_+(x):=\frac{f(x)}{g(x)}_x^{\mathrm{}}f(t)𝑑t.$$
(2.8)
The summation in (2.6) can be evaluated according to the Christoffel-Darboux formula, and the corresponding scaling limits are well known . To compute the scaled limit of the quantity $`F_{N1}`$ in (2.6), we first make use of results from the work to provide the explicit evaluation of the coefficients
$$(p_l,I_{})_2:=_{\mathrm{}}^{\mathrm{}}𝑑xf(x)p_l(x)_{\mathrm{}}^x𝑑tf(t),$$
applicable in all the classical cases (2.3). First we note
$$_{\mathrm{}}^xf(t)dt=\frac{1}{2}_{\mathrm{}}^{\mathrm{}}\mathrm{sgn}(xt)f(t)dt+\frac{1}{2}_{\mathrm{}}^{\mathrm{}}f(t)dt=:\stackrel{~}{\varphi }_0(x)+\stackrel{~}{s}_0,$$
(2.9)
which allows us to write
$$(p_l,I_{})_2=_{\mathrm{}}^{\mathrm{}}f(x)p_l(x)\stackrel{~}{\varphi }_0(x)𝑑x+\stackrel{~}{s}_0_{\mathrm{}}^{\mathrm{}}f(x)p_l(x)𝑑x.$$
(2.10)
Now results in give
$$\stackrel{~}{\varphi }_0(x)=\frac{1}{\gamma _0}\frac{g(x)}{f(x)}\underset{\nu =0}{\overset{\mathrm{}}{}}\underset{k=1}{\overset{\nu }{}}\left(\frac{\gamma _{2k1}}{\gamma _{2k}}\right)\frac{p_{2\nu +1}(x)}{(p_{2\nu +1},p_{2\nu +1})_2}$$
where
$$\gamma _k(p_k,p_k)_2=\{\begin{array}{cc}1,\hfill & \mathrm{Gaussian}\hfill \\ \frac{1}{2},\hfill & \mathrm{Laguerre}\hfill \\ \frac{1}{2}(2k+2+a+b),\hfill & \mathrm{Jacobi}\hfill \\ \alpha 1k,\hfill & \mathrm{Cauchy}\hfill \end{array}$$
(2.11)
This allows us to immediately evaluate the first term in (2.10).
It remains to compute the second term in (2.10). With the notation
$$\stackrel{~}{s}_l:=\frac{1}{2}_{\mathrm{}}^{\mathrm{}}f(x)p_l(x)𝑑x.$$
this term is given by $`2\stackrel{~}{s}_0\stackrel{~}{s}_l`$. Consider first the case $`l`$ even. With
$$\stackrel{~}{\varphi }_l(x):=\frac{1}{2}_{\mathrm{}}^{\mathrm{}}\mathrm{sgn}(xt)p_l(t)f(t)𝑑t$$
we know from that
$$\stackrel{~}{\varphi }_{2k}(x)=\frac{1}{\gamma _{2k}}\frac{1}{_{l=1}^k(\gamma _{2l1}/\gamma _{2l})}\frac{g(x)}{f(x)}\underset{\nu =k}{\overset{\mathrm{}}{}}\underset{l=1}{\overset{\nu }{}}\left(\frac{\gamma _{2l1}}{\gamma _{2l}}\right)\frac{p_{2\nu +1}(x)}{(p_{2\nu +1},p_{2\nu +1})_2}.$$
Forming the ratio $`\stackrel{~}{\varphi }_{2k}(x)/\stackrel{~}{\varphi }_{2k2}(x)`$ and taking the limit $`x\mathrm{}`$ shows that $`\stackrel{~}{s}_{2k}/\stackrel{~}{s}_{2k2}=\gamma _{2k2}/\gamma _{2k1}`$ and thus we have the evaluation
$$\stackrel{~}{s}_{2l}=\stackrel{~}{s}_0\underset{j=0}{\overset{l1}{}}\frac{\gamma _{2j}}{\gamma _{2j+1}}.$$
To evaluate $`\stackrel{~}{s}_l`$, $`l`$ odd, we recall the formula
$$\stackrel{~}{\varphi }_{2k+1}(x)\frac{\gamma _{2k1}}{\gamma _{2k}}\stackrel{~}{\varphi }_{2k1}(x)=\frac{1}{\gamma _{2k}}\frac{g(x)}{f(x)}\frac{p_{2k}(x)}{(p_{2k},p_{2k})_2}.$$
Taking the limit $`x\mathrm{}`$ implies $`\stackrel{~}{s}_{2k+1}=(\gamma _{2k1}/\gamma _{2k})\stackrel{~}{s}_{2k1}`$ and since $`\stackrel{~}{s}_1:=0`$ this gives
$$\stackrel{~}{s}_{2l+1}=0.$$
Thus the second term in (2.10) is fully determined.
Substituting the evaluation of $`(p_l,I_{})_2`$ obtained from the above working in (2.7) shows
$$F_{N1}(x)=\frac{1}{\gamma _0}\underset{\nu =[(N1)/2]}{\overset{\mathrm{}}{}}\frac{\underset{l=1}{\overset{\nu }{}}(\gamma _{2l1}/\gamma _{2l})}{(p_{2\nu +1},p_{2\nu +1})_2}p_{2\nu +1}(x)+2\stackrel{~}{s}_0^2\underset{l=[N/2]}{\overset{\mathrm{}}{}}\frac{\underset{j=0}{\overset{l1}{}}(\gamma _{2j}/\gamma _{2j+1})}{(p_{2l},p_{2l})_2}p_{2l}(x).$$
(2.12)
From this result we read off that
$$(p_{N1},F_{N1})_2=\{\begin{array}{cc}\frac{1}{\gamma _0}_{l=1}^{(N2)/2}(\gamma _{2l1}/\gamma _{2l})\hfill & N\mathrm{even}\hfill \\ 2\stackrel{~}{s}_0^2_{j=1}^{(N1)/2}(\gamma _{2j2}/\gamma _{2j1})\hfill & N\mathrm{odd}\hfill \end{array}$$
(2.13)
so all quantities in the expression (2.6) for $`\rho _{(k)}`$ are now known explicitly.
In the case of the ensemble odd$`(\mathrm{OE}_N(f)\mathrm{OE}_N(f))`$ the definition (2.7) of $`F_{N1}`$ has $`I_{}`$ replaced by $`I_+`$. Noting
$$_x^{\mathrm{}}f(t)𝑑t=\frac{1}{2}_{\mathrm{}}^{\mathrm{}}\mathrm{sgn}(xt)f(t)𝑑t+\frac{1}{2}_{\mathrm{}}^{\mathrm{}}f(t)𝑑t=\stackrel{~}{\varphi }_0(x)+\stackrel{~}{s}_0$$
which differs from (2.9) only in the sign of the first term, we see by revising the working which led from (2.9) to (2.12) the only modification needed to the formula (2.12) is that a minus sign be placed in front of the first term (and similarly in (2.13)).
### 2.1 Guassian ensemble at the soft edge
In the Gaussian case
$$f(x)=e^{x^2/2},g(x)=e^{x^2},p_l(x)=2^lH_l(x),(p_l,p_l)_2=\pi ^{1/2}2^ll!,$$
(2.14)
where $`H_l(x)`$ denotes the Hermite polynomial. The soft edge scaling is
$$x=(2N)^{1/2}+\frac{X}{2^{1/2}N^{1/6}}$$
(2.15)
so we seek to compute
$$\rho _{(k)}^{(\mathrm{GOE})^2\mathrm{s}}(X_1,\mathrm{},X_k):=\underset{N\mathrm{}}{lim}\left(\frac{1}{2^{1/2}N^{1/6}}\right)^k\rho _{(k)}^\mathrm{E}((2N)^{1/2}+\frac{X_1}{2^{1/2}N^{1/6}},\mathrm{},(2N)^{1/2}+\frac{X_k}{2^{1/2}N^{1/6}})$$
where E denotes the ensemble (1.23) and the r.h.s. is given by (2.6) with the substitutions (2.14).
Regarding the summation in (2.6), we know from the study of the GUE at the soft edge that
$$\underset{N\mathrm{}}{lim}\frac{1}{2^{1/2}N^{1/6}}\underset{l=0}{\overset{N2}{}}\frac{p_l(x)p_l(y)}{(p_l,p_l)_2}|_{\genfrac{}{}{0pt}{}{x=(2N)^{1/2}+X/2^{1/2}N^{1/6}}{y=(2N)^{1/2}+Y/2^{1/2}N^{1/6}}}=K^\mathrm{s}(X,Y),$$
where, with $`\mathrm{Ai}(x)`$ denoting the Airy function
$$K^\mathrm{s}(x,y):=\frac{\mathrm{Ai}(x)\mathrm{Ai}^{}(y)\mathrm{Ai}^{}(x)\mathrm{Ai}(y)}{xy}.$$
(2.16)
This is obtained using the Christoffel-Darboux summation formula and the asymptotic expansion
$$e^{x^2/2}H_n(x)=\pi ^{3/4}2^{n/2+1/4}(n!)^{1/2}n^{1/12}\left(\pi \mathrm{Ai}(u)+\mathrm{O}(n^{2/3})\right)$$
(2.17)
where $`x=(2n)^{1/2}u/2^{1/2}n^{1/6}`$.
It remains to compute the scaled limit of the term involving $`F_{N1}`$ in (2.6). Now substituting the values of $`\gamma _k`$ and $`(p_k,p_k)_2`$ from (2.11) and (2.14), and noting $`\stackrel{~}{s}_0^2=\pi /2`$, (2.12) reads
$$F_{N1}(x)=\pi ^{1/2}\underset{\nu =[(N1)/2]}{\overset{\mathrm{}}{}}\frac{\nu !}{(2\nu +1)!}H_{2\nu +1}(x)+\pi \underset{l=[N/2]}{\overset{\mathrm{}}{}}\frac{1}{2^{2l}l!}H_{2l}(x).$$
(2.18)
Next we want to combine this result with (2.13). For definiteness take $`N`$ to be even. Then we see that
$$\left(g(x)g(y)\right)^{1/2}\frac{p_{N1}(x)F_{N1}(y)}{(p_{N1},F_{N1})_2}=e^{x^2/2}2^{(N1)}H_{N1}(x)\left(A_1(y)+A_2(y)\right)$$
(2.19)
with
$`A_1(y)`$ $`:=`$ $`{\displaystyle \frac{e^{y^2/2}}{(N/21)!}}{\displaystyle \underset{\nu =0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(N/21+\nu )!}{(N1+2\nu )!}}H_{N1+2\nu }(y),`$
$`A_2(y)`$ $`:=`$ $`{\displaystyle \frac{\pi ^{1/2}e^{y^2/2}}{(N/21)!}}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{2^{N+2l}(N/2+l)!}}H_{N+2l}(y).`$ (2.20)
We remark that the summation defining $`A_1(y)`$ (with the lower terminal $`\nu =0`$ replaced by $`\nu =1`$) occurs in the study of the soft edge distribution at $`\beta =1`$ , and furthermore a procedure has been given to compute its asymptotic behaviour with the scaling (2.15), the key ingredient of which is the asymptotic expansion (2.17).
For the $`x`$-dependent terms in (2.19), (2.17) gives
$$e^{x^2/2}2^{(N1)}H_{N1}(x)|_{x=(2N)^{1/2}+X/2^{1/2}N^{1/6}}\pi ^{1/4}2^{(N1)/2+1/4}(N1)!^{1/2}N^{1/12}\mathrm{Ai}(X).$$
For $`A_1(y)`$, first note that for large $`N`$
$$\frac{(N/21+\nu )!(N1)!^{1/2}}{(N1+2\nu )!^{1/2}(N/21)!}2^\nu .$$
Then use of (2.17) with $`n=N+2\nu 1`$, $`uY2\nu /N^{1/3}`$ shows that the sum becomes the Riemann approximation to a definite integral, and thus
$$(N1)!^{1/2}A_1((2N)^{1/2}+Y/2^{1/2}N^{1/6})2^{(N1)/2}2^{1/4}\pi ^{1/4}N^{1/12}\frac{N^{1/3}}{2}_0^{\mathrm{}}\mathrm{Ai}(Yv)𝑑v.$$
Hence
$$\underset{N\mathrm{}}{lim}\frac{1}{2^{1/2}N^{1/6}}e^{x^2/2}2^{(N1)}H_{N1}(x)A_1(y)|_{\genfrac{}{}{0pt}{}{x=(2N)^{1/2}+X/2^{1/2}N^{1/6}}{y=(2N)^{1/2}+Y/2^{1/2}N^{1/6}}}=\frac{1}{2}\mathrm{Ai}(X)_0^{\mathrm{}}\mathrm{Ai}(Yv)𝑑v.$$
(2.21)
Similarly, for $`A_2(y)`$, noting that for large $`N`$
$$\frac{(N+2l)!^{1/2}}{2^{N+2l}(N/2+l)!}\frac{(N1)!^{1/2}}{(N/21)!}\frac{1}{(2\pi )^{1/2}}2^l,$$
and using (2.17) with $`n=N+2\nu `$, $`uY2\nu /N^{1/3}`$ we find
$$(N1)!^{1/2}A_2((2N)^{1/2}+Y/2^{1/2}N^{1/6})\frac{1}{\pi ^{1/2}}A_1((2N)^{1/2}+Y/2^{1/2}N^{1/6})$$
and hence
$$\underset{N\mathrm{}}{lim}\frac{1}{2^{1/2}N^{1/6}}e^{x^2/2}2^{(N1)}H_{N1}(x)A_2(y)|_{\genfrac{}{}{0pt}{}{x=(2N)^{1/2}+X/2^{1/2}N^{1/6}}{y=(2N)^{1/2}+Y/2^{1/2}N^{1/6}}}=\frac{1}{2}\mathrm{Ai}(X)_0^{\mathrm{}}\mathrm{Ai}(Yv)𝑑v.$$
(2.22)
The contributions (2.21) and (2.22) thus reinforce, so after adding to (2.16) we obtain
$$\rho _{(k)}^{(\mathrm{GOE})^2\mathrm{s}}(X_1,\mathrm{},X_k)=det\left[\left(K^\mathrm{s}(X_j,X_l)+\mathrm{Ai}(X_j)_0^{\mathrm{}}\mathrm{Ai}(X_lv)𝑑v\right)\right]_{j,l=1,\mathrm{},k}.$$
(2.23)
### 2.2 Laguerre ensemble at hard edge
In the Laguerre case
$$f(x)=x^{(a1)/2}e^{x/2},g(x)=x^ae^x,(x>0)$$
$$p_l(x)=(1)^ll!L_l^a(x),(p_l,p_l)_2=\mathrm{\Gamma }(l+1)\mathrm{\Gamma }(a+l+1),$$
(2.24)
where $`L_l^a`$ denotes the Laguerre polynomial. At the hard edge the appropriate scaling is specified by (1.6), so the task is to compute
$$\rho _{(k)}^{(\mathrm{LOE})^2\mathrm{h}}(X_1,\mathrm{},X_k;a):=\underset{N\mathrm{}}{lim}\left(\frac{1}{4N}\right)^k\rho _{(k)}^\mathrm{E}(\frac{X_1}{4N},\mathrm{},\frac{X_k}{4N})$$
(2.25)
where E denotes the ensemble (1.24) and $`\rho _{(k)}`$ on the r.h.s. is specified by (2.6).
The scaled limit of the summation in (2.6) at the hard edge occurs in the study of the LUE and is known . Thus making use of the Christoffel-Darboux formula and the large $`n`$ asymptotic formula
$$e^{y/2}y^{a/2}L_n^a(y)n^{a/2}J_a(2\sqrt{ny}).$$
(2.26)
one finds
$$\underset{N\mathrm{}}{lim}\frac{1}{4N}\underset{l=0}{\overset{N2}{}}\frac{p_l(x)p_l(y)}{(p_l,p_l)_2}|_{\genfrac{}{}{0pt}{}{x=X/4N}{y=Y/4N}}=K^\mathrm{h}(X,Y)$$
(2.27)
where $`K^\mathrm{h}`$ is specified by (1.7).
The first step in computing the term involving $`F_{N1}`$ in (2.6) is to substitute the formulas (2.24) in (2.12) modified so that there is a minus sign before the first term (recall the paragraph below (2.13). This gives
$`F_{N1}(x)`$ $`=`$ $`{\displaystyle \frac{2a!}{(a/2)!}}{\displaystyle \underset{\nu =[(N1)/2]}{\overset{\mathrm{}}{}}}{\displaystyle \frac{2^{2\nu }(a/2+\nu )!\nu !}{(a+2\nu +1)!}}L_{2\nu +1}^a(x)`$
$`+2^a{\displaystyle \frac{((a1)/2)!^2(a/2)!}{a!}}{\displaystyle \underset{l=[N/2]}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(2l)!}{2^{2l}l!(a/2+l)!}}L_{2l}^a(x).`$
We see from this formula and (2.13) (for definiteness in the latter we take $`N`$ to be even; recall that in this case a minus sign must be inserted) that
$`\left(g(x)g(y)\right)^{1/2}{\displaystyle \frac{p_{N1}(x)F_{N1}(y)}{(p_{N1},F_{N1})_2}}=(g(x))^{1/2}L_{N1}^a(x)\left({\displaystyle \frac{(N1)!}{2^{N2}((N2)/2)!(a/2+(N2)/2)!}}\right)`$ (2.28)
$`\times ({\displaystyle \underset{\nu =(N2)/2}{\overset{\mathrm{}}{}}}{\displaystyle \frac{2^{2\nu }(a/2+\nu )!\nu !}{(a+2\nu +1)!}}(g(y))^{1/2}L_{2\nu +1}^a(y)`$
$`+2^{a1}{\displaystyle \frac{((a1)/2)!^2(a/2)!^2}{a!^2}}{\displaystyle \underset{l=N/2}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(2l)!}{2^{2l}l!(a/2+l)!}}(g(y))^{1/2}L_{2l}^a(y))`$
$`:=`$ $`(g(x))^{1/2}L_{N1}^a(x){\displaystyle \frac{(N1)!}{2^{N2}((N2)/2)!(a/2+(N2)/2)!}}\left(B_1(y)+B_2(y)\right),`$
where $`B_1`$ and $`B_2`$ denote the two terms in the final brackets of the line before.
Consider the function $`B_1(y)`$. From Stirling’s formula we have
$$\frac{2^{2\nu }(a/2+\nu )!\nu !}{(a+2\nu +1)!}\pi ^{1/2}2^{(a+1)}\nu ^{a/21/2}.$$
Using this and (2.26) we that the sum defining $`B_1(y)`$ is the Riemann approximation to a definite integral and we find
$$B_1\left(\frac{Y}{4N}\right)\frac{\pi ^{1/2}}{2^{a/2+1}}\left(\frac{N}{2}\right)^{1/2}_1^{\mathrm{}}t^{1/2}J_a(\sqrt{tY})𝑑t=\frac{\pi ^{1/2}}{2^{a/2}}\left(\frac{N}{2}\right)^{1/2}\frac{1}{Y^{1/2}}_{Y^{1/2}}^{\mathrm{}}J_a(t)𝑑t.$$
(2.29)
Another application of Stirling’s formula shows
$$\frac{(N1)!}{2^{N2}((N2)/2)!(a/2+(N2)/2)!}2^{(a+1)/2}\pi ^{1/2}N^{(1a)/2}$$
so we see from further use of (2.26) together with (2.29) that
$`\underset{N\mathrm{}}{lim}(g(x))^{1/2}{\displaystyle \frac{1}{4N}}L_{N1}^a(x){\displaystyle \frac{(N1)!}{2^{N2}((N2)/2)!(a/2+(N2)/2)!}}B_1(y)|_{\genfrac{}{}{0pt}{}{x=X/4N}{y=Y/4N}}`$
$`={\displaystyle \frac{J_a(\sqrt{X})}{4\sqrt{Y}}}{\displaystyle _{Y^{1/2}}^{\mathrm{}}}J_a(t)𝑑t.`$ (2.30)
A similar analysis applied to $`B_2(y)`$ shows
$`\underset{N\mathrm{}}{lim}{\displaystyle \frac{1}{4N}}(g(x))^{1/2}L_{N1}^a(x){\displaystyle \frac{(N1)!}{2^{N2}((N2)/2)!(a/2+(N2)/2)!}}B_2(y)|_{\genfrac{}{}{0pt}{}{x=X/4N}{y=Y/4N}}`$
$`={\displaystyle \frac{J_a(\sqrt{X})}{4\sqrt{Y}}}{\displaystyle _{Y^{1/2}}^{\mathrm{}}}J_a(t)𝑑t,`$ (2.31)
which thus reinforces (2.2). Thus adding twice (2.2) to (2.27) we have
$$\rho _{(k)}^{(\mathrm{LOE})^2\mathrm{h}}(X_1,\mathrm{},X_k;(a1)/2)=det\left[\left(K^\mathrm{h}(X_j,X_l)+\frac{J_a(\sqrt{X_j})}{2\sqrt{X_l}}_{\sqrt{X_l}}^{\mathrm{}}J_a(t)𝑑t\right)\right]_{j,l=1,\mathrm{},k}.$$
(2.32)
## 3 Gap probabilities at the spectrum edge
### 3.1 The probability $`E_1^\mathrm{s}(0;(s,\mathrm{}))`$
It was remarked in the Introduction that the probability $`E_1^\mathrm{s}(0;(s,\mathrm{}))`$ has been computed in terms of a Painlevé II transcendent by Tracy and Widom . Expliclity let $`q_\mathrm{s}`$ denote the solution of the particular Painlevé II differential equation
$$(q_\mathrm{s})^{\prime \prime }=sq_\mathrm{s}+2(q_\mathrm{s})^3$$
(3.1)
subject to the boundary condition $`q_\mathrm{s}(s)\mathrm{Ai}(s)`$ as $`s\mathrm{}`$. Then it is shown in that
$$\left(E_1^\mathrm{s}(0;(s,\mathrm{}))\right)^2=E_2^\mathrm{s}(0;(s,\mathrm{}))\mathrm{exp}\left(_s^{\mathrm{}}q_\mathrm{s}(t)𝑑t\right)$$
(3.2)
where
$$E_2^\mathrm{s}(0;(s,\mathrm{}))=\mathrm{exp}\left(_s^{\mathrm{}}(ts)q_\mathrm{s}^2(t)𝑑t\right).$$
(3.3)
Here we will use the evaluation of the $`k`$-point distribution (2.23) to provide a simplified derivation of (3.2) while still following the essential strategy of .
By definition of (1.23) it follows that
$$\left(E_1^{\mathrm{GOE}_N}(0;(s,\mathrm{}))\right)^2=E^{\mathrm{even}(\mathrm{GOE}_N\mathrm{GOE}_N)}(0;(s,\mathrm{})),$$
which with the scaling (2.15) implies
$$\left(E_1^\mathrm{s}(0;(s,\mathrm{}))\right)^2=E^{(\mathrm{GOE})^2\mathrm{s}}(0;(s,\mathrm{})).$$
(3.4)
Now, recalling the determinant formula (2.23), we see from (1.15) and the text immediately below that $`E^{(\mathrm{GOE})^2s}`$ can be written as the determinant of a Fredholm integral operator. Thus
$$\left(E_1^\mathrm{s}(0;(s,\mathrm{}))\right)^2=det\left(1(K^\mathrm{s}+AB)\right)$$
(3.5)
where $`K^\mathrm{s}`$ is the integral operator on $`(s,\mathrm{})`$ with kernel (2.16) while $`A`$ is the operator which multiplies by $`\mathrm{Ai}(x)`$, while $`B`$ is the integral operator with kernel $`_0^{\mathrm{}}\mathrm{Ai}(yv)𝑑v`$.
Removing $`(1K^\mathrm{s})`$ as a factor from (3.5) and recalling
$$det(1K^\mathrm{s})=E_2^\mathrm{s}(0;(s,\mathrm{}))$$
we obtain
$`\left(E_1^\mathrm{s}(0;(s,\mathrm{}))\right)^2`$ $`=`$ $`E_2^\mathrm{s}(0;(s,\mathrm{}))det\left(1(1K^\mathrm{s})^1AB\right)`$ (3.6)
$`=`$ $`E_2^\mathrm{s}(0;(s,\mathrm{}))\left(1{\displaystyle _s^{\mathrm{}}}(1K^\mathrm{s})^1A[y]B(y)𝑑y\right)`$
where the second equality follows from the fact that $`(1K^\mathrm{s})^1A[y]`$ is the eigenfunction of the operator $`(1K^\mathrm{s})^1AB`$, so the eigenvalue is
$$_s^{\mathrm{}}(1K^\mathrm{s})^1A[y]B(y)𝑑y$$
Analogous to the notation of we put
$$\varphi ^\mathrm{s}(x):=A(x)=\mathrm{Ai}(x),Q^\mathrm{s}(x):=(1K^\mathrm{s})^1A[x]$$
so that
$$_s^{\mathrm{}}(1K^\mathrm{s})^1A[y]B(y)dy=_s^{\mathrm{}}dyQ^\mathrm{s}(y)_{\mathrm{}}^y\varphi ^\mathrm{s}(v)dv=:u_ϵ^\mathrm{s}$$
(3.7)
(the notation $`u_ϵ^\mathrm{s}`$ – without the superscript s – is used for an analogous quantity in ). Note from (3.6) that with the notation (3.7) we have
$$(E_1^\mathrm{s}(0;(s,\mathrm{}))^2=E_2^\mathrm{s}(0;(0;(s,\mathrm{}))(1u_ϵ^\mathrm{s}).$$
(3.8)
Following , our objective is to derive coupled differential equations for $`u_ϵ`$ and the quantity
$$q_ϵ^\mathrm{s}:=_s^{\mathrm{}}𝑑y\rho ^\mathrm{s}(s,y)_{\mathrm{}}^y\varphi ^\mathrm{s}(v)𝑑v$$
(3.9)
where $`\rho ^\mathrm{s}(x,y)`$ is the kernel of the operator $`(1K^\mathrm{s})^1`$. These equations will involve
$$Q^\mathrm{s}(s)=:q_\mathrm{s}=_s^{\mathrm{}}dy\rho ^\mathrm{s}(s,y)\varphi ^\mathrm{s}(y),$$
(3.10)
which in is shown to be the Painlevé II transcendent specified by the solution of (3.1), and their derivation relies on the formula
$$\frac{}{s}Q^\mathrm{s}(y)=q_\mathrm{s}\left(\delta ^+(ys)+\rho ^\mathrm{s}(s,y)\right),$$
(3.11)
where $`\delta ^+(ys)`$ is such that
$$_s^{\mathrm{}}\delta ^+(ys)f(y)𝑑y=f(s),$$
as well as the formula
$$\left(\frac{}{s}+\frac{}{x}+\frac{}{y}\right)\rho ^\mathrm{s}(x,y)=Q^\mathrm{s}(x)Q^\mathrm{s}(y).$$
(3.12)
Now, differentiating (3.7) with respect to $`s`$ we have
$$(u_ϵ^\mathrm{s})^{}=q_\mathrm{s}_{\mathrm{}}^s\varphi ^\mathrm{s}(v)𝑑v+_s^{\mathrm{}}𝑑y\left(\frac{}{s}Q^\mathrm{s}(y)\right)_{\mathrm{}}^y\varphi ^\mathrm{s}(v)𝑑v=q_\mathrm{s}q_ϵ^\mathrm{s}$$
(3.13)
where to obtain the second equality use has been make of (3.11) and the definition (3.9). We now seek a formula for $`(q_ϵ^\mathrm{s})^{}`$. Making use of (3.12) in (3.9) shows
$`(q_ϵ^\mathrm{s})^{}`$ $`=`$ $`{\displaystyle _s^{\mathrm{}}}𝑑y{\displaystyle \frac{}{y}}\rho ^\mathrm{s}(s,y){\displaystyle _{\mathrm{}}^y}\varphi ^\mathrm{s}(v)𝑑vq_\mathrm{s}{\displaystyle _s^{\mathrm{}}}𝑑yQ^\mathrm{s}(y){\displaystyle _{\mathrm{}}^y}\varphi ^\mathrm{s}(v)𝑑v`$ (3.14)
$`=`$ $`{\displaystyle _s^{\mathrm{}}}\rho ^\mathrm{s}(s,y)\varphi ^\mathrm{s}(y)𝑑yq_\mathrm{s}u_ϵ^\mathrm{s}`$
$`=`$ $`q_\mathrm{s}(1u_ϵ^\mathrm{s})`$
where the final equality follows from the definitions (3.7) and (3.10).
As $`q_\mathrm{s}`$ is known, the system of equations (3.13) and (3.14) fully determines $`u_ϵ^\mathrm{s}`$ and $`q_ϵ^\mathrm{s}`$ once boundary conditions are specified. Now $`Q^\mathrm{s}(y)`$ is smooth, so we see from (3.7) that
$$u_ϵ^\mathrm{s}0$$
(3.15)
as $`s\mathrm{}`$. On the other hand $`\rho (s,y)=\delta ^+(sy)+R(s,y)`$ where $`R(s,y)`$ is smooth, so for $`s\mathrm{}`$
$$q_ϵ^\mathrm{s}_{\mathrm{}}^{\mathrm{}}\varphi ^\mathrm{s}(v)𝑑v=_{\mathrm{}}^{\mathrm{}}\mathrm{Ai}(v)𝑑v=1.$$
(3.16)
The unique solution of the coupled equations (3.13) and (3.14) satisfying (3.15) and (3.16) is easily shown to be
$$u_ϵ^\mathrm{s}=1e^{\mu _\mathrm{s}},q_ϵ^\mathrm{s}=e^{\mu _\mathrm{s}}$$
(3.17)
where
$$\mu _\mathrm{s}:=_s^{\mathrm{}}q_\mathrm{s}(x)𝑑x.$$
Substituting the evaluation of $`u_ϵ^\mathrm{s}`$ from (3.17) in (3.8) reclaims (3.2), as desired.
### 3.2 The probability $`E_1^\mathrm{h}(0;(0,s);a)`$
All the steps leading to the rederivation of (3.2) given in the previous section have analogues for the probability $`E_1^\mathrm{h}(0;(0,s);a)`$ which lead to the evaluation (1.18).
First, the analogue of (3.4) is
$$\left(E_1^\mathrm{h}(0;(0,s);a)\right)^2=E^{(\mathrm{LOE})^2\mathrm{h}}(0;(0,s);a),$$
while use of the determinant formula (2.32) in (1.15) then gives
$$\left(E_1^\mathrm{h}(0;(0,s);(a1)/2)\right)^2=det\left(1(K^\mathrm{h}+CD)\right).$$
(3.18)
Here $`K^\mathrm{h}`$ is the integral operator on $`(0,s)`$ with kernel (1.7), while $`C`$ is the operator which multiplies by $`J_a(\sqrt{y})`$, while $`D`$ is the integral operator with kernel
$$\frac{1}{2\sqrt{y}}_\sqrt{y}^{\mathrm{}}J_a(t)𝑑t.$$
(3.19)
Recalling
$$det(1K^\mathrm{h})=E_2^\mathrm{h}(0;(0,s);a)$$
we see that analogous to (3.6), (3.18) can be rewritten
$$\left(E_1^\mathrm{h}(0;(0,s);(a1)/2)\right)^2=E_2^\mathrm{h}(0;(0,s);a)\left(1_0^s(1K^\mathrm{h})^1C[y]D(y)𝑑y\right).$$
(3.20)
Analogous to the notation of we put
$$\varphi ^\mathrm{h}(x):=C(x)=J_a(\sqrt{x}),Q^\mathrm{h}(y):=(1K^\mathrm{h})^1C[y].$$
After changing variables $`t=\sqrt{u}`$ in (3.19) we see that in terms of this notation
$$_0^s(1K^\mathrm{h})^1C[y]D(y)dy=\frac{1}{4}_0^sdyQ^\mathrm{h}(y)\frac{1}{\sqrt{y}}_y^{\mathrm{}}du\frac{1}{\sqrt{u}}\varphi ^\mathrm{h}(u)=:u_ϵ^\mathrm{h},$$
(3.21)
and in turn this latter notation used in (3.20) gives
$$\left(E_1^\mathrm{h}(0;(0,s);(a1)/2)\right)^2=E_1^\mathrm{h}(0;(0,s);a)(1u_ϵ^\mathrm{h}).$$
(3.22)
Now, with
$$Q^\mathrm{h}(s)=:q_\mathrm{h}=_0^sdy\rho ^\mathrm{h}(s,y)\varphi ^\mathrm{h}(y),$$
which in is shown to be the Painlevé V transcendent specified by the nonlinear equation (1.16), analogous to (3.11) we have
$$\frac{}{s}Q^\mathrm{h}(y)=q_\mathrm{h}\left(\delta ^+(ys)+\rho ^\mathrm{h}(s,y)\right).$$
Use of this formula in (3.21) then shows
$$(u_ϵ^\mathrm{h})^{}=\frac{1}{4}q_\mathrm{h}q_ϵ^\mathrm{h},q_ϵ^\mathrm{h}:=_0^s𝑑y\rho ^\mathrm{h}(s,y)\frac{1}{\sqrt{y}}_y^{\mathrm{}}𝑑u\frac{1}{\sqrt{u}}\varphi (u)$$
(3.23)
which is the analogue of (3.13).
Next we seek a formula for the derivative with respect to $`s`$ of $`q_ϵ^\mathrm{h}`$. For this purpose we note from that
$$x\frac{}{x}\rho ^\mathrm{h}(x,y)+s\frac{}{s}\rho ^\mathrm{h}(x,y)=\frac{}{y}\left(y\rho ^\mathrm{h}(x,y)\right)+\frac{1}{4}Q^\mathrm{h}(x)Q^\mathrm{h}(y)$$
(c.f. (3.12)). This formula applied to (3.23) shows
$`s(q_ϵ^\mathrm{h})^{}`$ $`=`$ $`{\displaystyle _0^s}dy({\displaystyle \frac{d}{dy}}(y\rho ^\mathrm{h}(s,y)){\displaystyle \frac{1}{\sqrt{y}}}{\displaystyle _y^{\mathrm{}}}du{\displaystyle \frac{1}{\sqrt{u}}}\varphi ^\mathrm{h}(u)+q_\mathrm{h}u_ϵ^\mathrm{h}`$ (3.24)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _0^s}𝑑y\rho ^\mathrm{h}(s,y){\displaystyle \frac{1}{\sqrt{y}}}{\displaystyle _y^{\mathrm{}}}𝑑u{\displaystyle \frac{1}{\sqrt{u}}}\varphi ^\mathrm{h}(u){\displaystyle _0^s}𝑑y\rho ^\mathrm{h}(s,y)\varphi ^\mathrm{h}(y)+q_\mathrm{h}u_ϵ^\mathrm{h}`$
$`=`$ $`{\displaystyle \frac{1}{2}}q_ϵ^\mathrm{h}q_\mathrm{h}(1u_ϵ^\mathrm{h})`$
The coupled equations (3.23) and (3.24) must be solved subject to the $`s0`$ boundary conditions
$$u_ϵ^\mathrm{h}0,\sqrt{s}q_ϵ^\mathrm{h}_0^{\mathrm{}}\frac{1}{\sqrt{u}}\varphi (u)𝑑u=2_0^{\mathrm{}}J_a(v)𝑑v=2.$$
(3.25)
The occurence of $`\sqrt{s}q_ϵ^\mathrm{h}`$ in (3.25) suggests we introduce
$$\stackrel{~}{q}_ϵ^\mathrm{h}:=\sqrt{s}q_ϵ^\mathrm{h}$$
in (3.23) and (3.24). Doing this gives the system of equations
$$\sqrt{s}(u_ϵ^\mathrm{h})^{}=\frac{1}{4}q_\mathrm{h}\stackrel{~}{q}_ϵ^\mathrm{h},\sqrt{s}(\stackrel{~}{q}_ϵ^\mathrm{h})^{}=q_\mathrm{h}(1u_ϵ^\mathrm{h}).$$
(3.26)
Introducing the new independent variable
$$\mu _\mathrm{h}:=_0^s\frac{1}{x^{1/2}}q_\mathrm{h}(x)𝑑x$$
we see that (3.26) reduces to the system with constant coefficients
$$\frac{d}{d\mu }u_ϵ^\mathrm{h}=\frac{1}{4}\stackrel{~}{q}_ϵ^\mathrm{h},\frac{d}{d\mu }\stackrel{~}{q}_ϵ^\mathrm{h}=(1u_ϵ^\mathrm{h}).$$
(3.27)
The solution satisfying (3.25) is
$$u_ϵ^\mathrm{h}=1e^{\frac{1}{2}\mu _\mathrm{h}},\stackrel{~}{q}_ϵ^\mathrm{h}=2e^{\frac{1}{2}\mu _\mathrm{h}}.$$
(3.28)
The stated formula (1.18) for $`E_1^\mathrm{h}(0;(0,s);(a1)/2)`$ now follows by substituting this evaluation of $`u_ϵ^\mathrm{h}`$ in (3.22).
## 4 Discussion
### 4.1 Special values of $`a`$
Edelman was the first person to obtain the exact evaluation of $`E_1^\mathrm{h}(0;(0,s);a)`$, albeit for two special values of $`a`$ only, namely $`a=1/2`$ and $`a=0`$. In terms of the scaling (1.6) the results of are
$`E_1^\mathrm{h}(0;(0,s);{\displaystyle \frac{1}{2}})`$ $`=`$ $`e^{(s/8+\sqrt{s}/2)}`$ (4.1)
$`E_1^\mathrm{h}(0;(0,s);0)`$ $`=`$ $`e^{s/8}.`$ (4.2)
Subsequently it was shown by the present author that $`E_1^\mathrm{h}(0;(0,s);a)`$ for $`a\text{ZZ}_0`$ can be expressed as a $`2a`$-dimensional integral. Explicitly
$$E_1^\mathrm{h}(0;(0,s);a)=Ce^{s/8}\left(\frac{1}{2\pi s^{1/2}}\right)^{2a}_{[\pi ,\pi ]^{2a}}\underset{j=1}{\overset{2a}{}}e^{s^{1/2}\mathrm{cos}\theta _j}e^{i\theta _j}\underset{1j<k2a}{}|e^{i\theta _k}e^{i\theta _j}|^4d\theta _1\mathrm{}d\theta _{2a}$$
(4.3)
where
$$C=\underset{j=1}{\overset{2a}{}}\frac{\mathrm{\Gamma }(3/2)\mathrm{\Gamma }(3/2+j/2)}{\mathrm{\Gamma }(1+j/2)}.$$
We remark that well known integration procedures (see e.g. ) allow this integral to be expressed as a Pfaffian. Such Pfaffian formulas, deduced in a different way, have been given in .
The formula (1.18) relates $`E_1^\mathrm{h}(0;(0,s);(a1)/2)`$ to $`E_2^\mathrm{h}(0;(0,s);a)`$, so it is appropriate to consider the evaluation of the latter for special values of $`a`$. Analogous to (4.3) we have that for $`a\text{ZZ}_0`$
$$E_2^\mathrm{h}(0;(0,s);a)=e^{s/4}\left(\frac{1}{2\pi }\right)^a\frac{1}{a!}_{[\pi ,\pi ]^a}\underset{j=1}{\overset{a}{}}e^{s^{1/2}\mathrm{cos}\theta _j}\underset{1j<ka}{}|e^{i\theta _k}e^{i\theta _j}|^2d\theta _1\mathrm{}d\theta _a$$
(4.4)
This integral can easily be written as a Toeplitz determinant, with the result
$$E_2^\mathrm{h}(0;(0,s);a)=e^{s/4}det\left[I_{jk}(s^{1/2})\right]_{j,k=1,\mathrm{},a}$$
(4.5)
where $`I_n(x)`$ denotes the Bessel function of purely imaginary argument. As an aside, it is interesting to note that the integral (4.4) is the generating function for the enumeration of various combinatorial objects, including quantities related to random permutations , random words and random walks .
The formula (1.17) gives
$$4\frac{d}{ds}\left(s\frac{d}{ds}\mathrm{log}E_2^\mathrm{h}(0;(0,s);a)\right)=\left(q_\mathrm{h}(s)\right)^2,$$
(4.6)
so (4.5) implies that for $`a\text{ZZ}_0`$, $`q_\mathrm{h}^2`$ can be expressed in terms of Bessel functions. The simplest case is $`a=0`$, when we have
$$E_2^\mathrm{h}(0;(0,s);a))=e^{s/4}.$$
(4.7)
Substituting in (4.6) gives
$$q_\mathrm{h}(s)=1$$
(4.8)
and substituting this in (1.18) we reclaim (4.1). The next simplest case is $`a=1`$ when we have
$$E_2^\mathrm{h}(0;(0,s);1)=e^{s/4}I_0(s^{1/2})$$
(4.9)
and so
$$\left(q_\mathrm{h}(s)\right)^2=14\frac{d}{ds}\left(s\frac{d}{ds}\mathrm{log}I_0(s^{1/2})\right).$$
For this to be consistent with (1.18) we see that the identity
$$I_0(s^{1/2})=\mathrm{exp}(\frac{1}{2}_0^s\frac{1}{\sqrt{t}}(14\frac{d}{dt}\left(t\frac{d}{dt}\mathrm{log}I_0(t^{1/2})\right))^{1/2}dt$$
must hold. This in turn is equivalent to the statement that $`y:=\mathrm{log}I_0(s^{1/2})`$ satisfies the nonlinear equation
$$4sy^{\prime \prime }+4s(y^{})^2+4y^{}1=0,$$
a fact which is readily verified using Bessel function identities.
Special evaluations are also known for $`E_4^\mathrm{h}(0;(0,s);a)`$ in the case $`a\text{ZZ}_0`$ . These evaluations are in terms of a certain generalized hypergeometric function based on Jack polynomials, while in the case $`a`$ even Pfaffian formulas are also known . In quoting from these results, one must be aware that $`E_4^\mathrm{h}`$ is defined starting with the ensemble LSE<sub>N/2</sub>, and scaling according to (1.6). This means that the results of require some rescaling of $`s`$. Doing this, we note from that the simplest cases are $`a=0`$ and $`a=1`$, when we have
$`E_4^\mathrm{h}(0;(0,s);0)`$ $`=`$ $`e^{s/8}`$ (4.10)
$`E_1^\mathrm{h}(0;(0,s);1)`$ $`=`$ $`e^{s/8}{}_{0}{}^{}F_{1}^{}({\displaystyle \frac{1}{2}};{\displaystyle \frac{s}{16}})=e^{s/8}\mathrm{cosh}{\displaystyle \frac{\sqrt{s}}{2}}.`$ (4.11)
The result (4.10) can only be related to (1.20) in the limit $`a1^{}`$, since for $`a1`$ $`E_2^\mathrm{h}(0;(0,s);a)=0`$. However, as the limiting forms of the quantities on the r.h.s. are not known we cannot readily check the consistency with (4.10). In contrast the consistency between (1.20) and (4.11) is immediate upon recalling (4.7) and (4.8).
### 4.2 Connection between $`E_\beta ^\mathrm{h}`$ and $`E_\beta ^\mathrm{s}`$
In previous articles it has been noted that for $`a\mathrm{}`$, after appropriate rescaling of the coordinates, the scaled $`k`$-point distribution function for the infinite Laguerre ensemble at the hard edge becomes equal to the scaled $`k`$-point distribution function for the infinite Gaussian ensemble at the soft edge. Note that in the symplectic case the scalings are done starting with the ensembles GSE<sub>N/2</sub> and LSE<sub>N/2</sub>. Let
$$a(\beta )=\{\begin{array}{cc}(a1)/2,\hfill & \beta =1\hfill \\ a,\hfill & \beta =2\hfill \\ a+1,\hfill & \beta =4\hfill \end{array}$$
Explicitly, it was checked that the scaled $`k`$-point distribution function for the infinite Laguerre ensemble at the hard edge, with parameter $`aa(\beta )`$ and after the rescaling of coordinates
$$xa^22a(a/2)^{1/3}x,$$
equals the soft edge distribution functions for the corresponding Gaussian ensemble results in the $`a\mathrm{}`$ limit. We must therefore have
$$\underset{a\mathrm{}}{lim}E_\beta ^\mathrm{h}(0;(0,a^22a(a/2)^{1/3}s);a(\beta ))=E_\beta ^\mathrm{s}(0;(0,\mathrm{})).$$
(4.12)
To verify (4.12), we first recall some additional results from . Write
$$\frac{1}{2}q_\mathrm{h}=\left((sR_\mathrm{h})^{}\right)^{1/2}$$
(4.13)
so that, after integrating by parts, (1.17) reads
$$E_2^\mathrm{h}(0;(0,s);a)=\mathrm{exp}\left(_0^sR_\mathrm{h}(t)𝑑t\right).$$
(4.14)
Then $`\sigma :=sR_\mathrm{h}(s)`$ is shown to satisfy the particular Painlevé III equation in $`\sigma `$ form (for an account of the latter see )
$$(s\sigma ^{\prime \prime })^2+\sigma ^{}(\sigma s\sigma ^{})(4\sigma ^{}1)a^2(\sigma ^{})^2=0.$$
(4.15)
Similarly let
$$q_\mathrm{s}(s)=(R_\mathrm{s}^{}(s))^{1/2}$$
(4.16)
so that (3.3) reads
$$E_2^\mathrm{s}(0;(s,\mathrm{}))=\mathrm{exp}\left(_s^{\mathrm{}}R_\mathrm{s}(t)\right).$$
(4.17)
It is shown in that $`R_\mathrm{s}`$ satisfies the particular Painlevé II equation in $`\sigma `$ form
$$(R_\mathrm{s}^{\prime \prime })^2+4R_\mathrm{s}^{}\left((R_\mathrm{s}^{})^2sR_\mathrm{s}^{}+R_\mathrm{s}\right)=0.$$
(4.18)
Substituting (4.14) and (4.17) in (4.12) with $`\beta =2`$ we deduce that the validity of the latter is equivalent to the statement
$$2a(a/2)^{1/3}R_\mathrm{h}(a^22a(a/2)^{1/3}t)\underset{a\mathrm{}}{}R_\mathrm{s}(t).$$
(4.19)
This can be verified by introducing the function
$$\stackrel{~}{\sigma }(s)=\frac{2a(a/2)^{1/3}}{a^2}\sigma \left(a^22a(a/2)^{1/3}s\right)2a(a/2)^{1/3}R_\mathrm{h}\left(a^22a(a/2)^{1/3}s\right)$$
into (4.15) and taking the limit $`a\mathrm{}`$. One finds the differential equation (4.18) results with $`R_\mathrm{s}=\stackrel{~}{\sigma }(s)`$.
In terms of (4.13), the evaluations (1.18) and (1.20) read
$`\left(E_1^\mathrm{h}(0;(0,s);(a1)/2)\right)^2`$ $`=`$ $`E_2^\mathrm{h}(0;(0,s);a)\mathrm{exp}\left({\displaystyle _0^s}{\displaystyle \frac{((tR_\mathrm{h}(t))^{})^{1/2}}{\sqrt{t}}}𝑑t\right)`$
$`\left(E_4^\mathrm{h}(0;(0,s);a+1)\right)^2`$ $`=`$ $`E_2^\mathrm{h}(0;(0,s);a)\mathrm{cosh}^2\left({\displaystyle \frac{1}{2}}{\displaystyle _0^s}{\displaystyle \frac{(tR_\mathrm{h}(t))^{})^{1/2}}{\sqrt{t}}}𝑑t\right).`$
Making use of (4.12) in the case $`\beta =2`$ (which has just been verified), and (4.19) together with (4.16), we see that
$`\underset{a\mathrm{}}{lim}\left(E_1^\mathrm{h}(0;(0,a^22a(a/2)^{1/3}s);(a1)/2)\right)^2`$ $`=`$ $`E_2^\mathrm{s}(0;(s,\mathrm{}))\mathrm{exp}\left({\displaystyle _s^{\mathrm{}}}q_\mathrm{s}(t)𝑑t\right)`$
$`\underset{a\mathrm{}}{lim}\left(E_4^\mathrm{h}(0;(0,a^22a(a/2)^{1/3}s);a+1)\right)^2`$ $`=`$ $`E_2^\mathrm{s}(0;(s,\mathrm{}))\mathrm{cosh}^2\left({\displaystyle \frac{1}{2}}{\displaystyle _s^{\mathrm{}}}q_\mathrm{s}(t)𝑑t\right).`$
We recognize the right hand sides in these expressions as $`E_\beta ^\mathrm{s}(0;(s,\mathrm{}))`$ for $`\beta =1`$ and 4 respectively (recall eq. (3.2); as noted in $`E_4^\mathrm{s}(0;(s,\mathrm{}))`$ can be deduced from $`E_1^\mathrm{s}(0;(s,\mathrm{}))`$ and $`E_2^\mathrm{s}(0;(s,\mathrm{}))`$ because of the validity of an inter-relationship analogous to (1.19)).
### Acknowledgements
This work was supported by the Australian Research Council.
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# TESTING CP VIOLATION IN THE HIGGS SECTOR AT MUON COLLIDERS⋆
## 1 Introduction
In spite of the fact that the Standard Model (SM) of electroweak interactions has been tested with very high precision, its scalar sector still evades experimental confirmation. In particular, it is an open question if the proper theory should contain one or more physical Higgs bosons. Since the CP nature of Higgs particles is a model dependent feature $`^\mathrm{?}`$ its determination would not only provide information concerning the mechanism of CP violation but would also restrict possible extensions of the SM of electroweak interactions and therefore reveal the structure of fundamental interactions beyond the SM. A muon collider with transversely polarized beams is the only place where CP properties of a second generation fermion Yukawa coupling can be probed. This is the subject of our more complete analysis $`^\mathrm{?}`$ which is summarized in this talk. We follow the line of our previous works $`^\mathrm{?}`$ where we have tried to unveil the CP-nature of Higgs bosons in a model-independent way.
## 2 Production of Higgs Bosons
The attractive possibility of s-channel Higgs boson production at a muon collider has been discussed before $`^\mathrm{?}`$ together with the possibility of the measurement of CP violation in the muon Yukawa couplings $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$. The latter is based on the fact that in any muon collider design $`^{\mathrm{?},\mathrm{?}}`$ there is a natural beam polarization of the order of 20% $`^\mathrm{?}`$ that allows for the rare possibility of direct Higgs boson production with known polarization of the initial state particles.
The cross section for the Higgs boson resonance production, $`\mu ^+\mu ^{}R`$, depends on the transverse $`P_T^\pm `$ and longitudinal $`P_L^\pm `$ beam polarizations and the $`\overline{\mu }r\mathrm{exp}(i\delta \gamma _5)\mu `$ muon Yukawa coupling in the following way:
$$\sigma _S(\zeta )=\sigma _S^0\left[1+P_L^+P_L^{}+P_T^+P_T^{}\mathrm{cos}(2\delta +\zeta )\right]$$
(1)
where $`\zeta `$ is the angle in the transverse plane between the beam polarizations and $`\sigma _S^0`$ is the unpolarized cross section. We stress that only the transverse polarization term is sensitive to $`\delta `$ of the muon Yukawa coupling. Since it is proportional to the product of the transverse polarizations, it is essential to have large $`P_T^+`$ and $`P_T^{}`$, as obtained by applying stronger cuts while selecting muons from the decaying pions initially produced (which, however, causes a reduction of luminosity). To compensate, a more intensive proton source or the ability to repack muon bunches will be needed. Another speculative option, would be high, up to 50%, polarization obtained by a phase-rotation technique $`^\mathrm{?}`$ which would lead to less luminosity reduction.
While varying $`\zeta `$, one can observe a maximum at $`\zeta =2\delta `$ and a minimum at $`\zeta =\pi 2\delta `$. Thus, studying $`\zeta `$ dependence is essential for resolving the $`\delta `$ value. A muon collider offers the unique possibility of a setup which in a natural way provides a scan over different $`\zeta `$ values. We will not discuss this option here. Our results will correspond to a configuration with four fixed $`\zeta `$ values: 0,90,180 and 270 degrees. Even though this cannot be accomplished experimentally, due to the spin precession in the accelerator ring, it can be well approximated by a simple but realistic setup $`^\mathrm{?}`$ that yields the same results as the fixed $`\zeta `$ analysis at the expense of 50% luminosity increase.
In order to illustrate the ability to reject different Higgs boson CP scenarios we can assume that the measured data is mimicked by the SM Higgs boson. For given luminosity $`L`$ and total polarization $`P`$ for each of the beams, we can place 1 and 3 $`\sigma `$ limits on the $`\delta `$ value for the observed resonance, assuming the $`\delta =0`$ SM is input. The limits for Higgs boson masses of 110 and 130 GeV are presented in table 1. For the expected yearly luminosity of $`L=150\mathrm{pb}^1`$, even several years of running at the natural 20% polarization would be insufficient for useful limits. However, 1$`\sigma `$ limits for the $`P=39\%`$ option (with reduced $`L`$) do give a rough indication of the CP nature of the resonance. 3$`\sigma `$ limits in 110-130 GeV mass range require either $`>40\%`$ polarization or $`<50\%`$ luminosity loss. (The requirements are less stringent for a 110 GeV Higgs boson.) We stress that there is no other way the measurement of the muon Yukawa $`\delta `$ can be done and that operation in the transverse polarization mode should not interfere with most of the other studies. For a heavier resonance, operation of a muon collider as an s-channel Higgs boson factory is justified only if the branching ratio $`BR(R\mu ^+\mu ^{})`$ is enhanced. Then, the analysis sketched above applies as well. Such enhancement arises in the Minimal Supersymmetric Standard Model (MSSM) at large $`\mathrm{tan}\beta `$. If the pseudoscalar mass is large ($`m_A>300\text{GeV}`$), the $`H`$ and $`A`$ masses will be similar. The increasing degeneracy with increasing $`\mathrm{tan}\beta `$ is illustrated for $`m_A=400\text{GeV}`$ in Fig. 1, assuming squark masses of $`1\text{TeV}`$ and no squark mixing and a beam energy spread of $`R=0.1\%`$. Since the total widths of the $`H`$ and $`A`$ are substantial ($`>1\text{GeV}`$) for the $`m_A`$ and $`\mathrm{tan}\beta `$ values being considered, it is not guaranteed that we will be able to separate the peaks. The figure shows that we are able to observe two separate peaks (the $`A`$ peak being at lower mass than the $`H`$ peak) for moderate $`\mathrm{tan}\beta <6`$. But, for higher $`\mathrm{tan}\beta `$ values the peaks begin to merge; for $`\mathrm{tan}\beta >8`$, $`|m_Hm_A|<11\text{GeV}`$ and one sees only a single merged peak. The picture changes if squark mixing is substantial; for instance, for $`m_A=300\text{GeV}`$, squark masses of $`1\text{TeV}`$ and large squark mixing ($`A_t=A_b=3\text{TeV}`$), the $`H`$ and $`A`$ peaks actually cross at $`\mathrm{tan}\beta 5`$. It would be crucial to distinguish such a case from a single CP-violating Higgs boson which may appear e.g. in the MSSM $`^\mathrm{?}`$. Table 2 illustrates the very distinct event number pattern as a function of $`\zeta `$ that would yield the needed discrimination. The event rate for any single, CP conserving or CP violating Higgs boson, has a minimum and maximum as a function of $`\zeta `$. In contrast, overlapping CP-even and CP-odd resonances result in a pattern independent of $`\zeta `$.
In Fig. 2 we plot for $`m_A=300\mathrm{GeV}`$ $`\mathrm{\Delta }\chi ^2`$ <sup>a</sup><sup>a</sup>aTo test a model A against B we introduce $`\mathrm{\Delta }\chi ^2=\mathrm{\Sigma }_i\frac{(N_i^AN_i^B)^2}{N_i^B}`$, where $`N_i^{A/B}`$ denotes the number of events in the $`i`$th bin calculated within the model $`A/B`$; see Grzadkowski, Gunion and Pliszka $`^\mathrm{?}`$ for details. The background generated by $`\gamma ^{}`$ and $`Z`$ exchange is taken into account. obtained for the four polarization-luminosity options (I)-(IV): (I) $`P=0.2`$, $`L=3.0\mathrm{fb}^1`$, (II) $`P=0.39`$, $`L=1.5\mathrm{fb}^1`$, (III) $`P=0.48`$, $`L=1.5\mathrm{fb}^1`$, (IV) $`P=0.45`$, $`L=3.0\mathrm{fb}^1`$. We emphasize that (for $`R=0.1\%`$) options (I) and (II) do not require over-design of the proton source. The $`\mathrm{\Delta }\chi ^2`$ plots show that good discrimination is obtained even for option (I) once $`\mathrm{tan}\beta >10`$. Option (II) would be needed for good discrimination if $`\mathrm{tan}\beta 5`$.
We have found that for simple MSSM test cases with $`m_A=300400`$ GeV and $`\mathrm{tan}\beta >8`$ (for which we cannot see separate resonance peaks) even natural 20% polarization will allow us to distinguish two overlapping resonances from any single one at more than the 3$`\sigma `$ level. Higher polarization will allow for a precise measurement of the relative contribution from the CP-even and the CP-odd component.
## 3 Summary and Conclusions
We have presented results of a realistic study of measuring the CP properties of the muon Yukawa couplings in Higgs boson production at a muon collider with transversely polarized beams. We have found that transverse polarization is essential for determining the CP nature of the muon Yukawa couplings. In particular, a collider with $`P40\%`$ and at least 50% of the original luminosity retained will ensure that the CP nature of the produced scalar resonance will be revealed.
## Acknowledgments
The author is grateful to the organizers of the XXXVth Rencontres de Moriond, on “Electroweak Interactions and Unified Theories” for creating a very warm and inspiring atmosphere during the meeting. He thanks J.F. Gunion and J. Pliszka for a collaboration upon which the presented talk was based and S. Geer, R. Raja and R. Rossmanith for helpful conversations on experimental issues. This work was supported in part by the U.S. Department of Energy, the U.C. Davis Institute for High Energy Physics, the State Committee for Scientific Research (Poland) grant No. 2 P03B 014 14 and by Maria Sklodowska-Curie Joint Fund II (Poland-USA) grant No. MEN/NSF-96-252.
## References
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# Optical properties and electronic structure of 𝛼'-Na1-xCaxV2O5
## I Introduction
$`\alpha ^{}`$-Na<sub>1-x</sub>Ca<sub>x</sub>V<sub>2</sub>O<sub>5</sub> belongs to the larger group of $`\alpha ^{}`$ vanadium pentoxides, with the chemical formula AV<sub>2</sub>O<sub>5</sub> (A= Li, Na, Ca, Mg, etc.) . Their structure is remarkably similar to that of the parent compound V<sub>2</sub>O<sub>5</sub>, which consists of layers of square pyramids of O surrounding a V<sup>+5</sup> ion. The layers are kept together via weak forces, which account for the easy cleavage of this oxide along (001). The basic building blocks forming the V<sub>2</sub>O<sub>5</sub> layers are parallel ladders of VO<sub>5</sub> pyramids. In AV<sub>2</sub>O<sub>5</sub> the A atoms enter the space between the layers and act as electron donors for the V<sub>2</sub>O<sub>5</sub> layers. In the case of $`\alpha ^{}`$-NaV<sub>2</sub>O<sub>5</sub>, every doped electron is shared between two V atoms. As a result the average valence of the V-ions corresponds to V<sup>+4.5</sup>. X-ray diffraction indicates that at room temperature all V-ions are in the same mixed valence state .
Partial substitution of the Na<sup>+</sup> with Ca<sup>2+</sup> leaves the $`\alpha ^{}`$ crystal structure intact, but alters the relative abundance of V<sup>4+</sup> and V<sup>5+</sup>, $`N^{4+}:N^{5+}=(1+x):(1x)`$. In this paper we report spectroscopic ellipsometry measurements on $`\alpha ^{}`$-Na<sub>1-x</sub>Ca<sub>x</sub>V<sub>2</sub>O<sub>5</sub> (x=0, 0.06, 0.15 and 0.20), in the energy range 0.8-4.5 eV. We employ the dependence of the optical spectra on the $`N^{4+}:N^{5+}`$ ratio to identify the main components in the optical spectra, which in turn we use to reveal the electronic structure of this material.
## II Details of sample preparation and experimental setup
The crystals (CR8, 45008, and 45010) had dimensions of approximately 2, 1 and 0.3 mm along the a, b, and c axes respectively. The samples 45008 and 45010 were prepared from NaVO<sub>3</sub> flux . In a first step a mixture of Na<sub>2</sub>CO<sub>3</sub> and V<sub>2</sub>O<sub>5</sub> is heated up to 550 C in air to form NaVO<sub>3</sub>. In a second step the NaVO<sub>3</sub> is mixed with VO<sub>2</sub> in the ratio of 8:1 and then heated up to 800 C in an evacuated quartz tube and cooled down at a rate of 1 K per hour. The excess NaVO<sub>3</sub> was dissolved in water. Then the doped samples were produced by substituting in the first step Na<sub>2</sub>CO<sub>3</sub> by CaCO<sub>3</sub>. The chemical composition of the samples has been determined using Energy Dispersive X-ray Fluoresence microprobe measurements. The results showed that the real Ca content of some samples was smaller that the nominal one (with a factor of 0.75), and that position dependent variations of the Na stoichiometry are below 2%. A standard spectroscopic ellipsometer was used to collect ellipsometric data in the range of 6000 to 35000 cm$`1`$ from the $`ab`$ planes of the crystals using two different crystal orientations, and to measure normal incidence reflectivity spectra of the $`ac`$ plane with the electric field vector along the $`c`$-direction.
## III Data collection and analyzes
We performed ellipsometric measurements on the (001) surfaces of the crystals both with the plane of incidence of the light along the a and the b axis. An angle of incidence $`\theta `$, of $`66^0`$, was used in all experiments. Ellipsometry provides directly the amplitude and phase of the ratio of the reflectivity coefficients of $`p`$\- and $`s`$-polarized light $`r_p(\omega )/r_s(\omega )`$. For an anisotropic crystal with the three optical axes ariented along the surface normal ($`p`$), perpendicular to the plane of incidence (s), and along the intersection of the plane of incidence and the surface ($`p`$), this ratio is related to the dielectric tensor elements along these three directions ($`ϵ_p`$, $`ϵ_s`$, and $`ϵ_p`$) according to the expression:
$$\begin{array}{c}\frac{r_p}{r_s}=\frac{\left[\sqrt{ϵ_pϵ_p}\mathrm{cos}\theta \sqrt{ϵ_p\mathrm{sin}^2\theta }\right]\left[\mathrm{cos}\theta +\sqrt{ϵ_s\mathrm{sin}^2\theta }\right]}{\left[\sqrt{ϵ_pϵ_p}\mathrm{cos}\theta +\sqrt{ϵ_p\mathrm{sin}^2\theta }\right]\left[\mathrm{cos}\theta \sqrt{ϵ_s\mathrm{sin}^2\theta }\right]}\hfill \end{array}$$
(1)
To extract the dielectric constant from the ellipsometric parameters we proceed in two steps: First the pseudo-dielectric functions along the optical axes are extracted from the ellipsometric data using the inversion formula
$$ϵ_p^{ps}=\mathrm{sin}^2\theta \left(1+\mathrm{tan}^2\theta \left(\frac{1r_p/r_s}{1+r_p/r_s}\right)^2\right)$$
(2)
For isotropic crystals this expression provides the dielectric function directly. The pseudo-dielectric function is close to the dielectric tensor elements along the intersection of the plane of incidence and the crystal surface. A biaxial crystal like $`\alpha ^{}`$-NaV<sub>2</sub>O<sub>5</sub> has three complex dielectric functions, $`ϵ_a`$, $`ϵ_b`$ and $`ϵ_c`$ along each optical axis, and an ellipsometric measurement involves all three tensor components of the dielectric matrix. In addition to the pseudo-dielectric functions displayed Fig.1a, $`ϵ_c(\omega )`$ is required. No $`ac`$ or $`bc`$ crystal-planes were available large enough to do ellipsometry with our setup. We therefore measured the $`c`$-axis reflectivity (Fig. 1b) of the $`bc`$-plane of the pristine material (sample CR3). The spectrum contains no (or very weak) absorption peaks in the measured frequency range, as it was reported earlier , providing a very reliable determination of the dielectric function $`ϵ_c`$ using Kramers-Kronig analysis. Due to the absence of strong resonances, $`ϵ_c`$ has a minor influence on the recorded ellipsometric spectra. In Fig. 1c the optical conductivity is displayed taking into account all corrections due to the anisotropy. We see that the conversion from $`ϵ^{ps}(\omega )`$ to $`ϵ(\omega )`$ has indeed a rather small effect on the spectra. In essence it leads to a factor 0.5 re-scaling of $`\sigma (\omega )`$.
The data are in general agreement with previous results using Kramers-Kroning analysis of reflectivity data. Along the a-direction we observe a peak at 0.9 eV with a shoulder at 1.4 eV, a peak at 3.3 eV and the slope of a peak above 4.2 eV, outside our spectral window. A similar blue-shifted sequence is observed along the b-direction. The 1eV peak drops rather sharply and extrapolates to zero at 0.7 eV. However, weak absorption has been observed within the entire far and mid-infrared range. The strong optical absorbtion within the entire visible spectrum causes the characteristic black appearance of this material.
The peak positions appear to be doping independent, but the striking observation is that the intensity of the peaks depends strongly on doping. In particular, the measured intensity of the 1 eV peak for the a axis is directly proportional to $`1x`$ (Fig. 3). The 1 eV peak for the b axis shows also a decrease upon doping.
## IV Main elements of the electronic structure
Before entering the interpretation of the data, we need to discuss in some more detail the main elements of the electronic structure of these compounds. The basic building block of the crystal structure of $`\alpha ^{}`$-NaV<sub>2</sub>O<sub>5</sub> is formed by VOV dimers. These dimers form the rungs of quasi one-dimensional ladders. The V-ions forming the rungs are bonded along the ladder direction via oxygen ions.
The backbone of the electronic structure is formed by the oxygen $`2p`$ and V$`3d`$ states. Photoelectron spectroscopy has provided crucial information on the occupied electronic levels: The oxygen $`2p`$ states have the lowest energy (the highest binding energy). They form a band about 4 eV wide, which is fully occupied. The occupied part of the V $`3d`$ states is located about 3 eV above the top of the oxygen bands. Due to ligand field splittings the V $`3d`$ manifold is spread over a range of at least 3 eV. The $`3d_{xy}`$ is of the V<sup>4+</sup> is occupied with one electron. The unoccupied $`d_{xz}`$ and $`d_{yz}`$ levels have an energy at least 1 eV higher. These in turn are located about 2 eV below the $`d_{x^2y^2}`$ and the $`d_{z^2}`$ levels. The relevance of the O 2p bands is that they provide a path for virtual hopping processes between the V-sites. The coupling between V-sites is through virtual hopping via $`\pi `$-bonded O $`2p_y`$ states, indicated schematically in Fig.5b. The effective hopping parameter between V-sites is $`t_{}0.3eV`$ on the same rung, and $`t_{}0.2eV`$ along the legs of the ladder. The number of electrons is one per pair of V-atoms. Approaching the ladders as a linear array of rungs, weakly coupled along the direction of the ladder, results in a model of electrons occupying a narrow band of states formed by the anti-symmetric combination of the two V 3d-states forming the rungs, hereafter referred to as V-V bonding levels.
Hence, we see that the basic building block are the pairs of V$`3d`$ states, together forming the rungs of the ladders. The essential charge and spin degrees of freedom of a single rung are identical to the Heitler-London model of the H$`{}_{}{}^{+}{}_{2}{}^{}`$ ion, with the V$`3d_{xy}`$ states playing the role of the H $`1s`$ states. The relevant Hamiltonion is
$$\begin{array}{c}H=t_{}_\sigma \left\{d_{L\sigma }^{}d_{R\sigma }+d_{R\sigma }^{}d_{L\sigma }\right\}\hfill \\ +\frac{\mathrm{\Delta }}{2}_\sigma \left\{n_{L\sigma }n_{R\sigma }\right\}+U\left\{n_Ln_L+n_Rn_R\right\}\hfill \end{array}$$
(3)
where $`d_{L(R),\sigma }^{}`$ creates an electron in the lefthand (righthand) $`d_{xy}`$ orbital on the rung, and $`\sigma `$ is the spin-index. The bias potential $`\mathrm{\Delta }`$ between the two V-sites accounts for a possible left/right charge imbalance. The most relevant states for the ground state are $`d_{L,xy}`$ and $`d_{R,xy}`$. Pure NaV<sub>2</sub>O<sub>5</sub> contains one electron per rung in the ground state. Important in the present discussion are the eigenstates and energies of a rung with 0, 1 or 2 electrons. The eigenstates and energies are listed in Table I.
In Fig. 4 the level diagram is displayed. In this representation $`N_e=0(2)`$ corresponds to the one electron removal(addition) states, for noninteracting electron picture indicated as the ”occupied” (”empty”) states. In Fig.5 we display the same information represented as the the one electron removal and addition spectral function. For non-interacting electrons this represents the occupied (left) and unoccupied (right) states.
In the absence of electron-electron interactions, here represented by the on-site Hubbard repulsion parameter $`U`$, the Fermi energy would be located in the middle of the bonding band, resulting in a metallic conductor. In Figs.4 and 5 we have adopted $`U=4`$ eV. The model now predicts a gap of order $`E_{CT}`$1 eV. The Fermi energy is located within this gap. The fact that these materials are insulating therefore is associated with the large on-site Hubbard interaction.
The excitation of an electron across the gap involves a change of occupancy of two of the rungs: The final state has one empty and one doubly occupied rung. It is important in this context, that the two electrons $`|^{3,1}LR>`$ are in a correlated state: In the limit $`U\mathrm{}`$ one electron is located on the lefthand V-atom and the other on the righthand V-atom.
## V Discussion of the experimental spectra
The 0.9 eV peak marks the fundamental gap of the optical spectrum. The interpretation of this peak is still subject of a scientific controversy. Several interpretations have been put forward
1. Transitions between linear combinations of V 3d<sub>xy</sub>-states of the two V-ions forming the rungs. In Refs. and even and odd combinations were considered. The 0.9 eV peak in $`\sigma _a(\omega )`$ (peak A) would then correspond to the transition from V-V bonding to antibonding combinations on the same rung (Fig.5b). In Ref. this model was extended to allow lop-sided linear combinations of the same orbitals, so that the 0.9 eV peak is a transition between left- and right-oriented linear combinations.
2. On-site d-d transitions between the crystal field split levels of the V-ions. Because V<sup>5+</sup> has no occupied 3d-levels, such processes involve the V<sup>4+</sup> ions.
3. Transitions between the on-rung V 3d$`{}_{}{}^{+}{}_{xy}{}^{}`$ bonding combination and final states of $`d_{xz}`$ and $`d_{yz}`$ character.
Optical transitions having values below 2 eV were also seen in V<sub>6</sub>O<sub>13</sub> and VO<sub>2</sub>. In V<sub>2</sub>O<sub>5</sub> they have very small intensities, and were attributed to defects . The last two assignments are motivated by the fact that in $`\alpha ^{}`$-Na<sub>1-x</sub>Ca<sub>x</sub>V<sub>2</sub>O<sub>5</sub> the optical selection rules allow on-site d-d transitions by virtue of the low point-symmetry at the V-sites. To determine which one of these assignments is true, we have measured the doping dependence of the 0.9 eV peak in $`\sigma _a(\omega )`$ in Ca-substituted $`\alpha ^{}`$-NaV<sub>2</sub>O<sub>5</sub> (Fig.2a). Because Ca is divalent, substituting Na with Ca has the effect of increasing the average density of V<sup>4+</sup> ions.
In the local d-d scenario the intensity of the on-site V 3d-3d transitions would be proportional to the density of V<sup>4+</sup> ions. As a result the intensity of on-site d-d transitions increases upon substituting Na with Ca. In Fig. 3 we display the experimentally observed doping dependence of the intensity together with the theoretical expectation within this scenario. Clearly the experimental intensity of the 0.9 eV peak in Fig. 2 behaves opposite to the expected behaviour of $`dd`$ transitions. This definitely rules out item number 2 of the above list.
This also rules out item number 3 presented in Ref.: if the transition from the V 3d$`{}_{}{}^{+}{}_{xy}{}^{}`$ bonding combination to the $`d_{xz}`$ and $`d_{yz}`$ orbitals involves mainly transitions among orbitals at the same site, the same argument as for item 2 applies. If it involves mainly transitions between molecular orbitals formed by different sites on the same rung, the transition from the V 3d$`{}_{}{}^{+}{}_{xy}{}^{}`$ bonding combination to the antibonding V 3d$`{}_{}{}^{}{}_{xy}{}^{}`$ would still be the dominant transition.
To explain the intensity decrease of about x% upon substituting x% of the Na ions with Ca (see Fig.3), let us have a look at the many-body nature of the ground state and the excited state properties of NaV<sub>2</sub>O<sub>5</sub>. The rungs which are occupied with two electrons due to the Ca-doping will be in the many-body ground state (see Table I) $`|^1\stackrel{~}{LR}`$. Due to the dipole selection rules, the optical excitation with the electric field along the rung-direction is $`|^1\stackrel{~}{LR}|^1\stackrel{~}{O}`$. The energy to make this excitation is $`\frac{1}{2}U+\sqrt{U^2/4+4t_{}^2}`$. Hence the effects of Ca doping are (i) to remove the peak at $`E_{CT}1eV`$ for the rungs receiving the extra electron, and (ii) to place a new peak at an energy $`U4eV`$. Hence, the observed x% decrease of intensity of the $`|B|A`$ transition peak for the x% Ca doped sample is in excellent quantitative agreement with the expected value.
Using the tightbinding f-sum rule $`\sigma (\omega )𝑑\omega =(ed/\mathrm{})^2\pi t_{}/(2V)d_{L\sigma }^{}d_{R\sigma }+\text{HC}`$ the intensity of the $`|^1\stackrel{~}{LR}|^1\stackrel{~}{O}`$ peak relative to the 0.9 eV peak of the singly occupied rungs is (assuming $`\mathrm{\Delta }=0`$ for a doubly occupied rung):
$$\frac{I(0)+I(2)}{2I(1)}=\frac{\alpha _K\beta _K}{uv}=\sqrt{\frac{1}{1+(U/4t)^2}}\frac{4t_{}}{U}$$
(4)
With the parameters relevant to NaV<sub>2</sub>O<sub>5</sub> this implies that the $`|^1\stackrel{~}{LR}|^1\stackrel{~}{O},|^1\stackrel{~}{E}`$ transitions have factor 2-4 smaller spectal weight than the $`|B|A`$ transition. Hence we conclude that only the assigment of item number 1 is consistent with our data: The 1eV peak in $`\sigma _a(\omega )`$ is the on-rung $`|^2B_\sigma >|^2A_\sigma >`$ transition with an excitation energy $`E_{CT}\sqrt{4t_{}^2+\mathrm{\Delta }^2}`$.
The 1.1 eV peak in $`\sigma _b(\omega )`$ (peak B) involves transitions between neighboring rungs along the ladder. In the non-interacting model ($`U=0`$) this would correspond to a Drude-Lorentz optical conductivity centered at $`\omega =0`$, with a spectral weight $`_0^{\mathrm{}}\sigma _b(\omega )𝑑\omega =(e/\mathrm{})^2t_{}\pi b(2ac)^1`$. As a result of the correlation gap in the density of states, indicated in Fig. 5a, the optically induced transfer of electrons between neighboring rungs results in a final state with one rung empty, and a neighboring rung doubly occupied, in other words, an electron hole pair consisting of a hole in the band below E<sub>F</sub>, and an electron in the empty state above E<sub>F</sub> indicated in Fig. 5a. This corresponds to the process
$$2|^2B_\sigma >|^{3,1}LR>+|0>$$
(5)
Note that the final state wavefunction is qualitatively different from the on-rung bonding-antibonding excitations considered above, even though the excitations are close in energy: it involves one rung with no electron, and a neighboring rung with one electron occupying each V-atom. The energy of this process is approximately $`2t_{}+\delta V`$, where $`\delta V`$ represents the increase in Coulomb interaction by bringing two electrons together on the same rung. Since the distance betweeen the electrons changes from about 5.0$`\AA `$ to 3.4$`\AA `$, and taking into account a screening factor $`ϵ6`$, we estimate that $`\delta V0.2eV`$.
This value of $`\delta V`$ corresponds closely to the difference in peakpositions along the $`a`$ and $`b`$ directions. According to this interpretation the absorption at 1.1 eV along $`b`$ corresponds to the creation of a free electron and hole, capable of carrying electrical currents. The on-rung excitation at 0.9 eV along the $`a`$-direction is a localized (charge neutral) excitation, in other words an exciton. In this case the energy of the exciton involves the states of a single electron only, whereas the free carrier states involve many-body interactions.
Doping with Ca creates doubly occupied rungs, whose ground energy is not $`2|^2B_\sigma >`$ but $`|^1\stackrel{~}{LR}`$. Consequently, the electrons on these rungs will not be involved in the processes of Eq. 5, thus decresing the intensity of the B peak upon doping, as seen from Fig. 2b.
The room temperature crystal structure has four V-atoms per unit cell, organized in ladders with up and down oriented apical oxygens alternating along the a-direction, resulting in a double degeneracy of the electronic states discussed above. The coupling between adjacent ladders lifts the degeneracy of these states, resulting in a ”Davidov” splitting of the peaks A and B. This can create the two additional ”shoulders” in $`\sigma (\omega )`$ at 1.4 and 1.7 eV for peak A and B respectively.
With ARPES it has been observed that an energy of 3 eV separates the V$`3d`$ band from the O$`2p`$. We therefore attribute the peaks at 3.3 eV in $`\sigma _a(\omega )`$ and the peak at 4 eV in $`\sigma _b(\omega )`$ to transitions of the type
$$|^2B_\sigma >|\underset{¯}{2p}^1LR>$$
(6)
where the $`\underset{¯}{2p}`$ hole is located on the oxygen on the same rung for peak A, and inbetween the rungs for peak B. This is further supported by previous optical measurements on V<sub>2</sub>O<sub>5</sub> showed a peak at about 3 eV. In V<sub>2</sub>O<sub>5</sub> all V-ions have a formal V$`3d^0`$ configuration, hence the 3 eV peak can not be attributed to d-d transitions. However, the O$`2p`$V$`3d`$ transitions should appear at approximately the same photon energy as in NaV<sub>2</sub>O<sub>5</sub>, which further supports our assigment of the 3 eV peak in NaV<sub>2</sub>O<sub>5</sub> to O$`2p`$V$`3d`$ transitions.
## VI Conclusions
In conclusion, we have measured the dielectric function along the a and b axes of Ca<sub>x</sub>Na<sub>1-x</sub>V<sub>2</sub>O<sub>5</sub> for x=0, 0.06, 0.15 and x=0.20. The 0.9 eV peak in $`\sigma _a(\omega )`$ was shown experimentally to be a bonding-antibonding transition inside the V<sub>2</sub>O rung and not a vanadium d-d transition due to crystal field splitting. We identified the 3.3eV peak in $`\sigma _a(\omega )`$ as the transition from the oxygen orbitals to the antibonding one. This strongly supports the notion, previously expressed in Refs that NaV<sub>2</sub>O<sub>5</sub> is an insulator due to a combination of three factors: A crystal field splitting, an on-site Hubbard interaction, and an on-rung bonding-antibonding splitting of the two V$`3d_{xy}`$ orbitals, each of which is large compared to the inter-rung hopping parameter.
## VII Acknowledgements
We like to thank H. Bron for his assistence with the chemical analysis of the crystals, and prof. J.T.M. de Hosson for making available the microprobe equipment. This investigation was supported by the Netherlands Foundation for Fundamental Research on Matter (FOM) with financial aid from the Nederlandse Organisatie voor Wetenschappelijk Onderzoek (NWO).
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# A catalogue of symbiotic stars
## 1 Introduction
Symbiotic stars are interacting binaries, in which an evolved giant transfers material to much hotter, compact companion. In a typical configuration, a symbiotic binary comprises a red giant transferring material to a white dwarf via a stellar wind.
Amongst the evidence for this predominant mass-transfer mechanism is the fact that ellipsoidal light variations, characteristic of tidally distorted stars, are rarely observed for symbiotic stars. Thus far, only two systems, T CrB (\[46 \]) and CI Cyg (\[m229 \]), are known to have the ellipsoidal light variations of a distorted giant. In some symbiotic systems, the red giant is replaced by a yellow giant or a carbon star, and the white dwarf by a main-sequence or neutron star.
Most symbiotic stars ($`80\%`$) contain a normal giant star and these, based on their near-IR colours (showing the presence of stellar photospheres, $`T_{\mathrm{eff}}`$ 3000 – 4000 k), are classified as S-type systems (stellar). The remainder contain Mira variables and their near-IR colours indicate temperatures of $`1000`$k, giving away the presence of warm dust shells; these are classified as D-type systems (dusty). The IR type seems to be dependent on the orbital separation of the components. For large separations (long periods), the cool star seems able to evolve to the Mira stage and produce a dust shell that enshrouds the system; for smaller separations (shorter periods), we deal with normal giants. For a detailed review of symbiotic stars, we refer the reader to \[358 \].
Two catalogues of symbiotic stars have been published. The first was by David Allen in 1984 (\[30 \]); it included 129 symbiotic stars and 15 possible symbiotic objects with a concise summary of available observational data, finding charts and optical spectra for the most of listed objects. The second catalogue was by Scott Kenyon in 1986 (\[31 \]); it included 133 symbiotic stars and 20 possible symbiotic objects, as well as tables describing selected observational properties of all the objects and a spectroscopic summary of a selected sample. Kenyon’s work also provides the reader with an excellent overview and bibliography of selected symbiotic stars.
Since 1986, a number of papers have presented surveys of large samples of symbiotic stars (e.g. \[7 ; vwds93 ; ibm94 ; sih95 ; 79 ; 3 ; 313 \]) and in-depth investigations (\[ii88 ; 37 ; ibesm93 ; sk94 ; go96 \] for AX Per alone). New stars have been included in the family of symbiotic stars each year and, at the same time, better data have been collected and better data analysis has been performed for a number of well-known symbiotic stars. The aim of this work is to present the symbiotic star research community with a comprehensive compilation of existing data collected from a number of astronomical journals, electronic databases and unpublished data resources. For many objects a new classification has been necessary: some have been confirmed as symbiotic stars; some have been rejected; some new objects have been added. Our catalogue lists 188 symbiotic stars and 28 objects suspected of being symbiotic stars.
## 2 Classification criteria
The optical spectra of symbiotic stars are characterized by the presence of absorption features and continuum, as appropriate for a late-type M giant (often a Mira or semi-regular, SR, variable), and strong nebular emission lines of Balmer H i, He ii and forbidden lines of \[O iii\], \[Ne iii\], \[Ne v\] and \[Fe vii\]. Some symbiotics – the yellow symbiotic stars – contain K (or even G) giants or bright giants. The spectra of many symbiotic systems also show two broad emission features at $`\lambda \mathrm{\hspace{0.17em}6825}\mathrm{\AA }`$ and $`\lambda \mathrm{\hspace{0.17em}7082}\mathrm{\AA }`$. These features have never been observed in any other astrophysical objects — only symbiotic stars with high-excitation nebulae. For many years there was no plausible identification for these lines, but \[sch89 \] pointed out that the $`\lambda \lambda `$ 6825, 7082 lines are probably due to Raman scattering of the O vi $`\lambda \lambda `$ 1032, 1038 resonance lines by neutral hydrogen.
To classify an object as symbiotic star we adopted the following criteria:
1. The presence of the absorption features of a late-type giant; in practice, these include (amongst others) TiO, $`\mathrm{H}_2\mathrm{O}`$, CO, CN and VO bands, as well as Ca i, Ca ii, Fe i and Na i absorption lines.
2. The presence of strong emission lines of H i and He i and either
* emission lines of ions with an ionization potential of at least 35 eV (e.g. \[O iii\]), or
* an A- or F-type continuum with additional shell absorption lines from H i, He i, and singly-ionized metals.
The latter corresponds to the appearance of a symbiotic star in outburst.
3. The presence of the $`\lambda `$ 6825 emission feature, even if no features of the cool star (e.g. TiO bands) are found.
Our adopted criteria represent a compromise: a collection of the classification criteria proposed in the past 70 years (see Kenyon 1986 for details), based on the examples of well-studied and widely accepted symbiotic objects. We believe that such an approach is appropriate, especially given that symbiotic stars are variables with timescales often exceeding a dozen years and that — as Kenyon very sensibly noted — “every known symbiotic star has, at one time or another, violated all the classification criteria invented”. For those who would prefer additional or different definitions, we give the highest ionization potential ever observed in the optical and UV (for objects that have been observed at least once with the International Ultraviolet Explorer — IUE). We also comment on all objects for which our classification may not seem obvious (e.g. V934 Her, which some readers may consider to be symbiotic, but which in our catalogue is classified as a suspected symbiotic star).
## 3 The catalogue
The main catalogue is presented in Table 1. This table includes collated data for all the symbiotic stars we know of. Note that a colon indicates an uncertain measurement or an estimate. Stars are ordered by right ascension (R. A.) for the epoch J2000.0. The content of each column is described below.
1: Symbiotic star catalogue number. A star symbol, if present here, means that there is a classification note and/or comment for the given object. We would still advise the use of the symbiotic (or suspected symbiotic) star name, as given in the second column of Tables 1 and 2, and not the object’s catalogue number.
2: Symbiotic star name. If, for a given object, a variable star name exists, then it was chosen; otherwise, the name used most often in the literature was adopted.
3,4: R. A. and declination (J2000.0), taken from radio VLA positions (\[m3 ; 9 ; 8 \]) if available, and if not from the SIMBAD database but corrected in a few cases were obvious mistakes have been spotted. If position was taken from somewhere else then comment is given in section 4.
5,6: Galactic coordinates (not included for extragalactic objects).
7,8: Magnitudes in the $`V`$ and $`K`$ filters, respectively. As most (if not all) symbiotics are variable, these values are arbitrary (usually the average of published measurements) just to give the general level of an object’s brightness.
9: IR type. If two IR types are given for one object, we supply references to both estimates in the notes.
10: Information on whether an IUE spectrum (or spectra) is (are) available (+) or not ($``$). The number of spectra can be readily obtained from SIMBAD (http://simbad.u-strasbg.fr/Simbad) and the spectra can be obtained from the IUE data archives (http://nssdca.gsfc.nasa.gov/ndads).
11: Information about whether an object was ever detected in X-rays. Plus (+) means a detection; minus ($``$) means that an object was observed but not detected and that only an upper limit is available. Most of the detections and upper limits came from ROSAT and were reported in \[275 \] but some have also been observed by Einstein (\[364 \]), EXOSAT and ASCA.
12: Highest ionization potential ever observed in the emission-line spectra of an object. The potential is given in electron volts (eV).
13: The symbiotic star catalogue number (repeated).
14: An estimate of the spectral type of the cool component, with references. Since the blue and visual spectral regions are often contaminated by the circumstellar nebula and/or the hot component, we have given priority to estimates made in the near-IR region and, in the case of multiple estimates, to those made at quiescence and/or near to inferior conjunction of the cool giant. The estimates based on the TiO bands are separated from those based on CO 2.3-$`\mu `$m bands by ‘/’. Also, if the cool component was reported to behave as a Mira (i.e. if Mira-like pulsations have been detected or the object’s position in the near-IR/IRAS colour diagram coincides with the region occupied by Mira variables) then it is noted in this column.
15: Radio observations of symbiotics. Detections or upper limits are given. In parentheses, the wavelength of observation is reported. If more than one detection has been reported, only one is included and the priority is given to the most extensive radio survey of symbiotics at 3.6 cm (\[7 \]). Other extensive surveys of symbiotic stars which were searched for radio detection include \[8 ; 11 ; 9 ; 10 ; m3 \].
16,17: IRAS fluxes at 12 and 25 $`\mu `$m. The fluxes are taken from pointed observations, if available, (\[13 \]) or from survey observations as listed in SIMBAD. If, for some object, there was no report of observations in either of the above two sources, but IRAS fluxes were available from somewhere else, then the reference to reported observations is given in the notes. The upper limits are marked with capital L. The IRAS number is listed in Table 8 (which contains different object names for the symbiotic and suspected symbiotic stars). If the number is not there then the reference to the reported observations is given in the section containing comments.
18: Major literature references to the object. A number indicates the reference number; abbreviations in parentheses mark the subject the reference was noted for: fc – finding chart, spc – optical spectrum, class – classification, parm – the latest or the most extensive and up-to-date discussion of an object.
In Table 2, we present data for objects suspected of being symbiotic stars. The order and content of the columns is exactly the same as in Table 1. The catalogue numbers of suspected symbiotic stars are preceded by the letter ‘s’ throughout the catalogue.
The next two tables include data on symbiotic and suspected symbiotic star orbits. In Table 3, we have put orbital photometric ephemerides, including information on the presence of eclipses if available, and references to every ephemeris estimate. In Table 4, the reader will find the orbital elements of twenty symbiotic stars as well as spectroscopic periods, radial velocity semi-amplitudes for the cool components, mass ratios, systemic velocities, eccentricities, times of inferior spectroscopic conjunctions of the giant, sizes of the giant orbits, mass functions and references to each orbital estimate.
In Table 5, we have collected the pulsation ephemerides for Miras in symbiotic and suspected symbiotic stars. Again, a reference to every estimate is given.
Table 6 includes known Hipparcos parallaxes for symbiotic stars.
Table 7 includes information on symbiotic and suspected symbiotic star flickering and outburst characteristics.
Tables 8 includes most of different names for symbiotic and suspected symbiotic stars. Symbiotic stars appear first, then suspected symbiotic stars follow. Objects are first listed by their catalogue number, then by the name (translated to SIMBAD nomenclature, if possible — the name by which the object is known in Table 1 or 2), then other names are given. The names are compatible with SIMBAD and general internet database nomenclature. In some cases, the catalogue name differs between Table 1 (or 2) and Table 8. This discrepancy is due the most commonly accepted name (Table 1 or 2) not following SIMBAD nomenclature (Table 8).
## 4 Classification notes and comments on particular objects
Symbiotic stars
004=SMC3 $`V`$ magnitude during outburst.
005=SMC N60 IR-type S –\[31 ; 35 \],D –\[30 \].
008=AX Per Incorrect coordinates given by \[30 ; 31 \].
009=V471 Per This star appears in previous symbiotic catalogs (\[30 ; 31 \]) as V741 Per. The correct name is V471 Per as given in the General Catalog of Variable Stars (\[357 \]).
010=o Ceti Cool component is Mira \[22 \] of type M2-7 III \[201 \], UV spectrum shows emission lines with ionization potentials up to 54.4 eV \[202 \] and in the optical spectrum there are emission lines of H i and He i \[204 \]. Preliminary orbit (orbital period = 400 yrs) \[202 \].
011=BD Cam Cool component is S giant of type S5.3 \[207 \]; UV spectra shows emission lines with ionization potential up to 77.5 eV \[317 \]. 24.76-day periodicity estimated from $`BVRI`$ photometry; pulsational origin has been suggested \[206 \].
016=UV Aur IR-type S –\[31 ; 10 \],D’ –\[30 ; 4 \].
017=V1261 Ori Cool component is S giant of type S4.1 \[211 \]; UV spectrum shows emission lines with ionization potential up to 77.5 eV \[212 \].
018=LMC1 IUE spectra described in \[35 \].
020=Sanduleak’s star In the optical spectrum, there is an emission feature at 6825 Å \[34 \]; moreover, there are emission lines with an ionization potential up to 108.8 eV \[34 ; 310 \] including lines of H i and Hei \[35 \]. The IUE spectra are described in \[35 \].
023=BX Mon IRAS data from \[13 \].
024=V694 Mon Object in permanent outburst \[293 \]; contains M3-5 giant \[223 ; m16 \]; optical spectrum shows emission lines of H i and He i with highly blueshifted ($`2000`$$`7000`$ km s<sup>-1</sup>) shell absorption \[m16 ; 223 ; m222 \] and emission lines of singly-ionised metals \[223 \] over an A-B type continuum \[223 \]. $`VK`$ magnitudes are appropriate for the outburst.
026=RX Pup Highly variable radio emission \[is94 \]. Nebula resolved at optical and radio wavelengths with a possible jet-like feature in the \[N ii\] line (\[not3 \] and references therein).
027=Hen 3$``$160 IRAS data from \[13 \].
028=AS 201 A spherical nebula detected in H<sub>α</sub> and \[N ii\] lines (\[not3 \] and references therein).
029=KM Vel Cool component is Mira \[22 ; m49 \] of M spectral type \[m1 \]; optical spectrum shows emission lines with ionization potential up to 41.0 eV \[m121 \] and emission lines of H i and He i \[121 \]. Finding chart in \[185 \] is incorrect and no other has been published.
032=SS73 29 IUE observations reported in \[66 \]. IRAS data from \[13 \].
033=SY Mus Spectropolarimetric orbit derived in \[not9 \].
034=BI Cru A bipolar nebula resolved in the optical (\[not3 \] and references therein) with the bipolar lobes and associated outflows perpendicular to the position angle of intrinsic scattering polarization \[65 \].
036=TX CVn Low ionization potential (IP$`{}_{\mathrm{max}}{}^{}=13.6`$ eV), but this is a confirmed symbiotic star (\[360 \]: combination spectrum of late B \+ early M, emission lines of H i and singly-ionised metals). Classification is also based on its light curve showing eruptions as in other symbiotics (with $`\mathrm{\Delta }m_{pg}`$ up to $`3^m`$). Since the 1970’s, the star is in permanent outburst with P-Cyg type spectrum.
038=Hen 3$``$828 IRAS data from \[13 \].
041=St 2$``$22 The SIMBAD database uses different name for this object: PN Sa 3$``$22.
043=V840 Cen IRAS data from \[118 \]. Finding chart available in \[237 \] where object is marked as star A \[186 \].
046=Hen 3$``$916 Finding chart in \[30 \] is wrong, object is 2mm ($`20\mathrm{"}`$) E of marked star \[52 \].
047=V704 Cen Cool component might be Mira \[22 \].
048=V852 Cen Cool component is Mira \[22 ; m49 \]; optical spectrum shows emission lines with ionization potential up to 100 eV \[57 ; 175 ; 3 \]; moreover, optical spectrum shows emission feature at 6825 Å \[m126 \] and emission lines of H i and He i \[3 \]. Bipolar nebula resolved in the optical (Southern Crab) (\[not3 \] and references therein).
050=V417 Cen An irregular nebula resolved at optical wavelengths (\[not3 \] and references therein).
055=HD 330036 This is a yellow symbiotic star; cool component is F5 giant or subgiant (\[62 \]). In UV, there are emission lines with ionization potential up to 77.5 eV \[62 \] and in the optical spectrum there are emission lines with ionization potential up to 54.4 eV \[3 ; 62 \]. IR-type D’ –\[\].
056=Hen 2$``$139 Only H i emission lines in spectrum according to \[30 \], but other emission lines (like \[O iii\]) are reported in \[m121 \].
058=AG Dra A secondary periodicity of $`355^\mathrm{d}`$ has been detected in the optical light curve and interpreted in terms of non-radial pulsation of the cool giant \[not1 \]. An orbital inclination, $`i120^{}`$, has been derived from spectropolarimetric observations \[not2 \].
060=V347 Nor An elliptical nebula resolved at optical wavelengths (\[not3 \] and references therein).
065=Hen 3$``$1213 IUE observations reported in \[90 \].
066=Hen 2$``$173 IRAS data from \[13 \].
067=Hen 2$``$176 IR-type S –\[3 \],D –\[30 ; 31 \].
068=KX TrA The finding chart in \[30 \] is wrong: the object is really 3mm ($`25\mathrm{"}`$) W of marked star, although tabulated coordinates are correct \[52 \].
071=CL Sco IRAS data from \[78 \].
073=V455 Sco An elliptical nebula possibly resolved in \[O iii\] (\[not3 \] and references therein).
074=Hen 3$``$1341 IUE observations reported in \[90 \]. Spectral signatures of collimated bipolar jets have been found during the 1999 outburst \[m230 \].
077=H 2$``$5 IR-type D –\[3 \],S –\[30 ; 31 \].
084=V2116 Oph Orbital period of 303.8 days is derived from the spin changes of the X-ray pulsar companion \[pbj99 ; cdd86 \].
088=M 1$``$21 $`VK`$ magnitudes – close and fainter companion also measured.
089=Hen$``$251 $`K`$-band spectrum is practically identical with that of the symbiotic Mira, RX Pup, as observed during the dust obscuration event, with strong dust continuum and weak CO 2.3-$`\mu `$m band \[ntt00 ; mik00 \].
092=RT Ser IRAS data from \[13 \].
093=AE Ara IRAS data from \[13 \].
094=SS73 96 IRAS data from \[13 \]. An axisymmetrical nebula resolved at radio wavelengths (\[not3 \] and references therein).
096=V2110 Oph IRAS data from \[13 \].
100=H 1$``$36 In \[5 \] there is an estimate of cool component spectral type M4-5 III based on TiO 7100Å band depth. However the spectrum of H 1$``$36 shown on their Fig.A1 does not show any absorption features or red continuum. A complex nebula resolved at optical and radio wavelengths (\[not3 \] and references therein). The only symbiotic star known to support an OH maser (\[ish94 \]).
101=RS Oph Bipolar nebula detected in radio range (\[not3 \] and references therein).
102=WRAY 16$``$312 IRAS and $`JHKL`$ colours confirm earlier suggestions \[m2 ; 22 \] that cool component of this system is a Mira \[\]. In the optical spectrum presented in \[30 \] there are emission lines with ionization potential up to 108.8 eV and moreover lines of H i and He i are present \[\]. IRAS data from \[13 \].
103=V4141 Sgr Classified as S-type in \[30 ; 4 \], but in the near-IR/IRAS colour diagrams it falls in the region occupied by symbiotic Miras \[3 \]. $`K`$-band spectrum shows strong CO 2.3-$`\mu `$m band consistent with an M6 giant \[ntt00 \]. Spectral type of cool component also estimated in \[179 ; 30 \] to be mid or late M.
105=AS 245 Classified as S-type in \[16 ; 3 \] but in the near-IR/IRAS colour diagrams it falls in the region occupied by symbiotic Miras \[\].
107=Bl 3$``$14 The finding chart in \[30 \] is good, but the coordinates are reported to disagree with the measured position: $`\alpha =17^\mathrm{h}52^\mathrm{m}06^\mathrm{s}.4,\delta =29^{}45^{}49^{\prime \prime }`$ (1950) \[52 \] (if this is right, our coordinates should also be corrected).
110=V745 Sco $`VK`$ magnitudes during decline from outburst \[253 \].
112=AS 255 IR-type S –\[30 ; 44 \],D –\[31 ; 43 \].
114=H2$``$34 Spectral type M5 is estimated by comparing ‘by eye’ the depths of TiO $`\lambda `$ 6180 and $`\lambda `$ 7100 $`\mathrm{\AA }`$ in the spectrum in Fig. 2 in \[3 \] with those of spectral standards.
115=SS73 117 IRAS data from \[13 \].
116=AS 269 This is a yellow symbiotic star, cool component is G-K giant \[3 ; m22 \]. In the optical spectrum there are emission lines with ionization potential up to 54.4 eV \[184 \].
118=SS73 122 IR-type D –\[13 \], others note only possible S type (\[30 ; 31 \]).
120=H 2$``$38 There was a report of a pulsational period of 433 days for this star in \[117 \], but this is a mistake and the reported number is the pulsation period of another symbiotic star: V366 Car (Hen 2$``$38). The spectral type of the cool component is estimated in \[16 \] to be M8.5.
122=Hen 3$``$1591 IR-type D –\[30 ; 4 \], S –\[31 ; m73 \].
124=Ve 2$``$57 Cool component is M star \[m16 \]. In the optical spectrum there are emission lines with ionization potential up to 35.1 eV or probably up to 54.4 eV \[m16 \].
125=AS 276 IR-type S –\[30 ; 44 \],D –\[31 \]. There is also a D’ classification in \[43 \], but it doesn’t look reliable.
128=V2506 Sgr IRAS data from \[108 \].
132=YY Her IRAS data from \[13 \].
133=V2756 Sgr Finding chart in \[185 \] is incorrect (\[m23 \]).
134=FG Ser $`K`$ magnitude during outburst. Coordinates taken from \[m233 \] – SIMBAD coordinates are not correct.
138=V4074 Sgr IUE observations reported in \[79 \].
139=V2905 Sgr IRAS data from \[78 \]. Spectral type of cool component also estimated in \[291 \] to be K/M.
146=V3811 Sgr Mis-identified in \[185 \] and in \[m42 \] (see \[m73 \]).
148=V3890 Sgr Cool component is M4-8 giant (\[m96 ; 261 ; 262 \]). In the optical spectrum there are emission lines with an ionization potential up to 361 eV \[261 \]. This object was earlier classified as recurrent nova with M type companion \[262 ; m96 \]. The spectrum is also presented in \[254 \].
156=FN Sgr IRAS data from \[78 \].
160=V1413 Aql Spectral type M4 estimated from the TiO $`\lambda `$ 7100 band depth as observed during mid-eclipse \[mik00a \].
162=Ap 3$``$1 Short description of optical spectrum is given in \[30 \]. The object was identified with the 2U 1907+2 X ray source \[m104 \] but so far there is no ROSAT detection, so this identification might not be correct.
166=BF Cyg IRAS data from \[13 \].
167=CH Cyg Complex nebula with jet-like features resolved at optical and radio wavelengths (\[not3 \] and references therein). Both the light curves and the radial velocity curves show multiple periodicities: a $`100^d`$ photometric period has been attributed to radial pulsation of the giant \[m107 \], while the nature of the secondary period of $`756^\mathrm{d}`$ also present in the radial velocity curve, is not clear \[not5 \]. There is controversy about whether the system is triple or binary \[93 \], and whether the symbiotic pair is the inner binary \[not6 \] or the white dwarf is on the longer orbit \[not7 ; not8 \].
169=HM Sge Mean $`K`$ magnitude during outburst. A complex nebula with possible jet-like features resolved at optical and radio wavelengths (\[not3 \] and references therein). The nebula is aligned with the binary orientation deduced from spectropolarimetry of the Raman scattered O vi lines \[m240 \].
170=Hen 3$``$1761 IRAS data from \[13 \]. IUE observations reported in \[90 \].
171=QW Sge IRAS data from \[13 \] although \[145 \] report no IRAS detection.
172=CI Cyg Coordinates from VLA observations \[m232 \].
174=V1016 Cyg A complex nebula with possible jet-like features resolved at optical and radio wavelengths (\[bang ; not3 \] and references therein).
176=PU Vul $`V`$ mag during the decline from outburst (XI 1994) \[153 \]. In \[m231 \] $`211^d`$ periodicity has been reported.
177=LT Del IRAS data from \[78 \]. Spectral type of cool component also estimated in \[m24 \] to be G5.
178=V1329 Cyg Spectral type of cool component also estimated in \[36 \] to be $`>`$M4.
180=ER Del Cool component is S star of type S5.5/2.5 \[267 \]. In the UV, there are emission lines with an ionization potential up to 47.9 eV and a strong UV continuum indicates the presence of a hot companion \[268 \]; in the optical spectrum there are emission lines of H i \[268 \].
181=V1329 Cyg The system inclination, $`i=86^{}\pm 2^{}`$, and the position angle of the orbital plane, $`11^{}`$, has been derived from spectropolarimetric studies. An extended nebulosity detected in the \[O iii\] $`\lambda 5007`$ line is aligned with the orbital plane \[159 \].
183=V407 Cyg IRAS data from \[118 \]. IR-type S –\[m3 \], D –\[18 \] and also there is D’ estimate in \[10 \]. In \[18 \] there is an estimate of the orbital period of 43 yrs.
184=StHA 190 In \[59 \] there is a suggestion, based on the IRAS ratio of F<sub>12</sub>/F<sub>25</sub>, that the cool component in this system is a Mira variable. Comparison of IRAS fluxes with diagnostic diagrams in \[13 \] shows that this object is among or close to D’ systems, and the $`VJHKL`$ colours are consistent with a G-K giant, so there is no reason to think that a Mira variable is present in this binary. The authors of \[59 \] argue that F<sub>12</sub>/F$`{}_{25}{}^{}>1.0`$ suggests the presence of a Mira but it may be merely the signature of dust around the system.
185=AG Peg $`VK`$ magnitudes during outburst. A complex nebula with possible bipolar structure detected at optical and radio wavelengths (\[not3 \] and references therein).
186=LL Cas The presence of the \[Fe vii\] 4892Å line reported in \[274 \] is not reliable because of the absence of the strongest \[Fe vii\] 6087Å iron line at that time. In \[274 \], there is a report of a possible pulsational period for the cool component of this system (286.6 days). This is a plausible explanation, as the spectrum taken at maximum light shows a more pronounced late-type continuum than the spectrum taken at minimum (see \[274 \]), indicating that the cool component is responsible for this variability. IR colours: J$`=8.90,H=8.03,K=7.55`$ \[m103 \]) with assumed modest amount of interstellar reddening (A<sub>K</sub>=0.2) give $`J_0=8.44,H_0=7.67,K_0=7.35`$ which corresponds to the colours of a normal giant in an S-type symbiotic star, although this might still be a Mira without an IR excess (like the Mira in R Aqr, which is another S-type symbiotic star).
187=Z And Spectral type of cool component also estimated in \[m90 \] to be $``$M5.2. An inclination of $`i=47\pm 12^{}`$ and an orbit orientation, $`\mathrm{\Omega }=72\pm 6^{}`$, derived from spectropolarimetry \[not12 \].
188=R Aqr The binary has been spatially resolved and a preliminary orbit (with a period of $`44`$ yrs) derived in \[hpl97 \]. The system is embedded in a complex bipolar nebula with jets (\[not3 \] and references therein). Only symbiotic star known to support H<sub>2</sub>O and SiO masers (\[ish94 ; m1 \]).
Suspected symbiotic stars
s01=RAW 1691 Carbon star \[199 \] + H<sub>α</sub> profile as for interacting binary star \[m18 \].
s02=\[BE74\] 583 Suspected in \[m139 \].
s03=StHA 55 Carbon star \[61 \] + with strong H i emission \[61 \] (too strong for single carbon star).
s04=GH Gem Suspected in \[178 ; 31 \].
s05=ZZ CMi This object was classified as symbiotic in \[48 ; 222 \]. We disagree with this classification because: i) colours are bluer at minimum \[316 \], the opposite to what is observed for symbiotics; the light curve looks more like a pulsational curve and not like a symbiotic light curve; ii) the spectrum presented in \[48 \] does not look like a symbiotic spectrum (e.g. H$`{}_{\gamma }{}^{}>\mathrm{H}_\beta `$) and is noisy (\[Ne iii\] line may not be present (so IP<sub>max</sub>=35.1 eV). However, this object contains a late-type star (though we do not know if the star is giant) and it displays an emission-line spectrum; also, the H<sub>α</sub> profile shown in \[222 \]) looks like a symbiotic star (for comparison see \[m18 \]). We therefore include this object as suspected symbiotic.
s06=NQ Gem Suspected in \[225 \]. Highly variable UV continuum with strong C iv\] emission and Si iii\]/C iii\] ratio similar to symbiotic stars. He ii 1640Å emission line has been detected in 1979 by IUE.
s07=WRAY 16$``$51 Probable presence of late-type star and emission-type spectrum (H i emission lines) \[179 \].
s08=Hen 3$``$653 Suspected in \[30 ; 291 \]: late-type star and emission-type spectrum (H i and He i emission lines).
s09=NSV 05572 Late-type giant and emission type-spectrum (H i emission lines).
s10=AE Cir Suspected in \[240 \]. Periods of 3900 and 100 days are mentioned in \[240 \] (based on visual photometric observations).
s11=V748 Cen Suspected in \[31 \]: M type giant \[177 ; 55 \] and emission-line spectrum (H i, Fe ii, \[Fe ii\], \[S ii\]) \[300 \] and UV excess.
s12=V345 Nor Suspected in \[m149 \]: M star \[245 \] and emission-line spectrum (H i, Fe ii) \[245 \]. Listed as N Nor 1985/2 in \[131 \].
s13=V934 Her Suspected in \[250 \]: M bright giant and UV emission lines with ionization potential up to 77.5 eV but no emission lines in optical spectrum and no short-wavelength continuum was found (the 1200–2800Å integrated flux $`<1.5\times 10^{14}`$ erg s<sup>-1</sup> cm<sup>-2</sup> Å<sup>-1</sup> at Earth) which excludes the presence of a hot white dwarf companion (although a neutron star is still possible).
s14=Hen 3$``$1383 Possible M type star \[m19 \] and emission-type spectrum (H i, He i) \[m16 \]. Nebula resolved at radio wavelengths (\[not3 \] and references therein).
s15=V503 Her Suspected in \[31 \]: M star \[m119 \] and blue excess in the optical spectra suggesting presence of hot companion.
s16=WSTB 19W032 Late type giant \[232 \] and emission-line spectrum: lines of H i, He i and others with ionization potential up to 35.1 eV. But this emission-line spectrum might not be physically connected with the giant \[232 \].
s17=WRAY 16$``$294 Suspected in \[3 \]: red continuum typical of reddened K giant and emission-line spectrum (H i, He i and others with ionization potential up to 35.1 eV). WRAY 16$``$294 appears as WRAY 16$``$296 in \[3 \].
s18=AS 241 Suspected in \[30 \]: M star \[3 \] and emission-line spectrum (H i, He i) \[3 \]. M6 spectral type of cool component and D’ IR type from \[43 \] are not reliable, as M6 does not agree with IR colours ($`JHK`$) and authors do not follow original definition of D’.
s19=DT Ser Considered as symbiotic in \[280 ; m122 \]: emission spectrum of H i, He i and other lines with ionization potential up to 54.4 eV plus G? \[280 \] or G2-K0 III-I \[m122 \] cool component. But there is a report of a G star 5<sup>′′</sup> from this object, so the cool component may not be connected physically with the source of the emission-line spectrum.
s20=V618 Sgr Presence of late-type component (TiO bands in optical spectrum \[279 \]) and emission-line spectrum (H i, Fe ii \[279 \]).
s21=AS 280 Suspected in \[3 \]: this object resembles a symbiotic star in outburst.
s22=AS 288 This object shows optical emission-line spectrum (H i, He i and others with ionization potential up to 54.4 eV) but no late-type component has been seen so far, however $`K`$ magnitude and IRAS fluxes compared to diagnostic diagrams in \[13 \] place this object among symbiotics (of IR type D), emission-line fluxes (\[O,iii\]4363, 5007Å, H<sub>β</sub>, H<sub>γ</sub>) compared to diagnostic diagrams in \[m5 \] place this object also among symbiotics (among IR type S but close to D-type objects).
s23=Hen 2$``$379 Cool component is G-K giant \[258 \] and there is emission-line spectrum: H i, He i and other lines with ionization potential up to 35.1 eV. But K giant might not be physically associated with nebula, which is source of the emission \[257 \]. Finding chart in \[185 \] is unclear as reported in \[m23 \].
s24=V335 Vul Suspected in \[287 \]: presence of carbon giant and optical emission-line spectrum (H i) displaying hot continuum in blue. We agree with this classification although this object might be a single pulsating star: (i) the carbon star might pulsate with period of 342 days \[m172 \] and then emission lines behave as for a Mira variable – they disappear near minimum light and show up again at maximum (see spectra in \[287 \]); (ii) H<sub>α</sub> is very narrow: 2Å (90 km s<sup>-1</sup>) and for a symbiotic star we would expect a width of about 300–500 km s<sup>-1</sup>; (iii) the Balmer decrement is different than that observed for symbiotic stars (it resembles that of Mira variable), although the authors of \[287 \] claim that the Balmer decrement resembles that of a symbiotic star.
s25=V850 Aql Probable presence of Mira \[m41 ; 357 \] or late-type star \[m22 \] in the centre of planetary nebula PK 037$``$6 2 (see note in \[357 \]) with emission-line spectrum (H i lines). In \[m22 \] and \[m41 \] there are notes that in \[m42 \] this object is classified as symbiotic, but this is not true and in \[m42 \] there are only IR colours for V850 Aql.
s26=Hen 2$``$442 Suspected in \[321 \]: TiO bands probably present, suggesting cool component \[m176 ; 321 \] and optical emission-line spectrum: H i, He i and other lines with ionization potential up to 100 eV. Hen 2$``$442 consists of two PN like objects: Hen 2$``$442A and Hen 2$``$442B \[m176 \] and values in catalogue are for the whole system. Symbiotic nature was suggested for Hen 2$``$442A.
s27=IRAS 19558+3333 Suspected in \[si94 \]: OH/IR star, based on IRAS colours, but without an OH maser, so a probable, extreme D-type system. Radio continuum emission implies a hot, ionising companion. Correct coordinates given here for precise radio emission (incorrect coordinates given by \[si94 \]).
s28=V627 Cas Suspected in \[m192 \]. Spectral type of cool component also estimated in \[m190 \] to be M2-4.
## 5 Comments on other objects not included in the catalogue
V1017 Sgr In some publications, regarded as symbiotic, probably after inclusion in Kenyon’s catalogue (\[31 \]), but this is not a symbiotic star. This object is a cataclysmic variable with orbital period of 5.7 days (\[129 \]).
CI Cam The optical counterpart of XTE J0421+560. Reported as symbiotic in \[357 \], possibly after suggestion in \[m237 \]. It is, however, a high-mass X-ray binary with a Be star donor (\[m238 ; m239 \]).
###### Acknowledgements.
This work has been funded by the KBN grant 2P03D02112, and also made use of the NASA Astrophysics Data System and SIMBAD database. We would like to thank Dr Estella Brandi and Dr Maciej Mikołajewski for many helpful comments. RJI acknowledges the award of a PPARC Advanced Fellowship. KB would like also to thank Dr Tomasz Bulik for help with preparation of this manuscript.
Table 1. Symbiotic stars
No. Name $`\alpha `$(2000) $`\delta `$(2000) l<sup>II</sup> b<sup>II</sup> $`V`$ $`K`$ IR IUE X IP<sub>max</sub> <sup>h</sup> <sup>m</sup> <sup>s</sup> ’ ” \[mag\] \[mag\] \[eV\] 001 SMC1 00 29 10.9 $``$74 57 38.9 16.2 S: + 114 002 SMC2 00 42 48.1 $``$74 42 00.0 16.2 S: + $``$ 114 003 EG And 00 44 37.1 +40 40 45.7 121.54 $``$22.17 7.1 2.6 S + + 100 004 SMC3 00 48 19.9 $``$73 31 54.9 15.5: S: + + 235 005 SMC N60 00 57 12.0 $``$74 13 00.0 16.8 13.0 S,D + $``$ 114 006 LIN 358 00 59 24.0 $``$75 04 59.9 15.2 11.4 S + + 114 007 SMC N73 01 04 42.0 $``$75 48 00.0 15.5 11.6 S + $``$ 114 008 AX Per 01 36 22.7 +54 15 02.5 129.53 $``$8.04 10.9 5.5 S + $``$ 109.3 009 V471 Per 01 58 49.6 +52 53 48.9 133.12 $``$8.64 13.0 9.8 D’ + $``$ 77.5 010 o Ceti 02 19 20.7 $``$02 58 39.5 167.76 $``$57.98 6.0 $``$2.7 + + 54.4 011 BD Cam 03 42 09.3 +63 13 00.5 140.84 +6.44 5.1 0.2 + 77.5 012 S32 04 37 45.0 $``$01 19 05.9 197.48 $``$30.04 13.5 S + + 114 013 LMC S154 04 51 50.2 $``$75 03 36.0 15.7 10.1 D + $``$ 114 014 LMC S147 04 54 04.6 $``$70 59 34.0 16.0 11.9 S + 114 015 LMC N19 05 03 24.0 $``$67 56 35.0 16.4 $``$ 114 016 UV Aur 05 21 48.8 +32 30 43.1 174.22 $``$2.35 8.5 2.1 S + $``$ 41.0 017 V1261 Ori 05 22 18.6 $``$08 39 58.0 210.63 $``$23.72 6.8 2.1 + + 77.5 018 LMC1 05 25 01.0 $``$62 28 46.9 15.9 9.9 D + 97.1 019 LMC N67 05 36 02.8 $``$64 43 23.9 15.9 11.4 S + $``$ 77.5 020 Sanduleak’s star 05 45 19.6 $``$71 16 09.9 16.9 13.0 D: + $``$ 114 021 LMC S63 05 48 44.1 $``$67 36 12.9 15.2 11.3 S + + 97.1 022 SMP LMC 94 05 54 10.3 $``$73 02 39.0 + 114 023 BX Mon 07 25 24.0 $``$03 36 00.0 220.04 +5.88 11.7 5.7 S + $``$ 54.4 024 V694 Mon 07 25 51.2 $``$07 44 08.0 223.76 +4.05 9.5 5.1 S + $``$ 24.6 025 WRAY 15$``$157 08 06 34.8 $``$28 31 57.0 246.60 +1.95 13.2 9.4 D’ $``$ 54.4 026 RX Pup 08 14 12.3 $``$41 42 29.0 258.50 $``$3.93 11.5 2.8 D + + 114 027 Hen 3$``$160 08 24 52.8 $``$51 28 36.0 267.68 $``$7.87 15.4 7.5 S $``$ 114 028 AS 201 08 31 42.9 $``$27 45 32.0 249.08 +6.97 11.8 9.9 D’ + $``$ 54.4 029 KM Vel 09 41 14.0 $``$49 22 47.0 274.19 +2.58 15.0 5.7 D $``$ 41.0 030 V366 Car 09 54 43.3 $``$57 18 52.4 280.81 $``$2.24 13.0 4.7 D + $``$ 114 031 Hen 3$``$461 10 39 08.5 $``$51 24 11.9 282.90 +6.25 12.3 3.9 S $``$ 100 032 SS73 29 11 08 27.4 $``$65 47 17.9 292.63 $``$5.00 14.1 10.6 S + $``$ 100 033 SY Mus 11 32 05.5 $``$65 25 08.0 294.80 $``$3.81 10.9 4.7 S + $``$ 114 034 BI Cru 12 23 26.0 $``$62 38 16.0 299.72 +0.06 11.1 5.0 D + $``$ 75.5 035 RT Cru 12 34 54.0 $``$64 33 54.0 301.16 $``$1.75 12.6 54.4 036 TX CVn 12 44 42.0 +36 45 50.6 130.93 +80.26 9.5 6.3 S + $``$ 13.6 037 Hen 2$``$87 12 45 20.1 $``$63 00 39.9 302.24 $``$0.15 15.5 6.0 S $``$ 114 038 Hen 3$``$828 12 50 58.0 $``$57 50 47.0 302.87 +5.03 13.4 7.1 S $``$ 114 039 SS73 38 12 51 26.2 $``$64 59 58.1 302.93 $``$2.13 14.5 6.1 D 108.8 040 Hen 3$``$863 13 07 43.8 $``$48 00 23.0 305.75 +14.78 11.8 8.5 S $``$ 54.4 041 St 2$``$22 13 14 30.0 $``$58 51 47.9 305.92 +3.87 8.5 S $``$ 114 042 CD$``$36 8436 13 16 01.6 $``$37 00 11.9 308.37 +25.61 11.2 5.7 S $``$ 54.4 043 V840 Cen 13 20 46.9 $``$55 50 35.9 307.07 +6.79 14.1 10.8 54.4 044 Hen 3$``$905 13 30 37.2 $``$57 58 18.0 308.12 +4.50 13.4 8.5 S 114 045 RW Hya 13 34 17.8 $``$25 22 52.1 314.99 +36.49 8.9 4.7 S + $``$ 77.5 046 Hen 3$``$916 13 35 28.9 $``$64 45 45.0 307.61 $``$2.29 12.9 7.8 S $``$ 114 047 V704 Cen 13 54 56.2 $``$58 27 16.9 311.17 +3.40 13.5 8.4 D $``$ 54.4 048 V852 Cen 14 11 52.1 $``$51 26 23.8 315.48 +9.46 14.0 6.9 D + $``$ 114 049 V835 Cen 14 14 09.4 $``$63 25 46.1 312.03 $``$2.03 12.9 5.0 D + $``$ 114 050 V417 Cen 14 15 55.9 $``$61 53 53.9 312.71 $``$0.64 12.2 D’ 35.1
Table 1. Symbiotic stars – cont.
No. Cool-star spectrum Radio IRAS IRAS References \[mJy\] F<sub>12</sub>\[Jy\] F<sub>25</sub>\[Jy\] 001 C3.2:,C(35 ; 34 ) 34 (fc,spc,class) 35 (parm) 002 K,G-K(35 ; 34 ) 34 (fc,spc,class) 35 (parm) 003 M3(ms99 ) 0.54(3.6cm) 4.5 1.25 30 (fc,class) 31 (spc) 26 (parm) 004 M0,K-M(34 ; 35 ) 34 (fc,spc,class) 197 (parm) 005 C3.3(35 ) 30 (fc,class) 35 (spc,parm) 006 mid K(35 ) 30 (fc,class) 35 (spc,parm) 007 K7(35 ) 30 (fc,class) 35 (spc,parm) 008 M6(ms99 ) 0.58(3.6cm) 0.32 0.10 30 (fc,class) 37 ; ibesm93 (spc,parm) 009 G5(ms99 ) 0.45(3.6cm) 1.53 2.70 30 (fc,class) 39 (spc,parm) 010 Mira,M2-7(22 ; 201 ) 4881 2261 204 (spc) 201 (parm) (class) 011 M3,S5.3(206 ; 207 ) 0.18(3.6cm) 40.95 10.82 207 (spc) 206 (parm) (class) 012 C1.1CH(42 ) m233 (fc) 61 (class) 42 (spc) 59 (parm) 013 C2.2(35 ) 208 (fc) 35 (class,spc,parm) 014 M1,K5(35 ; 33 ) 33 (fc,spc,class) 35 (parm) 015 M4(91 ) 96(6cm) 209 (fc) 91 (class,spc,parm) 016 Mira C8.1Je(41 ) $`<`$0.34(3.6cm) 69.41 20.64 30 (fc,spc,class) 40 (parm) 017 S4.1,M3(211 ; 212 ) $`<`$0.1(3.6cm) 7.98 2.01 (class) 211 (parm) 018 C4.3,C(35 ; 34 ) 34 (fc,spc,class) 35 (parm) 019 C3.2,C(35 ; 32 ) 296 (fc) 32 (spc,class) 35 (parm) 020 30 (fc) (class) 35 (spc,parm) 021 C2.1J(35 ) 30 (fc,class) 35 (spc,parm) 022 M(284 ) 285 (fc) m65 (spc,class) 023 M5(ms99 ) $`<`$0.07(3.6cm) 0.32 0.09L 30 (fc,class,spc) 343 (parm) 024 M6(ms99 ) 0.84 0.36: (class) 223 (spc,parm) 025 G5(ms99 ) 0.34(3.6cm) 1.00 1.42 30 (fc,class,spc) 026 Mira M5.5/M5(ms99 ; 16 ) 12.6$`÷`$69.7(3.6cm) 220.0 132.0 30 (fc,class) Metal99 (spc,parm) 027 M7(ms99 ; 16 ) $`<`$17(6cm) 0.12 0.03L 30 (fc,class,spc) 028 G5(ms99 ) 0.69(3.6cm) 1.00 2.81 30 (fc,class,spc) 54 (parm) 029 Mira,M(22 ; 1 ) 1.49(3.5cm) 11.84 7.81 233 ; 49 (spc) (class) m49 (parm) 030 Mira M6(ms99 ) 1.73(3.5cm) 9.26 4.28 30 (fc,class,spc) 3 (parm) 031 M6(ms99 ) 2.02 0.61 30 (fc) 291 (spc,class) 032 G/K(44 ,66 ) $`<`$16(6cm) 0.07 0.11L 30 (fc,class,spc) 66 (parm) 033 M5(ms99 ) $`<`$12(6cm) 0.94 0.43 30(fc,class,spc) 70(parm) 034 Mira M0-1(55 ) 5.11(3.5cm) 17.29 15.34 30 (fc,class,spc) 64 (spc) 65 (parm) 035 M4-5(6 ) 6 (spc,class,parm) 036 early M,K5.3(360 ; 36 ) $`<`$0.06(3.6cm) 1.07 0.33 31 (spc,class) 72 (parm) 037 M5.5/M7.5(ms99 ; 16 ) $`<`$12(6cm) 0.40 0.75: 30 (fc,class,spc) 3 (parm) 038 M6(ms99 ; 16 ) $`<`$16(6cm) 1.10 0.38 30 (fc,class,spc) 039 Mira C9(sl88 ) 5.12(3.5cm) 7.71 3.29 30 (fc,class,spc) 040 K4:(16 ) $`<`$12(6cm) 30 (fc,class,spc) 041 M4.5(ms99 ) 30 (fc,class) 3 (spc,parm) 042 M5.5(ms99 ) $`<`$0.24(6cm) 0.40 0.25L 30 (fc) 123 (spc,class,parm) 043 K5(186 ) 0.03L 0.22L 237 (fc) 186 (spc,class,parm) 044 M3(ms99 ) $`<`$13(6cm) 30 (fc,class,spc) 045 M2(ms99 ) 0.60(3.6cm) 0.88 0.45L 30 (fc,class,spc) 73 (parm) 046 M5/M6(ms99 ; 16 ) $`<`$20(6cm) 30 ; 52 (fc) 30 (class,spc) 047 M6.5(ms99 ) 1.43 1.16 30 (fc) 291 (spc) 123 (class,parm) 048 Mira(22 ) 3.31(3.5cm) 8.56 9.09 30 (fc,spc) (class) 175 (parm) 049 Mira $``$M5(55 ) 5.79(3.5cm) 32.45 26.05 30 (fc,class,spc) 80 (parm) 050 G8-K2,G2(m122 ; 239 ) 12.61 10.60 6 (spc,class) 239 (parm)
Table 1. Symbiotic stars – cont.
No. Name $`\alpha `$(2000) $`\delta `$(2000) l<sup>II</sup> b<sup>II</sup> $`V`$ $`K`$ IR IUE X IP<sub>max</sub> <sup>h</sup> <sup>m</sup> <sup>s</sup> ’ ” \[mag\] \[mag\] \[eV\] 051 BD$``$21 3873 14 16 34.3 $``$21 45 50.2 327.88 +36.95 10.7 7.2 S + $``$ 77.5 052 Hen 2$``$127 15 24 49.8 $``$51 49 52.6 325.53 +4.19 16.0 8.1 D $``$ 114 053 Hen 3$``$1092 15 47 10.6 $``$66 29 16.0 319.22 $``$9.35 13.5 7.8 S + $``$ 114 054 Hen 3$``$1103 15 48 28.5 $``$44 19 00.9 333.21 +7.90 13.0 8.4 S $``$ 114 055 HD 330036 15 51 15.9 $``$48 44 58.5 330.78 +4.15 11.0 7.6 D’ + $``$ 77.5 056 Hen 2$``$139 15 54 45.4 $``$55 29 36.9 326.91 $``$1.40 16.8 5.9 D $``$ 35.1 057 T CrB 15 59 30.1 +25 55 12.6 42.37 +48.16 10.1 4.8 S + + 77.5 058 AG Dra 16 01 40.5 +66 48 09.5 100.29 +40.97 9.1 6.2 S + + 114 059 WRAY 16$``$202 16 06 56.8 $``$49 26 39.0 332.28 +1.96 6.8 S $``$ 114 060 V347 Nor 16 14 00.3 $``$56 59 26.0 327.92 $``$4.30 16.6 5.0 D + $``$ 54.4 061 UKS Ce$``$1 16 15 29.2 $``$22 12 15.9 353.02 +20.25 15.0 11.3 S $``$ 54.4 062 QS Nor 16 21 07.9 $``$42 23 53.9 338.94 +5.36 13.3 8.1 S $``$ 35.1 063 WRAY 15$``$1470 16 23 21.6 $``$27 40 13.0 350.06 +15.24 12.9 7.8 S $``$ 54.4 064 Hen 2$``$171 16 34 03.9 $``$35 05 32.8 346.03 +8.55 14.8 6.7 D + $``$ 235 065 Hen 3$``$1213 16 35 15.3 $``$51 42 24.0 333.87 $``$2.81 10.4 6.7 S + $``$ 97.1 066 Hen 2$``$173 16 36 24.1 $``$39 51 52.4 342.77 +5.01 13.8 6.8 S $``$ 54.4 067 Hen 2$``$176 16 41 31.2 $``$45 13 04.7 339.39 +0.74 15 5.7 S,D $``$ 114 068 KX TrA 16 44 35.2 $``$62 37 14.0 326.41 $``$10.94 12.4 6.0 S + $``$ 114 069 AS 210 16 51 20.4 $``$26 00 27.1 355.51 +11.55 12.2 6.5 D + $``$ 114 070 HK Sco 16 54 39.0 $``$30 23 30.0 352.48 +8.27 13.5 7.9 S + $``$ 114 071 CL Sco 16 54 51.9 $``$30 37 18.0 352.32 +8.09 13.3 7.9 S + $``$ 54.9 072 MaC 1$``$3 17 01 27.9 $``$47 45 33.9 339.62 $``$3.52 18.2 D $``$ 114 073 V455 Sco 17 07 21.5 $``$34 05 21.0 351.17 +3.90 13.7 5.9 S $``$ 114 074 Hen 3$``$1341 17 08 36.6 $``$17 26 30.0 5.02 +13.39 12.5 7.6 S + $``$ 114 075 Hen 3$``$1342 17 08 55.0 $``$23 23 35.0 0.08 +9.92 12.7 8.5 S $``$ 114 076 AS 221 17 12 02.2 $``$32 34 24.9 352.98 +4.00 12.0 7.6 S $``$ 114 077 H 2$``$5 17 15 19.0 $``$31 34 05.9 354.20 +4.02 5.5 S,D $``$ 114 078 Sa 3$``$43 17 17 55.9 $``$30 01 48.0 355.79 +4.45 7.9 S $``$ 54.4 079 Draco C$``$1 17 19 57.6 +57 50 04.9 17.0 11.4 + + 77.5 080 Th 3$``$7 17 21 02.4 $``$29 22 59.4 356.70 +4.27 14.0 8.1 S $``$ 114 081 Th 3$``$17 17 27 31.9 $``$29 02 53.0 357.78 +3.27 14.0 8.2 S $``$ 54.4 082 Th 3$``$18 17 28 27.0 $``$28 38 34.0 358.23 +3.33 12.6 8.2 S $``$ 114 083 Hen 3$``$1410 17 29 06.0 $``$29 43 24.1 357.41 +2.61 12.8 8.5 S $``$ 114 084 V2116 Oph 17 32 03.0 $``$24 44 44.3 1.94 +4.79 19.0 8.1 S + 235 085 Th 3$``$29 17 32 28.0 $``$29 05 06.0 358.34 +2.35 17.0 7.0 S $``$ 35.1 086 Th 3$``$30 17 33 43.0 $``$28 07 18.0 359.30 +2.65 13.1 8.3 S $``$ 114 087 Th 3$``$31 17 34 26.5 $``$29 28 05.0 358.25 +1.78 13.6 7.6 S $``$ 114 088 M 1$``$21 17 34 17.2 $``$19 09 21.9 6.96 +7.36 13.7 7.2 S $``$ 114 089 Hen 2$``$251 17 35 22.3 $``$29 45 20.0 358.12 +1.46 15.4 7.3 D $``$ 114 090 Pt 1 17 38 49.8 $``$23 54 02.9 3.48 +3.94 15.0 8.6 S $``$ 114 091 K 6$``$6 17 39 18.3 $``$28 15 06.9 359.85 +1.54 16.5 S $``$ 114 092 RT Ser 17 39 51.8 $``$11 56 44.7 13.89 +9.97 15.0 7.0 S + $``$ 114 093 AE Ara 17 41 04.9 $``$47 03 20.9 344.00 $``$8.66 12.5 6.4 S + $``$ 77.5 094 SS73 96 17 41 28.3 $``$36 47 51.9 352.85 $``$3.38 15.2 6.4 S $``$ 100 095 UU Ser 17 42 37.9 $``$15 24 29.9 11.23 +7.62 15.3 9.1 S $``$ 114 096 V2110 Oph 17 43 33.3 $``$22 45 37.0 5.03 +3.62 19 7.9 D $``$ 114 097 V916 Sco 17 43 54.9 $``$36 03 24.9 353.73 $``$3.41 8.3 S $``$ 54.4 098 Hen 2$``$275 17 45 31.0 $``$38 39 47.9 351.67 $``$5.03 15.5 $`>`$9.3 S $``$ 114 099 V917 Sco 17 48 04.0 $``$36 08 17.9 354.10 $``$4.17 13.0 7.9 S $``$ 114 100 H 1$``$36 17 49 48.1 $``$37 01 27.9 353.51 $``$4.92 12.0 7.5 D + + 114
Table 1. Symbiotic stars – cont.
No. Cool-star spectrum Radio IRAS IRAS References \[mJy\] F<sub>12</sub>\[Jy\] F<sub>25</sub>\[Jy\] 051 K2(ms99 ) $`<`$0.22(3.6cm) 30 (fc,class,spc) 81 (parm) 052 Mira M5/M7(ms99 ; 16 ) 0.70(3.5cm) 0.49 0.29 30 (fc,class,spc) 3 (parm) 053 M5.5(ms99 ) 10(6cm) 30 (fc,class,spc) 3 (parm) 054 M3.5/M0(ms99 ; 16 ) $`<`$1(2cm) 30 (fc,class,spc) 055 G5(ms99 ) 34(2cm) 19.53 38.58 30 (fc,spc) (class) 59 (parm) 056 Mira,M6.5(ms99 ) 5.87 2.99 30 (fc,class,spc) 057 M4.5(ms99 ) $`<`$0.08(3.6cm) 0.68 0.20 30 (fc,class) 82 (spc) 46 (parm) 058 K2(ms99 ) 0.48(6cm) 0.22 0.74L 30 (fc,class) 39 (spc) 86 (parm) 059 M6(ms99 ) 30 (fc,class) 060 Mira M7/M8.5(ms99 ; 16 ) $`<`$12(6cm) 4.04 2.84 30 (fc,class,spc) 3 (parm) 061 C4.5Jch(42 ) $`<`$0.07(3.6cm) 30 (fc,spc,class) 42 (parm) 062 mid M(179 ) 185 (fc) 1 (spc) 3 (class) 291 (parm) 063 M3,M6(ms99 ; 5 ) $`<`$0.07(3.6cm) 30 (fc,class,spc) 064 Mira M6.5,M6(ms99 ; 5 ) 3.39(3.6cm) 7.47 4.58 30 (fc,class,spc) 89 (parm) 065 M2/K4(m19 ; 16 ) $`<`$12(6cm) 30 (fc,class,spc) 066 M4.5,M7(ms99 ; 5 ) 0.21(3.6cm) 0.20 0.05L 30 (fc,class,spc) 3 (parm) 067 Mira(m2 );M4/M7(ms99 ; 16 ) 4.05(3.6cm) 2.81 2.40 30 (fc,class,spc) 068 M6(ms99 ) $`<`$22(6cm) 0.26 0.25L 131 (fc) 30 (class,spc) 93 (parm) 069 Mira; C3(5 ) 4.56(3.6cm) 3.93 1.21 30 (fc,class,spc) 59 (parm) 070 M3.5(ms99 ) $`<`$0.07(3.6cm) 30 (fc,class,spc) 071 M5(ms99 ) 0.34(3.6cm) 0.04L 0.18L 30 (fc,class,spc) 95 (parm) 072 M2(3 ) 182 (fc) 3 (spc,class,parm) 073 M6.5,M6(ms99 ; 5 ) 1.99(3.6cm) 1.23 1.54 30 (fc,class,spc) 3 (parm) 074 M4/M0(ms99 ; 16 ) 1.04(3.6cm) 30 (fc,class,spc) 90 (parm) 075 M0(ms99 ) $`<`$0.08(3.6cm) 30 (fc,class,spc) 076 M7.5(ms99 ) 0.16(3.6cm) 3.32 3.46 30 (fc,class,spc) 3 (parm) 077 M5.5(ms99 ) $`<`$0.09(3.6cm) 1.20 0.67: 30 (fc,class,spc) 078 early M(30 ) $`<`$0.57(6cm) 30 (fc,class) 123 (spc,parm) 079 C1.2(31 ) $`<`$0.10(3.6cm) 30 (fc,class) 68 (spc) 66 (parm) 080 M5(5 ) 0.36(6cm) 30 (fc,class,spc) 3 (parm) 081 M3(3 ) $`<`$0.11(3.6cm) 30 (fc,class,spc) 3 (parm) 082 M3(5 ) $`<`$0.14(3.6cm) 30 (fc,class,spc) 3 (parm) 083 M1.5,M3(ms99 ; 3 ) 0.57(6cm) 30 (fc,class,spc) 3 (parm) 084 M5/M6:(25 ; 16 ) 0.06:(6cm) 24 (spc,class) 25 (fc,parm) 085 M3(3 ) $`<`$0.26(6cm) 6.02 5.62 1 (fc,class) 3 (spc) 086 M1(3 ) $`<`$0.11(3.6cm) 30 (fc,class,spc) 3 (parm) 087 M5(16 ) $`<`$0.10(3.6cm) 1.89 2.10 30 (fc,class,spc) 3 (parm) 088 M6/M2(ms99 ; 16 ) 0.51(3.6cm) 0.65 0.36: 30 (fc,class,spc) 3 (parm) 089 late M(ntt00 ) 7.33 7.46 185 (fc) 6 (spc) 3 (class,parm) 090 M3(ms99 ) $`<`$0.10(3.6cm) 30 (fc,class,spc) 97 (parm) 091 m25 (fc) 183 (spc) 3 (class,parm) 092 M6(ms99 ) 1.19(3.6cm) 0.12 0.05L 30 (fc,class,spc) 99 (parm) 093 M5.5/M2(ms99 ; 16 ) $`<`$18(2cm) 0.22 0.08L 30 (fc,class,spc) 100 (parm) 094 M0/M2(5 ; 16 ) 1.59(3.6cm) 0.43 0.35L 30 (fc,class,spc) 095 M4(3 ) $`<`$0.13(3.6cm) 30 (fc,class,spc) 3 (parm) 096 Mira M8(16 ) $`<`$0.19(3.6cm) 0.15 0.06L 30 (fc,class,spc) 097 M(31 ) $`<`$0.14(3.6cm) 30 (fc,class) 102 (parm) 098 M3(3 ) 1 (fc,spc,class) 3 (parm) 099 M7(16 ) $`<`$0.12(3.6cm) 30 (fc,class,spc) 100 Mira(104 );M4-5:(5 ) 65.3(3.6cm) 18.13L 28.16 30 (fc,class,spc) m68 (parm)
Table 1. Symbiotic stars – cont.
No. Name $`\alpha `$(2000) $`\delta `$(2000) l<sup>II</sup> b<sup>II</sup> $`V`$ $`K`$ IR IUE X IP<sub>max</sub> <sup>h</sup> <sup>m</sup> <sup>s</sup> ’ ” \[mag\] \[mag\] \[eV\] 101 RS Oph 17 50 13.2 $``$06 42 28.4 19.80 +10.37 11.5 6.5 S + + 114 102 WRAY 16$``$312 17 50 16.5 $``$30 57 41.1 358.79 $``$1.91 7.8 D + 108.8 103 V4141 Sgr 17 50 24.0 $``$19 53 42.0 8.31 +3.73 16.4 7.9 D,S $``$ 54.4 104 ALS 2 17 50 51.1 $``$17 47 56.0 10.18 +4.70 14.2 S 114 105 AS 245 17 51 00.8 $``$22 19 41.3 6.29 +2.38 11.0 7.2 S,D $``$ 114 106 Hen 2$``$294 17 51 38.2 $``$32 50 57.0 357.31 $``$3.12 8.4 S $``$ 114 107 Bl 3$``$14 17 52 25.9 $``$29 46 00.0 0.05 $``$1.70 14.3 8.7 S $``$ 114 108 Bl 3$``$6 17 52 57.0 $``$31 19 18.0 358.77 $``$2.59 17.0 S $``$ 114 109 Bl L 17 53 13.0 $``$30 18 00.0 359.68 $``$2.12 16.0 7.8 S $``$ 114 110 V745 Sco 17 55 22.2 $``$33 14 59.3 357.36 $``$4.00 17 8.4 S + $``$ 262 111 MaC 1$``$9 17 55 52.6 $``$14 06 49.3 13.99 +5.51 15.2 S $``$ 54.4 112 AS 255 17 57 08.7 $``$35 15 38.0 355.79 $``$5.32 12.5 8.4 S,D $``$ 54.4 113 V2416 Sgr 17 57 15.9 $``$21 41 35.4 7.57 +1.44 15.1 4.6 S + $``$ 114 114 H 2$``$34 17 58 28.0 $``$28 33 42.0 1.76 $``$2.23 18.0 $`>`$9.9 $``$ 114 115 SS73 117 18 02 22.9 $``$31 59 11.0 359.19 $``$4.66 12.5 7.7 S $``$ 114 116 AS 269 18 03 23.8 $``$32 42 21.9 358.67 $``$5.20 13.9 8.7 D’ $``$ 54.4 117 Ap 1$``$8 18 04 29.8 $``$28 21 28.0 2.59 $``$3.28 15.2 7.9 S $``$ 114 118 SS73 122 18 04 41.1 $``$27 09 32.8 3.66 $``$2.73 12.0 6.6 D $``$ 114 119 AS 270 18 05 33.6 $``$20 20 44.6 9.64 +0.37 13.1 5.5 S $``$ 54.4 120 H 2$``$38 18 06 01.0 $``$28 17 10.8 2.81 $``$3.54 13.8 6.7 D + + 114 121 SS73 129 18 07 05.7 $``$29 36 25.9 1.77 $``$4.39 12.5 8.0 S $``$ 114 122 Hen 3$``$1591 18 07 31.9 $``$25 53 44.0 5.07 $``$2.68 12.5 9.0 S,D’ + 54.4 123 V615 Sgr 18 07 49.6 $``$36 10 09.0 356.05 $``$7.66 13.3 7.6 S $``$ 54.4 124 Ve 2$``$57 18 08 24.0 $``$24 34 00.0 6.33 $``$2.20 12.5 7.4 S $``$ 35.1 125 AS 276 18 09 09.6 $``$41 13 26.0 351.64 $``$10.24 12.0 8.1 S,D $``$ 114 126 Ap 1$``$9 18 10 28.9 $``$28 07 41.0 3.43 $``$4.33 12.5 8.8 S $``$ 54.4 127 AS 281 18 10 47.5 $``$27 56 24.0 3.63 $``$4.30 12.9 7.0 S $``$ 114 128 V2506 Sgr 18 11 01.6 $``$28 32 39.9 3.12 $``$4.63 12.0 8.4 S $``$ 114 129 SS73 141 18 12 11.2 $``$33 10 41.0 359.13 $``$7.04 12.5 9.0 S $``$ 54.4 130 AS 289 18 12 22.0 $``$11 40 13.0 18.08 +3.20 12.1 5.0 S + $``$ 100 131 Y CrA 18 14 24.1 $``$42 51 00.0 350.60 $``$11.84 14.4 6.6 S + $``$ 114 132 YY Her 18 14 34.3 +20 59 20.0 48.14 +17.24 12.8 8.0 S + $``$ 109.3 133 V2756 Sgr 18 14 34.5 $``$29 49 22.4 2.36 $``$5.92 11.5 7.8 S $``$ 114 134 FG Ser 18 15 06.2 $``$00 18 57.6 28.5 +7.9 11.0 4.5 S + $``$ 100 135 HD 319167 18 15 24.5 $``$30 31 57.2 1.81 $``$6.41 12.5 7.5 S $``$ 54.4 136 Hen 2$``$374 18 15 36.4 $``$21 34 40.0 9.75 $``$2.23 12.0 6.5 S $``$ 114 137 Hen 2$``$376 18 15 46.0 $``$27 53 48.0 4.19 $``$5.24 14.1 S + 35.1 138 V4074 Sgr 18 16 05.5 $``$30 51 11.3 1.59 $``$6.69 11.5 7.7 + $``$ 235 139 V2905 Sgr 18 17 20.5 $``$28 09 51.0 4.11 $``$5.68 12.3 7.1 S $``$ 63.5 140 StHA 149 18 18 55.9 +27 26 12.0 54.81 +18.77 13.5 114 141 Hen 3$``$1674 18 20 19.2 $``$26 22 47.8 6.01 $``$5.43 12.5 7.7 S $``$ 114 142 AR Pav 18 20 27.8 $``$66 04 42.9 328.54 $``$21.60 10.0 7.2 S + $``$ 97.1 143 V3929 Sgr 18 20 58.8 $``$26 48 32.0 5.70 $``$5.76 12 7.4 D $``$ 114 144 V3804 Sgr 18 21 28.5 $``$31 32 03.9 1.52 $``$8.02 12.0 7.3 S $``$ 114 145 V443 Her 18 22 08.4 +23 27 20.0 51.23 +16.59 11.5 5.4 S + $``$ 97.1 146 V3811 Sgr 18 23 28.9 $``$21 53 08.9 10.34 $``$3.98 14 8.5 S $``$ 54.4 147 V4018 Sgr 18 25 27.0 $``$28 35 57.5 4.55 $``$7.46 11.4 7.5 S + $``$ 100 148 V3890 Sgr 18 30 43.0 $``$24 00 59.0 9.21 $``$6.44 $``$14 8.2 S + 361 149 V2601 Sgr 18 37 58.9 $``$22 42 07.9 11.14 $``$7.34 15 8.0 S $``$ 54.4 150 K 3$``$9 18 40 24.1 $``$08 43 46.6 23.91 $``$1.54 18.0 5.6 D $``$ 114
Table 1. Symbiotic stars – cont.
No. Cool-star spectrum Radio IRAS IRAS References \[mJy\] F<sub>12</sub>\[Jy\] F<sub>25</sub>\[Jy\] 101 M0-M2(105 ; am99 ) $`<`$0.33(3.6cm) 0.30 0.16 30 (fc,class,spc) 105 (parm) 102 Mira(22 ) 0.96(6cm) 5.0 3.1 30 (fc,spc) (class) 103 M6(ntt00 ) $`<`$0.17(6cm) 1.88 0.65 30 (fc) 123 (spc,class) 3 (parm) 104 M2(3 ) 1 (fc,spc,class) 3 (parm) 105 Mira(22 );M6(16 ) 0.87(6cm) 1.67 0.62 30 (fc,class,spc) 3 (parm) 106 M3(3 ) $`<`$0.11(3.6cm) 2.19 2.89 30 (fc,class,spc) 3 (parm) 107 M6(58 ) $`<`$0.12(3.6cm) 30 (fc,class,spc) 58 (parm) 108 M3(3 ) $`<`$0.20(6cm) 1.62 1.14 1 (fc,class) 3 (spc,parm) 109 M6.5(16 ) $`<`$0.10(3.6cm) 30 (fc,class,spc) 110 M6,M7(am99 ; 253 ) 253 (fc) 254 (spc,class) am99 (parm) 111 M2(3 ) 182 (fc) 3 (spc,class,parm) 112 K3(5 ) $`<`$0.10(3.6cm) 1.31L 1.58L 30 (fc,class,spc) 113 M6/M5(ms99 ; 16 ) 2.26(3.6cm) 3.80L 2.10 30 (fc,class,spc) 3 (parm) 114 M5:() 185(fc) 1 (spc,class) 115 M5/M6(5 ; 16 ) 0.25(3.6cm) 0.33 0.09 30 (fc,class,spc) 116 G5-K2(3 ) $`<`$1(2cm) 1.01 2.54 185 (fc) 3 (spc) (class) 117 M4/M0(3 ; 16 ) $`<`$0.10(3.6cm) 0.84 0.63 30 (fc,class,spc) 3 (parm) 118 Mira(22 );M7(5 ) 2.19(3.6cm) 1.44 1.00 30 (fc,class,spc) 3 (parm) 119 M5.5/M1(ms99 ; 16 ) 0.29(3.6cm) 1.67 5.98L 30 (fc,class,spc) 120 Mira(22 );M7(ms99 ) 4.50(3.6cm) 3.39 2.06 30(fc,class) m68 (spc,parm) 121 M0(ms99 ) $`<`$0.14(3.6cm) 0.74 1.42L 30 (fc,class,spc) 3 (parm) 122 G(16 ) 0.41(3.6cm) 14.19 11.47 30 (fc,class,spc) m4 (parm) 123 M5.5,M5(ms99 ; 5 ) $`<`$0.11(3.6cm) 0.87 0.54L 30 (fc,class,spc) 3 (parm) 124 M(m16 ) $`<`$0.25(6cm) 30 (fc,spc) (class) 125 M4.5(ms99 ) $`<`$0.14(3.6cm) 0.31L 0.37L 30 (fc,class,spc) 126 K5/K4(3 ; 16 ) $`<`$0.16(3.6cm) 1.13 0.70: 30 (fc,class,spc) 89 (parm) 127 M5(5 ) 0.24(3.6cm) 0.79 0.54: 30 (fc,class,spc) 3 (parm) 128 M5.5(ms99 ) 0.23(3.6cm) 0.21 0.18L 30 (fc,class,spc) 3 (parm) 129 M5(ms99 ) $`<`$0.10(3.6cm) 30 (fc,class,spc) 130 M3.5(ms99 ) 1.84(3.6cm) 1.16 0.73 30 (fc,class,spc) 131 M6/M5(ms99 ; 16 ) 0.31(3.6cm) 0.25L 0.25L 30 (fc,class,spc) 132 M4(ms99 ) $`<`$0.11(3.6cm) 0.10 0.04L 30 (fc,class,spc) 111 (parm) 133 M3(5 ) 0.17(3.6cm) 30 (fc,class,spc) 3 (parm) 134 M5.3,M5(36 ; ms99 ) 0.43(6cm) 1.00 0.42: 30 (fc,class,spc) Metal00 (parm) 135 M3(3 ; 5 ) 0.31(3.6cm) 30 (fc,class,spc) 3 (parm) 136 M5.5(ms99 ) 0.39(3.6cm) 1.28 0.71: 30 (fc,class,spc) 3 (parm) 137 M2(3 ) 1 (fc,spc) 6 (class) 3 (parm) 138 M4(117 ) $`<`$0.47(3.6cm) 14.6 65.1 30 (fc,class,spc) 117 (parm) 139 late M,M0(m130 ; 3 ) 0.07L 0.06L 3 (class,parm) 291 (spc) 140 M2(m138 ) 61 (spc,class) 141 M5(5 ) 0.28(3.6cm) 30 (fc,class,spc) 142 M5/M6(ms99 ; 16 ) $`<`$14(6cm) 0.14 0.05L 30 (fc,class,spc) 119 (parm) 143 Mira M1-2(5 ) 0.28(3.6cm) 4.68 4.37 30 (fc,class,spc) 208 (parm) 144 M5/M6(5 ; 16 ) 0.46(3.6cm) 0.48 0.50L 30 (fc) 291 (spc) 31 (class) 123 (parm) 145 M5.5(ms99 ) 0.23(3.6cm) 0.34 0.25L 30 (fc,class,spc) 125 (parm) 146 M3.5(ms99 ) $`<`$0.07(3.6cm) 5.40 2.20 30 (fc,class,spc) 147 M4(16 ) $`<`$0.11(3.6cm) 0.44 0.27L 30 (fc,class,spc) 126 (parm) 148 M5,M8(m96 ; 261 ) 131 (fc) 261 (spc) (class) am99 (parm) 149 M5(3 ) $`<`$0.12(3.6cm) 0.50 0.50L 30 (fc,class,spc) 3 (parm) 150 Mira(106 ; 3 );M3(3 ) 8.63(3.6cm) 8.09 4.08 185 (fc) 3 (spc,class) 106 ; 3 (parm)
Table 1. Symbiotic stars – cont.
No. Name $`\alpha `$(2000) $`\delta `$(2000) l<sup>II</sup> b<sup>II</sup> $`V`$ $`K`$ IR IUE X IP<sub>max</sub> <sup>h</sup> <sup>m</sup> <sup>s</sup> ’ ” \[mag\] \[mag\] \[eV\] 151 AS 316 18 42 33.0 $``$21 17 48.0 12.89 $``$7.67 12.0 7.8 S $``$ 114 152 DQ Ser 18 44 40.0 +05 03 29.9 36.67 +3.81 14.6 54.4 153 MWC 960 18 47 55.9 $``$20 05 49.9 14.53 $``$8.27 12.0 7.8 S + $``$ 114 154 AS 323 18 48 35.7 $``$06 41 08.9 26.66 $``$2.41 15.5 7.9 S $``$ 114 155 AS 327 18 53 16.9 $``$24 22 54.0 11.14 $``$11.23 11.8 8.5 S $``$ 114 156 FN Sgr 18 53 52.9 $``$18 59 42.0 16.15 $``$9.06 12.7 7.9 S + $``$ 114 157 Pe 2$``$16 18 54 10.0 $``$04 38 53.9 29.10 $``$2.72 16.0 8.1 S $``$ 114 158 CM Aql 19 03 35.2 $``$03 03 15.0 31.59 $``$4.09 13.6 7.7 S $``$ 54.4 159 V919 Sgr 19 03 46.0 $``$16 59 53.9 19.01 $``$10.32 12.2 7.3 S $``$ 54.4 160 V1413 Aql 19 03 51.6 +16 28 31.7 48.97 +4.77 14.0 7.5 S + $``$ 114 161 NSV 11776 19 09 55.8 $``$02 47 40.3 32.55 $``$5.39 13.6 114 162 Ap 3$``$1 19 10 36.1 +02 49 32.0 37.64 $``$2.97 15.4 8.6 S +? 114 163 MaC 1$``$17 19 12 57.3 $``$05 21 20.0 30.59 $``$7.22 15.6 S $``$ 63.5 164 V352 Aql 19 13 33.6 +02 18 14.0 37.51 $``$3.86 16.2 $`>`$9.6 S $``$ 114 165 ALS 1 19 16 16.2 $``$08 17 45.9 28.30 $``$9.27 14.8 S 114 166 BF Cyg 19 23 53.4 +29 40 25.1 62.93 +6.70 12.3 6.3 S + $``$ 97.1 167 CH Cyg 19 24 33.0 +50 14 29.1 81.86 +15.58 7.1 $``$0.7 S + + 77.5 168 StHA 164 19 28 40.9 $``$06 03 42.0 31.73 $``$11.03 13.6 114 169 HM Sge 19 41 57.1 +16 44 39.9 53.57 $``$3.15 17 4.3 D + + 205.1 170 Hen 3$``$1761 19 42 25.3 $``$68 07 35.3 327.67 $``$29.76 10.4 5.6 S + $``$ 100 171 QW Sge 19 45 49.6 +18 36 50.0 55.64 $``$3.02 12.8 7.1 S + $``$ 114 172 CI Cyg 19 50 11.8 +35 41 03.2 70.90 +4.74 11.0 4.5 S + $``$ 109.3 173 StHA 169 19 51 28.9 +46 23 06.0 80.38 +9.85 $`>`$13.5 54.4 174 V1016 Cyg 19 57 04.9 +39 49 33.9 75.17 +5.68 11.2 5.3 D + + 141.3 175 RR Tel 20 04 18.5 $``$55 43 33.1 342.16 $``$32.24 10.8 4.2 D + + 141.3 176 PU Vul 20 21 12.0 +21 34 41.9 62.58 $``$8.52 11.6 6.2 S + + 235 177 LT Del 20 35 57.3 +20 11 34.0 63.40 $``$12.15 13.1 9.4 S + $``$ 77.5 178 StHA 180 20 39 21.0 $``$05 16 59.9 40.89 $``$26.36 13.5 114 179 Hen 2$``$468 20 41 18.9 +34 44 52.0 75.94 $``$4.44 14.2 8.0 S $``$ 114 180 ER Del 20 42 46.4 +08 40 56.4 54.46 $``$20.00 10 + $``$ 47.9 181 V1329 Cyg 20 51 01.1 +35 34 51.2 77.80 $``$5.56 13.3 6.9 S + $``$ 141.3 182 CD$``$43 14304 21 00 06.3 $``$42 38 49.9 358.65 $``$41.10 10.8 7.6 S + + 114 183 V407 Cyg 21 02 13.0 +45 46 30.0 86.99 $``$0.49 14 3.3 D + $``$ 54.4 184 StHA 190 21 41 44.8 +02 43 54.4 58.41 $``$35.43 10.5 7.8 D’ + $``$ 47.9 185 AG Peg 21 51 01.9 +12 37 29.4 69.28 $``$30.89 9.0 3.9 S + + 100 186 LL Cas 23 09 20.1 +54 44 53.0 108.47 $``$5.23 15 7.6 S 54.4: 187 Z And 23 33 39.5 +48 49 05.4 109.98 $``$12.09 10.8 5.0 S + + 114 188 R Aqr 23 43 49.4 $``$15 17 04.2 66.52 $``$70.33 9.1 $``$0.9 D + + 331
Table 1. Symbiotic stars – cont.
No. Cool-star spectrum Radio IRAS IRAS References \[mJy\] F<sub>12</sub>\[Jy\] F<sub>25</sub>\[Jy\] 151 M4(ms99 ) 0.28(3.6cm) 30 (fc,class,spc) 3 (parm) 152 M3-5(280 ) 280 (fc,spc,class) 153 K7/M0(ms99 ; 16 ) $`<`$0.22(3.6cm) 30 (fc,class,spc) 66 (parm) 154 M3(3 ) $`<`$0.3(6cm) 185 (fc) 1 (spc,class) 3 (parm) 155 M4,M2(5 ; 135 ) 0.21(3.6cm) 30 (fc,class,spc) 3 (parm) 156 M3,M4(ms99 ; 5 ) 0.14(3.6cm) 0.16 0.06L 30 (fc,class,spc) 336 (parm) 157 M5(16 ) 0.14(3.6cm) 30 (fc,class,spc) 3 (parm) 158 M4/K3(h60 ; 16 ) 0.30(3.6cm) 30 (fc,class,spc) 159 M4.5/M1(ms99 ; 16 ) $`<`$0.08(3.6cm) 30 (fc,class,spc) 136 (parm) 160 M4/M5(16 ) 0.66(3.6cm) 1.68 0.69 30 (fc,class,spc) 137 (parm) 161 M7(ms99 ) 6 (spc,class) 162 M5,M7(ms99 ; 52 ) $`<`$0.07(3.6cm) 30 (fc,class) 163 M1(3 ) 182 (fc) 3 (class) 291 (spc,parm) 164 M3,M2(3 ; 186 ) $`<`$0.2(6cm) 0.46 1.01 131 (fc) 3 (class) 186 (spc,parm) 165 M3(3 ) 1 (fc,spc) 3 (class,parm) 166 M5(36 ; ms99 ) 0.66(3.6cm) 0.23 0.11 30 (fc,class) 87 (spc) 138 (parm) 167 M6.5,M7(36 ; ms99 ) 1.6(6cm) 565.0 191.0 30 (fc) 39 (spc) 31 (class) 139 (parm) 168 mid M(61 ) 61 (spc,class) 169 Mira,M7(ms99 ) 99.6(3.6cm) 106.0 75.5 30 (fc,class,spc) 143 (parm) 170 M5.5/M3(16 ; ms99 ) $`<`$12(6cm) 0.37 0.09 30 (fc,class,spc) 90 ; 330 (parm) 171 M5/M6(ms99 ; 16 ) 0.20(3.6cm) 0.14 0.05L 30 (fc,class,spc) 145 (parm) 172 M5.5(ms99 ) 2.47(3.6cm) 0.85 0.15 30 (fc,class) 146 (spc,parm) 173 M(61 ) 61 (spc,class) 174 Mira(22 );M7(ms99 ) 61.1(3.6cm) 42.88 34.17 30 (fc,class) 39 (spc) 133 (parm) 175 Mira(22 );M5/M6(16 ; ms99 ) 26.6(3.5cm) 19.8 15.8 30 (fc,class) m235 ; 150 (spc) 151 (parm) 176 M6,M6.5(ms99 ; 154 ) 0.50(3.6cm) 0.39 0.25L 30 (fc) 31 (class) 156 (spc) 155 (parm) 177 G6(158 ) $`<`$0.05(3.6cm) 0.03L 0.07 30 (fc,class,spc) 158 (parm) 178 mid M(61 ) 61 (spc,class) 179 late M(30 ) $`<`$0.04(3.6cm) 30 (fc,class) 134 (spc) 180 S5.5/2.5(267 ) $`<`$0.45(3.6) (class) 268 (parm) 181 M6,M6-7(ms99 ; 160 ) 1.81(3.6cm) 0.88 0.74 30 (fc,class) 39 (spc) 159 (parm) 182 K7/K3(ms99 ; 16 ) $`<`$0.15(3.6cm) 30 (fc,class,spc) 162 (parm) 183 Mira;M6(163 ) $`<`$0.06(3.6cm) 27.24 15.14 30 (fc,class) 163 (spc,parm) 184 G5(ms99 ) 1.84 1.65 61 (spc,class) 59 (parm) 185 M3,M4(36 ; ms99 ; m90 ) 8.15(6cm) 1.6 0.60 30 (fc,class,spc) 88 (parm) 186 M(m177 ) 185 (fc) 274 (spc,class) 187 M4.5(ms99 ) 1.74(3.6cm) 0.68 0.24 30 (fc,class) 39 (spc) 167 (parm) 188 Mira;M7,M8(169 ; ms99 ) 12.3(6cm) 1577.0 543.8 30 (fc,class) 169 (spc) 174 (parm)
Table 2. Suspected symbiotic stars
No. Name $`\alpha `$(2000) $`\delta `$(2000) l<sup>II</sup> b<sup>II</sup> $`V`$ $`K`$ IR IUE X IP<sub>max</sub> <sup>h</sup> <sup>m</sup> <sup>s</sup> ’ ” \[mag\] \[mag\] \[eV\] s01 RAW 1691 01 18 36.1 $``$72 42 24.0 15.3 12 $``$ 13.6 s02 \[BE74\] 583 05 26 54.0 $``$71 06 00.0 16.1 13.6 s03 StHA 55 05 46 42.0 +06 43 48.0 199.34 $``$11.12 13.5 13.6 s04 GH Gem 07 04 04.9 +12 02 12.0 203.57 +8.23 14.6 $`>`$9.7 $``$ s05 ZZ CMi 07 24 13.9 +08 53 51.7 208.64 +11.30 9.9 2.8 S 35.1 s06 NQ Gem 07 31 54.5 +24 30 12.5 194.63 +19.35 7.9 3.0 + $``$ 54.4 s07 WRAY 16$``$51 09 33 29.4 $``$46 34 49.0 271.35 +3.80 4.4 13.6 s08 Hen 3$``$653 11 25 32.5 $``$59 56 31.9 292.36 +1.16 12.5 5.4 S $``$ 29.6 s09 NSV 05572 12 21 52.5 $``$13 53 09.9 292.10 +48.36 15 13.6 s10 AE Cir 14 44 52.0 $``$69 23 35.9 312.67 $``$8.69 14.1 54.4 s11 V748 Cen 14 59 37.0 $``$33 25 23.9 331.51 +22.24 12.6 8.1 $``$ 13.6 s12 V345 Nor 16 06 44.3 $``$52 02 30.1 330.51 +0.05 11.4 $``$ 13.6 s13 V934 Her 17 06 34.5 +23 58 18.5 45.15 +32.99 7.8 3.3 + + 77.5 s14 Hen 3$``$1383 17 20 31.5 $``$33 09 55.7 353.53 +2.20 12.5 $``$ 24.6 s15 V503 Her 17 36 46.0 +23 18 18.0 47.00 +26.23 13.8 $``$ s16 WSTB 19W032 17 39 02.8 $``$28 56 35.0 359.24 +1.22 17.2 35.1 s17 WRAY 16$``$294 17 39 13.9 $``$25 38 06.0 2.06 +2.94 15.5 S 35.1 s18 AS 241 17 44 58.0 $``$38 18 12.9 351.92 $``$4.76 12.0 7.8 S $``$ 24.6 s19 DT Ser 18 01 52.0 $``$01 26 06.0 25.92 +10.34 15.4 54.4 s20 V618 Sgr 18 07 57.0 $``$36 29 35.9 355.77 $``$7.83 15.2 13.6 s21 AS 280 18 09 52.9 $``$33 19 41.9 358.77 $``$6.69 13.2 $`>`$9.4 S 54.4 s22 AS 288 18 12 48.0 $``$28 20 00.9 3.49 $``$4.87 8.4 D? 54.4 s23 Hen 2$``$379 18 16 17.4 $``$27 04 32.9 4.97 $``$4.96 12.5 9.3 35.1 s24 V335 Vul 19 23 14.2 +24 27 40.2 58.22 +4.40 11.8 13.6 s25 V850 Aql 19 23 34.6 +00 38 03.0 37.18 $``$6.85 5.0 S 13.6 s26 Hen 2$``$442 19 39 39.0 +26 30 42.0 61.80 +2.13 14 5.3 100 s27 IRAS 19558+3333 19 57 48.4 +33 41 15.9 69.98 +2.38 D? s28 V627 Cas 22 57 41.2 +58 49 14.9 108.66 $``$0.86 12.9 3.3 D 13.6
Table 2. Suspected symbiotic stars – cont.
No. Cool-star spectrum Radio IRAS IRAS References \[mJy\] F<sub>12</sub>\[Jy\] F<sub>25</sub>\[Jy\] s01 C(199 ) 198 (fc) 199 (spc) (class) s02 G-K(m139 ) 209 (fc) m139 (class) s03 C(61 ) 61 (spc) (class) s04 F2:(178 ) 178 (fc) 31 (class) s05 M6,M5(48 ; 221 ) 6.60 2.65 319 (fc) 48 (spc) (class) s06 C6.2(47 ) 3.28 0.84 226 (spc) 224 (class) 225 (parm) s07 S?,M?(179 ) 2.33 0.69 (class) s08 late M,M6,M(m16 ; 2 ; 291 ) 0.90 0.28L 30 (fc,class) 291 (spc) s09 late M,M9(56 ; 236 ) 56 (fc,spc) (class) s10 242 (fc) 240 (spc,class) s11 M3,M4(55 ; 177 ) 300 ; 55 (spc) 31 (class) 177 (parm) s12 M3-5(245 ) 131 (fc) (class) 245 (spc,parm) s13 M3(250 ) $`<`$0.8(20cm) 2.8 0.74 250 (fc,spc,class) m151 (parm) s14 Mep?(m19 ) 2.54(6cm) (class) s15 M2p(m119 ) 178 (fc) m119 (spc) 31 (class) s16 K0-M3(232 ) 21(6cm) 3.42 12.57 251 (fc) 232 (spc) (class) s17 K5(3 ) 3 (spc,class,parm) s18 M1(3 ) $`<`$0.21(6cm) 0.67L 0.28L 30 (fc,spc,class) s19 G?,G2-K0(280 ; m122 ) 280 (fc,spc) (class) s20 279 (spc,class) s21 185 (fc) 3 (class,spc) s22 26.4(6cm) 1.24 1.10 185 (fc) (class) s23 K2,G-K(258 ; m23 ) 9.0(6cm) 0.69L 1.67: 185 (fc) 259 (spc) (class) s24 Mira?,C(m172 ; 287 ) 1.79 0.62 m172 (fc) 287 (spc,class) s25 Mira?,late(m41 ; m22 ) 4(2cm) 5.96 2.66 185 (fc) (class) s26 6.0(6cm) 185 (fc) (class) m176 (spc,parm) s27 OH/IR (si94 ) 0.44 (3.6cm) 42.9 46.0 si94 (class,parm) s28 Mira M2,M4(m192 ; m228 ) $`<`$9(2cm) 189.4 170.4 308 (fc) m192 (spc,class) m190 (parm)
Table 3. Orbital photometric ephemerides for symbiotic and suspected symbiotic stars
No. Star Ephemeris Eclipse Ref. 003 EG And Min(UV) = JD 2445380 + 481$`\times `$E Yes \[337 \] 008 AX Per Min($`\mathrm{m}_{\mathrm{vis}}/\mathrm{m}_{\mathrm{pg}}`$) = JD 2436667 + 680.8$`\times `$E Yes \[37 \] 017 V1261 Ori Min(Vyb) = JD 2445494 + 642$`\times `$E Yes \[213 ; 318 \] 023 BX Mon Min($`\mathrm{m}_{\mathrm{pg}}`$) JD 2449530 + 1401$`\times `$E Yes \[343 \] 024 V694 Mon Min($`\mathrm{m}_{\mathrm{pg}}`$/B) = JD 2437455 + 1930$`\times `$E No \[344 \] 033 SY Mus Min($`\mathrm{m}_{\mathrm{vis}}`$) = JD 2450176 + 625$`\times `$E Yes \[345 \] 045 RW Hya Min(UV) = JD 2445072 + 370.2$`\times `$E Yes \[m11 \] 050 V417 Cen Min($`\mathrm{m}_{\mathrm{pg}}`$) = JD 2427655 + 245.7$`\times `$E ? \[239 \] 057 T CrB Min(UBV) = JD 2435687.6 + 227.67$`\times `$E No \[338 \] 058 AG Dra Max(U) = JD 2443886 + 554$`\times `$E No \[342 \] 071 CL Sco Min($`\mathrm{m}_{\mathrm{pg}}`$) = JD 2427020 + 624.7$`\times `$E No \[31 \] 092 RT Ser Max(B) = JD 2446600 + 4500$`\times `$E No \[99 \] 093 AE Ara Min($`\mathrm{m}_{\mathrm{vis}}`$) = JD 2448571 + 830$`\times `$E ? \[100 \] 132 YY Her Min(V) = JD 2448945 + 590$`\times `$E No \[111 \] 133 V2756 Sgr Min($`\mathrm{m}_{\mathrm{pg}}`$) = JD 2437485 + 243$`\times `$E Yes? \[107 \] 134 FG Ser Min(B) = JD 2448492 + 650$`\times `$E Yes \[113 \] 142 AR Pav Min($`\mathrm{m}_{\mathrm{vis}}/\mathrm{m}_{\mathrm{pg}}`$) = JD 2411265.6 + 604.46$`\times `$E Yes \[119 \] 145 V443 Her Min(U) = JD 2443660 + 594$`\times `$E No \[124 \] 149 V2601 Sgr Min($`\mathrm{m}_{\mathrm{pg}}`$) = JD 2429850 + 850$`\times `$E Yes? \[347 \] 156 FN Sgr Min($`\mathrm{m}_{\mathrm{vis}}`$) = JD 2441750 + 567.2$`\times `$E Yes \[336 \] 160 V1413 Aql Min(B) = JD 2446650 + 434.1$`\times `$E Yes \[350 \] 166 BF Cyg Min(V) = JD 2415058 + 756.8$`\times `$E Yes \[87 \] 167 CH Cyg Min(U) = JD 2446275 + 5700$`\times `$E Yes \[339 \] 172 CI Cyg Min($`\mathrm{m}_{\mathrm{pg}}`$/B) = JD 2442687.1 + 855.6$`\times `$E Yes \[164 \] 176 PU Vul Min($`\mathrm{m}_{\mathrm{vis}}`$) = JD 2444550 + 4900$`\times `$E Yes \[348 \] 177 LT Del Min(U) = JD 2445910 + 478.5$`\times `$E No \[341 \] 181 V1329 Cyg Min(V) = JD 2444890.0 + 956.5$`\times `$E Yes \[159 \] 185 AG Peg Max(V) = JD 2442710.1 + 816.5$`\times `$E No \[346 \] 187 Z And Min($`\mathrm{m}_{\mathrm{vis}}`$) = JD 2442666 + 758.8$`\times `$E No \[349 \] s11 V748 Cen Min(V) = JD 2441917 + 566.5$`\times `$E Yes \[300 \]
Table 4. Orbital elements for symbiotic binaries
No. Star $`P`$ $`K_g`$ $`qM_\mathrm{g}/M_\mathrm{h}`$ $`\gamma _0`$ $`e`$ T<sub>0</sub> $`a_g`$sin$`i`$ $`f(M)`$ Ref. \[days\] \[km/s\] \[km/s\] \[JD\] \[R\] \[M\] 003 EG And 482.6 7.3 $``$95 0 2450804<sup>(3)</sup> 70 0.020 \[373 \] 008 AX Per 680.8 7.5 2.5 $``$116.5 0 2445511.8 100 0.030 \[37 \] 011 BD Cam 596.21 8.5 $``$22.3 0.09 2442794<sup>(1)</sup> 99.7 0.037 \[205 \] 017 V1261 Ori 642 7.5 3.0 79.7 0.07 2446778<sup>(3)</sup> 95 0.028 \[333 \] 023 BX Mon 1401 4.3 6.7 29.1 0.49 2449530 104 0.0076 \[343 \] 1259 4.6 29.1 0.44 2449680 103 0.0092 \[373 \] 033 SY Mus 624.5 7.4 12.9 0 2449082<sup>(3)</sup> 91 0.026 \[70 \] 036 TX CVn 199 5.7 2.3 0.16 2445195.25<sup>(1)</sup> 22 0.004 \[72 \] 045 RW Hya 370.2 8.8 12.4 0 2445071.6 65 0.026 \[m11 \] 370.4 8.8 12.9 0 2449512 0.026 \[73 \] 051 BD$``$21 3873 281.6 10.6 203.9 0 2449087.3<sup>(3)</sup> 59 0.035 \[81 \] 057 T CrB 227.57 23.9 0.6 $``$27.8 0 2447918.6 107 0.322 \[373 ; 46 \] 058 AG Dra 554 5.1 $``$148.3 0 2446366.4 56 0.008 \[86 \] 101 RS Oph 455.7 16.7 $``$40.2 0 2450154.1<sup>(3)</sup> 150 0.221 \[373 \] 134 FG Ser 650 8.3 2.8 71.2 0 2448491 107 0.039 \[Metal00 \] 142 AR Pav 604.5 12.8 2 $``$68.2 0 2420331.3<sup>(2)</sup> 153 0.13 \[Qetal00 \] 145 V443 Her 594 3.2 $``$49.2 0 2446007.7 39 0.002 \[125 \] 167 CH Cyg 5700 4.9 $``$57.7 0.47 2445086 478 0.045 \[339 \] 756.0 2.6 $``$60.6 0 2446643.7 39 0.0014 \[353 \] 5292 4.8 0.06 2445592<sup>(1)</sup> 500 0.060 \[353 \] 172 CI Cyg 855.3 6.7 3 18.4 0 2445241.8 114 0.027 \[146 \] 853.8 6.7 15.0 0.11 2450426.4 112 0.026 \[373 \] 182 CD$``$43 14304 1448 4.4 27.6 0 2445929<sup>(3)</sup> 126 0.013 \[162 \] 1442 4.6 27.5 0.22 2445560<sup>(1)</sup> 128 0.014 \[162 \] 185 AG Peg 816.5 5.3 4 $``$15.9 0 2431667.5 84 0.012 \[356 \] 818.2 5.4 $``$15.9 0.11 2446812 87 0.0135 \[373 \] 187 Z And 758.8 6.8 $``$1.8 0 2445703.0 102 0.025 \[167 \] <sup>(1)</sup> T<sub>0</sub> – time of the passage through periastron
<sup>(2)</sup> assumed from photometric ephemeris (eclipse)
<sup>(3)</sup> time of maximum velocity
Table 5. Pulsations ephemerides for Miras in symbiotic and suspected symbiotic stars
No. Star T<sub>0</sub> \[JD\] P \[days\] Ref. 010 o Ceti Max(V)=2444839 331.96 \[357 \] 016 UV Aur Max($`\mathrm{m}_{\mathrm{pg}}`$)=2441062 395.42 \[357 \] 026 RX Pup Min(J)=2440214 578 \[Metal99 \] 029 KM Vel ? 370 \[m49 \] 030 V366 Car ? 433 \[m49 \] 034 BI Cru ? 280 \[63 \] 048 V852 Cen ? 400 \[m2 \] 049 V835 Cen ? 450 \[m49 \] 060 V347 Nor ? 370-380 \[22 \] 100 H 1$``$36 ? 450-500 \[m2 \] 169 HM Sge Max(J)=2446856 527 \[m108 \] 174 V1016 Cyg Min(K)=2444852 478 \[335 \] Max(K)=2445038 478 \[335 \] 175 RR Tel Max(J/K)=2442999 387 \[Fetal83 \] 183 V407 Cyg Max(B)=2429710 745 \[163 \] 188 R Aqr Max(V)=2442398 386.96 \[357 \] s25 V850 Aql Max(V)=2425888 320 \[357 \] s28 V627 Cas ? 466 \[m192 \]
Table 6. Hipparcos parallaxes for symbiotic stars \[365 \]
No. Star $`\pi `$ \[mas\] 003 EG And 1.48 $`\pm `$ 0.97 010 o Ceti 7.79 $`\pm `$ 1.07 011 BD Cam 6.27 $`\pm `$ 0.63 017 V1261 Ori 1.32 $`\pm `$ 0.99 057 T CrB $``$1.61 $`\pm `$ 0.63 058 AG Dra $``$1.72 $`\pm `$ 0.98 142 AR Pav 3.37 $`\pm `$ 2.47 167 CH Cyg 3.73 $`\pm `$ 0.85 172 CI Cyg $``$0.36 $`\pm `$ 1.58 185 AG Peg $``$0.30 $`\pm `$ 1.17 187 Z And 2.34 $`\pm `$ 2.91 188 R Aqr 5.07 $`\pm `$ 3.15
Table 7. Flickering and outburst characteristics of symbiotic and suspected symbiotic stars. SyN – symbiotic nova, SyRN – symbiotic recurrent nova, Z And – Z And type outburst.
No. Star Flick. Ref. Out. Type of Out. Ref. 002 SMC2 + \[34 \] 004 SMC3 + \[34 \] 008 AX Per + Z And \[37 \] 010 o Ceti + \[m236 \] 012 S32 +? \[m29 \] 013 LMC S154 +? \[208 \] 024 V694 Mon + \[m112 \] + \[293 \] 026 RX Pup + SyRN? \[Metal99 \] 035 RT Cru + \[6 \] + \[6 \] 036 TX CVn + Z And \[324 \] 043 V840 Cen + \[238 \] 050 V417 Cen +? \[239 \] 057 T CrB + \[366 \] + SyRN \[128 \] 058 AG Dra + Z And \[86 ; 370 \] 068 KX Tra + \[326 \] 074 Hen 3$``$1341 + \[m230 \] 084 V2116 Oph + \[96 \] 092 RT Ser + SyN \[327 ; 85 \] 093 AE Ara + Z And? \[100 \] 096 V2110 Oph +? \[328 \] 097 V916 Sco +? \[102 \] 099 V917 Sco + \[103 \] 101 RS Oph + \[m112 \] + SyRN \[31 \] 107 Bl 3$``$14 + \[58 \] 110 V745 Sco + SyRN \[253 \] 128 V2506 Sgr + \[131 \] 132 YY Her + Z And \[111 \] 134 FG Ser + Z And \[112 \] 138 V4074 Sgr + \[329 \] 139 V2905 Sgr +? \[298 \] 142 AR Pav + Z And \[368 \] 143 V3929 Sgr +? \[208 \] 148 V3890 Sgr + SyRN \[m96 \] 156 FN Sgr + Z And \[100 \] 158 CM Aql + Z And \[372 \] 160 V1413 Aql + Z And \[137 ; 369 \] 164 V352 Aql + \[288 \] 166 BF Cyg + Z And \[138 \] 167 CH Cyg + \[m112 \] + \[m106 \] 169 HM Sge + SyN \[85 \] 172 CI Cyg + Z And \[146 \] 174 V1016 Cyg + SyN \[85 \] 175 RR Tel + SyN \[85 \] 176 PU Vul + SyN \[153 ; 85 \] 177 LT Del + \[157 \] 181 V1329 Cyg + SyN \[85 \] 183 V407 Cyg + \[18 \] 185 AG Peg + SyN \[88 ; 85 \] 187 Z And + \[367 \] + Z And \[167 \] 188 R Aqr + \[371 \] s11 V748 Cen + \[117 \] s12 V345 Nor + \[131 \] s13 V934 Her + \[m151 \]
Table 8. Different names for symbiotic and suspected symbiotic stars
001=SMC 1=NAME SMC1=\[MH95\] 183
002=SMC 2=NAME SMC2
003=EG And=HD 4174=BD+39 167=SAO 36618=GCRV 403=HIC 3494=GEN# +1.00004174= AG+40 66=GC 880=DO 8473=GPM1 20=SKY# 1157=AGKR 609=IRC +40014=JP11 413= PPM 43262=HIP 3494=IRAS 00415+4024
004=SMC 3=NAME SMC3=RX J0048.4$``$7332
005=SMC N60=LHA 115$``$N 60=LIN 323=HV 1707(???)
007=SMC N73=LHA 115$``$N 73=LIN 445 a
008=AX Per=MWC 411=HV 5488=CSI+54$``$01331=GCRV 896=JP11 5465=IRAS 01331+5359
009=V471 Per=PN M 1$``$2=PK 133$``$08 1=PN VV 8=LS V +52 1=PN G133.1$``$08.6=PN ARO 116=CSI+52$``$01555=PN VV’ 11=IRAS 01555+5239
010=o Cet=MIRA=HD 14386=RAFGL 318=SKY# 3428=GC 2796=omi Cet=ADS 1778 AP= IRC +00030=YZ 93 562=MWC 35=BD$``$03 353=GEN# +1.00014386J= CCDM J02194$``$0258AP=PLX 477=GCRV 1301=JP11 625=CSI$``$03 353 1=DO 430= HIC 10826=SAO 129825=68 Cet=HR 681=UBV 21604=LTT 1179=TD1 1361= PPM 184482=JOY 1AP=IRAS 02168$``$0312
011=BD Cam=HD 22649=BD+62 597=HR 1105=SAO 12874=GC 4383=IRC +60125=FK4 129= UBV 3468=CSS 79=\[HFE83\] 244=GEN# +1.00022649=AG+63 277=N30 751= GCRV 2027=RAFGL 506=PPM 14446=HIP 17296=SKY# 5606=UBV M 9615=PLX 758= JP11 803=CSV 328=HIC 17296=S1\* 60=IRAS 03377+6303
012=S32=StHA 32
013=LMC S154=LHA 120$``$S 154
014=LMC S147=LHA 120$``$S 147=\[BE74\] 484
015=LMC N19=LHA 120$``$N 19=\[BE74\] 191=LMC B0503$``$6803=MC4(0503$``$680)= PMN J0503$``$6757
016=UV Aur=HD 34842=BD+32 957=ADS 3934 A=IRC +30110=SAO 57941=HV 3322=Case 9= GEN# +1.0003482A=CGCS 911=IDS 05153+3224 A=AG+32 505=JP11 1034=UBV M 10852= AN 58.1911=LEE 179=RAFGL 735=PPM 70251=HIC 25050=GCRV 3199= DO 11210=CSI+32 957 1=Fuen C 29=CCDM J05218+3231A=C\* 318=IRAS 05185+3227
017=V1261 Ori=HD 35155=BD$``$08 1099=RAFGL 736=SAO 132035=GCRV 56103=YZ 98 1455= HIC 25092=GEN# +1.00035155=GC 6602=HERZ 11764=CSS 133=HIP 25092=SKY# 8520= IRC $``$10086=PPM 187990=S1\* 98=IRAS 05199$``$0842
019=LMC N67=CH$``$95=LHA 120$``$N 67=\[HC88\] 95
021=LMC S63=LHA 120$``$S 63=HV 12671
022=SMP LMC 94=LHA 120$``$S 170=LM 2$``$44
023=BX Mon=AS 150=HV 10446=MHA 61$``$12=JP11 5448=SS73 4
024=V694 Mon=MWC 560=SS73 5=LS 391=CSI$``$07$``$07234=GSC 05396$``$01135= MHA 61$``$14=IRAS 07233$``$0737
025=WRAY 15$``$157=IRAS 08045$``$2823
026=RX Pup=HD 69190=WRAY 16$``$17=HV 3372=CPD$``$41 2287=Hen 3$``$138=AN 88.1914=SS73 8= JP11 1653=PK 258$``$03 1=CD$``$41 3911=IRAS 08124$``$4133
027=Hen 3$``$160=SS73 9=WRAY 15$``$208
028=AS 201=MHA 382$``$43=Hen 3$``$172=SCM 38=PK 249+06 1=PDS 32=PN SaSt 1$``$1= PN G249.0+06.9=IRAS 08296$``$2735
029=KM Vel=Hen 2$``$34=WRAY 16$``$56=PK 274+02 1=ESO 212$``$13=PN G274.1+02.5= SS73 14=IRAS 09394$``$4909
030=V366 Car=Hen 2$``$38=PK 280$``$02 1=PN SaSt 1$``$3=ESO 167$``$11=WRAY 16$``$62= IRAS 09530$``$5704
031=Hen 3$``$461=PK 283+06 2=IRAS 10370$``$5108
033=SY Mus=HD 100336=SS73 32=HV 3376=WRAY 15$``$824=TD1 15706=Hen 3$``$667= AN 118.1914=CPD$``$65 11298=CSI$``$65$``$11298=SPH 127=IRAS 11299$``$6508
034=BI Cru=Hen 3$``$782=SVS 1855=IRAS 12206$``$6221
035=RT Cru=HV 1245=AN 131.1906
036=TX CVn=BD+37 2318=UBV M 2779=PPM 76696=JP11 5415=CDS 866=FB 110= Case A-F 788=AG+37 1208=SAO 63173=GEN# +0.03702318=IRAS 12423+3702
037=Hen 2$``$87=SS73 36=PK 302$``$00 1=WRAY 16$``$119=ESO 95$``$16=IRAS 12423$``$6244
038=Hen 3$``$828=SS73 37=WRAY 15$``$1022
039=SS73 38=CD$``$64 665=CGCS 3295=DJM 390=IRAS 12483$``$6443
041=St 2$``$22=PN Sa 3$``$22=PK 305+03 1
042=CD$``$36 8436=NSV 06160=Hen 3$``$886=JP11 311=PK 308+25 6=IRAS 13131$``$3644
043=V840 Cen=LILLER’S OBJECT
044=Hen 3$``$905=SS73 40=WRAY 15$``$1108
045=RW Hya=HD 117970=MWC 412=CPD$``$24 5101=YZ 115 9978=SAO 181760=CD$``$24 10977= GCRV 8034=HV 3237=PPM 261808=GEN# +1.00117970=JP11 2404=AN 51.1910= FAUST 3820=IRAS 13315$``$2507
046=Hen 3$``$916=SS73$``$42=WRAY 15$``$1123
047=V704 Cen=Hen 2$``$101=PK 311+03 1=PN G311.1+03.4=WRAY 16$``$141=ESO 133$``$7= IRAS 13515$``$5812
048=V852 Cen=Hen 2$``$104=PK 315+09 1=ESO 221$``$31=PN G315.4+09.4=SS73 43= NAME SOUTHERN CRAB=WRAY 16$``$147=SCM 78=IRAS 14085$``$5112
049=V835 Cen=Hen 2$``$106=WRAY 16$``$148=PK 312$``$02 1=ESO 97$``$14=SCM 79=NSV 06587= IRAS 14103$``$6311
050=V417 Cen=HV 6516=IRAS 14122$``$6139
052=Hen 2$``$127=SS73 45=PK 325+04 2=ESO 224$``$2=SCM 94=WRAY 16$``$180= IRAS 15210$``$5139
053=Hen 3$``$1092=Hen 2$``$134=SS73 46=PK 319$``$09 1=WRAY 16$``$186=LS 3391=ESO 99$``$10= JP11 5230=CSI$``$66$``$15425
054=Hen 3$``$1103=SS73 47=WRAY 15$``$1359
055=HD 330036=BD$``$48 10371=PN Cn 1$``$1=PK 330+04 1=CD$``$48 10371=HIP 77662= ESO 225$``$1=HIC 77662=WRAY 15$``$1364=PN G330.7+04.1=IRAS 15476$``$4836
056=Hen 2$``$139=PK 326$``$01 1=WRAY 16$``$193=ESO 178$``$2=IRAS 15508$``$5520
057=T CrB=HD 143454=MWC 413=GC 21491=IDS 15553+2612 AB=DO 15377=PPM 104498= GEN# +1.00143454J=BD+26 2765=GCRV 9203=SAO 84129=NOVA CrB 1866= CCDM J15595+2555AB=AG+26 1536=HR 5958=SBC 558=NOVA CrB 1946=HIC 78322= HIP 78322=IRAS 15574+2603
058=AG Dra=AG+66 715=CDS 889=SAO 16931=HIC 78512=1ES 1601+66.9= BPS BS 16087$``$0012=BD+67 922=GCRV 9231=UBV 13635=GEN# +0.06700922= 2E 1601.3+6656=2RE J1601+664=HIP 78512=JP11 236=SVS 1155=PPM 19692= 2E 3573=2RE J160133+664805=IRAS 16013+6656
059=WRAY 16$``$202=PN Sa 3$``$33=PK 332+01 1
060=V347 Nor=Hen 2$``$147=WRAY 16$``$208=PK 327$``$04 1=ESO 178$``$13=PN G327.9$``$04.3= IRAS 16099$``$5651
061=UKS Ce$``$1=UKS$``$Ce1=UKS Ce1=UKS 1612$``$22.0
062=QS Nor=Hen 2$``$156=SS73 53=WRAY 15$``$1461=PK 338+05 2=SCM 111=ESO 331$``$2= CSV 2635
063=WRAY 15$``$1470=Hen 3$``$1187=SS73 55
064=Hen 2$``$171=SS73 59=WRAY 16$``$226=PK 346+08 1=PN G346.0+08.5=ESO 390$``$7= IRAS 16307$``$3459
065=Hen 3$``$1213=SS73 60=WRAY 15$``$1511
066=Hen 2$``$173=SS73 61=WRAY 15$``$1518=PK 342+05 1=ESO 331$``$7
067=Hen 2$``$176=PK 339+00 1=WRAY 16$``$230=ESO 277$``$5=IRAS 16379$``$4507
068=KX TrA=Hen 3$``$1242=PN Cn 1$``$2=Hen 2$``$177=PK 326$``$10 1=WRAY 16$``$233= ESO 101$``$10=JP11 5232=PN StWr 1$``$1=JP11 5233=IRAS 16401$``$6231
069=AS 210=MHA 276$``$52=Hen 3$``$1265=SS73 66=JP11 5249=WRAY 16$``$237=PK 355+11 1= IRAS 16482$``$2555
070=HK Sco=HV 4493=AS 212=MHA 71$``$6=Hen 3$``$1280=SS73 68=WRAY 15$``$1563
071=CL Sco=HV 4035=AS 213=MHA 71$``$5=Hen 3$``$1286=SS73 69=WRAY 15$``$1564= CD$``$30 13603
072=PN MaC 1$``$3=PK 339$``$03 1
073=V455 Sco=HV 7869=AS 217=MHA 71$``$15=Hen 3$``$1334=SS73 74=PN H 2$``$2= PK 351+03 1=ESO 392$``$3=WRAY 16$``$252=Haro 3$``$2=PN VV’ 162=IRAS 17040$``$3401
074=Hen 3$``$1341=NSV 08226=SS73 75
075=Hen 3$``$1342=SS73 77
076=AS 221=MHA 276$``$12=Hen 3$``$1348=SS73 79=PN H 2$``$4=PK 352+03 1=WRAY 15$``$1637= PN VV’ 166=Haro 3$``$4=ESO 392$``$4=IRAS 17087$``$3230
077=PN H 2$``$5=PK 354+04 2=PN VV’ 172=ESO 454$``$3=WRAY 15$``$1655=Haro 3$``$5= IRAS 17121$``$3131
078=PN Sa 3$``$43=PK 355+04 1
079=Draco C$``$1=Irwin Dra 22025=Stetson 3203=McClure C1=\[ALS82\] C1
080=PN Th 3$``$7=PK 356+04 3=PN ARO 249=ESO 454$``$12
081=PN Th 3$``$17=PK 357+03 3=WRAY 17$``$81=ESO 454$``$30=PN Sa 3$``$54
082=PN Th 3$``$18=PK 358+03 5=ESO 454$``$32=WRAY 17$``$83=Hen 2$``$228
083=Hen 3$``$1410=PN Th 3$``$20=PK 357+02 3=PN ARO 251=NSV 08805=WRAY 15$``$1714= ESO 454$``$37
084=V2116 Oph=GX 1+04=PK 001+04 1=4U 1728$``$24=1M 1728$``$247=H 1728$``$247= PN MaC 1$``$5=2S 1728$``$247=GX 2+05=2U 1728$``$24=3U 1728$``$24=3A 1728$``$247= 1H 1728$``$247
085=PN Th 3$``$29=WRAY 17$``$89=PN Sa 3$``$61=PK 358+02 3=ESO 455$``$15= IRAS 17299$``$2905
086=PN Th 3$``$30=PK 359+02 1=ESO 455$``$18=WRAY 17$``$90=Hen 2$``$243=PN Sa 3$``$63
087=PN Th 3$``$31=PK 358+01 2=ESO 455$``$19=Hen 2$``$245=WRAY 17$``$91= IRAS 17312$``$2926
088=PN M 1$``$21=SS73 90=Hen 2$``$247=PK 006+07 1=PN VV’ 216=PN VV 103= ESO 588$``$7=IRAS 17313$``$1909
089=Hen 2$``$251=PK 358+01 3=PN H 1$``$25=PN VV’ 218=Haro 2$``$25=WRAY 16$``$289= SS73 92=SCM 151=ESO 455$``$22=PN Bl A=IRAS 17321$``$2943
090=PN Pt 1=PK 003+03 3=EQ 1735$``$2352=ESO 520$``$12
091=PN K 6$``$6=WR 99=LuSt 1=PN G359.8+01.5=Terz V 2955
092=RT Ser=MWC 265=SS73 94=CSI$``$11$``$17371=NOVA Ser 1909=NOVA Ser 1917= GCRV 10203=AN 7.1917
093=AE Ara=HV 5491=MWC 591=Hen 3$``$1451=SS73 95=WRAY 15$``$1754=PN PC 18= PK 344$``$08 1=GSC 08347$``$01978=ESO 279$``$5
095=UU Ser=AS 237=SS73 98=AN 720.1936=HV 8771=CSV 2420
096=V2110 Oph=AS 239=MHA 79$``$52=Hen 3$``$1465=JP11 5250
097=V916 Sco=SSM 1=EQ 1740$``$360
098=Hen 2$``$275=WRAY 16$``$304=PK 351$``$05 1=ESO 334$``$8
099=V917 Sco=Hen 3$``$1481=SS73 103
100=PN H 1$``$36=PK 353$``$04 1=Hen 2$``$289=PN Sa 2$``$249=PN G353.5$``$04.9= ESO 393$``$31=PN VV’ 259=Haro 2$``$36=IRAS 17463$``$3700
101=RS Oph=HD 162214=MWC 414=BD$``$06 4661=NOVA Oph 1898=AN 20.1901= PPM 201101=GCRV 10316=HV 164=SS73 106=JP11 2898=NOVA Oph 1933= IRAS 17474$``$0641
102=WRAY 16$``$312=PN Sa 3$``$80=PK 358$``$01 2
103=V4141 Sgr=PN Th 4$``$4=PK 008+03 2=ESO 589$``$10=PN ARO 260=IRAS 17477$``$1948
104=ALS 2=\[ALS88\] 2=SS 324=EQ 174755.9$``$174710=PK 010+04 7
105=AS 245=MHA 359$``$110=Hen 3$``$1501=SS73 107=PN H 2$``$28=PK 006+02 2= ESO 589$``$13=PN VV’ 268=PN ARO 261=Haro 3$``$28=IRAS 17479$``$2218
106=Hen 2$``$294=PK 357$``$03 1=ESO 394$``$1=WRAY 16$``$318=IRAS 17483$``$3250
107=PN Bl 3$``$14=PK 000$``$01 4=ESO 455$``$55
108=PN Bl 3$``$6=PK 358$``$02 2=PN Sa 3$``$85=ESO 456$``$4=IRAS 17496$``$3120
109=PN Bl L=PK 359$``$02 1=PN Sa 3$``$86=ESO 456$``$6
110=V745 Sco=NOVA Sco 1937
111=PN MaC 1$``$9=PK 013+05 2
112=AS 255=MHA 363$``$45=Hen 3$``$1525=SS73 111=IRAS 17537$``$3513
113=V2416 Sgr=PN M 3$``$18=SS73 112=PK 007+01 2=ESO 589$``$24=PN VV’ 295= Hen 2$``$312=Ve 2$``$61=IRAS 17542$``$2142
114=PN H 2$``$34=PN Sa 3$``$105=PK 001$``$02 1=ESO 456$``$28=PN VV’ 304=Haro 3$``$34
116=AS 269=PN H 1$``$49=SS73 119=PK 358$``$05 2=Hen 2$``$331=MHA 304$``$52=ESO 394$``$23= WRAY 16$``$356=Haro 2$``$49=PN VV’ 323=SCM 178=IRAS 18001$``$3242
117=PN Ap 1$``$8=SS73 121=PK 002$``$03 1=ESO 456$``$47=IRAS 18013$``$2821
118=SS73 122=PN KFL 6=IRAS 18015$``$2709
119=AS 270=Hen 3$``$1581=SS73 126=IRAS 18026$``$2025
120=PN H 2$``$38=SS73 128=PK 002$``$03 4=Hen 2$``$343=WRAY 17$``$108=ESO 456$``$51= PN VV’ 342=Haro 3$``$38=IRAS 18028$``$2817
121=SS73 129=PK 001$``$04 3=T 17=PN KFL 8=IRAS 18038$``$2932
122=Hen 3$``$1591=SS73 132=T 53=NSV 10219=IRAS 18044$``$2558
123=V615 Sgr=Hen 2$``$349=PK 356$``$07 1=SS73 131=ESO 394$``$29=WRAY 15$``$1840= HV 7199=IRAS 18044$``$3610
124=Ve 2$``$57=SS 134=NSV 10241
125=AS 276=Hen 3$``$1595=SS73 135=MHA 363$``$7=IRAS 18058$``$4108
126=PN Ap 1$``$9=SS73 137=PK 003$``$04 2=Hen 2$``$356=WRAY 16$``$377=ESO 456$``$63= SCM 186=IRAS 18073$``$2753
127=AS 281=MHA 208$``$83=SS73 138=PN Ap 1$``$10=PK 003$``$04 1=WRAY 16$``$378= ESO 456$``$65=Hen 2$``$357=IRAS 18076$``$2757
128=V2506 Sgr=AS 282=MHA 304$``$113=Hen 2$``$358=SS73 139=PN Ap 1$``$11= PK 003$``$04 6=WRAY 16$``$379=ESO 456$``$66
129=SS73 141=PK 359$``$07 2=WRAY 16$``$384
130=AS 289=Hen 3$``$1627=SS73 143=F1$``$11=IRAS 18095$``$1140
131=Y CrA=HD 166813=Hen 3$``$1632=CSI$``$42$``$18107=SS73 144=HV 169= PK 350$``$11 1=AN 25.1901=IRAS 18110$``$4252
132=YY Her=AS 297=CSI+20$``$18124=GSC 01579$``$00381=AN 6.1919=JP11 5444= MHA 352$``$34
133=V2756 Sgr=AS 293=MHA 304$``$122=Hen 2$``$370=SS73 145=PK 002$``$05 1=ESO 456$``$79= WRAY 16$``$392
134=FG Ser=AS 296=D 143$``$2=SS73 148=SON 10363=JP11 5251=IRAS 18125$``$0019
135=HD 319167=PN CnMy 17=Hen 2$``$373=SS73 146=PK 001$``$06 1=WRAY 16$``$395= SCM 195=ESO 456$``$81
136=Hen 2$``$374=SS73 147=PK 009$``$02 1=ESO 590$``$10=IRAS 18126$``$2135
137=Hen 2$``$376=AS 294=NSV 10435=PK 004$``$05 2=MHA 304$``$123=WRAY 16$``$396= PN Sa 3$``$126=ESO 457$``$1
138=V4074 Sgr=AS 295B=Hen 3$``$1641=HIC 89526=HIP 89526
139=V2905 Sgr=AS 299=SS73 151=GEN# +6.20010299=MHA 208$``$92
141=Hen 3$``$1674=SS73 153=T 21=WRAY 15$``$1864=PN KFL 17
142=AR Pav=MWC 600=CPD$``$66 3307=GCRV 10756=HIC 89886=HV 7860=PPM 363277= HIP 89886=SBC 668=GSC 09080$``$00788=IRAS 18157$``$6609
143=V3929 Sgr=Hen 2$``$390=SS73 154=CSV 4026=PK 005$``$05 2=ESO 522$``$19= PN ARO 274=PN StWr 2$``$3=SCM 202=HV 9397=P 4629=IRAS 18178$``$2649
144=V3804 Sgr=AS 302=MHA 304$``$33=Hen 3$``$1676=SS73 155=IRAS 18182$``$3133
145=V443 Her=MWC 603=GCRV 68111=CSI+23$``$18201=JP11 5216=FB 171= IRAS 18200+2325
146=V3811 Sgr=Hen 2$``$396=ESO 590$``$19=PK 010$``$03 1=SS73 160=IRAS 18206$``$2157
147=V4018 Sgr=CD$``$28 14567=AS 304=Hen 3$``$1691=SS73 162=PN KFL 20= GSC 06869$``$00806=IRAS 18221$``$2837
148=V3890 Sgr=NOVA Sgr 1962=SS 390
149=V2601 Sgr=AS 313=MHA 208$``$51=SS73 171=IRAS 18349$``$2244
150=PN K 3$``$9=PN Sa 3$``$142=PK 023$``$01 1=IRAS 18376$``$0846
151=AS 316=MHA 208$``$58=Hen 2$``$417=SS73 172=PK 012$``$07 1=ESO 591$``$14
152=DQ Ser=CSV 4342
153=MWC 960=MHA 204$``$22=Hen 3$``$1726=SS73 174
154=AS 323=PN K 4$``$7=PK 026$``$02 2=MHA 369$``$39=PN ARO 292
155=AS 327=MHA 208$``$67=Hen 3$``$1730=SS73 176=PK 011$``$11 1=JP11 5253=WRAY 16$``$421
156=FN Sgr=AS 329=SS73 177=CSI$``$19$``$18509=NOVA Sgr 1925=GCRV 11342=JP11 5475
157=PN Pe 2$``$16=PK 029$``$02 1=PN Th 1$``$F=PN ARO 296=PN VV’ 464
159=V919 Sgr=AS 337=MHA 227$``$6=SS73 178=AN 237.1932
160=V1413 Aql=AS 338=MHA 305$``$6=Hen 3$``$1737=PN K 4$``$12=PK 048+04 1=SS 428= IRAS 19015+1625
162=PN Ap 3$``$1=PK 037$``$02 1=SCM 231=PN ARO 145
163=PN MaC 1$``$17=PK 030$``$07 1
164=V352 Aql=PN K 3$``$25=PK 037$``$03 3=PN ARO 306=AN 279.1931=IRAS 19111+0212
165=ALS 1=\[ALS88\] 1=EQ 191333.1$``$082308=PK 028$``$09 3
166=BF Cyg=MWC 315=LS II +29 5=JP11 5433=CSI+29$``$19219 2=GCRV 11847= AN 112.1914
167=CH Cyg=HD 182917=SAO 31632=BD+49 2999=GCRV 11865=JP11 3103=RAFGL 2383= HIC 95413=GEN# +1.00182917=AG+50 1370=DO 37228=YZ 50 6001=HIP 95413= SKY# 36122=GC 26820=CDS 1064=IRC +50294=AN 24.1924=PPM 37375= IRAS 19232+5008
169=HM Sge=PK 053$``$03 2=2E 1939.7+1637=SVS 2183=2E 4280=IRAS 19396+1637
170=Hen 3$``$1761=NSV 12264=IRAS 19373$``$6814
171=QW Sge=AS 360=MHA 80$``$5=Hen 3$``$1771=JP11 5254=NSV 12383
172=CI Cyg=MWC 415=HV 3625=CSI+35$``$19484=JP11 5434=HIP 97594=GCRV 12195= AN 10.1922=HIC 97594=IRAS 19483+3533
173=StHA 169=NSV 12466
174=V1016 Cyg=AS 373=CSI+39$``$19553=JP11 5437=2E 1955.3+3941=PK 075+05 1= JP11 5438=2E 4302=GCRV 70112=MHA 328$``$116=IRAS 19553+3941
175=RR Tel=Hen 3$``$1811=SKY# 37701=AN 166.1908=2E 4313=CSI$``$55$``$20003= GCRV 6924 E=HV 3181=2E 2000.3$``$5552=IRAS 20003$``$5552
176=PU Vul=NOVA Vul 1979=KUWANO-HONDA=IRAS 20189+2124
177=LT Del=Hen 2$``$467=PK 063$``$12 1=PN ARO 353=StHA 179
179=Hen 2$``$468=PK 075$``$04 1=PN ARO 355
180=ER Del=BD+08 4506=AG+08 2842=SVS 654
181=V1329 Cyg=HBV 475=PK 077$``$05 1=UBV M 46710=VES 248=IRAS 20492+3518
182=CD$``$43 14304=Hen 3$``$1924
183=V407 Cyg=AS 453=MHA 289$``$90=AN 148.1940=NOVA Cyg 1936
184=StHA 190=RJHA 120=IRAS 21392+0230
185=AG Peg=HD 207757=MWC 379=AGKR 19529=JP11 3412=TD1 28516=SKY# 41636= HIC 107848=GEN# +1.00207757=BD+11 4673=SAO 107436=DO 7622=HIP 107848= AG+12 2570=PLX 5279=GCRV 13724=YZ 12 8693=SBC 879=PPM 140717= IRAS 21486+1223
186=LL Cas=PN K 4$``$46=PK 108$``$05 1=AN 207.1940=PN ARO 379
187=Z And=HD 221650=MWC 416=PLX 5697=SAO 53146=2E 2331.6+4834=GEN# +1.00221650= BD+48 4093=GCRV 14773=HV 193=PPM 64386=2E 4735=AG+48 2087=JP11 3636= AN 41.1901=HIC 116287=HIP 116287=IRAS 23312+4832
188=R Aqr=HD 222800=SAO 165849=HR 8992=GCRV 14862=MWC 400=SKY# 44830= GC 32948=YZ 105 8733=PPM 242022=HIP 117054=GEN# +1.00222800=IRC $``$20642= BD$``$16 6352=RAFGL 3136=HIC 117054=IRAS 23412$``$1533
s01=RAW 1691=\[MA93\] 1858=LIN 521
s04=GH Gem=CSV 950=AN 241.1943
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# Three Dimensional Radiative Transfer
## 1 One Motivation: Structure Formation and the beginning of the bright ages
The atomic nuclei in the primordial gas (mostly hydrogen and helium) first (re)combined with electrons at a redshift $`z1000`$. From the study of absorption spectra of high redshift quasars we know that this then neutral gas must have been ionized prior to $`z=5`$. Most likely this reionization process was caused via photoionization by UV photons produced in proto galactic objects either by massive stars or by the accretion onto compact objects. The formation of these first objects in the universe and their potential impact on subsequent structure formation is a highly topical issue in physical cosmology to date. In our standard models of structure formation cosmological objects form via hierarchical build up from smaller pieces. The dynamics is controlled by gravity of a dominant cold dark matter (CDM) component. Baryons will fall into virializing CDM halos in which they may cool and possibly fragment to form stars. The lower limit on the masses of luminous objects that may be formed is determined by (1) the pressure of the primordial gas, which determines whether it can settle in the dark matter halo (2) the ability of the baryons to cool to collapse to stellar densities.
Both these issues depend sensitively on the presence of UV photons (see Abel and Haehnelt 1999, Haiman, Abel and Rees 1999, and references therein). Since the intergalactic medium (IGM) is initially optically thick to $`h\nu >13.6`$ eV photons ionization fronts will be formed around the first sources. Because we believe that the radiation is produced in structures condensed from the IGM by gravitational instability, the first UV photons will see a clumpy inhomogeneous IGM. As a consequence the time–varying ionized regions will have complex morphologies.
The above physical processes have prompted us to develop methods for the treatment of RT in three dimensional cosmological hydrodynamics. In the following we describe the conditions under which the cosmological RT becomes equivalent to classical RT. Then we go on to discuss some possible methods to solve the latter. In this contribution I will only give a brief overview of my and my collaborators work. However, note that also Razoumov and Scott (1999) offer a different approach. Gnedin (1999) chose to use dramatic simplifications constructed such as to mimic the expected effects of radiative transfer in order to study the reionization of intergalactic hydrogen.
## 2 Cosmological Radiative Transfer
The equation of cosmological radiative transfer in comoving coordinates (cosmological, not fluid) is:
$$\frac{1}{c}\frac{I_\nu }{t}+\frac{\widehat{n}I_\nu }{\overline{a}}\frac{H(t)}{c}(\nu \frac{I_\nu }{\nu }3I_\nu )=\eta _\nu \chi _\nu I_\nu $$
(1)
where $`I_\nu I(t,\stackrel{}{x},\stackrel{}{\mathrm{\Omega }},\nu )`$ is the monochromatic specific intensity of the radiation field, $`\widehat{n}`$ is a unit vector along the direction of propagation of the ray; $`H(t)\dot{a}/a`$ is the (time-dependent) Hubble constant, and $`\overline{a}\frac{1+z_{em}}{1+z}`$ is the ratio of cosmic scale factors between photon emission at frequency $`\nu `$ and the present time t. The remaining variables have their traditional meanings (e.g, Mihalas 1978.) Equation (1) will be recognized as the standard equation of radiative transfer with two modifications: the denominator $`\overline{a}`$ in the second term, which accounts for the changes in path length along the ray due to cosmic expansion, and the third term, which accounts for cosmological redshift and dilution.
One could, in principle, attempt to solve equation (1) directly for the intensity at every point given the emissivity $`\eta `$ and absorption coefficient $`\chi `$. However, the high dimensionality of the problem (three positions, two angles, one frequency and time $`=`$ 7D!) not to mention the high spatial and angular resolution needed in cosmological simulations would make this approach impractical for dynamic computations. Therefore we proceed through a sequence of well-motivated approximations which reduce the complexity to a tractable level.
### 2.1 Local quasi–static Approximation
We begin by eliminating the cosmological terms and factors. That we can do this can be understood on simple physical grounds. Before the universe is reionized, it is opaque to H and He Lyman continuum photons. Consequently, ionizing sources are local to scales of interest, and not at cosmological distances. In particular this means that the term multiplicative term $`\frac{H(t)}{c}`$ which is simply the reciprocal of the Hubble horizon at the time $`t`$ will ensure the cosmological term to be small as long as the opacity is much smaller than the horizon scale. If this is not the case and we have a simulation box size much smaller than the mean free path than we are at the limit where the cosmological terms will modify the boundary values but still not be important as the photons transverse the box. Therefore, setting $`\overline{a}1`$, equation (1) reduces to its standard, non-cosmological form:
$$\frac{1}{c}\frac{I_\nu }{t}+\widehat{n}I_\nu =\eta _\nu \chi _\nu I_\nu $$
(2)
where now $`\nu `$ is the instantaneous, comoving frequency.
Thinking of the special case of a point source that switches on instantaneously one realizes that initially the ionization front will always propagate at the speed of light. Eventually it slows down as the ’incoming flux’ of neutrals grows with the increasing ionization surface. This will ensure that eventually the light crossing time (1/(c times I-front distance)) will become much shorter than the timescales of change of the emissivities and absorption coefficient on the right hand side of equation (2). At this point an explicit integration can employ large time steps making the term $`\frac{1}{c}\frac{I_\nu }{t}`$ negligible. From this point of the evolution on it will suffice to solve the static classic equation of radiative transfer ,
$$\widehat{n}I_\nu =\eta _\nu \chi _\nu I_\nu .$$
(3)
## 3 Methods
### 3.1 Brute Force: Ray Tracing
Abel, Norman and Madau (1999) give a method that integrates this quasi–static approximation along rays casted from point–sources. That method has the particular advantage that it will ensure photon conservation independent of resolution by exploiting the known analytic solution of radiative transfer for a homogeneous slab.
Consider the simple case of only absorption. Then across a computational cell where we assume the density of absorbing material to be constant the outgoing photon number flux is simple $`e^\tau `$ times the incoming one. Hence the number of absorbed photons must be $`(1e^\tau )`$ times the incoming flux. So one can compute the number of photoionizations per second by adding all these $`(1e^\tau )`$ terms for the rays that pass this cell. Now by definition we ensure that the number of photoionizations will always equal the number of photons absorbed. As a consequence this method propagates ionization fronts at always the correct speed independent of resolution. This is a highly desirable feature of any method of multi-dimensional radiative transfer. This control of accuracy comes at high computational cost. In this method the $`1/r^2`$ drop in the photon flux in an optically thin region around the source is captured by the simple fact that many more rays traverse through cells near the source than cells further away. Obviously a large amount of computational time is wasted on computing the flux in such optically thin cell where it would simply be given by $`I(0)/r^2`$. This can be overcome as is discussed below. However, this method nevertheless can be used for a variety of realistic cases. This can be seen from Figure 1 and 2. Both of them employed the ray-tracing of Abel, Norman and Madau (1999) and are here shown as illustration of the practicality of this method.
### 3.2 Ionization Front tracking
Let us quickly side-track to point a simple way of solving a specific problem. If one is interested in the propagation of a R-type ionization front in a static medium it suffices to integrate the jump condition
$$n_{HI}(r)\frac{dR}{dt}=\frac{F(0)}{4\pi R^2}_0^R\alpha (T(r))n_{HII}(r)n_e(r)𝑑r,$$
(4)
where $`R`$, $`\alpha `$, and $`F(0)`$ denote the radius of the I-front, the recombination rate–coefficient, and the ionizing photon number luminosity, respectively. Dividing by $`n_{HI}`$ we can integrate equation (4) explicitly along rays<sup>1</sup><sup>1</sup>1where one uses the raytracing technique of Abel, Norman, and Madau (1999) and find the time at which the ionization fronts arrives at a given cell. Storing the arrival time in a 3D array allows one to investigate the time dependence morphology of ionization front by simply taking iso-surfaces on this array of arrival times. Such data is also interesting to compute how many ionizing photons are used to ionize a given volume in a static case, etc.. To get the full time evolution of the ionization front of one source on a 128<sup>3</sup> numerical grid requires $``$ one minute computation (wall clock) on a workstation.
### 3.3 Computer Graphics
In many fields as e.g. biomedical imaging interactive volume rendering of 3 dimensional data is highly desirable. A lot of effort went into designing fast algorithms that yield optical depths from a light source and to the observers eye. Not any such method will be suitable for application in astrophysics in particular one needs to worry about exact photon (energy) conservation. However, imagine one has a method that gives the optical depth to a source at every point in the computational volume. For the case of pure attenuation one then also knows the photon number flux (photons per second per area) everywhere from
$$\stackrel{}{F}(\stackrel{}{r})=\frac{F_{source}(0)}{|\stackrel{}{r}|^2}e^{\tau (\stackrel{}{r})}\frac{\stackrel{}{r}}{|\stackrel{}{r}|},$$
(5)
where $`\stackrel{}{r}`$ denotes the vector from the source to the point of interest. Now from the obvious ’discontinuity equation’:
$$\stackrel{}{F}(\stackrel{}{r})\dot{n_{HI}}=0$$
(6)
one can ensure the number of photoionizations $`\dot{n_{HI}}`$ to be computed self–consistently independently of resolution. Such a method has been implemented and tested by Abel and Welling (2000) and found to give speed ups in excess of a factor hundred as compared to Abel, Norman and Madau (1999) in the limit of large ionized regions.
### 3.4 Moment Methods
The methods presented above focus and the correct implementation of radiative transfer for point sources. However, ideally we also want to be able to treat regions of diffuse emission as it arises e.g. due to bremsstrahlung and recombination radiation. In Norman, Paschos and Abel (1998) we have outlined a possibility approach to treat point sources and diffuse radiation by means of a variable Eddington tensor formalism. Although we have significantly improved on some ingredients of this method (as e.g. we derived an analytic expression for the Eddington tensors in the pre–overlap stage) we have not succeeded as yet in constructing a stable implementation.
## 4 Concluding Remarks
For the applications to numerical cosmology some of the methods of 3D radiative transfer discussed above will have to be combined. I-front tracking is useful to initialize the environment of new sources. The methods drawn from Computer Graphics can be used to compute accurate boundary conditions for the moment methods that are the most promising in the limit of many sources. A number of interesting problems still will need to be solved before cosmological radiation hydrodynamics can become a standard tool for the study of the formation and evolution of structure in the universe. However, the existing techniques should be employed for the study of interstellar problems in which only few sources are of interest. Planetary Nebulae and HII regions are ideal candidates for such three-dimensional radiation magneto–hydrodynamic modeling.
###### Acknowledgements.
I am greatful to my collaborators Pascal Paschos, Mike Norman, Piero Madau, Aaron Sokasian, Lars Hernquist, and Joel Welling for all the fun we are having in devising these new approaches and learning the physics. Part of this work was supported by NASA ATP grants NAG5-4236 and NAG5-3923.
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# Shimura Curve Computations
## 1 Introduction
### 1.1 Why and how to compute with Shimura curves
The classical modular curves, associated to congruence subgroups of $`\mathrm{PSL}_2(𝐐)`$, have long held and repaid the interest of number theorists working theoretically as well as computationally. In the fundamental paper \[S2\] Shimura defined curves associated with other quaternion algebras other over totally real number fields in the same way that the classical curves are associated with the algebra $`M_2(𝐐)`$ of $`2\times 2`$ matrices over $`𝐐`$. These Shimura curves are now recognized as close analogues of the classical modular curves: almost every result involving the classical curves generalizes with some more work to Shimura curves, and indeed Shimura curves figure alongside classical ones in a key step in the recent proof of Fermat’s “last theorem” \[Ri\].
But computational work on Shimura curves lags far behind the extensive effort devoted to the classical modular curves. The 19th century pioneers investigated some arithmetic quotients of the upper half plane which we now recognize as Shimura curves (see for instance \[F1, F2\]) with the same enthusiasm that they applied to the $`\mathrm{PSL}_2(𝐐)`$ curves. But further inroads proved much harder for Shimura curves than for their classical counterparts. The $`\mathrm{PSL}_2(𝐐)`$ curves parametrize elliptic curves with some extra structure; the general elliptic curve has a simple explicit formula which lets one directly write down the first few modular curves and maps between them. (For instance, this is how Tate obtained the equations for the first few curves X$`{}_{1}{}^{}(N)`$ parametrizing elliptic curves with an $`N`$-torsion point; see for instance \[Kn, pp.145–148\].) Shimura showed that curves associated with other quaternion algebras also parametrize geometric objects, but considerably more complicated ones (abelian varieties with quaternionic endomorphisms); even in the first few cases beyond $`M_2(𝐐)`$, explicit formulas for these objects were obtained only recently \[HM\], and using such formulas to get at the Shimura curves seems a most daunting task. Moreover, most modern computations with modular curves (e.g. \[C, E5\]) sidestep the elliptic interpretation and instead rely heavily on $`q`$-expansions, i.e. on the curves’ cusps. But arithmetic subgroups of $`\mathrm{PSL}_2(𝐑)`$ other than those in $`\mathrm{PSL}_2(𝐐)`$ contain no parabolic elements, so their Shimura curves have no cusps, and thus any method that requires $`q`$-expansions must fail.
But while Shimura curves pose harder computational problems than classical modular curves, efficient solutions to these problems promise great benefits. These curves tempt the computational number theorist not just because, like challenging mountainpeaks, “they’re there”, but because of their remarkable properties, direct applications, and potential for suggesting new ideas for theoretical research. Some Shimura curves and natural maps between them provide some of the most interesting examples in the geometry of curves of low genus; for instance each of the five curves of genus $`g[2,14]`$ that attains the Hurwitz bound $`84(g1)`$ on the number of automorphisms of a curve in characteristic zero is a Shimura curve. Shimura curves, like classical and Drinfeld modular curves, reduce to curves over the finite field $`𝐅_{q^2}`$ of $`q^2`$ elements that attain the Drinfeld-Vlăduţ upper bound $`(q1+o(1))g`$ on the number of points of a curve of genus $`g\mathrm{}`$ over that field \[I3\]. Moreover, while all three flavors of modular curves include towers that can be given by explicit formulas and thus used to construct good error-correcting codes \[Go1, Go2, TVZ\], only the Shimura curves, precisely because of their lack of cusps, can give rise to totally unramified towers, which should simplify the computation of the codes; we gave formulas for several such towers in \[E6\]. Finally, the theory of modular curves indicates that CM (complex multiplication) points on Shimura curves, elliptic curves covered by them, and modular forms on them have number-theoretic significance. The ability to efficiently compute such objects should suggest new theoretical results and conjectures concerning the arithmetic of Shimura curves. For instance, the computations of CM points reported in this paper should suggest factorization formulas for the difference between the coordinates of two such points analogous to those of Gross and Zagier \[GZ\] for $`j`$-invariants of elliptic curves, much as the computation of CM values of the Weber modular functions suggested the formulas of \[YZ\]. Also, as in \[GS\], rational CM points on rational Shimura curves with only three elliptic points (i.e. coming from arithmetic triangle groups $`G_{p,q,r}`$) yield identities $`A+B=C`$ in coprime integers $`A,B,C`$ with many repeated factors; we list the factorizations here, though we found no example in which $`A,B,C`$ are perfect $`p,q,r`$-th powers, nor any new near-record ABC ratios. Finally, CM computations on Shimura curves may also make possible new Heegner-point constructions as in \[E4\].
So how do we carry out these computations? In a few cases (listed in \[JL\]), the extensive arithmetic theory of Shimura curves has been used to obtain explicit equations, deducing from the curves’ $`p`$-adic uniformizations Diophantine conditions on the coefficients of their equations stringent enough to determine them uniquely. But we are interested, not only in the equations, but in modular covers and maps between Shimura curves associated to the same quaternion algebra, and in CM points on those curves. The arithmetic methods may be able to provide this information, but so far no such computation seems to have been done. Our approach relies mostly on the uniformization of these curves qua Riemann surfaces by the hyperbolic plane, and uses almost no arithmetic. This approach is not fully satisfactory either; for instance it probably cannot be used in practice to exhibit all natural maps between Shimura curves of low genus. But it will provide equations for at least a hundred or so curves and maps not previously accessible, which include some of the most striking examples and should provide more than enough data to suggest further theoretical and computational work.
When a Shimura curve $`C`$ comes from an arithmetic subgroup of $`\mathrm{PSL}_2(𝐑)`$ contained in a triangle group $`G_{p,q,r}`$, the curve $`/G_{p,q,r}`$ has genus 0, and $`C`$ is a cover of that curve branched only above three points, so may be determined from the ramification data. (We noted in \[E5, p.48\] that this method was available also for classical modular curves comings from subgroups of $`\mathrm{PSL}_2(𝐙)G_{2,3,\mathrm{}}`$, though there better methods are available thanks to the cusp. Subgroups of $`\mathrm{PSL}_2(𝐑)`$ commensurate with<sup>1</sup><sup>1</sup>1 Recall that two subgroups $`H,K`$ of a group $`G`$ are said to be commensurate if $`HK`$ is a subgroup of finite index in both $`H`$ and $`K`$. but not contained in $`G_{p,q,r}`$ may be handled similarly via the common subgroup of finite index.) The identification of $`/G_{p,q,r}`$ with $`𝐏^1`$ is then given by a quotient of hypergeometric functions on $`𝐏^1`$, which for instance lets us compute the $`𝐏^1`$ coordinate of any CM point on $`C`$ as a complex number to high precision and thus recognize it at least putatively as an algebraic number.
Now it is known \[T\] that only nineteen commensurability classes of arithmetic subgroups of $`\mathrm{PSL}_2(𝐑)`$ contain a triangle group. These include some of the most interesting examples — for instance, congruence subgroups of arithmetic triangle groups account for several of the sporadic “arithmetically exceptional functions” (rational functions $`f(X)𝐐(X)`$ which permute $`𝐏^1(𝐅_p)`$ for infinitely many primes $`p`$) of \[Mü\]; but an approach that could only deal with those nineteen classes would be limited indeed. When there are more than three elliptic points, a new difficulty arises: even if $`C=/𝒢`$ still has genus 0, we must first determine the relative locations of the elliptic points, and to locate other CM points we must replace the hypergeometric functions to solutions of more general “Schwarzian differential equations” in the sense of \[I1\]. We do both by in effect using nontrivial elements of the “commensurator” of the group $`G\mathrm{PSL}_2(𝐑)`$, i.e. transformations in $`\mathrm{PSL}_2(𝐑)`$ which do not normalize $`G`$ but conjugate $`G`$ to a group commensurable with $`G`$. Ihara had already used these commensurators in \[I1\] theoretically to prove that both $`C`$ and its Schwarzian equation are defined over a number field, but this method has apparently not been actually used to compute such equations until now.
### 1.2 Overview of the paper
We begin with a review of the necessary definitions and facts on quaternion algebras and Shimura curves, drawn mostly from \[S2\] and \[V\]. We then give extended computational accounts of Shimura curves and their supersingular and rational CM points for the two simplest indefinite quaternion algebras over $`𝐐`$ beyond the classical case of the matrix algebra $`M_2(𝐐)`$, namely the quaternion algebras ramified at $`\{2,3\}`$ and $`\{2,5\}`$. In the final section we more briefly treat some other examples which illustrate features of our methods that do not arise in the $`\{2,3\}`$ and $`\{2,5\}`$ cases, and conclude with some open questions suggested by our computations that may point the way to further computational investigation on these curves.
### 1.3 Acknowledgements
Many thanks to B.H. Gross for introducing me to Shimura curves and for many enlightening conversations and clarifications on this fascinating topic. Thanks also to Serre for a beautiful course that introduced me to three-point covers of $`𝐏^1`$ among other things (\[Se\], see also \[Mat\]); to Ihara for alerting me to his work \[I1, I2\] on supersingular points on Shimura curves and their relation with the curves’ uniformization by the upper half-plane; and to C. McMullen for discussions of the uniformization of quotients of $``$ by general co-compact discrete subgroups of $`\mathrm{PSL}_2(𝐑)`$. A. Adler provided several references to the 19th-century literature, and C. Doran informed me of \[HM\]. Finally, I thank B. Poonen for reading and commenting on a draft of this paper, leading to considerable improvements of exposition in several places.
The numerical and symbolic computations reported here were carried out using the gp/pari and macsyma packages, except for (70), for which I thank Peter Müller as noted there.
This work was made possible in part by funding from the David and Lucile Packard Foundation.
## 2 Review of quaternion algebras over $`𝐐`$ and their Shimura curves
### 2.1 Quaternion algebras over $`𝐐`$; the arithmetic groups $`\mathrm{\Gamma }(1)`$ and $`\mathrm{\Gamma }^{}(1)`$
Let $`K`$ be a field of characteristic zero; for our purposes $`K`$ will always be a number field or, rarely, its localization, and usually the number field will be $`𝐐`$. A quaternion algebra over $`K`$ is a simple associative algebra $`𝖠`$ with unit, containing $`K`$, such that $`K`$ is the center of $`𝖠`$ and $`dim_K𝖠=4`$. Such an algebra has a conjugation $`a\overline{a}`$, which is a $`K`$-linear anti-involution (i.e. $`\overline{\overline{a}}=a`$ and $`\overline{a_1a_2}=\overline{a}_2\overline{a}_1`$ hold identically in $`𝖠`$) such that $`a=\overline{a}aK`$. The trace and norm are the additive and multiplicative maps from $`𝖠`$ to $`K`$ defined by
$$\mathrm{tr}(a)=a+\overline{a},\mathrm{N}(a)=a\overline{a}=\overline{a}a;$$
(1)
every $`a𝖠`$ satisfies its characteristic equation
$$a^2(\mathrm{tr}(a))a+\mathrm{N}(a)=0.$$
(2)
The most familiar example of a quaternion algebra is $`M_2(K)`$, the algebra of $`2\times 2`$ matrices over $`K`$, and if $`K`$ is algebraically closed then $`M_2(K)`$ is the only quaternion algebra over $`K`$ up to isomorphism. The other well-known example is the algebra of Hamilton quaternions over $`𝐑`$. In $`M_2(K)`$ the trace is the usual trace of a square matrix, so the conjugate of $`aM_2(K)`$ is $`\mathrm{tr}(a)I_{2\times 2}a`$, and the norm is just the determinant. Any quaternion algebra with zero divisors is isomorphic with $`M_2(K)`$. An equivalent criterion is that the algebra contain a nonzero element whose norm and trace both vanish. Now the trace-zero elements constitute a $`K`$-subspace of $`𝖠`$ of dimension 3, on which the norm is a homogeneous quadric; so the criterion states that $`𝖠M_2(K)`$ if and only if that quadric has nonzero $`K`$-rational points. The Hamilton quaternions have basis $`1,i,j,k`$ satisfying the familiar relations
$$i^2=j^2=k^2=1,ij=ji=k,jk=kj=i,ki=ik=j;$$
(3)
the conjugates of $`1,i,j,k`$ are $`1,i,j,k`$, so a Hamilton quaternion $`\alpha _1+\alpha _2i+\alpha _3j+\alpha _4k`$ has trace $`2\alpha _1`$ and norm $`\alpha _1^2+\alpha _2^2+\alpha _3^2+\alpha _4^2`$. Thus the Hamilton quaternions over $`K`$ are isomorphic with $`M_2(K)`$ if and only if $`1`$ is a sum of two squares in $`K`$.
In fact it is known that if $`K=𝐑`$ then every quaternion algebra over $`K`$ is isomorphic with either $`M_2(𝐑)`$ or the Hamilton quaternions. In general if $`K`$ is any local field of characteristic zero then there is up to isomorphism exactly one quaternion algebra over $`K`$ other than $`M_2(K)`$ — with the exception of the field of complex numbers, which being algebraically closed admits no quaternion algebras other than $`M_2(𝐂)`$. If $`𝖠`$ is a quaternion algebra over a number field $`K`$ then a finite or infinite place $`v`$ of $`K`$ is said to be ramified in $`𝖠`$ if $`𝖠K_v`$ is not isomorphic with $`M_2(K_v)`$. There can only be a finite number of ramified places, because a nondegenerate quadric over $`K`$ has nontrivial local zeros at all but finitely many places of $`K`$. A less trivial result (the case $`K=𝐐`$ is equivalent to Quadratic Reciprocity) is that the number of ramified places is always even, and to each finite set of places $`\mathrm{\Sigma }`$ of even cardinality containing no complex places there corresponds a unique (again up to isomorphism) quaternion algebra over $`K`$ ramified at those places and no others. In particular an everywhere unramified quaternion algebra over $`K`$ must be isomorphic with $`M_2(K)`$.
An order in a quaternion algebra over a number field (or a non-Archimedean local field) $`K`$ is a subring containing the ring $`O_K`$ of $`K`$-integers and having rank 4 over $`O_K`$. For instance $`M_2(O_K)`$ and $`O_K[i,j]`$ are orders in the matrix and quaternion algebras over $`K`$. Any order is contained in at least one maximal order, that is, in an order not properly contained in any other. Examples of maximal orders are $`M_2(O_K)M_2(K)`$ and the Hurwitz order $`𝐙[1,i,j,(1+i+j+k)/2]`$ in the Hamilton quaternions over $`𝐐`$. It is known that if $`K`$ has at least one Archimedean place at which $`𝖠`$ is not isomorphic with the Hamilton quaternions then all maximal orders are conjugate in $`𝖠`$.
Now let<sup>2</sup><sup>2</sup>2 Most of our examples, including the two that will occupy us in the next two sections, involve quaternion algebras over $`𝐐`$. In \[S2\] Shimura associated modular curves to a quaternion algebra over any totally real number field $`K`$ for which the algebra is ramified at all but one of the infinite places of $`K`$. Since the special case $`K=𝐐`$ accounts for most of our computations, and is somewhat easier to describe, we limit our discussion to quaternion algebras over $`𝐐`$ from here until section 5.3. At that point we briefly describe the situation for arbitrary $`K`$ before working out a couple of examples with $`[K:𝐐]>`$1. $`K=𝐐`$. A quaternion algebra $`𝖠/𝐐`$ is called definite or indefinite according as $`𝖠𝐑`$ is isomorphic with the Hamilton quaternions or $`M_2(𝐑)`$, i.e. according as the infinite place is ramified or unramified in $`𝖠`$. \[These names allude to the norm form on the trace-zero subspace of $`𝖠`$, which is definite in the former case, indefinite in the latter.\] We shall be concerned only with the indefinite case. Then $`\mathrm{\Sigma }`$ consists of an even number of finite primes. Fix such a $`\mathrm{\Sigma }`$ and the corresponding quaternion algebra $`𝖠`$. Let $`𝒪`$ be a maximal order in $`𝖠`$; since $`𝖠`$ is indefinite, all its maximal orders are conjugate, so choosing a different maximal order would not materially affect the constructions in the sequel.
Let $`𝒪_1^{}`$ be the group of units of norm 1 in $`𝒪`$. We then define the following arithmetic subgroups of $`𝖠^{}/𝐐^{}`$:
$`\mathrm{\Gamma }(1)`$ $`:=`$ $`𝒪_1^{}/\{\pm 1\},`$ (4)
$`\mathrm{\Gamma }^{}(1)`$ $`:=`$ $`\{[a]𝖠^{}/𝐐^{}:a𝒪=𝒪a,\mathrm{N}(a)>0\}.`$ (5)
\[In other words $`\mathrm{\Gamma }^{}(1)`$ is the normalizer of $`\mathrm{\Gamma }(1)`$ in the positive-norm subgroup of $`𝖠^{}/𝐐^{}`$. Takeuchi \[T\] calls these groups $`\mathrm{\Gamma }^{(1)}(𝖠,𝒪_1)`$ and $`\mathrm{\Gamma }^{()}(𝖠,𝒪_1)`$; we use $`\mathrm{\Gamma }(1)`$ to emphasize the analogy with the classical case of $`\mathrm{PSL}_2(𝐙)`$, which makes $`\mathrm{\Gamma }^{}(1)`$ a natural adaptation of Takeuchi’s notation. Vignéras \[V, p. 121ff.\] calls the same groups $`\mathrm{\Gamma }`$ and $`G`$, citing \[Mi\] for the structure of their quotient.\] As noted, $`\mathrm{\Gamma }(1)`$ is a normal subgroup of $`\mathrm{\Gamma }^{}(1)`$. In fact $`\mathrm{\Gamma }^{}(1)`$ consists of the classes mod $`𝐐^{}`$ of elements of $`𝒪`$ whose norm is $`_{p\mathrm{\Sigma }^{}}p`$ for some (possibly empty) subset $`\mathrm{\Sigma }^{}\mathrm{\Sigma }`$, and $`\mathrm{\Gamma }^{}(1)/\mathrm{\Gamma }(1)`$ is an elementary abelian 2-group with $`\mathrm{\#}\mathrm{\Sigma }`$ generators.
### 2.2 The Shimura modular curves $`𝒳(1)`$ and $`𝒳^{}(1)`$
The group $`\mathrm{\Gamma }(1)`$, and thus any other group commensurable with it such as $`\mathrm{\Gamma }^{}(1)`$, is a discrete subgroup of $`(A𝐑)_+^{}/𝐑^{}`$ (the subscript “$`+`$” indicating positive norm), with compact quotient unless $`\mathrm{\Sigma }=\mathrm{}`$, and of finite covolume even in that case. Since $`A𝐑M_2(𝐑)`$, the group $`(A𝐑)_+^{}/𝐑^{}`$ is isomorphic with $`\mathrm{PSL}_2(𝐑)`$ and thus with $`\mathrm{Aut}()`$, the group of automorphisms of the hyperbolic upper half plane
$$:=\{z𝐂:\mathrm{Im}(z)>0\}.$$
(6)
Explicitly, a unimodular matrix $`\pm (\genfrac{}{}{0pt}{}{ab}{cd})`$ acts on $``$ via the fractional linear transformation $`z(az+b)/(cz+d)`$. We may define the Shimura curves $`𝒳(1)`$ and $`𝒳^{}(1)`$ qua compact Riemann surfaces by
$$𝒳(1):=/\mathrm{\Gamma }(1),𝒳^{}(1):=/\mathrm{\Gamma }^{}(1).$$
(7)
\[More precisely, the Riemann surfaces are given by (7) unless $`\mathrm{\Sigma }=\mathrm{}`$, in which case the quotient only becomes compact upon adjoining a cusp.\] The hyperbolic area of these quotients of $``$ is given by the special case $`k=𝐐`$ of a formula of Shimizu \[S1, Appendix\], quoted in \[T, p.207\]. Using the normalization $`\pi ^1𝑑x𝑑y/y^2`$ for the hyperbolic area (with $`z=x+iy`$; this normalization gives an ideal triangle unit area), that formula is
$$\mathrm{Area}(𝒳(1))=\frac{1}{6}\underset{p\mathrm{\Sigma }}{}(p1),$$
(8)
from which
$$\mathrm{Area}(𝒳^{}(1))=\frac{1}{[\mathrm{\Gamma }^{}(1):\mathrm{\Gamma }(1)]}\mathrm{Area}(𝒳(1))=\frac{1}{6}\underset{p\mathrm{\Sigma }}{}\frac{p1}{2}.$$
(9)
It is known (see for instance Ch.IV:2,3 of \[V\] for the following facts) that, for any discrete subgroup $`\mathrm{\Gamma }\mathrm{PSL}_2(𝐑)`$ of finite covolume, the genus of $`/\mathrm{\Gamma }`$ is determined by its area together with with information on elements of finite order in $`\mathrm{\Gamma }`$. All finite subgroups of $`\mathrm{\Gamma }`$ are cyclic, and there are finitely many such subgroups up to conjugation in $`\mathrm{\Gamma }`$. There are finitely many points $`P_j`$ of $`/\mathrm{\Gamma }`$ with nontrivial stabilizer, and the stabilizers are the maximal nontrivial finite subgroups of $`\mathrm{\Gamma }`$ modulo conjugation in $`\mathrm{\Gamma }`$. If the order of the stabilizer of $`P_j`$ is $`e_j`$ then $`P_j`$ is said to be an “elliptic point of order $`e_j`$”. Then if $`/\mathrm{\Gamma }`$ is compact then its genus $`g=g(/\mathrm{\Gamma })`$ is given by
$$2g2=\mathrm{Area}(/\mathrm{\Gamma })\underset{j}{}(1\frac{1}{e_j}).$$
(10)
Moreover $`\mathrm{\Gamma }`$ has a presentation
$$\mathrm{\Gamma }=\alpha _1,\mathrm{},\alpha _g,\beta _1,\mathrm{},\beta _g,s_j|s_j^{e_j}=1,\underset{j}{}s_j\underset{i=1}{\overset{g}{}}[\alpha _i,\beta _i]=1,$$
(11)
in which $`s_j`$ generates the stabilizer of a preimage of $`P_j`$ in $``$ and rotates a neighborhood of that preimage by an angle $`2\pi /e_j`$ (i.e. has derivative $`e^{2\pi i/e_j}`$ at its fixed point), and $`[\alpha ,\beta ]`$ is the commutator $`\alpha \beta \alpha ^1\beta ^1`$. \[This group is sometimes called $`(g;e_1,\mathrm{},e_g)`$.\] If $`/\mathrm{\Gamma }`$ is not compact then we must subtract the number of cusps from the right-hand side of (10) and include a generator $`s_j`$ of $`\mathrm{\Gamma }`$ of infinite order for each cusp, namely a generator of the infinite cyclic stabilizer of the cusp. This generator is a “parabolic element” of $`\mathrm{PSL}_2(𝐑)`$, i.e. a fractional linear transformation with a single fixed point; there are two conjugacy classes of such elements in $`\mathrm{PSL}_2(𝐑)`$, and $`s_j`$ will be in the class of $`zz+1`$. We assign $`e_j=\mathrm{}`$ to a cusp. For both finite and infinite $`e_j`$, the trace and determinant of $`s_j`$ are related by
$$\mathrm{Tr}^2(s_j)=4\mathrm{cos}^2\frac{\pi }{e_j}det(s_j).$$
(12)
Since we are working in quaternion algebras over $`𝐐`$, this means that $`e_j\{2,3,4,6,\mathrm{}\}`$, and only $`2,3,\mathrm{}`$ are possible if $`\mathrm{\Gamma }\mathrm{\Gamma }(1)`$. Moreover $`e_j=\mathrm{}`$ occurs only in the classical case $`\mathrm{\Sigma }=\mathrm{}`$.
We shall need to numerically compute for several such $`\mathrm{\Gamma }`$ the identification of $`/\mathrm{\Gamma }`$ with an algebraic curve $`X/𝐂`$, i.e. to compute the coordinates on $`X`$ of a point corresponding to (the $`\mathrm{\Gamma }`$-orbit of) a given $`z`$, or inversely to obtain $`z`$ corresponding to a point with given coordinates. In fact the two directions are essentially equivalent, because if we can efficiently compute an isomorphism between two Riemann surfaces then we can compute its inverse almost as easily. For classical modular curves one usually uses $`q`$-expansions to go from $`z`$ to rational coordinates; but this method is not available for our groups $`\mathrm{\Gamma }`$, which have no parabolic ($`e_j=\mathrm{}`$) generator. We can, however, still go in the opposite direction, computing the map from $`X`$ to $`/\mathrm{\Gamma }`$ by solving differential equations on $`X`$. The key is that while the function $`z`$ on $`X`$ is not well defined due to the $`\mathrm{\Gamma }`$ ambiguity, its Schwarzian derivative is. In local coordinates the Schwarzian derivative of a nonconstant function $`z=z(\zeta )`$ is the meromorphic function defined by
$$S_\zeta (z):=4z^1z_{}^{}{}_{}{}^{1/2}\frac{d^2}{d\zeta ^2}\frac{z}{z_{}^{}{}_{}{}^{1/2}}=\frac{2z^{}z^{\prime \prime \prime }3z_{}^{\prime \prime }{}_{}{}^{2}}{z_{}^{}{}_{}{}^{2}}.$$
(13)
This vanishes if and only if $`z`$ is a fractional linear transformation of $`\zeta `$. Moreover it satisfies a nice “chain rule”: if $`\zeta `$ is in turn a function of $`\eta `$ then
$$S_\eta (z)=\left(\frac{d\zeta }{d\eta }\right)^2S_\zeta (z)+S_\eta (\zeta ).$$
(14)
Thus if we choose a coordinate $`\zeta `$ on $`X`$ then $`S_\zeta (z)`$ is the same for each lift of $`z`$ from $`/\mathrm{\Gamma }`$ to $``$, and thus gives a well-defined function on the complement in $`X`$ of the elliptic points; changing the coordinate from $`\zeta `$ to $`\eta `$ multiples this function by $`(d\zeta /d\eta )^2`$ and adds a term $`S_\eta (\zeta )`$ that vanishes if $`\zeta `$ is a fractional linear transformation of $`\eta `$. In particular if $`X`$ has genus 0 and we choose only rational coordinates (i.e. $`\eta ,\zeta `$ are rational functions of degree 1) then these terms $`S_\eta (\zeta )`$ always vanish and $`S_\zeta (z)d\zeta ^2`$ is a well-defined quadratic differential $`\sigma `$ on $`X`$. Near an elliptic point $`\zeta _0`$ of index $`e_j`$, the function $`z`$ has a branch point such that $`(zz_0)/(z\overline{z}_0)`$ is $`(\zeta \zeta _0)^{1/e_j}`$ times an analytic function; for such $`z`$ the Schwarzian derivative is still well-defined in a neighborhood of $`\zeta _0`$ but has a double pole there with leading term $`(1e_j^2)/(\zeta \zeta _0)^2`$ \[or $`(1e_j^2)/\zeta ^2`$ if $`\zeta _0=\mathrm{}`$ — note that this too has a double pole when multiplied by $`d\zeta ^2`$\]. So $`\sigma =S_\zeta (z)d\zeta ^2`$ is a rational quadratic differential on $`X`$, regular except for double poles of known residue at the elliptic points, and independent of the choice of rational coordinate when $`X`$ has genus 0. Knowing $`\sigma `$ we may recover $`z`$ from the differential equation
$$S_\zeta (z)=\sigma /d\zeta ^2,$$
(15)
which determines $`z`$ up to a fractional linear transformation over $`𝐂`$, and can then remove the ambiguity if we know at least three values of $`z`$ (e.g. at elliptic points, which are fixed points of known elements of $`\mathrm{\Gamma }`$).
Because $`S_\zeta (z)`$ is invariant under fractional linear transformations of $`z`$, the third-order nonlinear differential equation (15) can be linearized as follows (see e.g. \[I1, §1–5\]). Let $`(f_1,f_2)`$ be a basis for the solutions of the linear second-order equation
$$f^{\prime \prime }=af^{}+bf$$
(16)
for some functions $`a(\zeta ),b(\zeta )`$. Then $`z:=f_1/f_2`$ is determined up to fractional linear transformation, whence $`S_\zeta (z)`$ depends only on $`a,b`$ and not the choice of basis. In fact we find, using either of the equivalent definitions in (13), that
$$S_\zeta (f_1/f_2)=2\frac{da}{d\zeta }a^24b.$$
(17)
Thus if $`a`$ is any rational function and $`b=\sigma /4d\zeta ^2+a^{}/2a^2/4`$ then the solutions of (15), and thus a map from $`X`$ to $`/\mathrm{\Gamma }`$, are ratios of linearly independent pairs of solutions of (16). In the terminology of \[I1\], (16) is then a Schwarzian equation for $`/\mathrm{\Gamma }`$. We shall always choose $`a`$ so that $`ad\zeta `$ has at most simple poles at the elliptic points and no other poles; the Schwarzian equation then has regular singularities at the elliptic points and no other singularities. The most familiar example is the case that $`\mathrm{\Gamma }`$ is a triangle group, i.e. $`X`$ has genus 0 and three elliptic points (if $`g=0`$ there must be at least three elliptic points by (10)). In that case $`\sigma `$ is completely determined by its poles and residues: if two different $`\sigma `$’s were possible, their difference would be a nonzero quadratic differential on $`𝐏^1`$ with at most three simple poles, which is impossible. If we choose the coordinate on $`X`$ that puts the elliptic points at $`0,1,\mathrm{}`$, and require that $`a`$ be chosen of the form $`a=C_0/\zeta +C_1/(\zeta 1)`$ so that $`b`$ has only simple poles at $`0,1`$, then there are four choices for $`(C_0,C_1)`$, each giving rise to a hypergeometric equation upon multiplying (16) by $`\zeta (1\zeta )`$:
$$\zeta (1\zeta )f^{\prime \prime }=[(\alpha +\beta +1)\zeta \gamma ]f^{}+\alpha \beta f.$$
(18)
Here $`\alpha ,\beta ,\gamma `$ are related to the indices $`e_1,e_2,e_3`$ at $`\zeta =0,1,\mathrm{}`$ by
$$\frac{1}{e_1}=\pm (1\gamma ),\frac{1}{e_2}=\pm (\gamma \alpha \beta ),\frac{1}{e_3}=\pm (\alpha \beta );$$
(19)
then $`F(\alpha ,\beta ;\gamma ;\zeta )`$ and $`(1\zeta )^\gamma F(\alpha \gamma +1,\beta \gamma +1;2\gamma ;\zeta )`$ constitute a basis for the solutions of (16), where $`F={}_{2}{}^{}F_{1}^{}`$ is the hypergeometric function defined for $`|\zeta |<1`$ by
$$F(\alpha ,\beta ;\gamma ;\zeta ):=\underset{n=0}{\overset{\mathrm{}}{}}\left[\underset{k=0}{\overset{n1}{}}\frac{(\alpha +k)(\beta +k)}{(\gamma +k)}\right]\frac{\zeta ^n}{n!},$$
(20)
and by similar power series in neighborhoods of $`\zeta =1`$ and $`\zeta =\mathrm{}`$ (see for instance \[GR, 9.10 and 9.15\]). In general, knowing $`\sigma `$ we may construct and solve a Schwarzian equation in power series, albeit series less familiar than $`{}_{2}{}^{}F_{1}^{}`$, and numerically compute the map $`X/\mathrm{\Gamma }`$ as the quotient of two solutions. But once $`\mathrm{\Gamma }`$ is not a triangle group — that is, when $`X`$ has more than three elliptic points or positive genus — the elliptic points and their orders no longer determine $`\sigma `$ but only restrict it to an affine space of finite but positive dimension. In general it is a refractory problem to find the “accessory parameters” that tell where $`\sigma `$ lies in that space. If $`\mathrm{\Gamma }`$ is commensurable with a triangle group $`\mathrm{\Gamma }^{}`$ then we obtain $`\sigma `$ from the quadratic differential on $`/\mathrm{\Gamma }^{}`$ via the correspondence between that curve and $`/\mathrm{\Gamma }`$; but this only applies to Shimura curves associated with the nineteen quaternion algebras listed by Takeuchi in \[T\], including only two over $`𝐐`$, the matrix algebra and the algebra ramified at $`\{2,3\}`$. One of the advances in the present paper is the computation of $`\sigma `$ for some arithmetic groups not commensurable with any triangle group.
We now return to the Shimura curves $`𝒳(1)`$, $`𝒳^{}(1)`$ obtained from arithmetic groups $`\mathrm{\Gamma }=\mathrm{\Gamma }(1),\mathrm{\Gamma }^{}(1)`$. These curves also have a modular interpretation that gives them the structure of algebraic curves over $`𝐐`$. To begin with, $`𝒳(1)`$ is the modular curve for principally polarized abelian surfaces (ppas) $`A`$ with an embedding $`O\mathrm{End}(A)`$. (In the classical case $`𝒪=M_2(𝐙)`$, corresponding to $`\mathrm{\Sigma }=\mathrm{}`$, such an abelian surface is simply the square of an elliptic curve and we recover the familiar picture of modular curves parametrizing elliptic ones, but for nonempty $`\mathrm{\Sigma }`$ the surfaces $`A`$ are simple except for those associated to CM points on $`𝒳(1)`$; we shall say more about CM points later.) The periods of these surfaces satisfy a linear second-order differential equation which is a Schwarzian equation for $`/\mathrm{\Gamma }(1)`$, usually called a “Picard-Fuchs equation” in this context. \[This generalizes the expression for the periods of elliptic curves (a.k.a. “complete elliptic integrals”) as $`{}_{2}{}^{}F_{1}^{}`$ values, for which see e.g. \[GR, 8.113 1.\].\] The group $`\mathrm{\Gamma }^{}(1)/\mathrm{\Gamma }(1)`$ acts on $`𝒳(1)`$ with quotient curve $`𝒳^{}(1)`$. For each $`p\mathrm{\Sigma }`$ there is then an involution $`𝗐_p\mathrm{\Gamma }^{}(1)/\mathrm{\Gamma }(1)`$ associated to the class in $`\mathrm{\Gamma }^{}(1)/\mathrm{\Gamma }(1)`$ of elements of $`𝒪`$ of norm $`p`$, and these involutions commute with each other. (We chose the notation $`𝗐_p`$ to suggest an analogy with the Atkin-Lehner involutions $`w_l`$, which as we shall see have a more direct counterpart in our setting when $`l\mathrm{\Sigma }`$.)
In terms of abelian surfaces these involutions $`𝗐_p`$ of $`𝒳(1)`$ may be explained as follows. Let $`I_p𝒪`$ consist of the elements whose norm is divisible by $`p`$. Then $`I_p`$ is a two-sided prime ideal of $`𝒪`$, with $`𝒪/I_p𝐅_{p^2}`$ and $`I_p^2=p𝒪`$. Given an action of $`𝒪`$ on a ppas $`A`$, the kernel of $`I_p`$ is a subgroup of $`A`$ of size $`p^2`$ isotropic under the Weil pairing, so the quotient surface $`A^{}:=A/\mathrm{ker}I_p`$ is itself principally polarized. Moreover, since $`I_p`$ is a two-sided ideal, $`A^{}`$ inherits an action of $`𝒪`$. Thus if $`A`$ corresponds to some point $`P𝒳(1)`$ then $`A^{}`$ corresponds to a point $`P^{}𝒳(1)`$ determined algebraically by $`P`$; that is, we have an algebraic map $`𝗐_p:PP^{}`$ from $`𝒳(1)`$ to itself. Applying this construction to $`A^{}`$ yields $`A/\mathrm{ker}I_p^2=A/\mathrm{ker}p𝒪=A/\mathrm{ker}pA`$; thus $`𝗐_p(P^{})=P`$ and $`𝗐_p`$ is indeed an involution. The quotient curve $`𝒳^{}(1)`$ then parametrizes surfaces $`A`$ up to the identification of $`A`$ with $`A/\mathrm{ker}I`$ where $`I=_{p\mathrm{\Sigma }^{}}I_p=_{p\mathrm{\Sigma }^{}}I_p`$ for some $`\mathrm{\Sigma }^{}\mathrm{\Sigma }`$.
Since $`𝒳(1)`$, $`𝒳^{}(1)`$ have the structure of algebraic curves over $`𝐐`$, they can be regarded as curves over $`𝐑`$. Now a real structure on any Riemann surface is equivalent to an anti-holomorphic involution of the surface. For surfaces $`/\mathrm{\Gamma }`$ uniformized by the upper half-plane, we can give such an involution by choosing a group $`(\mathrm{\Gamma }:2)\mathrm{PGL}_2(𝐑)`$ containing $`\mathrm{\Gamma }`$ with index 2 such that $`(\mathrm{\Gamma }:2)\mathrm{PSL}_2(𝐑)`$. An element $`(\genfrac{}{}{0pt}{}{ab}{cd})𝐑^{}`$ of $`\mathrm{PGL}_2(𝐑)\mathrm{PSL}_2(𝐑)`$ (i.e. with $`adbc<0`$) acts on $``$ anti-holomorphically $`z(a\overline{z}+b)/(c\overline{z}+d)`$. Such a fractional conjugate-linear transformation has fixed points on $``$ if and only if $`a+d=0`$, in which case it is an involution and its fixed points constitute a hyperbolic line. Thus $`/\mathrm{\Gamma }`$, considered as a curve over $`𝐑`$ using $`\mathrm{\Gamma }:2`$, has real points if and only if $`(\mathrm{\Gamma }:2)\mathrm{\Gamma }`$ contains an involution of $``$. The real structures on $`𝒳(1)`$, $`𝒳^{}(1)`$ are defined by
$`(\mathrm{\Gamma }(1):2)`$ $`:=`$ $`𝒪^{}/\{\pm 1\},`$ (21)
$`(\mathrm{\Gamma }^{}(1):2)`$ $`:=`$ $`\{[a]𝖠^{}/𝐐^{}:a𝒪=𝒪a\}.`$ (22)
That is, compared with (4,5) we drop the condition that the norm be positive. If $`\mathrm{\Sigma }\mathrm{}`$ then $`𝒳(1)`$ has no real points, because if $`\mathrm{\Gamma }(1):2`$ contained an involution $`\pm a`$ then the characteristic equation of $`a`$ would be $`a^21=0`$ and $`𝖠`$ would contain the zero divisors $`a\pm 1`$. This is a special case of the result of \[S3\]. But $`𝒳^{}(1)`$ may have real points. For instance, we shall see that if $`\mathrm{\Sigma }=\{2,3\}`$ then $`\mathrm{\Gamma }^{}(1)`$ is isomorphic with the triangle group $`G_{2,4,6}`$. For general $`p,q,r`$ with<sup>3</sup><sup>3</sup>3 If $`1/p+1/q+1/r`$ equals or exceeds $`1`$, an analogous situation occurs with $``$ replaced by the complex plane or Riemann sphere. $`1/p+1/q+1/r<1`$ we can (and, if $`p,q,r`$ are distinct, can only) choose $`G_{p,q,r}:2`$ so that the real locus of $`/G_{p,q,r}`$ consists of three hyperbolic lines joining the three elliptic points in pairs, forming a hyperbolic triangle, with $`G_{p,q,r}:2`$ generated by hyperbolic reflections in the triangle’s sides; it is this triangle to which the term “triangle group” alludes.
### 2.3 The Shimura modular curves $`𝒳(N)`$ and $`𝒳^{}(N)`$ (with $`N`$ coprime to $`\mathrm{\Sigma }`$); the curves $`𝒳_0(N)`$ and $`𝒳_0^{}(N)`$ and their involution $`w_N`$
Now let $`l`$ be a prime not ramified in $`𝖠`$. Then $`𝖠𝐐_l`$ and $`𝒪𝐙_l`$ are isomorphic with $`M_2(𝐐_l)`$ and $`M_2(𝐙_l)`$ respectively. Thus $`(𝒪𝐐_l)_1^{}/\{\pm 1\}\mathrm{PSL}_2(𝐙_l)`$, with the subscript $`1`$ indicating the norm-1 subgroup as in (4). We can thus define congruence subgroups $`\mathrm{\Gamma }(l)`$, $`\mathrm{\Gamma }_1(l)`$, $`\mathrm{\Gamma }_0(l)`$ of $`\mathrm{\Gamma }(1)`$ just as in the classical case in which $`\mathrm{\Sigma }=\mathrm{}`$ and $`\mathrm{\Gamma }(1)=\mathrm{PSL}_2(𝐙)`$. For instance, $`\mathrm{\Gamma }(l)`$ is the normal subgroup
$$\{\pm a𝒪_+^{}/\{\pm 1\}:a1modl\}$$
(23)
of $`\mathrm{\Gamma }(1)`$, with $`\mathrm{\Gamma }(1)/\mathrm{\Gamma }(l)\mathrm{PSL}_2(𝐅_l)`$; once we choose an identification of the quotient group $`\mathrm{\Gamma }(1)/\mathrm{\Gamma }(l)`$ with $`\mathrm{PSL}_2(𝐅_l)`$ we may define $`\mathrm{\Gamma }_0(l)`$ as the preimage in $`\mathrm{\Gamma }(1)`$ of the upper triangular subgroup of $`\mathrm{PSL}_2(𝐅_l)`$. Likewise we have subgroups $`\mathrm{\Gamma }(l^r)`$, $`\mathrm{\Gamma }_0(l^r)`$ etc., and even $`\mathrm{\Gamma }(N)`$, $`\mathrm{\Gamma }_0(N)`$ for a positive integer $`N`$ not divisible by any of the primes of $`\mathrm{\Sigma }`$.
The quotients of $``$ by these subgroups of $`\mathrm{\Gamma }(1)`$ are then modular curves covering $`𝒳(1)`$, which we denote by $`𝒳(l)`$, $`𝒳_0(l)`$, etc. They parametrize ppas’s $`A`$ with an $`𝒪`$-action and extra structure: in the case of $`𝒳(N)`$, a choice of basis for the $`N`$-torsion points $`A[N]`$; in the case of $`𝒳_0(N)`$, a subgroup $`GA[N]`$ isomorphic with $`(𝐙/N)^2`$ and isotropic under the Weil pairing. In the latter case the surface $`A^{}=A/G`$ is itself principally polarized and inherits an action of $`𝒪`$ from $`A`$, and the image of $`A[N]`$ in $`A^{}`$ is again a subgroup $`G^{}(𝐙/N)^2`$ isotropic under the Weil pairing. Thus if we start from some point $`P`$ on $`𝒳_0(N)`$ and associate to it a pair $`(A,G)`$ we obtain a new pair $`(A^{},G^{})`$ of the same kind and a new point $`P^{}𝒳_0(N)`$ determined algebraically by $`P`$. Thus we have an algebraic map $`w_N:PP^{}`$ from $`𝒳_0(N)`$ to itself. As in the classical case — in which it is easy to see that the construction of $`A^{},G^{}`$ from $`A,G`$ amounts to (the square of) the familiar picture of cyclic subgroups and dual isogenies — this $`w_N`$ is an involution of $`𝒳_0(N)`$ that comes from a trace-zero element of $`𝖠`$ of norm $`N`$ whose image in $`𝖠^{}/𝐐^{}`$ is an involution normalizing $`\mathrm{\Gamma }_0(N)`$.
By abuse of terminology we shall say that a pair of points $`P,P^{}`$ on $`𝒳(1)`$ are “cyclically $`N`$-isogenous”<sup>4</sup><sup>4</sup>4 This qualifier “cyclically” is needed to exclude cases such as the multiplication-by-$`m`$ map, which as in the case of elliptic curves would count as an “$`m^2`$-isogeny” but not a cyclic one. if they correspond to ppas’s $`A,A^{}`$ with $`A^{}=A/G`$ as above, and call the quotient map $`AA/GA^{}`$ a “cyclic $`N`$-isogeny”. If we regard $`P,P^{}`$ as $`\mathrm{\Gamma }(1)`$-orbits in $``$ then they are cyclically $`N`$-isogenous iff a point in the first orbit is taken to a point in the second by some $`a𝒪`$ of norm $`N`$ such that $`ama^{}`$ for any $`a^{}𝒪`$ and $`m>1`$; since in that case $`\overline{a}`$ also satisfies this condition and acts on $``$ as the inverse of $`a`$, this relation on $`P,P^{}`$ is symmetric. Then $`𝒳_0(N)`$ parametrizes pairs of $`N`$-isogenous points on $`𝒳(1)`$, and $`w_N`$ exchanges the points in such a pair. The involutions $`𝗐_p`$ on $`𝒳(1)`$ lift to the curves $`𝒳(N)`$, $`𝒳_0(N)`$, etc., and commute with $`w_N`$ on $`𝒳_0(N)`$.
The larger group $`\mathrm{\Gamma }^{}(1)`$ likewise has congruence groups such as $`\mathrm{\Gamma }^{}(N)`$, $`\mathrm{\Gamma }_0^{}(N)`$, etc., which give rise to modular curves covering $`𝒳^{}(1)`$ called $`𝒳^{}(N)`$, $`𝒳_0^{}(N)`$, etc. The involution $`w_N`$ on $`𝒳_0(N)`$ descends to an involution on $`𝒳_0^{}(N)`$ which we shall also call $`w_N`$. We extend our abuse of terminology by saying that two points on $`𝒳^{}(1)`$ are “cyclically $`N`$-isogenous” if they lie under two $`N`$-isogenous points of $`𝒳(1)`$, and speak of “$`N`$-isogenies” between the equivalence classes of ppas’s parametrized by $`𝒳^{}(1)`$. One new feature of the congruence subgroups of $`\mathrm{\Gamma }^{}(1)`$ is that, while $`\mathrm{\Gamma }^{}(N)`$ is still normal in $`\mathrm{\Gamma }^{}(1)`$, the quotient group may be larger than $`\mathrm{PSL}_2(𝐙/N)`$, due to the presence of the $`𝗐_p`$. For instance if $`l\mathrm{\Sigma }`$ is prime then $`\mathrm{\Gamma }^{}(1)/\mathrm{\Gamma }^{}(l)`$ is $`\mathrm{PSL}_2(𝐅_l)`$ only if all the primes in $`\mathrm{\Sigma }`$ are squares modulo $`l`$; otherwise the quotient group is $`\mathrm{PGL}_2(𝐅_l)`$. In either case the index of $`\mathrm{\Gamma }_0^{}(N)`$ in $`\mathrm{\Gamma }^{}(1)`$, and thus also the degree of the cover $`𝒳_0^{}(N)/𝒳^{}(1)`$, is $`l+1`$.
Since these curves are all defined over $`𝐐`$, they can again be regarded as curves over $`𝐑`$ by a suitable choice of $`(\mathrm{\Gamma }:2)`$. For instance, if $`\mathrm{\Gamma }=\mathrm{\Gamma }(N)`$, $`\mathrm{\Gamma }_1(N)`$, $`\mathrm{\Gamma }_0(N)`$ we obtain $`(\mathrm{\Gamma }:2)`$ by adjoining $`a𝒪`$ of norm $`1`$ such that $`a(\genfrac{}{}{0pt}{}{10}{01})modN`$ under our identification of $`𝒪/N𝒪`$ with $`M_2(𝐙/N)`$. Note however that most of the automorphisms $`\mathrm{PSL}_2(𝐙/N)`$ of $`𝒳(N)`$ do not commute with $`(\genfrac{}{}{0pt}{}{10}{01})`$ and thus do not act on $`𝒳(N)`$ regarded as a real curve. Similar remarks apply to $`\mathrm{\Gamma }^{}(N)`$ etc.
Now fix a prime $`l\mathrm{\Sigma }`$ and consider the sequence of modular curves $`X_r=𝒳_0(l^r)`$ or $`X_r=𝒳_0^{}(l^r)`$ ($`r=0,1,2,\mathrm{}`$). The $`r`$-th curve parametrizes $`l^r`$-isogenies, which is to say sequences of $`l`$-isogenies
$$A_0A_1A_2\mathrm{}A_n$$
(24)
such that the composite isogeny $`A_{j1}A_{j+1}`$ is a cyclic $`l^2`$-isogeny for each $`j`$ with $`0<j<n`$. Thus for each $`m=0,1,\mathrm{},n`$ there are $`n+1m`$ maps $`\pi _j:X_nX_m`$ obtained by extracting for some $`j=0,1,\mathrm{},nm`$ the cyclic $`l^m`$-isogeny $`A_jA_{j+m}`$ from (24). Each of these maps has degree $`l^{nm}`$, unless $`m=0`$ when the degree is $`(l+1)l^{n1}`$. In particular we have a tower of maps
$$X_n\stackrel{\pi _0}{}X_{n1}\stackrel{\pi _0}{}X_{n2}\stackrel{\pi _0}{}\mathrm{}\stackrel{\pi _0}{}X_2\stackrel{\pi _0}{}X_1,$$
(25)
each map being of degree $`l`$. We observed in \[E6, Prop. 1\] that explicit formulas for $`X_1,X_2`$, together with their involutions $`w_l,w_{l^2}`$ and the map $`\pi _0:X_2X_1`$, suffice to exhibit the entire tower (25) explicitly: For $`n2`$ the product map
$$\pi =\pi _0\times \pi _1\times \pi _2\times \mathrm{}\times \pi _{n2}:X_nX_2^{n1}$$
(26)
is a 1:1 map from $`X_n`$ to the set of $`(P_1,P_2,\mathrm{},P_{n1})X_2^{n1}`$ such that
$$\pi _0\left(w_{l^2}(P_j)\right)=w_l\left(\pi _0(P_{j+1})\right)$$
(27)
for each $`j=1,2,\mathrm{},n2`$. Here we note that this information on $`X_1,X_2`$ is in turn determined by explicit formulas for $`X_0,X_1`$, together with the involution $`w_l`$ and the map $`\pi _0:X_1X_0`$. Indeed $`\pi _1:X_1X_0`$ is then $`\pi _0w_l`$, and the product map $`\pi _0\times \pi _1:X_2X_1^2`$ identifies $`X_2`$ with a curve in $`X_1^2`$ contained in the locus of
$$\{(Q_1,Q_2)X_1^2:\pi _1(Q_1)=\pi _0(Q_2)\},$$
(28)
which decomposes as the union of that curve with the graph of $`w_l`$.<sup>5</sup><sup>5</sup>5 This is where we use the hypothesis that $`l`$ is prime. The description of $`X_n`$ in (26,27) holds even for composite $`l`$, but the description of $`X_2`$ in terms of $`X_1`$ does not, because then (28) has other components. This determines $`X_2`$ and the projections $`\pi _j:X_2X_1`$ ($`j=0,1`$); the involution $`w_{l^2}`$ is
$$(Q_1,Q_2)(w_lQ_2,w_lQ_1).$$
(29)
Thus the equations we shall exhibit for certain choices of $`𝖠`$ and $`l`$ suffice to determine explicit formulas for towers of Shimura modular curves $`𝒳_0(l^r)`$, $`𝒳_0^{}(l^r)`$, towers whose reduction at any prime $`l^{}\mathrm{\Sigma }\{l\}`$ is known to be asymptotically optimal over the field of $`l_{}^{}{}_{}{}^{2}`$ elements \[I3, TVZ\].
### 2.4 Complex-multiplication (CM) and supersingular points on Shimura curves
Let $`F`$ be a quadratic imaginary field, and let $`O_F`$ be its ring of integers. Assume that none of the primes of $`\mathrm{\Sigma }`$ split in $`F`$. Then $`F`$ embeds in $`𝖠`$ (in many ways), and $`O_F`$ embeds in $`𝒪`$. For any embedding $`\iota :F𝖠`$, the image of $`F^{}`$ in $`𝖠^{}/𝐐^{}`$ then has a unique fixed point on $``$; the orbit of this point under $`\mathrm{\Gamma }(1)`$, or under any other congruence subgroup $`\mathrm{\Gamma }𝖠^{}/𝐐^{}`$, is then a CM point on the Shimura curve $`/\mathrm{\Gamma }`$. In particular, on $`𝒳(1)`$ such a point parametrizes a ppas with extra endomorphisms by $`\iota (F)𝒪`$. For instance if $`\iota (F)𝒪=\iota (O_F)`$ then this ppas is a product of elliptic curves each with complex multiplication by $`O_F`$ (but not in the product polarization). In general $`\iota ^1(\iota (F)𝒪)`$ is called the CM ring of the CM point on $`𝒳(1)`$. Embeddings conjugate by $`\mathrm{\Gamma }(1)`$ yield the same point on $`𝒳(1)`$, and for each order $`OF`$ there are finitely many embeddings up to conjugacy, and thus finitely many CM points on $`𝒳(1)`$ with CM ring $`O`$; in fact their number is just the class number of $`O`$. In \[S2\] Shimura already showed that all points with the same CM ring are Galois conjugate over $`𝐐`$, from which it follows that a CM point is rational if and only if its CM ring has unique factorization.
Thus far the description is completely analogous to the theory of complex multiplication for $`j`$-invariants of elliptic curves. But when $`\mathrm{\Sigma }\mathrm{}`$ a new phenomenon arises: CM points on the quotient curve $`𝒳^{}(1)`$ may be rational even when their preimages on $`𝒳(1)`$ are not. For instance, a point with CM ring $`O_F`$ is rational on $`𝒳^{}(1)`$ if and only if the class group of $`F`$ is generated by the classes of ideals $`IO_F`$ such that $`I^2`$ is the principal ideal $`(p)`$ for some rational prime $`p\mathrm{\Sigma }`$. This has the amusing consequence that when $`\mathrm{\Sigma }=\{2,3\}`$ the number of rational CM points on $`𝒳^{}(1)`$ is more than twice the number of rational CM points on the classical modular curve $`X(1)`$. \[Curiously, already in the classical setting $`X(1)`$ does not hold the record: it has 13 rational CM points, whilst $`X_0^{}(6)=X_0(6)/w_2,w_3`$ has 14. The reason again is fields $`F`$ with nontrivial class group generated by square roots of the ideals $`(2)`$ or $`(3)`$, though with a few small exceptions both 2 and 3 must ramify in $`F`$. In the $`𝒳^{}(1)`$ setting the primes of $`\mathrm{\Sigma }`$ are allowed to be inert as well, which makes the list considerably longer.\] In fact for each of the first four cases $`\mathrm{\Sigma }=\{2,3\},\{2,5\},\{2,7\},\{3,5\}`$ we find more rational CM points than on any classical modular curve.
A major aim of this paper is computation of the coordinates of these points. We must first list all possible $`O`$. The class number of $`O`$, and thus of $`F`$, must be a power of 2 no greater than $`2^{\mathrm{\#}\mathrm{\Sigma }}`$. In each of our cases, $`\mathrm{\#}\mathrm{\Sigma }=2`$, so $`F`$ has class number at most 4 and we may refer to the list of imaginary quadratic number fields with class group $`(𝐙/2)^r`$ ($`r=0,1,2`$), proved complete by Arno \[A\].<sup>6</sup><sup>6</sup>6 It might be possible to avoid that difficult proof for our application, since we are only concerned with fields whose class group is accounted for by ramified primes in a given set $`\mathrm{\Sigma }`$, and it may be possible to provably list them all using the arithmetic of CM points on either classical or Shimura modular curves, as in Heegner’s proof that $`𝐐(\sqrt{163})`$ is the last quadratic imaginary field of class number 1. Given $`F`$ we easily find all possible $`O`$, and imbed each into $`𝒪`$ by finding $`a𝒪`$ such that $`(a\overline{a})^2=disc(O)`$. This gives us the CM point on $``$. But we want its coordinates on the Shimura curve $`/\mathrm{\Gamma }^{}(1)`$ as rational numbers. Actually only one coordinate is needed because $`𝒳^{}(1)`$ has genus 0 for each of our four $`\mathrm{\Sigma }`$. We recover the coordinate as a real number using our Schwarzian uniformization of $`𝒳^{}(1)`$ by $``$. (Of course a coordinate on $`𝐏^1`$ is only defined up to $`\mathrm{PGL}_2(𝐐)`$, but in each case we choose a coordinate once and for all by specifying it on the CM points.) We then recognize that number as a rational number from its continued fraction expansion, and verify that the putative rational coordinate not only agrees with our computations to as many digits as we want but also satisfies various arithmetic conditions such as those described later in this section.
Of course this is not fully satisfactory; we do not know how to prove that, for instance, $`t=13^267^2109^2139^2157^2163/2^{10}5^611^617^6`$ (see Tables 1,2 below) is the CM point of discriminant $`163`$ on the curve $`𝒳^{}(1)`$ associated with the algebra ramified at $`\{2,3\}`$. But we can prove that above half of our numbers are correct, again using the modular curves $`𝒳^{}(l)`$ and their involutions $`w_l`$ for small $`l`$. This is because CM points behave well under isogenies: any point isogenous to a CM point is itself CM, and moreover a point on $`𝒳(1)`$ or $`𝒳^{}(1)`$ is CM if and only if it admits a cyclic $`d`$-isogeny to itself for some $`d>1`$. Once we have formulas for $`𝒳_0^{}(l)`$ and $`w_l`$ we may compute all points cyclically $`l`$-isogenous either with an already known CM points or with themselves. The discriminant of a new rational CM point can then be determined either by arithmetic tests or by identifying it with a real CM point to low precision.
The classical theory of supersingular points also largely carries over to the Shimura setting. We may use the fact that the ppas parametrized by a CM point has extra endomorphisms to define CM points of Shimura curves algebraically, and thus in any characteristic $`\mathrm{\Sigma }`$. In positive characteristic $`p\mathrm{\Sigma }`$, any CM point is defined over some finite field, and conversely every $`\overline{𝐅_p}`$-point of a Shimura curve is CM. All but finitely many of these parametrize ppas’s whose endomorphism ring has $`𝐙`$-rank 8; the exceptional points, all defined over $`𝐅_{p^2}`$, yield rank 16, and are called supersingular, all other $`\overline{𝐅_p}`$-points being ordinary. One may choose coordinates on $`𝒳(1)`$ (or $`𝒳^{}(1)`$) such that a CM point in characteristic zero reduces mod $`p`$ to a ordinary point if $`p`$ splits in the CM field, and to a supersingular point otherwise. Conversely each ordinary point mod $`p`$ lifts to a unique CM point (cf. \[D\] for the classical case). This means that if two CM points with different CM fields have the same reduction mod $`p`$, their common reduction is supersingular, and then as in \[GZ\] there is an upper bound on $`p`$ proportional to the product of the two CM discriminants. So for instance if $`𝒳^{}(1)𝐏^1`$ then the difference between the coordinates of two rational CM points is a product of small primes. This remains the case, for similar reasons, even for distinct CM points with the same CM field, and may be checked from the tables of rational CM points in this paper. The preimages of the supersingular points on modular covers such as $`𝒳_0(l)`$ yield enough $`𝐅_{p^2}`$-rational points on these curves to attain the Drinfeld-Vlăduţ bound \[I3\]; these curves are thus “asymptotically optimal” over $`𝐅_{p^2}`$. Asymptotically optimal curves over $`𝐅_{p^{2f}}`$ ($`f>1`$) likewise come from Shimura curves associated to quaternion algebras over totally real number fields with a prime of residue field $`𝐅_{p^f}`$.
In the case of residue field $`𝐅_p`$ (so in particular for quaternion algebras over $`𝐐`$) Ihara \[I2\] found a remarkable connection between the hyperbolic uniformization of a Shimura curve $`𝒳=/\mathrm{\Gamma }`$ and the supersingular points of its reduction mod $`p`$. We give his result in the case that $`𝒳`$ has genus 0, because we will only apply it to such curves and the result can be stated in an equivalent and elementary form (though the proof is still far from elementary). Since we are working over $`𝐅_p`$, we may identify any curve of genus 0 with $`𝐏^1`$, and choose a coordinate (degree-1 function) $`t`$ on $`𝐏^1`$ such that $`t=\mathrm{}`$ is an elliptic point. Let $`t_i`$ be the coordinates of the remaining elliptic points.
First, the hyperbolic area of the curve controls the number of points, which is approximately $`\frac{1}{2}(p+1)\mathrm{Area}(𝒳)`$ — “approximately” because $`\frac{1}{2}(p+1)\mathrm{Area}(𝒳)`$ is not the number of points but their total mass. The mass of a non-elliptic supersingular point is $`1`$, but an elliptic point with stabilizer $`G`$ has mass $`1/\mathrm{\#}G`$. If the elliptic point mod $`p`$ is the reduction of only one elliptic point on $`/\mathrm{\Gamma }`$ (which, for curves coming from quaternion algebras over $`𝐐`$, is always the case once $`p>3`$), then its stabilizer is $`𝐙/e𝐙`$ and its mass is $`1/e`$ where $`e`$ is the index of that elliptic point. \[The mass formula also holds for $`𝒳`$ of arbitrary genus, and for general residue fields provided $`p`$ is replaced by the size of the field.\] Let $`d`$ be the number of non-elliptic supersingular points, and choose a Schwarzian equation (16) with at most regular singularities at $`t=\mathrm{},t_i`$ and no other singularities. Then the supersingular points are determined uniquely by the condition that their $`t`$-coordinates are the roots of a polynomial $`P(t)`$ of degree $`d`$ such that for some $`r_i𝐐`$ the algebraic function $`_i(tt_i)^{r_i}P(t)`$ is a solution of the Schwarzian differential equation (16)! For instance \[I2, 4.3\], if $`\mathrm{\Gamma }`$ is a triangle group we may choose $`t_i=0,1`$, and then $`P(t)`$ is a finite hypergeometric series mod $`p`$.
Given $`t_0𝐐`$ we may then test whether $`t_0`$ is ordinary or supersingular mod $`p`$ for each small $`p`$. If $`t_0`$ is a CM point with CM field then its reduction is ordinary if $`p`$ splits in $`F`$, supersingular otherwise. When we have obtained $`t_0`$ as a good rational approximation to a rational CM point, but could not prove it correct, we checked for many $`p`$ whether $`t_0`$ is ordinary or supersingular mod $`p`$; when each prime behaves as expected from its behavior in $`F`$, we say that $`t_0`$ has “passed the supersingular test” modulo those primes $`p`$.
## 3 The case $`\mathrm{\Sigma }=\{2,3\}`$
### 3.1 The quaternion algebra and the curves $`𝒳(1)`$, $`𝒳^{}(1)`$
For this section we let $`𝖠`$ be the quaternion algebra ramified at $`\{2,3\}`$. This algebra is generated over $`𝐐`$ by elements $`b,c`$ satisfying
$$b^2=2,c^2=3,bc=cb.$$
(30)
The conjugation of $`𝖠`$ fixes 1 and takes $`b,c,bc`$ to $`b,c,bc`$; thus for any element $`\alpha =\alpha _1+\alpha _2b+\alpha _3c+\alpha _4bc𝖠`$ the conjugate and norm of $`\alpha `$ are given by
$$\overline{\alpha }=\alpha _1\alpha _2b\alpha _3c\alpha _4bc,\mathrm{N}(\alpha )=\alpha _1^22\alpha _2^2+3\alpha _3^26\alpha _4^2.$$
(31)
Since $`𝖠`$ is indefinite, all its maximal orders are conjugate; let $`𝒪`$ be the maximal order generated by $`b`$ and $`(1+c)/2`$. Then $`\mathrm{\Gamma }^{}(1)`$ contains $`\mathrm{\Gamma }(1)`$ with index $`2^{\mathrm{\#}\mathrm{\Sigma }}=4`$, and consists of the classes mod $`𝐐^{}`$ of elements of $`𝒪`$ of norm 1, 2, 3, or 6. In row II of Table 3 of \[T\] (p.208) we find that $`\mathrm{\Gamma }^{}(1)`$ is isomorphic with the triangle group
$$G_{2,4,6}:=s_2,s_4,s_6|s_2^2=s_4^4=s_6^6=s_2s_4s_6=1.$$
(32)
Indeed we find that $`\mathrm{\Gamma }^{}(1)`$ contains elements
$$s_2=[bc+2c],s_4=[(2+b)(1+c)],s_6=[3+c]$$
(33)
\[NB $`(2+b)(1+c),3+c2𝒪`$\] of orders $`2,4,6`$ with $`s_2s_4s_6=1`$. The subgroup of $`\mathrm{\Gamma }^{}(1)`$ generated by these elements is thus isomorphic with $`G_{2,4,6}`$. But a hyperbolic triangle group cannot be isomorphic with a proper subgroup (since the areas of the quotients of $``$ by the group and its subgroup are equal), so $`\mathrm{\Gamma }^{}(1)`$ is generated by $`s_2,s_4,s_6`$. Note that these generators have norms $`6,2,3`$ mod $`(𝐐^{})^2`$, and thus represent the three nontrivial cosets of $`\mathrm{\Gamma }^{}(1)`$ in $`𝒪^{}/\{\pm 1\}`$.
Since $`\mathrm{\Gamma }^{}(1)`$ is a triangle group, $`𝒳^{}(1)`$ is a curve of genus 0. Moreover $`𝒳^{}(1)`$ has $`𝐐`$-rational points (e.g. the three elliptic points, each of which must be rational because it is the only one of its index), so $`𝒳^{}(1)𝐏^1`$ over $`𝐐`$. Let $`t`$ be a rational coordinate on that curve (i.e. a rational function of degree 1). In general a rational coordinate on $`𝐏^1`$ is determined only up to the $`\mathrm{PGL}_2`$ action on $`𝐏^1`$, but can be specified uniquely by prescribing its values at three points. In our case $`𝒳^{}(1)`$ has three distinguished points, namely the elliptic points of orders $`2,4,6`$; we fix $`t`$ by requiring that it assume the values $`0,1,\mathrm{}`$ respectively at those three points.
None of $`s_2,s_4,s_6`$ is contained in $`\mathrm{\Gamma }(1)`$. Hence the $`(𝐙/2)^2`$ cover $`𝒳(1)/𝒳^{}(1)`$ is ramified at all three elliptic points. Thus $`s_2`$ lies under two points of $`𝒳(1)`$ with trivial stabilizer, while $`s_4`$ lies under two points of index 2 and $`s_6`$ under two points of index 3. By either the Riemann-Hurwitz formula or from (10) we see that $`𝒳(1)`$ has genus 0. This and the orders $`2,2,3,3`$ of the elliptic points do not completely specify $`\mathrm{\Gamma }(1)`$ up to conjugacy in $`\mathrm{PSL}_2(𝐑)`$: to do that we also need the cross-ratio of the four elliptic points. Fortunately this cross-ratio is determined by the existence of the cover $`𝒳(1)𝒳^{}(1)`$, or equivalently of an involution $`s_4`$ on $`𝒳(1)`$ that fixes the two order-2 points and switches the order-3 points. This forces the pairs of order-2 and order-3 points to have a cross-ratio of $`1`$, or to “divide each other harmonically” as the Greek geometers would say. The function field of $`𝒳(1)`$ is generated by the square roots of $`c_0t`$ and $`c_1(t1)`$ for some $`c_0,c_1𝐐^{}/𝐐_{}^{}{}_{}{}^{2}`$, but we do not yet know which multipliers $`c_0,c_1`$ are appropriate. If both $`c_0,c_1`$ were 1 then $`𝒳(1)`$ would be a rational curve with coordinate $`u`$ with $`t=((u^2+1)/2u)^2=1+((u^21)/2u)^2`$, the familiar parametrization of Pythagorean triples. The elliptic points of order 2 and 3 would then be at $`u=\pm 1`$ and $`u=0,\mathrm{}`$. However it will turn out that the correct choices are $`c_0=1,c_1=3`$, and thus that $`𝒳(1)`$ is the conic with equation
$$X^2+Y^2+3Z^2=0$$
(34)
and no rational points even over $`𝐑`$. \[That $`𝒳(1)`$ is the conic (34) is announced in \[Ku, p.279\] and attributed to Ihara; that there are no real points on the Shimura curve $`𝒳(1)`$ associated to any indefinite quaternion algebra over $`𝐐`$ other than $`M_2(𝐐)`$ was already shown by Shimura \[S3\]. The equation (34) for $`𝒳(1)`$ does not uniquely determine $`c_0,c_1`$, but the local methods of \[Ku\] could probably supply that information as well.\]
### 3.2 Shimura modular curves $`𝒳_0^{}(l)`$ and $`𝒳(l)`$ for $`l=5,7,13`$
Let $`l`$ be a prime other than the primes $`2,3`$ of $`\mathrm{\Sigma }`$. We determine the genus of the curve $`𝒳_0^{}(l)`$ using the formula (10). Being a cover of $`𝒳^{}(1)`$ of degree $`l+1`$, the curve $`𝒳_0^{}(l)`$ has normalized hyperbolic area $`(l+1)/12`$. It has $`1+(6/l)`$ elliptic points of order 2, $`1+(1/l)`$ elliptic points of order 4, and $`1+(3/l)`$ elliptic points of order 6. This is a consequence of our computation of $`s_2,s_4,s_6`$, which lift to elements of $`𝖠`$ that generate subfields isomorphic with $`𝐐(\sqrt{6})`$, $`𝐐(\sqrt{1})`$, and $`𝐐(\sqrt{3})`$. Actually the orders $`2,4,6`$ of the elliptic points suffice. Consider the images of $`s_2,s_4,s_6`$ in the Galois group ($`\mathrm{PGL}_2(𝐅_l)`$) of the cover $`𝒳_0^{}(l)/𝒳^{}(1)`$, and the cycle structures of their actions on the $`l+1`$ points of $`𝐏^1(𝐅_l)`$. These images $`\sigma _2,\sigma _4,\sigma _6`$ are group elements of order $`2,4,6`$. For 4 and 6, the order determines the conjugacy class, which joins as many of the points of $`𝐏^1(𝐅_l)`$ as possible in cycles of length 4 or 6 respectively and leaves any remaining points fixed; the number of fixed points is two or none according to the residue of $`l`$ mod 4 or 6. For $`\sigma _2`$ there are two conjugacy classes in $`\mathrm{PGL}_2(𝐅_l)`$, one with two fixed points and the other with none, but the choice is determined by the condition that the genus $`g(𝒳_0^{}(l))`$ be an integer, or equivalently by the requirement that the signs of $`\sigma _2,\sigma _4,\sigma _6`$ considered as permutations of $`𝐏^1(𝐅_l)`$ be consistent with $`s_2s_4s_6=1`$. We readily check that this means that the image of $`s_2`$ has two fixed points if and only if $`(6/l)=+1`$, as claimed. From (10) we conclude that
$$g(𝒳_0^{}(l))=\frac{1}{24}\left[l6\left(\frac{6}{l}\right)9\left(\frac{1}{l}\right)10\left(\frac{3}{l}\right)\right].$$
(35)
We tabulate this for $`l<50`$:
$`l`$ 5 7 11 13 17 19 23 29 31 37 41 43 47 $`g(𝒳_0^{}(l))`$ 0 0 1 0 1 1 2 1 1 1 2 2 3
It so happens that in the first seven cases $`g(𝒳_0^{}(l))`$ coincides with the genus of the classical modular curve X$`{}_{0}{}^{}(l)`$, but of course this cannot go on forever because the latter genus is $`l/12+O(1)`$ while the former is only $`l/24+O(1)`$, and indeed $`g(𝒳_0^{}(l))`$ is smaller for all $`l>23`$. Still, as with X$`{}_{0}{}^{}(l)`$, we find that $`𝒳_0^{}(l)`$ has genus 0 for $`l=5,7,13`$, but not for $`l=11`$ or any $`l>13`$. For the three genus-0 cases we shall use the ramification behavior of the cover $`𝒳_0^{}(l)/𝒳^{}(1)`$ to find an explicit rational function of degree $`l+1`$ on $`𝐏^1`$ that realizes that cover and determine the involution $`w_l`$.
Now for any $`l>3`$ the solution of $`\sigma _2\sigma _4\sigma _6=1`$ in elements $`\sigma _2,\sigma _4,\sigma _6`$ of orders $`2,4,6`$ in $`\mathrm{PGL}_2(𝐅_l)`$ is unique up to conjugation in that group. Thus we know from the general theory of \[Mat\] that the cover $`𝒳_0^{}(l)/𝒳^{}(1)`$ is determined by its Galois group and ramification data. Unfortunately the proof of this fact does not readily yield an efficient computation of the cover; for instance the Riemann existence theorem for Riemann surfaces is an essential ingredient. We use a method for finding the rational function $`t:𝒳_0^{}(l)𝒳^{}(1)`$ explicitly that amounts to solving for its coefficients, using the cycle structures of $`\sigma _2,\sigma _4,\sigma _6`$ to obtain algebraic conditions. In effect these conditions are the shape of the divisors $`(t)_0`$, $`(t)_1`$, $`(t)_{\mathrm{}}`$. But a rational function satisfying these conditions is not in general known to have the right Galois group: all we know is that the monodromy elements around $`0,1,\mathrm{}`$ have the right cycle structures in the symmetric group S<sub>l+1</sub>. Thus we obtain several candidate functions, only one of which has Galois group $`\mathrm{PGL}_2(𝐅_l)`$ (or $`\mathrm{PSL}_2(𝐅_l)`$ if $`l1mod24`$). Fortunately for $`l=5,7`$ we can exclude the impostors by inspection, and for $`l=13`$ the computation has already been done for us.
l=5. Here the cycle structures of $`s_2,s_4,s_6`$ are 2211, 411, 6. Curiously if the identity in the symmetric group $`S_6`$ is written as the product of three permutations $`\sigma _2,\sigma _4,\sigma _6`$ with these cycle structures then they can never generate all of $`S_6`$. This can be seen by considering their images $`\sigma _2^{},\sigma _4^{},\sigma _6^{}`$ under an outer automorphism of $`S_6`$: these have cycle structures 2211, 411, 321, and thus have too many cycles to generate a transitive subgroup (if two permutations of $`n`$ letters generate a transitive subgroup of $`S_n`$ then they and their product together have at most $`n+2`$ cycles). It turns out that the subgroup generated by $`\sigma _2^{},\sigma _4^{},\sigma _6^{}`$ can be either $`A_4\times S_2`$ or the point stabilizer $`S_5`$. In the former case $`\sigma _2,\sigma _4,\sigma _6`$ generate a transitive but imprimitive subgroup of $`S_6`$: the six letters are partitioned into three pairs, and the group consists of all permutations that respect this partition and permute the pairs cyclically. In the latter case $`\sigma _2,\sigma _4,\sigma _6`$ generate $`\mathrm{PGL}_2(𝐅_5)`$; this is the case we are interested in. In each of the two cases the triple $`(\sigma _2,\sigma _4,\sigma _6)`$ is determined uniquely up to conjugation in the subgroup of $`S_6`$ generated by the $`\sigma `$’s, each of which is in a rational conjugacy class in the sense of \[Mat\]. Thus each case corresponds to a unique degree-6 cover $`𝐏^1𝐏^1`$ defined over $`𝐐`$. We shall determine both covers.
Let $`t`$ be a rational function on $`𝐏^1`$ ramified only above $`t=0,1,\mathrm{}`$ with cycle structures 2211, 411, 6. Choose a rational coordinate $`x`$ on $`𝐏^1`$ such that $`x=\mathrm{}`$ is the sextuple pole of $`t`$ and $`x=0`$ is the quadruple zero of $`t1`$; this determines $`x`$ up to scaling. Then $`t`$ is a polynomial of degree 6 in $`x`$ with two double roots such that $`t1modx^4`$. The double roots are necessarily the roots of the quadratic polynomial $`x^3dt/dx`$. Thus $`t`$ is a polynomial of the form $`c_6x^6+c_5x^5+c_4x^4+1`$ divisible by $`6c_6x^2+5c_5x+4c_4`$. We readily compute that there are two possibilities for $`c_4,c_5,c_6`$ up to scaling $`(c_4,c_5,c_6)(\lambda ^4c_4,\lambda ^5c_5,\lambda ^6c_6)`$. One possibility gives $`t=2x^63x^4+1=(x^21)^2(2x^2+1)`$; being symmetric under $`xx`$ this must be the imprimitive solution. Thus the remaining possibility must give the $`\mathrm{PGL}_2(𝐅_5)`$ cover $`𝒳_0^{}(5)/𝒳^{}(1)`$. The following choice of scaling of $`x=x_5`$ seems simplest:
$$t=540x^6+324x^5+135x^4+1$$
(36)
$$=1+27x^4(20x^2+12x+5)=(15x^26x+1)(6x^2+3x+1)^2.$$
The elliptic points of order 2 and 4 on $`𝒳_0^{}(5)`$ are the simple zeros of $`t`$ and $`t1`$ respectively, i.e. the roots of $`15x^26x+1`$ and $`20x^2+12x+5`$. The involution $`w_5`$ switches each elliptic point with the other elliptic point of the same order; this suffices to determine $`w_5`$. The fact that two pairs of points on $`𝐏^1`$ switched by an involution of $`𝐏^1`$ determine the involution is well-known, but we have not found in the literature an explicit formula for doing this. Since we shall need this result on several occasion we give it in an Appendix as Proposition A. Using that formula (89), we find that
$$w_5(x)=\frac{4255x}{55+300x}.$$
(37)
l=7. This time $`s_2,s_4,s_6`$ have cycle structures 22211, 44, 611. Again there are several ways to get the identity permutation on 8 letters as a product of three permutations with these cycle structures, none of which generate the full symmetric group $`S_8`$. There are two ways to get the imprimitive group $`2^4:S_4`$; the corresponding covers are obtained from the $`S_4`$ cover $`t=4\xi ^33\xi ^4`$ by taking $`\xi =x^2+\xi _0`$ where $`\xi _0`$ is either root of the quadratic $`3\xi ^2+2\xi +1=(1t)/(\xi 1)^2`$. The remaining solution corresponds to our $`\mathrm{PGL}_2(𝐅_7)`$ cover.
To find that cover, let $`t`$ be a rational function on $`𝐏^1`$ ramified only above $`t=0,1,\mathrm{}`$ with cycle structures 2211, 411, 6, and choose a rational coordinate $`x`$ on $`𝐏^1`$ such that $`x=\mathrm{}`$ is the sextuple pole of $`t`$. This determines $`x`$ up to an affine linear transformation. Then there is a cubic polynomial $`P`$ and quadratic relatively prime polynomials $`Q_1,Q_2,Q_3`$ in $`x`$ such that $`t=P^2Q_1/Q_3=1+Q_2^4/Q_3`$, i.e. such that $`P^2Q_1Q_2^4`$ is quadratic. Equivalently, the Taylor expansion of $`Q_2^2/\sqrt{Q_1}`$ about $`x=\mathrm{}`$ should have vanishing $`x^1`$ and $`x^2`$ coefficients, and then $`R(x)`$ is obtained by truncating that Taylor expansion after its constant term. We assume without loss of generality that $`Q_1,Q_2`$ are monic. By translating $`x`$ (a.k.a. “completing the square”) we may assume that $`Q_1`$ is of the form $`x^2+\alpha `$. If the same were true of $`Q_2`$ then $`t`$ would be a rational function of $`x^2`$ and we would have an imprimitive cover. Thus the constant coefficient of $`Q_2`$ is nonzero, and by scaling $`x`$ we may take $`Q_2=x^2+x+\beta `$. We then set the $`x^1,x^2`$ coefficients of of $`Q_2^2/\sqrt{Q_1}`$ to zero, obtaining the equations
$$3\alpha ^28\alpha \beta +8\beta ^24\alpha =3\alpha ^24\alpha \beta =0.$$
(38)
Thus either $`\alpha =0`$ or $`\alpha =4\beta /3`$. The first option yields $`\beta =0`$ which fails because then $`Q_1,Q_2`$ have the common factor $`x`$. The second option yields $`\beta =0`$, which again fails for the same reason, but also $`\beta =2`$ which succeeds. Substituting $`(2x+1)/3`$ for $`x`$ to reduce the coefficients we then find:
$$t=\frac{(4x^2+4x+25)(2x^33x^2+12x2)^2}{108(7x^28x+37)}$$
(39)
$$=1\frac{(2x^2x+8)^4}{108(7x^28x+37)}.$$
The elliptic points of order 2 and 6 on $`𝒳_0^{}(7)`$ are respectively the simple zeros and poles of $`t`$, i.e. the roots of $`4x^2+4x+25`$ and $`7x^238x+7`$. The involution $`w_7`$ is again by the fact that it switches each elliptic point with the other elliptic point of the same order: it is
$$w_7(x)=\frac{1169x}{9+20x}.$$
(40)
l=13. Here the cycle structures are $`2^7`$, 44411, 6611. The computation of the degree-14 map is of course much more complicated than for the maps of degrees $`6,8`$ for $`l=5,7`$. Fortunately this computation was already done in \[MM, §4\] (a paper concerned not with Shimura modular curves but with examples of rigid $`\mathrm{PSL}_2(𝐅_p)`$ covers of the line). There we find that there is a coordinate $`x=x_{13}`$ on $`𝒳_0^{}(13)`$ for which
$$t=1\frac{27}{4}\frac{(x^2+36)(x^3+x^2+35x+27)^4}{(7x^2+2x+247)(x^2+39)^6}$$
(41)
$$=\frac{(x^750x^6+63x^55040x^4+783x^3168426x^26831x1864404)^2}{4(7x^2+2x+247)(x^2+39)^6}.$$
The elliptic points of order 4 and 6 on $`𝒳_0^{}(13)`$ are respectively the simple zeros and poles of $`t1`$, i.e. the roots of $`x^2+36`$ and $`7x^2+2x+247`$. Once more we use (89) to find the involution from the fact that it switches each elliptic point with the other elliptic point of the same order:
$$w_{13}(x)=\frac{5x+72}{2x5}.$$
(42)
From an equation for $`𝒳^{}(l)`$ and the rational map $`t`$ on that curve we recover $`𝒳_0(l)`$ by adjoining square roots of $`c_0t`$ and $`c_1(t1)`$. For each of our three cases $`l=5,7,13`$ the resulting curve has genus 1, and its Jacobian is an elliptic curve of conductor $`6l`$ — but only if we choose $`c_0,c_1`$ that give the correct quadratic twist. For $`l=5`$, $`l=7`$, $`l=13`$ it turns out that we must take a square root of $`3t(1t)`$, $`t`$, $`3(t1)`$ respectively. Fortunately these are consistent and we obtain $`c_0=1`$ and $`c_1=3`$ as promised. The resulting curves $`𝒳_0(5),𝒳_0(7),𝒳_0(13)`$ have no rational or even real points (because this is already true of the curve $`𝒳(1)`$ which they all cover); their Jacobians are the curves numbered 30F, 42C, 78B in the Antwerp tables in \[BK\] compiled by Tingley et al., and and 30-A6, 42-A3,78-A2 in Cremona \[C\].
### 3.3 Supersingular points on $`𝒳^{}(1)modl`$
We have noted that Ihara’s description of supersingular points on Shimura curves is particularly simple in the case of a triangle group: the non-elliptic supersingular points are roots of a hypergeometric polynomial, and the elliptic points are CM in characteristic zero so the Deuring test determines whether each one is supersingular or not.
In our case, The elliptic points $`t=0`$, $`t=1`$, $`t=\mathrm{}`$ are supersingular mod $`l`$ if and only iff $`l`$ is inert in $`𝐐(\sqrt{6})`$, $`𝐐(\sqrt{1})`$, $`𝐐(\sqrt{3})`$ respectively, i.e. iff $`6`$, $`1`$, $`3`$ is a quadratic nonresidue of $`l`$. Thus the status of all three elliptic points depends on $`lmod24`$, as shown in the next table:
$`l`$ mod 24 $`t`$ $`e`$ 1 5 7 11 13 17 19 23 0 2 $``$ $``$ $``$ $``$ 1 4 $``$ $``$ $``$ $``$ $`\mathrm{}`$ 6 $``$ $``$ $``$ $``$
(bullets mark elliptic points with supersingular reduction). This could also be obtained from the total mass $`(l+1)/24`$ of supersingular points, together with the fact that the contribution to this mass of the non-elliptic points is integral: in each column the table shows the unique subset of $`1/2,1/4,1/6`$ whose sum is congruent to $`(l+1)/24`$ mod 1. The hypergeometric polynomial whose roots are the non-elliptic supersingular points has degree $`l/24`$, and depends on $`lmod24`$ as follows:
$$\{\begin{array}{cc}F(\frac{1}{24},\frac{5}{24};\frac{1}{2};t),\hfill & \text{if }l1\text{ or 5 mod 24;}\hfill \\ F(\frac{7}{24},\frac{11}{24};\frac{1}{2};t),\hfill & \text{if }l7\text{ or 11 mod 24;}\hfill \\ F(\frac{13}{24},\frac{17}{24};\frac{3}{2};t),\hfill & \text{if }l13\text{ or 17 mod 24;}\hfill \\ F(\frac{19}{24},\frac{23}{24};\frac{3}{2};t),\hfill & \text{if }l19\text{ or 23 mod 24.}\hfill \end{array}$$
(43)
For example, for $`l=163(19mod24)`$ we find
$`F({\displaystyle \frac{19}{24}},{\displaystyle \frac{23}{24}};{\displaystyle \frac{3}{2}};t)`$ $`=`$ $`43t^6+89t^5+97t^4+52t^3+149t^2+132t+1`$ (44)
$`=(t+76)(t+78)(t+92)(t+127)(t^2+65t+74)`$
in characteristic 163, so the supersingular points mod 163 are $`0,1`$, and the roots of (44) in $`𝐅_{163^2}`$.
### 3.4 CM points on $`𝒳^{}(1)`$ via $`𝒳_0^{}(l)`$ and $`w_l`$
We noted already that the elliptic points $`t=0,1,\mathrm{}`$ on $`𝒳^{}(1)`$ are CM points, with discriminants $`3,4,24`$. Using our formulas for $`𝒳_0(l)`$ and $`w_l`$ ($`l=5,7,13`$) we can obtain fourteen further CM points: three points isogenous to one of the elliptic CM points, and eleven more points cyclically isogenous to themselves. This accounts for all but ten of the 27 rational CM points on $`𝒳^{}(1)`$. The discriminants of the three new points isogenous to $`t=1`$ or $`t=\mathrm{}`$ are determined by the isogenies’ degrees. The discriminants of the self-isogenous points can be surmised by testing them for supersingular reduction at small primes: in each case only one discriminant small enough to admit a self-isogeny of that degree has the correct quadratic character at the first few primes, which is then confirmed by extending the test to all primes up to 200.
On $`𝒳_0^{}(5)`$ the image of $`x_5=\mathrm{}`$ under $`w_5`$ is $`11/60`$, which yields the CM point $`t=152881/138240`$; likewise from $`w_5(0)=42/55`$ we recover the point $`421850521/1771561`$. These CM points are 5-isogenous with the elliptic points $`t=\mathrm{}`$, $`t=1`$ respectively, and thus have discriminants $`35^2`$ and $`45^2`$. Similarly on $`𝒳_0^{}(7)`$ we have $`w_7(\mathrm{})=9/20`$ at which $`t=1073152081/3024000000`$, a CM point 7-isogenous with $`t=\mathrm{}`$ and thus of discriminant $`37^2`$.
For each of $`l=5,7,13`$ the two fixed points of $`w_l`$ on $`𝒳_0^{}(l)`$ are rational and yields two new CM points of discriminants $`cl`$ for some factors $`c`$ of $`24`$. For $`𝒳_0^{}(5)`$ these fixed points are $`x_5=3/5`$ and $`x_5=7/30`$, at which $`t=2312/125`$ and $`t=5776/3375`$ respectively; these CM points have discriminants $`40`$, $`120`$ by the supersingular test. For $`𝒳_0^{}(7)`$ we find $`x_7=2`$ and $`x_7=29/10`$, and thus $`t=169/27`$, $`t=701784/15625`$ of discriminants $`84`$, $`168`$ divisible by $`7`$. For $`𝒳_0^{}(13)`$ the fixed points $`x_{13}=9`$, $`x_{13}=4`$ yield $`t=6877/15625`$ and $`t=27008742384/27680640625`$, with discriminants $`52=413`$ and $`312=2413`$.
Each of these new CM points admits an $`l`$-isogeny to itself. By solving the equation $`t(x_l)=t(w_l(x_l))`$ we find the remaining such points; those not accounted for by fixed points of $`w_l`$ admit two self-isogenies of degree $`l`$, and correspond to a quadratic pair of $`x_l`$ values over $`𝐐(t)`$. As it happens all the $`t`$’s thus obtained are rational with the exception of a quadratic pair coming from the quartic $`167x_{13}^460x_{13}^3+12138x_{13}^21980x_{13}+221607=0`$. Those points are: from $`𝒳_0^{}(5)`$, the known $`t=1`$, $`t=169/25`$, and the new $`t=1377/1024`$, $`t=3211/1024`$ of discriminants $`51`$, $`19`$; from $`𝒳_0^{}(7)`$, the CM points $`t=0`$, $`152881/138240`$, $`3211/1024`$, $`2312/125`$, $`6877/15625`$ seen already, but also $`t=13689/15625`$ of discriminant $`132`$; and from $`𝒳_0^{}(13)`$, seven of the CM points already known and also the two new values $`t=21250987/16000000`$, $`15545888/20796875`$ of discriminants $`43`$, $`88`$.
### 3.5 Numerical computation of CM points on $`𝒳^{}(1)`$
If we could obtain equations for the modular cover of $`𝒳^{}(1)`$ by the elliptic curve $`𝒳^{}(11)`$, $`𝒳^{}(17)`$ or $`𝒳^{}(19)`$ we could similarly find a few more rational CM points on $`𝒳^{}(1)`$. But we do not know how to find these covers, let alone the cover $`𝒳^{}(l)`$ for $`l`$ large enough to get at the rational CM point of discriminant $`163`$; moreover, some applications may require irrational CM points of even higher discriminants. We thus want a uniform way of computing the CM points of any given discriminant as an algebraic irrationality. We come close to this by finding these points and their algebraic conjugates as real (or, in the irrational case, complex) numbers to high precision, and then using continued fractions to recognize their elementary symmetric functions as rational numbers. We say that this “comes close” to solving the problem because, unlike the case of the classical modular functions such as $`j`$, we do not know a priori how much precision is required, since the CM values are generally not integers, nor is an effective bound known on their height. However, even when we cannot prove that our results are correct using an isogeny of low degree, we are quite confident that the rational numbers we exist are correct because they not only match their numerical approximations to many digits but also pass all the supersingularity tests we tried as well as the condition that differences between pairs of CM values are products of small primes as in \[GZ\].
To do this we must be able to compute numerically the rational function $`t:/\mathrm{\Gamma }^{}(1)\stackrel{}{}𝐏^1`$. Equivalently, we need to associate to each $`t𝐏^1`$ a representative of its corresponding $`\mathrm{\Gamma }^{}(1)`$-orbit in $``$. We noted already that this is done, up to a fractional linear transformation over $`𝐂`$, by the quotient of two hypergeometric functions in $`t`$. To fix the transformation we need images of three points, and we naturally choose the elliptic points $`t=0,1,\mathrm{}`$. These go to fixed points of $`s_2,s_4,s_6\mathrm{\Gamma }^{}(1)`$, and to find those fixed points we need an explicit action of $`\mathrm{\Gamma }^{}(1)`$ on $``$.
To obtain such an action we must imbed that group into $`\mathrm{Aut}()=\mathrm{PSL}_2(𝐑)`$. Equivalently, we must choose an identification of $`𝖠𝐑`$ with the algebra $`M_2(𝐑)`$ of $`2\times 2`$ real matrices. Having done this, to obtain the action of some $`g\mathrm{\Gamma }^{}(1)𝖠^{}/𝐐^{}`$ on $``$ we will choose a representative of $`g`$ in $`𝖠^{}`$, identify this representative with an invertible matrix $`(\genfrac{}{}{0pt}{}{ab}{cd})`$ of positive determinant, and let $`g`$ act on $`z`$ by $`z(az+b)/(cz+d)`$. Identifying $`𝖠𝐑`$ with $`M_2(𝐑)`$ is in turn tantamount to solving (30) in $`M_2(𝐑)`$.
We choose the following solution:
$$b:=\left(\begin{array}{cc}\sqrt{2}& 0\\ 0& \sqrt{2}\end{array}\right),c:=\left(\begin{array}{cc}0& \sqrt{3}\\ \sqrt{3}& 0\end{array}\right).$$
(45)
The elliptic points are then the $`\mathrm{\Gamma }^{}(1)`$ orbits of the fixed points in the upper half-plane of $`s_2,s_4,s_6`$, that is, of
$$P_2:=(1+\sqrt{2})i,P_4:=\frac{1+\sqrt{2}}{\sqrt{3}}(1+\sqrt{2}i),P_6:=i.$$
(46)
Thus for $`|t|<1`$ the point on $`/\mathrm{\Gamma }^{}(1)`$ which maps to $`t`$ is the $`\mathrm{\Gamma }^{}(1)`$ orbit of $`z`$ near $`P_2`$ such that
$$(zP_2)/(z\overline{P}_2)=F_1(t)/F_2(t)$$
(47)
for some solutions $`F_1,F_2`$ of the hypergeometric equation (18). Since the fractional linear transformation $`z(zP_2)/(z\overline{P}_2)`$ takes the hyperbolic lines $`\overline{P_2P_4}`$ and $`\overline{P_2P_6}`$ to straight lines through the origin, $`F_2`$ must be a power series in $`t`$, and $`F_1`$ is such a power series multiplied by $`\sqrt{t}`$; that is,
$$(zP_2)/(z\overline{P}_2)=Ct^{1/2}F(\frac{13}{24},\frac{17}{24},\frac{3}{2},t)/F(\frac{1}{24},\frac{5}{24},\frac{1}{2},t)).$$
(48)
for some nonzero constant $`C`$. We evaluate $`C`$ by taking $`t=1`$ in (48). Then $`z=P_4`$, which determines the left-hand side, while the identity \[GR, 9.122\]
$$F(a,b;c;1)=\frac{\mathrm{\Gamma }(c)\mathrm{\Gamma }(cab)}{\mathrm{\Gamma }(ca)\mathrm{\Gamma }(cb)}$$
(49)
gives us the coefficient of $`C`$ in the right-hand side in terms of gamma functions. We find $`C=(.314837\mathrm{})i/(2.472571\mathrm{})=(.128545\mathrm{})i`$. Likewise we obtain convergent power series for computing $`z`$ in neighborhoods of $`t=1`$ and $`t=\mathrm{}`$.
Now let $`D`$ be the discriminant of an order $`O_D`$ in a quadratic imaginary field $`𝐐(\sqrt{D})`$ such that $`O_D`$ has a maximal embedding in $`𝒪`$ (i.e. an embedding such that $`O_D=(O_D𝐐)𝒪`$) and the embedding is unique up to conjugation in $`\mathrm{\Gamma }^{}(1)`$. Then there is a unique, and therefore rational, CM point on $`𝒳^{}(1)`$ of discriminant $`D`$. Being rational, the point is real, and thus can be found on one of the three hyperbolic line segments $`\overline{P_2P_4}`$, $`\overline{P_2P_6}`$, $`\overline{P_4P_6}`$. It is thus the fixed point of a positive integer combination, with coprime coefficients, of two of the elliptic elements $`s_2=bc+2c`$, $`s_4=(2+b)(1+c)/2`$, $`s_6=(3+c)/2`$ with fixed points $`P_2,P_4,P_6`$. In each case a short search finds the appropriate linear combination and thus the fixed point $`z`$. Using (48) or the analogous formulas near $`t=1`$, $`t=\mathrm{}`$ we then solve for $`t`$ as a real number with sufficient accuracy (60 decimals was more than enough) to recover it as a rational number from its continued-fraction expansion.
### 3.6 Tables of rational CM points on $`𝒳^{}(1)`$
There are 27 rational CM points on $`𝒳^{}(1)`$. We write the discriminant $`D`$ of each of them as $`D_0D_1`$ where $`D_0|24`$ and $`D_1`$ is coprime to $`6`$. In Table 1 we give, for each $`|D|=D_0D_1`$, the integers $`A,B`$ with $`B0`$ such that $`(A:B)`$ is the $`t`$-coordinate of a CM point of discriminant $`D`$. In the last column of this table we indicate whether the point was obtained algebraically (via an isogeny of degree 5, 7, or 13) and thus proved correct, or only computed numerically. The CM points are listed in order of increasing height $`\mathrm{max}(|A|,B)`$.
In Table 2 we give, for each except the first three cases, the factorizations of $`|A|,B,|C|`$ where $`C=AB`$, and also the associated “$`ABC`$ ratio” \[E1\] defined by $`r=\mathrm{log}N(ABC)/\mathrm{log}\mathrm{max}(|A|,B,|C|)`$. As expected, the $`A,B,C`$ values are “almost” perfect squares, sixth powers, and fourth powers respectively: a prime at which at which the valuation of $`A,B,C`$ is not divisible by 2, 6, 4 resp. is either 2, 3, or the unique prime in $`D_1`$. When $`D_1>1`$ its unique prime factor is listed at the end of the $`|A|`$, $`B`$, or $`|C|`$ factorization in which it appears; otherwise the prime factors are listed in increasing order.
In the factorization of the difference between the last two $`t=A/B`$ values in this table, the primes not accounted for by common factors in the last two rows of the table are 79, 127, 271, 907, 2287, 2971, 3547, each occurring once.
## 4 The case $`\mathrm{\Sigma }=\{2,5\}`$
### 4.1 The quaternion algebra and the curves $`𝒳(1)`$, $`𝒳^{}(1)`$
For this section we let $`𝖠`$ be the quaternion algebra ramified at $`\{2,5\}`$. This time $`𝖠`$ is generated over $`𝐐`$ by elements $`b,e`$ satisfying
$$b^2+2=e^25=be+eb=0,$$
(50)
and the conjugate and norm of an element $`\alpha =\alpha _1+\alpha _2b+\alpha _3e+\alpha _4be𝖠`$ are
$$\overline{\alpha }=\alpha _1\alpha _2b\alpha _3e\alpha _4be,\alpha \overline{\alpha }=\overline{\alpha }\alpha =\alpha _1^2+2\alpha _2^25\alpha _3^210\alpha _4^2.$$
(51)
The elements $`b`$ and $`(1+e)/2`$ generate a maximal order, which we use for $`𝒪`$.
By (9), the curve $`𝒳^{}(1)`$ has hyperbolic area $`1/6`$. Since the algebra $`𝖠`$ is not among the nineteen algebras listed in \[T\] that produce arithmetic triangle groups, $`𝒳^{}(1)`$ must have at least four elliptic points. On the other hand, by (10) a curve of area as small as $`1/6`$ cannot have more than four elliptic points, and if it has exactly four then their orders must be $`2,2,2,3`$. Indeed we find in $`\mathrm{\Gamma }^{}(1)`$ the elements of finite order
$$s_2=[b],s_2^{}=[2e+5bbe],s_2^{\prime \prime }=[5bbe],s_3=[2be1]$$
(52)
\[NB $`2e+5bbe,5bbe,2be12𝒪`$\] of orders $`2,2,2,3`$ with $`s_2^{}s_2^{}s_2^{\prime \prime }s_3^{}=1`$. As in the case of the $`G_{2,4,6}`$ we conclude that here $`\mathrm{\Gamma }^{}(1)`$ has the presentation
$$s_2^{},s_2^{},s_2^{\prime \prime },s_3^{}|s_2^2=s_{2}^{}{}_{}{}^{2}=s_{2}^{\prime \prime }{}_{}{}^{2}=s_3^3=s_2^{}s_2^{}s_2^{\prime \prime }s_3^{}=1.$$
(53)
Of the four generators only $`s_3`$ is in $`\mathrm{\Gamma }(1)`$; thus the $`(𝐙/2)^2`$ cover $`𝒳(1)/𝒳^{}(1)`$ is ramified at the elliptic points of order 2. Therefore $`𝒳(1)`$ is a rational curve with four elliptic points of order 3, and $`\mathrm{\Gamma }^(1)`$ is generated by four 3-cycles whose product is the identity, for example by $`s_3`$ and its conjugates by $`s_2,s_2^{},s_2^{\prime \prime }`$. (The genus and number of elliptic points of $`𝒳(1),𝒳^{}(1)`$, but not the generators of $`\mathrm{\Gamma }(1),\mathrm{\Gamma }^{}(1)`$, are already tabulated in \[V, Ch.IV:2\].)
### 4.2 Shimura modular curves $`𝒳_0^{}(l)`$, in particular $`𝒳_0^{}(3)`$
The elliptic elements $`s_3^{},s_2^{},s_2^{},s_2^{\prime \prime }`$ have discriminants $`3,8,20,40`$. Thus the curve $`𝒳_0^{}(l)`$ has genus
$$g(𝒳_0^{}(l))=\frac{1}{12}\left[l4\left(\frac{3}{l}\right)3\left(\frac{2}{l}\right)3\left(\frac{5}{l}\right)3\left(\frac{10}{l}\right)\right].$$
(54)
Again we tabulate this for $`l<50`$:
$`l`$ 3 7 11 13 17 19 23 29 31 37 41 43 47 $`g(𝒳_0^{}(l))`$ 0 0 1 1 2 1 2 3 3 3 3 3 4
Since $`g(𝒳_0^{}(l))(l13)/12`$, the cases $`l=3,7`$ of genus 0 occurring in this table are the only ones. We next find an explicit rational functions of degree 4 on $`𝐏^1`$ that realizes the cover $`𝒳_0^{}(3)/𝒳_0^{}(1)`$, and determine the involution $`w_3`$.
The curve $`𝒳_0^{}(3)`$ is a degree-4 cover of $`𝒳^{}(1)`$ with Galois group $`\mathrm{PGL}_2(𝐅_3)`$ and cycle structures 31, 211, 211, 22 over the elliptic points $`P_3^{},P_2^{},P_2^{},P_2^{\prime \prime }`$. Thus there are coordinates $`\tau ,x`$ on $`𝒳^{}(1)`$, $`𝒳_0^{}(3)`$ such that $`\tau (x)=(x^2c)^2/(x1)^3`$ for some $`c`$. To determine the parameter $`c`$, we use the fact that $`w_3`$ fixes the simple pole $`x=\mathrm{}`$ and takes each simple preimage of the 211 points $`P_2^{},P_2^{}`$ to the other simple preimage of the same point. That is,
$$(x^2c)^1(x1)^4\frac{dx}{dt}=x^24x+3c$$
(55)
must have distinct roots $`x_i`$ ($`i=1,2`$) that yield quadratic polynomials
$$\frac{(x1)^3(\tau (x)\tau (x_i))}{(xx_i)^2}$$
(56)
with the same $`x`$ coefficient. We find that this happens only for $`c=5/3`$, i.e. that $`\tau =(3x^2+5)^2/9(x1)^3`$. For future use it will prove convenient to use
$$t=\frac{6^3}{9\tau +8}=\frac{(6x6)^3}{(x+1)^2(9x^210x+17)},$$
(57)
with $`w_3(x)=\frac{10}{9}x`$. \[Smaller coefficients can be obtained by letting $`x=1+2/x^{}`$, $`\tau =2t^{}/9`$, when $`t^{}=(2x^2+3x^{}+3)^2/x^{}`$ and $`w_3(x^{})=9x^{}/(4x^{}+9)`$. But our choice of $`x`$ will simplify the computation of the Schwarzian equation, while the choice of $`t`$ will turn out to be the correct one 3-adically.\] The elliptic points are then $`P_3:t=0`$, $`P_2^{\prime \prime }:t=27`$, and $`P_2,P_2^{}:t=\mathrm{},2`$. In fact the information so far does not exclude the possibility that the pole of $`t`$ might be at $`P_2^{}`$ instead of $`P_2`$; that in fact $`t(P_2)=\mathrm{},t(P_2^{})=2`$ and not the other way around can be seen from the order of the elliptic points on the real locus of $`𝒳^{}(1)`$, or (once we compute the Schwarzian equation) checked using the supersingular test.
### 4.3 CM points on $`𝒳^{}(1)`$ via $`𝒳_0^{}(3)`$ and $`w_3`$
From $`w_3`$ we obtain five further CM points. Three of these are 3-isogenous to known elliptic points: $`w_3`$ takes the triple zero $`x=1`$ of $`t`$ to $`x=1/9`$, which gives us $`t=192/25`$, the point 3-isogenous to $`P_3`$ with discriminant $`27`$; likewise $`w_3`$ takes the double root $`x=5`$ and double pole $`x=1`$ of $`t2`$ to $`x=35/9,19/9`$ and thus to $`t=2662/169`$ and $`t=125/147`$, the points 3-isogenous to $`t=2`$ and $`t=\mathrm{}`$ and thus (once these points are identified with $`P_2^{}`$ and $`P_2`$) of discriminants $`180`$ and $`72`$. One new CM point comes from the other fixed point $`x=5/9`$ of $`w_3`$, which yields $`t=27/49`$ of discriminant $`120`$. Finally the remaining solutions of $`t(x)=t(w_3(x))`$ are the roots of $`9x^210x+65`$; the resulting CM point $`t=64/7`$, with two 3-isogenies to itself, turns out to have discriminant $`35`$.
### 4.4 The Schwarzian equation on $`𝒳^{}(1)`$
We can take the Schwarzian equation on $`𝒳^{}(1)`$ to be of the form
$$t(t2)(t27)f^{\prime \prime }+(At^2+Bt+C)f^{}+(Dt+E)=0.$$
(58)
The coefficients $`A,B,C,D`$ are then forced by the indices of the elliptic points. Near $`t=0`$, the solutions of (58) must be generated by functions with leading terms $`1`$ and $`t^{1/3}`$; near $`t=2`$ ($`t=27`$), by functions with leading terms 1 and $`(t2)^{1/2}`$ (resp. $`(t27)^{1/2}`$); and at infinity, by functions with leading terms $`t^e`$ and $`t^{e1/2}`$ for some $`e`$. The conditions at the three finite singular points $`t=0,2,27`$ determine the value of the $`f^{}`$ coefficient at those points, and thus yield $`A,B,C`$, which turn out to be $`5/3,203/6,36`$. Then $`e,e+1/2`$ must be roots of an “indicial equation” $`e^22e/3+D=0`$, so $`e=1/12`$ and $`D=7/144`$.
Thus (58) becomes
$$t(t2)(t27)f^{\prime \prime }+\frac{10t^2203t+216}{6}f^{}+(\frac{7t}{144}+E)=0.$$
(59)
To determine the “accessory parameter” $`E`$, we again use the cover $`𝒳_0^{}(3)/𝒳^{}(1)`$ and the involution $`w_3`$. A Schwarzian equation for $`𝒳_0^{}(3)`$ is obtained by substituting $`t=(6x6)^3/(x+1)^2(9x^210x+17)`$ in (59). The resulting equation will not yet display the $`w_3`$ symmetry, because it will have a spurious singular point at the double pole $`x=1`$ of $`t(x)`$. To remove this singularity we consider not $`f(t(x))`$ but
$$g(x):=(x+1)^{1/6}f(t(x)).$$
(60)
The factor $`(x+1)^{1/6}`$ is also singular at $`x=\mathrm{}`$, but that is already an elliptic point of $`𝒳_0^{}(3)`$ and a fixed point of $`w_3`$. Let $`x=u+5/9`$, so $`w_3`$ is simply $`uu`$. Then we find that the differential equation satisfied by $`g`$ is
$$4(81u^2+20)(81u^2+128)^2g^{\prime \prime }+108u(81u^2+128)(405u^2+424)g^{}$$
(61)
$$+(3^{11}u^4163296u^2+170496+72(18E+7)(9u4)(81u^2+128))g=0.$$
Clearly this has the desired symmetry if and only if $`18E+7=0`$. Thus the Schwarzian equation is
$$t(t2)(t27)f^{\prime \prime }+\frac{10t^2203t+216}{6}f^{}+(\frac{7t}{144}\frac{7}{18})=0.$$
(62)
### 4.5 Numerical computation of CM points on $`𝒳^{}(1)`$
We can now expand a basis of solutions of (62) in power series about each singular point $`t=0,2,27,\mathrm{}`$ (using inverse powers of $`t\frac{27}{2}`$ for the expansion about $`\mathrm{}`$ to assure convergence for real $`t[0,27]`$). As with the $`\mathrm{\Sigma }=\{2,3\}`$ case we need to identify $`𝖠𝐑`$ with $`M_2(𝐑)`$, and use the solution
$$b:=\left(\begin{array}{cc}0& \sqrt{2}\\ \sqrt{2}& 0\end{array}\right),e:=\left(\begin{array}{cc}\sqrt{5}& 0\\ 0& \sqrt{5}\end{array}\right).$$
(63)
of (50), analogous to (30). We want to proceed as we did for $`\mathrm{\Sigma }=\{2,3\}`$, but there is still one obstacle to computing, for given $`t_0𝐑`$, the point on the hyperbolic quadrilateral formed by the fixed points of $`s_2^{},s_2^{},s_2^{\prime \prime },s_3^{}`$ at which $`t=t_0`$. In the $`\mathrm{\Sigma }=\{2,3\}`$ case, the solutions of the Schwarzian equation were combinations of hypergeometric functions, whose value at $`1`$ is known. This let us determine two solutions whose ratio gives the desired map to $``$. But here $`\mathrm{\Gamma }^{}(1)`$ is not a triangle group, so our basic solutions of (62 are more complicated power series and we do not know a priori their values at the neighboring singular points. In general this obstacle can be overcome by noting that for each nonsingular $`t_0𝐑`$ its image in $``$ can be computed from the power-series expansions about either of its neighbors and using the condition that the two computations agree for several choices of $`t_0`$ to determine the maps to $``$. In our case we instead removed the obstacle using the non-elliptic CM points computed in the previous section. For example, we used the fact that $`t_0=125/147`$ is the CM point of discriminant $`72`$, and thus maps to the unique fixed point in $``$ of $`(9b+4ebe)/2`$, to determine the correct ratio of power series about $`t=0`$ and $`t=2`$. Two or three such points suffice to determine the four ratios needed to compute our map $`𝐑`$ to arbitrary accuracy; since we actually had five non-elliptic CM points, we used the extra points for consistency checks, and then used the resulting formulas to numerically compute the $`t`$-coordinates of the remaining CM points.
There are 21 rational CM points on $`𝒳^{}(1)`$. We write the discriminant $`D`$ of each of them as $`D_0D_1`$ where $`D_0|40`$ and $`D_1`$ is coprime to 10. Table 3 is organized in the same way as Table 1: we give, for each $`|D|=D_0D_1`$, the integers $`A,B`$ with $`B0`$ such that $`(A:B)`$ is the $`t`$-coordinate of a CM point of discriminant $`D`$. The last column identifies with a “Y” the nine points obtained algebraically from the computation of $`𝒳_0^{}(3)`$ and $`w_3`$. Some but not all of the remaining twelve points would move from “N” to “Y” if we also had the equations for the degree-8 map $`𝒳_0^{}(7)𝒳^{}(1)`$ and the involution $`w_7`$ on $`𝒳_0^{}(7)`$.
It will be seen that the factor $`3^3`$ in our normalization (57) of $`t`$ was needed<sup>7</sup><sup>7</sup>7 On the other hand the factor $`2^3`$ in (57) was a matter of convenience, to make the four elliptic points integral. to make $`t`$ a good coordinate 3-adically: $`3`$ splits in the CM field iff $`t`$ is not a multiple of $`3`$.
In Table 4 we give the factorizations of $`|A|,B,|A2B|,|A27B|`$; as expected, $`|A|`$ is always “almost” a perfect cube, and $`B,|A2B|,|A27B|`$ “almost” a perfect square, any exceptional primes other than 2 or 5 being the unique prime in $`D_1`$, which if it occurs is listed at the end of its respective factorization.
## 5 Further examples and problems
Our treatment here is briefer because most of the ideas and methods of the previous sections apply here with little change. Thus we only describe new features that did not arise for the algebras ramified at $`\{2,3\}`$ and $`\{2,5\}`$, and exhibit the final results of our computations of modular curves and CM points.
### 5.1 The case $`\mathrm{\Sigma }=\{2,7\}`$
We generate $`𝖠`$ by elements $`b,g`$ with
$$b^2+2=g^27=bg+gb=0,$$
(64)
and a maximal order $`𝒪`$ by $`𝐙[b,g]`$ together with $`(1+b+g)/2`$ (and $`b(1+g)/2`$). By (9), the curve $`𝒳^{}(1)`$ has hyperbolic area $`1/4`$. Since $`\mathrm{\Gamma }^{}(1)`$ is not a triangle group (again by \[T\]), we again conclude by (10) that $`𝒳^{}(1)`$ has exactly four elliptic points, this time of orders $`2,2,2,4`$. We find in $`\mathrm{\Gamma }^{}(1)`$ the elements of finite order
$$s_2=[b],s_2^{}=[7b2gbg],s_2^{\prime \prime }=[7b+2gbg],s_4=[1+2b+g]$$
(65)
\[NB $`7b\pm 2gbg2𝒪`$\] of orders $`2,2,2,4`$ with $`s_2^{}s_2^{}s_2^{\prime \prime }s_4^{}=1`$, and conclude that $`s_2^{},s_2^{},s_2^{\prime \prime },s_4^{}`$ generate $`\mathrm{\Gamma }^{}(1)`$ with relations determined by $`s_2^2=s_{2}^{}{}_{}{}^{2}=s_{2}^{\prime \prime }{}_{}{}^{2}=s_4^4=s_2s_2^{}s_2^{\prime \prime }s_4=1`$. None of these is in $`\mathrm{\Gamma }^{}(1)`$: the representatives $`b,1+2b+g`$ of $`s_2,s_4`$ have norm $`2`$, while $`s_2^{},s_2^{\prime \prime }`$ have representatives $`(7b\pm 2gbg)/2`$ of norm $`14`$. The discriminants of $`s_4,s_2,s_2^{},s_2^{\prime \prime }`$ are $`4,8,56,56`$; note that $`56`$ is not among the “idoneal” discriminants (discriminants of imaginary quadratic fields with class group $`(𝐙/2)^r`$), and thus that the elliptic fixed points $`P_2^{},P_2^{\prime \prime }`$ of $`s_2^{},s_2^{\prime \prime }`$ are quadratic conjugates on $`𝒳^{}(1)`$.
Again we use the involution $`w_3`$ on the modular curve $`𝒳_0^{}(3)`$ to simultaneously determine the relative position of the elliptic points $`P_4,P_2,P_2^{},P_2^{\prime \prime }`$ on $`𝒳^{}(1)`$ and the modular cover $`𝒳_0^{}(3)𝒳^{}(1)`$, and then to obtain a Schwarzian equation on $`𝒳^{}(1)`$. Clearly $`P_4`$ is completely ramified in $`𝒳_0^{}(3)`$. Since $`8`$ and $`56`$ are quadratic residues of $`3`$, each of $`P_2,P_2^{},P_2^{\prime \prime }`$ has ramification type 211. Thus $`𝒳_0^{}(3)`$ is a rational curve with six elliptic points all of index 2, and we may choose coordinates $`t,x`$ on $`𝒳^{}(1),𝒳_0^{}(3)`$ such that $`t(P_4)=\mathrm{}`$, $`t(P_2)=0`$, and $`x=\mathrm{}`$, $`x=0`$ at the quadruple pole and double zero respectively of $`t`$.
We next determine the action of $`w_3`$ on the elliptic points of $`𝒳_0^{}(3)`$. Necessarily the simple preimages of $`P_2`$ parametrize two 3-isogenies from $`P_2`$ to itself. On the other hand the simple preimages of $`P_2^{}`$ parametrize two 3-isogenies from that point to $`P_2^{\prime \prime }`$ and vice versa, because the squares of the primes above $`3`$ in $`𝐐(\sqrt{14})`$ are not principal. Therefore $`w_3`$ exchanges the simple preimages of $`P_2`$ but takes each of the two simple points above $`P_2^{}`$ to one above $`P_2^{\prime \prime }`$ and vice versa.
So again we have a one-parameter family of degree-4 functions on $`𝐏^1`$, and a single condition in the existence of the involution $`w_3`$; but this time it turns out that there are (up to scaling the coordinates $`t,x`$) two ways to satisfy this condition:
$$t=\frac{1}{3}(x^4+4x^3+6x^2),w_3(x)=\frac{1x}{1+x},P_2^{},P_2^{\prime \prime }:t^23t+3=0$$
(66)
and
$$t=\frac{1}{27}(x^4+2x^3+9x^2),w_3(x)=\frac{52x}{2+x},P_2^{},P_2^{\prime \prime }:16t^2+13t+8=0.$$
(67)
How to choose the correct one? We could consider the next modular curve $`𝒳_0^{}(5)`$ and its involution to obtain a new condition that would be satisfied by only one of (66,67). Fortunately we can circumvent this laborious calculation by noting that the Fuchsian group associated with (66) is commensurable with a triangle group, since its three elliptic points of index 2 are the roots of $`(1t)^3=1`$ and are thus permuted by a 3-cycle that fixes the fourth elliptic point $`t=\mathrm{}`$. The quotient by that 3-cycle is a curve parametrized by $`(1t)^3`$ with elliptic points of order $`2,3,12`$ at $`1,0,\mathrm{}`$. But by \[T\] there is no triangle group commensurable with an arithmetic subgroup of $`𝖠^{}/𝐐^{}`$; indeed we find there that $`G_{2,3,12}`$ is associated with the quaternion algebra over $`𝐐(\sqrt{3})`$ ramified at the prime above 2 and at one of the infinite places of that number field.<sup>8</sup><sup>8</sup>8 See \[T\], table 3, row IV. In terms of that algebra $`𝖠^{}`$, the triangle group $`G_{2,3,12}`$ is $`\mathrm{\Gamma }^{}(1)`$; the index-3 normal subgroup whose quotient curve is parametrized by the $`t`$ of (66) is the normalized in $`\mathrm{\Gamma }^{}(1)`$ of $`\{[a]𝒪^{}/\{\pm 1\}:a1modI_2\}`$; and the intersection of this group with $`\mathrm{\Gamma }_0^{}(3)`$ yields as quotient curve the $`𝐏^1`$ with coordinate $`x`$. Therefore (67) is the correct choice. Alternatively, we could have noticed that since $`𝒳(1)`$ is a $`(𝐙/2)^2`$ cover of $`𝒳^{}(1)`$ ramified at all four elliptic points, it has genus 1, and then used the condition that this curve’s Jacobian have conductor 14 to exclude (66). The function field of $`𝒳^{}(1)`$ is obtained by adjoining square roots of $`c_0t`$ and $`c_1(16t^2+13t+8)`$ for some $`c_0,c_1`$; for the Jacobian to have the correct conductor we must have $`c_0c_1=1`$ mod squares. The double cover of $`𝒳_0^{}(3)`$ obtained by adjoining $`\sqrt{c_1(16t^2+13t+8)}`$ also has genus 1, and so must have Jacobian of conductor at most 42; this happens only when $`c_1=1`$ mod squares, the Jacobian being the elliptic curve 42-A3 (42C). The curve $`𝒳(1)`$ then has the equation
$$y^2=16s^4+13s^28(t=s^2),$$
(68)
and its Jacobian is the elliptic curve 14-A2 (14D). Kurihara had already obtained in \[Ku\] an equation birational with (68). Let $`\mathrm{\Gamma }_0^{}(3^r)`$ be the group intermediate between $`\mathrm{\Gamma }_0(3^r)`$ and $`\mathrm{\Gamma }_0^{}(3^r)`$ consisting of the elements of norm 1 or 7 mod $`𝐐_{}^{}{}_{}{}^{2}`$. Then the corresponding curves $`𝒳_0^{}(3^r)`$ ($`r>0`$) of genus $`3^{r1}+1`$ are obtained from $`𝒳_0^{}(3^r)`$ by extracting a square root of $`t(16t^2+13t+8)`$, and constitute an unramified tower of curves over the genus-2 curve
$$𝒳_0^{}(3):y^2=3(4x^6+12x^5+75x^4+50x^3+255x^2288x+648)$$
(69)
whose reductions are asymptotically optimal over $`𝐅_{l^2}`$ ($`l2,3,7)`$ with each step in the tower being a cyclic cubic extension. (Of course when we consider only reductions to curves over $`𝐅_{l^2}`$ the factor of 3 in (69) may be suppressed.)
Using $`w_3`$ we may again find the coordinates of several non-elliptic CM points: $`t=4/3`$ and $`t=75/16`$ of discriminants $`36`$ and $`72`$, i.e. the points 3-isogenous to $`P_4`$ and $`P_2`$, other than $`P_4,P_2`$ themselves; $`t=4/9`$ and $`t=200/9`$ of discriminants $`84`$ and $`168`$, coming from the fixed points $`x=1`$ and $`x=5`$ of $`w_3`$; and the points $`t=1`$, $`t=5`$ of discriminants $`11`$ and $`35`$, coming from the remaining solutions of $`t(x)=t(w_3(x))`$ and each with two 3-isogenies to itself.
Even once the relative position of the elliptic points is known, the computation of the cover $`𝒳_0^{}(5)/𝒳^{}(1)`$ is not a trivial matter; I thank Peter Müller for performing this computation using J.-C. Faugere’s Gröbner basis package GB. It turns out that there are eight PGL$`{}_{2}{}^{}(𝐅_5)`$ covers consistent with the ramification of which only one is defined over $`𝐐`$:
$$t=\frac{(256x^3+224x^2+232x+217)^2}{50000(x^2+1)},w_5(x)=\frac{247x}{7+24x}.$$
(70)
This yields the CM points of discriminants $`11`$, $`35`$, $`36`$, $`84`$ already known from $`w_3`$, and new points of discriminants $`91`$, $`100`$, $`280`$.
This accounts for eleven of the nineteen rational CM points on $`𝒳^{}(1)`$; the remaining ones were computed numerically as we did for the $`\mathrm{\Sigma }=\{2,5\}`$ curve. We used the Schwarzian equation
$$t(16t^2+13t+8)f^{\prime \prime }+(24t^2+13t+4)f^{}+\left(\frac{3}{4}t+\frac{3}{16}\right)f=0,$$
(71)
for which the “accessory parameter” $`3/16`$ was again determined by pulling back to $`𝒳_0^{}(3)`$ and imposing the condition of symmetry under $`w_3`$. We tabulate the coordinates $`t=A/B`$ and factorizations for all nineteen points in Table 5.
We see that $`t`$ is also a good coordinate 3-adically: a point of $`𝒳^{}(1)`$ is supersingular at 3 iff the denominator of its $`t`$-coordinate is a multiple of 3. (It is supersingular at 5 iff $`5|t`$.)
### 5.2 The case $`\mathrm{\Sigma }=\{3,5\}`$
Here the area of $`𝒳^{}(1)`$ is $`1/3`$. This again is small enough to show that there are only four elliptic points, but leaves two possibilities for their indices: 2,2,2,6 or 2,2,3,3. It turns out that the first of these is correct. This fact is contained in the table of \[V, Ch.IV:2\]; it can also be checked as we did in the cases $`\mathrm{\Sigma }=\{2,p\}`$ ($`p=3,5,7`$) by exhibiting appropriate elliptic elements of $`\mathrm{\Gamma }^{}(1)`$ — which we need to do anyway to compute the CM points. We chose to write write $`𝒪=𝐙[\frac{1}{2}1+c,e]`$ with
$$c^2+3=e^25=ce+ec=0,$$
(72)
and found the elliptic elements
$$s_2=[4c3e],s_2^{}=[5c3ece],s_2^{\prime \prime }=[20c9e7ce],s_6=[3+c]$$
(73)
\[NB $`20c9e7ce,3+c2𝒪`$\] of orders $`2,2,2,6`$ with $`s_2^{}s_2^{}s_2^{\prime \prime }s_6^{}=1`$. The corresponding elliptic points $`P_2^{},P_2^{},P_2^{\prime \prime },P_6^{}`$ have CM discriminants $`3,12,15,60`$. For the first time we have a curve $`𝒳_0^{}(2)`$, and here it turns out that the elliptic points $`P_2^{}`$ is not ramified in the cover $`𝒳_0^{}(2)/𝒳^{}(1)`$: it admits two 2-isogenies to itself, and one to $`P^{\prime \prime }`$. Of the remaining elliptic points, $`P_6`$ is totally ramified, and each of $`P_2,P_2^{\prime \prime }`$ has one simple and one double preimage. So we may choose coordinates $`x,t`$ on $`𝒳_0^{}(2)`$ and $`𝒳^{}(1)`$ such that $`t=x(x3)^2/4`$, with $`t(P_6)=\mathrm{}`$, $`t(P_2)=0`$, $`t(P_2^{\prime \prime })=1`$. To determine $`t(P_2^{})`$ we use the involution $`w_2`$, which switches $`x=\mathrm{}`$ (the triple pole) with $`x=0`$ (the simple zero), $`x=4`$ (the simple preimage of $`P_2^{\prime \prime }`$) with one of the preimages $`x_1`$ of $`P_2^{}`$ (the one parametrizing the isogeny from $`P_2^{}`$ to $`P_2^{\prime \prime }`$), and the other two preimages of $`P_2^{}`$ with each other. Then $`w_2`$ is $`x4x_1/x`$, so the product of the roots of $`(t(x_1)t(x))/(xx_1)`$ is $`4x_1`$. Thus
$$x(x3)^24t(P_2^{})=(xx_1)(x^2+ax+4x_1)$$
(74)
for some $`a`$. Equating $`x^2`$ coefficients yields $`a=x_16`$, and equating the coefficients of $`x`$ we find $`9=10x_1^{}x_1^2`$. Thus $`x_1=1`$ or $`x_1=9`$; but the first would give us $`t(P_2^{})=1=t(P_2^{\prime \prime })`$ which is impossible. Thus $`x_1=9`$ and $`t(P_2^{})=81`$, with $`w_2(x)=36/x`$. This lets us find six further rational CM points, of discriminants $`7,28,40,48,120,240`$; we can also solve for the accessory parameter $`1/2`$ in the Schwarzian equation
$$t(t1)(t81)f^{\prime \prime }+\left(\frac{3}{2}t^282t+\frac{81}{2}\right)f^{}+\left(\frac{1}{18}t\frac{1}{2}\right)f=0,$$
(75)
and use it to compute the remaining twelve rational CM points numerically. We tabulate the coordinates $`t=A/B`$ and factorizations for the twenty-two rational CM points on $`𝒳^{}(1)`$ in Table 6.
An equivalent coordinate that is also good 2-adically is $`(t1)/4`$, which is supersingular at 2 iff its denominator is even.
The elliptic curve $`𝒳(1)`$ is obtained from $`𝒳^{}(1)`$ by extracting square roots of $`At`$ and $`B(t1)(t81)`$ for some $`A,B𝐐^{}/𝐐_{}^{}{}_{}{}^{2}`$. Using the condition that the Jacobian of $`𝒳(1)`$, and any elliptic curve occurring in the Jacobian of $`𝒳_0(2)`$, have conductor at most 15 and 30 respectively, we find $`A=B=3`$. Then $`𝒳(1)`$ has equation
$$y^2=(3s^2+1)(s^2+27)$$
(76)
(with $`t=3s^2`$) and Jacobian isomorphic with elliptic curve 15C (15-A1); the curve intermediate between $`𝒳^{}(2)`$ and $`𝒳_0(2)`$ whose function field is obtained from $`𝐐(𝒳^{}(2))`$ by adjoining $`\sqrt{3(t1)(t81)}`$ has equation
$$y^2=3(x^410x^3+33x^2360x+1296)$$
(77)
and Jacobian 30C (30-A3). Fundamental domains for $`\mathrm{\Gamma }^{}(1)`$ and $`\mathrm{\Gamma }(1)`$, computed by Michon \[Mi\] and drawn by C. Léger, can be found in \[V, pp.123–127\]; an equation for $`𝒳(1)`$ birational with (76) is reported in the table of \[JL, p.235\].
### 5.3 The triangle group $`G_{2,3,7}`$ as an arithmetic group
It is well-known that the minimal quotient area of a discrete subgroup of $`\mathrm{Aut}()`$ $`=\mathrm{PSL}_2(𝐑)`$ is $`1/42`$, attained only by the triangle group $`G_{2,3,7}`$, and that the Riemann surfaces $`/\mathrm{\Gamma }`$ with $`\mathrm{\Gamma }`$ a proper normal subgroup of finite index in $`G_{2,3,7}`$ are precisely the curves of genus $`g>1`$ whose number of automorphisms attains Hurwitz’s upper bound $`84(g1)`$. Shimura observed in \[S2\] that this group is arithmetic.<sup>9</sup><sup>9</sup>9 Actually this fact is due to Fricke \[F1, F2\], over a century ago; but Fricke could not relate $`G_{2,3,7}`$ to a quaternion algebra because the arithmetic of quaternion algebras had yet to be developed. Indeed, let $`K`$ be the totally real cubic field $`𝐐(\mathrm{cos}2\pi /7)`$ of minimal discriminant $`49`$, and let $`𝖠`$ be a quaternion algebra over $`K`$ ramified at two of the three real places and at no finite primes of $`K`$. Now for any totally real number field of degree $`n>1`$ over $`𝐐`$, and any quaternion algebra over that field ramified at $`n1`$ of its real places, the group $`\mathrm{\Gamma }(1)`$ of norm-1 elements of a maximal order embeds as a discrete subgroup of $`\mathrm{PSL}_2(𝐑)=\mathrm{Aut}()`$, with $`/\mathrm{\Gamma }`$ of finite area given by Shimizu’s formula
$$\mathrm{Area}(𝒳(1))=\frac{d_K^{3/2}\zeta _K(2)}{4^{n1}\pi ^{2n}}\underset{\mathrm{}\mathrm{\Sigma }}{}(\mathrm{N}\mathrm{}1)\left[=\frac{(1)^n}{2^{n2}}\zeta _K(1)\underset{\mathrm{}\mathrm{\Sigma }}{}(\mathrm{N}\mathrm{}1)\right]$$
(78)
(from which we obtained (8) by taking $`K=𝐐`$). Thus, in our case of $`K=𝐐(\mathrm{cos}2\pi /7)`$, $`\mathrm{\Sigma }=\{\mathrm{},\mathrm{}^{}\}`$, the area of $`/\mathrm{\Gamma }(1)`$ is $`1/42`$, so $`\mathrm{\Gamma }(1)`$ must be isomorphic with $`G_{2,3,7}`$. From this Shimura deduced \[S2, p.83\] that for any proper ideal $`IO_K`$ his curve $`𝒳(I)=/\mathrm{\Gamma }(I)`$ attains the Hurwitz bound. For instance, if $`I`$ is the prime ideal $`\mathrm{}_7^{}`$ above the totally ramified prime $`7`$ of $`𝐐`$ then $`𝒳(\mathrm{}_7^{})`$ is the Klein curve of genus 3 with automorphism group $`\mathrm{PSL}_2(𝐅_7)`$ of order 168. The next-smallest example is the ideal $`\mathrm{}_8^{}`$ above the inert prime 2, which yields a curve of genus 7 with automorphism group \[P\]SL$`{}_{2}{}^{}(𝐅_8)`$ of order 504. This curve is also described by Shimura as a “known curve”, and indeed it first appears in \[F3\]; an equivalent curve was studied in detail only a few years before Shimura by Macbeath \[Mac\], who does not cite Fricke, and the identification of Macbeath’s curve with Fricke’s and with Shimura’s $`𝒳(\mathrm{}_8^{})`$ may first have been observed by Serre in a 24.vii.1990 letter to Abhyankar. At any rate, we obtain towers $`\{𝒳(\mathrm{}_7^r)\}_{r>0}`$, $`\{𝒳(\mathrm{}_8^r)\}_{r>0}`$ of unramified abelian extensions which are asymptotically optimal over the quadratic extensions of residue fields<sup>10</sup><sup>10</sup>10 That is, over the fields of size $`p^2`$ for primes $`p=7`$ or $`p\pm 1mod7`$, and $`p^6`$ for other primes $`p`$. of $`K`$ other than $`𝐅_{49}`$ and $`𝐅_{64}`$ respectively, which are involved in the class field towers of exponents $`7,2`$ of the Klein and Macbeath curves over those fields.
These towers are the Galois closures of the covers of $`𝒳(1)`$ by $`𝒳_0(\mathrm{}_7^r)`$, $`𝒳_0(\mathrm{}_8^r)`$, which again may be obtained from the curves $`𝒳_0(\mathrm{}_7^{})`$, $`𝒳_0(\mathrm{}_8^{})`$ together with their involutions. It turns out that these curves both have genus 0 (indeed the corresponding arithmetic subgroups $`\mathrm{\Gamma }_0(\mathrm{}_7^{})`$, $`\mathrm{\Gamma }_0(\mathrm{}_8^{})`$ of $`\mathrm{\Gamma }(1)`$ are the triangle groups $`G_{3,3,7}`$, $`G_{2,7,7}`$ in \[T, class X\]). The cover $`𝒳_0(\mathrm{}_7^{})/𝒳(1)`$ has the same ramification data as the degree-8 cover of classical modular curves $`\mathrm{X}_0(7)/\mathrm{X}(1)`$, and is thus given by the same rational function
$`t`$ $`=`$ $`{\displaystyle \frac{(x_7^48x_7^318x_7^288x_7^{}+1409)^2}{2^{13}3^3(9x_7^{})}}`$
$`=`$ $`1+{\displaystyle \frac{(x_7^28x_7^{}5)^3(x_7^2+8x_7^{}+43)}{2^{13}3^3(9x_7^{})}}`$
(with the elliptic points of orders $`2,3,7`$ at $`t=0,1,\mathrm{}`$, i.e. $`t`$ corresponds to $`112^3j`$). The involution is different, though: it still switches the two simple zeros $`x_7^{}=4\pm \sqrt{27}`$ of $`t1`$, but it takes the simple pole $`x_7^{}=0`$ to itself instead of the septuple pole at $`x_7^{}=\mathrm{}`$. Using (89) again we find
$$w_\mathrm{}_7^{}(x_7)=\frac{19x_7^{}+711}{13x_7^{}19}.$$
(80)
For the degree-9 cover $`𝒳_0(\mathrm{}_8^{})/𝒳(1)`$ we find
$`t`$ $`=`$ $`{\displaystyle \frac{(1x_8)(2x_8^4+4x_8^3+18x_8^2+14x_8^{}+25)^2}{27(4x_8^2+5x_8^{}+23)}}`$
$`=`$ $`1{\displaystyle \frac{4(x_8^3+x_8^2+5x_8^{}1)^3}{27(4x_8^2+5x_8^{}+23)}},`$
with the involution fixing the simple zero $`x_8=1`$ and switching the simple poles, i.e.
$$w_\mathrm{}_8^{}(x_8)=\frac{5119x_8^{}}{19+13x_8^{}}.$$
(82)
Note that all of these covers and involutions have rational coefficients even though a priori they are only known to be defined over $`K`$. This is possible because $`K`$ is a normal extension of $`𝐐`$ and the primes $`\mathrm{}_7^{}`$, $`\mathrm{}_8^{}`$ used to define our curves and maps are Galois-invariant. To each of the three real places of $`K`$ corresponds a quaternion algebra ramified only at the other two places, and thus a Shimura curve $`𝒳(1)`$ with three elliptic points $`P_2,P_3,P_7`$ to which we may assign coordinates $`0,1,\mathrm{}`$. Then Gal$`(K/𝐐)`$ permutes these three curves; since we have chosen rational coordinates for the three distinguished points, any point on or cover of $`𝒳(1)`$ defined by a Galois-invariant construction must be fixed by this action of Galois and so be defined over $`𝐐`$. The same applies to each of the triangle groups $`G_{p,q,r}`$ associated with quaternion algebras over number fields $`F`$ properly containing $`𝐐`$, which can be found in cases III through XIX of Takeuchi’s list \[T\]. In each case, $`F`$ is Galois over $`𝐐`$, and the finite ramified places of the quaternion algebra are Galois-invariant. Moreover, even when $`G_{p,q,r}`$ is not $`\mathrm{\Gamma }(1)`$, it is still related with $`\mathrm{\Gamma }(1)`$ by a Galois-invariant construction (such as intersection with $`\mathrm{\Gamma }_0(\mathrm{})`$ or adjoining $`𝗐_{\mathrm{}}`$ or $`w_{\mathrm{}}`$ for a Galois-invariant prime $`\mathrm{}`$ of $`F`$). At least one of the triangle groups in each commensurability class has distinct indices $`p,q,r`$, whose corresponding elliptic points may be unambiguously identified with $`0,1,\mathrm{}`$; this yields a model of the curve $`/G_{p,q,r}`$, and thus of all its commensurable triangle curves, that is defined over $`𝐐`$.
This discussion bears also on CM points on $`𝒳(1)`$. There are many CM points on $`𝒳(1)`$ rational over $`K`$, but only seven of those are $`𝐐`$-rational: a CM point defined over $`𝐐`$ must come from a CM field $`K^{}`$ which is Galois not only over $`K`$ but over $`𝐐`$. Thus $`K^{}`$ is the compositum of $`K`$ with an imaginary quadratic field, which must have unique factorization. We check that of the nine such fields only five retain unique factorization when composed with $`K`$. One of these, $`𝐐(\sqrt{7})`$, yields the cyclotomic field $`𝐐(e^{2\pi i/7})`$, whose ring of integers is the CM ring for the elliptic point $`P_7:t=\mathrm{}`$; two subrings still have unique factorization and yield CM points $`\mathrm{}_7^{}`$\- and $`\mathrm{}_8^{}`$-isogenous to that elliptic point, which again are not only $`K`$\- but even $`𝐐`$-rational thanks to the Galois invariance of $`\mathrm{}_7^{}`$, $`\mathrm{}_8^{}`$. The other four cases are the fields of discriminant $`3,4,8,11`$, which yield one rational CM point each. The first two are the elliptic points $`P_3,P_2:t=1,0`$. To find the coordinates of the CM point of discriminant $`8`$, and of the two points isogenous with $`P_7`$, we may use the involutions (80,82) on $`𝒳_0(\mathrm{}_7^{})`$ and $`𝒳_0(\mathrm{}_8^{})`$. On $`𝒳_0(\mathrm{}_7^{})`$, the involution takes $`x_7^{}=\mathrm{}`$ to $`19/13`$, yielding the point $`t=3593763963/4015905088`$ $`\mathrm{}_7^{}`$-isogenous with $`P_7`$ on $`𝒳(1)`$; on $`𝒳_0(\mathrm{}_8^{})`$ the involution takes $`x_8^{}=\mathrm{}`$ to $`19/13`$, yielding the point $`t=47439942003/8031810176`$ $`\mathrm{}_8^{}`$-isogenous with $`P_7`$. On the latter curve, the second fixed point of the involution (besides $`x_8^{}=1`$) is $`x_8^{}=51/13`$, which yields the CM point $`t=1092830632334/1694209959`$ of discriminant $`8`$. The two points isogenous with $`P_7`$ also arise from the second fixed point of $`w_\mathrm{}_7^{}`$ and a further solution of $`t(x_8^{})=t(w_\mathrm{}_8^{}(x_8^{}))`$.
This still leaves the problem of locating the CM point of discriminant $`11`$. We found it numerically using quotients of hypergeometric functions as we did for $`G_{2,4,6}`$. Let $`c=2\mathrm{cos}2\pi /7`$, so $`c`$ is the unique positive root of $`c^3+c^22c1`$. Consider the quaternion algebra over $`K`$ generated by $`i,j`$ with
$$i^2=j^2=c,ij=ji.$$
(83)
This is ramified at the two other real place of $`K`$, in which $`c`$ maps to the negative reals $`2\mathrm{cos}4\pi /7`$ and $`2\mathrm{cos}6\pi /7`$, but not at the place with $`c=2\mathrm{cos}2\pi /7`$; since $`c`$ is a unit, neither is this algebra ramified at any finite place with the possible exception of $`\mathrm{}_8^{}`$, which we exclude using the fact that the set of ramified places has even cardinality. Thus $`K(i,j)`$ is indeed our algebra $`𝖠`$. A maximal order $`𝒪`$ is obtained from $`O_K[i,j]`$ by adjoining the integral element $`(1+ci+(c^2+c+1)j)/2`$. Then $`𝒪^{}`$ contains the elements
$$g_2:=ij/c,g_3:=\frac{1}{2}(1+(c^22)j+(3c^2)ij),$$
(84)
$$g_7:=\frac{1}{2}(c^2+c1+(2c^2)i+(c^2+c2)ij)$$
of norm 1, with $`g_2^2=g_3^3=g_7^7=1`$ and $`g_2=g_7g_3`$. Thus the images of $`g_2,g_3,g_7`$ in $`\mathrm{\Gamma }(1)`$ are elliptic elements that generate that group. A short search finds the linear combination $`(2c^2)g_3+(c^2+c)g_7𝒪`$ of discriminant $`11`$; computing its fixed point in $``$ and solving for $`t`$ to high precision (150 decimals, which turned out to be overkill), we obtain a real number whose continued fraction matches that of
$$\frac{88983265401189332631297917}{45974167834557869095293}=\frac{7^343^2127^2139^2207^2659^211}{3^313^783^7},$$
(85)
with numerator and denominator differing by $`2^929^341^3167^3281^3`$. Having also checked that this number differs from the $`t`$-coordinates of the three non-elliptic CM points by products of small ($`<10^4`$) primes,<sup>11</sup><sup>11</sup>11 If $`10^4`$ does not seem small, remember that the factorizations are really over $`K`$, not $`𝐐`$; the largest inert prime that occurs is 19, and the split primes are really primes of $`K`$ of norm at most comparable with that of 19. and that it passes the supersingular test, we are quite confident that (85) is in fact the $`t`$-coordinate of the CM point of discriminant $`11`$.
### 5.4 An irrational example: the algebras over $`𝐐[\tau ]/(\tau ^34\tau +2)`$ with $`\mathrm{\Sigma }=\{\mathrm{}_i,\mathrm{}_j\}`$
While our examples so far have all been defined over $`𝐐`$, this is not generally the case for Shimura curves associated to a quaternion algebra over a totally number field $`K`$ properly containing $`𝐐`$. For instance, $`K`$ may not be a Galois extension of $`𝐐`$; or, $`K`$ may be Galois, but the set of finite ramified places may fails to be Galois-stable; or, even if that set is Galois-stable, the congruence conditions on the subgroup of $`𝖠^{}/K^{}`$ may not be Galois-invariant, and the resulting curve would not be defined over $`𝐐`$ even though $`𝒳(1)`$ would be. In each case different real embeddings of the field yield different arithmetic subgroups of $`\mathrm{PSL}_2(𝐑)`$ and thus different quotient curves. We give here what is probably the simplest example: a curve $`𝒳(1)`$ associated to a quaternion algebra with no finite ramified places over a totally real cubic field which is not Galois over $`𝐐`$. While the curve has genus 0, no degree-1 rational function on it takes $`𝐐`$-rational values at all four of its elliptic points, and the towers of modular curves over this $`𝒳(1)`$ are defined over $`K`$ but not over $`𝐐`$.
Let $`K`$ be the cubic field $`𝐐[\tau ]/(\tau ^34\tau +2)`$ and discriminant $`148=2^237`$, which is minimal for a totally real non-Galois field. Let $`𝖠/K`$ be a quaternion algebra ramified at two of the three real places and at no finite primes of $`K`$. Using gp/pari to compute $`\zeta _K(2)`$, we find from Shimizu’s formula (78) that the associated Shimura curve $`𝒳(1)=𝒳^{}(1)`$ has hyperbolic area $`.16666\mathrm{}`$; thus the area is $`1/6`$ and, since $`𝖠`$ is not in Takeuchi’s list, the curve $`𝒳(1)`$ has genus 0 and four elliptic points, one of order 3 and three of order 2. The order-3 point $`P_3`$ has discriminant $`3`$ as expected, but the order-2 points are a bit more interesting: their CM field is $`K(i)`$, but the ring of integers of that field is not $`O_K[i]`$! Note that the rational prime 2 is totally ramified in $`K`$, being the cube of the prime $`(\tau )`$; thus $`(1+i)/\tau `$ is an algebraic integer, and we readily check that it generates the integers of $`K(i)`$ over $`O_K`$. One of the elliptic points, call it $`P_2`$, has CM ring $`O_K[(1+i)/\tau ]`$ and discriminant $`4/\tau ^2`$; of its three $`(\tau )`$-isogenous points, one is $`P_2`$ itself, and the others are the remaining elliptic points $`P_2^{},P_2^{\prime \prime }`$, with CM ring $`O_K[i]`$ of discriminant $`4`$.
Thus the modular curve $`𝒳_0((\tau ))`$ is a degree-3 cover of $`𝒳(1)`$ unramified above the elliptic point $`P_2`$, and ramified above the other three elliptic points with type 3 for $`P_3`$ and 21 for $`P_2^{},P_2^{\prime \prime }`$. This determines the cover up to $`\overline{K}`$-isomorphism — the curve $`𝒳_0((\tau ))`$ has genus 0, and we can choose coordinates $`x`$ on that curve and $`t`$ on $`𝒳(1)`$ such that $`t(P_3)=\mathrm{}`$ and $`t=x^33cx`$ for some $`c0`$ — but not the location of the unramified point $`P_2`$ relative to the other three elliptic points. To determine that we once again use the involution, this time $`w_{(\tau )}`$, of $`𝒳_0((\tau ))`$: this involution fixes the point above $`P_2`$ corresponding to its self-isogeny, and pairs the other two preimages of $`P_2`$ with the simple preimages of $`P_2^{},P_2^{\prime \prime }`$. We find that there are three ways to satisfy this condition:
$$t=x^33(\tau ^23)x,P_2:t=1300188\tau 351\tau ^2,P_2^{},P_2^{\prime \prime }:t=\pm 2(\tau ^23)^{3/2},$$
(86)
and its Galois conjugates. The correct choice is determined by the condition that the Shimura curves must be fixed by the involution of the Galois closure of $`K/𝐐`$ that switches the two real embeddings of $`K`$ that ramify $`𝖠`$: the image of $`\tau `$ under the the third (split) embedding must be used in (86). We find that the simple and double preimages of $`P_2^{},P_2^{\prime \prime }`$ are $`x=\pm 2\sqrt{a^23}`$, $`\sqrt{a^23}`$, and the preimages of $`P_2`$ are $`122\tau 3\tau ^2`$ (fixed by $`w_{(\tau )}`$) and $`(12+2\tau +3\tau ^2\pm (3a^212)\sqrt{a^23})/2`$. From this we recover as usual the tower of curves $`𝒳_0((\tau )^r)`$, whose reductions at primes of $`K`$ other than $`\tau `$ are asymptotically optimal over the quadratic extensions of the primes’ residue fields, and which in this case is a tower of double (whence cyclic) covers unramified above the genus-3 curve $`𝒳_0((\tau )^4)`$ and thus involved in that curve’s class-field tower.
### 5.5 Open problems
#### 5.5.1 Computing modular curves and covers.
Given a nonempty even set $`\mathrm{\Sigma }`$ of rational primes, and thus a quaternion algebra $`𝖠/𝐐`$, how to compute the curve $`𝒳^{}(1)`$ together with its Schwarzian equation and modular covers such as $`𝒳(1)`$ and $`𝒳_0^{}(l)`$? Even in the simplest case $`\mathrm{\Sigma }=\{2,3\}`$ where $`\mathrm{\Gamma }^{}(1)`$ is a triangle group and all the covers $`𝒳_0^{}(l)/𝒳^{}(1)`$ are in principle determined by their ramifications, finding those covers seems at present a difficult problem once $`l`$ gets much larger than the few primes we have dealt with here. This is the case even when $`l`$ is still small enough that $`𝒳_0^{}(l)`$ has genus small enough, say $`g5`$, that the curve should have a simple model in projective space. For instance, according to 35 the curve $`𝒳_0^{}(73)`$ has genus 1. Thus its Jacobian is an elliptic curve; moreover it must be one of the six elliptic curves of conductor $`673`$ tabulated in \[C\]. Which one of those curves it is, and which principal homogeneous space of that curve is isomorphic with $`𝒳_0^{}(73)`$, can probably be decided by local methods such as those of \[Ku\]; indeed such a computation was made for $`𝒳_0(11)`$ in D. Roberts’ thesis \[Ro\]. But that still leaves the problem of finding the degree-74 map on that curve which realizes the modular cover $`𝒳_0^{}(73)𝒳^{}(1)`$. For classical modular curves (i.e. with $`\mathrm{\Sigma }=\mathrm{}`$) of comparable and even somewhat higher levels, the equations and covers can be obtained via $`q`$-expansions as explained in \[E5\]; but what can we do here in the absence of cusps and thus of $`q`$-expansions? Can we do anything at all once the primes in $`\mathrm{\Sigma }`$ are large or numerous enough to even defeat the methods of the present paper for computing $`𝒳^{}(1)`$ and the location of the elliptic points on this curve? Again this happens while the genus of $`𝒳^{}(1)`$ is still small; for instance it seems already a difficult problem to locate the elliptic points on all curves $`𝒳^{}(1)`$ of genus zero and determine their Schwarzian equations, let alone find equations for all curves $`𝒳^{}(1)`$ of genus 1, 2, or 3. By \[I2\] the existence of the involutions $`w_l`$ on $`𝒳_0^{}(l)`$ always suffices in principle to answer these questions, but the computations needed to actually do this become difficult very quickly; it seems that a perspicuous way to handle these computations, or a new and more efficient approach, is called for. The reader will note that so far we have said nothing about computing with modular forms on Shimura curves. Not only is this an intriguing question in its own right, but solving it may also allow more efficient computation of Shimura curves and the natural maps between them, as happens in the classical modular setting.
In another direction, we ask: is there a prescription, analogous to (27), for towers of Shimura curves whose levels are powers of a ramified prime of the algebra? For a concrete example (from case III of \[T\]), let $`𝖠`$ be the quaternion algebra over $`𝐐(\sqrt{2})`$ with $`\mathrm{\Sigma }=\{\mathrm{}_1,\mathrm{}_2\}`$, where $`\mathrm{}_1`$ is one of the two Archimedean places and $`\mathrm{}_2`$ is the prime ideal $`(\sqrt{2})`$ above $`2`$; let $`𝒪𝖠`$ be a maximal order, $`I=I_\mathrm{}_2𝒪`$ the ideal of elements whose norm is a multiple of $`\sqrt{2}`$, and
$$\mathrm{\Gamma }_n=\{[a]O_1^{}/\{\pm 1\}:a1modI^n\}$$
(87)
for $`n=0,1,2,\mathrm{}`$ . Then $`\mathrm{\Gamma }_{n+1}`$ is a normal subgroup of $`\mathrm{\Gamma }_n`$ with index $`3,2^2,2`$ according as $`n=0`$, $`n`$ is odd, or $`n`$ is even and positive. Consulting \[T\], we find that $`\mathrm{\Gamma }_0,\mathrm{\Gamma }_1`$ are the triangle groups $`G_{3,3,4}`$ and $`G_{4,4,4}`$. Let $`X_n`$ be the Shimura curve $`/\mathrm{\Gamma }_n`$, which parametrizes principally polarized abelian fourfolds with endomorphisms by $`𝖠`$ and complete level-$`I^n`$ structure. Then $`\{X_n\}_{n>0}`$ is a tower of $`𝐙/2`$ or $`(𝐙/2)^2`$ covers, unramified above the curve $`X_3`$. Moreover, $`X_n`$ has genus zero for $`n=0,1,2`$, while $`X_3`$ is isomorphic with the curve $`y^2=x^5x`$ of genus 2 with maximal automorphism group. The reduction of this tower at any prime $`\mathrm{}\mathrm{}_2`$ of $`𝐐(\sqrt{2})`$ is asymptotically optimal over the quadratic extension of the residue field of $`\mathrm{}`$. So we ask for explicit recursive equations for the curves in this tower. Note that unlike the tower (25), this one does not seem to offer a $`w_l`$ or $`𝗐_\mathrm{}_2`$ shortcut.
#### 5.5.2 CM points.
Once we have found a Shimura modular curve together with a Schwarzian equation, we have seen how to compute the coordinates of CM points on the curve, at least as real or complex numbers to arbitrary precision. But this still leaves many theoretical and computational questions open. For instance, what form does the Gross-Zagier formula \[GZ\] for the difference between $`j`$-invariants of elliptic curves take in the context of Shimura curves such as $`𝒳_0^{}(1)`$ or $`𝒳(1)`$? Note that a factorization theorem would also yield a rigorous proof that our tabulated rational coordinates of CM points are correct. Our tables also suggest that at least for rational CM points the heights increase more or less regularly with $`D_1`$; can this be explained and generalized to CM points of degree $`>1`$? For CM points on the classical modular curve X$`(1)`$ this is easy: a CM $`j`$-invariant is an algebraic integer, and its size depends on how close the corresponding point of $`/\mathrm{PSL}_2(𝐙)`$ is to the cusp; so for instance if $`𝐐(\sqrt{D})`$ has class number 1 then the CM $`j`$-invariant of discriminant $`D`$ is a rational integer of absolute value $`\mathrm{exp}(\pi \sqrt{D})+O(1)`$. But such a simple explanation probably cannot work for Shimura curves which have neither cusps nor integrality of CM points. Within a commensurability class of Shimura curves (i.e. given the quaternion algebra $`𝖠`$), the height is inversely proportional to the area of the curve; does this remain true in some sense when $`𝖠`$ is varied?
As a special case we might ask: how does the minimal polynomial of a CM point of discriminant $`D`$ factor modulo the primes contained in $`D_1`$? That the minimal polynomials for CM $`j`$-invariants are almost squares modulo prime factors of the discriminant was a key component of our results on supersingular reduction of elliptic curves \[E2, E3\]; analogous results on Shimura curves may likewise yield a proof that, for instance, for every $`t𝐐`$ there are infinitely many primes $`p`$ such that the point on the $`(2,4,6)`$ curve with coordinate $`t`$ reduces to a supersingular point mod $`p`$.
#### 5.5.3 Enumeration and arithmetic of covers.
When an arithmetic subgroup of $`\mathrm{PSL}_2(𝐑)`$ is commensurable with a triangle group $`G=G_{p,q,r}`$, as was the case for the $`\mathrm{\Sigma }=\{2,3\}`$ algebra, any modular cover $`/G^{}`$ of $`/G`$ (for $`G^{}G`$ a congruence subgroup) is ramified above only three points on the genus-0 curve $`/G`$. We readily obtain the ramification data, which leave only finitely many possibilities for the cover. We noted that, even when there is only one such cover, actually finding it can be far from straightforward; but much is known about covers of $`𝐏^1`$ ramified at three points — for instance, the number of such covers with given Galois group and ramification can be computed by solving equations in the group (see \[Mat\]), and the cover is known \[Be\] to have good reduction at each prime not dividing the size of the group. But when $`G`$, and any group commensurable with it, has positive genus or more than three elliptic points, we were forced to introduce additional information about the cover, namely the existence of an involution exchanging certain preimages of the branch points. In the examples we gave here (and in several others to be detailed in future work) this was enough to uniquely determine the cover $`/G^{}/G`$. But there is as yet no general theory that predicts the number of solutions of this kind of covering problem. The arithmetic of the solutions is even more mysterious: recall for instance that in our final example the cubic field $`𝐐[\tau ]/(\tau ^34\tau +2)`$ emerged out of conditions on the cover $`𝒳_0((\tau ))/𝒳(1)`$ in which that field, and even its ramified prime $`37`$, are nowhere to be seen.
## 6 Appendix: Involutions of $`𝐏^1`$
We collect some facts concerning involutions of the projective line over a field of characteristic other than 2. We do this from a representation-theoretic point of view, in the spirit of \[FH\]. That is, we identify a pair of points $`t_i=(x_i:y_i)`$ $`(i=1,2)`$ of $`𝐏^1`$ with a binary quadric, i.e. a one-dimensional space of homogeneous quadratic polynomials $`Q(X,Y)=AX^2+2BXY+CY^2`$, namely the polynomials vanishing at the two points; we regard the three-dimensional space $`V_3`$ of all such polynomials $`AX^2+2BXY+CY^2`$ as a representation of the group SL<sub>2</sub> acting on $`𝐏^1`$ by unimodular linear transformations of $`(X,Y)`$.
An invertible linear transformation of a two-dimensional vector space $`V_2`$ over any field yields an involution of the projective line $`𝐏^1=𝐏(V_2^{})`$ if and only if it is not proportional to the identity and its trace vanishes (the first condition being necessary only in characteristic 2). Over an algebraically closed field of characteristic other than 2, every involution of $`𝐏^1`$ has two fixed points, and any two points are equivalent under the action of $`\mathrm{PSL}_2`$ on $`𝐏^1`$. It is clear that the only involution fixing $`0,\mathrm{}`$ is $`tt`$; it follows that any pair of points determines a unique involution fixing those two points. Explicitly, if $`B^2AC`$, the involution fixing the distinct roots of $`AX^2+2BXY+CY^2`$ is $`(X:Y)(BX+CY:AXBY)`$. Note that the 2-transitivity of $`\mathrm{PSL}_2`$ on $`𝐏^1`$ also means that this group acts transitively on the complement in the projective plane $`𝐏V_3`$ of the conic $`B^2=AC`$ (and also acts transitively on that conic); indeed it is well-known that $`\mathrm{PSL}_2`$ is just the special orthogonal group for the discriminant quadric $`B^2AC`$ on $`V_3`$.
Now let $`Q_1,Q_2V_3`$ be two polynomials without a common zero. Then there is a unique involution of $`𝐏^1`$ switching the roots of $`Q_1`$ and also of $`Q_2`$. (If $`Q_i`$ has a double zero the condition on $`Q_i`$ means that its zero is a fixed point of the involution.) This can be seen by using the automorphism group $`\mathrm{Aut}(𝐏^1)=`$PGL<sub>2</sub> to map $`Q_i`$ to $`XY`$ or $`Y^2`$ and noting that the involutions that switch $`t=0`$ with $`\mathrm{}`$ are $`ta/t`$ for nonzero $`a`$, while the involutions fixing $`t=\mathrm{}`$ are $`tat`$ for arbitrary $`a`$. As before, we regard the involution determined in this way by $`Q_1,Q_2`$ as an element of $`𝐏V_3`$. This yields an algebraic map $`f`$ from (an open set in) $`𝐏V_3\times 𝐏V_3`$, parametrizing $`Q_1,Q_2`$ without common zeros, to $`𝐏V_3`$. We next determine this map explicitly.
First we note that this map is covariant under the action of $`\mathrm{PSL}_2`$: we have $`f(gQ_1,gQ_2)=g(f(Q_1,Q_2))`$ for any $`g\mathrm{PSL}_2`$. Next we show that $`f`$ has degree 1 in each factor. Using the action of $`\mathrm{PSL}_2`$ it is enough to show that if $`Q_1=XY`$ or $`Y^2`$ then $`f`$ is linear as a function of $`Q_2=AX^2+2BXY+CY^2`$. In the first case, the involution is $`tC/At`$ and its fixed points are the roots of $`AX^2CY^2`$. In the second case, the involution is $`t(2B/A)t`$ with fixed points $`t=\mathrm{}`$ and $`t=B/A`$, i.e. the roots of $`AXY+BY^2`$. In either case the coefficients of $`f(Q_1,Q_2)`$ are indeed linear in $`A,B,C`$.
But it turns out that these two conditions completely determine $`f`$: there is up to scaling a unique $`\mathrm{PSL}_2`$-covariant bilinear map from $`V_3\times V_3`$ to $`V_3`$; equivalently, $`V_3`$ occurs exactly once in the representation $`V_3V_3`$ of $`\mathrm{PSL}_2`$. In fact it is known (see e.g. \[FH, §11.2\]) that $`V_3V_3`$ decomposes as $`V_1V_3V_5`$, where $`V_1`$ is the trivial representation and $`V_5`$ is the space of homogeneous polynomials of degree 4 in $`X,Y`$. The factor $`V_3`$ is particularly easy to see, because it is just the antisymmetric part $`^2V_3`$ of $`V_3V_3`$. Now the next-to-highest exterior power $`^{dimV1}V`$ of any finite-dimensional vector space $`V`$ is canonically isomorphic with $`(detV)V^{}`$, where $`detV`$ is the top exterior power $`^{dimV}V`$. Taking $`V=V_3`$, we see that $`detV_3`$ is the trivial representation of $`\mathrm{PSL}_2`$. Moreover, thanks to the invariant quadric $`B^2AC`$ we know that $`V_3`$ is self-dual as a $`\mathrm{PSL}_2`$ representation. Unwinding the resulting identification $`^2V_3\stackrel{}{}V_3^{}\stackrel{}{}V_3`$, we find:
Proposition A. Let $`Q_i=A_iX^2+2B_iXY+C_iY^2`$ ($`i=1,2`$) be two polynomials in $`V_3`$ without a common zero. Then the unique involution of $`𝐏^1`$ switching the roots of $`Q_1`$ and also of $`Q_2`$ is the involution whose fixed points are the roots of
$$(A_1B_2A_2B_1)X^2+(A_1C_2A_2C_1)XY+(B_1C_2B_2C_1)Y^2,$$
(88)
i.e. the fractional linear transformation
$$t\frac{(A_1C_2A_2C_1)t+2(B_1C_2B_2C_1)}{2(B_1A_2B_2A_1)t+(C_1A_2C_2A_1)}.$$
(89)
Proof: The coordinates of $`Q_1Q_2`$ for the basis of $`V_3^{}`$ dual to $`(X^2,2XY,Y^2)`$ are ($`B_1C_2B_2C_1`$, $`A_2C_1A_1C_2`$, $`A_1B_2A_2B_1`$). To identify $`V_3^{}`$ with $`V_3`$ we need a $`\mathrm{PSL}_2`$-invariant element of $`V_3^2`$. We could get this invariant from the invariant quadric $`B^2ACV_3^2`$, but it is easy enough to exhibit it directly: it is
$$X^2Y^2\frac{1}{2}\mathrm{\hspace{0.17em}2}XY2XY+Y^2X^2,$$
(90)
the generator of the kernel of the multiplication map Sym$`{}_{}{}^{2}(V_3)V_5`$. The resulting isomorphism from $`V_3^{}`$ to $`V_3`$ takes the dual basis of $`(X^2,2XY,Y^2)`$ to $`(Y^2,XY,X^2)`$, and thus takes $`Q_1Q_2`$ to (88) as claimed. $`\mathrm{}`$
Of course this is not the only way to obtain (89). A more “geometrical” approach (which ultimately amounts to the same thing) is to regard $`𝐏^1`$ as a conic in $`𝐏^2`$. Then involutions of $`𝐏^1`$ correspond to points $`p𝐏^2`$ not on the conic: the involution associated with $`p`$ takes any point $`q`$ of the conic to the second point of intersection of the line $`pq`$ with the conic. Of course the fixed points are then the points $`q`$ such that $`pq`$ is tangent to the conic at $`q`$. Given $`Q_1,Q_2`$ we obtain for $`i=1,2`$ the secant of the conic through the roots of $`Q_i`$, and then $`p`$ is the intersection of those secants.
From either of the two approaches we readily deduce
Corollary B. Let $`Q_i=A_iX^2+2B_iXY+C_iY^2`$ ($`i=1,2,3`$) be three polynomials in $`V_3`$ without a common zero. Then there is an involution of $`𝐏^1`$ switching the roots of $`Q_i`$ for each $`i`$ if and only if the determinant
$$\left|\begin{array}{ccc}A_1& B_1& C_1\\ A_2& B_2& C_2\\ A_3& B_3& C_3\end{array}\right|$$
(91)
vanishes.
As an additional check on the formula (88), we may compute that the discriminant of that quadratic polynomial is exactly the resolvent
$$det\left(\begin{array}{cccc}A_1& 2B_1& C_1& 0\\ 0& A_1& 2B_1& C_1\\ A_2& 2B_2& C_2& 0\\ 0& A_2& 2B_2& C_2\end{array}\right)$$
(92)
of $`Q_1,Q_2`$ which vanishes if and only if these two polynomials have a common zero.
Corrigendum
While I was preparing my paper (call it \[SCC\]) I did not have David Roberts’ thesis \[Ro\] to hand. Roberts has now kindly provided a copy of \[Ro\]; it turns out that the second-hand information I had to rely on concerning the contents of this thesis was wrong in several details, requiring specific corrections as follows.
I ask several times for a formula for the factorization of the differences between the coordinates of CM points on Shimura curves analogous to the formulas proved Gross and Zagier for X$`{}_{0}{}^{}(1)`$ \[GZ\] and obtained experimentally by Yui and Zagier for singular values of the Weber functions \[YZ\]. Roberts had already answered this question in principle, at least for the $`𝒳_0`$ curves, by obtaining an arithmetic intersection formula for CM points on these curves (section 6.5 of \[Ro\]). It would still take some work to extract from it (say) an explicit factorization of the difference between two such points on a Shimura curve of genus zero, but the technical framework now exists.
Roberts identified arithmetically most of the elliptic curves of conductor 60 or less which arise as Jacobians of Shimura curves for quaternion algebras over Q \[Ro, §7.4\]. This extends considerably the list of previously computed Shimura curves, and includes most that arise in \[SCC\]. Contrary to a statement in the first paragraph of \[SCC, §5.5\], the list in Roberts’ thesis does not, however, include the curve of conductor 66 which arises as the Jacobian of $`𝒳_0^{}(11)`$ \[ $`=X_{11,6}`$ in Roberts’ notation\], though it would be easy to obtain from his methods.
In \[SCC,§5.2\], the sentence preceding (76) needs some explanation: conductor “at most” 15 and 30 rather than exactly? In fact Jac($`𝒳(1)`$) is known a priori to have conductor exactly 15, but Jac($`𝒳_0(2)`$) has factors of conductor 15 \[namely Jac($`𝒳(1)`$), from “oldforms” on $`𝒳_0(2)`$\] as well as 30; if there were no newforms at all on $`𝒳_0(2)`$ its Jacobian would consist entirely of curves of conductor 15. \[Cf. the case of the classical modular curve X$`{}_{0}{}^{}(22)`$, whose Jacobian is isogenous with the square of X$`{}_{0}{}^{}(11)`$.\] For that matter, the fact that the conductors are thus bounded at all needs a reference. There are various ways of doing this; surely the easiest (given existing work) is to cite \[Ro\] for results finding (factors isogenous to) the Jacobians of Shimura curves inside the Jacobians of classical modular curves X$`{}_{0}{}^{}(N)`$, and using the known results about elliptic curves occurring in $`J_0(N)`$.
Thanks again to David Roberts for bringing these matters to my attention.
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# Vacuum Energy Density Fluctuations in Minkowski and Casimir States via Smeared Quantum Fields and Point Separation
## I Introduction
Recent years saw the beginning of serious studies of the fluctuations of the energy momentum tensor (EMT) $`\widehat{T}_{\mu \nu }`$ of quantum fields in spacetimes with boundaries (such as Casimir effect ) , nontrivial topology (such as imaginary time thermal field theory) or nonzero curvature (such as the Einstein universe) . A natural measure of the strength of fluctuations is $`\chi `$ , the ratio of the variance $`\mathrm{\Delta }\rho ^2`$ of fluctuations in the energy density (expectation value of the $`\widehat{\rho }^2`$ operator minus the square of the mean $`\widehat{\rho }`$ taken with respect to some quantum state) to its mean-squared (square of the expectation value of $`\widehat{\rho }`$):
$$\chi \frac{\widehat{\rho }^2\widehat{\rho }^2}{\widehat{\rho }^2}\frac{\mathrm{\Delta }\rho ^2}{\overline{\rho }^2}$$
(1)
Alternatively, we can use the quantity introduced by Kuo and Ford
$$\mathrm{\Delta }\frac{\widehat{\rho }^2\widehat{\rho }^2}{\widehat{\rho }^2}=\frac{\chi }{\chi +1}$$
(2)
Assuming a positive definite variance $`\mathrm{\Delta }\rho ^20`$, then $`0\chi \mathrm{}`$ and $`0\mathrm{\Delta }1`$ always, with $`\mathrm{\Delta }1`$ falling in the classical domain. Kuo and Ford (KF) displayed a number of quantum states (vacuum plus 2 particle state, squeezed vacuum and Casimir vacuum) with respect to which the expectation value of the energy momentum tensor (00 component) gives rise to negative local energy density. For these states $`\mathrm{\Delta }`$ is of order unity. From this result they drew the implications, amongst other interesting inferences, that semiclassical gravity (SCG) based on the semiclassical Einstein equation
$$G_{\mu \nu }=8\pi G\widehat{T}_{\mu \nu }$$
(3)
(where $`G_{\mu \nu }`$ is the Einstein tensor and $`G`$ the Newton gravitational constant) could become invalid under these conditions. The validity of semiclassical gravity in the face of fluctuations of quantum fields as source is an important issue which has caught the attention of many authors . We hold a different viewpoint on this issue from KF, which we hope to clarify in this study.
In this series of papers, we would like to examine more closely this and related issues of SCG, such as the regularization in the energy momentum tensor of quantum fields and fluctuations of spacetime metric. In this paper we study a free field in flat space and spacetimes with boundaries. Later papers deal with curved spacetimes depicting the early universe (quantum fluctuations and structure formation) and black holes (horizon fluctuations and dynamical backreactions).
As explained elsewhere by one of us, when fluctuations of the energy momentum tensors are included as source for the dynamics of spacetime, these problems are best discussed in the larger context of stochastic semiclassical gravity (SSG) program based on Einstein-Langevin type of equations , which is the proper framework to address Planck scale physics.
There are two groups of interrelated issues in quantum field theory in flat (ordinary QFT) or curved spacetimes (QFTCST), or semiclassical gravity (SCG –where the background spacetime dynamics is determined by the backreaction of the mean value of quantum fields): one pertaining to quantum fields and the other to spacetimes. We discuss the first set relating to the fluctuations of the EMT over its mean values with respect to the vacuum state. It strikes us as no great surprise that states which are more quantum (e.g., squeezed states) in nature than classical (e.g., coherent states) may lead to large fluctuations comparable to the mean in the energy density. This can be seen even in the ratio of expectation values of moments of the displacement operators in simple quantum harmonic oscillators . Such a condition exists peacefully with the underlying spacetime at least at the low energy of today’s universe. We don’t see sufficient ground to question the validity of SCG at energy below the Planck energy when the spacetime is depictable by a manifold structure, approximated locally by the Minkowski space. Besides, the cases studied in Kuo and Ford as well as many others are of a test-field nature, where backreaction is not considered. (So KF’s criterion pertains more to QFTCST than to SCG, where in the former the central issue is compatibility, which is a weaker condition than consistency in the latter.)
To assess this situation we aim at calculating the variance of fluctuations to mean-squared ratio of a quantum field for the simplest case of Minkowski spacetime i.e., for ordinary quantum field theory. We find that $`\mathrm{\Delta }=2/5`$. This is a clear-cut counter-example to the claim of KF, since $`\mathrm{\Delta }=O(1)`$ holds also for Minkowski space, where SCG is known to be valid at large scales. We view this situation as arising from the quantum nature of the vacuum state and saying little about the compatibility of the field source with the spacetime the quantum field lives in. In contrast, our view on this issue is that one should refer to a scale (of interaction or for probing accuracy) when measuring the validity of SCG. The conventional belief is that when reaching the Planck scale from below, QFTCST will break down because, amongst other things happening, graviton production at that energy will become significant so as to render the classical background spacetime unstable, and the mean value of quantum field taken as a source for the Einstein equation becomes inadequate. To address this issues as well as the issue of the spatial extent where negative energy density can exist, we view it necessary to introduce a scale in the spacetime regions where quantum fields are defined to monitor how the mean value and the fluctuations of the energy momentum tensor change. Point separation is an ideal method to adopt for this purpose.
In conventional field theories the stress tensor built from the product of a pair of field operators evaluated at a single point in the spacetime manifold is, strictly speaking, ill-defined. The point separation scheme was introduced as a method of regularization of the energy momentum tensor for quantum fields in curved spacetime. In this scheme, one introduces an artificial separation of the single point $`x`$ to a pair of closely separated points $`x`$ and $`x^{}`$. The problematic terms involving field products such as $`\widehat{\varphi }(x)^2`$ becomes $`\widehat{\varphi }(x)\widehat{\varphi }(x^{})`$, whose expectation value is well defined. One then brings the two points back (taking the coincidence limit) to identify the divergences present, which will then be removed (regularization) or moved (by renormalizing the coupling constants), thereby obtaining a well-defined, finite stress tensor at a single point. In this context point separation was introduced as a trick for identifying the ultraviolet divergences in a covariant manner. One of us in the development of the stochastic semiclassical gravity program has maintained the view (and as we will expound further in later papers ) that instead of being used as a mere technical device in QFTCST, this method has much greater physical content. We prefer to view the operator valued EM bi-tensor $`\widehat{T}_{\mu \nu }(x,y)`$ and EM two point function $`\widehat{T}_{\mu \nu }(x,y)`$ as the more fundamental objects in a more basic theory of spacetime and matter which has the point-defined quantum field theory as a low energy limit. This allows for spacetime to acquire an extended structure at sub-Planckian scale.<sup>*</sup><sup>*</sup>*Viewing the two points here as the end points of an open string gives what one of us called a ‘skeletal representation’ of string theory, where the internal excitation modes are suppressed. Related to the so-called dipole approximation of string theory in recent literature it should also carry features of non-commutative geometry of spacetime possible at the Planck scale.
For our stated purpose above, there is another way to introduce a scale in the quantum field theory, i.e., by introducing a (spatial) smearing function $`f(𝐱)`$ to define smeared field operators $`\widehat{\varphi }_t(f_𝐱)`$. In this paper we shall construct a scheme to encompass both aspects, by defining field operators at two separated points (connected by distance $`r`$) and using a Gaussian smearing function (with variance $`\sigma ^2`$). We derive expressions for the EM bi-tensor operator, its mean and its fluctuations as functions of $`r,\sigma `$, for a massless scalar field in both the Minkowski and the Casimir spacetimes. The interesting result we find is that while both the vacuum expectation value and the fluctuations of energy density grow as $`\sigma 0`$, the ratio of the variance of the fluctuations to its mean-squared remains a constant $`\chi _d`$ ($`d`$ is the spatial dimension of spacetime) which is independent of $`\sigma `$. The measure $`\mathrm{\Delta }_d`$ ($`=\chi _d/(\chi _d+1)`$) depends on the dimension of space and is of the order unity. It varies only slightly for spacetimes with boundary or nontrivial topology. For example $`\mathrm{\Delta }`$ for Minkowski is $`2/5`$, while for Casimir is $`6/7`$ (cf, from ). Add to this our prior result for the Einstein Universe, $`\mathrm{\Delta }=111/112`$, independent of curvature, and that for hot flat space , we see a pattern emerging.
These results allow us to address two interrelated issues: 1) Fluctuations of the energy density and validity of semiclassical gravity, and 2) The spatial extent where negative energy density can exist. For Issue 1) we see that i) the fluctuations of the energy density as well as its mean both increase with decreasing distance (or probing scale), while ii) the ratio of the variance of the fluctuations in EMT to its mean-squared is of the order unity. We view the first but not the second feature as linked to the question of the validity of SCG –the case for Minkowski spacetime alone is sufficient to testify to the fallacy of Kuo and Ford’s criterion. The second feature represents something quite different, pertaining more to the quantum nature of the vacuum state than to the validity of SCG.
For Issue 2) it is well known that negative energy density exists in Casimir geometry, moving mirrors, black holes and worm holes. Proposals have also been conjured to use the negative energy density for the design of time machines . Our results (Figures 1, 2) provide an explicit scale dependence of the regularized vacuum energy density $`\rho _{L,reg}`$ and its fluctuations $`\mathrm{\Delta }_{L,reg}`$ , specifically $`\sigma /L`$, the ratio of the smearing length (field scale) to that of the Casimir length (geometry scale). For example, Fig. 2 shows that only for $`\sigma /L<0.24`$ is $`\rho _{L,reg}<0`$. Recall $`\sigma `$ gives the spatial extent the field is probed or smeared. Ordinary pointwise quantum field theory which probes the field only at a point does not carry information about the spatial extent where negative energy density sustains. These results have direct implications on wormhole physics (and time machines, if one gets really serious about these fictions ). If L is the scale characterizing the size (‘throat’) of the wormhole where one thinks negative energy density might prevail, and designers of ‘time machines’ wish to exploit for ‘time-travel’ , our result provides a limit on the size of the probe (spaceship in the case of time-travel) in ratio to L where such conditions may exist. It could also provide a quantum field-theoretical bound on the probability of spontaneous creation of baby universes from quantum field energy fluctuations.
An equally weighty issue brought to light in this study is 3) the meaning of regularization in the face of EMT fluctuations. Since we have the point-separated expressions of the EMT and its fluctuations we can study how they change as a function of separation or smearing. In particular we can see how divergences arise at the coincidence limit. Whether certain cross terms containing divergences have physical meaning is a question raised by the recent studies of Wu and Ford . We can use these calculations to examine these issues and ask the broader question of what exactly regularization entails, where divergences arise and how they are to be treated. The consideration of divergences in the fluctuations of EMT requires a more sophisticated rationale and reveals a deeper layer of issues pertaining to effective versus more fundamental theories. If we view ordinary quantum field theory defined at points as a low energy limit of a theory of spacetime involving extended structures (such as string theory), then these results would shed light on their meaning and inter-connections.
In this paper we will discuss three aspects of quantum field theory in curved spacetime in the light of fluctuations of quantum stress energy : 1) Fluctuation to mean ratio of vacuum energy density and the validity of SCG; 2) The point separated results of the mean and the fluctuations of the energy density for two states: the Minkowski case which has no scale present (massless field) and the Casimir case which has a scale present (the separation of the plates); 3) The circumstances when and how divergences appear and the meaning of regularization in point-defined field theories versus theories defined at separated points and/or smeard fields. A summary of the main points of 1) has been given in . In this paper we give details of the calculations and discuss the regularization issue 3), while leaving the issue 2) of the spatial extent of negative energy density and its implications for quantum effects of worm holes, baby universes and time travel to a future investigation. In Sec. 2 we define the smeared field operators and their products defined at separated points. In Sec. 3 we construct from these the smeared energy density and its fluctuations, and calculate the ratio of the fluctuations to the mean for a flat space (Minkowski geometry). In Sec. 4 we analyze the case for a Casimir geometry of one periodic spatial dimension. In Sec. 5 we consider the point-separation calculation in Minkowski space, comparing the different results obtained from taking the coincidence limit from temporal versus spatial directions. Finally in Sec. 6 we summarize the major results and discuss the meaning of our finding in relation to the issues raised above.
## II Smeared Field Operators at Separated Points
Since the field operator in conventional point-defined quantum field theory is an operator-valued distribution, products of field operators at a point become problematic. This parallels the problem with defining the square of a delta function $`\delta ^2(x)`$. Distributions are defined via their integral against a test function: they live in the space dual to the test function space. By going from the field operator $`\widehat{\varphi }(x)`$ to its integral against a test function, $`\widehat{\varphi }(f)=\widehat{\varphi }f`$, we can now readily consider products.
When we take the test functions to be spatial Gaussians, we are smearing the field operator over a finite spatial region. Physically we see smearing as representing the necessarily finite extent of an observer’s probe, or the intrinsic limit of resolution in carrying out a measurement at a low energy (compared to Planck scale). In contrast to the ordinary point-defined quantum field theory, where ultraviolet divergences occur in the energy momentum tensor, smeared fields give no ultraviolet divergence. This is because smearing is equivalent to a regularization scheme which imparts an exponential suppression to the high momentum modes and restricts the contribution of the high frequency modes in the mode sum.
With this in mind, we start by defining the spatially smeared field operator
$$\widehat{\varphi }_t(f_𝐱)=\widehat{\varphi }(t,𝐱^{})f_𝐱(𝐱^{})𝑑𝐱^{}$$
(4)
where $`f_𝐱(𝐱^{})`$ is a suitably smooth function. With this, the two point operator becomes
$$\left(\widehat{\varphi }_t(f_𝐱)\right)^2=\widehat{\varphi }(t,𝐱^{})\widehat{\varphi }(t,𝐱^{\prime \prime })f_𝐱(𝐱^{})f_𝐱(𝐱^{\prime \prime })𝑑𝐱𝑑𝐱^{}$$
(5)
which is now finite. In terms of the vacuum $`|0`$ ($`\widehat{a}_𝐤|0=0`$, for all $`𝐤`$) we have the usual mode expansion
$$\widehat{\varphi }(t_1,𝐱_1)=𝑑\mu \left(𝐤_1\right)\left(\widehat{a}_{𝐤_1}u_{𝐤_1}(t_1,𝐱_1)+\widehat{a}_{𝐤_1}^{}u_{𝐤_1}^{}(t_1,𝐱_1)\right)$$
(6)
with
$$u_{𝐤_1}(t_1,𝐱_1)=N_{k_1}e^{i\left(𝐤_1𝐱_1t_1\omega _1\right)},\omega _1=\left|𝐤_1\right|,$$
(7)
where the integration measure $`𝑑\mu \left(𝐤_1\right)`$ and the normalization constants $`N_{𝐤_1}`$ are given for a Minkowski and Casimir spaces by (17) and (IV) respectively.
In this work, we use a Gaussian smearing function
$$f_{𝐱_0}(𝐱)=\left(\frac{1}{4\pi \sigma ^2}\right)^{\frac{d}{2}}e^{\left(\frac{𝐱_0𝐱_1}{2\sigma }\right)^2}$$
(8)
with the properties $`f_{𝐱_0}\left(𝐱^{}\right)𝑑𝐱^{}=1`$, $`𝐱^{}f_{𝐱_0}\left(𝐱^{}\right)𝑑𝐱^{}=𝐱_0`$ and $`|𝐱^{}|^2f_{𝐱_0}\left(𝐱^{}\right)𝑑𝐱^{}=2d\sigma ^2+|𝐱_0|^2`$. Using
$`{\displaystyle u_{𝐤_1}(t,𝐱)f_{𝐱_1}(𝐱)𝑑𝐱}`$ $`=`$ $`N_{k_1}e^{it\omega _1}{\displaystyle \underset{i=1}{\overset{d}{}}}\left({\displaystyle \frac{1}{2\sqrt{\pi }\sigma }}{\displaystyle e^{+ik_{1i}x_i\left(\frac{x_{1i}x_i}{2\sigma }\right)^2}𝑑x_i}\right)`$ (9)
$`=`$ $`N_{k_1}e^{itw+i𝐤_1𝐱_1\sigma ^2k_{1}^{}{}_{}{}^{2}}`$ (10)
we get the smeared field operator
$$\widehat{\varphi }_{t_1}\left(f_{𝐱_1}\right)=𝑑\mu \left(𝐤_1\right)N_{k_1}e^{i𝐤_1𝐱_1\sigma ^2k_{1}^{}{}_{}{}^{2}it_1\omega _1}\left(e^{2i𝐤_1𝐱_1}\widehat{a}_{𝐤_1}+e^{2it_1\omega _1}\widehat{a}_{𝐤_1}^{}\right)$$
(11)
and their product:
$`\widehat{\varphi }_{t_1}\left(f_{𝐱_1}\right)\widehat{\varphi }_{t_2}\left(f_{𝐱_2}\right)`$ $`=`$ $`{\displaystyle 𝑑\mu (𝐤_1,𝐤_2)N_{k_1}N_{k_2}e^{\sigma ^2\left(k_{1}^{}{}_{}{}^{2}+k_{2}^{}{}_{}{}^{2}\right)i\left(𝐤_1𝐱_1+𝐤_2𝐱_2+t_1\omega _1+t_2\omega _2\right)}}`$ (14)
$`(e^{2i\left(𝐤_1𝐱_1+𝐤_2𝐱_2\right)}\widehat{a}_{𝐤_1}\widehat{a}_{𝐤_2}+e^{2i\left(𝐤_1𝐱_1+t_2\omega _2\right)}\widehat{a}_{𝐤_1}\widehat{a}_{𝐤_2}^{}+e^{2i\left(𝐤_2𝐱_2+t_1\omega _1\right)}\widehat{a}_{𝐤_1}^{}\widehat{a}_{𝐤_2}`$
$`+e^{2i\left(t_1\omega _1+t_2\omega _2\right)}\widehat{a}_{𝐤_1}^{}\widehat{a}_{𝐤_2}^{}).`$
The smearing has introduced a factor $`e^{\sigma ^2k^2}`$ for each of the momenta. This factor acts as the regulator for the mode sums; in this sense the high momenta that led to the divergences have been controlled. We also note the $`\sigma 0`$ limit amounts to the relaxation of the regulator and the point field theory version of the field operator is recovered.
### A Coincident Smeared Green Function
We can let the two spatial points come together, $`𝐱_2=𝐱_1=0`$, and get the coincident limit of the Green function
$`G\left(\mathrm{\Delta }t\right)`$ $`=`$ $`0\left|\widehat{\varphi }_{t_1}\left(f_{𝐱_1}\right)\widehat{\varphi }_{t_2}\left(f_{𝐱_1}\right)\right|0`$ (15)
$`=`$ $`{\displaystyle 𝑑\mu \left(𝐤_1\right)N_{k_1}^2e^{2\sigma ^2k_1^2+i\mathrm{\Delta }t\omega _1}}`$ (16)
For flat space the normalization and integration measure are
$$N_{k_1}^2=\frac{1}{2^{d+1}\pi ^d\omega _1}\mathrm{and}𝑑\mu (𝐤_1)=\frac{2\pi ^{\frac{d}{2}}}{\mathrm{\Gamma }\left(\frac{d}{2}\right)}_0^{\mathrm{}}k_1^{d1}𝑑k_1$$
(17)
and the Green function is
$`G(\mathrm{\Delta }t)`$ $`=`$ $`{\displaystyle \frac{1}{2^{\frac{3}{2}\left(d+1\right)}\pi ^{\frac{d}{2}}\sigma ^d\mathrm{\Gamma }\left(\frac{d}{2}\right)}}\{2\sigma \mathrm{\Gamma }\left({\displaystyle \frac{d1}{2}}\right){}_{1}{}^{}F_{1}^{}({\displaystyle \frac{d1}{2}};{\displaystyle \frac{1}{2}};{\displaystyle \frac{\mathrm{\Delta }t^2}{8\sigma ^2}})`$ (20)
$`+i\sqrt{2}\mathrm{\Delta }t\mathrm{\Gamma }\left({\displaystyle \frac{d}{2}}\right){}_{1}{}^{}F_{1}^{}({\displaystyle \frac{d}{2}};{\displaystyle \frac{3}{2}};{\displaystyle \frac{\mathrm{\Delta }t^2}{8\sigma ^2}})\}`$
$`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{d1}{2}\right)}{2^{\frac{3d+1}{2}}\pi ^{\frac{d}{2}}\sigma ^{d1}\mathrm{\Gamma }\left(\frac{d}{2}\right)}}\left(1+{\displaystyle \frac{i\mathrm{\Delta }t\mathrm{\Gamma }\left(\frac{d}{2}\right)}{\sqrt{2}\sigma \mathrm{\Gamma }\left(\frac{d1}{2}\right)}}{\displaystyle \frac{\left(d1\right)\mathrm{\Delta }t^2}{8\sigma ^2}}\right)+O\left(\mathrm{\Delta }t^3\right),`$ (22)
which we see is finite for this spatial coincident limit, and is finite for the $`\mathrm{\Delta }t0`$ limit as well. For spatial dimension $`d=3`$, the smeared Green function is
$`G(\mathrm{\Delta }t)`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2\sigma ^2}}{\displaystyle \frac{\mathrm{\Delta }t\left(\mathrm{Erfi}(\frac{\mathrm{\Delta }t}{2\sqrt{2}\sigma })i\right)}{32\sqrt{2}e^{\frac{\mathrm{\Delta }t^2}{8\sigma ^2}}\pi ^{\frac{3}{2}}\sigma ^3}}`$ (24)
$`=`$ $`{\displaystyle \frac{1}{16\pi ^2\sigma ^2}}\left(1+{\displaystyle \frac{i\mathrm{\Delta }t\sqrt{\pi }}{2\sqrt{2}\sigma }}{\displaystyle \frac{\mathrm{\Delta }t^2}{4\sigma ^2}}\right)+O\left(\mathrm{\Delta }t^3\right)`$ (26)
### B Point-Separated Smeared Energy Density Operator
For a classical scalar function, the energy density is
$$\rho (t_1,𝐱_1)=\frac{1}{2}\left(\left(_{t_1}\varphi \right)^2+\left(\stackrel{}{}\varphi \right)^2\right)$$
(27)
We cannot go directly to the quantum field case since the energy density has pairs of field operators evaluated at the same point. We can however define an energy density operator which contains smeared field operators at separated points. Then by taking the coincident ($`𝐱_1𝐱_2`$) limit, we will be using the smearing to regularize the energy density, while if we take the zero smearing width ($`\sigma 0`$) limit, we are using point separation regularization.
Point separation consists of symmetrically splitting the operator product as
$$\widehat{\varphi }(t_1,𝐱_1)^2\frac{1}{2}\left(\widehat{\varphi }(t_1,𝐱_1)\widehat{\varphi }(t_2,𝐱_2)+\widehat{\varphi }(t_2,𝐱_2)\widehat{\varphi }(t_1,𝐱_1)\right).$$
(28)
Products of derivatives of the field operator are symmetrically split according to, e.g., time derivatives become
$$\left(_{t_1}\widehat{\varphi }(t_1,𝐱_1)\right)^2\frac{1}{2}\left(\left(_{t_1}\widehat{\varphi }(t_1,𝐱_1)\right)\left(_{t_2}\widehat{\varphi }(t_2,𝐱_2)\right)+\left(_{t_2}\widehat{\varphi }(t_2,𝐱_2)\right)\left(_{t_1}\widehat{\varphi }(t_1,𝐱_1)\right)\right)$$
(29)
We introduce the smeared field operator derivatives
$`\left(_{t_1}\widehat{\varphi }_{t_1}\right)\left(f_{𝐱_1}\right)`$ $`=`$ $`{\displaystyle \left(_{t_1}\widehat{\varphi }(t_1,𝐱^{})\right)f_{𝐱_1}\left(𝐱^{}\right)𝑑𝐱^{}}`$ (31)
$`=`$ $`i{\displaystyle 𝑑\mu (𝐤_1)N_{k_1}\omega _1e^{i𝐤_1𝐱_1\sigma ^2k_{1}^{}{}_{}{}^{2}it_1\omega _1}\left(e^{2it_1\omega _1}\widehat{a}_{𝐤_1}^{}e^{2i𝐤_1𝐱_1}\widehat{a}_{𝐤_1}\right)}`$ (32)
$`\left(\stackrel{}{}_{𝐱_1}\widehat{\varphi }_{t_1}\right)\left(f_{𝐱_1}\right)`$ $`=`$ $`{\displaystyle \left(\stackrel{}{}_𝐱^{}\widehat{\varphi }(t_1,𝐱^{})\right)f_{𝐱_1}\left(𝐱^{}\right)𝑑𝐱^{}}`$ (34)
$`=`$ $`i{\displaystyle 𝑑\mu (𝐤_1)𝐤_1N_{k_1}e^{i𝐤_1𝐱_1\sigma ^2k_{1}^{}{}_{}{}^{2}it_1\omega _1}\left(e^{2it_1\omega _1}\widehat{a}_{𝐤_1}^{}e^{2i𝐤_1𝐱_1}\widehat{a}_{𝐤_1}\right)}`$ (35)
and use the symmetric splitting to define the point separated smeared energy density operator
$`\widehat{\rho }(t_1,𝐱_1;t_2,𝐱_2;\sigma )`$ $`=`$ $`{\displaystyle \frac{1}{4}}\{\left(\left(_{t_1}\widehat{\varphi }_{t_1}\right)\left(f_{𝐱_1}\right)\right)\left(\left(_{t_2}\widehat{\varphi }_{t_2}\right)\left(f_{𝐱_2}\right)\right)+\left(\left(_{t_2}\widehat{\varphi }_{t_2}\right)\left(f_{𝐱_2}\right)\right)\left(\left(_{t_1}\widehat{\varphi }_{t_1}\right)\left(f_{𝐱_1}\right)\right)`$ (37)
$`+\left(\left(\stackrel{}{}_{𝐱_1}\widehat{\varphi }_{t_1}\right)\left(f_{𝐱_1}\right)\right)\left(\left(\stackrel{}{}_{𝐱_2}\widehat{\varphi }_{t_2}\right)\left(f_{𝐱_2}\right)\right)+\left(\left(\stackrel{}{}_{𝐱_2}\widehat{\varphi }_{t_2}\right)\left(f_{𝐱_2}\right)\right)\left(\left(\stackrel{}{}_{𝐱_1}\widehat{\varphi }_{t_1}\right)\left(f_{𝐱_1}\right)\right)\}`$
$`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle 𝑑\mu (𝐤_1)𝑑\mu (𝐤_2)N_{k_1}N_{k_2}\left(𝐤_1𝐤_2+\omega _1\omega _2\right)}`$ (43)
$`\times e^{i\left(𝐤_1𝐱_1+𝐤_2𝐱_2\right)i\left(t_1\omega _1+t_2\omega _2\right)\sigma ^2\left(k_{1}^{}{}_{}{}^{2}+k_{2}^{}{}_{}{}^{2}\right)}`$
$`\times (e^{2i𝐤_1𝐱_1+2i𝐤_2𝐱_2}\widehat{a}_{𝐤_1}\widehat{a}_{𝐤_2}e^{2i𝐤_1𝐱_1+2it_2\omega _2}\widehat{a}_{𝐤_1}\widehat{a}_{𝐤_2}^{}`$
$`+e^{2i𝐤_1𝐱_1+2i𝐤_2𝐱_2}\widehat{a}_{𝐤_2}\widehat{a}_{𝐤_1}e^{2i𝐤_2𝐱_2+2it_1\omega _1}\widehat{a}_{𝐤_2}\widehat{a}_{𝐤_1}^{}`$
$`e^{2i𝐤_2𝐱_2+2it_1\omega _1}\widehat{a}_{𝐤_1}^{}\widehat{a}_{𝐤_2}+e^{2it_1\omega _1+2it_2\omega _2}\widehat{a}_{𝐤_1}^{}\widehat{a}_{𝐤_2}^{}`$
$`e^{2i𝐤_1𝐱_1+2it_2\omega _2}\widehat{a}_{𝐤_2}^{}\widehat{a}_{𝐤_1}+e^{2it_1\omega _1+2it_2\omega _2}\widehat{a}_{𝐤_2}^{}\widehat{a}_{𝐤_1}^{})`$
Its vacuum expectation value is
$`\rho (t_1,𝐱_1;t_2,𝐱_2;\sigma )`$ $`=`$ $`0\left|\widehat{\rho }_t(t_1,𝐱_1;t_2,𝐱_2;\sigma )\right|0`$ (44)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle 𝑑\mu \left(𝐤_1\right)N_{k_1}^2\omega _1^2e^{2\sigma ^2k_1^2}}`$ (46)
$`\times \left(e^{i\left(𝐤_1\left(𝐱_1𝐱_2\right)\left(t_1t_2\right)\omega _1\right)}+e^{i\left(𝐤_1\left(𝐱_1𝐱_2\right)\left(t_1t_2\right)\omega _1\right)}\right)`$
We obtain the smeared vacuum energy density by taking the coincidence
$$\rho \left(\sigma \right)=𝑑\mu \left(𝐤_1\right)N_{k_1}^2\omega _1^2e^{2\sigma ^2k_1^2}$$
(47)
while the point separation expression is obtained from taking the zero smearing-width limit:
$$\rho (t_1,𝐱_1;t_2,𝐱_2)=𝑑\mu (𝐤_1)N_{k_1}^2\omega _1^2\mathrm{cos}(𝐤_1\left(𝐱_1𝐱_2\right)\left(t_1t_2\right)\omega _1)$$
(48)
### C Point-Separated Smeared Energy Density Correlation Function
We now consider the point separated vacuum correlation function for the energy density operator:
$`\mathrm{\Delta }\rho ^2(t_1,𝐱_1,t_1^{},𝐱_1^{};t_2,𝐱_2,t_2^{},𝐱_2^{};\sigma )`$ $`=`$ $`0\left|\widehat{\rho }(t_1,𝐱_1;t_1^{},𝐱_1^{};\sigma )\widehat{\rho }(t_2,𝐱_2;t_2^{},𝐱_2^{};\sigma )\right|0`$ (50)
$`\rho (t_1,𝐱_1;t_1^{},𝐱_1^{};\sigma )\rho (t_2,𝐱_2;t_2^{},𝐱_2^{};\sigma )`$
With this definition, the vacuum energy density correlation function is
$`\mathrm{\Delta }\rho ^2(t_1,𝐱_1;t_2,𝐱_2)`$ $``$ $`0\left|\widehat{\rho }(t_1,𝐱_1)\widehat{\rho }(t_2,𝐱_2)\right|00\left|\widehat{\rho }(t_1,𝐱_1)\right|00\left|\widehat{\rho }(t_2,𝐱_2)\right|0`$ (51)
$`=`$ $`\mathrm{\Delta }\rho ^2\left(t_1,𝐱_1,t_1,𝐱_1;t_2,𝐱_2,t_2,𝐱_2;\sigma =0\right)`$ (52)
Since the divergences present in $`0\left|\widehat{\rho }(t_1,𝐱_1)\widehat{\rho }(t_2,𝐱_2)\right|0`$ for $`(t_2,𝐱_2)(t_1,𝐱_1)`$ are canceled by those due to $`0\left|\widehat{\rho }(t_1,𝐱_1)\right|0`$ and $`0\left|\widehat{\rho }(t_2,𝐱_2)\right|0`$, we can assume $`(t_1^{},𝐱_1^{})=(t_1,𝐱_1)`$ and $`(t_2^{},𝐱_2^{})=(t_2,𝐱_2)`$ from the start. This will be confirmed during the computation of the vacuum expectation value.
First we consider just the square of the energy density operator; its expectation value is
$`0\left|\widehat{\rho }^2\right|0`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle }d\mu (𝐤_1,𝐤_2)N_{k_1}^2N_{k_2}^2e^{2\sigma ^2\left(k_1^2+k_2^2\right)}\{`$ (57)
$`(𝐤_1𝐤_2+\omega _1\omega _2)^2(e^{i\left(𝐤_1\left(𝐱_1𝐱_2^{}\right)+𝐤_2\left(𝐱_1^{}𝐱_2\right)\right)i\left(\left(t_1t_2^{}\right)\omega _1+\left(t_1^{}t_2\right)\omega _2\right)}`$
$`+e^{i\left(𝐤_1\left(𝐱_1𝐱_2\right)+𝐤_2\left(𝐱_1^{}𝐱_2^{}\right)\right)i\left(\left(t_1t_2\right)\omega _1+\left(t_1^{}t_2^{}\right)\omega _2\right)})`$
$`+[\omega _1^2(e^{i\left(𝐤_1\left(𝐱_1𝐱_2\right)\left(t_1t_1^{}\right)\omega _1\right)}+e^{i\left(𝐤_1\left(𝐱_1𝐱_1^{}\right)\left(t_1t_1^{}\right)\omega _1\right)})`$
$`\times \omega _2^2(e^{i\left(𝐤_2\left(𝐱_2𝐱_4\right)\left(t_2t_2^{}\right)\omega _2\right)}+e^{i\left(𝐤_2\left(𝐱_2𝐱_2^{}\right)\left(t_2t_2^{}\right)\omega _2\right)})]\}`$
By comparing the last two lines of the above expression with Eq.(46), we see this is but $`\rho (t_1,𝐱_1;t_1^{},𝐱_1^{};\sigma )\rho (t_2,𝐱_2;t_2^{},𝐱_2^{};\sigma )`$. Thus, the remainder is the desired expression for $`\mathrm{\Delta }\rho ^2(t_1,𝐱_1,t_1^{},𝐱_1^{};t_2,𝐱_2,t_2^{},𝐱_2^{};\sigma )`$. Even for the case $`\sigma 0`$, this expression is finite for $`(t_1^{},𝐱_1^{})(t_1,𝐱_1)`$ and $`(t_2^{},𝐱_2^{})(t_2,𝐱_2)`$, as long as $`(t_1,𝐱_1)(t_2,𝐱_2)`$. Letting $`(t,𝐱)=(t_2,𝐱_2)(t_1,𝐱_1)`$, our results for the energy density \[from (46)\] and its correlation function \[from here\] are
$$\rho (t,𝐱;\sigma )=𝑑\mu \left(𝐤\right)N_k^2\omega ^2e^{2\sigma ^2k^2}\mathrm{cos}(𝐱𝐤t\omega )$$
(58)
$$\mathrm{\Delta }\rho ^2(t,𝐱;\sigma )=\frac{1}{2}𝑑\mu (𝐤_1,𝐤_2)N_{k_1}^2N_{k_2}^2\left(𝐤_1𝐤_2+\omega _1\omega _2\right)^2e^{2\sigma ^2\left(k_1^2+k_2^2\right)i𝐱\left(𝐤_1+𝐤_2\right)+it\left(\omega _1+\omega _2\right)}$$
(59)
## III Smeared-Field Energy Density and Fluctuations in Minkowski Space
We consider a Minkowski space $`R^1\times R^d`$ with $`d`$-spatial dimensions. For this space the mode density is
$$𝑑\mu (𝐤)=_0^{\mathrm{}}k^{d1}𝑑k_{S^{d1}}𝑑\mathrm{\Omega }_{d1}\mathrm{with}_{S^{d1}}𝑑\mathrm{\Omega }_{d1}=\frac{2\pi ^{\frac{d}{2}}}{\mathrm{\Gamma }\left(\frac{d}{2}\right)}$$
(60)
and the mode function normalization constant is $`N_{k_1}=1/\sqrt{2^{d+1}\pi ^d\omega _1}`$. We introduce the angle between two momenta in phase space, $`\gamma `$, via
$$𝐤_1𝐤_2=k_1k_2\mathrm{cos}(\gamma )=\omega _1\omega _2\mathrm{cos}(\gamma ).$$
(61)
The averages of the cosine and cosine squared of this angle over a pair of unit spheres are
$`{\displaystyle _{S^{d1}}}𝑑\mathrm{\Omega }_1{\displaystyle _{S^{d1}}}𝑑\mathrm{\Omega }_2\mathrm{cos}(\gamma )`$ $`=`$ $`0`$ (63)
$`{\displaystyle _{S^{d1}}}𝑑\mathrm{\Omega }_1{\displaystyle _{S^{d1}}}𝑑\mathrm{\Omega }_2\mathrm{cos}^2(\gamma )`$ $`=`$ $`{\displaystyle \frac{4\pi ^d}{d\mathrm{\Gamma }\left(\frac{d}{2}\right)^2}}.`$ (64)
The smeared energy density (47) becomes
$`\rho (\sigma )`$ $`=`$ $`{\displaystyle \frac{1}{2^d\pi ^{\frac{d}{2}}\mathrm{\Gamma }\left(\frac{d}{2}\right)}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{k_{1}^{}{}_{}{}^{d}}{e^{2\sigma ^2k_{1}^{}{}_{}{}^{2}}}}𝑑k_1`$ (65)
$`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{d+1}{2}\right)}{2^{\frac{3\left(d+1\right)}{2}}\pi ^{\frac{d}{2}}\sigma ^{d+1}\mathrm{\Gamma }\left(\frac{d}{2}\right)}}`$ (66)
For the fluctuations of the smeared energy density operator, we evaluate (59) for this space and find
$`\mathrm{\Delta }\rho ^2(\sigma )`$ $`=`$ $`{\displaystyle \frac{1}{2^{(2d+3)}\pi ^{2d}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle _{S^{d1}}}{\displaystyle _{S^{d1}}}{\displaystyle \frac{\left(1+\mathrm{cos}(\gamma )\right)^2k_1^dk_2^d}{e^{2\sigma ^2\left(k_1^2+k_2^2\right)}}}𝑑\mathrm{\Omega }_1𝑑\mathrm{\Omega }_2𝑑k_1𝑑k_2`$ (67)
$`=`$ $`{\displaystyle \frac{\left(d+1\right)\mathrm{\Gamma }\left(\frac{d+1}{2}\right)^2}{2^{(3d+4)}d\pi ^d\sigma ^{2\left(d+1\right)}\mathrm{\Gamma }\left(\frac{d}{2}\right)^2}}`$ (68)
Defining the (dimension-dependent) constant
$$\chi _d\frac{1+d}{2d},$$
(69)
we write the smeared fluctuations in terms of the square of the smeared energy density
$$\mathrm{\Delta }\rho ^2(\sigma )=\chi _d\rho (\sigma )^2$$
(70)
We introduce the dimensionless measure of fluctuations
$$\mathrm{\Delta }=\left|1\frac{\rho ^2}{\rho ^2}\right|=\left|\frac{\mathrm{\Delta }\rho ^2}{\mathrm{\Delta }\rho ^2+\rho ^2}\right|=\frac{\chi _d}{\chi _d+1}$$
(71)
and for Minkowski space we have
$$\mathrm{\Delta }_{\mathrm{Minkowski}}(d)=\frac{1+d}{1+3d}$$
(72)
which has the particular values
$`\begin{array}{ccccc}& & & & \\ d& 1& 3& 5& \mathrm{}\\ & & & & \\ \mathrm{\Delta }_{\mathrm{Minkowski}}& \frac{1}{2}& \frac{2}{5}& \frac{3}{8}& \frac{1}{3}\end{array}`$
## IV Smeared-Field in Casimir Topology
The Casimir topology is obtained from a flat space (with $`d`$ spatial dimensions, i.e., $`R^1\times R^d`$ ) by imposing periodicity $`L`$ in one of its spatial dimensions, say, $`z`$, thus endowing it with a $`R^1\times R^{d1}\times S^1`$ topology. We decompose $`𝐤`$ into a component along the periodic dimension and call the remaining components $`𝐤_{}`$:
$`𝐤`$ $`=`$ $`(𝐤_{},{\displaystyle \frac{2\pi n}{L}})=(𝐤_{},ln),l2\pi /L`$ (74)
$`\omega _1`$ $`=`$ $`\sqrt{k_1^2+l^2n_{1}^{}{}_{}{}^{2}}`$ (75)
The normalization and momentum measure are
$`{\displaystyle 𝑑\mu (𝐤)}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}k^{d2}𝑑k{\displaystyle _{S^{d2}}}𝑑\mathrm{\Omega }_{d2}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}`$ (76)
$`N_{k_1}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2^dL\pi ^{d1}\omega _1}}}`$ (77)
With this, the energy density (47) becomes
$`\rho _L\left(\sigma \right)`$ $`=`$ $`{\displaystyle \frac{l}{2^d\pi ^{\frac{d+1}{2}}\mathrm{\Gamma }\left(\frac{d1}{2}\right)}}{\displaystyle \underset{n_1=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle _0^{\mathrm{}}}k_1^{d2}\left(k_1^2+l^2n_1^2\right)^{\frac{1}{2}}e^{2\sigma ^2\left(k_1^2+l^2n_1^2\right)}𝑑k_1`$ (78)
$`=`$ $`{\displaystyle \frac{l}{2^d\pi ^{\frac{d+1}{2}}\mathrm{\Gamma }\left(\frac{d1}{2}\right)}}{\displaystyle \underset{n_1=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{k_{1}^{}{}_{}{}^{d}}{\sqrt{k_1^2+l^2n_1^2}}}e^{2\sigma ^2\left(k_1^2+l^2n_1^2\right)}𝑑k_1`$ (80)
$`+{\displaystyle \frac{l}{2^d\pi ^{\frac{d+1}{2}}\mathrm{\Gamma }\left(\frac{d1}{2}\right)}}{\displaystyle \underset{n_1=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{l^2k_{1}^{}{}_{}{}^{2+d}n_1^2}{\sqrt{k_1^2+l^2n_1^2}}}e^{2\sigma ^2\left(k_1^2+l^2n_1^2\right)}𝑑k_1`$
Using the analysis of the Appendix A we write this as the sum of the two smeared Green function derivatives
$`\rho _L\left(\sigma \right)`$ $`=`$ $`0_L\left|\left(\left(_{}\varphi _t\right)\left(f_𝐱\right)\right)^2\right|0_L+0_L\left|\left(\left(_z\varphi _t\right)\left(f_𝐱\right)\right)^2\right|0_L`$ (81)
$`=`$ $`G_L(\sigma )_{,x_{}x_{}}+G_L(\sigma )_{,zz}`$ (82)
where $`|0_L`$ is the Casimir vacuum.
### A Regularized Casimir Energy Density
Since $`G_L(\sigma )_{,i}=G_{L,i}^{\mathrm{div}}+G_{L,i}^{\mathrm{fin}}`$ ($`i=x_{}x_{}`$ or $`zz`$) we see how to split the smeared energy density into a $`\sigma 0`$ divergent term and the finite contribution:
$$\rho _L\left(\sigma \right)=\rho _L^{\mathrm{div}}+\rho _L^{\mathrm{fin}}$$
(83)
where
$`\rho _L^{\mathrm{div}}`$ $`=`$ $`G_{L,x_{}x_{}}^{\mathrm{div}}+G_{L,zz}^{\mathrm{div}}`$ (84)
$`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{d+1}{2}\right)}{2^{\frac{3\left(d+1\right)}{2}}\pi ^{\frac{d}{2}}\sigma ^{d+1}\mathrm{\Gamma }\left(\frac{d}{2}\right)}}`$ (85)
$`=`$ $`\rho \left(\sigma \right)`$ (86)
and
$`\rho _L^{\mathrm{fin}}`$ $`=`$ $`G_{L,x_{}x_{}}^{\mathrm{fin}}+G_{L,zz}^{\mathrm{fin}}`$ (87)
$`=`$ $`{\displaystyle \frac{d\mathrm{\Gamma }\left(\frac{d}{2}\right)\mathrm{\Gamma }\left(\frac{d}{2}\right)}{(4\pi )^{(d+3)/2}l^{d+1}}}{\displaystyle \underset{p=1}{\overset{\mathrm{}}{}}}\left(1\right)^p(2l)^{2p}p\left(2p1\right)^2\sigma ^{2\left(p1\right)}{\displaystyle \frac{B_{2p+d1}}{2p+d1}}{\displaystyle \frac{\left(2p3\right)!!}{\left(2p\right)!}}{\displaystyle \frac{\mathrm{\Gamma }\left(p\frac{1}{2}\right)}{\mathrm{\Gamma }\left(p+\frac{d}{2}\right)}}`$ (88)
With this we define the regularized energy density
$`\rho _{L,\mathrm{reg}}`$ $``$ $`\underset{\sigma 0}{lim}\left(\rho _L\left(\sigma \right)\rho \left(\sigma \right)\right)`$ (89)
$`=`$ $`{\displaystyle \frac{d\pi ^{\frac{d}{2}}B_{d+1}\mathrm{\Gamma }\left(\frac{d}{2}\right)\mathrm{\Gamma }\left(\frac{d}{2}\right)}{2\left(d+1\right)L^{d+1}\mathrm{\Gamma }\left(\frac{d}{2}+1\right)}}`$ (90)
and get the usual results
$`\begin{array}{cccc}& & & \\ d& 1& 3& 5\\ & & & \\ \rho _{L,\mathrm{reg}}& \frac{\pi }{6L^2}& \frac{\pi ^2}{90L^4}& \frac{2\pi ^3}{945L^6}\end{array}`$
### B Casimir energy density fluctuations
For the $`d`$-dimensional Casimir geometry, the fluctuations are
$`\mathrm{\Delta }\rho _L^2\left(\sigma \right)`$ $`=`$ $`{\displaystyle \frac{l^2}{2^{2d+3}\pi ^{2d}}}{\displaystyle \underset{n_1=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n_2=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle _0^{\mathrm{}}}k_1^{d2}𝑑k_1{\displaystyle _0^{\mathrm{}}}k_2^{d2}𝑑k_2{\displaystyle _{S^{d2}}}𝑑\mathrm{\Omega }_1{\displaystyle _{S^{d2}}}𝑑\mathrm{\Omega }_2`$ (92)
$`\times {\displaystyle \frac{e^{2\sigma ^2\left(\omega _1^2+\omega _2^2\right)}}{\omega _1\omega _2}}\left(\mathrm{cos}(\gamma )k_1k_2+l^2n_1n_2+\omega _1\omega _2\right)^2`$
$`=`$ $`{\displaystyle \frac{l^2}{2^{2d}\pi ^{d+1}\mathrm{\Gamma }\left(\frac{d1}{2}\right)^2}}{\displaystyle \underset{n_1,n_2=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{\mathrm{}}}𝑑k_1𝑑k_2{\displaystyle \frac{k_1^{d+2}e^{2\sigma ^2\left(k_1^2+l^2n_1^2\right)}}{\sqrt{k_1^2+l^2n_1^2}}}{\displaystyle \frac{k_2^{d+2}e^{2\sigma ^2\left(k_2^2+l^2n_2^2\right)}}{\sqrt{k_2^2+l^2n_2^2}}}`$ (94)
$`\times \left({\displaystyle \frac{dk_1^2k_2^2}{2\left(d1\right)}}+{\displaystyle \frac{l^2}{2}}\left(k_2^2n_1^2+k_1^2n_2^2\right)+l^4n_1^2n_2^2\right)`$
We write this expression in terms of products of the Green functions derivatives used above:
$$\mathrm{\Delta }\rho _L^2\left(\sigma \right)=\frac{d\left(G_L(\sigma )_{,x_{}x_{}}\right)^2}{2\left(d1\right)}+G_L(\sigma )_{,zz}\left(G_L(\sigma )_{,x_{}x_{}}+G_L(\sigma )_{,zz}\right)$$
(95)
We can split $`\mathrm{\Delta }\rho _L^2\left(\sigma \right)`$ into three general terms
$$\mathrm{\Delta }\rho _L^2\left(\sigma \right)=\mathrm{\Delta }\rho _L^{2,\mathrm{div}}+\mathrm{\Delta }\rho _L^{2,\mathrm{cross}}+\mathrm{\Delta }\rho _L^{2,\mathrm{fin}}$$
(96)
The first term contains only the divergent parts of the Green functions while the last term contains only the finite parts. This is similar to the split we used for the smeared energy density above. What is new here is the middle term $`\mathrm{\Delta }\rho _L^{2,\mathrm{cross}}`$. This comes about from the products of the divergent part of one Green function and the finite part of the other. That this term arises for computations of the energy density fluctuations is a generic feature. We will discuss in greater detail the meaning of this term later.
The results of Appendix A give
$`\mathrm{\Delta }\rho _L^{2,\mathrm{div}}`$ $`=`$ $`{\displaystyle \frac{d\left(G_{L,x_{}x_{}}^{\mathrm{div}}\right)^2}{2\left(d1\right)}}+G_{L,zz}^{\mathrm{div}}\left(G_{L,x_{}x_{}}^{\mathrm{div}}+H_2^{\mathrm{div}}\right)`$ (98)
$`=`$ $`{\displaystyle \frac{\left(d+1\right)\mathrm{\Gamma }\left(\frac{d+1}{2}\right)^2}{d\mathrm{\hspace{0.17em}2}^{3d+4}\pi ^d\sigma ^{2\left(d+1\right)}\mathrm{\Gamma }\left(\frac{d}{2}\right)^2}}`$ (99)
$`=`$ $`\chi _d\left(\rho _L^{\mathrm{div}}\right)^2=\chi _d\left(\rho \left(\sigma \right)\right)^2`$ (100)
$`\mathrm{\Delta }\rho _L^{2,\mathrm{cross}}`$ $`=`$ $`{\displaystyle \frac{d}{d1}}G_{L,x_{}x_{}}^{\mathrm{div}}G_{L,x_{}x_{}}^{\mathrm{fin}}+\left(2G_{L,zz}^{\mathrm{div}}G_{L,zz}^{\mathrm{fin}}+G_{L,x_{}x_{}}^{\mathrm{div}}G_{L,zz}^{\mathrm{fin}}+G_{L,zz}^{\mathrm{div}}G_{L,x_{}x_{}}^{\mathrm{fin}}\right)`$ (101)
$`=`$ $`{\displaystyle \frac{\left(d+1\right)\mathrm{\Gamma }\left(\frac{d+1}{2}\right)}{2^{\frac{3\left(d+1\right)}{2}}d\pi ^{\frac{d}{2}}\sigma ^{d+1}\mathrm{\Gamma }\left(\frac{d}{2}\right)}}\left(G_{L,x_{}x_{}}^{\mathrm{fin}}+G_{L,zz}^{\mathrm{fin}}\right)`$ (102)
$`=`$ $`2\chi _d\rho _L^{\mathrm{div}}\rho _L^{\mathrm{fin}}`$ (103)
$`\mathrm{\Delta }\rho _L^{2,\mathrm{fin}}`$ $`=`$ $`{\displaystyle \frac{d\left(G_{L,x_{}x_{}}^{\mathrm{fin}}\right)^2}{2\left(d1\right)}}+G_{L,zz}^{\mathrm{fin}}\left(G_{L,x_{}x_{}}^{\mathrm{fin}}+G_{L,zz}^{\mathrm{fin}}\right)`$ (104)
$`=`$ $`{\displaystyle \frac{d^2l^{2(d+1)}}{2^{2d+7}\pi ^{d+3}}}\mathrm{\Gamma }\left({\displaystyle \frac{d}{2}}\right)^2\mathrm{\Gamma }\left({\displaystyle \frac{d}{2}}\right)^2`$ (108)
$`\times {\displaystyle \underset{p,q=1}{\overset{\mathrm{}}{}}}((1)^{p+q}\mathrm{\hspace{0.17em}2}^{2\left(p+q\right)}l^{2\left(p+q\right)}pq(2p1)(2q1)\sigma ^{2\left(p+q2\right)}`$
$`\times \left(4+d^26q+d\left(2p+2q3\right)+2p\left(4q3\right)\right)`$
$`\times {\displaystyle \frac{B_{2p+d1}B_{2q+d1}\left(2p3\right)!!\left(2q3\right)!!}{\left(2p+d1\right)\left(2q+d1\right)\left(2p\right)!\left(2q\right)!}}{\displaystyle \frac{\mathrm{\Gamma }\left(p\frac{1}{2}\right)\mathrm{\Gamma }\left(q\frac{1}{2}\right)}{\mathrm{\Gamma }\left(p+\frac{d}{2}\right)\mathrm{\Gamma }\left(q+\frac{d}{2}\right)}})`$
$`\stackrel{\sigma 0}{}`$ $`{\displaystyle \frac{d^3\pi ^dB_{d+1}^2\mathrm{\Gamma }\left(\frac{d}{2}\right)^2\mathrm{\Gamma }\left(\frac{d}{2}\right)^2}{8\left(d+1\right)L^{2\left(d+1\right)}\mathrm{\Gamma }\left(1+\frac{d}{2}\right)^2}}`$ (109)
$`=`$ $`{\displaystyle \frac{d\left(d+1\right)}{2}}\left(\rho _{L,\mathrm{reg}}\right)^2`$ (110)
From this we see the divergent and cross terms can be related to the smeared energy density via
$$\mathrm{\Delta }\rho _L^{2,\mathrm{div}}+\mathrm{\Delta }\rho _L^{2,\mathrm{cross}}=\chi _d\left\{\left(\rho _L^{\mathrm{div}}\right)^2+2\rho _L^{\mathrm{div}}\rho _L^{\mathrm{fin}}\right\}$$
(111)
where $`\chi _d`$ is the function that relates the fluctuations of the energy density to the mean energy density when the boundaries are not present, i.e., Minkowski space. This leads us to interpret these terms as due the vacuum fluctuations that are always present. With this in mind, we define the regularized fluctuations of the energy density
$`\mathrm{\Delta }\rho _{L,\mathrm{reg}}^2`$ $`=`$ $`\underset{\sigma 0}{lim}\left(\mathrm{\Delta }\rho _L^2\left(f\right)\chi _d\left\{\left(\rho _L^{\mathrm{div}}\right)^2+2\rho _L^{\mathrm{div}}\rho _L^{\mathrm{fin}}\right\}\right)`$ (112)
$`=`$ $`\chi _{d,L}\left(\rho _{L,\mathrm{reg}}\right)^2`$ (113)
where
$$\chi _{d,L}\frac{d\left(d+1\right)}{2}.$$
(114)
We also define a regularized version of the dimensionless measure $`\mathrm{\Delta }`$:
$$\mathrm{\Delta }_{L,\mathrm{Reg}}\frac{\mathrm{\Delta }\rho _{L,\mathrm{Reg}}^2}{\mathrm{\Delta }\rho _{L,\mathrm{Reg}}^2+\left(\rho _{L,\mathrm{Reg}}\right)^2}=\frac{d\left(d+1\right)}{2+d+d^2}$$
(115)
and note the values:
$`\begin{array}{ccccc}& & & & \\ d& 1& 3& 5& \mathrm{}\\ & & & & \\ \mathrm{\Delta }_{L,\mathrm{Reg}}& \frac{1}{2}& \frac{6}{7}& \frac{15}{16}& 1\end{array}`$
Following the procedures described in Appendix B, we have made two plots, Fig. 1 of $`\mathrm{\Delta }(\sigma ,L)`$ and $`\mathrm{\Delta }_{L,Reg}`$ versus $`\sigma /L`$, (which we call $`\sigma ^{}`$ here for short); and Fig. 2 of $`\rho _{L,Reg}`$ and $`\sqrt{\mathrm{\Delta }\rho _{L,Reg}^2}`$ versus $`\sigma ^{}`$. The range of $`\sigma ^{}`$ is limited to $`0.4`$ because going any further would make the meaning of a local energy density ill-defined, as the smearing of the field extends to the Casimir boundary in space. (We believe this infrared limit also carry important physical meaning in reference to the structure of spacetime, it is outside the focus of this paper.)
Let us ponder on the meaning they convey. In Fig. 1, we first note that both curves are of the order unity. But the behavior of $`\mathrm{\Delta }`$ (recall that the energy density fluctuations thus defined include the cross term along with the finite part and the state independent divergent part) is relatively insensitive to the smearing width, whereas $`\mathrm{\Delta }_{L,Reg}`$, which measures only the finite part of the energy density fluctuations to the mean has more structure. In particular, it saturates its upper bound of 1 around $`\sigma ^{}=0.24`$. Note that if one adheres to the KF criterion one would say that semiclassical gravity fails, but all that is happening here is that $`\rho _{L,\mathrm{Reg}}=0`$ while $`\mathrm{\Delta }\rho _{L,\mathrm{Reg}}^2`$ shows no special feature. The real difference between these two functions is the cross term, which is responsible for their markedly different structure and behavior. We have more to say about what to make of the cross term in the last section, which should be contrasted with the opinion of Wu and Ford on its physical significance. In Fig. 2, the main feature to notice is that the regularized energy density crosses from negative to positive values at around $`\sigma ^{}=0.24`$ . The negative Casimir energy density calculated in a point-wise field theory which corresponds to small ranges of $`\sigma ^{}`$ is expected, and is usually taken to signify the quantum nature of the Casimir state. As $`\sigma ^{}`$ increases we are averaging the field operator over a larger region, and thus sampling the field theory from the ultraviolet all the way to the infrared region. At large $`\sigma ^{}`$ finite size effect begins to set in. The difference and relation of these two effects are explained in : Casimir effect arises from summing up the quantum fluctuations of ALL modes (as altered by the boundary), with no insignificant short wavelength contributions, whereas finite size effect has dominant contributions from the LONGEST wavelength modes, and thus reflect the large scale behavior. As the smearing moves from a small scale to the far boundary of space, the behavior of the system is expected to shift from a Casimir-dominated to a finite size-dominated effect. This could be the underlying reason in the crossover behavior of $`\rho _{L,Reg}`$.
## V Point Separated Energy Density and Fluctuations in Minkowski Space
We return now to the Minkowski space and consider the point-separated expressions for the energy density and its fluctuations. This will lead to some interesting new observations about point-separation and maybe even the extended structure of spacetime. Consider the point separated energy density (58) with $`\sigma =0`$ in a Minkowski space with $`d`$-spatial dimensions
$`\rho (t,𝐱)`$ $`=`$ $`{\displaystyle 𝑑\mu (𝐤)N_k^2\omega ^2\mathrm{cos}(𝐱𝐤t\omega )}`$ (116)
$`=`$ $`{\displaystyle \frac{1}{2^d\pi ^{\frac{d+1}{2}}\mathrm{\Gamma }\left(\frac{d1}{2}\right)}}{\displaystyle _0^{\mathrm{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{cos}\left(xk_xt\sqrt{k_{}^2+k_x^2}\right)k_{}^{d2}\sqrt{k_{}^2+k_x^2}𝑑k_x𝑑k_{}`$ (117)
where we take $`𝐱=𝐱_1𝐱_2=x\widehat{x}`$ and decompose $`𝐤=(k_x,𝐤_{})`$ into one component along $`\widehat{x}`$ and two perpendicular to $`\widehat{x}`$. We change variables to $`k_x=k\mathrm{cos}\varphi `$ and $`k_{}=|𝐤_{}|=k\mathrm{sin}\varphi `$ and evaluate
$`\rho (t,x)`$ $`=`$ $`{\displaystyle \frac{1}{2^d\pi ^{\frac{d+1}{2}}\mathrm{\Gamma }\left(\frac{d1}{2}\right)}}{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^\pi }k^d\mathrm{cos}(k\left(tx\mathrm{cos}\varphi \right))\mathrm{sin}^{d+2}\varphi d\varphi dk`$ (118)
$`=`$ $`{\displaystyle \frac{1}{2^{\frac{d}{2}+1}\pi ^{\frac{d}{2}}}}{\displaystyle _0^{\mathrm{}}}k^d\left(k^2x^2\right)^{\frac{2d}{4}}J_{\frac{d}{2}1}\left(\sqrt{k^2x^2}\right)\mathrm{cos}(kt)𝑑k`$ (119)
$`=`$ $`{\displaystyle \frac{1}{2^{\frac{d}{2}+1}\pi ^{\frac{d}{2}}}}\mathrm{}{\displaystyle _0^{\mathrm{}}}k^d\left(k^2x^2\right)^{\frac{2d}{4}}J_{\frac{d}{2}1}\left(\sqrt{k^2x^2}\right)e^{ϵk+ikt}𝑑k`$ (120)
A small positive imaginary part $`iϵ`$ is added to $`t`$ to guarantee the $`k`$ integral convergences. Continuing,
$`=`$ $`{\displaystyle \frac{1}{2\pi ^{\frac{d+1}{2}}}}\mathrm{}\left\{{\displaystyle \frac{\left(d\left(ϵit\right)^2x^2\right)}{\left(ϵit\right)^{d+3}}}\left(1+{\displaystyle \frac{x^2}{\left(ϵit\right)^2}}\right)^{\frac{d+3}{2}}\mathrm{\Gamma }\left({\displaystyle \frac{d+1}{2}}\right)\right\}`$ (121)
$`=`$ $`{\displaystyle \frac{1}{2\pi ^{\frac{d+1}{2}}}}{\displaystyle \frac{\left(dt^2+x^2\right)}{\left(t^2x^2\right)^{\frac{d+3}{2}}}}\mathrm{\Gamma }\left({\displaystyle \frac{d+1}{2}}\right)\mathrm{sin}({\displaystyle \frac{d\pi }{2}})`$ (122)
The $`\mathrm{sin}\left(d\pi /2\right)`$ factor shows this result only holds for odd $`d`$. Restricting ourselves to odd $`d`$’s, the final result for the point separated energy density in Minkowski space is
$$\rho (t,x)=\frac{\left(1\right)^{\frac{d1}{2}}\mathrm{\Gamma }\left(\frac{d+1}{2}\right)}{2\pi ^{\frac{d+1}{2}}}\frac{\left(dt^2+x^2\right)}{\left(t^2x^2\right)^{\frac{d+3}{2}}}$$
(123)
Now we consider the correlation function (59). By identifying the two point function derivatives
$`G_{x_{}x_{}}(t,x)`$ $`=`$ $`0\left|_x_{}^2\left(\widehat{\varphi }(t_1,𝐱_1)\widehat{\varphi }(t_2,𝐱_2)\right)\right|0`$ (125)
$`=`$ $`{\displaystyle \frac{1}{2^d\pi ^{\frac{d+1}{2}}\mathrm{\Gamma }\left(\frac{d1}{2}\right)}}{\displaystyle _0^{\mathrm{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{k_{}^de^{i\left(xk_xt\sqrt{k_{}^2+k_x^2}\right)}}{\sqrt{k_{}^2+k_x^2}}}𝑑k_x𝑑k_{}`$ (126)
$`G_{xx}(t,x)`$ $`=`$ $`0\left|{\displaystyle \frac{^2}{x^2}}\left(\widehat{\varphi }(t_1,𝐱_1)\widehat{\varphi }(t_2,𝐱_2)\right)\right|0`$ (128)
$`=`$ $`{\displaystyle \frac{1}{2^d\pi ^{\frac{d+1}{2}}\mathrm{\Gamma }\left(\frac{d1}{2}\right)}}{\displaystyle _0^{\mathrm{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{k_{}^{d2}k_x^2}{\sqrt{k_{}^2+k_x^2}}}e^{i\left(xk_xt\sqrt{k_{}^2+k_x^2}\right)}𝑑k_x𝑑k_{}`$ (129)
$`G_{tx}(t,x)`$ $`=`$ $`0\left|{\displaystyle \frac{^2}{tx}}\left(\widehat{\varphi }(t_1,𝐱_1)\widehat{\varphi }(t_2,𝐱_2)\right)\right|0`$ (131)
$`=`$ $`{\displaystyle \frac{1}{2^d\pi ^{\frac{d+1}{2}}\mathrm{\Gamma }\left(\frac{d1}{2}\right)}}{\displaystyle _0^{\mathrm{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}k_{}^{d2}k_xe^{i\left(xk_xt\sqrt{k_{}^2+k_x^2}\right)}𝑑k_x𝑑k_{}`$ (132)
we write the energy density correlation function as
$$\mathrm{\Delta }\rho ^2(t,x)=\frac{dG_{x_{}x_{}}^2(t,x)}{2\left(d1\right)}+G_{x_{}x_{}}(t,x)G_{xx}(t,x)+G_{xx}^2(t,x)+G_{tx}^2(t,x)$$
(133)
We proceed with the evaluation of the Green functions in a similar manner as we did for the point separated energy density above and obtain
$`G_{x_{}x_{}}(t,x)`$ $`=`$ $`{\displaystyle \frac{(1)^{\frac{d+1}{2}}\mathrm{\Gamma }\left(\frac{d+1}{2}\right)}{2\pi ^{\frac{d+1}{2}}}}{\displaystyle \frac{(d1)}{\left(t^2x^2\right)^{\frac{d+1}{2}}}}`$ (135)
$`G_{xx}(t,x)`$ $`=`$ $`{\displaystyle \frac{(1)^{\frac{d+1}{2}}\mathrm{\Gamma }\left(\frac{d+1}{2}\right)}{2\pi ^{\frac{d+1}{2}}}}{\displaystyle \frac{\left(t^2+dx^2\right)}{\left(t^2x^2\right)^{\frac{d+3}{2}}}}`$ (136)
$`G_{tx}(t,x)`$ $`=`$ $`{\displaystyle \frac{(1)^{\frac{d+1}{2}}\mathrm{\Gamma }\left(\frac{d+1}{2}\right)}{2\pi ^{\frac{d+1}{2}}}}{\displaystyle \frac{(d+1)tx}{\left(x^2t^2\right)^{\frac{d+3}{2}}}}`$ (137)
With these results, correlation function is
$$\mathrm{\Delta }\rho ^2(t,x)=\frac{\mathrm{\Gamma }\left(\frac{d+1}{2}\right)^2}{\pi ^{d+1}}\left(\frac{4t^2x^2+d\left(t^2+x^2\right)^2}{\left(t^2x^2\right)^{d+3}}\right)$$
(138)
The constant $`\chi _d`$ is now a function of the temporal and spatial separation,
$$\chi _d(t,x)=\frac{d+1}{2}\left(\frac{4t^2x^2+d\left(t^2+x^2\right)^2}{\left(dt^2+x^2\right)^2}\right),$$
(139)
and we write the correlation function in terms of the square of the point separated energy density
$$\mathrm{\Delta }\rho ^2(t,x)=\chi _d(t,x)\left(\rho (t,x)\right)^2$$
(140)
Our dimensionless measure is also a function of the separation:
$$\mathrm{\Delta }(t,x)=\frac{\left(d+1\right)\left(d(t^2+x^2)^2+4t^2x^2\right)}{\left(d+1\right)\left(d(t^2+x^2)^2+4t^2x^2\right)+2(dt^2+x^2)^2}$$
(141)
To extract physical meaning out of this for a point-wise quantum field theory, we have to work in the $`(t,x)0`$ limit (recall $`t=t_1t_2`$, $`𝐱=𝐱_1𝐱_\mathrm{𝟐}=x\widehat{x}`$), for only then $`\rho (t,x)0\left|\widehat{\rho }\right|0`$. With this in mind, we parameterize the direction dependence via $`t=r\mathrm{sin}\theta `$ and $`x=r\mathrm{cos}\theta `$ (note this is only a parameterization, the imaginary time $`\tau `$ and $`x`$ shown below would carry physical meaning in the Euclidean sense). In the $`r0`$ limit we have
$$\mathrm{\Delta }(\theta )=\frac{\left(d+1\right)\left(1+2d\mathrm{cos}(4\theta )\right)}{\left(d+1\right)\left(1+2d\mathrm{cos}(4\theta )\right)+4(\mathrm{cos}^2(\theta )+d\mathrm{sin}^2(\theta ))^2}$$
(142)
which, as expected is finite. Taking the limit along the $`t`$axis ($`\theta =\pi /2`$), we get
$$\mathrm{\Delta }(t,x=0)=\frac{1+d}{1+3d}=\mathrm{\Delta }_{\mathrm{Minkowski}}$$
(143)
On the other hand, taking the limit along the spatial direction:
$$\mathrm{\Delta }(t=0,x)=\frac{d\left(d+1\right)}{2+d+d^2}=\mathrm{\Delta }_{L,\mathrm{Reg}}$$
(144)
We also approach this problem in another way. Since both the point separated energy density and the correlation function have a direction dependence, we “average” over the direction. We take the hyper-spherical averaging procedure. This involves first Wick rotating to imaginary time $`(ti\tau )`$. Then we take the hyper-spherical average in the Euclidean geometry and then Wick rotate back to Minkowski space. For the energy density
$$\rho _E(\tau ,x)=\frac{\mathrm{\Gamma }\left(\frac{d+1}{2}\right)}{2\pi ^{\frac{d+1}{2}}}\frac{\left(d\tau ^2x^2\right)}{\left(\tau ^2+x^2\right)^{\frac{d+3}{2}}}$$
(145)
Now expressing $`\tau =r\mathrm{sin}\theta `$ and $`x=r\mathrm{cos}\theta `$ we do the averaging
$`\rho _E(r)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{2\pi }}\rho _E(r\mathrm{sin}\theta ,r\mathrm{cos}\theta )𝑑\theta `$ (146)
$`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{d+1}{2}\right)}{4\pi ^{\frac{d+3}{2}}r^{d+1}}}{\displaystyle _0^{2\pi }}(d\mathrm{sin}(\theta )^2\mathrm{cos}(\theta )^2)d\theta `$ (147)
$`=`$ $`{\displaystyle \frac{(d1)\mathrm{\Gamma }\left(\frac{d+1}{2}\right)}{4\pi ^{\frac{d1}{2}}r^{d+1}}}`$ (148)
W do the same for the correlation function:
$`\mathrm{\Delta }\rho _E^2(r)`$ $`=`$ $`{\displaystyle \frac{(d+1)\mathrm{\Gamma }\left(\frac{d+1}{2}\right)^2}{32\pi ^{d+2}r^{2(d+1)}}}{\displaystyle _0^{2\pi }}\left(d1+\left(d+1\right)\mathrm{cos}(4\theta )\right)𝑑\theta `$ (149)
$`=`$ $`{\displaystyle \frac{\left(d^21\right)\mathrm{\Gamma }\left(\frac{d+1}{2}\right)^2}{32\pi ^{d+1}r^{2\left(d+1\right)}}}`$ (150)
With these results, we have
$$\chi _{d,\mathrm{Avg}}=\frac{1+d}{2\left(d1\right)}\mathrm{and}\mathrm{\Delta }_{d,\mathrm{Avg}}=\frac{1+d}{3d1}$$
(151)
independent of whether or not we Wick rotate back to Minkowski space. Also,
$`\begin{array}{ccccc}& & & & \\ d& 1& 3& 5& \mathrm{}\\ & & & & \\ \mathrm{\Delta }_{\mathrm{Avg}}& 1& \frac{1}{2}& \frac{3}{7}& \frac{1}{3}\end{array}`$
It is interesting to observe that the first set of results depend on the direction the two points come together, and changes if one averages over all directions. This feature of point-separation is known, but it could also reveal some properties of possible extended structure of the underlying spacetime.
## VI Discussions
Let us ponder on the implication of these findings pertaining to a) fluctuations to mean ratio and the validity of semiclassical gravity b) the dependence of fluctuations on both the intrinsic scale (defined by smearing or point-separation) and the extrinsic scale (such as the Casimir or finite temperature periodicity) c) the treatment of divergences and meaning of regularization
### A Fluctuation to Mean ratio and Validity of SCG
If we adopt the criterion of Kuo and Ford that the variance of the fluctuation relative to the mean-squared (vev taken with respect to the ordinary Minkowskian vacuum) being of the order unity be an indicator of the failure of SCG, then all spacetimes studied above would indiscriminately fall into that category, and SCG fails wholesale, regardless of the scale these physical quantities are probed. This contradicts with the common expectation that SCG is valid at scales below Planck energy. We believe the criterion for the validity or failure of a theory should depend on the range or the energy probed. Our findings here suggest that this is indeed the case: Both the mean (the vev of EMT) with respect to the Minkowski vacuum) AND the fluctuations of EMT increase as the scale deceases. As one probes into an increasingly finer scale or higher energy the expectation value of EMT grows in value and the induced metric fluctuations become important, leading to the failure of SCG. A generic scale for this to happen is the Planck length. At such energy densities and above, particle creation from the quantum field vacuum would become copious and their backreaction on the background spacetime would become important . Fluctuations in the quantum field EMT entails these quantum processes. The induced metric fluctuations renders the smooth manifold structure of spacetime inadequate, spacetime foams including topological transitions begin to appear and SCG no longer can provide an adequate description of these dominant processes. This picture first conjured by Wheeler is consistent with the common notion adopted in SCG, and we believe it is a valid one.
### B Dependence of fluctuations on intrinsic and extrinsic scales
In the previous section we have presented some detailed analysis on the results of our calculations for the fluctuations of the energy density for the separate cases of Minkowski and Casimir states. Let us now look at the bigger picture and see if we can capture the essence of these results with some general qualitative arguments. We want to see if there is any simple reason behind the following results we obtained:
a) $`\mathrm{\Delta }=O(1)`$
b) the specific numeric value of $`\mathrm{\Delta }`$ for the different cases
c) why $`\mathrm{\Delta }`$ for the Minkowski space from the coincidence limit of taking a spatial point separation is identical to the Casimir case at the coincidence limit (6/7) and identical to the hot flat space result (2/5) from taking the coincidence limit of a temporal point separation?
Point a) has also been shown by earlier calculations , and our understanding is that this is true only for states of quantum nature, including the vacuum and certain squeezed states, but probably not true for states of a more classical nature like the coherent state. We also emphasized that this result should not be used as a criterion for the validity of semiclassical gravity.
For point b), we can trace back the calculation of the fluctuations (second moment) of the energy momentum tensor in ratio to its mean (first moment) to the integral of the term containing the inner product of two momenta $`𝐤_1𝐤_2`$ summed over all participating modes. The modes contributing to this are different for different geometries, e.g., Minkowski versus Casimir boundary –for the Einstein universe this enters as 3j symbols – and could account for the difference in the numerical values of $`\mathrm{\Delta }`$ for the different cases.
For point c) the difference of results between taking the coincidence limit of a spatial versus a temporal point separation is well-known in QFTCST. The case of temporal split involves integration of three spatial dimensions while the case of spatial split involve integration of two remaining spatial and one temporal dimension, which would give different results. The calculation of fluctuations involves the second moment of the field and is in this regard similar to what enters into the calculation of moments of inertia for rotating objects. We suspect that the difference between the temporal and the spatial results is similar, to the extent this analogy holds, to the difference in the moment of inertia of the same object but taken with respect to different axes of rotation.
It may appear surprising, as we felt when we first obtained these results, that in a Minkowski calculation the result of Casimir geometry or thermal field should appear, as both cases involve a scale – the former in the spatial dimension and the latter in the (imaginary) temporal dimension. But if we note that the results for Casimir geometry or thermal field are obtained at the coincidence (ultraviolet) limit, when the scale (infrared) of the problem does not intercede in any major way, then the main components of the calculations for these two cases would be similar to the two cases of taking coincident limit in the spatial and temporal directions in Minkowski space. All of these cases are effectively devoid of scale as far as the pointwise field theory is concerned. As soon as we depart from this limit the effect of the presence of a scale shows up. The point-separated or field-smeared results for the Casimir calculation in Sec. 4 shows clearly that the boundary scale enters in a major way and the result for the fluctuations and the ratio are different from those obtained at the coincident limit. For other cases where a scale enters intrinsically in the problem such as that of a massive or non-conformally coupled field it would show a similar effect in these regards as the present cases (of Casimir and thermal field) where a periodicity condition exists (in the spatial and temporal directions respectively). We expect a similar strong disparity between point-coincident and point-separated cases. The field theory changes its nature in a fundamental and physical way when this limit is taken. This brings us to the second major issue brought out in this investigation, i.e., the appearance of divergences and the meaning of regularization in the light of a point-separated versus a point-defined quantum field theory.
### C Regularization in the Fluctuations of EMT
From our calculations, the smeared energy density fluctuations for the Casimir topology has the form
$$\mathrm{\Delta }\rho _L^2(\sigma )=\mathrm{\Delta }\rho _L^{\mathrm{div}}+\mathrm{\Delta }\rho _L^{\mathrm{cross}}+\mathrm{\Delta }\rho _L^{\mathrm{fin}}$$
(152)
with
$`\mathrm{\Delta }\rho _L^{\mathrm{div}}`$ $`=`$ $`\chi _d\left(\rho _L^{\mathrm{div}}\right)^2=\chi _d(\rho (\sigma ))^2`$ (154)
$`\mathrm{\Delta }\rho _L^{\mathrm{cross}}`$ $`=`$ $`2\chi _d\rho _L^{\mathrm{div}}\rho _L^{\mathrm{fin}}`$ (155)
$`\mathrm{\Delta }\rho _L^{\mathrm{fin}}`$ $`=`$ $`2\chi _{d,L}\left(\rho _L^{\mathrm{fin}}\right)^2+\mathrm{terms}\mathrm{that}\mathrm{vanish}\mathrm{as}\sigma 0`$ (156)
where $`\chi _d`$ is the ratio between the fluctuations for Minkowski space and the square of the corresponding energy density: $`\mathrm{\Delta }\rho ^2=\chi _d(\rho (\sigma ))^2`$. Our results show that $`\mathrm{\Delta }\rho _L^2(\sigma )`$ diverges as the width $`\sigma `$ of the smearing function shrinks to zero with contributions from the truly divergent and the cross terms. We also note that the divergent term $`\mathrm{\Delta }\rho ^{\mathrm{div}}`$ is state independent, in the sense that it is independent of $`L`$, while the cross term $`\mathrm{\Delta }\rho ^{\mathrm{cross}}`$ is state dependent, as is the finite term $`\mathrm{\Delta }\rho ^{\mathrm{fin}}`$.
If we want to ask about the strength of fluctuations of the energy density, the relevant quantity to study is the energy density correlation function $`H(x,y)=\widehat{\rho }(x)\widehat{\rho }(y)\widehat{\rho }(x)\widehat{\rho }(y)`$. It is finite at $`xy`$ for a linear quantum theory (this happens since the divergences for $`\widehat{\rho }(x)\widehat{\rho }(y)`$ are exactly the same as the product $`\widehat{\rho }(x)\widehat{\rho }(y)`$), but diverges as $`yx`$, corresponding to the coincident or unsmeared limit $`\sigma 0`$.
To define a procedure for rendering our expression for $`\mathrm{\Delta }\rho _L^2(\sigma )`$ finite, one can see that there exists choices – which means ambiguities in the regularization scheme. Three possibilities present themselves: The first is to just drop the state independent $`\mathrm{\Delta }\rho ^{\mathrm{div}}`$. This is easily seen to fail since we are left with the divergences from the cross term. The second is to neglect all terms that diverge as $`\sigma 0`$. This is too rash a move since $`\mathrm{\Delta }\rho ^{\mathrm{cross}}`$ has, along with its divergent parts, ones that are finite in the $`\sigma 0`$ limit. This comes about since it is of the form $`\rho _L^{\mathrm{div}}\rho _L^{\mathrm{fin}}`$ and the negative powers of $`\sigma `$ present in $`\rho _L^{\mathrm{div}}`$ will cancel out against the positive powers in $`\rho _L^{\mathrm{fin}}`$. Besides, they yield results in disagreement with earlier results using well-tested methods such as normal ordering in flat space and zeta-function regularization in curved space .
The third choice is the one we have used in this paper. For the energy density, we can think of regularization as computing the contribution “above and beyond” the Minkowski vacuum contribution. Same for regularizing the fluctuations. So we need to first determine for Minkowski space vacuum how the fluctuations of the energy density are related to the vacuum energy density $`\mathrm{\Delta }\rho ^2=\chi (\rho )^2`$. This we obtained for finite smearing. For Casimir topology the sum of the divergent and cross terms take the form
$$\mathrm{\Delta }\rho _L^{\mathrm{div}}+\mathrm{\Delta }\rho _L^{\mathrm{cross}}=\chi \left\{\left(\rho _L^{\mathrm{div}}\right)^2+2\rho _L^{\mathrm{div}}\rho _L^{\mathrm{fin}}\right\}=\chi \left\{\left(\rho _L\right)^2\left(\rho _L^{\mathrm{fin}}\right)^2\right\}$$
(157)
where $`\chi `$ is the ratio derived for Minkowski vacuum. We take this to represent the (state dependent) vacuum contribution. What we find interesting is that to regularize the smeared energy density fluctuations, a state dependent subtraction must be used. With this, just the $`\sigma 0`$ limit of the finite part $`\mathrm{\Delta }\rho _L^{\mathrm{fin}}`$ is identified as the regularized fluctuations $`\mathrm{\Delta }\rho _{L,\mathrm{Reg}}^2`$. The ratio $`\chi _L`$ thus obtained gives exactly the same result as derived by Kou and Ford for $`d=3`$ via normal ordering and by ourselves for arbitrary $`d`$ via the $`\zeta `$-function.
That this procedure is the one to follow can be seen by considering the problem from the point separation method. For this method, the energy density expectation value is defined as the $`x^{}x`$ limit of
$$\rho (x,x^{})=𝒟_{x,x^{}}G(x,x^{})$$
(158)
for the suitable Green function $`G(x,x^{})`$ and $`𝒟_{x,x^{}}`$ is a second order differential operator. (For the more general stress tensor, details are reviewed in .) In the limit $`x^{}x`$, $`G(x,x^{})`$ is divergent. The Green function is regularized by subtracting from it a Hadamard form $`G^L(x,x^{})`$: $`G_{\mathrm{Reg}}(x,x^{})=G(x,x^{})G^L(x,x^{})`$ . With this, the regularized energy density can be obtained
$$\rho _{\mathrm{Reg}}(x)=\underset{x^{}x}{lim}\left(𝒟_{x,x^{}}G_{\mathrm{Reg}}(x,x^{})\right)$$
(159)
Or, re-arranging terms, we can define the divergent and finite pieces as
$$G^{\mathrm{div}}(x,x^{})=G^L(x,x^{}),G^{\mathrm{fin}}(x,x^{})=G_{\mathrm{Reg}}(x,x^{})=G(x,x^{})G^L(x,x^{})$$
(160)
and
$$\rho (x,x^{})=\rho ^{\mathrm{div}}(x,x^{})+\rho ^{\mathrm{fin}}(x,x^{})$$
(161)
$`\rho ^{\mathrm{div}}(x,x^{})=𝒟_{x,x^{}}G^{\mathrm{div}}(x,x^{})\mathrm{and}\rho ^{\mathrm{fin}}(x,x^{})=𝒟_{x,x^{}}G^{\mathrm{fin}}(x,x^{})`$
so that $`\rho _{\mathrm{Reg}}(x)=lim_{x^{}x}\rho ^{\mathrm{fin}}(x,x^{})`$, which corresponds to the $`\sigma 0`$ limit in our computation of the Casimir energy density.
Now turning to the fluctuations, we have the point separated expression for the correlation function
$$H(x,y)=\underset{x^{}x}{lim}\underset{y^{}y}{lim}𝒟_{x,x^{}}𝒟_{y,y^{}}G(x,x^{},y,y^{})$$
(162)
where $`G(x,x^{},y,y^{})`$ is the suitable four point function. For linear theories we use Wick’s Theorem to express this in terms of products of Green functions $`G(x,x^{},y,y^{})=G(x,y)G(x^{},y^{})+\mathrm{permutations}\mathrm{of}(x,x^{},y,y^{})`$. Excluded from the permutations is $`G(x,x^{})G(y,y^{})`$. (Details are in , which includes correct identifications of needed permutations and Green functions.) The general form is
$$H(x,y)=\underset{x^{}x}{lim}\underset{y^{}y}{lim}𝒟_{x,x^{}}𝒟_{y,y^{}}G(x,y)G(x^{},y^{})+\mathrm{permutations}$$
(163)
The $`(x^{},y^{})(x,y)`$ limits are only retained to keep track of which derivatives act on which Green functions, but we can see there are no divergences for $`yx`$. However, to get the point-wise fluctuations of the energy density, the divergences from $`lim_{yx}G(x,y)`$ will present a problem. Splitting the Green function into its finite and divergent pieces, we can recognize terms leading to those we found for $`\mathrm{\Delta }\rho _L^2(\sigma )`$:
$$H(x,y)=H^{\mathrm{div}}(x,y)+H^{\mathrm{cross}}(x,y)+H^{\mathrm{fin}}(x,y)$$
(164)
where
$`H^{\mathrm{div}}(x,y)`$ $`=`$ $`\underset{x^{}x}{lim}\underset{y^{}y}{lim}𝒟_{x,x^{}}𝒟_{y,y^{}}G^{\mathrm{div}}(x,y)G^{\mathrm{div}}(x^{},y^{})`$ (166)
$`H^{\mathrm{cross}}(x,y)`$ $`=`$ $`2\underset{x^{}x}{lim}\underset{y^{}y}{lim}𝒟_{x,x^{}}𝒟_{y,y^{}}G^{\mathrm{div}}(x,y)G^{\mathrm{fin}}(x^{},y^{})`$ (167)
$`H^{\mathrm{fin}}(x,y)`$ $`=`$ $`\underset{x^{}x}{lim}\underset{y^{}y}{lim}𝒟_{x,x^{}}𝒟_{y,y^{}}G^{\mathrm{fin}}(x,y)G^{\mathrm{fin}}(x^{},y^{})`$ (168)
plus permutations. Thus we see the origin of both the divergent and cross terms. When the un-regularized Green function is used, we must get a cross term, along with the expected divergent term. If the fluctuations of the energy density is regularized via point separation, i.e. $`G(x,x^{})`$ is replaced by $`G_{\mathrm{Reg}}(x,x^{})=G^{\mathrm{fin}}(x,y)`$, then we should do the same replacement for the fluctuations. When this is done, it is only the finite part above that will be left and we can define the point-wise fluctuations as
$$\mathrm{\Delta }\rho _{\mathrm{Reg}}^2=\underset{yx}{lim}H^{\mathrm{fin}}(x,y)$$
(169)
The parallel with the smeared-field derivation presented here can be seen when the analysis of $`G_L(\sigma )_{,zz}`$ and $`G_L(\sigma )_{,x_{}x_{}}`$ in the Appendix is considered. There it is shown they are derivatives of Green functions and can be separated into state-independent divergent part and state-dependent finite contribution: $`G_L(\sigma )_{,i}=G_{L,i}^{\mathrm{div}}+G_{L,i}^{\mathrm{fin}}`$, same as the split hereby shown for the Green function.
When analyzing the energy density fluctuations, discarding the divergent piece is the same as subtracting from the Green function its divergent part. If this is done, we also no longer have the cross term, just as viewing the problem from the point separation method outlined above. We feel this makes it problematic to analyze the cross term without also including the divergent term. At the same time, regularization of the fluctuations involving the subtraction of state dependent terms as realized in this calculation raises new issues on regularization which merits further investigations.
To end this discussion, we venture one philosophical point we find resounding throughout all the cases studied here. It has to do with the meaning of a point-defined versus a point-separated field theory, the former we take as an effective theory coarse-grained from the latter, the point-separated theory reflecting a finer level of spacetime structure. It bears on the meaning of regularization, not just at the level of technical procedures, but related to finding an effective description and matching with physics observed at a coarser scale or lower energy.
In particular, we feel that finding a finite energy momentum tensor (and its fluctuations as we do here) which occupied the center of attention in the research of quantum field theory in curved spacetime in the 70’s is only a small part of a much larger and richer structure of theories of fields and spacetimes . We come to understand that whatever regularization method one uses to get these finite parts in a point-wise field theory should not be viewed as universally imparting meaning beyond its specified function, i.e., to identify the divergent pieces and provide a prescription for their removal. We believe the extended structure of spacetime (e.g., via point-separation or smearing ) and the field theory defined therein has its own much fuller meaning beyond just reproducing the well-recognized result in ordinary quantum field theory as we take the point-wise or coincident limit. In this way of thinking, the divergence- causing terms are only ‘bad’ when they are forced to a point-wise limit, because of our present inability to observe or resolve otherwise . If we accord them with the full right of existence beyond this limit, and acknowledge that their misbehavior is really due to our own inability to cope, we will be rewarded with the discovery of new physical phenomena and ideas of a more intricate world. (Maybe this is just another way to appreciate the already well-heeded paths of string theory.)
Acknowledgement We thank Professor Larry Ford for interesting discussions, especially on the meaning of the cross term, Dr. Alpan Raval for useful comments, especially on the generic nature of the vacuum fluctuations to the mean, and Prof. Raphael Sorkin for discussions on the relevance of our results to worm hole physics. A few researchers whom we have met have voiced their doubts to us on the Kuo-Ford criterion, and indicated a view similar to ours as expounded here. We would like to thank Prof. Paul Anderson and Prof. Ted Jacobson for conversations of this nature. This work is supported in part by NSF grant PHY98-00967
## A Evaluation of Two Point Functions for Casimir Topology
We want to compute the smeared derivatives of the field operator two-point functions for the Casimir geometry:
$$G_L(\sigma )_{,x_{}x_{}}=0_L\left|\left(\left(_{}\varphi \right)\left(f_𝐱\right)\right)^2\right|0_L\mathrm{and}G_L(\sigma )_{,zz}=0_L\left|\left(\left(_z\varphi \right)\left(f_𝐱\right)\right)^2\right|0_L,$$
(A1)
where $`|0_L`$ is the Casimir vacuum. Performing the differentiation and taking the vacuum expectation values, we need the integrals and sums
$`G_L(\sigma )_{,x_{}x_{}}`$ $`=`$ $`{\displaystyle \frac{l}{2^d\pi ^{\frac{d+1}{2}}\mathrm{\Gamma }\left(\frac{d1}{2}\right)}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{k^d}{\sqrt{k^2+l^2n^2}}}e^{2\left(k^2+l^2n^2\right)\sigma ^2}𝑑k`$ (A3)
$`G_L(\sigma )_{,zz}`$ $`=`$ $`{\displaystyle \frac{l}{2^d\pi ^{\frac{d+1}{2}}\mathrm{\Gamma }\left(\frac{d1}{2}\right)}}{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{k^{d+2}l^2n^2}{\sqrt{k^2+l^2n^2}}}e^{2\left(k^2+l^2n^2\right)\sigma ^2}𝑑k`$ (A5)
With the definitions of the functions
$`F_{x_{}x_{}}(n)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{2k^d}{\sqrt{k^2+l^2n^2}}}e^{2\left(k^2+l^2n^2\right)\sigma ^2}𝑑k`$ (A7)
$`F_{zz}(n)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{2k^{d2}l^2n^2}{\sqrt{k^2+l^2n^2}}}e^{2\left(k^2+l^2n^2\right)\sigma ^2}𝑑k`$ (A9)
we use the Euler-Maclauren sum formula to re-arrange the terms to the more useful form ($`i=x_{}x_{}`$ or $`zz`$):
$`G_L(\sigma )_{,i}`$ $`=`$ $`{\displaystyle \frac{l}{2^d\pi ^{\frac{d+1}{2}}\mathrm{\Gamma }\left(\frac{d1}{2}\right)}}\underset{N\mathrm{}}{lim}\left({\displaystyle \frac{1}{2}}F_i(0)+{\displaystyle \underset{n=1}{\overset{N}{}}}F_i(n)\right)`$ (A10)
$`=`$ $`{\displaystyle \frac{l}{2^d\pi ^{\frac{d+1}{2}}\mathrm{\Gamma }\left(\frac{d1}{2}\right)}}\underset{N\mathrm{}}{lim}\left({\displaystyle _0^N}F_i(n)𝑑n+{\displaystyle \frac{1}{2}}F_i(N)+{\displaystyle \underset{p=1}{\overset{q}{}}}{\displaystyle \frac{B_{2p}}{(2p)!}}\left(F_i^{(2p1)}(N)F_i^{(2p1)}(0)\right)\right)`$ (A11)
As we will show, $`F_i(N)`$ vanishes exponentially with $`N`$ so that $`F_i(N)`$ and $`F_i^{(2p1)}(N)`$ give no contributions to the final result and we are left with
$$G_L(\sigma )_{,i}=\frac{l}{2^d\pi ^{\frac{d+1}{2}}\mathrm{\Gamma }\left(\frac{d1}{2}\right)}\left(_0^{\mathrm{}}F_i(n)𝑑n\underset{p=1}{\overset{q}{}}\frac{B_{2p}}{(2p)!}F_i^{(2p1)}(0)\right)$$
(A12)
This re-arrangement of the terms has allowed us to separate the expectation values into terms that diverge as $`\sigma 0`$ and those that are finite in this limit: $`G_L(\sigma )_{,i}=G_{L,i}^{\mathrm{div}}+G_{L,i}^{\mathrm{fin}}`$ with
$`G_{L,i}^{\mathrm{div}}`$ $`=`$ $`{\displaystyle \frac{l}{2^d\pi ^{\frac{d+1}{2}}\mathrm{\Gamma }\left(\frac{d1}{2}\right)}}{\displaystyle _0^{\mathrm{}}}F_i(n)𝑑n`$ (A14)
$`G_{L,i}^{\mathrm{fin}}`$ $`=`$ $`{\displaystyle \frac{l}{2^d\pi ^{\frac{d+1}{2}}\mathrm{\Gamma }\left(\frac{d1}{2}\right)}}{\displaystyle \underset{p=1}{\overset{q}{}}}{\displaystyle \frac{B_{2p}}{(2p)!}}F_i^{(2p1)}(0)`$ (A16)
The $`k`$ integrations give the explicit form of the functions $`F_i(n)`$:
$`F_{x_{}x_{}}(n)`$ $`=`$ $`{\displaystyle \frac{l^dn^d\mathrm{\Gamma }\left(\frac{d}{2}\right)\mathrm{\Gamma }\left(\frac{d+1}{2}\right)}{\sqrt{\pi }}}{}_{1}{}^{}F_{1}^{}({\displaystyle \frac{1}{2}};1+{\displaystyle \frac{d}{2}};2l^2n^2\sigma ^2)`$ (A19)
$`+{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{d}{2}\right)}{2^{\frac{d}{2}}\sigma ^d}}{}_{1}{}^{}F_{1}^{}({\displaystyle \frac{1d}{2}};1{\displaystyle \frac{d}{2}};2l^2n^2\sigma ^2)`$
and
$`F_{zz}(n)`$ $`=`$ $`{\displaystyle \frac{l^dn^d\mathrm{\Gamma }\left(1\frac{d}{2}\right)\mathrm{\Gamma }\left(\frac{d1}{2}\right)}{\sqrt{\pi }}}{}_{1}{}^{}F_{1}^{}({\displaystyle \frac{1}{2}};{\displaystyle \frac{d}{2}};2l^2n^2\sigma ^2)`$ (A21)
$`+{\displaystyle \frac{l^2n^2}{2^{\frac{d}{2}1}\sigma ^{d2}}}\mathrm{\Gamma }\left({\displaystyle \frac{d}{2}}1\right){}_{1}{}^{}F_{1}^{}({\displaystyle \frac{3}{2}}{\displaystyle \frac{d}{2}};2{\displaystyle \frac{d}{2}};2l^2n^2\sigma ^2)`$
with $`{}_{1}{}^{}F_{1}^{}(a;b;z)`$ the Kummer confluent hypergeometric function.
We carry out the $`n`$ integrations and obtain the divergent parts of the expectation values
$`G_{L,x_{}x_{}}^{\mathrm{div}}`$ $`=`$ $`{\displaystyle \frac{\left(d1\right)\mathrm{\Gamma }\left(\frac{d+1}{2}\right)}{2^{\frac{3\left(d+1\right)}{2}}d\pi ^{\frac{d}{2}}\mathrm{\Gamma }\left(\frac{d}{2}\right)\sigma ^{d+1}}}`$ (A23)
$`G_{L,zz}^{\mathrm{div}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{d+1}{2}\right)}{2^{\frac{3\left(d+1\right)}{2}}d\pi ^{\frac{d}{2}}\mathrm{\Gamma }\left(\frac{d}{2}\right)\sigma ^{d+1}}}`$ (A24)
Turning to the finite contribution, we need the general form of
$$H=\frac{d^{2p1}}{dn^{2p1}}\left(An^\beta {}_{1}{}^{}F_{1}^{}(\alpha ;\gamma ;2l^2n^2\sigma ^2)\right)|_{n=0}$$
(A25)
To this end, we make use of the relation
$$\frac{d^pg\left(an^2\right)}{dn^p}|_{n=0}=\{\begin{array}{cc}(p1)!!\left(2a\right)^{\frac{p}{2}}g^{\left(\frac{p}{2}\right)}(0);\hfill & p\mathrm{even}\hfill \\ 0;\hfill & p\mathrm{odd}\hfill \end{array}$$
(A26)
along with
$`{\displaystyle \frac{d{}_{1}{}^{}F_{1}^{}(\alpha ,\gamma ;z)}{dz}}`$ $`=`$ $`{\displaystyle \frac{\alpha }{\gamma }}{}_{1}{}^{}F_{1}^{}(\alpha +1,\gamma +1;z)`$ (A27)
$`{\displaystyle \frac{d^p{}_{1}{}^{}F_{1}^{}(\alpha ,\gamma ;z)}{dz^p}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(\alpha +p)\mathrm{\Gamma }(\gamma )}{\mathrm{\Gamma }(\alpha )\mathrm{\Gamma }(\gamma +p)}}{}_{1}{}^{}F_{1}^{}(\alpha +p,\gamma +p;z)`$ (A28)
For $`\gamma `$ not a negative integer, $`{}_{1}{}^{}F_{1}^{}(\alpha ,\gamma ;0)=1`$. These results lead to
$`H`$ $`=`$ $`{\displaystyle \frac{A\left(2p1\right)!}{\mathrm{\Gamma }\left(2p\beta \right)}}\left({\displaystyle \frac{d^{2p1\beta }{}_{1}{}^{}F_{1}^{}(\alpha ;\gamma ;2l^2n^2\sigma ^2)}{dn^{2p1\beta }}}\right)|_{n=0}`$ (A29)
$`=`$ $`\{\begin{array}{cc}\frac{A\left(2p1\right)!}{\mathrm{\Gamma }\left(2p\beta \right)}2^{2p\beta 1}\left(\left(l^2\sigma ^2\right)\right)^{\frac{2p\beta 1}{2}}\left(2p\beta 2\right)!!\left(\frac{d^{\frac{2p\beta 1}{2}}{}_{1}{}^{}F_{1}^{}(\alpha ;\gamma ;z)}{dz^{\frac{2p\beta 1}{2}}}\right)|_{z=0};\hfill & \beta \mathrm{odd}\hfill \\ 0;\hfill & \beta \mathrm{even}\hfill \end{array}`$ (A32)
Staying with $`\beta `$ odd, the final result is
$$\frac{A\left(1+2p\right)!}{\mathrm{\Gamma }\left(\beta +2p\right)}2^{2p\beta 1}\left(\left(l^2\sigma ^2\right)\right)^{\frac{2p\beta 1}{2}}\left(2p\beta 2\right)!!\frac{\mathrm{\Gamma }\left(\gamma \right)\mathrm{\Gamma }\left(\frac{1}{2}+\alpha \frac{\beta }{2}+p\right)}{\mathrm{\Gamma }\left(\alpha \right)\mathrm{\Gamma }\left(\frac{1}{2}\frac{\beta }{2}+\gamma +p\right)}$$
(A33)
To determine the finite contributions to the smeared Green function derivatives, we use
$`\begin{array}{cccccc}& & A& \beta & \alpha & \gamma \\ F_{x_{}x_{}}& \mathrm{term}\mathrm{\hspace{0.33em}1}& l^d\mathrm{\Gamma }\left(\frac{d}{2}\right)\mathrm{\Gamma }\left(\frac{d+1}{2}\right)/\sqrt{\pi }& d& \frac{1}{2}& 1+\frac{d}{2}\\ F_{x_{}x_{}}& \mathrm{term}\mathrm{\hspace{0.33em}2}& 2^{\frac{d}{2}}\sigma ^d\mathrm{\Gamma }\left(\frac{d}{2}\right)& 0& \frac{1d}{2}& 1\frac{d}{2}\\ F_{zz}& \mathrm{term}\mathrm{\hspace{0.33em}1}& l^d\mathrm{\Gamma }\left(1\frac{d}{2}\right)\mathrm{\Gamma }\left(\frac{d1}{2}\right)/\sqrt{\pi }& d& \frac{1}{2}& \frac{d}{2}\\ F_{zz}& \mathrm{term}\mathrm{\hspace{0.33em}2}& 2^{1\frac{d}{2}}l^2\sigma ^{2d}\mathrm{\Gamma }\left(1+\frac{d}{2}\right)& 2& \frac{3}{2}\frac{d}{2}& 2\frac{d}{2}\end{array}`$
For $`d`$ odd, only the first terms of $`F_{x_{}x_{}}`$ and $`F_{zz}`$ contribute. For $`d`$ even, the situation involves more analysis, since for the second terms, even though $`\beta `$ is even, $`\gamma `$ is a negative integer. This implies $`{}_{1}{}^{}F_{1}^{}(\alpha ,\gamma ,2l^2n^2\sigma ^2)`$ divergent structure needs to be considered as well. For $`d`$ odd,
$$G_{L,x_{}x_{}}^{\mathrm{fin}}=\frac{l^{d1}d\left(d1\right)\mathrm{\Gamma }\left(\frac{d}{2}\right)\mathrm{\Gamma }\left(\frac{d}{2}\right)}{2^{d+3}\pi ^{\frac{d+3}{2}}}\underset{p=1}{}\left(4l^2\right)^p\sigma ^{2\left(p1\right)}p\left(2p1\right)\frac{B_{2p+d1}\left(2p3\right)!!\mathrm{\Gamma }\left(p\frac{1}{2}\right)}{\left(2p+d1\right)\left(2p\right)!\mathrm{\Gamma }\left(p+\frac{d}{2}\right)}$$
(A35)
and
$$G_{L,zz}^{\mathrm{fin}}=\frac{l^{d1}d\mathrm{\Gamma }\left(\frac{d}{2}\right)\mathrm{\Gamma }\left(\frac{d}{2}\right)}{2^{d+3}\pi ^{\frac{d+3}{2}}}\underset{p=1}{\overset{\mathrm{}}{}}\left(4l^2\right)^p\sigma ^{2\left(p1\right)}p\left(2p1\right)\left(2p+d2\right)\frac{B_{2p+d1}\left(2p3\right)!!\mathrm{\Gamma }\left(p\frac{1}{2}\right)}{\left(2p+d1\right)\left(2p\right)!\mathrm{\Gamma }\left(p+\frac{d}{2}\right)}$$
(A36)
## B Plotting Smeared Casimir Results
In this appendix we outline how $`\mathrm{\Delta }_L`$ and $`\mathrm{\Delta }_{L,\mathrm{Reg}}`$ are manipulated so they can be plotted as a function of $`\sigma `$, the smearing width. We know $`G_L(\sigma )_{,i}=G_{L,i}^{\mathrm{div}}+G_{L,i}^{\mathrm{fin}}`$ and $`\rho =\rho ^{\mathrm{div}}+\rho ^{\mathrm{fin}}`$, where
$$\rho ^{\mathrm{div}}=\frac{1}{32\pi ^2\sigma ^4},G_{L,x_{}x_{}}^{\mathrm{div}}=\frac{1}{48\pi ^2\sigma ^4}=\frac{2}{3}\rho ^{\mathrm{div}}\mathrm{and}G_{L,zz}^{\mathrm{div}}=\frac{1}{96\pi ^2\sigma ^4}=\frac{1}{3}\rho ^{\mathrm{div}}$$
(B1)
We write $`G_L(\sigma )_{,i}=F_i(0)+2_{n=1}^{\mathrm{}}F_i(n)`$ and $`\rho _L(\sigma )=F_\rho (0)+2_{n=1}^{\mathrm{}}F_\rho (n)`$ with
$`F_1`$ $`=`$ $`{\displaystyle \frac{n}{8e^{8n^2\pi ^2\sigma ^2}\sigma ^2}}+{\displaystyle \frac{\mathrm{Erfc}\left(2\sqrt{2}n\pi \sigma \right)}{32\sqrt{2\pi }\sigma ^3}}{\displaystyle \frac{n^2\pi ^{\frac{3}{2}}\mathrm{Erfc}\left(2\sqrt{2}n\pi \sigma \right)}{2\sqrt{2}\sigma }}`$ (B2)
$`F_2`$ $`=`$ $`{\displaystyle \frac{n^2\pi ^{\frac{3}{2}}\mathrm{Erfc}\left(2\sqrt{2}n\pi \sigma \right)}{2\sqrt{2}\sigma }}`$ (B3)
$`F_\rho `$ $`=`$ $`{\displaystyle \frac{n}{8e^{8n^2\pi ^2\sigma ^2}\sigma ^2}}+{\displaystyle \frac{\mathrm{Erfc}\left(2\sqrt{2}n\pi \sigma \right)}{32\sqrt{2\pi }\sigma ^3}}`$ (B4)
The problem here is that each of the terms in the sum over $`n`$ diverge as $`\sigma 0`$. By defining $`\stackrel{~}{F}_i(n)=F_i(n)/X_i^{\mathrm{div}}`$ ($`i=x_{}x_{},zz,\rho `$ and $`X_i=G_{L,x_{}x_{}}^{\mathrm{div}},G_{L,zz}^{\mathrm{div}},\rho _L^{\mathrm{div}}`$) and $`\stackrel{~}{X}_i=2_{}^{}{}_{n=0}{}^{\mathrm{}}\stackrel{~}{F}_i(n)1`$, where $`^{}`$ has a factor of $`\frac{1}{2}`$ for $`n=0`$, then
$$X_i=X_i^{\mathrm{div}}\stackrel{~}{X}_i$$
(B5)
The defined functions are
$`\stackrel{~}{F}_{x_{}x_{}}(\sigma ,n)`$ $`=`$ $`{\displaystyle \frac{6n\pi ^2\sigma ^2}{e^{8n^2\pi ^2\sigma ^2}}}+{\displaystyle \frac{3\pi ^{\frac{3}{2}}\sigma \mathrm{Erfc}\left(2\sqrt{2}n\pi \sigma \right)}{2\sqrt{2}}}12\sqrt{2}n^2\pi ^{\frac{7}{2}}\sigma ^3\mathrm{Erfc}\left(2\sqrt{2}n\pi \sigma \right)`$ (B6)
$`\stackrel{~}{F}_{zz}(\sigma ,n)`$ $`=`$ $`24\sqrt{2}n^2\pi ^{\frac{7}{2}}\sigma ^3\mathrm{Erfc}\left(2\sqrt{2}n\pi \sigma \right)`$ (B7)
$`\stackrel{~}{F}_\rho (\sigma ,n)`$ $`=`$ $`{\displaystyle \frac{4n\pi ^2\sigma ^2}{e^{8n^2\pi ^2\sigma ^2}}}+{\displaystyle \frac{\pi ^{\frac{3}{2}}\sigma \mathrm{Erfc}\left(2\sqrt{2}n\pi \sigma \right)}{\sqrt{2}}}`$ (B8)
Each of these new functions to be summed over are now finite as $`\sigma 0`$. We have divided out the known divergent part.
Now we can turn to the smeared fluctuations of the energy density. First, for the sum of the divergent and cross terms we have
$$\mathrm{\Delta }\rho _L^{\mathrm{div}}+\mathrm{\Delta }\rho _L^{\mathrm{cross}}=\frac{2}{3}\left((\rho _L^{\mathrm{div}})^2+2\rho _L^{\mathrm{div}}\rho _L^{\mathrm{fin}}\right)=\frac{2}{3}\left(1+2\stackrel{~}{\rho }_L\right)\left(\rho _L^{\mathrm{div}}\right)^2$$
(B9)
We can clearly see how this has factored out the divergent coefficient. For the finite term:
$`\mathrm{\Delta }\rho _L^{\mathrm{fin}}={\displaystyle \frac{(\rho _L^{\mathrm{div}})^2}{9}}\left(3\stackrel{~}{G}_{L,x_{}x_{}}^2+2\stackrel{~}{G}_{L,x_{}x_{}}\stackrel{~}{G}_{L,zz}+\stackrel{~}{G}_{L,zz}^{}{}_{}{}^{2}\right)`$ (B10)
Taking together, the fluctuations can be written as
$$\mathrm{\Delta }\rho _L^2=\frac{(\rho _L^{\mathrm{div}})^2}{9}\left(6+3\stackrel{~}{G}_{L,x_{}x_{}}^{}{}_{}{}^{2}+2\stackrel{~}{G}_{L,x_{}x_{}}\stackrel{~}{G}_{L,zz}+\stackrel{~}{G}_{L,zz}^{}{}_{}{}^{2}+12\stackrel{~}{\rho }_L\right)$$
(B11)
From this, we get the dimensionless measure
$$\mathrm{\Delta }(\sigma ,L)=\frac{6+3\stackrel{~}{G}_{L,x_{}x_{}}^{}{}_{}{}^{2}+2\stackrel{~}{G}_{L,x_{}x_{}}\stackrel{~}{G}_{L,zz}+\stackrel{~}{G}_{L,zz}^{}{}_{}{}^{2}+12\stackrel{~}{\rho }_L}{15+3\stackrel{~}{G}_{L,x_{}x_{}}^{}{}_{}{}^{2}+2\stackrel{~}{G}_{L,x_{}x_{}}\stackrel{~}{G}_{L,zz}+\stackrel{~}{G}_{L,zz}^{}{}_{}{}^{2}+30\stackrel{~}{\rho }_L+9\stackrel{~}{\rho }_{L}^{}{}_{}{}^{2}}$$
(B12)
Also, considering just the finite terms
$$\mathrm{\Delta }_{L,\mathrm{Reg}}(\sigma ,L)=\frac{\mathrm{\Delta }\rho _{L,\mathrm{Reg}}}{\mathrm{\Delta }\rho _{L,\mathrm{Reg}}+\left(\rho _{L,\mathrm{Reg}}\right)^2}=\frac{3\stackrel{~}{G}_{L,x_{}x_{}}^{}{}_{}{}^{2}+2\stackrel{~}{G}_{L,x_{}x_{}}\stackrel{~}{G}_{L,zz}+\stackrel{~}{G}_{L,zz}^{}{}_{}{}^{2}}{3\stackrel{~}{G}_{L,x_{}x_{}}^{}{}_{}{}^{2}+2\stackrel{~}{G}_{L,x_{}x_{}}\stackrel{~}{G}_{L,zz}+\stackrel{~}{G}_{L,zz}^{}{}_{}{}^{2}+9\stackrel{~}{\rho }_{L}^{}{}_{}{}^{2}}$$
(B13)
We can now numerically evaluate the above ratios. Plots of $`\mathrm{\Delta }(\sigma ,L)`$ and $`\mathrm{\Delta }_{L,\mathrm{Reg}}(\sigma ,L)`$ as a function of $`\sigma /L`$ are presented in Figure 1. Figure 2 presents plots of $`\rho _L^{\mathrm{fin}}(\sigma ,L)`$ and $`\sqrt{\mathrm{\Delta }\rho _L^{\mathrm{fin}}(\sigma ,L)}`$.
We need to worry about the error for $`\sigma 1`$($`=L`$). Considering only the periodic $`z`$ direction, this was smeared with the function
$$f(z,\sigma )=\frac{1}{\sqrt{2\pi }\sigma }\mathrm{exp}\left(\frac{z^2}{2\sigma ^2}\right)$$
(B14)
For an error estimate, we use $`f(1,\sigma )`$, for as long as this is small, then the Gaussian smearing function does not detect the periodicity. At $`\sigma =0.4`$, this error is only 4%.
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# Spin Gap of 𝑆=1/2 Heisenberg Model on Distorted Diamond Chain
Recently, Ishii et al.$`^{\text{?}\text{)}}`$ measured the magnetic susceptibility $`\chi `$ for $`\mathrm{Cu}_3\mathrm{Cl}_6(\mathrm{H}_2\mathrm{O})_22\mathrm{H}_8\mathrm{C}_4\mathrm{SO}_2`$, which is considered to be a quasi-one-dimensional material consisting of $`S=\frac{1}{2}`$ trimer spin chains. The result indicates that $`\chi `$ vanishes in the low temperature limit. They also measured the magnetization process for this material, and showed that there is a plateau of zero magnetization below the critical field $`H_c3.9`$ T. From these experimental results, they concluded that the ground state is a singlet state with spin gap. The spin gap $`\mathrm{\Delta }`$ is estimated as $`\mathrm{\Delta }/k_\mathrm{B}5.2`$ K from the value of $`H_c`$.
The proposed Hamiltonian$`^{\text{?)}}`$ representing a spin chain in this material is given by
$`H`$ $`=`$ $`J_1{\displaystyle \underset{j}{}}\left(𝑺_{3j1}𝑺_{3j}+𝑺_{3j}𝑺_{3j+1}\right)`$ (1)
$`+`$ $`J_2{\displaystyle \underset{j}{}}𝑺_{3j+1}𝑺_{3j+2}`$
$`+`$ $`J_3{\displaystyle \underset{j}{}}\left(𝑺_{3j2}𝑺_{3j}+𝑺_{3j}𝑺_{3j+2}\right),`$
where $`𝑺_j`$ is the $`S=\frac{1}{2}`$ spin on site $`j`$. Three spins $`𝑺_{3j1}`$, $`𝑺_{3j}`$ and $`𝑺_{3j+1}`$ form a trimer. The lattice structure is shown in Fig. 1. Three kinds of exchange constants $`J_1`$, $`J_2`$ and $`J_3`$ are inferred to be positive and to satisfy the relation $`J_1>J_2,J_3`$ from the lattice parameters of the material .$`^{\text{?}\text{)}}`$ Hereafter, we use the unit of $`J_1=1`$.
The symmetric case of $`J_3=J_1(=1)`$ has been studied by Takano et al.$`^{\text{?}\text{)}}`$ and the system has been called the diamond chain. They almost exactly showed that there exist three phases in the parameter space; the ferrimagnetic phase for $`J_2<0.909`$, the tetramer-dimer (TD) singlet phase for $`0.909<J_2<2`$ and the dimer-monomer (DM) singlet phase for $`J_2>2`$. The TD phase is a disordered phase with spin gap which originates from frustration among exchange interactions, while the DM phase is a spin fluid phase without spin gap. Okamoto et al.$`^{\text{?}\text{)}}`$ studied the general case of $`J_31`$; i. e. the distorted diamond chain. The three phases develop in the $`J_2`$-$`J_3`$ plane. They numerically determined the phase boundaries. Also Tonegawa et al.$`^{\text{?}\text{)}}`$ numerically studied the magnetization process and showed plateaux for $`\frac{1}{3}`$ and $`\frac{2}{3}`$ of the saturation field.
In this article, we estimate the values of the spin gap by the numerical diagonalization. Then we produce a contour map in the $`J_2`$-$`J_3`$ parameter space. The contour map represents an overall feature of the gapped phase of the $`S=\frac{1}{2}`$ Heisenberg model on the distorted diamond chain. When further experimental information on $`\mathrm{Cu}_3\mathrm{Cl}_6(\mathrm{H}_2\mathrm{O})_22\mathrm{H}_8\mathrm{C}_4\mathrm{SO}_2`$ is given, the contour map will be useful to determine the values of the exchange constants for the real material.
We first calculate the spin gap $`\mathrm{\Delta }_L`$ for finite chains with system size $`L`$. The spin gap $`\mathrm{\Delta }_{\mathrm{}}`$ in the thermodynamic limit is evaluated by extrapolation. We assume the size dependence of $`\mathrm{\Delta }_L`$ as
$$\mathrm{\Delta }_L=\mathrm{\Delta }_{\mathrm{}}+\frac{c_1}{L}+\frac{c_2}{L^2}$$
(2)
with constants $`c_1`$ and $`c_2`$. The numerical diagonalization has been done for $`L=`$12, 18 and 24 under the periodic boundary condition. We determine $`c_1`$, $`c_2`$ and $`\mathrm{\Delta }_{\mathrm{}}`$ by fitting. In Fig. 2, we show $`\mathrm{\Delta }_L`$ as a function of $`L`$ for several values of $`J_3`$ at $`J_2=1`$. For $`J_3=0`$, the estimated value of $`\mathrm{\Delta }_{\mathrm{}}`$ is about 0.002 and is close to zero; the nonzero value is interpreted as an extrapolation error .$`^{\text{?}\text{)}}`$ For $`0<J_3\stackrel{<}{_{}}0.4`$, the true value of the spin gap is very small or may be regarded as zero, since the estimated values are less than 0.002 and are within the extrapolation error. For $`J_3\stackrel{>}{_{}}0.6`$, the figure shows that the system has a finite spin gap. For $`J_3=0.5`$, $`\mathrm{\Delta }_{\mathrm{}}`$ is 0.0034, which is small but seems to be finite. This is consistent with the result of Okamoto et al. that the spin gap opens at the critical value $`J_3^c0.35`$ for $`J_2=1`$ .$`^{\text{?)}}`$
In general, the spin gap in a dimer phase is exponentially small near the phase boundary to a spin fluid phase. Hence it is difficult to estimate $`\mathrm{\Delta }_{\mathrm{}}`$ near the boundary in the present case. To overcome this difficulty, we assume that $`J_3`$ dependence of the spin gap is given by
$$D(J_3)=a_1\sqrt{J_3J_3^c}\mathrm{exp}\left(\frac{a_2}{J_3J_3^c}\right)$$
(3)
for $`J_3J_3^c`$,$`^{\text{?}\text{)}}`$ where $`a_1`$ and $`a_2`$ are constants. We have the values of $`J_3^c`$ by inspecting the phase diagram of Okamoto et al. ;$`^{\text{?)}}`$ e. g. $`J_3^c`$ = 0.374, 0.354 and 0.460 for $`J_2`$ = 0.7, 1.0 and 1.4, respectively. We carry out the fitting of the extrapolation data $`\mathrm{\Delta }_{\mathrm{}}`$ to eq. (3) and determine $`a_1`$ and $`a_2`$. Figure 3 represents the fitting function $`D(J_3)`$ and the extrapolation data. We find that the extrapolation data are well reproduced by eq. (3) for $`\mathrm{\Delta }_{\mathrm{}}\stackrel{>}{_{}}0.02`$. Hence the function form in eq. (3) is reliable. We use eq. (3) to estimate the spin gap for $`\mathrm{\Delta }_{\mathrm{}}\stackrel{<}{_{}}0.02`$ near the critical value $`J_3^c`$. For example, the spin gap is estimated as $`1.0\times 10^2`$ at $`J_30.57`$, $`1.0\times 10^3`$ at $`J_30.50`$, $`1.0\times 10^4`$ at $`J_30.47`$ and $`1.0\times 10^5`$ at $`J_30.45`$ for $`J_2=1.0`$.
Using these results, we draw contour lines of the spin gap in the $`J_2`$-$`J_3`$ plain. The resultant contour map is shown in Fig. 4. We have calculated $`\mathrm{\Delta }_{\mathrm{}}`$ at the discrete positions ($`J_2`$, $`J_3`$) with $`J_2`$ = 0.7, 0.8, …, 2.0 and $`J_3`$ = 0.5, 0.55, …, 1.0. For $`\mathrm{\Delta }_{\mathrm{}}>0.02`$, the positions of solid circles are determined by the linear interpolation among the spin gaps $`\mathrm{\Delta }_{\mathrm{}}`$ at the discrete positions. For $`\mathrm{\Delta }_{\mathrm{}}<0.02`$, the positions of open circles are determined by using $`D(J_3)`$ (eq. (3)) instead of $`\mathrm{\Delta }_{\mathrm{}}`$.
The temperature dependence of the experimental magnetic susceptibility has a broad peak at $``$70 K. It suggests that the energy scale of the characteristic exchange constant $`J_1`$ is larger than 70 K. Here we consider a case of $`J_1`$ being 100 K as an example.$`^{\text{?}\text{)}}`$ In this case, we have $`\mathrm{\Delta }_{\mathrm{}}0.05`$ according to the observed spin gap $``$5 K. Then $`J_2`$ and $`J_3`$ are limited to values close to the contour line of $`\mathrm{\Delta }_{\mathrm{}}=0.05`$ and of $`J_2<1`$.
One of the authors (K. T.) would like to thank H. Tanaka and M. Ishii for explaining their experimental results, and K. Okamoto for discussion. This work is partially supported by the Grant-in-Aid for Scientific Research from the Ministry of Education, Science, Sports and Culture, Japan.
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# Quantum Amplitude Amplification and Estimation
## 1 Introduction
Quantum computing is a field at the junction of theoretical modern physics and theoretical computer science. Practical experiments involving a few quantum bits have been successfully performed, and much progress has been achieved in quantum information theory, quantum error correction and fault tolerant quantum computation. Although we are still far from having desktop quantum computers in our offices, the quantum computational paradigm could soon be more than mere theoretical exercise.
The discovery by Peter Shor of a polynomial-time quantum algorithm for factoring and computing discrete logarithms was a major milestone in the history of quantum computing. Another significant result is Lov Grover’s quantum search algorithm . Grover’s algorithm does not solve NP–complete problems in polynomial time, but the wide range of its applications more than compensates for this.
In this paper, we generalize Grover’s algorithm in a variety of directions. Consider a problem that is characterized by a Boolean function $`\chi (x,y)`$ in the sense that $`y`$ is a good solution to instance $`x`$ if and only if $`\chi (x,y)=1`$. (There could be more than one good solution to a given instance.) If we have a probabilistic algorithm $`𝒫`$ that outputs a guess $`𝒫(x)`$ on input $`x`$, we can call $`𝒫`$ and $`\chi `$ repeatedly until a solution to instance $`x`$ is found. If $`\chi (x,𝒫(x))=1`$ with probability $`p_x>0`$, we expect to repeat this process $`1/p_x`$ times on the average. Consider now the case when we have a quantum algorithm $`𝒜`$ instead of the probabilistic algorithm. Assume $`𝒜`$ makes no measurements: instead of a classical answer, it produces quantum superposition $`|\mathrm{\Psi }_x`$ when run on input $`x`$. Let $`a_x`$ denote the probability that $`|\mathrm{\Psi }_x`$, if measured, would be a good solution. If we repeat the process of running $`𝒜`$ on $`x`$, measuring the output, and using $`\chi `$ to check the validity of the result, we shall expect to repeat $`1/a_x`$ times on the average before a solution is found. This is no better than the classical probabilistic paradigm.
In Section 2, we describe a more efficient approach to this problem, which we call amplitude amplification. Intuitively, the probabilistic paradigm increases the probability of success roughly by a constant on each iteration; by contrast, amplitude amplification increases the amplitude of success roughly by a constant on each iteration. Because amplitudes correspond to square roots of probabilities, it suffices to repeat the amplitude amplification process approximately $`1/\sqrt{a_x}`$ times to achieve success with overwhelming probability. For simplicity, we assume in the rest of this paper that there is a single instance for which we seek a good solution, which allows us to dispense with input $`x`$, but the generalization to the paradigm outlined above is straightforward. Grover’s original database searching quantum algorithm is a special case of this process, in which $`\chi `$ is given by a function $`f:\{0,1,\mathrm{},N1\}\{0,1\}`$ for which we are promised that there exists a unique $`x_0`$ such that $`f(x_0)=1`$. If we use the Fourier transform as quantum algorithm $`𝒜`$—or more simply the Walsh–Hadamard transform in case $`N`$ is a power of 2—an equal superposition of all possible $`x`$’s is produced, whose success probability would be $`1/N`$ if measured. Classical repetition would succeed after an expected number $`N`$ of evaluations of $`f`$. Amplitude amplification corresponds to Grover’s algorithm: it succeeds after approximately $`\sqrt{N}`$ evaluations of the function.
We generalize this result further to the case when the probability of success $`a`$ of algorithm $`𝒜`$ is not known ahead of time: it remains sufficient to evaluate $`𝒜`$ and $`\chi `$ an expected number of times that is proportional to $`1/\sqrt{a}`$. Moreover, in the case $`a`$ is known ahead of time, we give two different techniques that are guaranteed to find a good solution after a number of iterations that is proportional to $`1/\sqrt{a}`$ in the worst case.
It can be proven that Grover’s algorithm goes quadratically faster than any possible classical algorithm when function $`f`$ is given as a black box. However, it is usually the case in practice that information is known about $`f`$ that allows us to solve the problem much more efficiently than by exhaustive search. The use of classical heuristics, in particular, will often yield a solution significantly more efficiently than straight quantum amplitude amplification would. In Section 3, we consider a broad class of classical heuristics and show how to apply amplitude amplification to obtain quadratic speedup compared to any such heuristic.
Finally, Section 4 addresses the question of estimating the success probability $`a`$ of quantum algorithm $`𝒜`$. We call this process amplitude estimation. As a special case of our main result (Theorem 12), an estimate for $`a`$ is obtained after any number $`M`$ of iterations which is within $`2\pi \sqrt{a(1a)}/M+\pi ^2/M^2`$ of the correct value with probability at least $`8/\pi ^2`$, where one iteration consists of running algorithm $`𝒜`$ once forwards and once backwards, and of computing function $`\chi `$ once. As an application of this technique, we show how to approximately count the number of $`x`$ such that $`f(x)=1`$ given a function $`f:\{0,1,\mathrm{},N1\}\{0,1\}`$. If the correct answer is $`t>0`$, it suffices to compute the function $`\sqrt{N}`$ times to obtain an estimate roughly within $`\sqrt{t}`$ of the correct answer. A number of evaluations of $`f`$ proportional to $`\frac{1}{\epsilon }\sqrt{N/t}`$ yields a result that is likely to be within $`\epsilon t`$ of the correct answer. (We can do slightly better in case $`\epsilon `$ is not fixed.) If it is known ahead of time that the correct answer is either $`t=0`$ or $`t=t_0`$ for some fixed $`t_0`$, we can determine which is the case with certainty using a number of evaluations of $`f`$ proportional to $`\sqrt{N/t_0}`$. If we have no prior knowledge about $`t`$, the exact count can be obtained with high probability after a number of evaluations of $`f`$ that is proportional to $`\sqrt{t(Nt)}`$ when $`0<t<N`$ and $`\sqrt{N}`$ otherwise. Most of these results are optimal.
We assume in this paper that the reader is familiar with basic notions of quantum computing.
## 2 Quantum amplitude amplification
Suppose we have a classical randomized algorithm that succeeds with some probability $`p`$. If we repeat the algorithm, say, $`j`$ times, then our probability of success increases to roughly $`jp`$ (assuming $`jp1`$). Intuitively, we can think of this strategy as each additional run of the given algorithm boosting the probability of success by an additive amount of roughly $`p`$.
A quantum analogue of boosting the probability of success would be to boost the amplitude of being in a certain subspace of a Hilbert space. The general concept of amplifying the amplitude of a subspace was discovered by Brassard and Høyer as a generalization of the boosting technique applied by Grover in his original quantum searching paper . Following and , we refer to their idea as amplitude amplification and detail the ingredients below.
Let $``$ denote the Hilbert space representing the state space of a quantum system. Every Boolean function $`\chi :\{0,1\}`$ induces a partition of $``$ into a direct sum of two subspaces, a good subspace and a bad subspace. The good subspace is the subspace spanned by the set of basis states $`|x`$ for which $`\chi (x)=1`$, and the bad subspace is its orthogonal complement in $``$. We say that the elements of the good subspace are good, and that the elements of the bad subspace are bad.
Every pure state $`|\mathrm{{\rm Y}}`$ in $``$ has a unique decomposition as $`|\mathrm{{\rm Y}}=|\mathrm{{\rm Y}}_1+|\mathrm{{\rm Y}}_0`$, where $`|\mathrm{{\rm Y}}_1`$ denotes the projection onto the good subspace, and $`|\mathrm{{\rm Y}}_0`$ denotes the projection onto the bad subspace. Let $`a_\mathrm{{\rm Y}}=\mathrm{{\rm Y}}_1|\mathrm{{\rm Y}}_1`$ denote the probability that measuring $`|\mathrm{{\rm Y}}`$ produces a good state, and similarly, let $`b_\mathrm{{\rm Y}}=\mathrm{{\rm Y}}_0|\mathrm{{\rm Y}}_0`$. Since $`|\mathrm{{\rm Y}}_1`$ and $`|\mathrm{{\rm Y}}_0`$ are orthogonal, we have $`a_\mathrm{{\rm Y}}+b_\mathrm{{\rm Y}}=1`$.
Let $`𝒜`$ be any quantum algorithm that acts on $``$ and uses no measurements. Let $`|\mathrm{\Psi }=𝒜|0`$ denote the state obtained by applying $`𝒜`$ to the initial zero state. The amplification process is realized by repeatedly applying the following unitary operator on the state $`|\mathrm{\Psi }`$,
$$𝐐=𝐐(𝒜,\chi )=𝒜𝐒_0𝒜^1𝐒_\chi .$$
(1)
Here, the operator $`𝐒_\chi `$ conditionally changes the sign of the amplitudes of the good states,
$$|x\{\begin{array}{cc}|x\hfill & \text{if }\chi (x)=1\hfill \\ |x\hfill & \text{if }\chi (x)=0\text{,}\hfill \end{array}$$
while the operator $`𝐒_0`$ changes the sign of the amplitude if and only if the state is the zero state $`|0`$. The operator $`𝐐`$ is well-defined since we assume that $`𝒜`$ uses no measurements and, therefore, $`𝒜`$ has an inverse.
The usefulness of operator $`𝐐`$ stems from its simple action on the subspace $`_\mathrm{\Psi }`$ spanned by the vectors $`|\mathrm{\Psi }_1`$ and $`|\mathrm{\Psi }_0`$.
###### Lemma 1
We have that
$`𝐐|\mathrm{\Psi }_1`$ $`=(12\mathrm{a})|\mathrm{\Psi }_12\mathrm{a}|\mathrm{\Psi }_0`$
$`𝐐|\mathrm{\Psi }_0`$ $`=\mathrm{\hspace{0.17em}2}(1a)|\mathrm{\Psi }_1+(12a)|\mathrm{\Psi }_0,`$
where $`a=\mathrm{\Psi }_1|\mathrm{\Psi }_1`$.
It follows that the subspace $`_\mathrm{\Psi }`$ is stable under the action of $`𝐐`$, a property that was first observed by Brassard and Høyer and rediscovered by Grover .
Suppose $`0<a<1`$. Then $`_\mathrm{\Psi }`$ is a subspace of dimension 2, and otherwise $`_\mathrm{\Psi }`$ has dimension 1. The action of $`𝐐`$ on $`_\mathrm{\Psi }`$ is also realized by the operator
$$𝐔_\mathrm{\Psi }𝐔_{\mathrm{\Psi }_0},$$
(2)
which is composed of 2 reflections. The first operator, $`𝐔_{\mathrm{\Psi }_0}=𝐈\frac{2}{1a}|\mathrm{\Psi }_0\mathrm{\Psi }_0|`$, implements a reflection through the ray spanned by the vector $`|\mathrm{\Psi }_0`$, while the second operator $`𝐔_\mathrm{\Psi }=𝐈2|\mathrm{\Psi }\mathrm{\Psi }|`$ implements a reflection through the ray spanned by the vector $`|\mathrm{\Psi }`$.
Consider the orthogonal complement $`_\mathrm{\Psi }^{}`$ of $`_\mathrm{\Psi }`$ in $``$. Since the operator $`𝒜𝐒_0𝒜^1`$ acts as the identity on $`_\mathrm{\Psi }^{}`$, operator $`𝐐`$ acts as $`𝐒_\chi `$ on $`_\mathrm{\Psi }^{}`$. Thus, $`𝐐^2`$ acts as the identity on $`_\mathrm{\Psi }^{}`$, and every eigenvector of $`𝐐`$ in $`_\mathrm{\Psi }^{}`$ has eigenvalue $`+1`$ or $`1`$. It follows that to understand the action of $`𝐐`$ on an arbitrary initial vector $`|\mathrm{{\rm Y}}`$ in $``$, it suffices to consider the action of $`𝐐`$ on the projection of $`|\mathrm{{\rm Y}}`$ onto $`_\mathrm{\Psi }`$.
Since operator $`𝐐`$ is unitary, the subspace $`_\mathrm{\Psi }`$ has an orthonormal basis consisting of two eigenvectors of $`𝐐`$,
$$|\mathrm{\Psi }_\pm =\frac{1}{\sqrt{2}}\left(\frac{1}{\sqrt{a}}|\mathrm{\Psi }_1\pm \frac{ı}{\sqrt{1a}}|\mathrm{\Psi }_0\right),$$
(3)
provided $`0<a<1`$, where $`ı=\sqrt{1}`$ denotes the principal square root of $`1`$. The corresponding eigenvalues are
$$\lambda _\pm =e^{\pm ı2\theta _a},$$
(4)
where the angle $`\theta _a`$ is defined so that
$$\mathrm{sin}^2(\theta _a)=a=\mathrm{\Psi }_1|\mathrm{\Psi }_1$$
(5)
and $`0\theta _a\pi /2`$.
We use operator $`𝐐`$ to boost the success probability $`a`$ of the quantum algorithm $`𝒜`$. First, express $`|\mathrm{\Psi }=𝒜|0`$ in the eigenvector basis,
$$𝒜|0=|\mathrm{\Psi }=\frac{ı}{\sqrt{2}}\left(e^{ı\theta _a}|\mathrm{\Psi }_+e^{ı\theta _a}|\mathrm{\Psi }_{}\right).$$
(6)
It is now immediate that after $`j`$ applications of operator $`𝐐`$, the state is
$`𝐐^j|\mathrm{\Psi }`$ $`={\displaystyle \frac{ı}{\sqrt{2}}}\left(e^{(2j+1)ı\theta _a}|\mathrm{\Psi }_+e^{(2j+1)ı\theta _a}|\mathrm{\Psi }_{}\right)`$ (7)
$`={\displaystyle \frac{1}{\sqrt{a}}}\mathrm{sin}((2j+1)\theta _a)|\mathrm{\Psi }_1+{\displaystyle \frac{1}{\sqrt{1a}}}\mathrm{cos}((2j+1)\theta _a)|\mathrm{\Psi }_0.`$ (8)
It follows that if $`0<a<1`$ and if we compute $`𝐐^m|\mathrm{\Psi }`$ for some integer $`m0`$, then a final measurement will produce a good state with probability equal to $`\mathrm{sin}^2((2m+1)\theta _a)`$.
If the initial success probability $`a`$ is either 0 or 1, then the subspace $`_\mathrm{\Psi }`$ spanned by $`|\mathrm{\Psi }_1`$ and $`|\mathrm{\Psi }_0`$ has dimension 1 only, but the conclusion remains the same: If we measure the system after $`m`$ rounds of amplitude amplification, then the outcome is good with probability $`\mathrm{sin}^2((2m+1)\theta _a)`$, where the angle $`\theta _a`$ is defined so that Equation 5 is satisfied and so that $`0\theta _a\pi /2`$.
Therefore, assuming $`a>0`$, to obtain a high probability of success, we want to choose integer $`m`$ such that $`\mathrm{sin}^2((2m+1)\theta _a)`$ is close to 1. Unfortunately, our ability to choose $`m`$ wisely depends on our knowledge about $`\theta _a`$, which itself depends on $`a`$. The two extreme cases are when we know the exact value of $`a`$, and when we have no prior knowledge about $`a`$ whatsoever.
Suppose the value of $`a`$ is known. If $`a>0`$, then by letting $`m=\pi /4\theta _a`$, we have that $`\mathrm{sin}^2((2m+1)\theta _a)1a`$, as shown in . The next theorem is immediate.
###### Theorem 2 (Quadratic speedup)
Let $`𝒜`$ be any quantum algorithm that uses no measurements, and let $`\chi :\{0,1\}`$ be any Boolean function. Let $`a`$ the initial success probability of $`𝒜`$. Suppose $`a>0`$, and set $`m=\pi /4\theta _a`$, where $`\theta _a`$ is defined so that $`\mathrm{sin}^2(\theta _a)=a`$ and $`0<\theta _a\pi /2`$. Then, if we compute $`𝐐^m𝒜|0`$ and measure the system, the outcome is good with probability at least $`\mathrm{max}(1a,a)`$.
Note that any implementation of algorithm $`𝐐^m𝒜|0`$ requires that the value of $`a`$ is known so that the value of $`m`$ can be computed. We refer to Theorem 2 as a quadratic speedup, or the square-root running-time result. The reason for this is that if an algorithm $`𝒜`$ has success probability $`a>0`$, then after an expected number of $`1/a`$ applications of $`𝒜`$, we will find a good solution. Applying the above theorem reduces this to an expected number of at most $`(2m+1)/\mathrm{max}(1a,a)\mathrm{\Theta }(\frac{1}{\sqrt{a}})`$ applications of $`𝒜`$ and $`𝒜^1`$.
As an application of Theorem 2, consider the search problem in which we are given a Boolean function $`f:\{0,1,\mathrm{},N1\}\{0,1\}`$ satisfying the promise that there exists a unique $`x_0\{0,1,\mathrm{},N1\}`$ on which $`f`$ takes value 1, and we are asked to find $`x_0`$. If $`f`$ is given as a black box, then on a classical computer, we need to evaluate $`f`$ on an expected number of roughly half the elements of the domain in order to determine $`x_0`$.
By contrast, Grover discovered a quantum algorithm that only requires an expected number of evaluations of $`f`$ in the order of $`\sqrt{N}`$. In terms of amplitude amplification, Grover’s algorithm reads as follows: Let $`\chi =f`$, and let $`𝒜=𝐖`$ be the Walsh-Hadamard transform on $`n`$ qubits that maps the initial zero state $`|0`$ to $`\frac{1}{\sqrt{N}}_{x=0}^{N1}|x`$, an equally-weighted superposition of all $`N=2^n`$ elements in the domain of $`f`$. Then the operator $`𝐐=𝒜𝐒_0𝒜^1𝐒_\chi `$ is equal to the iterate $`\mathrm{𝐖𝐒}_0\mathrm{𝐖𝐒}_f`$ applied by Grover in his searching paper . The initial success probability $`a`$ of $`𝒜`$ is exactly $`1/N`$, and if we measure after $`m=\pi /4\theta _a`$ iterations of $`𝐐`$, the probability of measuring $`x_0`$ is lower bounded by $`11/N`$ .
Now, suppose that the value of $`a`$ is not known. In Section 4, we discuss techniques for finding an estimate of $`a`$, whereafter one then can apply a weakened version of Theorem 2 in which the exact value of $`a`$ is replaced by an estimate of it. Another idea is to try to find a good solution without prior computation of an estimate of $`a`$. Within that approach, by adapting the ideas in Section 6 in we can still obtain a quadratic speedup.
###### Theorem 3 (Quadratic speedup without knowing $`a`$)
There exists a quantum algorithm QSearch with the following property. Let $`𝒜`$ be any quantum algorithm that uses no measurements, and let $`\chi :\{0,1\}`$ be any Boolean function. Let $`a`$ denote the initial success probability of $`𝒜`$. Algorithm QSearch finds a good solution using an expected number of applications of $`𝒜`$ and $`𝒜^1`$ which are in $`\mathrm{\Theta }(\frac{1}{\sqrt{a}})`$ if $`a>0`$, and otherwise runs forever.
The algorithm in the above theorem utilizes the given quantum algorithm $`𝒜`$ as a subroutine and the operator $`𝐐`$. The complete algorithm is as follows:
* Algorithm( $`\text{QSearch}(𝒜,\chi )`$ )
1. Set $`l=0`$ and let $`c`$ be any constant such that $`1<c<2`$.
2. Increase $`l`$ by 1 and set $`M=c^l`$.
3. Apply $`𝒜`$ on the initial state $`|0`$, and measure the system. If the outcome $`|z`$ is good, that is, if $`\chi (z)=1`$, then output $`z`$ and stop.
4. Initialize a register of appropriate size to the state $`𝒜|0`$.
5. Pick an integer $`j`$ between 1 and $`M`$ uniformly at random.
6. Apply $`𝐐^j`$ to the register, where $`𝐐=𝐐(𝒜,\chi )`$.
7. Measure the register. If the outcome $`|z`$ is good, then output $`z`$ and stop. Otherwise, go to step 2.
The intuition behind this algorithm is as follows. In a 2-dimensional real vector space, if we pick a unit vector $`(x,y)=(\mathrm{cos}(),\mathrm{sin}())`$ uniformly at random then the expected value of $`y^2`$ is $`1/2`$. Consider Equation 8. If we pick $`j`$ at random between 1 and $`M`$ for some integer $`M`$ such that $`M\theta _a`$ is larger than, say, $`100\pi `$, then we have a good approximation to a random unit vector, and we will succeed with probability close to $`1/2`$.
To turn this intuition into an algorithm, the only obstacle left is that we do not know the value of $`\theta _a`$, and hence do not know an appropriate value for $`M`$. However, we can overcome this by using exponentially increasing values of $`M`$, an idea similar to the one used in “exponential searching” (which is a term that does not refer to the running time of the method, but rather to an exponentially increasing growth of the size of the search space).
The correctness of algorithm QSearch is immediate and thus to prove the theorem, it suffices to show that the expected number of applications of $`𝒜`$ and $`𝒜^1`$ is in the order of $`1/\sqrt{a}`$. This can be proven by essentially the same techniques applied in the proof of Theorem 3 in and we therefore only give a very brief sketch of the proof.
On the one hand, if the initial success probability $`a`$ is at least $`3/4`$, then step 3 ensures that we soon will measure a good solution. On the other hand, if $`0<a<3/4`$ then, for any given value of $`M`$, the probability of measuring a good solution in step 7 is lower bounded by
$$\frac{1}{2}\left(1\frac{1}{2M\sqrt{a}}\right).$$
(9)
Let $`c_0>0`$ be such that $`c=2(1c_0)`$ and let $`M_0=1/(2c_0\sqrt{a})`$. The expected number of applications of $`𝒜`$ is upper bounded by $`T_1+T_2`$, where $`T_1`$ denotes the maximum number of applications of $`𝒜`$ the algorithm uses while $`M<M_0`$, and where $`T_2`$ denotes the expected number of applications of $`𝒜`$ the algorithm uses while $`MM_0`$. Clearly $`T_1O(M_0)=O(\frac{1}{\sqrt{a}})`$ and we now show that $`T_2O(\frac{1}{\sqrt{a}})`$ as well.
For all $`MM_0`$, the measurement in step 7 yields a good solution with probability at least $`\frac{1}{2}(1c_0)`$, and hence it fails to yield a good solution with probability at most $`p_0=\frac{1}{2}(1+c_0)`$. Thus for all $`i0`$, with probability at most $`p_0^i`$, we have that $`MM_0c^i`$ at some point after step 2 while running the algorithm. Hence $`T_2`$ is at most on the order of $`_{i0}M_0(cp_0)^i`$ which is in $`O(M_0)`$ since $`cp_0<1`$. The total expected number of applications of $`𝒜`$ is thus in $`O(M_0)`$, which is $`O(\frac{1}{\sqrt{a}})`$.
For the lower bound, if $`M`$ were in $`o\left(\frac{1}{\sqrt{a}}\right)`$, then the probability that we measure a good solution in step 7 would be vanishingly small. This completes our sketch of the proof of Theorem 3.
### 2.1 Quantum de-randomization when the success <br>probability is known
We now consider the situation where the success probability $`a`$ of the quantum algorithm $`𝒜`$ is known. If $`a=0`$ or $`a=1`$, then amplitude amplification will not change the success probability, so in the rest of this section, we assume that $`0<a<1`$. Theorem 2 allows us to boost the probability of success to at least $`\mathrm{max}(1a,a)`$. A natural question to ask is whether it is possible to improve this to certainty, still given the value of $`a`$. It turns out that the answer is positive. This is unlike classical computers, where no such general de-randomization technique is known. We now describe 2 optimal methods for obtaining this, but other approaches are possible.
The first method is by applying amplitude amplification, not on the original algorithm $`𝒜`$, but on a slightly modified version of it. By Equation 8, if we measure the state $`𝐐^m𝒜|0`$, then the outcome is good with probability $`\mathrm{sin}^2((2m+1)\theta _a)`$. In particular, if $`\stackrel{~}{m}=\pi /4\theta _a1/2`$ happens to be an integer, then we would succeed with certainty after $`\stackrel{~}{m}`$ applications of $`𝐐`$. In general, $`\overline{m}=\stackrel{~}{m}`$ iterations is a fraction of 1 iteration too many, but we can compensate for that by choosing $`\overline{\theta _a}=\pi /(4\overline{m}+2)`$, an angle slightly smaller than $`\theta _a`$. Any quantum algorithm that succeeds with probability $`\overline{a}`$ such that $`\mathrm{sin}^2(\overline{\theta _a})=\overline{a}`$, will succeed with certainty after $`\overline{m}`$ iterations of amplitude amplification. Given $`𝒜`$ and its initial success probability $`a`$, it is easy to construct a new quantum algorithm that succeeds with probability $`\overline{a}a`$: Let $``$ denote the quantum algorithm that takes a single qubit in the initial state $`|0`$ and rotates it to the superposition $`\sqrt{1\overline{a}/a}|0+\sqrt{\overline{a}/a}|1`$. Apply both $`𝒜`$ and $``$, and define a good solution as one in which $`𝒜`$ produces a good solution, and the outcome of $``$ is the state $`|1`$. Theorem 4 follows.
###### Theorem 4 (Quadratic speedup with known $`a`$)
Let $`𝒜`$ be any quantum algorithm that uses no measurements, and let $`\chi :\{0,1\}`$ be any Boolean function. There exists a quantum algorithm that given the initial success probability $`a>0`$ of $`𝒜`$, finds a good solution with certainty using a number of applications of $`𝒜`$ and $`𝒜^1`$ which is in $`\mathrm{\Theta }(\frac{1}{\sqrt{a}})`$ in the worst case.
The second method to obtain success probability 1 requires a generalization of operator $`𝐐`$. Given angles $`0\varphi ,\phi <2\pi `$, redefine $`𝐐`$ as follows,
$$𝐐=𝐐(𝒜,\chi ,\varphi ,\phi )=𝒜𝐒_0(\varphi )𝒜^1𝐒_\chi (\phi ).$$
(10)
Here, the operator $`𝐒_\chi (\phi )`$ is the natural generalization of the $`𝐒_\chi `$ operator,
$$|x\{\begin{array}{cc}e^{ı\phi }|x\hfill & \text{if }\chi (x)=1\hfill \\ |x\hfill & \text{if }\chi (x)=0\text{.}\hfill \end{array}$$
Similarly, the operator $`𝐒_0(\varphi )`$ multiplies the amplitude by a factor of $`e^{ı\varphi }`$ if and only if the state is the zero state $`|0`$. The action of operator $`𝐐(𝒜,\chi ,\varphi ,\phi )`$ is also realized by applying an operator that is composed of two pseudo-reflections: the operator $`𝒜𝐒_0(\varphi )𝒜^1`$ and the operator $`𝐒_\chi (\phi )`$.
The next lemma shows that the subspace $`_\mathrm{\Psi }`$ spanned by $`|\mathrm{\Psi }_1`$ and $`|\mathrm{\Psi }_0`$ is stable under the action of $`𝐐`$, just as in the special case $`𝐐(𝒜,\chi ,\pi ,\pi )`$ studied above.
###### Lemma 5
Let $`𝐐=𝐐(𝒜,\chi ,\varphi ,\phi )`$. Then
$`𝐐|\mathrm{\Psi }_1`$ $`=e^{ı\phi }((1e^{ı\varphi })a1)|\mathrm{\Psi }_1+\mathrm{e}^{\mathrm{ı}\mathrm{\phi }}(1\mathrm{e}^{\mathrm{ı}\mathrm{\varphi }})\mathrm{a}|\mathrm{\Psi }_0`$
$`𝐐|\mathrm{\Psi }_0`$ $`=(1\mathrm{e}^{\mathrm{ı}\mathrm{\varphi }})(1\mathrm{a})|\mathrm{\Psi }_1((1e^{ı\varphi })a+e^{ı\varphi })|\mathrm{\Psi }_0,`$
where $`a=\mathrm{\Psi }_1|\mathrm{\Psi }_1`$.
Let $`\stackrel{~}{m}=\pi /4\theta _a1/2`$, and suppose that $`\stackrel{~}{m}`$ is not an integer. In the second method to obtain a good solution with certainty, we also apply $`\stackrel{~}{m}`$ iterations of amplitude amplification, but now we slow down the speed of the very last iteration only, as opposed to of all iterations as in the first method. For the case $`\stackrel{~}{m}<1`$, this second method has also been suggested by Chi and Kim . We start by applying the operator $`𝐐(𝒜,\chi ,\varphi ,\phi )`$ with $`\varphi =\phi =\pi `$ a number of $`\stackrel{~}{m}`$ times to the initial state $`|\mathrm{\Psi }=𝒜|0`$. By Equation 8, this produces the superposition
$$\frac{1}{\sqrt{a}}\mathrm{sin}\left((2\stackrel{~}{m}+1)\theta _a\right)|\mathrm{\Psi }_1+\frac{1}{\sqrt{1a}}\mathrm{cos}\left((2\stackrel{~}{m}+1)\theta _a\right)|\mathrm{\Psi }_0.$$
Then, we apply operator $`𝐐`$ one more time, but now using angles $`\varphi `$ and $`\phi `$, both between $`0`$ and $`2\pi `$, satisfying
$$\begin{array}{c}e^{ı\phi }(1e^{ı\varphi })\sqrt{a}\mathrm{sin}\left((2\stackrel{~}{m}+1)\theta _a\right)\hfill \\ \hfill =((1e^{ı\varphi })a+e^{ı\varphi })\frac{1}{\sqrt{1a}}\mathrm{cos}\left((2\stackrel{~}{m}+1)\theta _a\right).\end{array}$$
(11)
By Lemma 5, this ensures that the resulting superposition has inner product zero with $`|\mathrm{\Psi }_0`$, and thus a subsequent measurement will yield a good solution with certainty.
The problem of choosing $`\varphi ,\phi `$ such that Equation 11 holds is equivalent to requiring that
$$\mathrm{cot}\left((2\stackrel{~}{m}+1)\theta _a\right)=e^{ı\phi }\mathrm{sin}(2\theta _a)\left(\mathrm{cos}(2\theta _a)+ı\mathrm{cot}(\varphi /2)\right)^1.$$
(12)
By appropriate choices of $`\varphi `$ and $`\phi `$, the right hand side of Equation 12 can be made equal to any nonzero complex number of norm at most $`\mathrm{tan}(2\theta _a)`$. Thus, since the left hand side of this equation is equal to some real number smaller than $`\mathrm{tan}(2\theta _a)`$, there exist $`\varphi ,\phi `$ such that Equation 12 is satisfied, and hence also such that the expression in Equation 11 vanishes. In conclusion, applying $`𝐐(𝒜,\chi ,\varphi ,\phi )`$ with such $`\varphi ,\phi `$ at the very last iteration allows us to measure a good solution with certainty.
## 3 Heuristics
As explained in the previous section, using the amplitude amplification technique to search for a solution to a search problem, one obtains a quadratic speedup compared to a brute force search. For many problems, however, good heuristics are known for which the expected running time, when applied to a “real-life” problem, is in $`o(\sqrt{N})`$, where $`N`$ is the size of the search space. This fact would make amplitude amplification much less useful unless a quantum computer is somehow able to take advantage of these classical heuristics. In this section we concentrate on a large family of classical heuristics that can be applied to search problems. We show how these heuristics can be incorporated into the general amplitude amplification process.
By a heuristic, we mean a probabilistic algorithm, running in polynomial time, that outputs what one is searching for with some non-negligible probability.
Suppose we have a family $``$ of functions such that each $`f`$ is of the form $`f:X\{0,1\}`$. For a given function $`f`$ we seek an input $`xX`$ such that $`f(x)=1`$. A heuristic is a function $`G:\times RX`$, for an appropriate finite set $`R`$. The heuristic $`G`$ uses a random seed $`rR`$ to generate a guess for an $`x`$ such that $`f(x)=1`$. For every function $`f`$, let $`t_f=|\{xXf(x)=1\}|`$, the number of good inputs $`x`$, and let $`h_f=|\{rRf(G(f,r))=1\}|`$, the number of good seeds. We say that the heuristic is efficient for a given $`f`$ if $`h_f/|R|>t_f/|X|`$, that is, if using $`G`$ and a random seed to generate inputs to $`f`$ succeeds with a higher probability than directly guessing inputs to $`f`$ uniformly at random. The heuristic is good in general if
$$\text{E}_{}\left(\frac{h_f}{|R|}\right)>\text{E}_{}\left(\frac{t_f}{|X|}\right).$$
Here $`\text{E}_{}`$ denotes the expectation over all $`f`$ according to some fixed distribution. Note that for some $`f`$, $`h_f`$ might be small but repeated uses of the heuristic, with seeds uniformly chosen in $`R`$, will increase the probability of finding a solution.
###### Theorem 6
Let $`\{ff:X\{0,1\}\}`$ be a family of Boolean functions and $`𝒟`$ be a probability distribution over $``$. If on a classical computer, using heuristic $`G:\times RX`$, one finds $`x_0X`$ such that $`f(x_0)=1`$ for random $`f`$ taken from distribution $`D`$ in expected time $`T`$ then using a quantum computer, a solution can be found in expected time in $`O(\sqrt{T})`$.
* Proof
A simple solution to this problem is to embed the classical heuristic $`G`$ into the function used in the algorithm QSearch. Let $`\chi (r)=f(G(f,r))`$ and $`x=G(f,\text{QSearch}(𝐖,\chi ))`$, so that $`f(x)=1`$. By Theorem 3, for each function $`f`$, we have an expected running time in $`\mathrm{\Theta }(\sqrt{|R|/h_f})`$. Let $`P_f`$ denote the probability that $`f`$ occurs. Then $`_fP_f=1`$, and we have that the expected running time is in the order of $`_f\sqrt{|R|/h_f}P_f`$, which can be rewritten as
$$\underset{f}{}\sqrt{\frac{|R|}{h_f}P_f}\sqrt{P_f}\left(\underset{f}{}\frac{|R|}{h_f}P_f\right)^{1/2}\left(\underset{f}{}P_f\right)^{1/2}=\left(\underset{f}{}\frac{|R|}{h_f}P_f\right)^{1/2}$$
by Cauchy–Schwarz’s inequality. $``$$``$
An alternative way to prove Theorem 6 is to incorporate the heuristic into the operator $`𝒜`$ and do a minor modification to $`f`$. Let $`𝒜`$ be the quantum implementation of $`G`$. It is required that the operator $`𝒜`$ be unitary, but clearly in general the classical heuristic does not need to be reversible. As usual in quantum algorithms one will need first to modify the heuristic $`G:\times RX`$ to make it reversible, which can be done efficiently using standard techniques . We obtain a reversible function $`G_f^{}:R\times \mathrm{𝟎}R\times X`$. Let $`𝒜`$ be the natural unitary operation implementing $`G_f^{}`$ and let us modify $`\chi `$ (the good set membership function) to consider only the second part of the register, that is $`\chi ((r,x))=1`$ if and only if $`f(x)=1`$. We then have that $`a=h_f/|R|`$ and by Theorem 3, for each function $`f`$, we have an expected running time in $`\mathrm{\Theta }(\sqrt{|R|/h_f})`$. The rest of the reasoning is similar. This alternative technique shows, using a simple example, the usefulness of the general scheme of amplitude amplification described in the preceding section, although it is clear that from a computational point of view this is strictly equivalent to the technique given in the earlier proof of the theorem.
## 4 Quantum amplitude estimation
Section 2 dealt in a very general way with combinatorial search problems, namely, given a Boolean function $`f:X\{0,1\}`$ find an $`xX`$ such that $`f(x)=1`$. In this section, we deal with the related problem of estimating $`t=|\{xXf(x)=1\}|`$, the number of inputs on which $`f`$ takes the value 1.
We can describe this counting problem in terms of amplitude estimation. Using the notation of Section 2, given a unitary transformation $`𝒜`$ and a Boolean function $`\chi `$, let $`|\mathrm{\Psi }=𝒜|0`$. Write $`|\mathrm{\Psi }=|\mathrm{\Psi }_1+|\mathrm{\Psi }_0`$ as a superposition of the good and bad components of $`|\mathrm{\Psi }`$. Then amplitude estimation is the problem of estimating $`a=\mathrm{\Psi }_1|\mathrm{\Psi }_1`$, the probability that a measurement of $`|\mathrm{\Psi }`$ yields a good state.
The problem of estimating $`t=|\{xXf(x)=1\}|`$ can be formulated in these terms as follows. For simplicity, we take $`X=\{0,1,\mathrm{},N1\}`$. If $`N`$ is a power of 2, then we set $`\chi =f`$ and $`𝒜=𝐖`$. If $`N`$ is not a power of 2, we set $`\chi =f`$ and $`𝒜=𝐅_N`$, the quantum Fourier transform which, for every integer $`M1`$, is defined by
$$𝐅_M:|x\frac{1}{\sqrt{M}}\underset{y=0}{\overset{M1}{}}e^{2\pi ıxy/M}|y(0x<M).$$
(13)
Then in both cases we have $`a=t/N`$, and thus an estimate for $`a`$ directly translates into an estimate for $`t`$.
To estimate $`a`$, we make good use of the properties of operator $`𝐐=𝒜𝐒_0𝒜^1𝐒_f`$. By Equation 8 in Section 2, we have that the amplitudes of $`|\mathrm{\Psi }_1`$ and $`|\mathrm{\Psi }_0`$ as functions of the number of applications of $`𝐐`$, are sinusoidal functions, both of period $`\frac{\pi }{\theta _a}`$. Recall that $`0\theta _a\pi /2`$ and $`a=\mathrm{sin}^2(\theta _a)`$, and thus an estimate for $`\theta _a`$ also gives an estimate for $`a`$.
To estimate this period, it is a natural approach to apply Fourier analysis like Shor does for a classical function in his factoring algorithm. This approach can also be viewed as an eigenvalue estimation and is best analysed in the basis of eigenvectors of the operator at hand . By Equation 4, the eigenvalues of $`𝐐`$ on the subspace spanned by $`|\mathrm{\Psi }_1`$ and $`|\mathrm{\Psi }_0`$ are $`\lambda _+=e^{ı2\theta _a}`$ and $`\lambda _{}=e^{ı2\theta _a}`$. Thus we can estimate $`a`$ simply by estimating one of these two eigenvalues. Errors in our estimate $`\stackrel{~}{\theta }_a`$ for $`\theta _a`$ translate into errors in our estimate $`\stackrel{~}{a}=\mathrm{sin}^2(\stackrel{~}{\theta }_a)`$ for $`a`$, as described in the next lemma.
###### Lemma 7
Let $`a=\mathrm{sin}^2(\theta _a)`$ and $`\stackrel{~}{a}=\mathrm{sin}^2(\stackrel{~}{\theta }_a)`$ with $`0\theta _a,\stackrel{~}{\theta }_a2\pi `$ then
$$\left|\stackrel{~}{\theta }_a\theta _a\right|\epsilon |\stackrel{~}{a}a|2\epsilon \sqrt{a(1a)}+\epsilon ^2.$$
* Proof
For $`\epsilon 0`$, using standard trigonometric identities, we obtain
$`\mathrm{sin}^2(\theta _a+\epsilon )\mathrm{sin}^2(\theta _a)`$ $`=`$ $`\sqrt{a(1a)}\mathrm{sin}(2\epsilon )+(12a)\mathrm{sin}^2(\epsilon )\text{ and}`$
$`\mathrm{sin}^2(\theta _a)\mathrm{sin}^2(\theta _a\epsilon )`$ $`=`$ $`\sqrt{a(1a)}\mathrm{sin}(2\epsilon )+(2a1)\mathrm{sin}^2(\epsilon ).`$
The inequality follows directly. $``$$``$
We want to estimate one of the eigenvalues of $`𝐐`$. For this purpose, we utilize the following operator $`\mathrm{\Lambda }`$. For any positive integer $`M`$ and any unitary operator $`𝐔`$, the operator $`\mathrm{\Lambda }_M(𝐔)`$ is defined by
$$|j|y|j(𝐔^j|y)(0j<M).$$
(14)
Note that if $`|\mathrm{\Phi }`$ is an eigenvector of $`𝐔`$ with eigenvalue $`e^{2\pi ı\omega }`$, then $`\mathrm{\Lambda }_M(𝐔)`$ maps $`|j`$$`|\mathrm{\Phi }`$ to $`e^{2\pi ı\omega j}|j|\mathrm{\Phi }`$.
###### Definition 8
For any integer $`M>0`$ and real number $`0\omega <1`$, let
$$|𝒮_M(\omega )=\frac{1}{\sqrt{M}}\underset{y=0}{\overset{M1}{}}e^{2\pi ı\omega y}|y.$$
We then have, for all $`0xM1`$
$$𝐅_M|x=|𝒮_M(x/M).$$
The state $`|𝒮_M(\omega )`$ encodes the angle $`2\pi \omega `$ ($`0\omega <1`$) in the phases of an equally weighted superposition of all basis states. Different angles have different encodings, and the overlap between $`|𝒮_M(\omega _0)`$ and $`|𝒮_M(\omega _1)`$ is a measure for the distance between the two angles $`\omega _0`$ and $`\omega _1`$.
###### Definition 9
For any two real numbers $`\omega _0,\omega _1`$, let $`d(\omega _0,\omega _1)=\mathrm{min}_z\{|z+\omega _1\omega _0|\}`$.
Thus $`2\pi d(\omega _0,\omega _1)`$ is the length of the shortest arc on the unit circle going from $`e^{2\pi ı\omega _0}`$ to $`e^{2\pi ı\omega _1}`$.
###### Lemma 10
For $`0\omega _0<1`$ and $`0\omega _1<1`$ let $`\mathrm{\Delta }=d(\omega _0,\omega _1)`$. If $`\mathrm{\Delta }=0`$ we have $`\left|𝒮_M(\omega _0)|𝒮_M(\omega _1)\right|^2=1`$. Otherwise
$$\left|𝒮_M(\omega _0)|𝒮_M(\omega _1)\right|^2=\frac{\mathrm{sin}^2(M\mathrm{\Delta }\pi )}{M^2\mathrm{sin}^2(\mathrm{\Delta }\pi )}.$$
* Proof
$`\left|𝒮_M(\omega _0)|𝒮_M(\omega _1)\right|^2`$ $`=`$ $`\left|\left({\displaystyle \frac{1}{\sqrt{M}}}{\displaystyle \underset{y=0}{\overset{M1}{}}}e^{2\pi ı\omega _0y}y|\right)\left({\displaystyle \frac{1}{\sqrt{M}}}{\displaystyle \underset{y=0}{\overset{M1}{}}}e^{2\pi ı\omega _1y}|y\right)\right|^2`$
$`=`$ $`{\displaystyle \frac{1}{M^2}}\left|{\displaystyle \underset{y=0}{\overset{M1}{}}}e^{2\pi ı\mathrm{\Delta }y}\right|^2={\displaystyle \frac{\mathrm{sin}^2(M\mathrm{\Delta }\pi )}{M^2\mathrm{sin}^2(\mathrm{\Delta }\pi )}}.`$
$``$$``$
Consider the problem of estimating $`\omega `$ where $`0\omega <1`$, given the state $`|𝒮_M(\omega )`$. If $`\omega =x/M`$ for some integer $`0x<M`$, then $`𝐅_M^1|𝒮_M(x/M)=|x`$ by definition, and thus we have a perfect phase estimator. If $`M\omega `$ is not an integer, then observing $`𝐅_M^1|𝒮_M(\omega )`$ still provides a good estimation of $`\omega `$, as shown in the following theorem.
###### Theorem 11
Let $`X`$ be the discrete random variable corresponding to the classical result of measuring $`𝐅_M^1|𝒮_M(\omega )`$ in the computational basis. If $`M\omega `$ is an integer then $`\mathrm{Prob}(X=M\omega )=1`$. Otherwise, letting $`\mathrm{\Delta }=d(\omega ,x/M)`$,
$$\mathrm{Prob}(X=x)=\frac{\mathrm{sin}^2(M\mathrm{\Delta }\pi )}{M^2\mathrm{sin}^2(\mathrm{\Delta }\pi )}\frac{1}{(2M\mathrm{\Delta })^2}.$$
For any $`k>1`$ we also have
$$\mathrm{Prob}\left(d(X/M,\omega )k/M\right)1\frac{1}{2(k1)}$$
and, in the case $`k=1`$ and $`M>2`$,
$$\mathrm{Prob}\left(d(X/M,\omega )1/M\right)\frac{8}{\pi ^2}.$$
* Proof
Clearly
$`\mathrm{Prob}(X=x)`$ $`=\left|x|𝐅^1|𝒮_M(\omega )\right|^2`$
$`=\left|(𝐅|x)^{}|𝒮_M(\omega )\right|^2`$
$`=\left|𝒮_M(x/M)|𝒮_M(\omega )\right|^2`$
thus using Lemma 10 we directly obtain the first part of the theorem. We use this fact to prove the next part of the theorem.
$`\mathrm{Prob}\left(d(X/M,\omega )k/M\right)`$ $`=`$ $`1\mathrm{Prob}(d(X/M,\omega )>k/M)`$
$``$ $`12{\displaystyle \underset{j=k}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{4M^2(\frac{j}{M})^2}}`$
$``$ $`1{\displaystyle \frac{1}{2(k1)}}.`$
For the last part, we use the fact that for $`M>2`$, the given expression attains its minimum at $`\mathrm{\Delta }=1/(2M)`$ in the range $`0\mathrm{\Delta }1/M`$.
$`\mathrm{Prob}\left(d(X/M,\omega )1/M\right)`$ $`=`$ $`\mathrm{Prob}(X=M\omega )+\mathrm{Prob}(X=M\omega )`$
$`=`$ $`{\displaystyle \frac{\mathrm{sin}^2(M\mathrm{\Delta }\pi )}{M^2\mathrm{sin}^2(\mathrm{\Delta }\pi )}}+{\displaystyle \frac{\mathrm{sin}^2(M(\frac{1}{M}\mathrm{\Delta })\pi )}{M^2\mathrm{sin}^2((\frac{1}{M}\mathrm{\Delta })\pi )}}`$
$``$ $`{\displaystyle \frac{8}{\pi ^2}}.`$
$``$$``$
The following algorithm computes an estimate for $`a`$, via an estimate for $`\theta _a`$.
* Algorithm( $`\text{Est\_Amp}(𝒜,\chi ,M)`$ )
1. Initialize two registers of appropriate sizes to the state $`|0𝒜|0`$.
2. Apply $`𝐅_M`$ to the first register.
3. Apply $`\mathrm{\Lambda }_M(𝐐)`$ where $`𝐐=𝒜𝐒_0𝒜^1𝐒_\chi `$.
4. Apply $`𝐅_M^1`$ to the first register.
5. Measure the first register and denote the outcome $`|y`$.
6. Output $`\stackrel{~}{a}=\mathrm{sin}^2(\pi \frac{y}{M})`$.
Steps 1 to 5 are illustrated on Figure 1. This algorithm can also be summarized, following the approach in , as the unitary transformation
$$\left((𝐅_M^1𝐈)\mathrm{\Lambda }_M(𝐐)(𝐅_M𝐈)\right)$$
applied on state $`|0𝒜|0`$, followed by a measurement of the first register and classical post-processing of the outcome. In practice, we could choose $`M`$ to be a power of 2, which would allow us to use a Walsh–Hadamard transform instead of a Fourier transform in step 2.
###### Theorem 12 (Amplitude Estimation)
For any positive integer $`k`$, the algorithm $`\text{Est\_Amp}(𝒜,\chi ,M)`$ outputs $`\stackrel{~}{a}`$ $`(0\stackrel{~}{a}1)`$ such that
$$\left|\stackrel{~}{a}a\right|\mathrm{\hspace{0.33em}2}\pi k\frac{\sqrt{a(1a)}}{M}+k^2\frac{\pi ^2}{M^2}$$
with probability at least $`\frac{8}{\pi ^2}`$ when $`k=1`$ and with probability greater than $`1\frac{1}{2(k1)}`$ for $`k2`$. It uses exactly $`M`$ evaluations of $`f`$. If $`a=0`$ then $`\stackrel{~}{a}=0`$ with certainty, and if $`a=1`$ and $`M`$ is even, then $`\stackrel{~}{a}=1`$ with certainty.
* Proof
After step 1, by Equation 6, we have state
$`|0𝒜|0`$ $`=`$ $`{\displaystyle \frac{ı}{\sqrt{2}}}|0\left(e^{ı\theta _a}|\mathrm{\Psi }_+e^{ı\theta _a}|\mathrm{\Psi }_{}\right).`$
After step 2, ignoring global phase, we have
$$\frac{1}{\sqrt{2M}}\underset{j=0}{\overset{M1}{}}|j\left(e^{ı\theta _a}|\mathrm{\Psi }_+e^{ı\theta _a}|\mathrm{\Psi }_{}\right)$$
and after applying $`\mathrm{\Lambda }_M(𝐐)`$ we have
$`{\displaystyle \frac{1}{\sqrt{2M}}}{\displaystyle \underset{j=0}{\overset{M1}{}}}|j\left(e^{ı\theta _a}e^{2ıj\theta _a}|\mathrm{\Psi }_+e^{ı\theta _a}e^{2ıj\theta _a}|\mathrm{\Psi }_{}\right)`$
$`=`$ $`{\displaystyle \frac{e^{ı\theta _a}}{\sqrt{2M}}}{\displaystyle \underset{j=0}{\overset{M1}{}}}e^{2ıj\theta _a}|j|\mathrm{\Psi }_+{\displaystyle \frac{e^{ı\theta _a}}{\sqrt{2M}}}{\displaystyle \underset{j=0}{\overset{M1}{}}}e^{2ıj\theta _a}|j|\mathrm{\Psi }_{}`$
$`=`$ $`{\displaystyle \frac{e^{ı\theta _a}}{\sqrt{2}}}|𝒮_M(\frac{\theta _a}{\pi })|\mathrm{\Psi }_+{\displaystyle \frac{e^{ı\theta _a}}{\sqrt{2}}}|𝒮_M(1\frac{\theta _a}{\pi })|\mathrm{\Psi }_{}.`$
We then apply $`𝐅_M^1`$ to the first register and measure it in the computational basis.
The rest of the proof follows from Theorem 11. Tracing out the second register in the eigenvector basis, we see that the first register is in an equally weighted mixture of $`𝐅_M^1|𝒮_M(\frac{\theta _a}{\pi })`$ and $`𝐅_M^1|𝒮_M(1\frac{\theta _a}{\pi })`$. Thus the measured value $`|y`$ is the result of measuring either the state $`𝐅_M^1|𝒮_M(\frac{\theta _a}{\pi })`$ or the state $`𝐅_M^1|𝒮_M(1\frac{\theta _a}{\pi })`$. The probability of measuring $`|y`$ given the state $`𝐅_M^1|𝒮_M(1\frac{\theta _a}{\pi })`$ is equal to the probability of measuring $`|My`$ given the state $`𝐅_M^1|𝒮_M(\frac{\theta _a}{\pi })`$. Since $`\mathrm{sin}^2\left(\pi \frac{(My)}{M}\right)=\mathrm{sin}^2\left(\pi \frac{y}{M}\right)`$, we can assume we measured $`|y`$ given the state $`𝐅_M^1|𝒮_M(\frac{\theta _a}{\pi })`$ and $`\stackrel{~}{\theta }_a=\pi \frac{y}{M}`$ estimates $`\theta _a`$ as described in Theorem 11. Thus we obtain bounds on $`d(\stackrel{~}{\theta }_a,\theta _a)`$ that translate, using Lemma 7, into the appropriate bounds on $`|\stackrel{~}{a}a|`$. $``$$``$
A straightforward application of this algorithm is to approximately count the number of solutions $`t`$ to $`f(x)=1`$. To do this we simply set $`𝒜=𝐖`$ if $`N`$ is a power of 2, or in general $`𝒜=𝐅_N`$ or any other transformation that maps $`|0`$ to $`\frac{1}{\sqrt{N}}_{j=0}^{N1}|j`$. Setting $`\chi =f`$, we then have $`a=\mathrm{\Psi }_1|\mathrm{\Psi }_1=t/N`$, which suggests the following algorithm.
* Algorithm( $`\text{Count}(f,M)`$ )
1. Output $`t^{}=N\times \text{Est\_Amp}(𝐅_N,f,M)`$.
By Theorem 12, we obtain the following.
###### Theorem 13 (Counting)
For any positive integers $`M`$ and $`k`$, and any Boolean function $`f:\{0,1,\mathrm{},N1\}\{0,1\}`$, the algorithm Count$`(f,M)`$ outputs an estimate $`t^{}`$ to $`t=|f^1(1)|`$ such that
$$\left|t^{}t\right|\mathrm{\hspace{0.33em}2}\pi k\frac{\sqrt{t(Nt)}}{M}+\pi ^2k^2\frac{N}{M^2}$$
with probability at least $`8/\pi ^2`$ when $`k=1`$, and with probability greater than $`1\frac{1}{2(k1)}`$ for $`k2`$. If $`t=0`$ then $`t^{}=0`$ with certainty, and if $`t=N`$ and $`M`$ is even, then $`t^{}=N`$ with certainty.
Note that Count$`(f,M)`$ outputs a real number. In the following counting algorithms we will wish to output an integer, and therefore we will round off the output of Count to an integer. To assure that the rounding off can be done efficiently<sup>1</sup><sup>1</sup>1For example, if $`t^{}+\frac{1}{2}`$ is super-exponentially close to an integer $`n`$ we may not be able to decide efficiently if $`t^{}`$ is closer to $`n`$ or $`n1`$. we will round off to an integer $`\stackrel{~}{t}`$ satisfying $`\left|\stackrel{~}{t}\text{Count}(f,M)\right|\frac{2}{3}`$.
If we want to estimate $`t`$ within a few standard deviations, we can apply algorithm Count with $`M=\sqrt{N}`$.
###### Corollary 14
Given a Boolean function $`f:\{0,1,\mathrm{},N1\}\{0,1\}`$ with $`t`$ defined as above, rounding off the output of $`\text{Count}(f,\sqrt{N})`$ gives an estimate $`\stackrel{~}{t}`$ such that
$$\left|\stackrel{~}{t}t\right|<\mathrm{\hspace{0.33em}2}\pi \sqrt{\frac{t(Nt)}{N}}+11$$
(15)
with probability at least $`8/\pi ^2`$ and requires exactly $`\sqrt{N}`$ evaluations of $`f`$.
We now look at the case of estimating $`t`$ with some relative error, also referred to as approximately counting $`t`$ with accuracy $`\epsilon `$. For this we require the following crucial observation about the output $`t^{}`$ of algorithm $`\text{Count}(f,L)`$. Namely $`t^{}`$ is likely to be equal to zero if and only if $`Lo(\sqrt{N/t})`$. Thus, we can find a rough estimate of $`\sqrt{N/t}`$ simply by running algorithm $`\text{Count}(f,L)`$ with exponentially increasing values of $`L`$ until we obtain a non-zero output. Having this rough estimate $`L`$ of $`\sqrt{N/t}`$ we can then apply Theorem 13 with $`M`$ in the order of $`\frac{1}{\epsilon }L`$ to find an estimate $`\stackrel{~}{t}`$ of $`t`$ with the required accuracy. The precise algorithm is as follows.
* Algorithm( $`\text{Basic\_Approx\_Count}(f,\epsilon )`$ )
1. Start with $`\mathrm{}=0`$.
2. Increase $`\mathrm{}`$ by 1.
3. Set $`t^{}=`$ Count$`(f,2^{\mathrm{}})`$.
4. If $`t^{}=0`$ and $`2^{\mathrm{}}<2\sqrt{N}`$ then go to step 2.
5. Set $`M=\frac{20\pi ^2}{\epsilon }2^{\mathrm{}}`$.
6. Set $`t^{}=\text{Count}(f,M)`$.
7. Output an integer $`\stackrel{~}{t}`$ satisfying $`\left|\stackrel{~}{t}t^{}\right|\frac{2}{3}`$.
###### Theorem 15
Given a Boolean function $`f`$ with $`N`$ and $`t`$ defined as above, and any $`0<\epsilon 1`$, $`\text{Basic\_Approx\_Count}(f,\epsilon )`$ outputs an estimate $`\stackrel{~}{t}`$ such that
$$\left|\stackrel{~}{t}t\right|\epsilon t$$
with probability at least $`\frac{2}{3}`$, using an expected number of evaluations of $`f`$ which is in $`\mathrm{\Theta }\left(\frac{1}{\epsilon }\sqrt{N/t}\right)`$. If $`t=0`$, the algorithm outputs $`\stackrel{~}{t}=t`$ with certainty and $`f`$ is evaluated a number of times in $`\mathrm{\Theta }\left(\sqrt{N}\right)`$.
* Proof
When $`t=0`$, the analysis is straightforward. For $`t>0`$, let $`\theta `$ denote $`\theta _{t/N}`$ and $`m=\mathrm{log}_2(\frac{1}{5\theta })`$. From Theorem 11 we have that the probability that step 3 outputs Count$`(f,2^{\mathrm{}})=0`$ for $`\mathrm{}=1,2,\mathrm{},m`$ is
$$\underset{\mathrm{}=1}{\overset{m}{}}\frac{\mathrm{sin}^2(2^{\mathrm{}}\theta )}{2^2\mathrm{}\mathrm{sin}^2(\theta )}\underset{\mathrm{}=1}{\overset{m}{}}\mathrm{cos}^2(2^{\mathrm{}}\theta )=\frac{\mathrm{sin}^2(2^{m+1}\theta )}{2^{2m}\mathrm{sin}^2(2\theta )}\mathrm{cos}^2\left(\frac{2}{5}\right).$$
The previous inequalities are obtained by using the fact that $`\mathrm{sin}(M\theta )M\mathrm{sin}(\theta )\mathrm{cos}(M\theta )`$ for any $`M0`$ and $`0M\theta <\frac{\pi }{2}`$, which can be readily seen by considering the Taylor expansion of $`\mathrm{tan}(x)`$ at $`x=M\theta `$.
Now assuming step 3 has outputted $`0`$ at least $`m`$ times (note that $`2^m\frac{1}{5\theta }\frac{1}{5}\sqrt{N/t}<2\sqrt{N}`$), after step 5 we have $`M\frac{20\pi ^2}{\epsilon }2^{m+1}\frac{4\pi ^2}{\epsilon \theta }`$ and by Theorem 13 (and the fact that $`\theta \frac{\pi }{2}\mathrm{sin}(\theta )=\frac{\pi }{2}\sqrt{t/N}`$) the probability that Count$`(f,M)`$ outputs an integer $`t^{}`$ satisfying $`|t^{}t|\frac{\epsilon }{4}t+\frac{\epsilon ^2}{64}t`$ is at least $`8/\pi ^2`$. Let us suppose this is the case. If $`\epsilon t<1`$, then $`|\stackrel{~}{t}t|<1`$ and, since $`\stackrel{~}{t}`$ and $`t`$ are both integers, we must have $`t=\stackrel{~}{t}`$. If $`\epsilon t1`$, then rounding off $`t^{}`$ to $`\stackrel{~}{t}`$ introduces an error of at most $`\frac{2}{3}\frac{2\epsilon }{3}t`$, making the total error at most $`\frac{\epsilon }{4}t+\frac{\epsilon ^2}{64}t+\frac{2\epsilon }{3}t<\epsilon t`$. Therefore the overall probability of outputting an estimate with error at most $`\epsilon t`$ is at least $`\mathrm{cos}^2\left(\frac{2}{5}\right)\times (8/\pi ^2)>\frac{2}{3}`$.
To upper bound the number of applications of $`f`$, note that by Theorem 13, for any integer $`L18\pi \sqrt{N/t}`$, the probability that Count$`(f,L)`$ outputs 0 is less than $`1/4`$. Thus the expected value of $`M`$ at step 6 is in $`\mathrm{\Theta }(\frac{1}{\epsilon }\sqrt{N/t})`$. $``$$``$
We remark that in algorithm Basic\_Approx\_Count, we could alternatively to steps 1 to 4 use algorithm QSearch of Section 2, provided we have QSearch also output its final value of $`M`$. In this case, we would use (a multiple of) that value as our rough estimate of $`\sqrt{N/t}`$, instead of using the final value of $`2^{\mathrm{}}`$ found in step 4 of Basic\_Approx\_Count.
Algorithm Basic\_Approx\_Count is optimal for any fixed $`\epsilon `$, but not in general. In Appendix A we give an optimal algorithm, while we now present two simple optimal algorithms for counting the number of solutions exactly. That is, we now consider the problem of determining the exact value of $`t=|f^1(1)|`$. In the special case that we are given a nonzero integer $`t_0`$ and promised that either $`t=0`$ or $`t=t_0`$, then we can determine which is the case with certainty using a number of evaluations of $`f`$ in $`O(\sqrt{N/t_0})`$. This is an easy corollary of Theorem 4 and we state it without proof.
###### Theorem 16
Let $`f:\{0,1,\mathrm{},N1\}\{0,1\}`$ be a given Boolean function such that the cardinality of the preimage of 1 is either 0 or $`t_0`$. Then there exists a quantum algorithm that determines with certainty which is the case using a number of evaluations of $`f`$ which is in $`\mathrm{\Theta }\left(\sqrt{N/t_0}\right)`$, and in the latter case, also outputs a random element of $`f^1(1)`$.
For the general case in which we do not have any prior knowledge about $`t`$, we offer the following algorithm.
* Algorithm( $`\text{Exact\_Count}(f)`$ )
1. Set $`t_1^{}=\text{Count}(f,14\pi \sqrt{N})`$ and $`t_2^{}=\text{Count}(f,14\pi \sqrt{N})`$.
2. Let $`M_i=30\sqrt{(t_i^{}+1)(Nt_i^{}+1)}`$ for $`i=1,2`$.
3. Set $`M=\mathrm{min}\{M_1,M_2\}`$.
4. Set $`t^{}=\text{Count}(f,M)`$.
5. Output an integer $`\stackrel{~}{t}`$ satisfying $`\left|\stackrel{~}{t}t^{}\right|\frac{2}{3}`$.
The main idea of this algorithm is the same as that of algorithm Basic\_Approx\_Count. First we find a rough estimate $`t_r^{}`$ of $`t`$, and then we run algorithm Count$`(f,M)`$ with a value of $`M`$ that depends on $`t_r^{}`$. By Theorem 13, if we set $`M`$ to be in the order of $`\sqrt{t_r^{}(Nt_r^{})}`$, then the output $`t^{}=\text{Count}(f,M)`$ is likely to be so that $`|t^{}t|<\frac{1}{3}`$, in which case $`\stackrel{~}{t}=t`$.
###### Theorem 17
Given a Boolean function $`f`$ with $`N`$ and $`t`$ defined as above, algorithm Exact\_Count requires an expected number of evaluations of $`f`$ which is in $`\mathrm{\Theta }(\sqrt{(t+1)(Nt+1)})`$ and outputs an estimate $`\stackrel{~}{t}`$ which equals $`t`$ with probability at least $`\frac{2}{3}`$ using space only linear in $`\mathrm{log}(N)`$.
* Proof
Apply Theorem 13 with $`k=7`$. For each $`i=1,2`$, with probability greater than $`\frac{11}{12}`$, outcome $`t_i^{}`$ satisfies $`\left|t_i^{}t\right|<\sqrt{\frac{t(Nt)}{N}}+1/4`$, in which case we also have that $`\sqrt{t(Nt)}\frac{\sqrt{2}}{30}M_i`$. Thus, with probability greater than $`\left(\frac{11}{12}\right)^2`$, we have
$$\frac{\sqrt{t(Nt)}}{M}\frac{\sqrt{2}}{30}.$$
Suppose this is the case. Then by Theorem 13, with probability at least $`8/\pi ^2`$,
$$|t^{}t|\frac{2\pi \sqrt{2}}{30}+\frac{4\pi ^2}{30^2}<\frac{1}{3}$$
and consequently
$$|\stackrel{~}{t}t|<1.$$
Hence, with probability at least $`\left(\frac{11}{12}\right)^2\times 8/\pi ^2>\frac{2}{3}`$, we have $`\stackrel{~}{t}=t`$.
The number of applications of $`f`$ is $`214\pi \sqrt{N}+M`$. Consider the expected value of $`M_i`$ for $`i=1,2`$. Since
$$\sqrt{(t_i^{}+1)(Nt_i^{}+1)}\sqrt{(t+1)(Nt+1)}+\sqrt{N|t_i^{}t|}$$
for any $`0t_i^{},tN`$, we just need to upper bound the expected value of $`\sqrt{N|t_i^{}t|}`$. By Theorem 13, for any $`k2`$,
$$|t_i^{}t|k\sqrt{\frac{t(Nt)}{N}}+k^2$$
with probability at least $`1\frac{1}{k}`$. Hence $`M_i`$ is less than
$$30(1+k)\left(\sqrt{(t+1)(Nt+1)}+\sqrt{N}\right)+1$$
(16)
with probability at least $`1\frac{1}{k}`$.
In particular, the minimum of $`M_1`$ and $`M_2`$ is greater than the expression given in Equation 16 with probability at most $`\frac{1}{k^2}`$. Since any positive random variable $`Z`$ satisfying $`\text{Prob}(Z>k)\frac{1}{k^2}`$ has expectation upper bounded by a constant, the expected value of $`M`$ is in $`O\left(\sqrt{(t+1)(Nt+1)}\right)`$. $``$$``$
It follows from Theorem 4.10 of that any quantum algorithm capable of deciding with high probability whether or not a function $`f:\{0,1,\mathrm{},N1\}\{0,1\}`$ is such that $`|f^1(1)|t`$, given some $`0<t<N`$, must query $`f`$ a number of times which is at least in $`\mathrm{\Omega }\left(\sqrt{(t+1)(Nt+1)}\right)`$ times. Therefore, our exact counting algorithm is optimal up to a constant factor.
Note also that successive applications of Grover’s algorithm in which we strike out the solutions as they are found will also provide an algorithm to perform exact counting. In order to obtain a constant probability of success, if the algorithm fails to return a new element, one must do more than a constant number of trials. In particular, repeating until we get $`\mathrm{log}(N)`$ failures will provide an overall constant probability of success. Unfortunately, the number of applications of $`f`$ is then in $`O\left(\sqrt{tN}+\mathrm{log}(N)\sqrt{N/t}\right)`$ and the cost in terms of additional quantum memory is prohibitive, that is in $`\mathrm{\Theta }(t)`$.
## 5 Concluding remarks
Let $`f:\{0,1,\mathrm{},N1\}\{0,1\}`$ be a function provided as a black box, in the sense that the only knowledge available about $`f`$ is given by evaluating it on arbitrary points in its domain. We are interested in the number of times that $`f`$ must be evaluated to achieve certain goals, and this number is our measure of efficiency. Grover’s algorithm can find the $`x_0`$ such that $`f(x_0)=1`$ quadratically faster in the expected sense than the best possible classical algorithm provided the solution is known to be unique . We have generalized Grover’s algorithm in several directions.
* The quadratic speedup remains when the solution is not unique, even if the number of solutions is not known ahead of time.
* If the number of solutions is known (and nonzero), we can find one quadratically faster in the worst case than would be possible classically even in the expected case.
* If the number $`t`$ of solutions is known to be either 0 or $`t_0`$, we can tell which is the case with certainty, and exhibit a solution if $`t>0`$, in a time in $`O(\sqrt{N/t_0})`$ in the worst case. By contrast, the best classical algorithm would need $`Nt_0+1`$ queries in the worst case. This is much better than a quadratic speedup when $`t_0`$ is large.
* The quadratic speedup remains in a variety of settings that are not constrained to the black-box model: even if additional information about $`f`$ can be used to design efficient classical heuristics, we can still find solutions quadratically faster on a quantum computer, provided the heuristic falls under the broad scope of our technique.
* We give efficient quantum algorithms to estimate the number of solutions in a variety of error models. In all cases, our quantum algorithms are proven optimal, up to a multiplicative constant, among all possible quantum algorithms. In most cases, our quantum algorithms are known to be quadratically faster than the best possible classical algorithm. In the case of counting the number of solutions up to relative error $`\epsilon `$, our optimal quantum algorithm is quadratically faster than the best known classical algorithm for fixed $`\epsilon `$, but in fact it is better than that when $`\epsilon `$ is not a constant. Since we do not believe that a super-quadratic quantum improvement for a non-promise black-box problem is possible, we conjecture that there exists a classical algorithm that uses a number of queries in $`O(\mathrm{min}\{M^2,N\})`$, where $`M=\sqrt{\frac{N}{\epsilon t+1}}+\frac{\sqrt{t(Nt)}}{\epsilon t+1}`$ is proportional to the number of queries required by our optimal quantum algorithm. This conjecture is further supported by the fact that we can easily find a good estimate for $`M^2`$, without prior knowledge of $`t`$, using a number of classical queries in $`O(\frac{1}{\epsilon }+\frac{N}{t+1})`$.
* We can amplify efficiently the success probability not only of classical search algorithms, but also of quantum algorithms. More precisely, if a quantum algorithm can output an $`x`$ that has probability $`a>0`$ of being such that $`f(x)=1`$, then a solution can be found after evaluating $`f`$ an expected number of time in $`O(1/\sqrt{a})`$. If the value of $`a`$ is known, a solution can be found after evaluating $`f`$ a number of time in $`O(1/\sqrt{a})`$ even in the worst case. We call this process *amplitude amplification*. Again, this is quadratically faster than would be possible if the quantum search algorithm were available as a black box to a classical algorithm.
* Finally, we provide a general technique, known as *amplitude estimation*, to estimate efficiently the success probability $`a`$ of quantum search algorithms. This is the natural quantum generalization of the above-mentioned technique to estimate the number of classical solutions to the equation $`f(x)=1`$.
The following table summarizes the number of applications of the given function $`f`$ in the quantum algorithms presented in this paper. The table also compares the quantum complexities with the classical complexities of these problems, when the latter are known. Any lower bounds indicated (implicit in the use of the “$`\mathrm{\Theta }`$” notation) correspond to those in the black-box model of computation. In the case of the efficiency of quantum counting with accuracy $`\epsilon `$, we refer to the algorithm given below in the Appendix.
| Problem | Quantum Complexity | Classical Complexity |
| --- | --- | --- |
| Decision | $`\mathrm{\Theta }(\sqrt{N/(t+1)})`$ | $`\mathrm{\Theta }(N/(t+1))`$ |
| Searching | $`\mathrm{\Theta }(\sqrt{N/(t+1)})`$ | $`\mathrm{\Theta }(N/(t+1))`$ |
| Counting with error $`\sqrt{t}`$ | $`\mathrm{\Theta }(\sqrt{N})`$ | |
| Counting with accuracy $`\epsilon `$ | $`\mathrm{\Theta }\left(\sqrt{\frac{N}{\epsilon t+1}}+\frac{\sqrt{t(Nt)}}{\epsilon t+1}\right)`$ | $`O(\frac{1}{\epsilon ^2}N/(t+1))`$ |
| Exact counting | $`\mathrm{\Theta }\left(\sqrt{(t+1)(Nt+1)}\right)`$ | $`\mathrm{\Theta }(N)`$ |
We leave as open the problem of finding a quantum algorithm that exploits the structure of some searching or counting problem in a genuinely quantum way. By this, we mean in a way that is not equivalent to applying amplitude amplification or amplitude estimation to a classical heuristic. Note that Shor’s factoring algorithm does this in the different context of integer factorization.
## Acknowledgements
We are grateful to Joan Boyar, Harry Buhrman, Artur Ekert, Ashwin Nayak, Jeff Shallitt, Barbara Terhal and Ronald de Wolf for helpful discussions.
## Appendix A Tight Algorithm for Approximate <br>Counting
Here we combine the ideas of algorithms Basic\_Approx\_Count and Exact\_Count to obtain an optimal algorithm for approximately counting. That this algorithm is optimal follows readily from Corollary 1.2 and Theorem 1.13 of Nayak and Wu .
###### Theorem 18
Given a Boolean function $`f`$ with $`N`$ and $`t`$ defined as above, and any $`\epsilon `$ such that $`\frac{1}{3N}<\epsilon 1`$, the following algorithm $`\text{Approx\_Count}(f,\epsilon )`$ outputs an estimate $`\stackrel{~}{t}`$ such that
$$\left|\stackrel{~}{t}t\right|\epsilon t$$
with probability at least $`\frac{2}{3}`$, using an expected number of evaluations of $`f`$ in the order of
$$S=\sqrt{\frac{N}{\epsilon t+1}}+\frac{\sqrt{t(Nt)}}{\epsilon t+1}.$$
If $`t=0`$ or $`t=N`$, the algorithm outputs $`\stackrel{~}{t}=t`$ with certainty.
We assume that $`\epsilon N>1/3`$, since otherwise approximately counting with accuracy $`\epsilon `$ reduces to exact counting. Set
$$S^{}=\mathrm{min}\{\frac{1}{\sqrt{\epsilon }}\sqrt{\frac{N}{t}}\left(1+\sqrt{\frac{Nt}{\epsilon N}}\right),\sqrt{(t+1)(Nt+1)}\}$$
(17)
and note that $`S^{}\mathrm{\Theta }(S)`$ where $`S`$ is defined as in Theorem 18. The algorithm works by finding approximate values for each of the different terms in Equation 17. The general outline of the algorithm is as follows.
* Algorithm( $`\text{Approx\_Count}(f,\epsilon )`$ )
1. Find integer $`L_1`$ approximating $`\sqrt{N/(t+1)}`$.
2. Find integer $`L_2`$ approximating $`\sqrt{(Nt)/(\epsilon N)}`$.
3. Set $`M_1=\frac{1}{\sqrt{\epsilon }}L_1(1+L_2)`$.
4. If $`M_1>\sqrt{N}`$ then find integer $`M_2`$ approximating $`\sqrt{(t+1)(Nt+1)}`$. If $`M_1\sqrt{N}`$ then set $`M_2=\mathrm{}`$.
5. Set $`M=\mathrm{min}\{M_1,M_2\}`$.
6. Set $`t^{}=\text{Count}(f,10\pi M)`$.
7. Output an integer $`\stackrel{~}{t}`$ satisfying $`\left|\stackrel{~}{t}t^{}\right|\frac{2}{3}`$.
* Proof
To find $`L_1`$, we run steps 1 to 4 of algorithm Basic\_Approx\_Count and then set $`L_1=9\pi \times 2^l`$. A proof analogous to that of Theorem 15 gives that
+ $`L_1>\sqrt{N/(t+1)}`$ with probability at least $`0.95`$, and
+ the expected value of $`L_1`$ is in $`\mathrm{\Theta }\left(\sqrt{N/(t+1)}\right)`$.
This requires a number of evaluations of $`f`$ which is in $`\mathrm{\Theta }(L_1)`$ , and thus, the expected number of evaluations of $`f`$ so far is in $`O(S^{})`$.
In step 2, for some constant $`c`$ to be determined below, we use $`2\frac{c}{\sqrt{\epsilon }}`$ evaluations of $`f`$ to find integer $`L_2`$ satisfying
+ $`L_2>\sqrt{(Nt)/(\epsilon N)}`$ with probability at least $`0.95`$, and
+ the expected value of $`L_2`$ is in $`O\left(\sqrt{(Nt+1)/(\epsilon N)}\right)`$.
Since $`Nt=|f^1(0)|`$, finding such $`L_2`$ boils down to estimating, with accuracy in $`\mathrm{\Theta }(\sqrt{\epsilon })`$, the square root of the probability that $`f`$ takes the value 0 on a random point in its domain. Or equivalently, the probability that $`\neg f`$ takes the value 1, where $`\neg f=1f`$. Suppose for some constant $`c`$, we run $`\text{Count}(\neg f,\frac{c}{\sqrt{\epsilon }})`$ twice with outputs $`\stackrel{~}{r}_1`$ and $`\stackrel{~}{r}_2`$. By Theorem 13, each output $`\stackrel{~}{r}_i`$ ($`i=1,2`$) satisfies that
$$\left|\sqrt{\frac{\stackrel{~}{r}_i}{\epsilon N}}\sqrt{\frac{Nt}{\epsilon N}}\right|\sqrt{\frac{2\pi k}{c}}\sqrt[\text{4}]{\frac{Nt}{\epsilon N}}+\frac{\pi k}{c}$$
with probability at least $`1\frac{1}{2(k1)}`$ for every $`k2`$. It follows that $`\stackrel{~}{r}=\mathrm{min}\{\sqrt{\stackrel{~}{r}_1/(\epsilon N)},\sqrt{\stackrel{~}{r}_2/(\epsilon N)}\}`$ has expected value in $`O\left(\sqrt{(Nt+1)/(\epsilon N)}\right)`$. Setting $`k=21`$, $`c=8\pi k`$, and $`L_2=2\stackrel{~}{r}+1`$, ensures that $`L_2`$ satisfies the two properties mentioned above. The number of evaluations of $`f`$ in step 2 is in $`\mathrm{\Theta }(\frac{1}{\sqrt{\epsilon }})`$ which is in $`O(S^{})`$.
In step 3, we set $`M_1=\frac{1}{\sqrt{\epsilon }}L_1(1+L_2)`$. Note that
+ $`M_1>\frac{1}{\sqrt{\epsilon }}\sqrt{\frac{N}{t+1}}\left(1+\sqrt{\frac{Nt}{\epsilon N}}\right)`$ with probability at least $`0.95^2`$, and
+ the expected value of $`M_1`$ is in the order of $`\frac{1}{\sqrt{\epsilon }}\sqrt{\frac{N}{t+1}}\left(1+\sqrt{\frac{Nt+1}{\epsilon N}}\right)`$.
In step 4, analogously to algorithm Exact\_Count, a number of evaluations of $`f`$ in $`\mathrm{\Theta }(\sqrt{N})`$ suffices to find an integer $`M_2`$ such that
+ $`M_2>\sqrt{(t+1)(Nt+1)}`$ with probability at least $`0.95`$, and
+ the expected value of $`M_2`$ is in $`\mathrm{\Theta }\left(\sqrt{(t+1)(Nt+1)}\right)`$.
Fortunately, since $`\sqrt{(t+1)(Nt+1)}\sqrt{N}`$, we shall only need $`M_2`$ if $`M_1>\sqrt{N}`$. We obtain that, after step 5,
+ $`M`$ is greater than
$$\mathrm{min}\{\frac{1}{\sqrt{\epsilon }}\sqrt{\frac{N}{t+1}}\left(1+\sqrt{\frac{Nt}{\epsilon N}}\right),\sqrt{(t+1)(Nt+1)}\}$$
with probability at least $`0.95^3>0.85`$, and
+ the expected value of $`M`$ is in $`O(S^{})`$.
To derive this latter statement, we use the fact that the expected value of the minimum of two random variables is at most the minimum of their expectation.
Finally, by Theorem 13, applying algorithm $`\text{Count}(f,10\pi M)`$ given such an $`M`$, produces an estimate $`t^{}`$ of $`t`$ such that $`|t^{}t|\frac{\epsilon t}{3}`$ (which implies that $`|\stackrel{~}{t}t|\epsilon t`$) with probability at least $`8/\pi ^2`$. Hence our overall success probability is at least $`0.85\times 8/\pi ^2>2/3`$, and the expected number of evaluations of $`f`$ is in $`O(S^{})`$. $``$$``$
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# 1 Introduction
## 1 Introduction
One of the remarkable insights of orbifold string theory is an indication of the existence of a new cohomology theory of orbifolds containing so-called twisted sectors as the contribution of singularities. Mathematically, such an orbifold cohomology theory has been constructed by Chen-Ruan \[CR\]. Author believes that there is a ”stringy” geometry and topology of orbifolds of which orbifold cohomology is its core. One aspect of this new geometry and topology is the twisted orbifold cohomology and its relation to discrete torison. Let me first explain their physical origin. Physicists usually work over a global quotient $`X=Y/G`$ only, where $`G`$ is a finite group acting smoothly on $`Y`$. A discrete torsion is a cohomology class $`\alpha H^2(G,U(1))`$. Physically, a discrete torsion counts the freedom to choose a phase factor to weight path integral over each twisted sector without destroying the consistency of string theory. For each $`\alpha `$, Vafa-Witten \[VW\] constructed the twisted orbifold cohomology group $`H_{orb,\alpha }^{}(X/G,𝐂)`$.
Mathematically, Vafa-Witten suggested that discrete torsion and twisted orbifold cohomology is connected to the problem of desingularizations. Recall that there are two methods to remove singularities, resolution or deformation. Both play important roles in the theory of Calabi-Yau 3-folds. One can obtain a smooth manifold $`Y`$ from an orbifold $`X`$ by using a combination of resolution and deformation. We call $`Y`$ a desingularization of $`X`$. In string theory, we also require the resolution to be a crepant resolution. It is known that a desingualization may not exist in dimension higher than three. In this case, we allow our desingularization to be an orbifold. As we mentioned in \[CR\], physicists predicted that ordinary orbifold cohomology group is the same as ordinary cohomology group of its crepant resolution. Vafa-Witten suggested that discrete torsion is a parameter for deformation. Furthermore, the cohomology of the desingularization is the twisted orbifold cohomology of discrete torsion plus possible contributions of exceptional loci of small resolution. A small resolution is a special kind of resolution whose exceptional loci is of codimension 2 or more. However, this proposal immediately ran into trouble because there are many more desingularizations than the number of discrete torsions. For example, D. Joyce \[JO\] constructed five different desingularizations of $`T^6/𝐙_4`$ while $`H^2(𝐙_4,U(1))=0`$. To count these ”missing” desingularizations seems to be a serious problem. On the another hand, it is well-known that most orbifolds (even Calabi-Yau orbifolds) are not global quotients. Therefore, it is also necessary to develop the theory over general orbifolds.
We will address both problems in this paper. First, we introduce the notion of inner local system $``$ for arbitrary orbifold. A local system is defined as an assignment of a flat (orbifold)-line bundle $`L_{(g)}`$ to each sector $`X_{(g)}`$ satisfying certain compatibility condition (See definition 2.1). Such a compatibility condition is designed in such a way that Poincare duality and cup product in ordinary orbifold cohomology survive the process of twisting. Then, twisted orbifold cohomology $`H_{orb}^{}(X,)`$ is defined as orbifold cohomology with value in the inner local system (See Definition 2.2). We will demonstrate that our inner local systems count all the examples constructed by D. Joyce. The author believes that the inner local system is a more fundamental notion than the discrete torsion. Then, we can formulate following mathematical conjecture: Suppose that $`X`$ is a Calabi-Yau Gorenstein orbifold. For every desingularization, we can associate a inner local systems such that as additive groups the ordinary orbifold cohomology of desigularization is the sum of twisted orbifold cohomology and contributions from exceptional loci of small resolution.
Our next goal is to determine appropriate notion of discrete torsion for general orbifold. Let $`X`$ be an arbitrary almost complex orbifold. The author’s key observation is that we should use the orbifold fundamental group $`\pi _1^{orb}(X)`$ (See definition 2.1) to replace $`G`$. Then a discrete torsion of $`X`$ is defined as a cohomology class $`\alpha H^2(\pi _1^{orb}(X),U(1))`$. Note that if $`X=Y/G`$ is a global quotient, there is a short exact sequence
$$1\pi _1(Y)\pi _1^{orb}(X)G1.$$
$`(1.1)`$
It induces a homomorphism $`H^2(G,U(1))H^2(\pi _1^{orb}(X),U(1))`$. Hence a discrete torsion in the sense of Vafa-Witten induces the discrete torsion in this paper. In fact, we can do better, we can define a local discrete torsion for each connected component of singular loci. A global discrete torsion is defined as an assignment of a local discrete torsion to a connected component of singular loci. Then, the link between discrete torsion and twisted orbifold cohomology is the theorem that a global discrete torsion induces an inner local system and hence define a twisted orbifold cohomology. However, we want to emphasis that not every inner local system comes from discrete torsion (See example 5.3).
We will introduce inner local system and twisted orbifold cohomology ring in section 2. The section 3 is devoted to discrete torsion. The relation between discrete torsion and local system is discussed in section 4. Finally, some examples are computed in section 5. This paper can be viewed as a sequel to \[CR\]. Since many constructions are similar, we will follow the notations of \[CR\] and be sketchy in the details. The author strongly encourages readers to read \[CR\] first before reading this paper.
This paper was completed while author was visiting Caltech. He would like to thank R. Pandharipande and the Caltech math department for financial support and hospitality. The author would like to thank E. Zaslow for bringing his attention to Vafa-Witten’s paper and E. Witten for explaining to him \[VW\]. He would like also thank A. Adem for many valuable discussions about group cohomology.
## 2 Local system and twisted orbifold cohomology
### 2.1 Review of ordinary orbifold cohomology
Suppose that $`X`$ is an orbifold. By the definition, $`X`$ is a topological space with a system of orbifold charts (uniformizing system). Namely, every point $`pX`$ has a system of orbifold chart of the form $`U_p/G_p`$ where $`U_p`$ is a smooth manifold and $`G_p`$ is a finite group acting on $`U_p`$ fixing $`p`$. $`G_p`$ is called a local group. Note that the action of $`G_p`$ does not have to be effective. If it does, we call it a reduced orbifold. We use $`(U_p,G_p)`$ to denote the chart. The patching condition is follows: if $`qU_p/G_pU_r/G_r`$, there is an orbifold chart $`U_q/G_qU_p/G_pU_r/G_r`$. Moreover, the inclusion map $`i:U_q/G_qU_p/G_p`$ can be lifted to a smooth map
$$\stackrel{~}{i}_{pq}:U_qU_p$$
$`(2.1)`$
and an injective homomorphism
$$i_{\mathrm{\#},pq}:G_qG_p$$
$`(2.2)`$
such that $`\stackrel{~}{i}`$ is $`i_\mathrm{\#}`$-equivariant.
$$i_{pq}=(\stackrel{~}{i},i_\mathrm{\#}):(U_q,G_q)(U_p,G_p)$$
is called an injection. A different lifting differs from $`\stackrel{~}{i}`$ by the action of an element of $`G_p`$. Moreover, $`i_\mathrm{\#}`$ differs by the conjugation of the same element. We say that the corresponding injections are equivalent.
Therefore, for any $`gG_q`$, the conjugacy class $`(i_\mathrm{\#}(g))_{G_p}`$ is well-defined. We define an equivalence relation $`(g)_{G_q}(i_\mathrm{\#}(g))_{G_p}`$. Let $`T_1`$ be the set of equivalence classes. By abusing the notation, we use $`(g)`$ to denote the equivalence class to which $`(g)_{G_q}`$ belongs to. For each $`(g)T_1`$, we can define a sector
$$X_{(g)}=\{(x,(g^{})_{G_x})|g^{}G_x,(g^{})_{G_x}(g)\}.$$
$`(2.3)`$
It was shown in \[CR\] that $`X_{(g)}`$ is an orbifold. It is the common convention that we call $`X_{(g)}`$ for $`g1`$ a twisted sector and $`X_{(1)}`$ a nontwisted sector. Once we define sectors, we diagonalize the action of $`g`$ in $`T_xX_{(g)}`$ for each $`xX_{(g)}`$. Suppose that $`g=diag(e^{\frac{2\pi im_1}{m}},\mathrm{},e^{\frac{2\pi im_n}{m}})`$, where $`m`$ is the order of $`g`$ and $`0\frac{m_i}{m}<1`$. Then we define the degree shifting number $`\iota _{(g)}=_i\frac{m_i}{m}`$. One can show that $`\iota _{(g)}`$ is independent of $`xX_{(g)}`$. The ordinary orbifold cohomology is defined as
$$H_{orb}^{}(X,𝐂)=_{(g)T_1}H^{2\iota _{(g)}}(X_{(g)},𝐂).$$
$`(2.4)`$
There is a diffeomorphism $`I:X_{(g)}X_{(g^{1)}}`$ defined by $`(x,(g))(x,(g^1))`$. Poincare paring $`<>_{orb}`$ of orbifold cohomology is defined as the direct sum of
$$<>_{orb}^{(g)}:H^{d2\iota _{(g)}}(X_{(g)},𝐂)H^{2nd2\iota _{(g^1)}}(X_{(g^1)},𝐂)𝐂$$
$`(2.5)`$
defined by
$$<\alpha ,\beta >_{orb}^{(g)}=_{X_{(g)}}\alpha I^{}\beta .$$
$`(2.6)`$
Next, we consider cup product. We first construct a moduli space (see \[CR\] (section 4.1)).
$$X_3=\{(x,(g_1,g_2,g_3)_{G_x})|g_iG_x,g_1g_2g_3=1\}$$
$`(2.7)`$
$`X_3`$ is an orbifold. Let $`𝐠=(g_1,\mathrm{},g_k)`$ with $`g_iG_q`$. By abusing the notation, we simply say $`𝐠G_q`$. We define the equivalence relation $`(𝐠)_{G_q}(i_\mathrm{\#}(𝐠))_{G_q}`$. Let $`T_k`$ be the set of equivalence class and use $`(𝐠)`$ to denote the equivalence class such that $`(𝐠)_{G_q}(𝐠)`$. We will use $`T_k^oT_k`$ to denote the set of equivalence classes of $`(𝐠)`$ such that $`𝐠=(g_1,\mathrm{},g_k)`$ with $`g_1\mathrm{}g_k=1`$. It was proved in \[CR\] that
$$X_{(𝐠)}=\{(x,(𝐠^{})_{G_x})|𝐠^{}G_x,(𝐠^{})=(𝐠)\}.$$
$`(2.8)`$
is an orbifold. One can check that
$$X_3=\underset{(𝐠)T_3^o}{}X_{(𝐠)},$$
$`(2.9)`$
Then, for each $`(𝐠)T_3^o`$, we can define evaluation maps
$$e_i:X_{(𝐠)}X_{(g)}$$
by
$$e_i(x,(g_1,g_2,g_3)_{G_x})=(x,(g_i)_{G_x}).$$
Furthermore, there is an obstruction bundle $`E`$ (see Lemma 4.2.2 \[CR\]). Then, we can define a three-point function
$$<\alpha ,\beta ,\gamma >_{orb}=_{X_{(𝐠)}}e_1^{}\alpha e_2^{}\beta e_3^{}\gamma e(E),$$
$`(2.10)`$
for any $`\alpha H^{p\iota _{(g_1)}}(X_{(g_1)},𝐂),\beta H^{q\iota _{(g_2)}}(X_{(g_2)},𝐂),\gamma H^{2npq\iota _{(g_3)}}(X_{(g_3)},𝐂)`$. Once the three point function is defined, the cup product is defined by the equation
$$<\alpha _{orb}\beta ,\gamma >_{orb}=<\alpha ,\beta ,\gamma >_{orb}.$$
$`(2.11)`$
for arbitrary $`\gamma `$.
There is a Dolbeaut version of orbifold cohomology ring (Dolbeaut orbifold cohomology ring) with identical construction. We refer reader to \[CR\] for detail.
Next, we consider the bundle over each sector and its pull-back. By \[CR1\], this is a very subtle problem and one has to be careful. We will give explicit construction in our case and refer reader to general theory in \[CR1\].
Now, let’s examine the orbifold structure of twisted sectors more carefully. Suppose that $`pX_{(g)}`$ and an orbifold chart of $`pX`$ is $`U_p/G_p`$. By \[CR1\](Lemma 3.1.1), an orbifold chart of $`X_{(g)}`$ can be described as follows. Choose a representative of $`(g)_{G_p}`$, say $`g_p`$. Then, a local orbifold chart of $`p`$ is $`U_{g_p}/C(g_p)`$, where $`U_{g_p}`$ is the fixed point loci of $`g_p`$ and $`C(g_p)`$ is the centralizer. In general, $`C(g_p)`$ may not acts freely on $`U_{g_p}`$. The patching map of $`X_{(g)}`$ is defined in the same way. More generally, $`X_{(𝐠)}`$ also has a structure of a orbifold given by $`(U_{𝐠_p},C(𝐠_p))`$, where $`U_{𝐠_p},C(𝐠_p)`$ are the fixed point loci and centralizer of $`𝐠_p`$. We denote the corresponding patching map by injection $`i_{𝐠,pq}=(\stackrel{~}{i}_{pq,(𝐠)},i_{\mathrm{\#},pq,(𝐠)})`$.
Now, an orbifold-bundle $`f:EX_{(𝐠)}`$ is a continuous map between topological space such that $`E`$ has a structure of orbifold as follows. Suppose that $`pX_{(𝐠)}`$. $`E`$ is covered by chart of the form $`(U_{𝐠_p}\times 𝐑^n,C(𝐠_p))`$ such that the restriction of $`f`$ is the projection $`U_{𝐠_p}\times 𝐑^nU_{𝐠_p}`$ equivariant under $`C(𝐠_p)`$. For any injection $`i:(U_{𝐠_q},C(𝐠_q))(U_{𝐠_p},C(𝐠_p))`$, there is an injection between charts $`(U_{𝐠_q}\times 𝐑^n,G_{𝐠_q})(U_{𝐠_p}\times 𝐑^n,G_{𝐠_p})`$ given by $`(\stackrel{~}{i}_{pq,(𝐠)}\times g_i,i_{\mathrm{\#},pq,(𝐠)}),`$ where $`g_i:U_{𝐠_q}Aut(𝐑^n)`$ satisfies the condition $`g_{ji}=g_jg_i`$. The last condition is to ensure that the equivalent injections on $`X_{(𝐠)}`$ implies the equivalence between corresponding injections on total space of bundle.
Next, we consider evaluation map $`e_i:X_{(𝐠)}X_{(g)}`$. A crucial observation is that we can choose the orbifold charts of $`X_{(𝐠)},X_{(g)}`$ silmontaneously. Suppose that $`pX_{(𝐠)}`$. We first choose $`𝐠_p`$. Then, it gives a natural choice for $`g_{i,p}`$. Similarly, an injection $`i_{pq,(𝐠)}`$ between the charts of $`X_{(𝐠)}`$ gives a natural choice of injection $`\lambda (i_{pq,(𝐠)})=i_{pq,(g_i)}`$ with the property $`\lambda (ji)=\lambda (j)\lambda (i)`$. The evaluation map is interpreted as an inclusion
$$e_{i,p}:U_{𝐠_p}U_{g_{i,p}},e_{\mathrm{\#},i,p}:C(𝐠_p)C(g_{i,p}).$$
$`(2.11)`$
Following \[CR1\], we say that $`e_i`$ is a good map. By \[CR1\], if $`EX_{(g_i)}`$ is a orbifold-bundle, $`e^{}EX_{(𝐠)}`$ is a orbifold-bundle.
Note that the direct sum and tensor product of orbifold-bundles is still a orbifold bundle. Moreover, all the differential geometric constructions such as differential form, connection and curvature work over orbifold-bundle.
### 2.2 Inner Local system and twisted orbifold cohomology ring
Now, we introduce the notion of inner local system for orbifold. Definition 2.1: Suppose that $`X`$ is an orbifold (almost complex or not). An inner local system $`=\{L_{(g)}\}_{gT_1}`$ is an assignment of a flat complex orbifold line bundle over
$$L_{(g)}X_{(g)}$$
to each sector $`X_{(g)}`$ satisfying the compatibility condition
$`L_{(1)}=1`$ is trivial.
$`I^{}L_{(g^1)}=L_{(g)}^1.`$
Over each $`X_{(𝐠)}`$ with $`(𝐠)T_3^o`$, $`_ie_i^{}L_{(g_i)}=1`$.
If $`X`$ is a complex orbifold, we assume that $`L_{(g)}`$ is holomorphic. Definition 2.2: Given an inner local system $``$, we define the twisted orbifold cohomology
$$H_{orb}^{}(X,)=_{(g)}H^{2\iota _{(g)}}(X_{(g)},L_{(g)}).$$
Definition 2.3: Suppose that $`X`$ is a closed complex orbifold and $``$ is an inner local system. We define Dolbeault cohomology groups
$$H_{orb}^{p,q}(X,)=_{(g)}H^{p\iota _{(g)},q\iota _{(g)}}(X_{(g)};L_{(g)}).$$
$`(2.12)`$
Proposition 2.4: If $`X`$ is a Kahler orbifold, we have Hodge decomposition
$$H_{orb}^k(X,)=_{k=p+q}H_{orb}^{p,q}(X,).$$
$`(2.13)`$
Proof: Note that each sector $`X_{(g)}`$ is a Kähler orbifold. The proposition follows by applying the ordinary Hodge theorem with twisted coefficients to each sector $`X_{(g)}`$. $`\mathrm{}`$
The property (2) of Definition 2.1 can be used to show that Poincare pairing defined in (2.6) can be adopted to twisted orbifold cohomology. Definition (Poincaré duality) 2.5: Suppose that $`X`$ is a $`2n`$-dimensional closed almost complex orbifold. We define a pairing
$$<>_{orb,}:H_{orb}^d(X,)H_{orb}^{2nd}(X,)𝐂.$$
$`(2.14)`$
as the direct sum of
$$<>_{orb,}^{(g)}:H^{d2\iota _{(g)}}(X_{(g)},L_{(g)})H^{2nd2\iota _{(g^1)}}(X_{(g^1)},L_{(g^1)})𝐂$$
$`(2.15)`$
defined by
$$<\alpha ,\beta >_{orb,}^{(g)}=_{X_{(g)}}\alpha I^{}\beta .$$
$`(2.16)`$
Note that $`L_{(g)}I^{}L_{(g^1)}=1`$. Hence the integral (2.6) makes sense. Theorem 2.6: The pairing $`<>_{orb,}`$ is nondegenerate. Proof: The proof follows from ordinary Poincare duality on $`X_{(g)}`$ with twisted coefficient.
There is also a version of Poincaré duality for twisted Dolbeault cohomology. Suppose that $`X`$ is a closed complex orbifold of complex dimension $`n`$. Then $`X_{(g)}`$ is a closed complex orbifold. Definition 2.7: We define a pairing
$$<>_{orb,}:H_{orb}^{p,q}(X,)H_{orb}^{np,nq}(X,)𝐂.$$
$`(2.17)`$
as the direct sum of
$$<>_{orb,}^{(g)}:H^{p\iota _{(g)},q\iota _{(g)}}(X_{(g)},L_{(g)})H^{np\iota _{(g^1)},nq\iota _{(g^1)}}(X_{(g^1)},L_{(g^1)})𝐂$$
$`(2.18)`$
defined by
$$<\alpha ,\beta >_{orb,}^{(g)}=_{X_{(g)}}\alpha I^{}\beta .$$
$`(2.19)`$
Theorem 2.8: The pairing (2.17) is nondegenerate.
The property (3) of Definition 2.1 shows that the integral (2.10) makes sense for twisted orbifold cohomology classes. The same construction of \[CR\] goes through without change. We can define a twisted orbifold product $`_{orb,}`$ in the same fashion. The same proof as in \[CR\] yields
Theorem 2.9: Let $`X`$ be a closed almost complex orbifold with almost complex structure $`J`$ and $`dim_𝐂X=n`$. There is a cup product $`_{orb,}:H_{orb}^p(X;)\times H_{orb}^q(X;)H_{orb}^{p+q}(X;)`$ for any $`0p,q2n`$ such that $`p+q2n`$, which has the following properties:
1. The total twisted orbifold cohomology group $`H_{orb}^{}(X;)=_{0d2n}H_{orb}^d(X;)`$ is a ring with unit $`e_X^0H_{orb}^0(X;)`$ under $`_{orb,}`$, where $`e_X^0`$ is the Poincareé dual to the fundamental class of nontwisted sector.
2. Restricted to each $`H_{orb}^d(X;)\times H_{orb}^{2nd}(X;)H_{orb}^{2n}(X;)=H^{2n}(X,𝐂)`$,
$$_X\alpha _{orb,}\beta =<\alpha ,\beta >_{orb,}.$$
$`(2.20)`$
3. The cup product $`_{orb,}`$ is invariant under deformations of $`J`$.
4. When $`X`$ is of integral degree shifting numbers, the total twisted orbifold cohomology group $`H_{orb}^{}(X;)`$ is integrally graded, and we have supercommutativity
$$\alpha _1_{orb,}\alpha _2=(1)^{\mathrm{deg}\alpha _1\mathrm{deg}\alpha _2}\alpha _2_{orb,}\alpha _1.$$
5. Restricted to the nontwisted sectors, i.e., the ordinary cohomology $`H^{}(X;𝐂)`$, the cup product $`_{orb,}`$ equals the ordinary cup product on $`X`$.
6. $`_{orb,}`$ is associative.
Similarly, we also have a holomorphic version.
Theorem 2.10: Let $`X`$ be an n-dimensional closed complex orbifold with complex structure $`J`$. The orbifold cup product
$$_{orb,}:H_{orb}^{p,q}(X;)\times H_{orb}^{p^{},q^{}}(X;)H_{orb}^{p+p^{},q+q^{}}(X;)$$
has the following properties:
1. The total orbifold Dolbeault cohomology group is a ring with unit $`e_X^0H_{orb}^{0,0}(X;)`$ under $`_{orb,}`$, where $`e_X^0`$ is the class represented by the equal one constant function on $`X`$.
2. Restricted to each $`H_{orb}^{p,q}(X;)\times H_{orb}^{np,nq}(X;)H_{orb}^{n,n}(X;)`$, the integral $`_X\alpha _{orb,}\beta `$ equals the Poincare pairing $`<\alpha ,\beta >_{orb,}`$.
3. The cup product $`_{orb,}`$ is invariant under the deformation of $`J`$.
4. When $`X`$ is of integral degree shifting numbers, the total twisted orbifold Dolbeault cohomology group of $`X`$ is integrally graded, and we have supercommutativity
$$\alpha _1_{orb,}\alpha _2=(1)^{\mathrm{deg}\alpha _1\mathrm{deg}\alpha _2}\alpha _2_{orb,}\alpha _1.$$
5. Restricted to the nontwisted sectors, i.e., the ordinary Dolbeault cohomologies $`H^,(X;𝐂)`$, the cup product $`_{orb,}`$ equals the ordinary wedge product on $`X`$.
6. The cup product is associative.
7. When $`X`$ is Kähler, the cup product $`_{orb,}`$ coincides with the twisted orbifold cup product over the twisted orbifold cohomology groups $`H_{orb}^{}(X;)`$ under the relation
$$H_{orb}^r(X;)𝐂=_{p+q=r}H_{orb}^{p,q}(X;).$$
Remark 2.11: If $`X`$ is open, we can define usual twisted orbifold cohomology $`H_{orb}^{}(X,)`$ and twisted orbifold cohomology with compact support $`H_{orb,c}^{}(X,)`$ in the same fashion. The Poincare paring should be understood as the paring between $`H_{orb}^d(X,)`$ and $`H_{orb,c}^{2nd}(X,)`$.
## 3 Orbifold fundamental group and discrete torsion
First, we recall the definition of orbifold fundamental group. Definition 3.1: A smooth map $`f:YX`$ is an orbifold cover iff (1) each $`pY`$ has a neighborhood $`U_p/G_p`$ such that the restriction of $`f`$ to $`U_p/G_p`$ is isomorphic to a map $`U_p/G_pU_p/\mathrm{\Gamma }`$ such that $`G_p\mathrm{\Gamma }`$ is a subgroup. (2) Each $`qX`$ has a neighborhood $`U_q/G_q`$ for which each component of $`f^1(U_q/G_q)`$ is isomorphic to $`U_q/\mathrm{\Gamma }^{}`$ such that $`\mathrm{\Gamma }^{}G_q`$ is a subgroup. An orbifold universal cover $`f:YX`$ of $`X`$ has the property: (i) $`Y`$ is connected; (ii)if $`f^{}:Y^{}X`$ is an orbifold cover, then there exists an orbifold cover $`h:YY^{}`$ such that $`f=f^{}h`$. If $`Y`$ exists, we call $`Y`$ the orbifold universal cover of $`X`$ and the group of deck translations the orbifold fundamental group $`\pi _1^{orb}(X)`$ of $`X`$. By Thurston \[T\], an orbifold universal cover exists. It is clear from the definition that the orbifold universal cover is unique. Suppose that $`f:YX`$ is an orbifold universal cover. Then
$$f:Yf^1(\mathrm{\Sigma }X)X\mathrm{\Sigma }X$$
$`(3.1)`$
is an honest cover with $`G=\pi _1^{orb}(X)`$ as covering group, where $`\mathrm{\Sigma }`$ is the singular loci of $`X`$. Therefore, $`X=Y/G`$ and there is a surjective homomorphism
$$p_f:\pi _1(X\mathrm{\Sigma }X)G.$$
$`(3.2)`$
In general, (3.1) is not a universal covering. Hence, $`p_f`$ is not an isomorphism.
Remark 3.2: Suppose that $`X=Z/G`$ for an orbifold $`Z`$ and $`Y`$ is the orbifold universal cover of $`Z`$. By the definition, $`Y`$ is an orbifold universal cover of $`X`$. It is clear that there is a short exact sequence
$$1\pi _1(Z)\pi ^{orb}(X)G1.$$
$`(3.3)`$
Example 3.3: Consider the Kummer surface $`T^4/\tau `$ where $`\tau `$ is the involution
$$\tau (e^{it_1},e^{it_2},e^{it_3},e^{it_4})=(e^{it_1},e^{it_2},e^{it_3},e^{it_4}).$$
$`(3.4)`$
The universal cover is $`𝐑^4`$. The group $`G`$ of deck translations is generated by translations $`\lambda _i`$ by an integral point and the involution
$$\tau :(t_1,t_2,t_3,t_4)(t_1,t_2,t_3,t_4).$$
It is easy to check that
$$G=\{\lambda _i(i=1,2,3,4),\tau |\tau ^2=1,\tau \lambda _i=\lambda _i^1\tau ,\}$$
$`(3.5)`$
where $`\lambda _i`$ represents translation and $`\tau `$ represents involution. Example 3.4: Let $`T^6=𝐑^6/\mathrm{\Gamma }`$ where $`\mathrm{\Gamma }`$ is the lattice of integral points. Consider $`𝐙_2^2`$ acting on $`T^6`$ lifted to an action on $`𝐑^6`$ as
$$\sigma _1(t_1,t_2,t_3,t_4,t_5,t_6)=(t_1,t_2,t_3,t_4,t_5,t_6)$$
$$\sigma _2(t_1,t_2,t_3,t_4,t_5,t_6)=(t_1,t_2,t_3,t_4,t_5,t_6)$$
$$\sigma _3(t_1,t_2,t_3,t_4,t_5,t_6)=(t_1,t_2,t_3,t_4,t_5,t_6).$$
This example was considered by Vafa-Witten \[VW\]. The orbifold fundamental group
$$\pi _1^{orb}(T^6/𝐙_2^2)=\{\tau _i(1i6),\sigma _j(1j3)|$$
$$\sigma _i^2=1,\sigma _1\tau _i=\tau _i^1\sigma _1(i5,6),\sigma _2\tau _i=\tau _i^1\sigma _2(i3,4),\sigma _3\tau _i=\tau _i^1\sigma _3(i1,2)\}.$$
$`(3.6)`$
The following example was taken from \[SC\] Example 3.5: Consider the orbifold Riemann surface $`\mathrm{\Sigma }_g`$ of genus $`g`$ and $`n`$-orbifold points $`𝐳=(x_1,\mathrm{},x_n)`$ with orders $`k_1,\mathrm{},k_n`$. Then,
$$\pi _1^{orb}(\mathrm{\Sigma }_g)=\{\lambda _i(i2g),\sigma _i(in)|\sigma _1\mathrm{}\sigma _n\underset{i}{}[\lambda _{2i1},\lambda _{2i}]=1,\sigma _i^{k_i}=1\},$$
$`(3.7)`$
where $`\lambda _i`$ are the generators of $`\pi _1(\mathrm{\Sigma }_g)`$ and $`\sigma _i`$ are the generators of $`\mathrm{\Sigma }_g𝐳`$ represented by a loop around each orbifold point.
Note that $`\pi _1^{orb}(\mathrm{\Sigma }_g)`$ is just $`\pi _1(\mathrm{\Sigma }_g𝐳)`$ modulo by the relation $`\sigma _i^{k_i}=1`$. This suggests that one can first take the cover of $`\mathrm{\Sigma }_g𝐳`$ induced by $`\pi _1^{orb}(\mathrm{\Sigma })`$. The relation $`\sigma _i^{k_i}=1`$ implies that the preimage of the punctured disc around $`x_i`$ is a punctured disc. Then we can fill in the center point to obtain the orbifold universal cover. Definition 3.6: Suppose that $`SX`$ is a connected component of singular loci. A local discrete torsion $`\alpha _S`$ at $`S`$ is defined as a cohomology class $`\alpha _SH^2(\pi _1^{orb}(U(S)),U(1))=H^2(B\pi _1^{orb}(U(S)),U(1))`$, where $`U(S)`$ is a small open neighborhood of $`S`$. A global discrete torsion $`\alpha =\{\alpha _S\}`$ is an assignment of a local discrete torsion to each connected component of singular loci. If $`X=Z/G`$ for a finite group $`G`$, by Remark 3.2, there is a surjective homomorphism
$$\pi :\pi _1^{orb}(X)G.$$
$`\pi `$ induces a homomorphism
$$\pi ^{}:H^2(G,U(1))H^2(\pi _1^{orb}(X),U(1)).$$
$`(3.8)`$
Hence, an element of $`H^2(G,U(1))`$ induces a discrete torsion of $`X`$.
They are many ways to define $`H^2(G,U(1))`$. The definition $`H^2(G,U(1))=H^2(BG,U(1))`$ is a very useful definition for computation since we can use algebro-topological machinery. However, we can also take the original definition in terms of cocycles. A 2-cocycle is a map $`\alpha :G\times GU(1)`$ satisfying
$$\alpha _{g,1}=\alpha _{1,g}=1,\alpha _{g,hk}\alpha _{h,k}=\alpha _{g,h}\alpha _{gh,k},$$
$`(3.9)`$
for any $`g,h,kG`$. We denote the set of two cocycles by $`Z^2(G,U(1))`$. For any map $`\rho :GU(1)`$ with $`\rho _1=1`$, its coboundary is defined by formula
$$(\delta \rho )_{g,h}=\rho _g\rho _h\rho _{gh}^1.$$
$`(3.10)`$
Let $`B^2(G,U(1))`$ be the set of coboundaries. Then, $`H^2(G,U(1))=Z^2(G,U(1))/B^2(G,U(1))`$. $`H^2(G,U(1))`$ naturally appears in many important places of mathematics. For example, it classifies the group extension of $`G`$ by $`U(1)`$. If we have a unitary projective representation of $`G`$, it naturally induces a class of $`H^2(G,U(1))`$. In many instances, this class completely classifies the projective unitary representation. In fact, it is in this context that discrete torsion arises in orbifold string theory.
Definition 3.7: For each 2-cocycle $`\alpha `$, we define its phase
$$\gamma (\alpha )_{g,h}=\alpha _{g,h}\alpha _{h,g}^1.$$
$`(3.11)`$
It is clear that $`\gamma (\alpha )_{g,g}=1,\gamma (\alpha )_{g,h}=\gamma (\alpha )_{h,g}^1.`$ Lemma 3.8: Suppose that $`gh=hg,gk=kg.`$ Then
$`\gamma (\delta \rho )_{g,h}=1`$.
$`\gamma (\alpha )_{g,hk}=\gamma (\alpha )_{g,h}\gamma (\alpha )_{g,k}.`$
The (2) implies $`L_g^\alpha =\gamma _{g,}:C(g)U(1)`$ is a group homomorphism. We call $`L_g^\alpha `$ a $`\alpha `$-twisted character. Proof: (1) is obvious. For (2),
$$\begin{array}{ccc}\gamma (\alpha )_{g,hk}\hfill & =\hfill & \alpha _{g,hk}\alpha _{hk,g}^1\hfill \\ & =\hfill & \alpha _{g,hk}\alpha _{gh,k}^1\alpha _{hg,k}\alpha _{h,gk}^1\alpha _{h,kg}\alpha _{hk,g}^1\hfill \\ & =\hfill & \alpha _{g,h}\alpha _{h,k}^1\alpha _{g,k}\alpha _{h,g}^1\alpha _{h,k}\alpha _{k,g}^1\hfill \\ & =\hfill & \gamma (\alpha )_{g,h}\gamma (\alpha )_{g,k}\hfill \end{array}$$
Next, we calculate discrete torsions for some groups. We first consider the case of finite abelian group $`G`$. In this case $`H^i(G,𝐐)=0`$ for $`i0`$. The exact sequence
$$0𝐙𝐂𝐂^{}1$$
implies that $`H^2(G,U(1))=H^2(G,𝐂^{})=H^3(G,𝐙)`$. By universal coefficient theorem, $`H^3(G,𝐙)=H_2(G,𝐙).`$ Example 3.9 $`G=𝐙/n\times 𝐙/m`$: Notes that $`H^2(G,U(1))=H_2(G,𝐙)=𝐙/n𝐙/m=Z_{gcd(n,m)}`$. In this case, one can write down the phase of discrete torsion explicitly \[VW\]. Let $`\xi (\zeta )`$ be $`n(m)`$-root of unity. Any element of $`𝐙/n\times 𝐙/m`$ can be written as $`(\xi ^a,\zeta ^b)`$. Let $`p=gcd(n,m)`$. The phase of a discrete torsion can be written as
$$\gamma _{(\xi ^a,\zeta ^b),(\xi ^a^{},\zeta ^b^{})}=\omega _p^{m(ab^{}ba^{})}$$
with $`\omega _p=e^{2\pi i/p},m=1,\mathrm{},p.`$ There are $`p`$-phases for $`p`$-discrete torsions. It is trivial to generalize this construction to an arbitrary finite abelian group.
## 4 Discrete torsion and local system
Suppose that $`f:YX`$ is the orbifold universal cover and $`G`$ is the orbifold fundamental group which acts on $`Y`$ such that $`X=Y/G`$. Suppose $`X_{(g)}`$ is a sector (twisted or nontwisted) of $`X`$. For any $`qX`$, choose an orbifold chart $`U_q/G_q`$ satisfying Definition 3.1. A component of $`f^1(U_q/G_q)`$ is of the form $`U_q/\mathrm{\Gamma }^{}`$ for $`\mathrm{\Gamma }^{}G_q`$. It is clear that $`G_q/\mathrm{\Gamma }^{}`$ is a subgroup of the orbifold fundamental group. Therefore, we obtain a group homomorphism
$$\varphi _q:G_q\pi _1^{orb}(X).$$
$`(4.1)`$
It is easy to check that a different choice of component of $`f^1(U_q/G_q)`$ or a different choice of $`qX_{(g)}`$ induces a homomorphism differing by a conjugation. Therefore, there is a unique map from the conjugacy classes of $`G_q`$ to the conjugacy classes of $`\pi _1^{orb}(X)`$. Definition 4.1: We call $`X_{(g)}`$ a dormant sector if $`\varphi _p(g)=1`$. If $`X_{(g)}`$ is a dormant sector, we define $`L_{(g)}=1`$. It will not receive any correction from discrete torsion. Non-dormant sectors are of the form $`Y_g/C(g)`$, where $`Y_g\mathrm{}`$ is the fixed point loci of $`1g\pi _1^{orb}(X)`$. $`Y_g`$ is a smooth suborbifold of $`Y`$. It is clear that $`Y_{h^1gh}`$ is diffeomorphic to $`Y_g`$ by the action of $`h`$. By abusing the notation, we denote the twisted sector $`Y_g/C(g)`$ by $`X_{(g)}`$, where $`C(g)`$ is the centralizer of $`g`$.
Let $`\alpha `$ be a global discrete torsion. Suppose that $`S`$ is the connected component of singular loci containing the image of $`X_{(g)}`$ in $`X`$. We choose a small open neighborhood $`U(S)`$ of $`S`$ and suppose that local discrete torsion is $`\alpha _S`$. We replace $`X`$ by $`U(S)`$ in above construction. We can use $`L_g^\alpha `$ to define a flat complex orbifold line-bundle
$$L_g=Y_g\times _{L_g^\alpha }𝐂$$
over $`X_{(g)}`$. Lemma 4.2:
$`L_{tgt^1}`$ is isomorphic to $`L_g`$ by the map
$$t\times Id:Y_g\times 𝐂Y_{tgt^1}\times 𝐂.$$
$`(4.3)`$
Hence, we can denote $`L_g`$ by $`L_{(g)}`$.
$`L_{(g)}^1=L_{(g^1)}.`$
When we restrict to $`X_{(g_1,\mathrm{},g_k)}=Y_{g_1}\mathrm{}Y_{g_k}/C(g_1,\mathrm{},g_k)`$, $`L_{(g_1,\mathrm{},g_k)}=L_{(g_1)}\mathrm{}L_{(g_k)}`$, where $`L_{(g_1,\mathrm{},g_k)}=Y_{g_1}\mathrm{}Y_{g_k}\times _{\gamma _{g_1\mathrm{}g_k}}𝐂.`$
Proof: Recall that there is an isomorphism
$$t_\mathrm{\#}:C(g)C(tgt^1)$$
given by $`t_\mathrm{\#}(h)=tht^1`$. The map
$$t:Y_gX_{tgt^1}$$
is $`t_\mathrm{\#}`$-equivariant. By Lemma 3.8, $`\gamma _{tgt^1}(tht^1)=\gamma _g(h)`$ for $`hC(g)`$. Then,
$$(t\times Id)(hx,\gamma (h)(v))=(thx,\gamma _g(h)(v))=(tht^1tx,\gamma _{tgt^1}(tht^1)(v)).$$
$`(4.4)`$
Then we take the quotient by $`C(g),C(tgt^1)`$ respectively to get an isomorphism between $`L_g,L_{tgt^1}`$. (2) and (3) follow from the fact that for any $`hC(g_1,\mathrm{},g_k)`$,
$$\gamma (\alpha )_{g_1\mathrm{}g_k,h}=\gamma (\alpha )_{h,g_1\mathrm{}g_k}^1=\gamma (\alpha )_{h,g_1}^1\mathrm{}\gamma (\alpha )_{h,g_k}^1=\gamma (\alpha )_{g_1,h}\mathrm{}\gamma (\alpha )_{g_k,h}.$$
$`(4.5)`$
Theorem 4.3: $`_\alpha =\{L_{(g)}\}_{(g)T_1}`$ is an inner local system of $`X`$. Proof: Property (1) is obvious. The property (2) follows from Lemma 4.2. Let’s prove property (3). Consider the image $`𝐠^{}=(g_1^{},g_2^{},g_3^{})`$ of $`𝐠`$ in $`\pi _1^{orb}(X)`$ under the homomorphism (4.1). Then, we still have $`g_1^{}g_2^{}g_3^{}=1`$. There are three possibilities, (i) $`g_1^{}=g_2^{}=g_3^{}=1`$ and there is nothing to prove in this case; (ii) $`g_3^{}=1,g_2^{}=(g^{})_1^1`$ is nontrivial; (iii) $`g_1^{},g_2^{},g_3^{}`$ are all nontrivial.
For the second case, let $`g=g_1^{}`$. We have the following factorization
$$e_1\times e_2\times e_3:X_{(𝐠)}X_{(g_1,g_2)}\times X_{(g_3)}X_{(g_1)}\times X_{(g_2)}\times X_{(g_3)}.$$
However, $`X_{(g_1,g_2)}=Y_gY_{(g^1)}/C(g,g^1)=Y_g/C(g).`$ Moreover, over $`X_{g_1,g_2}`$
$$e_1^{}L_{(g)}e_2^{}L_{(g^1)}=L_{(g)}I^{}L_{(g^1)}=1.$$
$`(4.6)`$
In the third case, $`X_{(𝐠)}=Y_{g_1}Y_{g_2}Y_{g_3}/C(g_1,g_2,g_3)`$. The proof follows from Lemma 4.2 (3). Definition 4.4: Suppose that $`\alpha `$ is a global discrete torsion. We define the twisted orbifold cohomology $`H_{orb,\alpha }^{}(X,𝐂)=H_{orb}^{}(X,_\alpha ).`$
## 5 Examples
Only a few examples of global quotients have been computed by physicists \[VW\] \[D\]. It is still a very important problem to develop general machinery to compute discrete torsion and twisted orbifold cohomology. Here we compute five examples. First two have nontrivial discrete torsion. One is a global quotient and another one is a non-global quotient. The second example has the phenomenon that the most of twisted sectors are dormant sectors. The third one is Joyce example, where there is no nontrivial discrete torsion. However, there are nontrivial local systems. We will compute twisted orbifold cohomology given by nontrivial local systems to match Joyce’s desingularizations. Orbifold cohomology is strongly intertwine with group theory. We demonstrate it in last two examples. Example 5.1 $`T^4/𝐙_2\times 𝐙_2`$: Here, $`T^4=𝐂^2/`$, where $``$ is the lattice of integral points. Suppose that $`g,h`$ are generators of the first and the second factor of $`𝐙_2\times 𝐙_2`$. The action of $`𝐙_2\times 𝐙_2`$ on $`T^4`$ is defined as
$$g(z_1,z_2)=(z_1,z_2),h(z_1,z_2)=(z_1,z_2).$$
$`(5.1)`$
The fixed point locus of $`g`$ is 4 copies of $`T^2`$. When we divide it by the remaining action generated by $`h`$, we obtain twisted sectors consisting of 4 copies of $`S^2`$. The degree shifting number for these twisted sectors is $`\frac{1}{2}`$. For the same reason, the fixed point locus of $`h`$ give twisted sectors consisting of $`4`$ copies of $`S^2`$ with degree shifting number $`\frac{1}{2}`$. The fixed point locus of $`gh`$ is 16 points, which are fixed by the whole group. The degree shifting number of the 16 points is $`1`$. An easy calculation shows that nontwisted sector contributes one generator to degree 0, 4 orbifold cohomology and two generators to degree 2 orbifold cohomology and no other. Using this information, we can compute the ordinary orbifold cohomology group
$$b_{orb}^0=b_{orb}^4=1,b_{orb}^1=b_{orb}^3=8,b_{orb}^2=18.$$
$`(5.2)`$
By example 2.10, $`H^2(𝐙_2\times 𝐙_2,U(1))=𝐙_2`$. By Remark 2.2, the nontrivial generator of $`H^2(𝐙_2\times 𝐙_2,U(1))`$ induces a discrete torsion $`\alpha `$. Next, we compute the twisted orbifold cohomology $`H_{orb,\alpha }^{}(T^4/𝐙_2\times 𝐙_2,𝐂)`$. Note that $`\gamma (\alpha )_{gh,g}=\gamma (\alpha )_{gh,h}=1`$. Hence, the flat orbifold-bundles over the twisted sectors given by 16 fixed points of $`gh`$ are nontrivial. Therefore, they contribute nothing to twisted orbifold cohomology. For two dimensional twisted sectors, let’s consider a component of fixed point locus of $`g`$. By the previous description, it is $`T^2`$. $`h`$ acts on $`T^2`$. Then the twisted sector $`S^2=T^2/\{h\}`$. We observe that the flat orbifold line bundle over $`S^2`$ is constructed as $`L=T^2\times _{L_g^\alpha }𝐂`$. Hence $`H^{}(S^2,L)`$ is isomorphic to the space of invariant cohomology of $`T^2`$ under the action of $`h`$ twisted by $`\gamma (\alpha )_g`$ as $`h(\beta )=\gamma (\alpha )_{g,h}h^{}\beta `$. By example 2.10, $`\gamma (\alpha )_{g,h}=1`$. The invariant cohomology is $`H^1(T^2,𝐂)`$. Using the degree shifting number to shift up its degree, we obtain the twisted orbifold cohomology
$$b_{orb,\alpha }^0=b_{orb,\alpha }^4=1,b_{orb,\alpha }^1=b_{orb,\alpha }^3=0,b_{orb,\alpha }^2=18.$$
$`(5.3)`$
Example 5.2 $`WP(2,2d_1)\times WP(2,2d_2)`$ ($`d_1,d_2>1,(d_1,d_2)=1`$): Here, $`WP(2,2d)`$ is the weighted projective space of weighted $`(2,2d)`$. $`WP(2,2d_1)\times WP(2,2d_2)`$ is not a global quotient unless $`d_1=d_2=1`$. In fact, its orbifold universal cover is $`WP(1,d_1)\times WP(1,d_2)`$ and $`WP(2,2d_1)\times WP(2,2d_2)=WP(1,d_1)\times WP(1,d_2)/𝐙_2\times 𝐙_2`$. Hence, the orbifold fundamental group is $`𝐙_2\times 𝐙_2`$. Therefore, there is a nontrivial discrete torsion $`\alpha H^2(𝐙_2\times 𝐙_2,U(1))`$.
Next, we describe the twisted sectors. Suppose that $`p=[0,1],q=[1,0]WP(1,d_1)`$. We also use $`p,q`$ to denote its image in $`WP(2,2d_1)`$. We use $`p^{},q^{}`$ to denote the corresponding points in $`WP(1,d_2),WP(2,2d_2)`$. $`\{p\}\times WP(2,2d_2),\{p^{}\}\times WP(2,2d_1)`$ give rise to two twisted sectors with degree shifting number $`\frac{1}{2}`$. $`\{q\}\times WP(2,2d_2),\{q^{}\}\times WP(2,2d_1)`$ give rise to $`2d_11,2d_21`$ many twisted sectors with degree shifting numbers $`\frac{i}{2d_1},\frac{j}{2d_2}`$ for $`1i2d_11,1j2d_21`$. $`\{p\}\times \{p^{}\}`$ give rise to a twisted sector with degree shifting number $`1`$. $`\{p\}\times \{q^{}\}`$ give rise to $`2d_21`$-many twisted sectors with degree shifting numbers $`\frac{1}{2}+\frac{i}{2d_2}`$ for $`1i2d_21`$. $`\{q\}\times \{p^{}\}`$ give rise to $`2d_11`$-many twisted sectors with degree shifting numbers $`\frac{1}{2}+\frac{i}{2d_1}`$ for $`1i2d_11`$. $`\{q\}\times \{q^{}\}`$ give rise to $`4d_1d_21`$-many twisted sectors with degree shifting numbers $`\frac{i}{2d_1}+\frac{j}{2d_2}`$ for all $`i,j`$ except $`(i,j)=(0,0)`$. Using this information, we can write down ordinary orbifold cohomology
$$b_{orb}^0=b_{orb}^4=1,b_{orb}^1=b_{orb}^3=6,b_{orb}^2=6$$
$$b_{orb}^{\frac{i}{d_1}}=b_{orb}^{\frac{i}{d_2}}=1,b_{orb}^{1+\frac{i}{d_1}}=b_{orb}^{1+\frac{i}{d_2}}=3,b_{orb}^{2+\frac{i}{d_1}}=b_{orb}^{2+\frac{i}{d_2}}=2,1id_11,1jd_21$$
$$b_{orb}^{\frac{i}{d_1}+\frac{j}{d_2}}=1,0i2d_11,0j2d_2,(i,j)(0,0),(d_1,d_2).$$
$`(5.4)`$
Next, we compute $`H_{orb,\alpha }^{}`$. In this example, the most of twisted sectors are dormant sectors. To find nondormant sectors, recall that $`WP(2,2d_1)\times WP(2,2d_2)=WP(1,d_1)\times WP(1,d_2)/Z_2\times Z_2`$. Let $`g`$ be the generator of the first factor and $`h`$ be the generator of the second factor. The fixed points of $`g`$ is $`\{p,q\}\times WP(1,d_2)`$. We have two nondormant sectors obtained by modulo the remaining action generated by $`h`$. However, $`\gamma (\alpha )_{g,h}=1`$. There is no invariant cohomology of $`WP(1,d_2)`$ under the action of $`h`$ twisted by $`L_g^\alpha `$. Hence, these two nondormant twisted sectors give no contribution to twisted orbifold cohomology. Their degree shifting numbers are 1. For the same reason, $`WP(1,d_1)\times \{p^{},q^{}\}/g`$ gives no contribution to twisted orbifold cohomology. The fixed point locus of $`gh`$ consists of 4 points which give 4 nondormant sectors. Again, their degree shifting numbers are 1. As we saw in last example, their flat orbifold bundles are nontrivial. Hence, they give no contribution to twisted orbifold cohomology. Therefore, the twisted orbifold cohomology is
$$b_{orb,\alpha }^0=b_{orb,\alpha }^4=1,b_{orb,\alpha }^1=b_{orb,\alpha }^3=2,b_{orb,\alpha }^2=2$$
$$b_{orb,\alpha }^{\frac{i}{d_1}}=b_{orb,\alpha }^{\frac{i}{d_2}}=1,b_{orb,\alpha }^{1+\frac{i}{d_1}}=b_{orb,\alpha }^{1+\frac{i}{d_2}}=3,b_{orb}^{2+\frac{i}{d_1}}=b_{orb,\alpha }^{2+\frac{i}{d_2}}=2,1id_11,1jd_21$$
$$b_{orb,\alpha }^{\frac{i}{d_1}+\frac{j}{d_2}}=1,0i2d_11,0j2d_2,(i,j)(0,0),(d_1,d_2).$$
$`(5.5)`$
Example 5.3 $`T^6/𝐙_4`$: Here, $`T^6=𝐂^3/`$, where $``$ is the lattice of integral points. The generator of $`𝐙_4`$ acts on $`T^6`$ as
$$\kappa :(z_1,z_2,z_3)(z_1,iz_2,iz_3).$$
$`(5.6)`$
This example has been studied by D. Joyce \[JO\], where he constructed five different desingularizations. However, there is no discrete torsion in the case which induces nontrivial orbifold cohomology.
First all, the nontwisted sector contributes one generator to $`H_{orb}^{0,0},H_{orb}^{3,3}`$, 5 generators to $`H_{orb}^{1,1},H_{orb}^{2,2}`$ and 2 generator to $`H_{orb}^{2,1},H_{orb}^{1,2}`$ The fixed point loci of $`\kappa ,\kappa ^3`$ are 16-points
$$\{(z_1,z_2,z_3)+:z_1\{0,\frac{1}{2},\frac{i}{2},\frac{1}{2}+\frac{i}{2}\},z_2,z_3\{0,\frac{1}{2}+\frac{i}{2}\}.$$
These points are fixed by $`Z_4`$. Therefore, they generate 32-twisted sectors in which 16 corresponds to the conjugacy class $`(\kappa )`$ and 16 corresponds to the conjugacy class $`(\kappa ^3)`$. The sector with conjugacy class $`(\kappa )`$ has degree shifting number 1. The sector with conjugacy class $`(\kappa ^3)`$ has degree shifting number 2.
The fixed point loci of $`\kappa ^2`$ is 16 copies of $`T^2`$, given by
$$\{(z_1,z_2,z_3)+:z_1𝐂,z_2,z_3\{0,\frac{1}{2},\frac{i}{2},\frac{1}{2}+\frac{i}{2}\}\}$$
Twelve of the 16-copies of $`T^2`$ fixed by $`\kappa ^2`$ are identified in pairs by the action of $`\kappa `$, and these contribute 6 copies of $`T^2`$ to the singular set of $`T^6/Z_4`$. On the remaining 4 copies $`\kappa `$ acts as $`1`$, so these contribute 4 copies of $`S^2=T^2/\{\pm 1\}`$ to singular set. The degree shifting number of these 2-dimensional twisted sectors is 1.
Next, we construct inner local systems. We start with two dimensional twisted sectors. Since $`\kappa ^2=\kappa ^2`$, the condition (2) of Definition 2.1 tells us that the flat orbifold line bundle $`L`$ over two dimensional sectors has the property $`L^2=1`$. Now, we assign trivial line bundle to all $`T^2`$-sectors and $`k(k=0,1,2,3,4)`$-many $`S^2=T^2/\{\pm 1\}`$-sectors. For the remaining $`S^2=T^2/\{\pm 1\}`$-sectors, we assign a flat orbifold line bundle $`T^2\times 𝐂/\{\pm 1\}`$. For the zero dimensional sectors, they are all points of two dimensional sectors. If we assign a trivial bundle on a two dimensional sector, we just assign trivial bundle to its point sectors. For these two dimensional sectors with nontrivial flat line bundle, we need to be careful to choose the flat orbifold line bundle on its point sectors to ensure the condition (3) of Definition 2.1. Suppose that $`\mathrm{\Sigma }`$ is one of 2-dimensional sectors supporting nontrivial flat orbifold line bundle. It contains 4 singular points which generate the point sectors. Let $`x`$ be one of 4-points. $`x`$ generates two sectors given by the conjugacy classes $`(\kappa ),(\kappa ^3)`$. For condition (3), we have to consider the conjugacy class of triple $`(g_1,g_2,g_3)`$ with $`g_1g_2g_3=1`$. The only nontrivial choices are $`(𝐠)=(\kappa ,\kappa ,\kappa ^2),(\kappa ^2,\kappa ^3,\kappa ^3)`$. The corresponding components of $`X_{(𝐠)}`$ are exactly the these singular points. Clearly, $`x`$ is a fixed by the whole group $`Z_4`$. The orbifold line bundle is given by the action of $`𝐙_4`$ on $`𝐂`$. Consider the component of $`X_{(𝐠)}`$ generated by $`x`$. The pull-back of flat orbifold line bundle from 2-dimensional sector ($`(\kappa ^2)`$-sector) is given by the action $`\kappa v=1`$. A moment of thought tells us that for sectors $`(\kappa ),(\kappa ^3)`$, we should assign a flat orbifold line bundle given by the action of $`𝐙_4`$ on $`𝐂`$ as $`\kappa v=iv.`$ It is easy to check that for above choices the condition (3) is satisfied for $`X_{(𝐠)}`$. Therefore, the twisted sectors given by $`(x,(\kappa )),(x,(\kappa ^3))`$ give no contribution to twisted orbifold cohomology. Suppose that the resulting local system is $`_k`$. For the sectors with trivial line bundle, they contribute $`6+k`$ generators to $`H_{orb}^{1,1},H_{orb}^{2,2}`$ and 6 generators to $`H_{orb}^{2,1},H_{orb}^{1,2}`$. Its point sectors contribute $`4k`$ generators to $`H_{orb}^{1,1},H_{orb}^{2,2}`$. The remaining sectors contribute $`4k`$ generators to $`H_{orb}^{2,1},H_{orb}^{1,2}`$. Its point sectors give no contribution. Moreover, the nontwisted sector contributes
$$h^{0,0}=h^{3,3}=2,h^{1,1}=5.$$
In summary, we obtain
$$dimH_{orb}^{0,0}(T^6/Z_4,_k)=dimH_{orb}^{3,3}(T^6/Z_4,_k)=1,dimH_{orb}^{1,1}(T^6/Z_4,_k)=dimH_{orb}^{2,2}(T^6/Z_4,_k)=11+5k,$$
$$dimH_{orb}^{1,2}(T^6/Z_4,_k)=dimH_{orb}^{2,1}(T^6/Z_4,_k)=12k$$
$`(5.7)`$
Our calculation matches the betti numbers of Joyce’s desingularizations.
The orbifold cohomology ring of following examples have been computed in \[CR\]. Here, we compute their twisted version. Example 5.4: Let’s consider the case that $`X`$ is a point with a trivial group action of $`G`$. Suppose that $`\alpha H^2(G,U(1))`$ is a discrete torsion. We want to compute $`H_{orb,\alpha }^{}(X,𝐂)`$. The twisted sector $`X_{(g)}`$ is a point with a group $`C(g)`$. It is obvious that $`H^0(X_{(g)},L_g^\alpha )=0`$ unless $`L_g^\alpha =1`$. Recall that a conjugacy class $`(g)`$ is $`\alpha `$-regular iff $`L_g^\alpha =1`$. Hence, only $`\alpha `$-regular class will contribute. Therefore, torbifold cohomology is generated by $`\alpha `$-regular conjugacy classes of elements of $`G`$. All the degree shifting numbers are zero. By the same argument as nontwisted case, as a ring, $`H_{orb,\alpha }^{}(X,𝐂)`$ is the center of twisted group algebra $`𝐂_\alpha [G]`$.
Example 5.5: Suppose that $`GSL(n,𝐂)`$ is a finite subgroup. Then, $`𝐂^n/G`$ is an orbifold. Suppose that $`\alpha H^2(G,U(1))`$ is a discrete torsion. For any $`gG`$, the fixed point set $`X_g`$ is a vector subspace and $`X_{(g)}=X_g/C(g)`$. By the definition, $`L_{(g)}=X_g\times \gamma (\alpha )_g𝐂`$. Therefore, $`H^{}(X_{(g)},L_{(g)})`$ is the subspace of $`H^{}(X_g,𝐂)`$ invariant under twisted action of $`C(g)`$
$$h\beta =\gamma (\alpha )_g(h)h^{}\beta $$
$`(5.8)`$
for any $`hC(g),\beta H^{}(X_g,𝐂)`$. However, $`H^i(X_g,𝐂)=0`$ for $`i1`$. Moreover, if $`\gamma (\alpha )_g`$ is nontrivial, $`H^0(X_g,L_{(g)})=0`$. Therefore, $`H_{orb}^{p,q}=0`$ for $`pq`$ and $`H_{orb}^{p,p}`$ is a vector space generated by conjugacy class of $`\alpha `$-regular elements $`g`$ with $`\iota _{(g)}=p`$. Therefore, we have a natural decomposition
$$H_{orb,\alpha }^{}(X,𝐂)=Z[𝐂_\alpha [G])=\underset{p}{}H_p,$$
$`(5.9)`$
where $`H_p`$ is generated by conjugacy classes of $`\alpha `$-regular elements $`g`$ with $`\iota _{(g)}=p`$. The ring structure is also easy to describe. Let $`x_{(g)}`$ be generator corresponding to zero cohomology class of twisted sector $`X_{(g)}`$ such that $`g`$ is $`\alpha `$-regular. The cup product is completely same as the nontwisted case except we replace conjugacy class by $`\alpha `$-conjugacy class. Let me sketch the calculation. As we showed in previous example, the multiplication of conjugacy classes can be described in terms of center of twisted group algebra $`Z(𝐂_\alpha [G])`$. But we have further restrictions in this case. It is clear
$$X_{(h_1,h_2,(h_1h_2)^1)}=X_{h_1}X_{h_2}/C(h_1,h_2).$$
To have nonzero invariant, we require that
$$\iota _{(h_1h_2)}=\iota _{(h_1)}+\iota _{(h_2)}.$$
$`(5.10)`$
Then, we need to compute
$$_{X_{h_1}X_{h_2}/C(h_1,h_2)}e_3^{}(vol_c(X_{h_1h_2}/C(h_1h_2)))e(E),$$
$`(5.11)`$
where $`vol_c(X_{h_1h_2}/C(h_1h_2))`$ is the compact supported top form with volume one. However,
$$X_{h_1}X_{h_2}/X_{h_1h_2}$$
is a submanifold. (5.11) is zero unless
$$X_{h_1}X_{h_2}=X_{h_1h_2}.$$
$`(5.12)`$
In this case, we call $`(h_1,h_2)`$ transverse. In this case, it is clear that obstruction bundle is trivial. Suppose that $`d_{h_1,h_2}`$ is the order of finite cover $`X_{h_1h_2}/C(h_1,h_2)X_{h_1h_2}/C(h_1h_2)`$. Then, the integral is $`d_{h_1,h_2}`$. Let
$$I_{g_1,g_2}=\{(h_1,h_2);h_i(g_i),\iota _{(h_1)}+\iota _{(h_2)}=\iota _{(h_1h_2)},(h_1,h_2)transverse,(h_1h_2)\alpha regular\}.$$
$`(5.13)`$
Then,
$$x_{(g_1)}x_{(g_2)}=\underset{(h_1,h_2)I_{g_1,g_2}}{}d_{h_1,h_2}x_{(h_1h_2)}.$$
$`(5.14)`$
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# 1 Localization of Light
## 1 Localization of Light
In some cases a medium may possess an electromagnetic band gap where a severe depressions of the photon density of states occurs. One such case is rather well known – This is the appearance of the polariton band gap due to the interaction of light with a dense medium \[1-3\]. Another possibility has been recently found when a spectral photon gap develops because of periodicity of dielectric structures . The latter are called photonic band–gap materials. In such dielectric superlattices, strong localization of photons happens \[4-6\]. The same effect of light localization arises if a resonance atom, with a frequency inside the polariton band gap, is doped into a dispersive medium \[7-9\]. If the atomic resonance frequency lies near the gap, a polariton–atom bound state appears with an eigenfrequency lying within the gap. The appearance of this bound state results in a significant suppression of spontaneous emission, that is, in localization of light. This means that an atom, in the stationary state, has a finite probability to be in the excited state, provided that it was excited at the initial moment of time.
To explain in simple parlance what does mean the localization of light, let us consider a two–level atom whose population difference is described by an operator $`\sigma ^z(t)`$. It is convenient to introduce the excitation function
$$\eta (t)\frac{1}{2}\left(1+s(t)\right),s(t)\sigma ^z(t),$$
in which the angle brackets $`\mathrm{}`$ imply a quatum–mechanical averaging. The atom is excited when the population difference $`s=1`$, i.e., the excitation function $`\eta =1`$. When the atom is not excited, then $`s=1`$ and $`\eta =0`$. There exists a finite probability to find the atom excited, in the stationary state, if and only if
$$\underset{t\mathrm{}}{lim}\eta (t)0,\underset{t\mathrm{}}{lim}s(t)=\zeta 1.$$
This is exactly what one means under the localization of light. The reason for this localization is the suppression of spontaneous emission, if the atom was initially excited. If the suppression is absolute, then, starting with $`s(0)=1`$, one comes to $`\zeta =1`$, that is, the atom does not become deexcited, keeping an absorbed photon forever. Vice versa, if an atom, with a frequency inside the polariton band gap, was not excited at the initial time $`t=0`$, that is, $`s(0)=1`$, then it should remain not excited in the stationary state, as $`t\mathrm{}`$, since there are no photons in the gap, which could excite the atom. The corresponding dynamics of the population difference can be described by the equation
$$\frac{ds}{dt}=\gamma _1(s\zeta ),\zeta s(0),$$
whose evident solution $`s(t)=\zeta s(0)`$ demonstrates the suppression of spontaneous emission.
The situation becomes more complicated when a collection of identical impurity atoms, with a transition frequency in the polariton band gap, is incorporated into the medium. If the spacing between the admixture atoms is much smaller than the transition wavelength, then the electromagnetic coupling of atoms leads to the formation of a photonic impurity band within the polarization band gap . If the density of the admixture is high, the polariton gap can be destroyed at all. Electromagnetic field can propagate in the impurity band formed by collective interactions of atoms, and coherent radiation becomes possible. If the spacing between the admixture atoms would be much larger than the transition wavelength, then the propagation band could not be formed and the atomic radiation would be prohibited. In this way, only coherent interactions can overcome the suppression of emission caused by the localization of light.
The time evolution of spontaneous emission near the edge of a photonic band gap has been considered for a simple concentrated Dicke model, where the radiation wavelength is assumed to be much larger than not only the interatomic spacing but the whole system. This model, evidently, is equivalent to a single–atom model with atomic variables factored by the number of atoms.
The aim of the present communication is to suggest a more realistic, though yet solvable, approach to describing coherent emission of admixture atoms placed in a medium with localizing light. Clearly, a more realistic approach is, at the same time, more and even much more complicated. Therefore, it would be impossible in the frame of this communication to expound it in whole. The main attention here is paid to the the formulation of the problem, with a brief survey of physical picture, and some new results are announced.
## 2 Formulation of Problem
Consider a system of $`N`$ resonance two–level radiators enumerated by the index $`i=1,2,\mathrm{},N`$. These can be atoms, molecules, nuclei, or quantum dots. For short, let us call them atoms. Their Hamiltonian is
$$\widehat{H}_a=\frac{1}{2}\underset{i=1}{\overset{N}{}}\omega _0(1+\sigma _i^z),$$
(1)
where $`\omega _0`$ is a transition frequency $`(\mathrm{}1)`$ and $`\sigma _i^z`$ is a population difference operator. The electromagnetic field Hamiltonian has the general form
$$\widehat{H}_f=\frac{1}{8\pi }[\stackrel{2}{\stackrel{}{E}}(\stackrel{}{r})+\stackrel{2}{\stackrel{}{H}}(\stackrel{}{r})]d\stackrel{}{r},$$
(2)
with the electric field $`\stackrel{}{E}`$ and magnetic field $`\stackrel{}{H}\stackrel{}{}\times \stackrel{}{A}`$, where $`\stackrel{}{A}`$ is the vector–potential satisfying the Coulomb gauge condition $`\stackrel{}{}\stackrel{}{A}=0`$. The interaction between the atoms and field, in the dipole approximation, is given by
$$\widehat{H}_{af}=\frac{1}{c}\underset{i=1}{\overset{N}{}}\underset{a}{\overset{}{J}}(\underset{i}{\overset{}{r}})\stackrel{}{A}(\underset{i}{\overset{}{r}}),$$
(3)
with the transition current
$$\underset{a}{\overset{}{J}}(\underset{i}{\overset{}{r}})=i\omega _0(\sigma _i^+\stackrel{}{\stackrel{}{d}}\sigma _i^{}\stackrel{}{d}),$$
(4)
in which $`\sigma _i^\pm `$ is a raising or lowering operator, respectively; and $`\stackrel{}{d}`$, a transition dipole. As usual , the relativistic term $`\stackrel{2}{\stackrel{}{A}}/c^2`$ is neglected. The system of radiators defined by Eqs.(1) to (4) forms the basis for the standard consideration of collective processes in the emission of light .
The case we are considering here is aggravated by the fact that the resonance atoms are not in empty space but are inserted as admixtures into a medium. The latter can be modeled in different ways, with the main requirement that a band gap should appear resulting in the localization of light for a single atom. Photonic band–gap materials can be described by periodic superstructures of scatterers . A frequency gap for propagating electromagnetic modes exists also in many natural dielectrics and semiconductors . For instance, the polariton effect is well developed in such semiconductors as $`CuCl,CuBr,CdS,CdSe,ZnSe,GaAs,GaSb,InAs,AlAs,SiC`$, and in some semiconductor microstructures including quantum dots, wells, and wires . A frequency gap for light propagation inside dense media is known to appear when the latter contain excitations being in resonance with the frequency of light . For example, the matter could be presented as an ensemble of two–level radiators. Then, these radiators together with the admixture atoms would form a kind of a resonance two–component system \[17-19\]. The polariton gap in dispersive dense media arises because of the interaction of light with some gapful elementary excitations, like excitons or optical phonons . Hence, a medium can be modeled by an ensemble of oscillators, possessing an optical branch, which represent optical–type collective excitations of the medium. A Hamiltonian of such collective excitations has, in general, the form
$$\widehat{H}_m=\underset{j=1}{\overset{N^{}}{}}\frac{\underset{j}{\overset{2}{\stackrel{}{p}}}}{2m}+\frac{1}{2}\underset{ij}{\overset{N^{}}{}}\underset{\alpha \beta }{\overset{3}{}}D_{ij}^{\alpha \beta }u_i^\alpha u_j^\beta ,$$
(5)
where $`N^{}`$ is the number of lattice sites; $`\underset{i}{\overset{}{p}}`$ and $`\underset{i}{\overset{}{u}}`$ are momentum and displacement operators, respectively; $`D_{ij}^{\alpha \beta }`$ is a dynamical matrix. The interaction of the matter excitations with electromagnetic field is described by the term
$$\widehat{H}_{mf}=\frac{1}{c}\underset{j=1}{\overset{N^{}}{}}\underset{m}{\overset{}{J}}(\underset{j}{\overset{}{r}})\stackrel{}{A}(\underset{j}{\overset{}{r}}),$$
(6)
in which
$$\underset{m}{\overset{}{J}}(\underset{j}{\overset{}{r}})=\frac{e}{m}\underset{j}{\overset{}{p}}$$
(7)
is a local current at the point $`\underset{j}{\overset{}{r}}`$ of the medium.
In this way, the total Hamiltonian of the considered system is
$$\widehat{H}=\widehat{H}_a+\widehat{H}_f+\widehat{H}_{af}+\widehat{H}_m+\widehat{H}_{mf},$$
(8)
consisting of the atom Hamiltonian (1), field term (2), atom–field interaction (3), medium Hamiltonian (5), and of the medium–field interaction (6).
The stationary states of the system with Hamiltonian (8) can be studied as follows. One expands the vector potential
$$\stackrel{}{A}=\underset{k\nu }{}\left(\frac{2\pi c}{kV}\right)^{1/2}(a_{k\nu }\underset{k\nu }{\overset{}{e}}e^{i\stackrel{}{k}\stackrel{}{r}}+a_{k\nu }^{}\underset{k\nu }{\overset{}{\stackrel{}{e}}}e^{i\stackrel{}{k}\stackrel{}{r}})$$
(9)
in plane waves, with $`\underset{k\nu }{\overset{}{e}}`$ being a polarization vector; $`k|\stackrel{}{k}|;\stackrel{}{k}`$, a wave vector; $`\nu =1,2`$; $`V`$, volume; and $`a_{k\nu }`$ being a photon operator indexed by the wave vector $`\stackrel{}{k}`$ and the polarization–branch number $`\nu `$. The displacement and momentum operators of the medium are also expanded in plane waves:
$$\underset{j}{\overset{}{u}}=\underset{ks}{}\left(2mN^{}\omega _{ks}\right)^{1/2}(b_{ks}+b_{ks}^{})\underset{ks}{\overset{}{e}}e^{i\stackrel{}{k}\underset{j}{\overset{}{r}}},$$
$$\underset{j}{\overset{}{p}}=i\underset{ks}{}\left(\frac{m\omega _{ks}}{2N^{}}\right)^{1/2}(b_{ks}b_{ks}^{})\underset{ks}{\overset{}{e}}e^{i\stackrel{}{k}\underset{j}{\overset{}{r}}},$$
(10)
where $`\underset{ks}{\overset{}{e}}`$ is a corresponding polarization vector; $`s=1,2,3`$; and $`\omega _{ks}`$ is the spectrum of collective excitations defined by the eigenvalue problem
$$\frac{1}{mN^{}}\underset{ij}{\overset{N^{}}{}}\underset{\beta =1}{\overset{3}{}}D_{ij}^{\alpha \beta }e^{i\stackrel{}{k}\underset{ij}{\overset{}{r}}}e_{ks}^\beta =\omega _{ks}^2e_{ks}^\alpha ,$$
(11)
with $`\underset{ij}{\overset{}{r}}=\underset{i}{\overset{}{r}}\underset{j}{\overset{}{r}}`$. The frequency and polarization vectors are assumed to be even functions of the wave vector, $`\omega _{ks}=\omega _{ks}`$, $`\underset{ks}{\overset{}{e}}=\underset{ks}{\overset{}{e}}`$. The destruction, $`b_{ks}`$, and creation, $`b_{ks}^{}`$, operators of collective oscillations satisfy the Bose commutation relations. With expansion (9), the field Hamiltonian (2) takes the known simple form
$$\widehat{H}_f=\underset{k\nu }{}ck\left(a_{k\nu }^{}a_{k\nu }+\frac{1}{2}\right).$$
(12)
And the matter Hamiltonian (5), by means of (10), becomes
$$\widehat{H}_m=\underset{ks}{}\omega _{ks}\left(b_{ks}^{}b_{ks}+\frac{1}{2}\right).$$
(13)
Recall that an optical–type spectrum $`\omega _{ks}`$ is to be assumed in (13). Introducing polariton operators that are linear combinations of the photon operators $`a_{k\nu }`$ and of the boson operators $`b_{k\nu }`$, one can, in some cases , diagonalize the sum of the Hamiltonians $`\widehat{H}_f+\widehat{H}_m+\widehat{H}_{mf}`$, obtaining a diagonal polariton Hamiltonian. For example, a detailed description of this procedure of diagonalization can be found in Ref., where a uniform and isotropic model is considered and several simplifications, in line with the Heitler–London and resonance approximation, are involved. In these approximations, one neglects the counter–rotating terms and two–boson transitions, whose influence, similarly to two–atom transitions , can become important only far from the resonance.
When there is only one admixture atom in the medium, that is $`N=1`$, then the stationary states of Hamiltonian (8) can be found resorting to the uniform, isotropic, and resonance approximations. If the atomic resonance frequency lies near the gap, there appears a polariton–atom bound state with an eigenfrequency lying within the gap. This means that light is localized at the atom. As a result of the appearance of this localized bound state, a significant suppression of spontaneous emission occurs . The behavior of many–polariton states is a little more diverse. Those states containing an even number of polaritons correspond to solitons that can propagate within the gap, while the states with an odd number of polaritons represent solitons that are pinned to the atom forming a many–polariton bound state. The latter state, similarly to the single–polariton bound state, also depicts the localization of light .
When two identical two–level atoms are placed in a frequency dispersive medium whose polariton spectrum has a gap, then the polariton–atom bound state lying within the polaritonic gap splits into a doublet due to an effective atom–atom interaction . The spontaneous emission can exist only if the resonance frequency of these two atoms lies in the polariton continuous spectrum. And if the resonance frequency lies within the gap, then two discrete modes represent a doublet of bound polariton–atom states, for which spontaneous emission is practically completely suppressed .
A qualitatively different situation develops when many resonance admixture atoms are placed in the medium. Stationary states for this case have been studied for a one–dimensional atomic chain incorporated in a uniform and isotropic system . The nearest–neighbor approximation has been used. In the case of spatially correlated atoms, with a frequency inside the gap, a polariton–impurity band is formed within the polaritonic gap. Then polaritons can propagate in this impurity band and the atomic chain provides a waveguide for the radiation field. Even in the nearest–neighbor approximation, when, for a chain, the interaction of only three atoms is effectively taken into account, the width of the impurity band, normalized with respect to the polariton gap, is $`0.14`$ . Since this width is, roughly speaking, proportional to $`N1`$, the collective interaction of about $`10`$ atoms should fill the whole polariton gap.
The energy–momentum representation employed for studying stationary states is not convenient for considering space–time dynamics of collective processes. For the latter purpose, to our mind, it is more appropriate to remain in real space and time. One may write the Heisenberg evolution equations for the operators of the problem. Then, formally solving the Maxwell equations, one may exclude the field variables (see details in ). After that, one comes to the equations for the atomic variables,
$$\frac{d\sigma _i^{}}{dt}=(i\omega _0+\gamma _2)\sigma _i^{}+\sigma _i^z\stackrel{}{\stackrel{}{d}}\underset{i}{\overset{}{D}}+$$
$$+ik_0^2\sigma _i^z\stackrel{}{d}\underset{j(i)}{\overset{N}{}}\frac{1}{r_{ij}}[\sigma _j^+(t\frac{r_{ij}}{c})\stackrel{}{\stackrel{}{d}}\sigma _j^{}(t\frac{r_{ij}}{c})\stackrel{}{d}],$$
(14)
$$\frac{d\sigma _i^z}{dt}=\gamma _1(\sigma _i^z\zeta )2(\sigma _i^+\stackrel{}{\stackrel{}{d}}+\sigma _i^{}\stackrel{}{d})\underset{i}{\overset{}{D}}$$
$$2ik_0^2(\sigma _i^+\stackrel{}{\stackrel{}{d}}+\sigma _i^{}\stackrel{}{d})\underset{j(i)}{\overset{N}{}}\frac{1}{r_{ij}}[\sigma _j^+(t\frac{r_{ij}}{c})\stackrel{}{\stackrel{}{d}}\sigma _j^{}(t\frac{r_{ij}}{c})\stackrel{}{d}],$$
(15)
where $`\gamma _1`$ and $`\gamma _2`$ are the longitudinal and transverse relaxation parameters, respectively; $`k_0\omega _0/c`$; $`r_{ij}|\underset{ij}{\overset{}{r}}|`$; and
$$\underset{i}{\overset{}{D}}(t)=\frac{k_0}{c}\underset{j(i)}{\overset{N^{}}{}}\frac{1}{r_{ij}}\underset{m}{\overset{}{J}}(\underset{j}{\overset{}{r}},t\frac{r_{ij}}{c})$$
is an electric field produced by the medium. Assuming that the single–atom stationary state corresponds to localized light, we put $`\zeta =\sigma _i^z(0)`$, where $`\mathrm{}`$ means the statistical averaging over an initial state. Equations (14) and (15) are complemented by initial conditions
$$u_0=\sigma _i^{}(0),s_0=\sigma _i^z(0).$$
## 3 New Results
The notion of light localization is relatively recent. This is why the main part of this communication has been devoted to the description of the related physical picture and to the formulation of the basic equations that would permit one to consider the time development of collective phenomena for a system of admixture atoms in a medium with localized light. The limited frames of this communication do not allow to expound in detail the way of solving the basic equations (14) and (15) and the analysis of the corresponding solutions. This will be done in a separate publication. Here, we only can briefly delineate the scheme of solving Eqs.(14) and (15) and present some fresh results.
From the operator equations (14) and (15), we pass to the equations for the related statistical averages. The pair correlation functions are decoupled in the semiclassical approximation. The retardation is treated in the quasirelativistic approximation as in Refs. \[26-28\]. The system of nonlinear differential equations is solved by means of the scale separation approach . The resulting physical picture depends on the values of the following parameters: initial conditions $`u_0`$ and $`s_0`$; the coupling parameter of coherent atomic interactions,
$$g=\frac{k_0^3d_0^2}{\gamma _2}\underset{j(i)}{\overset{N}{}}\frac{\mathrm{sin}(k_0r_{ij})}{k_0r_{ij}};$$
(16)
the effective parameter of coupling between the admixture atoms and matter,
$$\alpha =|e^{\mathrm{\Gamma }t}_0^te^{(i\mathrm{\Omega }+\mathrm{\Gamma })\tau }\stackrel{}{\stackrel{}{d}}\stackrel{}{D}(\tau )d\tau |^2,$$
(17)
where the double brackets $`\mathrm{}`$ imply the statistical and time averaging, and
$$\mathrm{\Omega }=\omega _0+\mathrm{\Delta }_Ls,\mathrm{\Gamma }=\gamma _2(1gs)$$
(18)
are an effective frequency and attenuation, with
$$\mathrm{\Delta }_L=k_0^3d_0^2\underset{j(i)}{\overset{N}{}}\frac{\mathrm{cos}(k_0r_{ij})}{k_0r_{ij}},s=\sigma _i^z;$$
and the critical atom–matter coupling parameter
$$\alpha _c=\frac{(1gs_0)^2+4g^2|u_0|^2}{4g^2s_0^2}.$$
(19)
The overall physical picture for the case $`g1`$ and $`\alpha \alpha _c`$ is as follows. If the admixture atoms at the initial time $`t=0`$ are excited, then after a delay time $`t_0T_1`$, defined by the values of the above parameters, a coherent burst occurs. Then, a series of coherent bursts follows, separated from each other by the periods of practically no radiations. This series of bursts lasts for the time of several $`T_1`$. Finally, dynamics tends to a stationary state with the excitation function
$$\underset{t\mathrm{}}{lim}\eta (t)=\frac{1}{2}\left(1+\frac{1}{g}\right).$$
(20)
Remembering the definition in Section 1, we see that the limit (20) exhibits a partial localization of light.
Acknowledgements
I am grateful for discussions to M.R. Singh and W. Lau. A Senior Fellowship from the University of Western Ontario, Canada, is appreciated.
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# The Transition in the Two-Dimensional Step Model: A Kosterlitz-Thouless Transition in Disguise
\[
## Abstract
Evidence for a Kosterlitz-Thouless transition in the 2D step model is obtained from Monte Carlo determinations of the helicity modulus. It is argued that the free energy of a single vortex at the center of the system depends logarithmically on the system size in spite of the fact that the spin interaction is not harmonic for small differences in the spin angles. We conclude that this is the reason for the Kosterlitz-Thouless transition in the 2D step model and that the harmonic spin interaction not is a necessary requirement.
\]
The phase transition in two-dimensional (2D) XY models is known to take place through the vortex unbinding mechanism due to Kosterlitz and Thouless (KT). From the principles of universality one expects this transition to remain the same independent of details of the system as e.g. the underlying lattice structure. The precise spin interaction potential $`U(\varphi )`$, where $`\varphi `$ is the angle difference between neighboring spins, is not supposed to be essential either. $`U(\varphi )`$ is however required to be periodic in $`2\pi `$ and it seems always to have been presumed that the interaction in addition has to be harmonic for small $`\varphi `$. The harmonicity for small $`\varphi `$ has to do with the energetics for vortex formation. With a harmonic potential the energy for a single vortex in a $`L\times L`$ lattice goes as $`\mathrm{ln}L`$, and in the classical argument by Kosterlitz and Thouless this property is crucial for the transition.
The subject of the present Letter is the 2D step model which is an XY model with a spin interaction that has no harmonic component. $`U(\varphi )`$ is instead a step-like function
$$U(\varphi )=J\mathrm{sign}(\mathrm{cos}\varphi )$$
(1)
Since this potential is flat around $`\varphi =0`$ there is no $`\mathrm{ln}L`$-dependence of the energy for a single vortex. This energy is instead independent of system size and one would therefore expect a non-vanishing density of free vortices at all finite temperatures, and consequently no phase transition. Against that background the evidence from simulations for a transition were very intriguing. The first clear evidence for a transition was obtained from a Monte Carlo (MC) study of the susceptibility and the specific heat. The increase of the susceptibility with lattice size was considered to suggest a phase transition at $`T1.1`$. Later simulations also provided evidence that the transition actually is in the same universality class as the harmonic XY models. The similarity of the behavior close to the transition of an harmonic XY model and the step model also led these authors to question the vortex unbinding as the mechanism behind the KT transition. How could a transition driven by vortices be altogether insensitive to the very different energy cost for vortices in the two models?
In the present Letter we address the question of the necessity of an harmonic spin potential for the KT transition by examining the behavior of the 2D step model. We first calculate the helicity modulus and show that the behavior of this quantity gives strong support for a KT transition. We then demonstrate that the cost in *free energy* for a single vortex at the center of the system in fact goes as $`\mathrm{ln}L`$. It is this feature that stabilizes the low temperature phase. Finally, we refine the arguments to obtain quantitatively satisfactory estimates.
The helicity modulus $`\mathrm{{\rm Y}}`$, is a convenient quantity for the study of KT transitions due to its universal value $`2T/\pi `$ at the transition, and the known form of the approach to this universal value with $`L`$. The usual procedure in MC simulations is to determine $`\mathrm{{\rm Y}}`$ from a correlation function which involves some derivatives of $`U(\varphi )`$. Clearly, with a step-like potential the derivatives of the potential cannot be calculated and this expression cannot be used. A way out is to instead start from the defining expression for the helicity modulus
$$\mathrm{{\rm Y}}=\frac{^2F}{\mathrm{\Delta }^2}|_{\mathrm{\Delta }=0},$$
(2)
and perform the simulations with fluctuating twist boundary conditions. In these simulations one collects a histogram of the total twist $`P(\mathrm{\Delta })`$. Since the probability for a certain twist is related to the free energy through $`P(\mathrm{\Delta })e^{F(\mathrm{\Delta })/T}`$ Eq. (2) becomes
$$\mathrm{{\rm Y}}=T\frac{^2\mathrm{ln}P}{\mathrm{\Delta }^2}|_{\mathrm{\Delta }=0}.$$
(3)
The simulations are done with twist variables in the two directions, $`\mathrm{\Delta }_x`$ and $`\mathrm{\Delta }_y`$, which beside the spin variables $`\theta _i`$ are updated with the Metropolis algorithm. With $`𝐫_{ij}`$ a unit vector between nearest neighbors and $`𝚫=(\mathrm{\Delta }_x,\mathrm{\Delta }_y)`$, the Hamiltonian may be written
$$H=\underset{ij}{}U\left(\theta _i\theta _j\frac{1}{L}𝐫_{ij}𝚫\right)=\underset{ij}{}U(\varphi _{ij}).$$
To get a good acceptance ratio for the twist update it is necessary to make use of $`L`$ different twist variables in each direction, with $`\mathrm{\Delta }_x=_{k=1}^L\mathrm{\Delta }_x^{(k)}`$, (and similarly in the $`y`$ direction) where $`k`$ is the column (row) number. In our simulations, which for convenience were for a $`O(256)`$ model, we used the potential $`U(\varphi )=J`$ for 129 angle differences $`[\pi /2`$, $`\pi /2]`$ and $`+J`$ for the remaining 127. This choice is not expected to be important for the transition properties, but gives a slight shift of the transition temperature as compared to the potential of Eq. (1). The length of the runs were typically $`5\times 10^8/L`$ sweeps through the lattice.
In Fig. 1 we show the histogram $`P(\mathrm{\Delta })`$ from Monte Carlo simulations at $`T=0.05J`$. Since the histogram is peaked around $`\mathrm{\Delta }=0`$ the figure immediately gives evidence for a low temperature phase with a finite stiffness. To determine $`\mathrm{{\rm Y}}`$ we fit a quadratic curve to $`\mathrm{ln}P`$ for $`|\mathrm{\Delta }/\pi |<1/3`$ and obtain $`\mathrm{{\rm Y}}/T=0.789(4)`$ (where the given error is one standard deviation). Note that this is slightly larger than the universal value $`2/\pi 0.637`$, which is required for a stable low temperature phase.
An important feature of the step model is the gap in excitation energies; there are no excitations with energy $`<2J`$. At $`TJ`$ the system is therefore at all times in one of its numerous ground states which means that the histogram $`P(\mathrm{\Delta })`$ is independent of temperature. From Eq. (3) then follows a linear temperature dependence for $`\mathrm{{\rm Y}}`$. This is in contrast to harmonic XY models for which $`\mathrm{{\rm Y}}/J`$ approaches unity in the low-temperature limit.
The temperature dependence of $`\mathrm{{\rm Y}}`$ is shown in Fig. 2 for several system sizes together with the dashed line for the universal jump condition $`2T/\pi `$. We note that the curves start out linearly at low temperatures, become size-dependent at $`T/J0.75`$, cross the universal line and then drop down to zero. Beside the unusual linear temperature-dependence at low $`T`$ this behavior is just as in an ordinary harmonic XY model and therefore precisely what one would expect for a KT transition.
To determine the KT temperature we make use of the finite size dependence of $`\mathrm{{\rm Y}}`$. We follow the procedure in Ref. of first fitting our MC data for $`\mathrm{{\rm Y}}`$ from a narrow temperature interval to second order polynomials in $`T`$, one for each $`L`$, and then fit the data to the expression
$$\mathrm{{\rm Y}}_L(T_{\mathrm{KT}})=\frac{2T_{\mathrm{KT}}}{\pi }\left(1+\frac{1}{A+2\mathrm{ln}L}\right).$$
(4)
which amounts to adjusting $`T_{\mathrm{KT}}`$ and $`A`$ to get the best possible fit. Using data for $`L16`$ we obtained $`T_{\mathrm{KT}}=0.765(6)`$. Fig. 3 illustrates the good fit of the data at the transition temperature to the line from Eq. (4). We consider the above MC data to be strong evidence for a KT transition. This is in agreement with the conclusion in Ref. that the step model is in the same universality class as the harmonic XY model.
We now propose an analytical analysis in order to understand this behavior. We focus on the properties in the low temperature phase where the angular differences are restricted to the low energy region, $`|\varphi |\pi /2`$. A central idea in the present letter is to note that, while in the harmonic XY model the spinwave-vortex interaction only is a smooth perturbation, in the step model this interaction leads to a dramatic and crucial effect: while the energy of a system with a single vortex fixed in the center of the system is finite, the *free energy* of this system grows as $`\mathrm{ln}L`$. This is due to the change in entropy of spinwave fluctuations for the configuration with the fixed vortex, as compared to the vortex free case. Note that this is the entropy associated with a *fixed* vortex, not the positional entropy associated with a free vortex’s variable location.
To demonstrate the existence of this spinwave entropy we consider the configuration of spins in Fig. 4a where we slightly reorganize the spins and delete the links in the radial direction. We will return below to the approximation involved in this step.
The condition for having a positive vortex in the center of the system in Fig. 4a is that the phase rotates by $`2\pi `$ along each of the circles. We therefore introduce $`\mathrm{\Phi }_r=\varphi `$ along the circle with radius $`r`$, which is a sum of $`2\pi r`$ values and let $`\mathrm{\Omega }_r(\mathrm{\Phi }_r)`$ denote the number of possible combinations of the $`\varphi `$:s at distance $`r`$ as a function of $`\mathrm{\Phi }_r`$, which is defined only for $`\mathrm{\Phi }_r=2\pi n`$ (with integer $`n`$). The probability for having a (positive) vortex is then determined by the product
$$P_{\mathrm{v}ort}\underset{r=1}{\overset{L}{}}\frac{\mathrm{\Omega }_r(2\pi )}{\mathrm{\Omega }_r(0)}$$
(5)
The fraction $`\mathrm{\Omega }_r(2\pi )/\mathrm{\Omega }_r(0)`$ may be calculated if we temporarily open up a closed path that makes up a circle. Since we want the same number of links along this path we need one more spin variable at one of the endpoints, which we take to be independent of the other endpoint. $`\mathrm{\Phi }_r`$ then becomes a sum of $`2\pi r`$ independent variables $`\varphi `$. With the average of $`\varphi `$ being equal to zero and its variance given by $`\sigma ^2`$ the distribution of $`\mathrm{\Phi }_r`$ for this open path becomes a Gaussian with width $`2\pi r\sigma ^2`$:
$$\mathrm{\Omega }_r^{\mathrm{o}pen}(\mathrm{\Phi }_r)\mathrm{exp}\left(\frac{\mathrm{\Phi }_r^2}{2\pi r\times 2\sigma ^2}\right)$$
(6)
We now make use of the fact that the number of possible configurations for the open and closed paths are the same if the spins at the endpoints of the open path are equal. Since this condition is equal to having $`\mathrm{\Phi }_r=2\pi n`$ we conclude that $`\mathrm{\Omega }_r(2\pi n)=\mathrm{\Omega }_r^{\mathrm{o}pen}(2\pi n)`$. From Eqs. (5) and (6) the probability for a vortex then becomes
$$P_{\mathrm{v}ort}=\mathrm{exp}\left(\frac{\pi }{\sigma ^2}\underset{r=1}{\overset{L}{}}\frac{1}{r}\right)\mathrm{exp}\left(\frac{\pi }{\sigma ^2}\mathrm{ln}L\right),$$
for the entropy of a fixed vortex we obtain $`S_{\mathrm{v}ort}=\mathrm{ln}P_{\mathrm{v}ort}=\frac{\pi }{\sigma ^2}\mathrm{ln}L`$, and the free energy for a vortex at a fixed position in a system of size $`L`$ finally becomes
$$F_{\mathrm{v}ort}=TS_{\mathrm{v}ort}=T\frac{\pi }{\sigma ^2}\mathrm{ln}L.$$
(7)
For the harmonic XY model the well-known argument for the phase transition gives the free energy for having a vortex at any of the $`L^2`$ positions as
$$\mathrm{\Delta }F=\left(\pi J2T\right)\mathrm{ln}L,$$
and the transition takes place at $`T/J=2/\pi `$. In the step model the corresponding expression becomes
$$\mathrm{\Delta }F=F_{\mathrm{v}ort}2T\mathrm{ln}L=T\left(\frac{\pi }{\sigma ^2}2\right)\mathrm{ln}L,$$
where the temperature, at first sight only appears as a prefactor. However, there is a hidden temperature dependence in the variance $`\sigma ^2`$. At low enough temperatures the $`\varphi _{ij}`$ are restricted to the interval $`[\pi /2,\pi /2]`$ but with increasing temperature the $`\varphi _{ij}`$ will more often take values outside this interval, and $`\sigma ^2`$ will increase. Therefore, if $`\mathrm{\Delta }F`$ is positive at low temperatures it will turn negative at some finite temperature and this will give the transition. But if $`\sigma ^2>\pi /2`$ already at low temperatures there will be no transition.
To see how the local restrictions on the angle differences give rise to the non-zero helicity modulus we now turn to a rectangular geometry and consider the determination of $`\mathrm{{\rm Y}}`$ from the distribution of $`\mathrm{\Delta }_x`$ and $`\mathrm{\Delta }_y`$. One point with examining the distribution of the twist is to give predictions that are easy to compare with MC simulations. In the simplest approximation we again delete all links in the perpendicular direction, as in Fig. 4b. For a single row the number of configurations consistent with a certain total twist becomes $`\mathrm{\Omega }_{\mathrm{r}ow}(\mathrm{\Delta })\mathrm{exp}\left(\mathrm{\Delta }^2/2L\sigma ^2\right)`$, in analogy with Eq. (6), and the number of configurations for the whole system with $`L`$ rows becomes
$$\mathrm{\Omega }(\mathrm{\Delta })=\left[\mathrm{\Omega }_{\mathrm{r}ow}(\mathrm{\Delta })\right]^L\mathrm{exp}\left(\mathrm{\Delta }^2/2\sigma ^2\right).$$
For the free energy we then arrive at $`F(\mathrm{\Delta })=T\mathrm{\Delta }^2/2\sigma ^2`$ and with Eq. (2) the helicity modulus becomes
$$\mathrm{{\rm Y}}=T/\sigma ^2.$$
(8)
In the absence of perpendicular links and at low temperatures, the $`\varphi _{ij}`$ have a rectangular distribution, and from elementary integrals one finds $`\sigma ^2=\pi ^2/120.822`$. Through Eq. (8) this gives $`\mathrm{{\rm Y}}/T1.22`$ which is about 50% larger than $`\mathrm{{\rm Y}}/T0.789`$ from Monte Carlo simulations, cf. Fig. 1. A better estimate will be obtained below by including some links in the perpendicular direction.
Comparing Eqs. (7) and (8) we find $`F_{\mathrm{v}ort}(L)=\mathrm{{\rm Y}}\pi \mathrm{ln}L`$ which is the same relation as in the harmonic model. This shows that our two different calculations are equivalent which is a consequence of using the same approximation in both cases, i.e. neglecting all links perpendicular to the direction of interest.
We now discuss the assumption used above, namely that the qualitatively correct behavior in a certain direction may be obtained even though one neglects the perpendicular links. This assumption is true only if the relative reduction of the number of allowed configurations obtained by introducing perpendicular links is essentially independent of the total twist in the direction of interest. We argue that this is a plausible assumption by considering two sets of configurations: 1) the set of all twist free configurations and 2) the set of configurations with a twist $`\mathrm{\Delta }_x`$. There is then a transformation $`\varphi _{ij}+\mathrm{\Delta }/L\varphi _{ij}`$ on all the horizontal links that transforms each member in the twist free set into a corresponding one in the twisted set. Since this transformation not affects the angle difference at the perpendicular links, it follows that the effect of the perpendicular links will be to exclude the same number of configurations in both these sets and this suggests that the relative reduction due to the perpendicular links will be independent of $`\mathrm{\Delta }`$. However, this argument only serves to make our assumption a reasonable one; it is not conclusive. There is nothing that guarantees that the relative reduction of the *allowed* configurations (with $`|\varphi _{ij}|<\pi /2`$ for all horizontal links) will be the same for the two different sets. That a certain member of the twist free set is allowed doesn’t imply that the corresponding member in the twisted set is allowed too.
Even though we have argued that the perpendicular links may be neglected in a qualitative discussion, they have to be included in order to get reasonable quantitative estimates, since they do affect the variance $`\sigma ^2`$. To include some perpendicular links we consider the configuration in Fig. 5a, where the perpendicular links have been deleted at every third row only. The approach is then to integrate out the upper and lower rows of spins in each of these triple rows to give $`L/3`$ one-dimensional rows, cf. Fig. 5b. These integrations may be done analytically or by Monte Carlo simulations on that geometry. The integrations also give correlations to the neighboring and next neighboring links, $`\varphi \varphi ^{}`$ and $`\varphi \varphi ^{\prime \prime }`$ (cf. Fig. 6), which together give an effective variance $`\sigma _{\mathrm{e}ff}^2`$:
$$\sigma _{\mathrm{e}ff}^2=\frac{1}{L}\left(\underset{i=1}{\overset{L}{}}\varphi _i\right)^2\sigma ^2+\varphi \varphi ^{}+\varphi \varphi ^{\prime \prime }.$$
(9)
With $`L/3`$ rows, as in Fig. 5b, the expression for the helicity modulus becomes $`\mathrm{{\rm Y}}=T/(3\sigma _{\mathrm{e}ff}^2)`$. From our MC simulations on the geometry of Fig. 5a we get $`\sigma ^2=0.672`$, $`\varphi \varphi ^{}=0.120`$, and $`\varphi \varphi ^{\prime \prime }=0.026`$. The first two of these numbers are easily obtained by integrating analytically with symbolic software. With Eq. (9) this gives $`\sigma _{\mathrm{e}ff}^2=0.526`$ and $`\mathrm{{\rm Y}}/T=0.633`$ which is less than 20% off $`\mathrm{{\rm Y}}/T=0.789`$ obtained above.
To conclude we have presented evidence from simulations that the 2D step model actually undergoes a KT transition. We have argued that the reason for the stability of the low temperature phase against the formation of free vortices is the $`\mathrm{ln}L`$-dependence of the free energy for a vortex at a fixed position. From these results we are led to the conclusion that the harmonic spin interaction is not a necessary condition for a KT transition in a 2D spin model, and that the KT transition is a more general phenomenon than has so far been recognized.
The authors thank Prof. P. Minnhagen and Prof. S. Teitel for critical reading of the manuscript. Financial support from the Swedish Natural Science Research Council through Contract No. E-EG 10376-312 is gratefully acknowledged.
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# First order transition from ferro- to antiferromagnetism in CeFe2 based pseudobinary alloys
## Abstract
We present results of ac susceptibility measurements highlighting the presence of thermal hysteresis and phase coexistence across the ferro-to antiferromagnetic transition in various CeFe<sub>2</sub> based pseudobinary systems. These results indicate that the ferro-to antiferromagnetic transition in these systems is first order in nature.
The C15-Laves phase compound CeFe<sub>2</sub> retains its identity amongst the other members of RFe<sub>2</sub> family (where R=Y,Zr and heavy rare earth elements). First, magnetic moment of CeFe<sub>2</sub> per formula unit ($`2.4\mu _B`$) is distinctly smaller than that found in other RFe<sub>2</sub> compounds. Second,its Curie temperature T<sub>C</sub>($``$ 235 K) is relatively small in comparison to the T<sub>C</sub> of other RFe<sub>2</sub> compounds .However, short range magnetic order is detected in its paramagnetic state even in the temperature regime upto four times T<sub>C</sub> .All these aspects drew the attention of the experimentalists during last thirty years, and amongst other things the role of Ce in the magnetic properties of CeFe<sub>2</sub> has been a subject of both theoretical and experimental investigations ;this in turn led to the discovery of newer interesting properties . The most recent neutron measurements on single crystal sample of pure CeFe<sub>2</sub> have now revealed the presence of low temperature antiferromagnetic fluctuation in this otherwise ferromagnetic compound . From the study of doped-CeFe<sub>2</sub> it is already known for quite sometime that the ferromagnetism of CeFe<sub>2</sub> is quite fragile in nature and a low temperature antiferromagnetic state can be established easily with small amount of doping with elements like Al,Co,Ru,Ir,Re,Os . It should be noted, however, that the destabilization of ferromagnetism in CeFe<sub>2</sub> is not a simple disorder induced one, since the doping with other elements like Ni,Mn,Rh,Pd leads to simple dilution of ferromagnetism .
Most of the early experimental activities in CeFe<sub>2</sub> were focussed to establish the exact nature of the low temperature magnetic phase, whether it is a re-entrant spin-glass or an antiferromagnetic state and except in few cases not much emphasis was given on the exact nature of this phase transtion. With the antiferromagnetic nature of the low temperature state being more or less established , in the present work we shall specifically address the question –what is the nature of this ferro- to antiferromagnetic transition? While there exists no complete theory (to our knowledge) to explain the interesting magnetic properties of CeFe<sub>2</sub>, a phenomenological model dealing with itinerant electron systems has often been invoked to explain the para-to ferro- to antiferromagnetic transition in the doped-psuedobinary alloys of CeFe<sub>2</sub>.This phenomenological model of Moriya and Usami predictd that the ferro- to antiferromagnetic transition would be a first order transition, while para- to ferromagnetic transition would be a second order transition . With our high resolution ac-susceptibility measurement across this ferro- to antiferromagnetic transition in two doped samples of CeFe<sub>2</sub>, we shall report characteristics which are typically associated with a first order transition.On the other hand the higher temperature para- to ferromagnetic transition can be characterized as a standard second order phase transtion. We believe that such a clear cut characterization of the various phase transitions in CeFe<sub>2</sub> based pseudobinaries is necessary, either for an appropraite extension of Moriya-Usami model or for the development of newer theory for the proper understanding of the magnetic properties of CeFe<sub>2</sub>.
Two samples–Ce(Fe,5%Ir)<sub>2</sub> and Ce(Fe,7%Ru)<sub>2</sub>–used in the present study were prepared by argon arc melting from metals of at least 99.99% purity.Details of sample preparation, heat treatment and characterisation can be found in Ref. 10.The same samples have earlier been used in some other studies .
The AC susceptibility setup consists of a coil system having a primary solenoid and two oppositely wound secondaries each consisting of 1500 turns. The coil is dipped in liquid nitrogen to ensure that the temperature of the coil remains constant during the entire experiment to avoid drifts in the value of the applied field. The sample is mounted in a double walled quartz insert and its temperature is raised by heating the exchange gas by a heater wound on a seperate teflon mounting.A temperature controller (Lake Shore– DRC-91CA) is used for controlling the temperature.A copper-constantan thermocouple is used in differential mode to monitor the small temperature lag between the sample and the sensor.The sinusoidal output of a lock in amplifier (Stanford Research–SR830) is given to a voltage to current convertor which drives the current through the coil to generate the neccessary ac magnetic field. The signal from the pickup coil which is proportional to the susceptibility is measured by the same lock in amplifier. The field and frequency values were 4 Oe rms and 333Hz respectively.
Fig.(1) shows the AC susceptibility ($`\chi `$) for both Ce(Fe,5% Ir)<sub>2</sub> and Ce(Fe,7% Ru)<sub>2</sub> as a function of temperature (T). The para- to ferromagnetic transition is characterized by a sharp increase in susceptibility ($`\chi `$) with the decrease in T at T$`{}_{Curie}{}^{}`$185K in the 5% Ir doped sample and T$`{}_{Curie}{}^{}`$165K in the 7% Ru doped sample. Below T<sub>Curie</sub> the susceptibility more or less flattens out for both the samples, before decreasing sharply at around 135K in Ce(Fe,5% Ir)<sub>2</sub> and at around 125K in Ce(Fe,7% Ru)<sub>2</sub>. This low temperature decrease in $`\chi `$ has earlier been taken as a signature of ferro- to antiferromagnetic transition , and the estimated transition temperatures (T<sub>N</sub>) from our present study agree well with the existing literature .
Our aim now is to find out the exact nature of these two magnetic transitions oberved in CeFe<sub>2</sub>-based pseudobinaries. Experimentally, the indication of a first order transition usually comes via a hysteretic behaviour of various properties, not necessarily thermodynamic ones. As an example, the first indication of a first order melting transition from elastic solid to vortex liquid in a vortex matter came via distinct hysteresis observed in transport property measurements . The confirmatory tests of the first order nature of a transition ofcourse involve the detection of discontinuous change in thermodynamic observables and the estimation of latent heat, and this has subsequently been achieved for vortex lattice melting in vortex matter . There also exists a less rigorous class of experimental tests which invloves the study of phase inhomogeneity and phase coexistence across a first order transtion. This kind of experiment has also come out to be pretty informative for the melting transtion as well as ordered solid to disordered solid transition in vortex matter. In our present study we shall use hysteresis and phase coexistence to investigate the nature of the magnetic transitions in CeFe<sub>2</sub> based systems; our observable will be ac susceptibility ($`\chi `$).
In order to observe a hysteresis in the transition, if any, we have chosen to sweep the temperature at a slow rate (0.006K/sec typical and slower when needed) instead of stabilizing at each temperature. This was done to ensure that the temperature is varied unidirectionally during both the heating and cooling cycles. The signal was measured at a temperature interval of 0.2K. The time constant of the low pass filter of the lock in amplifier was chosen such that the temperature changes negligibly(compared to our temperature step) within a time interval of 10 times the time constant. The temperature difference between the sensor and the sample, as monitored by the differential thermocouple, was always less than 1% of the sensor temperature and is used to obtain the correct value of the sample temperature.
First, we show the effect of temperature cycling on the para- to ferromagnetic transition in fig(2).In case of Ce(Fe,5% Ir)<sub>2</sub> the transition is reversible within an error of 0.15K to 0.2K.In case of Ce(Fe,7% Ru)<sub>2</sub> the reversibility is even better. The lack of hysteresis in para- to ferromagnetic transition within an error bar smaller than our temperature step, is indicative of a second order phase transition.
We then focus our attention on the ferro- to antiferromagnetic transition which has been shown to be associated with a structural distortion from cubic to rhombohedral , hinting towards a first order transition.The same protocol of sweeping the temperature and measuring the signal at closely spaced temperature values is followed during this measurement also.
Fig(3) shows the result of our measurements on both 5% Ir and 7% Ru doped CeFe<sub>2</sub> samples. Both the samples show a distinct thermal hysteresis in the ac-susceptibility across the ferro- to antiferromagnetic transition.The width of the hysteresis is about 2K which is well beyond the error in our measurements.
To study the phase coexistence we use the technique of minor hysteresis loop (MHL). We first define the “envelope curve” as the curve enclosing the thermally hysteretic susceptibility beteween the lower and higher temperature reversible region (see Fig.3). We can draw a MHL during the heating cycle i.e. start heating and increase T from the lower temperature reversible (antiferromagnetic) region and then reverse the direction of temperature before reaching the higher temperature reversible (ferromagnetic) region. We can also draw a MHL in the cooling cycle i.e. start cooling from the reversible ferromagnetic region and reverse the direction of temperature before reaching the lower temperature reversible antiferromagnetic region. If the heating is reversed at sufficiently ‘low’ temperatures the minor loop does not coincide with the cooling part of the ‘envelope curve’. Here in the lower part of the hysteretic regime the high temperature ferromagnetic phase phase is not formed in a sufficient quantity; so when the temperature is decreased the curve does not fall on the cooling part of the envelope curve which represents the curve along which the high temperature phase is supercooled. The MHL’s initiated from temperatures well inside the hysteretic regime coincide with the cooling part of envelope curve indicating that the high temperature phase has formed in a sufficient quantity. In Fig. 4 and 5 we present some representative MHLs both for the Ce(Fe,5%Ir)<sub>2</sub> and Ce(Fe,7%Ru)<sub>2</sub> alloys. We have drawn similar MHLs from the cooling branch of the enevelope curve, which are not shown here for the sake of clarity and conciseness. We have reproduced this behaviour of MHLs over many experimental cycles. The presence of these MHLs clearly suggest the existence of phase coexistence across the ferro- to antiferromagnetic transition. Had there been no phase coexistence we would have followed the cooling part of the envelope curve reversibly on increasing T.Very similar minor hysteresis loop technique has been used to study the phase coexistence associated with a first order metal-insulator transitions in NdNiO<sub>3</sub> .
It should be noted here that the pinning of solitons (domain walls) by lattice defects can also give rise to a thermal hysteresis in magnetic measurements. However, the observed thermal hysteresis in our present study is confined to a relatively narrow temperature window and this argues against such a possibility.
In conclusion we have shown that the ferro- to antiferromagnetic transition in the compounds Ce(Fe,5%Ir)<sub>2</sub> and Ce(Fe,7%Ru)<sub>2</sub> is accompanied by distinct thermal hysteresis as well as signatures of phase-coexistence. We argue that these observations are indicative of the first order nature of the concerned phase transition. The higher tempeature para- to ferromagnetic transition appears to be a typical second order phase transition. These results would support the applicability of Moriya-Usami’s model in explaining the double magnetic transitions in various CeFe<sub>2</sub> based pseudobinary systems. A calorimetric study is now required to confirm the conjecture that this ferro- to antiferromagnetic transition is first order in nature. However, it should be noted that in the case of small latent heats it might be difficult to distinguish a first order transition through calorimetric studies ; in such cases the observed hysteresis and phase coexistence would remain a useful tool for identification of a first order transition.
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# Models for the interpretation of CaT and the blue Spectral Indices in Elliptical Nuclei
## 1 Introduction
The study of chemical abundances in elliptical galaxies has traditionally been performed through the analysis of absorption features usually present in their spectra (see the recent review by Henry & Worthey HW (1999)). The observation of such indices in the blue spectral range —in particular the so called Lick indices— has been a very fruitful tool to interpret the physical properties of elliptical galaxies and globular clusters, both assumed to consist of old stellar populations. There are many articles compiling observational data for some of these indices (e.g. Trager et al. trager (1998) and references therein), specially in Mg<sub>2</sub> and $`\mathrm{Fe}=(\mathrm{Fe5270}+\mathrm{Fe5335})/2`$. Some works have also measured other indices in the same spectral region such as Mgb, NaD and H$`\beta `$.
Evolutionary synthesis models are the tool most frequently used to interpret observed spectra. There are a large number of different models (see Leitherer et al. leitherer (1998), and references therein ) which have become available thanks to the development of theoretical isochrones, computed for a wide range of ages and metallicities. An additional basic input for these models is an atlas of stellar spectra (empirical or theoretical) which provides the spectral energy distribution of each elemental area of the Hertzsprung-Russell Diagram (HRD). If the spectral resolution of the available stellar atlas is good enough, line-strength indices can be measured directly in the final spectrum resulting from the calculation (see Vazdekis V99 (1999)).
When this is not the case, or when the stellar atlas consists of atmosphere models, empirical calibrations of line-strength indices (also known as fitting functions) must be incorporated into the models. These fitting functions are obtained by observing a large sample of stars covering the widest available range of the basic atmospheric stellar parameters (effective temperature T<sub>eff</sub>, surface gravity $`\mathrm{log}g`$, and metallicity —usually parameterized by $`[\mathrm{Fe}/\mathrm{H}]`$—; some authors also include relative abundances $`[\mathrm{X}/\mathrm{Fe}]`$ parameters to introduce elemental ratios different from solar). Among the most employed sets of fitting functions are those provided by the Lick group (Gorgas et al. Gorgas93 (1993); Worthey et al. WFGB94 (1994) —hereafter WFGB94), and those of the Marseille group (Idiart & Freitas-Pacheco idiart95 (1995), Borges et al. BIFT95 (1995); hereafter BIFT95). Examples of evolutionary synthesis models, in which blue spectral indices for single stellar populations (SSP) of different ages and metallicities are computed, are those of Worthey (W94 (1994), hereafter W94), Casuso et al. (casuso (1996)), Bressan et al. (BCT96 (1996), hereafter BCT96), Vazdekis et al. (vazdekis96 (1996) hereafter VCPB96), and Kurth et al. (KFF99 (1999), hereafter KFF99). All these works have employed the polynomial functions of WFGB94. On the other hand, BIFT95 have made use of their own set of fitting functions, which have also been employed in the models of Tantalo et al. (tantalo98 (1998)). Most of these synthesis models give estimates for the blue-yellow line-strength indices, such as Mgb, Mg<sub>2</sub>, Fe5270 and Fe5335 (sometimes only $`\mathrm{Fe}`$), NaD and H$`\beta `$.
An important result is obtained from the study of the locus of data in the plane Mg<sub>2</sub>–Fe (where Fe means an iron index such as Fe5270, Fe5335 or $`\mathrm{Fe}`$). The correlation followed by globular cluster data is adequately reproduced by synthetic models of spectral indices applied to old stellar populations of low metallicities, a result which is not unexpected, since most of the poor-metal stars used to calibrate the spectral index dependence on metallicity are members of these globular clusters. This correlation is steeper than that found for elliptical galaxy nuclei which cannot be fitted by the models even by using the oldest and more metal-rich stellar populations (Burstein et al. B84 (1984); Gorgas et al. GEA90 (1990); Worthey et al. W92 (1992); Davies et al. davies93 (1993); Carollo & Danzinger CD94a (1994),1994b ; Fisher et al. fisher (1996); Vazdekis et al. vazdekis97 (1997)). In fact, elliptical galaxies are located below the lines in the mentioned diagrams. The usual explanation states that old elliptical galaxies formed stars very quickly in the past, after the production of large quantities of magnesium and other elements by massive stars (through the ejection of metals by Type II supernovae, SNe), but before the bulk of iron production, which is mainly synthesized by Type I supernovae resulting from the evolution of low-mass stars. The iron-peak elements appear at least 1 Gyr later than the $`\alpha `$-elements in the interstellar medium. This result limits the star formation duration to less than 1 Gyr, after the start of the process
Therefore, the so-called over-abundance of Magnesium over Iron is actually an under-abundance of Iron, in terms of the absolute values of total abundances, and since Calcium is also an $`\alpha `$-element, it should be expected that Ca indices follow the Mg<sub>2</sub> behavior: if CaT and Mg<sub>2</sub> indices were directly related to the abundances of Calcium and Magnesium, and both elements were mostly produced by Type II SNe, one should expect that a large Mg enrichment would also imply a large proportion of Ca in comparison with the Iron abundance, implying \[Ca/Fe\] $`>0`$, too. On the contrary, if models are not able to reproduce the observational data, a new explanation should be proposed. However, elliptical galaxies data seem to follow the model predictions in the Ca4455–$`\mathrm{Fe}`$ plane (Worthey W98 (1998)) thus implying a $`[\mathrm{Ca}/\mathrm{Fe}]=0`$. Since this result is not well understood, here we propose the use of the Calcium Triplet index at $`\lambda 8600`$ Å, CaT ($`\lambda 8542+8662`$ Å), to test the predictions of theoretical models against observational data in the plane CaT-$`\mathrm{Fe}`$. This point will be discussed in detail in the following sections.
It is important to stress that model predictions used to compare the variation of the CaT with other spectral features in the blue spectral region should be obtained with the same computational technique and inputs, i.e., the same models with identical stellar tracks and atlases of Stellar Energy Distributions, SEDs. For this reason, in this paper we will present index predictions obtained with a revised version of the evolutionary synthesis models already presented by García-Vargas et al. (Paper I (1998), hereinafter Paper I). There we modeled the equivalent width of the two main lines of the CaII Triplet ($`\lambda \lambda 8542,8662`$ Å), following the index definition given by Díaz et al. (DTT (1989), hereafter DTT), for SSPs with ages ranging from 1 Myr to 17 Gyr, and for 4 different metallicities Z=0.004, 0.008, 0.02 and 0.05. An important conclusion derived from Paper I is that the CaT index is almost constant with age, and only dependent on metallicity for ages older than 1 Gyr, when the IR flux is dominated essentially by giants. This result indicates that the CaT is a potential tool, in conjunction with other age-sensitive indices such as H$`\beta `$, to confront the well-known age-metallicity degeneracy problem in old populations. In fact, both indices produce a quasi-orthogonal grid of constant age and metallicity lines (see Fig. 7 in Paper I). In the revised version of the models employed in this work, we have included the computation of the most common indices in the blue spectral region, following the same strategy as that employed in Paper I for the near-IR indices.
In order to check our model results about CaT in that work, we compared the predicted indices with the globular cluster data (see paper I), and obtained a dependence of this CaT for the oldest stellar populations on the metallicity similar to that estimated from those data. Unfortunately, there were just a few CaT observations in elliptical galaxies to compare with the model results. In this new piece of work we have compiled data for a sample of elliptical galaxies, for which both the CaT and Lick indices are available in the literature. This will allow us to compare the predictions of the new models with the indices measured in both blue and near-IR spectral regions.
This paper is organized in the following way: in Section 2 we give a description of the evolutionary synthesis model, with special attention to changes introduced with respect to Paper I, and we discuss the criteria followed to select the fitting functions. The comparison of models with data is shown in Section 3. A discussion is performed in Section 4 and finally, our conclusions are presented in Section 5.
## 2 The evolutionary synthesis model
### 2.1 Model description
The evolutionary synthesis models presented in this paper have been computed using the technique already described in paper I. The total mass of every SSP is $`1\times 10^9`$ M with a Salpeter-type IMF, $`\mathrm{\Phi }(m)=m^\alpha `$, with $`\alpha =2.35`$, from m$`{}_{\mathrm{low}}{}^{}=0.6`$ M to $`m_{\mathrm{up}}=100`$ M. We have followed the passive evolution of a single–burst stellar population (SSP) through ages from 4 Myr to 20 Gyr (with a logarithmic age step of 0.1 dex until 10 Gyr, and 0.02 dex afterwards). This wide range is useful to analyze composite populations, as it occurs in spiral disks (see Mollá et al. molla (2000)) or the non-negligible possibility —non treated here— of having recent star formation over-imposed over an older stellar population (see Pellerin & Robert, pellerin (1999)). In particular it might be used for studying, in the near future, starburst galaxies where the star formation bursts were provoked by radial flows in spiral disks which pushed the gas towards their centers. Thus, this work also is useful to check that this model gives reasonable spectral indices for SSP before applying it to other more complex stellar populations.
The spectral energy distribution for each SSP is synthesized by adding the spectra corresponding to all the points in the theoretical HR diagram, taken from Bressan et al. (bressan93 (1993)) . The spectrum associated to each point was selected using the closest atmosphere model in the stellar parameter space, and properly luminosity scaled ( see García-Vargas et al. GVBD (1995)). Atmosphere models were taken from Lejeune et al. lejeune97 (1997), lejeune98 (1998)), who provide an extensive and homogeneous grid of low-resolution theoretical flux distributions for a large range in T<sub>eff</sub> (from 2000 K to 50000 K), gravity ($`1.02<\mathrm{log}g<5.50`$), and metallicity ($`5.0[\mathrm{M}/\mathrm{H}]+1.0`$), by including M dwarf model spectra. We chose their corrected fluxes, color-calibrated flux distributions, constructed through empirical effective temperature-color relations. These spectra span from 91 Å to 16000 nm, with an average spectral resolution of 20 Å in the optical range. In order to provide predictions for the youngest stellar populations, we supplemented the above data with hot stars (T$`{}_{\mathrm{eff}}{}^{}>50000`$ K) from Clegg & Middlemass (clegg (1987)). These last stars do not exist in old stellar populations but they will be necessary in spiral disk models, where young stars usually contribute to the total flux.
The present models have allowed us to compute the evolution of several line-strength indices: Na i, Mg i and CaT in the near-IR, and Mgb, Mg<sub>2</sub>, Fe5270, Fe5335, NaD and H$`\beta `$ in the blue band. The index definitions given by DTT have been used for the CaT and Mg i indices, whereas the Na i index was computed as in Zhou (Z91 (1991)). The index definitions for the blue indices (Lick/IDS system) can be found in Trager et al. (trager (1998)). Although in the comparison with observational data we are only using CaT, Mg<sub>2</sub>, $`\mathrm{Fe}`$ and H$`\beta `$, the calculation of the additional indices allow to us to check whether our blue model predictions are compatible with those obtained by other authors. In addition, the Mg i and Na i indices in the near-infrared are predictions which may be useful in the comparison with high-resolution spectral data in low velocity dispersion objects.
To calculate the synthetic line-strength indices, as a function of age and metallicity, we have followed the procedure explained in paper I. The Lick/IDS indices have been computed using two different sets of fitting functions: (1) the polynomial functions from WFGB94, and (2) the fitting functions from BIFT95. Both sets of formulae give the behavior of each index as a function of the main stellar parameters: gravity, effective temperature, and metallicity or iron abundance \[Fe/H\]. The BIFT95 functions allow one to include the dependence on $`\alpha `$–element abundance ratios, although in this paper we have assumed $`[\alpha /\mathrm{Fe}]=0.0`$.
Although the CaT was already synthesized in paper I, we have decided to recompute it here in order to consider the differences in the $`m_{\mathrm{low}}`$ (we have employed 0.6 M, instead of 0.8 M), and in the atmosphere models, which now include M dwarf spectra. Thus, this recalculation provides us with a set of indices obtained with the same computational technique. The only difference between the blue and the near-IR indices presented in this paper resides on the stellar libraries used to derive the corresponding indices values. The inclusion of the CaT in the models has been performed by using the predictions, based on model atmospheres, from Jørgensen et al. (JCJ92 (1992), hereafter JCJ92). These authors give the equivalent width of the two strongest calcium lines as a function of the stellar effective temperature, surface gravity, and calcium abundance, and when these functions are used for the DTT stellar library, the predictions are in excellent agreement with the data. Thus, we prefer to use JCJ92 models, which revealed a complex behavior of the calcium lines, instead of relations from DTT, who only gave empirical relations between CaT and the stellar parameters, as a first component analysis, in order to explain the parametric behavior of the star sample.
It is very important to highlight that, since JCJ92 predictions are only valid in the temperature range from 4000 to 6600 K, we have extrapolated their fitting functions for cooler stars (see also paper I). Although this is always dangerous, we are confident in the generic trend of JCJ92 results with Teff, because these relations, extrapolated for Teff $`<4000K`$, were already compared with M stars data in Paper I: a decreasing of CaT with Teff was found in both cases. To strengthen our confidence on this point, we show in Fig. 1 the CaT-Teff relations obtained with the JCJ92 extrapolated functions for different gravities (solid lines), and compare them to CaT data for the cool star samples from Zhou (Z91 (1991)), Mallik (mallik (1997)) and Zhu et al. (zhu (1999)). The two first sample data are shown with symbols representing the stellar gravity. The last set data are solid dots. The dotted line is the least squares fit for these data, given by the equation:
$$CaT=3.2510^3\times Teff4.10$$
(1)
The dependence of the CaT on Teff has a similar slope in both, empirical and theoretical, cases. Data are limited by the extrapolation lines for $`\mathrm{log}g=0.5`$ and $`\mathrm{log}g=2`$. Those with lower gravities seem to be located higher in the graph than those with higher gravities which are lower, by following the same trend than theoretical models. Maybe the dependence of this theoretical function on gravity does not exactly reproduce the observed one for cool stars, but this point cannot be estimated with the small number of stars with known gravities used for our fit. The influence of using one or the other kind of dependence (empirical equation (1) vs JCJ92 functions) for cool stars on model results will be analyzed in the following section. In summary, this being the main source of uncertainty of the present models, we attempt to use the most adequate solution until more confident fitting functions in this spectral range are available.
For Mg i index, the empirical calibration obtained from DTT have been used, whereas in the case of the Na i index we have followed Equations 3 and 5 from Zhou (Z91 (1991)) who calibrated this index as a function of gravity, stellar abundance and (R-I) Johnson-system color for spectral types from G0 to M3. For the coolest stars we assign a value of 1Å, following data shown by Z91 in their Figure 6. The (R-I) color has been taken from isochrones.
### 2.2 Model Results
Table 1 shows the line-strength predictions for three different metallicities and six age steps, using WFGB94 fitting functions. We have selected these ages since the discussion in this paper is focused on old populations in elliptical galaxies. (The whole set of results for six metallicities and ages from $`\mathrm{log}age=6.60`$ to $`\mathrm{log}age=10.30`$ is available in electronic format). <sup>1</sup><sup>1</sup>1Complete set of Table 2 is only available in electronic form at CDS via anonymous ftp to cdsarc.u-strasbg.fr (130.79.128.5) or via http://cdsweb.u-strasbg.fr/Abstract.html, or upon request to authors
The results for the CaT index are very similar to those from paper I. As a comparison, for an age of 13 Gyrs and abundances Z=0.004, 0.008, 0.02 and 0.05, the present models predict the following values: 4.58, 5.91, 7.08 and 8.30 Å, respectively. The corresponding indices in paper I were: 4.65, 5.95, 7.32 and 9.28 Å. These differences are lower than 10%, implying that the new low-mass limit does not have an important effect on the CaT values, although we note a larger difference (1 Å) for Z=0.05 probably due to the effect of including M dwarf spectra in the atlas, which is larger on the metal-rich stellar populations.
These results are also similar to those obtained by other authors, such as Vazdekis et al. (vazdekis96 (1996)), Idiart et al. (idiart96 (1996), hereafter ITFP96), Mayya(mayya (1997)) or even the most recent ones from Schiavon et al. (schiavon (1999), hereafter SBB99). The comparison with other models is made in Table 2 where results for two SSP of 1 Gyr and 13 Gyr old, with different metallicities, are shown. In all cases, the CaT is independent of age and very dependent on metallicity, except for the results from VCPB96 for Z$`=0.05`$, which are smaller than the solar abundance ones.
We represent the time evolution of the results of our model and other models for solar metallicities in Fig. 2, in order to determine whether our conclusion about the independence of CaT on age is a general behavior (Mayya’s model is not shown because this work is only dedicated to stellar populations younger than 1 Gyr). We see that our model gives an almost constant index for old stellar populations in agreement with BCPB96 and SBB99, while models from ITFP97, who used their own stellar library and fitting functions, give CaT values slightly increasing with age ( In order to compare with the same kind of definitions we have converted the CaT values from these authors to the system DTT89). When comparing our model with BCPB96’s results, (the only authors who calculated red (CaT and Mg i) and blue spectral indices with the same model) we observe that our values of CaT are ($``$ 1 Å) smaller. This is due to the different calibration used to calculate the stellar CaT indices: they use DTT while we use the theoretical equations given by JCJ92. The difference between both models is almost the same offset found between the two grids which were computed in Paper I with these two function sets. Mayya (mayya (1997)) also use JCJ92, at least in part, reaching basically the same behavior with the age for stellar populations younger than 1 Gyr. SBB99 also obtain no dependence on age, although the absolute values for CaT are not given in their work. In conclusion, with the present knowledge about this index we are confident in the behavior goodness for this index.
This independence of age does not depend on the extrapolation performed for the coolest stars: when the empirical equation (1) is used, instead of the theoretical equations from JCJ92, we also obtain a constancy of CaT with age, although with a smaller value (6.76 Åvs 7.10 Å– differences lower than 5 %– for solar abundance and $`\mathrm{log}age=10.20`$). This behavior is explained when the isochrones information is taken into account. When the stellar populations become older their stars are cooler, which decreases their CaT values following Fig. 1. At the same time the influence of these stars on the total flux of the stellar population increases with age. Both effects compensate each other, resulting in a contribution of these cool stars to the flux in the CaT index which remains constant with age.
We have also checked our model results for the blue spectral indices by comparing them with those derived by other authors. Table 3 shows that our predictions are similar to those from W94 and BCT96 when polynomial fitting functions from WFGB94 are used. A comparison of the age evolution of the Mg<sub>2</sub> index is shown in Fig. 3. The similarity with BCT96 is not surprising since their input atmosphere models and isochrones are not very different from those employed in this work. Our indices are also close to those from KFF99 calculated with Padova isochrones, although these authors followed a different method to assign the stellar models in the HR diagram. VCPB96 obtain larger values of Mg<sub>2</sub> and lower values of H$`\beta `$. The recent estimations of V99, which are obtained without fitting functions by measuring in the final high-resolution synthetic spectra, are also similar to all models although slightly lower.
The agreement between different authors is clear, with the exception of BIFT95, which predicts lower values than all the other works. Note that, even using the same fitting functions as in BIFT95, we are not able to reproduce their line-strength predictions. This not-well understood discrepancy between the BIFT95 results and our results (and those of other works) leads us to choose our predictions obtained with WFGB94 fitting functions for the comparison with observational data.
## 3 Data Analysis
### 3.1 Comparison of Data with Models
With the aim of addressing the behavior of calcium concerning the $`\alpha `$-element overabundance problem, we start by comparing our model predictions for Mg<sub>2</sub> and CaT with the available data. The set of data is shown in Table 4. This index was measured by Terlevich et al. (TDT90 (1990), hereafter TDT90) for a sample of galaxies, with the aim of finding a method to discriminate between active —i.e. with strong star formation— and normal galaxies. Delisle & Hardy (DH (1992), hereafter DH) also estimated the CaT equivalent width for a reduced sample of elliptical and spiral galaxies, but they used an index definition with different bandpasses. In our compilation there are three sources of IR data: 1) the data from TDT90 ; 2) the set taken from DH. These authors have kindly provided us their spectra in order to calculate again the calcium triplet CaT with the same definitions and windows than the first set; 3) data unpublished from Gorgas et al. (private communication, hereafter GCGV). All these data are in the same measure system. From these sources, we have selected only the galaxies for which blue spectral indices are also available in the literature, following references of column (11). All of them are in the Lick system, therefore we have compiled a uniform set of data. In Column (12) we also show values for Mg<sub>2</sub> as estimated by Davies et al. (davies87 (1987)). These values are systematically lower than those more recent of Column (5).
The measurement of equivalent widths is related to their internal velocity dispersion through the corresponding broadening of the spectral lines which affect both the continuum level and the lines. The effect of the broadening is to decrease the measured EW’s (see DTT89, and their fig. 3, where this point was carefully explained). Therefore, the raw data have been corrected, by adding a broadening correction to all of them, and this is how they are shown in the figures. These corrections are lower than the errors estimated in most of data: for the DTT89 sample, they are $`0.2\AA `$, well within the error bars.
We have already discussed that if magnesium and calcium were produced by the same type of massive stars, they should also trace the same mean abundance, in which case the \[$`\alpha `$/Fe\] overabundance found in E galaxies when studying the Mg lines should also be present in the analysis of the Ca features. We would expect to find all the observational points within the model lines.
This comparison is graphically shown in Fig. 4, where the lines correspond to the model results presented in Table 1: solid lines are isoabundance predictions (Z=0.05, 0.02 and 0.008, from top to bottom), whereas dashed lines indicate isoages values (ranging from 2 to 16 Gyr, from left to right). Interestingly, model predictions in this diagram indicate that Mg<sub>2</sub> and CaT do not exhibit a strong age–metallicity degeneracy, which means that this index–index plane could be useful for the overabundance study. It is clear from the figure that, for single populations older than 2 Gyr, the CaT index is roughly constant for a given abundance, as we already explained in the previous section, whereas Mg<sub>2</sub> is sensitive to both age and metallicity.
When the data compiled in Table 4 are plotted in this diagram (using different symbols to distinguish distinct CaT sources, as explained in the plot key), it is apparent that the observed Mg<sub>2</sub> indices spread beyond the model lines. This effect is equivalent to that found by Worthey et al. (W92 (1992)) in the Mg<sub>2</sub>$`\mathrm{Fe}`$ plane. In this figure, we use solid dots to represent the values of Column (12), while the open symbols are those of Column (5). Both kinds of points represent the same galaxies with different estimations of Mg<sub>2</sub>. It is clear that an uncertainty range appears to be involved in the calculation of the index: there exists an offset between Davies et al. (davies87 (1987)) and data from other authors, mostly Gonzalez (gonzalez (1993)), models being closer to the observations when the first set is used.
The observed CaT values indicate that the calcium abundance is about solar. Note that the location of data is not dependent on the CaT source (although TDT90 indices show a large abundance scatter, ranging from Z=0.02 to Z=0.05, but some of them are active galaxies).
Another piece of information came from the study of the CaT–$`\mathrm{Fe}`$ plane, shown in Fig. 5. Model predictions exhibit a larger degeneracy than that observed in Fig. 4. It is important to note, however, that CaT and $`\mathrm{Fe}`$ are not completely degenerate due to the dependence of the iron indices on age, almost negligible for the CaT. This result can be easily quantified using the metal sensitivity parameter ($`\mathrm{\Delta }\mathrm{log}\mathrm{age}/\mathrm{\Delta }\mathrm{log}\mathrm{Z}`$) defined by Worthey (WFGB94 (1994)): for CaT is 134, while for Mg2 this value is $`1.8`$ and for Fe5270 is $``$ 2.3. It is clear from this calculation that the CaT is an excellent abundance indicator, highly surpassing Fe5709, the iron index most sensitive to metallicity ($`\mathrm{\Delta }\mathrm{log}\mathrm{age}/\mathrm{\Delta }\mathrm{log}\mathrm{Z}=6.5`$) in the optical range.
When the observational data set is included in the above plane (see Fig. 5), it is clear that model predictions reproduce the location of most elliptical galaxy nuclei. In addition the abundances read from these models are mostly solar, for ages between 4–16 Gyr (note that this result is logically the same derived from Fig. 4 since metallicity is derived from the same index, i.e. CaT). There are also some outliers, but most of them correspond to data with large error bars in CaT and/or in $`\mathrm{Fe}`$. It is clear that most elliptical galaxies fall within the model predictions grid.
This diagram confirms the previous trend obtained when comparing $`\mathrm{Fe}`$, with Ca4455 (e.g. Worthey W98 (1998)). We conclude then that the simultaneous analysis of Figs. 4 and 5 adds further weight to the idea that calcium follows iron instead of magnesium, reinforcing the observational evidence that calcium and magnesium behave differently in elliptical galaxies (see section 4 for a discussion).
In order to check the possible influence of the extrapolation used on this diagram, we have also shown, as dotted lines, in Fig. 4 and Fig. 5 the results obtained when the empirical equation (1) is used for the stars cooler than 4000.The CaT values are lower but our former conclusions hold: there are a large number of points falling out of the diagram in Fig. 4 while in Fig. 5 some data now fall out but also in the opposite direction.
Other possible error source comes from the uncertainties in the RGB temperature which may be 200 degrees cooler or hotter. We have estimated the effect of a reduction of 200 K in the effective temperature for stars in the RGB on our indices predictions, by indicating these variations in each figure as an arrow located in a corner of each diagram. This shift is not sufficient to reach the data region in Fig. 4, even using the empirical calibration for the coolest stars.
### 3.2 Disentangling age and metallicity: the CaT-H$`\beta `$ plane
One of the main problems to understand stellar populations in early–type galaxies is how to disentangle age and metallicity effects. Gonzalez (gonzalez (1993)) showed that the combination of the H$`\beta `$ line-strength with indices like Mg<sub>2</sub> or $`\mathrm{Fe}`$ could break the degeneracy. It is well known that H$`\beta `$ is affected by nebular emission. To overcome this problem, higher-order Balmer lines, like H$`\gamma `$, have been proposed as a powerful alternative (Jones and Worthey Jones95 (1995); Vazdekis and Arimoto V99b (1999)). Unfortunately, not many accurate data on this feature have been presented in previous works (but see Kuntschner and Davies KD (1998)). Since the aim of this work is not to derive absolute ages and metallicities but to gather information from the relative trends in the index–index diagrams, we have preferred to use the classical H$`\beta `$ index, although we note that some particular galaxies may be partially affected by an emission component (see notes in Table 4).
The new model predictions in the H$`\beta `$–CaT plane are displayed in Fig. 6. The simultaneous computation of blue and near-IR indices performed in this work allows us to confirm the previous finding of Paper I, already discussed in the introduction to this paper, concerning the orthogonality of the CaT–H$`\beta `$ diagram and the advantage of this diagram to separate age and metallicity effects. Almost all the data points fall inside the new model predictions, within their error bars. Once again we represent the results obtained through the empirical fit for CaT as dotted lines and an arrow is included in the graph to indicate the possible shift of points if the RGB effective temperature is reduce 200 K.
Taking the models literally, the above diagram indicates that the stellar populations exhibit roughly solar abundances for \[Ca/H\]. This is naturally the same result obtained from Figs. 4 and 5. In addition, the derived ages range from 4 to 16 Gyr, except for NGC 221, with an age of $`2`$ Gyr, which is not too far from the result derived by Vazdekis & Arimoto (V99b (1999)), who found that the better fit of a synthetic spectrum to this galaxy is obtained for ages in the range from 2.5 to 5 Gyr. Obviously, these results are not new since they rely on the H$`\beta `$ indices.
An important problem associated with the use of the Balmer absorption features is that their continuum bandpasses usually include metallic lines. Another important drawback, specially in the case of H$`\beta `$, is that it is well established that a large fraction ($`50`$ %) of early-type galaxies exhibit Balmer emission lines at some extent. Therefore, the determination of ages is not totally safe.
## 4 Discussion
As we have already mentioned, a ratio $`[\mathrm{Mg}/\mathrm{Fe}]>0`$ can be explained by taking into account that both elements are synthesized in the interior of stars covering a different range of stellar masses. Invoking the same reasoning for calcium, it can be concluded that this element is generated in another type of stars than those synthesizing Mg. In this sense, Worthey (W98 (1998)) has suggested that Type II SNe could have two flavors, just like Type I and II SNe produce iron and magnesium at different times in the evolution of a galaxy.
It is important to note that the ratio \[$`\alpha `$/Fe\] is high because the iron abundance \[Fe/H\] is low, i.e., an iron deficiency translates into an over-abundance for alpha-elements. In particular, if one assumes \[$`\alpha `$/Fe\]=+0.4 for a given total abundance Z, it means that, for the total abundance Z, $`\mathrm{log}Z/Z_{}`$ is approximately the same and \[Fe/H\] is $`0.4`$ dex smaller. Isochrones do not change very much when alpha-elements are enhanced because abundances of the principal elements are also enhanced; but the iron abundance is decreased (see Tantalo et al. tantalo98 (1998)).
The overabundance in $`\alpha `$-elements, or underabundance in \[Fe/H\], which can be seen in the models, occurs because the bulk of the stars are created before the iron is produced. It may be seen when the star formation is low, usually in low-metallicities regions. Thus, the ratio \[Ca/Fe\] is over-solar in metal-poor stars (Wallerstein, wallerstein (1962); Hartmann & Gehren, Hartmann (1988); Zhao & Magain, zhao (1990); Gratton & Sneden, gratton (1991)), such as the \[Mg/Fe\] is. This overabundance decreases until solar ratios are reached when the iron abundance increases. For disk dwarfs, Edvardsson et al. (edvardsson (1993)). determined the calcium abundances in the range $`1.0<[\mathrm{Fe}/\mathrm{H}]<0.2`$ dex, finding that \[Ca/Fe\] $`0.25`$ dex at \[Fe/H\]$`=1.0`$, and then decreases to solar values just following the magnesium. This kind of behavior may also occurs for high metallicities if the time scale for the formation of the bulk of stars is short enough to create stars before the iron appearance, $`1`$ Gyr. McWilliam & Rich (mcwilliam (1994)) showed, for the Galactic bulge stars, an overabundance in \[Ca/Fe\]$`0.3`$, slightly lower than those of magnesium and silicon.
Chemical evolution models, including recent theoretical calculations of yields (Timmes et al. timmes (1995)), reproduce the Galactic abundances for the so called intermediate alpha-elements (Mg, Si, Ca and Na) as the result of the evolution of the massive stars. (We must note that it is necessary to increase the magnesium yield by a factor of 2 in order to obtain the solar abundance of Mg). Thus the bulk of the star formation should take a time which is shorter than $`1`$ Gyr (the minimum time that has to elapse before the low-mass stars eject the iron in the Type I SNe explosions) in regions where \[$`\alpha /\mathrm{Fe}]>0.0`$, and longer for the systems where \[Mg/Fe\]=0.0. This result is well reproduced by chemical evolution models (after adequately fitting the data in the solar neighborhood): in Fig. 5 of Molla & Ferrini (MF (1995)) we can see a similar behavior for Mg, Ca and Si when stars are formed in a short time scale ($`0.8`$ Gyr).
The fact that $`[\mathrm{Mg}/\mathrm{Fe}]>0`$ in elliptical galaxy nuclei implies that star formation may extend over a time shorter than 1 Gyr. If we now introduce the observed \[Ca/Fe\] ratio into this reasoning, and we assume that both calcium and magnesium are produced by similar Type II SNe (although with different masses), we can constrain even more the duration of the star formation episode in these regions. In Fig. 7 we represent, simultaneously, the fraction of ejected mass of <sup>24</sup>Mg and <sup>40</sup>Ca over the total yield for each element, taken from Woosley & Weaver (WW95 (1996)), as a function of the stellar mass (lower x-axis) and the mean stellar life (upper x-axis, in logarithmic scale). It is clear from this figure that calcium is produced by stars with mass in the range 12–30 M, while magnesium is generated in stars of 20–40 M. This is in agreement with the more recent metallicity-dependent yields from Portinari et al. (portinari (1998)), who show in their Fig. 4 the ejecta of each element as a function of the CO-core, indicating that calcium proceeds from lower mass stars than those which produce magnesium. Considering the lifetimes of the stars responsible for the bulk production of Ca and Mg, we find these values peak at $`8.9`$ and $`6.3`$ Myr, respectively. Therefore, even with a short difference between both lifetimes ($`2.6`$ Myr), there exists a time period at which the magnesium has already been ejected by the most massive stars but the bulk of the calcium has not yet been released. After this small delay, the calcium appears in the ISM, and stars formed afterwards will incorporate both elements, Mg and Ca, showing \[Mg/Fe\] and \[Ca/Fe\] larger than zero and a ratio \[Mg/Ca\] tending to solar.
The consequence of this reasoning is important: if in elliptical galaxy nuclei calcium is not over-abundant, while magnesium is, their stellar populations could have been formed, at least locally, in a very short burst, lasting a timescale smaller than a few Myr. <sup>2</sup><sup>2</sup>2A word of caution must be said when considering the speed of the star formation in a stellar system. The quoted times always refer to local times, i.e. those needed to produce the bulk of the star formation locally, although the required times to form the whole system, for instance a galaxy, could be much larger. If the star formation took more than this time, the calcium would have been incorporated into the new stars, and one should expect to observe $`[\mathrm{Mg}/\mathrm{Ca}]0`$, which does not seem to be the case. This is in agreement with the dissipative theory of galaxy formation where the collapse time scale must be shorter for more massive galaxies by producing a correlation between the time scale of the star formation in these galaxies and the total mass. This is supported by Richer et al. (richer (1998)), who claim that \[O/H\] varies systematically with the velocity dispersion for spheroidals, bulges and ellipticals, by relating the gravitational potential well with this effect and by concluding that the time scale governs \[O/Fe\] or \[Mg/Fe\]. The result implies a short time scale for star formation in more massive elliptical and bulges, supporting the old idea that ellipticals are simpler than low mass galaxies from their star formation history.
However, the deduced constraint for the local star formation time might seem to be too short. Then a new scenario, which could include the two Type II SNe flavors, must be invoked. This effect can also be enhanced (or even dominated) by a flatter initial mass function due to special conditions (dense environment or high metallicity). If an IMF is biased towards massive stars, the magnesium produced or yield in these elliptical nuclei will be larger, allowing one to find overabundances of magnesium and also of other elements ejected by the most massive stars such as oxygen, but not of those ejected by less massive stars. An IMF universal and constant is still a matter of discussion (Larson larson (1998), Padoan et al. padoan97 (1997), Meyer et al. meyer (1999), Chiosi et al. chiosi (1998)), but some recent works claim that it must be constant (Wyse wyse (1997), Tsujimoto et al. tsujimoto (1999), Chiappini et al. CMP (1999)).
One point that may enlighten this problem proceeds from the nucleosynthetic yield calculations: by using the existing yields in a chemical evolution model, it is not possible to reproduce the solar abundance \[Mg/H\]. It is necessary to add magnesium to the stellar production to reach the estimated solar value. The quantity of ejected magnesium must be a factor of $``$2 greater that the present yield, implying that Fig. 7 must change. If the deficiency in magnesium might be accounted for by the production of the most massive stars (M$`>40\mathrm{M}_{}`$), we might solve two problems: chemical evolution models would reproduce the solar abundances and the difference between stellar lifetimes for stars producing Mg and Ca would be large enough to allow the formation of stellar populations with an overabundance \[Mg/Ca\] in a sufficient time.
Another possible solution to this problem, besides the actual ratio \[Mg/Fe\] non solar, may be related with the measures of Mg<sub>2</sub> as we see by comparing columns (12) and (5) from Table 4. For some galaxies there exist data from different authors and there are differences as large as 0.05, that is, 15 %. Following Goudfrooij & Emsellem (Goudfrooij (1996)) the possible emission lines may affect the measures of absorption lines. They estimated that the index Mgb may be artificially enhanced by 0.4-0.1 Å and Mg<sub>2</sub> by 0.03 mag, due to the \[N i\] emission-line doublet at 5199 Å. Taking into account that 50 % of giant elliptical galaxies exhibit H$`\alpha +`$ \[N ii\] emission, maybe the actual Mg<sub>2</sub> values must be reduced. If we reduce the Mg<sub>2</sub> data by 15 %, most of them might be almost reproduced by models, as we show in Fig. 4 where the two sets of data (open vs full) are represented.
In this context, a surprising result has been obtained by Origlia et al. (origlia (1997)). These authors have measured the strength of the infrared absorption line at 1.59 $`\mu `$m, which is primarily sensitive to the total silicon abundance. These authors used this absorption to estimate the ratio \[Si/Fe\] necessary to reproduce their EW, by comparing a sample of elliptical galaxies and globular clusters with models. They found that silicon could be enhanced by about 0.5 dex in both kinds of stellar systems. Since Si and Ca proceed from stars with similar masses (Portinari et al.portinari (1998)), both should be overabundant at the same time, contrary to our and their findings and to our explanations for the existence of an overabundance in \[Mg/Ca\]. Thus, the question remains unclear until more precise calculations of nucleosynthesis yields and observations of alpha-elements become available.
## 5 Conclusions
We have developed an evolutionary synthesis model with which we have produced a grid of models for SSP at 6 metallicities and a wide age range. This code is able to predict indices in the blue-visible spectral range, the classical Lick indices and indices in the near-IR such CaT, Mg i and Na i at the same time.
We have carefully analyzed the behavior of this index for the coolest stars (Teff $`<`$ 4000), given by some samples of data available in the literature, obtaining the generic trend of CaT decreasing with effective temperature, in agreement with the extrapolated JCJ92’s theoretical functions. Therefore, although systematic effects may still be present in the theoretical predictions, mainly due to potential errors in the extrapolation of these JCJ92 equations, we use the most adequate solution until more reliable fitting functions become available.
We have compiled a set of data from the literature for galaxies for which both kind of indices, blue-visible and near-IR, had been observed, and we have compared their predictions with data.
We have used our results to study the relationship between different indices for old stellar populations by producing diagnostic diagrams in which observed data and models can be plotted to determine the basic physical properties of the dominant stellar population in these galaxies. We find that most of the galaxies with known data for Mg2, CaT and $`\mathrm{Fe}`$ remain in the CaT – Mg<sub>2</sub> plane at the place of Z=0.02, while they seem to have overabundances of Mg. This conclusion is in agreement with other data of Ca4300 in the blue region found by Worthey (W98 (1998)), but raises the question of how it is possible to have solar metallicity ratio for the calcium element and over-solar ratio for Mg abundance, while both are $`\alpha `$-element produced by the same type of massive stars.
If we accept the assumption of a relative abundance \[Mg/X\]$`>0`$, adopted to explain this kind of diagram and that it is due to a short time scale for the star formation, and we apply the same argument for the \[Mg/Ca\], this would imply that the star formation time scale in elliptical galaxy nuclei must be shorter than $`510`$ Myr. Otherwise, we should not find a discrepancy between the data and the models in the later diagnostic, where both indices proceed from the same kind of alpha-elements, which is not the case. We suggest that an update of the nucleosynthesis yields of Mg, increasing the production of Mg for the more massive stars, may solve this problem, by extending the elapsed time between the production of Mg and of Ca and, in consequence, the time scale for the star formation. An alternative explanation might be an IMF biased towards the massive stars ($`\mathrm{M}>40\mathrm{M}_{}`$) in the early phases of star formation in elliptical galaxies.
We must keep in mind that the emission over these spectra also may affect Mg<sub>2</sub> data: they may be reduced by 5 % if a careful analysis is done before obtaining the spectral indices Mgb and Mg2.
We use the orthogonal diagram CaT – H$`\beta `$ to date elliptical galaxies and to determinate their abundances, reaching the conclusion that elliptical galaxies are nearby solar in their abundances of calcium, in the same way as for iron, and that their ages range between 8 -16 Gyr.
A large campaign of observations in the near-infrared to estimate the CaT index, e.g. in the same set of galaxies given in Davies et al. (1987), would be very useful to date them and to determinate their metallicities/abundances in a clear way.
###### Acknowledgements.
We thank J. Gorgas and N. Cardiel for the permission to use their CaT data before publication and for fruitful discussions. We are grateful to Sonya Delisle and Eduardo Hardy who kindly sent us their spectra in the red bandpass. We thank the referee, Guy Worthey, for his useful comments for the improvement of this paper. We also thank A. I. Díaz for suggestions in the final version of this work and J. Gea Banacloche for the help in the english version. M.M. thanks the Université Laval (Quebec) for the nice period, during which part of this work was done. M.M. has been supported by a post-doctoral fellow of the Spanish Ministerio de Educación y Cultura. This work has made use of the NASA Astrophysics Data System.
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# REFERENCES
Comment on ”Spin relaxation in quantum Hall systems”.
In the recent publication the authors have considered the spin relaxation in a 2D quantum Hall system for the filling factor $`\nu 1`$. The interest to this system is determined by the fact that the spin relaxation here is due to the collective excitations, the so-called spin-excitons. Thus, correct description requires taking into account the Coulomb interaction and here we deal with the delocalized states. The authors considered only one spin flip mechanism among the three possible spin-orbit related ones. This direct spin-phonon coupling is described by the following term in the Hamiltonian
$$\widehat{}_{so}=\frac{1}{2}V_0\widehat{𝝈}\widehat{𝝋};\widehat{\phi }_x=\frac{1}{2}\{u_{xy},\widehat{p}_y\}_+,$$
(1)
where $`\widehat{𝐩}`$ is 2D electron momentum operator, the z-axis coincides with the normal to the 2D plane, $`\{,\}_+`$ denotes the anticommutator, $`u_{ij}`$ is the lattice strain tensor due to the acoustic phonons, and $`V_0=810^7cm/s`$. The authors came to the conclusion that the spin relaxation time due to this mechanism is quite short: around $`10^{10}`$ s at B=10 T (for GaAs) which is much shorter than the typical time ($`10^5`$ s) obtained in Ref. while considering the spin relaxation of 2D electrons in a quantizing magnetic field without Coulomb interaction and for the same spin-phonon coupling Eq(1). The authors related this fact to the presence of the Coulomb interaction and argued that now they are able to explain the earlier experimental data by M.Dobers et al. I show that their conclusion about the value of the spin-flip time is wrong and have deduced the correct time which is by several orders of magnitude longer. I also discuss the admixture mechanism of the spin-orbit interaction. The authors obtained the following expression for the spin- relaxation rate:
$$\frac{1}{\tau _{so}}=\underset{𝐐,i=x,y}{}\frac{\pi }{2\mathrm{}}\lambda ^i(𝐐)^2\delta (E_{ex}(q)\mathrm{}\omega _Q),$$
(2)
where $`𝐐=(𝐪,Q_z)`$ and $`\omega _Q`$ are the phonon wave vector and the dispersion law, $`E_{ex}(q)`$ is the 2D dispersion law for the spin-excitons, $`𝝀\mathbf{(}𝐐\mathbf{)}`$ is the constant which couples the projected spin density of the electrons and phonon creation operator in the electron-phonon part of the second-quantized Hamiltonian. Using Eq.(1), we get: $`\lambda ^i(𝐐)^2=(\mathrm{}^3V_0^2/32\rho sQ)q^2q_i^2\mathrm{exp}(q^2l_B^2/2)\mathrm{\Lambda }^2`$, where $`\mathrm{\Lambda }(Q_z)=𝑑z\chi _0^2(z)\mathrm{exp}(iQ_zz)`$, $`s`$ is the sound velocity, $`\rho `$ the crystal mass density, $`l_B`$ the magnetic length and $`\chi _0`$ is the wave function in the z- direction. Note that parameter $`x_0=(\mathrm{}s/ϵ_cl_B)^2ms^2/E_B510^4`$ in GaAs, here $`ϵ_c=e^2/\kappa l_B,E_B`$ are the Coulomb and Bohr energies. Then the main contribution in Eq.(2) comes from $`Q_zϵ_c/\mathrm{}sq1/l_B`$ and we obtain: $`1/\tau _{so}(0.48\sqrt{2}/4\pi ^{3/2})(V_0/s)^2(\mathrm{}/\rho l_B^5)x_0^{1/2}`$. The estimate for GaAs at B=10 T gives $`\tau _{so}410^5s`$ (we assumed $`\mathrm{\Lambda }=1`$ which can only overestimate the rate).
It is much more essential, that with taking into account the spin-independent interaction with the phonons the energy relaxation time is much shorter than the spin relaxation time and realistic situation corresponds to the spin relaxation of the 2D electron system which has the lowest possible energy. Then the characteristic energy transferred to the phonon during the spin-flip transition is much smaller than $`ϵ_c`$ which appeared in Eq.(2). Now it is determined by Zeeman energy $`\mathrm{\Delta }`$ or temperature $`T`$. (It remains unclear why the authors who assumed the temperature of the electron system to be much smaller than $`\mathrm{\Delta }`$, obtained Eq.(2)). Two different physical situations are possible depending on the degree of excitation of the spin system (i.e. the number of the spin-excitons created in the 2D system). If this number is macroscopically large (exceeds some critical value $`N_cT`$), then the dominating channel of the spin relaxation is the annihilation of two ”zero” (i.e. with $`q=0`$) excitons from the condensate with simultaneous generation of a ”nonzero” exciton and phonon. The interaction with the phonons was spin-independent but the spin-orbit term in the Hamiltonian $`\beta (\widehat{\sigma }_x\widehat{p}_x+\widehat{\sigma }_y\widehat{p}_y)`$ which is due to the absence of the inversion symmetry in the crystal cell leads to an admixture of the state with the opposite spin and allows spin-flip transition . The nonexponential relaxation of the $`S_z`$ component is described by exactly the same equation as that obtained in Ref. but with time (for the interaction with piezo-phonons) $`1/\tau _0(\mathrm{\Delta }Ms^2M\beta ^2/\mathrm{}^3\omega _c^2)((eh_{14})^2/\mathrm{}s^3\rho )`$, which is written for $`T,Ms^2\mathrm{\Delta }`$. Here $`h_{14}`$ is the piezomodulus and $`M`$ the exciton mass. This time does not depend on the magnetic field and is $`10^5÷10^6s`$ depending on $`\beta `$. The mechanism described by Eq.(1) gives a much smaller contribution to the rate. In the case of $`NN_c`$ or near equilibrium when the main process is a direct recombination of the ”nonzero” excitons, the spin relaxation rate is proportional to the temperature and can be also relatively small (time is longer than $`10^6`$ s for $`T<1K`$) .
Alexander V. Khaetskii,
Institute of Microelectronics Technology, Russian Academy of Sciences, 142432, Chernogolovka, Russia
PACS numbers: 73.40.Hm; 71.35.-y; 76.20.+q
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# Untitled Document
A simple method of determining the Hubble constant <sup>1</sup><sup>1</sup>1This essay received an “honorable mention” in the 1999 Essay Competition of the Gravity Research Foundation.
Yi-Ping Qin<sup>1,2,3,4</sup>
<sup>1</sup> Yunnan Observatory, Chinese Academy of Sciences, Kunming, Yunnan 650011, P. R. China; E-mail: ypqin@public.km.yn.cn
<sup>2</sup> National Astronomical Observatories, Chinese Academy of Sciences
<sup>3</sup> Chinese Academy of Science-Peking University joint Beijing Astrophysical Center
<sup>4</sup> Yunnan Astrophysics Center
Abstract
> Bidirectional relativistic proper motions of radio components of nearby extragalactic sources give a strong constraint on the determination of the Hubble constant $`H_0`$. Under the assumption that the real velocity of radio components of extragalactic sources is not less than that of Galactic sources, the value of $`H_0`$ can be estimated at a high level of accuracy. The assumption is reasonable due to the general belief that the activity in the core of galaxies must be more powerful than that of stars. This method is simple and with only one uncertainty — the real velocity of components. This uncertainty is related to the value of the real velocity of componenets of Galactic sources and the latter is always well-determined (note that the determination is independent of $`H_0`$ and the distance of Galactic sources can be directly measured at a rather high level of accuracy). Hopefully the method will play an important role in future research to fix the value of $`H_0`$. With the data of the three sources available so far and the assumption that the real velocity of componenets of at least one of the sources is not less than a known velocity of componenets of a Galactic source, the constant is estimated to be within $`27.08kms^1Mpc^1<H_053.15kms^1Mpc^1`$ with this method.
The real value of the Hubble constant has been a hot topic for a long time. In the past few years, some advances have been achieved (for more details, see Trimble and Aschwanden 1999, Trimble and McFadden 1998). An exciting result of observation by HST led to $`H_0=80\pm 17kms^1Mpc^1`$ (Freedman et al. 1994). However, with the same data, some people got a much smaller value, e. g., $`H_0=40kms^1Mpc^1`$ (Sandage et al. 1994). Among the many methods, that taking the peak brightness of Type Ia supernovae as a distance indicator is generally used. It gave a small value of the constant, $`H_0=61\pm 10kms^1Mpc^1`$, by Brach (1992).
Bidirectional relativistic motions of extragalactic radio sources can be used to estimate the constant (Marscher and Broderick 1982). Recently there were some sound works using this method published. The most successful one was done by Taylor et al. (1997). However, the method they used is somewhat complicated. And it does not tell how the uncertainties of the measurements used affect the estimation of $`H_0`$, and what one should do when the data of several sources are available.
In the following, we illustrate a simple method of determining the Hubble constant using bidirectional relativistic proper motions of extragalactic radio sources.
The apparent transverse velocities of components of a source along an axis at an angle $`\theta `$ to the line of sight at a velocity $`\beta c`$ can be expressed as (Rees 1966, 1967)
$$(\beta _{app})_a\frac{\mu _aD_L}{c(1+z)}=\frac{\beta \mathrm{sin}\theta }{1\beta \mathrm{cos}\theta },$$
(1)
$$(\beta _{app})_r\frac{\mu _rD_L}{c(1+z)}=\frac{\beta \mathrm{sin}\theta }{1+\beta \mathrm{cos}\theta },$$
(2)
where $`a`$ and $`r`$ stand for the motions of the approaching and receding components, with $`\mu _a`$ and $`\mu _r`$ being the corresponding proper motions, respectively.
These two equations give
$$\frac{D_L}{c(1+z)}=\frac{1}{2\mu _a\mu _r}\sqrt{\beta ^2(\mu _a+\mu _r)^2(\mu _a\mu _r)^2}.$$
(3)
For small redshift sources, the following is maintained for all kinds of the universe within the framework of the FRW cosmology,
$$\frac{D_L}{1+z}\frac{cz}{H_0},$$
(4)
where $`H_0`$ is the Hubble constant of the universe.
From Equation (3) one has
$$H_0\frac{2\mu _a\mu _rz}{\sqrt{\beta ^2(\mu _a+\mu _r)^2(\mu _a\mu _r)^2}}.$$
(5)
It shows that, for a nearby extragalactic source with measured values of $`z`$, $`\mu _a`$ and $`\mu _r`$, when $`\beta `$ is known, the Hubble constant would be well determined. This relation provides a very strong constrain on the determination of $`H_0`$.
From Equation (5), the law of $`\beta <1`$ leads to
$$H_0>z\sqrt{\mu _a\mu _r}$$
(6)
for any sources. Therefore, among many values of the lower limit of $`H_0`$, calculated from various sources, the largest one would be the closest value to the limit. If $`\mu _a`$ and $`\mu _r`$ are presented in the form $`\mu \pm \mu `$ for a source, then the lower limit of $`H_0`$ determined by the source should be
$$H_{0,\mathrm{min}}=z\sqrt{(\mu _a\mu _a)(\mu _r\mu _r)},$$
(7)
with $`H_{0,\mathrm{min}}`$ satisfying
$$H_{0,\mathrm{min}}<H_0.$$
(8)
In this way, the largest value of $`H_{0,\mathrm{min}}`$ among those calculated from various sources should be taken as the best estimation of the lower limit of $`H_0`$.
Let
$$\alpha 1\beta .$$
(9)
Since $`0\beta <1`$, then $`0<\alpha 1`$. For a small value of $`\alpha `$, Equation (5) gives
$$H_0z\sqrt{\mu _a\mu _r}[1+\frac{(\mu _a+\mu _r)^2}{4\mu _a\mu _r}\alpha ].$$
(10)
Considering the case where $`\alpha `$ is known to the extent of $`\alpha \alpha _{\mathrm{max}}`$ for a source, the upper limit of the Hubble constant would be determined by the source in the way
$$H_0H_{0,\mathrm{max}},$$
(11)
where
$$H_{0,\mathrm{max}}=z\sqrt{(\mu _a+\mu _a)(\mu _r+\mu _r)}[1+\frac{(\mu _a+\mu _a+\mu _r+\mu _r)^2}{4(\mu _a+\mu _a)(\mu _r+\mu _r)}\alpha _{\mathrm{max}}].$$
(12)
In determination of the range of $`H_0`$, there is a reasonable demand that the ranges of $`H_0`$ estimated from different sources should be overlapped. This demand is consistent with the above principle of choosing the lower limit of $`H_0`$. When determining the upper limit of $`H_0`$ from (12), the requirement can be realized by adopting various values of $`\alpha _{\mathrm{max}}`$ for different sources.
If among these sources, there is at least one source satisfying $`\alpha \alpha _{\mathrm{max}}`$ for a given $`\alpha _{\mathrm{max}}`$, the largest value of $`H_{0,\mathrm{max}}`$ calculated with this value of $`\alpha _{\mathrm{max}}`$ for all the sources should be taken as the best estimation of the upper limit of $`H_0`$. In this way, while the given value of $`\alpha _{\mathrm{max}}`$ is assigned to the source of the largest value of $`H_{0,\mathrm{max}}`$, some larger values of $`\alpha _{\mathrm{max}}`$ should be assigned to other sources, so that the estimated ranges of $`H_0`$ may be overlapped.
Since the value of $`\alpha _{\mathrm{max}}`$ for Galactic sources can be calculated at a rather high level of accuracy, that for extragalactic sources can then be well settled by assuming that it would not be less than that for Galactic sources. This assumption is reasonable due to the general belief that the activity in the core of galaxies must be more powerful than that of stars.
Recently, several Galactic sources were found to have bidirectional relativistic proper motions of radio components. Among them, the largest and well calculated value of $`\beta `$ is $`0.92\pm 0.02`$ for GRO J1655-40 (Hjellming and Rupen 1995). This corresponds to $`\beta _{\mathrm{min}}=0.9`$ and $`\alpha _{\mathrm{max}}=0.1`$ for the source. Therefore, at present, it is reasonable to take $`\alpha _{\mathrm{max}}=0.1`$ for extragalactic sources according to the assumption.
Till now, there are only a few extragalactic sources with measured values of $`\mu _a`$ and $`\mu _r`$ found in literature. Excluding those with high redshifts or uncertain values of proper motions, there are only three sources suitable for our study. They are: 1146+596 (NGC 3894), $`z=0.01085`$, $`\mu _a=0.26\pm 0.05masyr^1`$ and $`\mu _r=0.19\pm 0.05masyr^1`$ (Taylor et al. 1998); 0316+413 (3C 84), $`z=0.0172`$, $`\mu _a=0.58\pm 0.12masyr^1`$ and $`\mu _r0.28masyr^1`$ (Marr et al. 1989, Vermeulen et al. 1994); 1946+708, $`z=0.101`$, $`\mu _a=0.117\pm 0.020masyr^1`$ and $`\mu _r=0.053\pm 0.020masyr^1`$ (Taylor and Vermeulen 1997).
For the lower limit of $`H_0`$, the first and the third sources give $`H_{0,\mathrm{min}}=8.82kms^1Mpc^1`$ and $`27.08kms^1Mpc^1`$, respectively from (7), while the second source gives no lower limit of $`H_0`$. According to the above principle of choosing $`H_{0,\mathrm{min}}`$, the best estimation of the lower limit of $`H_0`$ from these data should be $`H_{0,\mathrm{min}}=27.08kms^1Mpc^1`$.
For the given value of $`\alpha _{\mathrm{max}}=0.1`$, the three sources give $`H_{0,\mathrm{max}}=15.45kms^1Mpc^1`$, $`40.51kms^1Mpc^1`$, and $`53.15kms^1Mpc^1`$, respectively from (12). Assuming that there is at least one source satisfying $`\alpha \alpha _{\mathrm{max}}`$ for $`\alpha _{\mathrm{max}}=0.1`$, then according to the above requirement we choose $`H_{0,\mathrm{max}}=53.15kms^1Mpc^1`$ as the best estimation of the upper limit of $`H_0`$.
Therefore, we obtain the range of $`27.08kms^1Mpc^1<H_053.15kms^1Mpc^1`$ for the Hubble constant from the data of the three sources.
In practice, the observable bidirectional relativistic proper motions of radio components of a source are always those moving almost perpendicular to the line of sight, and then the values of their $`\mu _a`$ and $`\mu _r`$ are close (see, e.g., Taylor et al. 1998). Therefore $`\frac{(\mu _a+\mu _a+\mu _r+\mu _r)^2}{4(\mu _a+\mu _a)(\mu _r+\mu _r)}1`$ and $`\frac{(\mu _a+\mu _a+\mu _r+\mu _r)^2}{4(\mu _a+\mu _a)(\mu _r+\mu _r)}\alpha _{\mathrm{max}}0.1`$ for $`\alpha _{\mathrm{max}}=0.1`$. It shows that taking $`\alpha _{\mathrm{max}}=0.1`$ will only produce about $`10\%`$ uncertainty for $`H_{0,\mathrm{max}}`$. When more such sources have been observed, the expected value of $`\alpha _{\mathrm{max}}`$ and the corresponding uncertainty will be smaller.
This method depends on only one assumption and it concerns only one uncertainty — the real velocity of components. This uncertainty concerns the value of the real velocity of components of Galactic sources and the latter is always well determined (note that the determination is independent of $`H_0`$ and the distance of Galactic sources can be directly measured at a rather high level of accuracy). The assumption is weak due to the general belief that the activity in the core of galaxies must be more powerful than that of stars. Also, the method is simple. To determine $`H_0`$ at a high level, one only needs to measure $`\mu _a`$ and $`\mu _r`$ at a high level of accuracy and finds a sufficient number of such sources (say, 10 or more). The method is then hopeful to play an important role for finally fixing the value of the Hubble constant in future researches.
ACKNOWLEDGEMENTS
It is a pleasure to thank G. B. Taylor, R. C. Vermeulen and Y. H. Zhang for kindly supplying necessary materials. This work was supported by the United Laboratory of Optical Astronomy, CAS, the Natural Science Foundation of China, and the Natural Science Foundation of Yunnan.
> REFERENCES
>
> Brach, D. 1992, Astrophys. J., 392, 35
>
> Freedman, W. L. et al. 1994, Nature, 371, 757
>
> Hjellming, R. M., and Rupen, M. P. 1995, Nature., 375, 464
>
> Marr, J. M., Backer, D. C., and Wright, M. C. H. et al. 1989, Astrophys. J., 337, 671
>
> Marscher, A. P., and Broderick, J. J. 1982, In Extragalactic Radio Sources, D. S. Heeschen and C. M. Wade, eds. (Dordrecht: Reidel), p. 359
>
> Rees, M. J. 1966, Nature, 211, 468
>
> Rees, M. J. 1967, Mon. Not. R. Astron. Soc., 289, 945
>
> Sandage, A. et al. 1994, Astrophys. J., 423, L13
>
> Taylor, G. B., and Vermeulen, R. C. 1997, Astrophys. J., 485, L9
>
> Taylor, G. B., Worbel, J. M., and Vermeulen, R. C. 1998, Astrophys. J., 498, 619
>
> Trimble, V., and Aschwanden, M. 1999, Pub. Astron. Soc. Pacific, 111, 385
>
> Trimble, V., and Mcfadden, L.-A. 1998, Pub. Astron. Soc. Pacific, 110, 223
>
> Vermeulen, R. C., Readhead, A. C. S., and Backer, D. C. 1994, Astrophys. J., 430, L41
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# Rigorous Effective Field Theory Study on Pion Form Factor
##
The process $`e^+e^{}\pi ^+\pi ^{}`$ at energies lower than the chiral symmetry spontaneously breaking scale contains very important information on low energy hadron dynamics. It was an active subjuct and studied continually during past fifty years. Experimentally, the effects of the strong interaction in process of $`e^+e^{}`$ annihilation is obvious to provides a large enhancement to production of pions in vector meson resonance region\[1-5\]. Theorectically, however, the problem was not studied by using a rigorous effective field theory(EFT) of QCD yet. Although at very low energy the chiral perturbative theory(ChPT) is a rigorou EFT of QCD, it in principle can not predict physics at vector meson resonance region. From viewpoint of quantum field theory, a rigorous theoretical study on $`e^+e^{}\pi ^+\pi ^{}`$ cross section and $`l=1,I=1`$ $`\pi \pi `$ phase shift provided by an EFT of QCD must satisfy the following requirements: 1) Some fundamental principles, such as symmetry and unitarity, must be satisfied in this EFT. 2) The experimental data of $`l=1,I=1`$ $`\pi \pi `$ scattering phase shift implies that width of $`\rho `$-meson $`\mathrm{\Gamma }_\rho `$ is transitional momentum-dependent, and must vanish at $`q^2=0`$(where $`q^2`$ denotes four-momentuma squre of off-shell $`\rho `$). The momentum-dependence of $`\mathrm{\Gamma }_\rho (q^2)`$ should be predicted by the EFT itself instead of being fitted by experiment. 3) This EFT must provide an effective method to evaluated error bar of this current calculation, i.e., the next order contribution should be able to be calculated. The purpose of this present paper is to provide a rigorous EFT study on pion form factor and $`I=1`$, P-wave $`\pi \pi `$ phase shift below 1GeV. In other words, all requirements mentioned above will be met in the study of this paper.
In some recent refrences, the authors have studied $`e^+e^{}\pi ^+\pi ^{}`$ cross section and $`l=1,I=1`$ phase shift at vector meson resonance region by using some very simple phenomenological models. These models are constructed in the intermediate energy region using some phenomenology considerations, such as vector meson dominant(VMD) and universal coupling. In principle, each of them can capture some leading order effects of low energy EFT of QCD and they are classfied by different symmetry realization for vector meson fields. However, so far, the low energy effective lagrangians including vector meson resonances are only up to $`O(p^4)`$ which are not enough for the physics at vector meson mass scale, and can not successfully evaluate very important one-loop effects of pseudoscalar mesons which corresponds to the next to leading order of $`N_c^1`$ expansion. Hence, these phenomenological models are not of rigorous EFT, and they can not provide any rigorous theoretical predictions on low energy hardon physics. This bad shortage can be overcome by using the EFT in ref., in which we constructed a consistent chiral constituent quark model(ChCQM) with lowest vector meson resonances and element Goldstone bosons. In this formalism we can capture all order information of chiral perturbative expansion and one-loop effects of pseudoscalar mesons.
In chiral limit, ChCQM is parameterized by the following chiral constituent quark lagrangian
$`_\chi `$ $`=`$ $`i\overline{q}(/+/\mathrm{\Gamma }+g__A/\mathrm{\Delta }\gamma _5i/V)qm\overline{q}q+{\displaystyle \frac{F^2}{16}}<_\mu U^\mu U^{}>+{\displaystyle \frac{1}{4}}m_0^2<V_\mu V^\mu >.`$ (1)
Here $`<\mathrm{}>`$ denotes trace in SU(3) flavour space, $`\overline{q}=(\overline{q}_u,\overline{q}_d,\overline{q}_s)`$ are constituent quark fields. $`V_\mu `$ denotes vector meson octet and singlet,
$$V_\mu (x)=\lambda 𝐕_\mu =\sqrt{2}\left(\begin{array}{ccc}\frac{\rho _\mu ^0}{\sqrt{2}}+\frac{\omega _\mu }{\sqrt{2}}& \rho _\mu ^+& K_\mu ^+\\ \rho _\mu ^{}& \frac{\rho _\mu ^0}{\sqrt{2}}+\frac{\omega _\mu }{\sqrt{2}}& K_\mu ^0\\ K_\mu ^{}& \overline{K}_\mu ^0& \varphi _\mu \end{array}\right).$$
(2)
The $`3\times 3`$ anti-Hermian matrices $`\mathrm{\Delta }_\mu `$ and $`\mathrm{\Gamma }_\mu `$ are defined as follows,
$`\mathrm{\Delta }_\mu `$ $`=`$ $`{\displaystyle \frac{1}{2}}\{\xi ^{}(_\mu ir_\mu )\xi \xi (_\mu il_\mu )\xi ^{}\},`$ (3)
$`\mathrm{\Gamma }_\mu `$ $`=`$ $`{\displaystyle \frac{1}{2}}\{\xi ^{}(_\mu ir_\mu )\xi +\xi (_\mu il_\mu )\xi ^{}\},`$ (4)
and covariant derivative are defined as follows
$`_\mu U`$ $`=`$ $`_\mu Uir_\mu U+iUl_\mu =2\xi \mathrm{\Delta }_\mu \xi ,`$ (5)
$`_\mu U^{}`$ $`=`$ $`_\mu U^{}il_\mu U^{}+iU^{}r_\mu =2\xi ^{}\mathrm{\Delta }_\mu \xi ^{},`$ (6)
where $`l_\mu =v_\mu +a_\mu `$ and $`r_\mu =v_\mu a_\mu `$ are linear combinations of external vector field $`v_\mu `$ and axial-vector field $`a_\mu `$, $`\xi `$ associates with non-linear realization of spontanoeusly broken global chiral symmetry introduced by Weinberg. This realization is obtained by specifying the action of global chiral group $`G=SU(3)_L\times SU(3)_R`$ on element $`\xi (\mathrm{\Phi })`$ of the coset space $`G/SU(3)__V`$:
$$\xi (\mathrm{\Phi })g_R\xi (\mathrm{\Phi })h^{}(\mathrm{\Phi })=h(\mathrm{\Phi })\xi (\mathrm{\Phi })g_L^{},g_L,g_RG,h(\mathrm{\Phi })H=SU(3)__V.$$
(7)
Explicit form of $`\xi (\mathrm{\Phi })`$ is usual taken
$$\xi (\mathrm{\Phi })=\mathrm{exp}\{i\lambda ^a\mathrm{\Phi }^a(x)/2\},U(\mathrm{\Phi })=\xi ^2(\mathrm{\Phi }),$$
(8)
where the Goldstone boson $`\mathrm{\Phi }^a`$ are treated as pseudoscalar meson octet. In ref. we have shown that the lagrangian( 1) is invariant under $`G_{\mathrm{global}}\times G_{\mathrm{local}}`$. The axial coupling constant $`g_A=0.75`$ is fitted by $`\beta `$-decay of neutron, and constituent quark mass $`m=480`$MeV is fitted by low energy limit of the model. It has been also illustrated that the value of $`g_A`$ has included effects of intermediate axial-vector meson resonances exchanges at low energy.
The EFT describing low energy meson interaction can be deduced via loop effects of constituent quarks. From viewpoint of symmetry, at leading order of vector mesons coupling to pseudoscalar mesons, the effective lagrangian is equivalent to WCCWZ lagrangian given by Brise. In terms of this EFT, we found that the chiral perturbative expansion converge slowly at vector meson energy scale. Thus the high order contributions of chiral perturbative expansion play important role at this energy scale. Phenomenologically, this model provides excellent theoretical predictions on $`\rho `$-physics and on $`\omega `$-physics. In this present paper, we focus our attention on vector-photon, vector-$`\pi \pi `$ and photon-$`\pi \pi `$ vertices. These relevant vertices have been calculated in ref. which including all order effects of the chiral perturbative expansion and one-loop contribution of pseudoscalar mesons. The “direct” photon-$`\pi \pi `$ coupling and vector-photon coupling vertices read,
$`_{\gamma \pi \pi }^c`$ $`=`$ $`{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}\overline{F}_\pi (q^2)A_\mu (q)[\pi ^+(x)^\mu \pi ^{}(x)^\mu \pi ^+(x)\pi ^{}(x)]},`$ (9)
$`_{\rho \gamma }^c`$ $`=`$ $`{\displaystyle \frac{1}{2}}e{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}b_{\rho \gamma }(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\rho ^{0\mu }(q)A^\nu (x)},`$ (10)
$`_{\omega \gamma }^c`$ $`=`$ $`{\displaystyle \frac{1}{6}}e{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}b_{\rho \gamma }(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q)A^\nu (x)},`$ (11)
where the super-srcipt “c” denotes these “complete” effective couplings which have contained one-loop effects of pseudoscalar mesons so that the “form factors”, $`\overline{F}_\pi (q^2)`$, $`b_{\rho \gamma }(q^2)`$ etc., are not real function. These “form factors” are given as follows,
$`\overline{F}_\pi (q^2)`$ $`=`$ $`1+{\displaystyle \frac{q^2b_\gamma (q^2)}{1+\mathrm{\Sigma }(q^2)}},b_{\rho \gamma }(q^2)={\displaystyle \frac{A(q^2)}{g(1+\zeta )}}f_\pi ^2b(q^2){\displaystyle \frac{\mathrm{\Sigma }_0(q^2)}{1+2\zeta }}[1+{\displaystyle \frac{q^2b_\gamma (q^2)}{1+\mathrm{\Sigma }(q^2)}}],`$ (12)
$`b_\gamma (q^2)`$ $`=`$ $`{\displaystyle \frac{gb(q^2)}{2(1+3\zeta )}}{\displaystyle \frac{1}{16\pi ^2f_\pi ^2}}\{\lambda +{\displaystyle _0^1}dxx(1x)\mathrm{ln}[(1{\displaystyle \frac{x(1x)p^2}{m__K^2}})({\displaystyle \frac{x(1x)p^2}{m__K^2}})^2]`$ (14)
$`{\displaystyle \frac{2}{3}}i\pi \theta (p^24m_\pi ^2)\}{\displaystyle \frac{C(q^2)\mathrm{\Sigma }_0(q^2)}{(1+11\zeta /3)}},`$
$`b(q^2)`$ $`=`$ $`{\displaystyle \frac{1}{gf_\pi ^2}}[A(q^2)+g_A^2B(q^2)],C(q^2)={\displaystyle \frac{A(q^2)+2g__A^2B(q^2)}{2f_\pi ^2}},`$ (15)
$`A(q^2)`$ $`=`$ $`g^2{\displaystyle \frac{N_c}{\pi ^2}}{\displaystyle _0^1}𝑑tt(1t)\mathrm{ln}(1{\displaystyle \frac{t(1t)q^2}{m^2}}),`$ (16)
$`B(q^2)`$ $`=`$ $`g^2+{\displaystyle \frac{N_c}{2\pi ^2}}{\displaystyle _0^1}dt_1t_1{\displaystyle _0^1}dt_2(1t_1t_2)[1+{\displaystyle \frac{m^2}{m^2t_1(1t_1)(1t_2)q^2}}`$ (18)
$`+\mathrm{ln}(1{\displaystyle \frac{t_1(1t_1)(1t_2)q^2}{m^2}})],`$
$`\mathrm{\Sigma }_0(q^2)`$ $`=`$ $`{\displaystyle \frac{2}{f_\pi ^2}}[2\mathrm{\Sigma }_\pi (q^2)\mathrm{\Sigma }_K(q^2)],\mathrm{\Sigma }(q^2)=[1+{\displaystyle \frac{q^2C(q^2)}{1+11\zeta /3}}]\mathrm{\Sigma }_0(q^2),`$ (19)
$`\mathrm{\Sigma }_K(q^2)`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi )^2}}\{\lambda (m__K^2{\displaystyle \frac{q^2}{6}})+{\displaystyle _0^1}𝑑t[m__K^2t(1t)q^2]\mathrm{ln}(1{\displaystyle \frac{t(1t)q^2}{m__K^2}})\},`$ (20)
$`\mathrm{\Sigma }_\pi (q^2)`$ $`=`$ $`{\displaystyle \frac{q^2}{(4\pi )^2}}\{{\displaystyle \frac{\lambda }{6}}+{\displaystyle _0^1}𝑑tt(1t)\mathrm{ln}{\displaystyle \frac{t(1t)q^2}{m__K^2}}{\displaystyle \frac{i}{6}}\pi \theta (q^24m_\pi ^2)\},`$ (21)
where $`f_\pi =185`$MeV is decay constant of pion, $`g`$ and $`\lambda `$ are constants which absorb the logarithmic divergence from constituent quark loops and the quadratic divergence from meson loops respectively,
$`g^2`$ $`=`$ $`{\displaystyle \frac{8}{3}}{\displaystyle \frac{N_c}{(4\pi )^{D/2}}}({\displaystyle \frac{\mu ^2}{m^2}})^{ϵ/2}\mathrm{\Gamma }(2{\displaystyle \frac{D}{2}}),`$ (22)
$`\lambda `$ $`=`$ $`({\displaystyle \frac{4\pi \mu ^2}{m__K^2}})^{ϵ/2}\mathrm{\Gamma }(2{\displaystyle \frac{D}{2}}),\zeta ={\displaystyle \frac{2\lambda }{(4\pi )^2}}{\displaystyle \frac{m__K^2}{f_\pi ^2}}.`$ (23)
In ref., $`\lambda =2/3`$ has been fitted by Zweig rule.
Traditionally, VMD assumes that all photon-hadron coupling is mediated by vector mesons. However, from an empirical or symmetrical point of view, one has a freedom, that a non-resonant background is alowed. For instance, in process of $`e^+e^{}\pi ^+\pi ^{}`$, since $`\pi ^+\pi ^{}`$ can consist of a vector-isovector system whose quantum numbers are same to $`\rho ^0`$, experiment or symmetry can not divide contribution of “direct” photon-$`\pi \pi `$ coupling from one from photon$`\rho ^0\pi ^+\pi ^{}`$. Therefore, in general, the traditional VMD is a strong assumption. From eq. (9), we can see that the “direct” photon-$`\pi \pi `$ coupling indeed exists in this EFT. The same problem is also questioned in isospin breaking decay $`\omega \pi ^+\pi ^{}`$, which is dominated by $`\rho ^0`$ exchange, but have a contirbution from “direct” $`\omega \pi ^+\pi ^{}`$ coupling yet.
The “complete” $`\rho \pi \pi `$ vertex reads
$$_{\rho \pi \pi }^c=\frac{d^4q}{(2\pi )^4}e^{iqx}g_{\rho \pi \pi }(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\rho ^{0\mu }(q)[\pi ^+(x)^\mu \pi ^{}(x)^\mu \pi ^+(x)\pi ^{}(x)],$$
(24)
with
$$g_{\rho \pi \pi }(q^2)=\frac{b(q^2)}{(1+2\zeta )(1+\mathrm{\Sigma }(q^2))}.$$
(25)
At leading order of vector meson coupling, the VMD vertex and $`\rho \pi \pi `$ vertex read respectively
$`_{\rho \gamma }^c`$ $`=`$ $`{\displaystyle \frac{1}{2}}eg{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}(q^2\delta _{\mu \nu }q_\mu q_\nu )\rho ^{0\mu }(q)A^\nu (x)},`$ (26)
$`_{\rho \pi \pi }^c`$ $`=`$ $`{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}\frac{1}{gf_\pi ^2}[g^2+g__A^2(\frac{N_c}{3\pi ^2}g^2)](q^2\delta _{\mu \nu }q_\mu q_\nu )\rho ^{0\mu }(q)[\pi ^+(x)^\mu \pi ^{}(x)^\mu \pi ^+(x)\pi ^{}(x)]}.`$ (27)
A gauge-like argument suggests that the $`\rho `$ couples to all hadrons with the same strength(universality). It is formulated by the first KSRF sum rule
$`g_{\rho \gamma }={\displaystyle \frac{1}{2}}f_\pi ^2g_{\rho \pi \pi }.`$ (28)
However, experimentally, the first KSRF sum rule is observed to be not quite exact. It can be naturally understood the universal coupling is a conclusion only working at leading order of vector meson coupling, and high order contribution will correct it. From eq. (26) we can see that the first KSRF sum rule is strictly satified when $`g=\pi ^1`$ for $`N_c=3`$. Thus $`g=\pi ^1`$ is a favorite choice. In addition, it has been shown in ref. how high order correction breaks the first KSRF sum rules. Besides of the parameters $`g_A`$, $`m`$, $`g`$ and $`\lambda `$, these are no other adjustable free parameters. So that this EFT will provide powerful prediction on low energy meson physics. For example, the theoretical prediction of on-shell decay width of $`\rho ^0e^+e^{}`$ is $`7.0`$MeV, which agree with experimental data, $`6.77\pm 0.32`$MeV, very well.
In this EFT, the $`\rho `$-resonance propagator(Breit-Wigner formula) can be naturally derived due to unitarity of the model instead of input,
$`\mathrm{\Delta }_{\mu \nu }^{(\rho )}(q^2)={\displaystyle \frac{i\delta _{\mu \nu }}{q^2\stackrel{~}{m}_\rho ^2+i\sqrt{q^2}\mathrm{\Gamma }_\rho (q^2)}},`$ (29)
where we have included only that part of the propagator which survives when coupled to conserved currents, $`\stackrel{~}{m_\rho }`$ is the (real valued) mass parameter and $`\mathrm{\Gamma }_\rho (q^2)`$ is the momentum-dependent width,
$`\mathrm{\Gamma }_\rho (q^2)={\displaystyle \frac{f_\pi ^2b^2(q^2)q^4}{2(1+2\zeta )^2}}\sqrt{q^2}\mathrm{Im}\{{\displaystyle \frac{\mathrm{\Sigma }_0(q^2)}{1+\mathrm{\Sigma }(q^2)}}\}={\displaystyle \frac{|g_{\rho \pi \pi }(q^2)|^2q^4}{48\pi }}\sqrt{q^2}(1{\displaystyle \frac{4m_\pi ^2}{q^2}})^{3/2}.`$ (30)
Numerically, the on-shell width $`\mathrm{\Gamma }_\rho =\mathrm{\Gamma }_\rho (q^2=m_\rho ^2)=146`$MeV, which agree with data very well.
Because the width in $`\rho `$-resonance (possessing a complex pole) propagator (29) is momentum-dependent, it must be addressed that the mass parameter $`\stackrel{~}{m}_\rho `$ is not the physical mass $`m_\rho =770`$MeV. Let us interpret this point briefly. Empirically, the physical mass of resonance is defined as position of pole(real value) in relevant scattering cross section, or theoretically, it should be defined as real part of complex pole possessed by resonance. It is well-known that the width of $`\rho `$-resonance is generated by pion loops. For a simple VMD model, the leading order of $`\rho \pi \pi `$ coupling is independent of $`q^2`$. Thus one has $`\mathrm{\Gamma }_\rho ^{(VMD)}(q^2)\sqrt{q^2}`$, and due to equation
$`q^2\stackrel{~}{m}_\rho ^2+i{\displaystyle \frac{\mathrm{\Gamma }_\rho }{m_\rho }}q^2=0,`$ (31)
we obtain the complex pole of $`\rho `$-resonance is $`q^2=m_\rho ^2(1iϵ+O(ϵ^2))`$ with $`ϵ=\mathrm{\Gamma }_\rho /m_\rho 0.19`$. The result yields $`\stackrel{~}{m}_\rho =m_\rho \sqrt{1+ϵ^2}=784`$MeV. In particular, in the EFT used by this present paper, $`\rho \pi \pi `$ coupling is proportional to $`q^2`$ at least. Hence one has $`\mathrm{\Gamma }_\rho (q^2)q^4\sqrt{q^2}`$ at least, and complex pole equation
$`q^2\stackrel{~}{m}_\rho ^2+i{\displaystyle \frac{\mathrm{\Gamma }_\rho }{m_\rho ^5}}q^6=0.`$ (32)
It yields $`\stackrel{~}{m}_\rho =m_\rho \sqrt{1+3ϵ^2}=810`$MeV, which poses a significant correction.
The above discussions imply that: 1) For resonance with large width, the mass parameter in its propagator is different from its physical mass. The correction is proportional to the ratio of resonant width to physical mass. 2) The mass in the orignal effective lagrangian only emerges as a parameter instead of a physical quantity measured directly in experiment. 3) The choice of mass parameter is relied on the choice of model. But the physical quantity must be independent of this choice.
Since in our result all hadronic couplings include all order information of the chiral perturbative expansion and one-loop effects of pseudoscalar mesons, the momentum-dependence of $`\mathrm{\Gamma }_\rho (q^2)`$ is very complicate. It is difficult to determine $`\stackrel{~}{m}_\rho `$ via the above method. Note that it is welcome that all vector meson resonances degenrate into a universal mass parameter $`m__V`$ at chiral limit and large $`N_c`$ limit. A reliable method is to determine $`m__V`$ via input mass of $`\omega `$-resonance (since $`\mathrm{\Gamma }_\omega m_\omega `$, $`\stackrel{~}{m}_\omega `$ is almost equal to $`m_\omega `$ and hereafter we do not distingusih them). Then $`\stackrel{~}{m}_\rho `$ can be obtained via dynamical calculation provide by this EFT. In general, the splitting between $`\stackrel{~}{m}_{\rho ^0}`$ and $`m_\omega `$ is caused by three sources: $`\rho ^0\omega `$ mixing, electromagnetic effects dur to VMD and one-loop effects of pseudoscalar mesons. Up to next to leading order $`N_c^1`$ expansion, the momentum-dependent $`\rho ^0\omega `$ mixing has been derived in ref.,
$$_{\omega \rho }^c=\frac{d^4q}{(2\pi )^4}e^{iqx}\mathrm{\Theta }_{\omega \rho }(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q)\rho ^{0\nu }(x),$$
(33)
where
$`\mathrm{\Theta }_{\omega \rho }(q^2)`$ $`=`$ $`{\displaystyle \frac{N_c}{6\pi ^2}}{\displaystyle \frac{m_um_d}{m}}\{g^2h_0(q^2)(1{\displaystyle \frac{4}{3}}\zeta )+q^2b(q^2)s(q^2)[\mathrm{\Sigma }_K(q^2){\displaystyle \frac{\mathrm{\Sigma }_\pi (q^2)}{1+\mathrm{\Xi }(q^2)}}]\}`$ (35)
$`+{\displaystyle \frac{\alpha _{\mathrm{e}.\mathrm{m}.}\pi }{3}}b_{\rho \gamma }^2(q^2)+O((a_1(m_um_d)+a_2\alpha _{\mathrm{e}.\mathrm{m}.})^2).`$
The function $`b(q^2)`$, $`b_{\rho \gamma }(q^2)`$, $`\mathrm{\Sigma }_K(q^2)`$ and $`\mathrm{\Sigma }_\pi (q^2)`$ are given in eq.(12), and $`s(q^2)`$, $`h_0(q^2)`$ and $`\mathrm{\Xi }(q^2)`$ are of follows
$`s(q^2)`$ $`=`$ $`{\displaystyle \frac{4}{gf_\pi ^2}}[h_0(q^2)+{\displaystyle \frac{3}{4}}g_A^2(h_1(q^2){\displaystyle \frac{h_2(q^2)}{2}})],`$ (36)
$`h_0(q^2)`$ $`=`$ $`{\displaystyle _0^1}𝑑t{\displaystyle \frac{6t(1t)}{1t(1t)q^2/m^2}},`$ (37)
$`h_1(q^2)`$ $`=`$ $`{\displaystyle _0^1}𝑑t_1t_1^2{\displaystyle _0^1}𝑑t_2(1t_2){\displaystyle \frac{32t_1^2t_2(1+2t_1)(1t_2)q^2/m^2}{[1t_1^2t_2(1t_2)q^2/m^2]^2}},`$ (38)
$`h_2(q^2)`$ $`=`$ $`{\displaystyle _0^1}𝑑t_1t_1^2{\displaystyle _0^1}𝑑t_2(1t_2){\displaystyle \frac{4(1t_1)[34t_1^2t_2(1t_2)q^2/m^2]}{[1t_1^2t_2(1t_2)q^2/m^2]^2}},`$ (39)
$`\mathrm{\Xi }(q^2)`$ $`=`$ $`4f_\pi ^2(1+{\displaystyle \frac{q^2N_c}{4\pi ^2f_\pi ^2}})\mathrm{\Sigma }_\pi (q^2).`$ (40)
In addition, it should be also noticed that the $`\pi `$-loop correction to $`m_\omega `$ is suppressed by isospin conservation. Thus one-loop correction to $`m_\omega `$ is dominated by $`K`$-meson. At tree level, the $`\omega KK`$ coupling reads
$`_{\omega KK}`$ $`=`$ $`{\displaystyle \frac{i}{2}}{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}b(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q)}`$ (41)
$`\times `$ $`\{[K^+(x)^\nu K^{}(x)^\nu K^+(x)K^{}(x)]+[K^0(x)^\nu \overline{K}^0(x)^\nu K^0(x)\overline{K}^0(x)]\}.`$ (42)
Then the physical mass of $`\omega `$-meson are
$`m_\omega ^2`$ $`=`$ $`m__V^2+\mathrm{Re}\{{\displaystyle \frac{q^4\mathrm{\Theta }_{\omega \rho }(q^2)}{q^2\stackrel{~}{m}_\rho ^2+i\sqrt{q^2}\mathrm{\Gamma }_\rho (q^2)}}+{\displaystyle \frac{\pi \alpha _{\mathrm{e}.\mathrm{m}.}}{9}}q^2b_{\rho \gamma }^2(q^2)\}|_{q^2=m_\omega ^2}`$ (44)
$`{\displaystyle \frac{q^4b^2(q^2)\mathrm{\Sigma }_K(q^2)}{(1+2\zeta )^2\{1+2f_\pi ^2[1+\frac{q^2C(q^2)}{1+11\zeta /3}]\mathrm{\Sigma }_K(q^2)\}}}|{\displaystyle \genfrac{}{}{0pt}{}{}{q^2=m_\omega ^2}}.`$
Input $`m_\omega =782`$MeV, one has $`m__V=785.8`$MeV. Similarly, the mass parameter of $`\rho `$-meson are
$`\stackrel{~}{m}_\rho ^2`$ $`=`$ $`m__V^2+\mathrm{Re}\{{\displaystyle \frac{q^4\mathrm{\Theta }_{\omega \rho }(q^2)}{q^2m_\omega ^2+im_\omega \mathrm{\Gamma }_\omega }}+\pi \alpha _{\mathrm{e}.\mathrm{m}.}q^2b_{\rho \gamma }^2(q^2)+{\displaystyle \frac{f_\pi ^2b^2(q^2)q^6\mathrm{\Sigma }_0(q^2)}{2(1+2\zeta )^2(1+\mathrm{\Sigma }(q^2))}}\}|_{q^2=m_\rho ^2}.`$ (45)
Using the above value of $`m__V`$, we have $`\stackrel{~}{m}_\rho =803.1`$MeV which is indeed significantly different from the physical mass $`m_\rho =770`$MeV. Success of this prediction will be checked in the following by localizing the position of pole in cross section of $`e^+e^{}\pi ^+\pi ^{}`$. Furthermore, the detail calculation shows that, the $`\rho ^0\omega `$ only makes $`\stackrel{~}{m}_\rho `$ shift $`0.25`$MeV, the VMD effects and one-loop effects of pseudoscalar mesons make $`\stackrel{~}{m}_\rho `$ shift $`+3.45`$MeV and $`+14.1`$MeV respectively.
For working out full shape of $`e^+e^{}\pi ^+\pi ^{}`$ cross section, the $`\omega \pi \pi `$ coupling is needed. It has been derived in ref., in which not only the coupling via $`\rho `$ exchange, but also the “direct” coupling are included,
$`_{\omega \pi \pi }^c=i{\displaystyle \frac{d^4q}{(2\pi )^4}e^{iqx}g_{\omega \pi \pi }(q^2)(q^2\delta _{\mu \nu }q_\mu q_\nu )\omega ^\mu (q)[\pi ^+(x)^\nu \pi ^{}(x)^\nu \pi ^+(x)\pi ^{}(x)]},`$ (46)
where
$`g_{\omega \pi \pi }(q^2)={\displaystyle \frac{q^2\mathrm{\Theta }_{\omega \rho }g_{\rho \pi \pi }(q^2)}{q^2\stackrel{~}{m}_\rho ^2+i\sqrt{q^2}\mathrm{\Gamma }_\rho (q^2)}}g_{\omega \pi \pi }^{(0)}(q^2),`$ (47)
with “direct” coupling strength
$`g_{\omega \pi \pi }^{(0)}(q^2)`$ $`=`$ $`{\displaystyle \frac{N_c}{12\pi ^2}}{\displaystyle \frac{m_um_d}{m}}\{s(q^2)({\displaystyle \frac{1}{1+\mathrm{\Xi }(q^2)}}{\displaystyle \frac{10}{3}}\zeta )`$ (49)
$`6f_\pi ^2\mathrm{\Sigma }_K(q^2)[8m^2f_\pi ^2b(q^2){\displaystyle \frac{s(q^2)}{3}}(1+{\displaystyle \frac{q^2N_c}{4\pi ^2f_\pi ^2}})]\}+{\displaystyle \frac{2\alpha _{\mathrm{e}.\mathrm{m}.}\pi }{3}}\overline{F}_\pi (q^2)b_{\rho \gamma }(q^2).`$
Eqs. (9), (24) and (46) lead to the electromagnetic form factor of pion as follow
$`F_\pi (q^2)=1+{\displaystyle \frac{q^2b_\gamma (q^2)}{1+\mathrm{\Sigma }(q^2)}}{\displaystyle \frac{q^4b_{\rho \gamma }(q^2)g_{\rho \pi \pi }(q^2)}{2(q^2\stackrel{~}{m}_\rho ^2+i\sqrt{q^2}\mathrm{\Gamma }_\rho (q^2))}}{\displaystyle \frac{q^4b_{\rho \gamma }(q^2)g_{\omega \pi \pi }(q^2)}{6(q^2m_\omega ^2+i\sqrt{q^2}\mathrm{\Gamma }_\omega )}}.`$ (50)
Here due to narrow width of $`\omega `$, we ignore the momentum-dependence of $`\mathrm{\Gamma }_\omega `$. In this form factor, we can see that the contributions of resonance exchange accompany $`q^4`$ factor. Due to this reason, some authors declared that the pion form factor in WCCWZ EFT exhibits an unphysical high energy behaviour($`\mu >m_\rho `$). However, this conclusion is wrong. It is caused by their wrong result for momentum-dependence of $`\mathrm{\Gamma }_\rho (q^2)`$ which is fitted by experimental instead of by dynamical prediction. In fact, since $`\sqrt{q^2}\mathrm{\Gamma }_\rho (q^2)`$ is proportional to $`q^6`$ at least, we do not need to worry that the form factor has a bad high energy behaviour. We can also see that there is a moment-dependent non-resonant contribution. It together with the contribution of resonance exchange determined the high energy behaviour of the factor. The cross-section for $`e^+e^{}\pi ^+\pi ^{}`$ is given by(neglecting the electron mass)
$`\sigma ={\displaystyle \frac{\pi \alpha _{\mathrm{e}.\mathrm{m}.}^2}{3}}{\displaystyle \frac{(q^24m_\pi ^2)^{3/2}}{(q^2)^{5/2}}}|F_\pi (q^2)|^2.`$ (51)
From defination of function $`A(q^2)`$ and $`B(q^2)`$ in eq. (12) we can see this EFT is unitary only for $`q^2<4m^2`$. Thus the effective prediction should be below $`m_{\pi \pi }<2m=960`$MeV. The result is shown in fig. 1. We can see the prediction agree with data well. Especially, the theoretical prediction in vector meson energy region agree with data excellently. Although the mass parameter $`\stackrel{~}{m}_\rho =803.1`$MeV in $`\rho `$ propagator is larger than physical mass, the position of pole is localized in $`\sqrt{q^2}=772`$MeV which is just the physical mass of $`\rho `$. It strongly supports our above discussion and dynamical calculation. It also implies that we must carefully distinguish the physical mass difference of $`\rho ^0`$ and $`\omega `$ from the diffrence of mass parameter in effective lagrangian.
Let us give some further remarks on pion form factor (50). From eqs. (12), (25) and (47) we can see that, in eq. (50), $`b_\gamma (q^2)`$, $`b_{\rho \gamma }(q^2)`$, etc., are all complex function instead of real function. It is caused by one-loop effects of pions. Thus the expression (50) can be rewritten as follow
$`F_\pi (q^2)=1+q^2a_1(q^2)e^{i\varphi _1(q^2)}{\displaystyle \frac{q^4a_2(q^2)e^{i\varphi _2(q^2)}}{2(q^2\stackrel{~}{m}_\rho ^2+i\sqrt{q^2}\mathrm{\Gamma }_\rho (q^2))}}{\displaystyle \frac{q^4a_3(q^2)e^{i\varphi _3(q^2)}}{6(q^2m_\omega ^2+i\sqrt{q^2}\mathrm{\Gamma }_\omega )}}.`$ (52)
Here $`a_i(q^2)(i=1,2,3)`$ are three real function and $`\varphi _i(q^2)(i=1,2,3)`$ are three momentum-dependent phases. In particular, $`\varphi _3(q^2=m_\omega ^2)`$, so called Orsay phase, has been extracted from data as $`100125`$ degrees. Our theorectical prediction is $`\varphi _3(q^2=m_\omega ^2)=116.5`$ degrees. However, so far, the phases $`\varphi _1(q^2)`$ and $`\varphi _2(q^2)`$ are not reported in any literatures. These momentum-dependent phases indicate that the dynamics including loop effects of pseudoscalar mesons is different from one only in tree level. In fig.2, we given theorectical curves of $`\varphi _i(q^2)`$. They are indeed nontrivial.
Obviously, $`F_\pi (q^2)`$ is an analytic function in the complex $`q^2`$ plane, with a branch cut along the real axis beginning at the two-pion threshold, $`q^2=4m_\pi ^2`$. Time-reversal invariance and the unitarity of the $`S`$-matrix requires that the phase of the form factor be that of $`l=1,I=1`$ $`\pi \pi `$ scattering. This last emerges as $`\pi \pi `$ scattering in the relevant channel is very nearly elastic from threshold through $`q^2(m_\pi +m_\omega )^2`$. In this region of $`q^2`$, then, the form factor is related to the $`l=1,I=1`$ $`\pi \pi `$ phase shift, $`\delta _1^1`$, via
$`F_\pi (q^2)=e^{2i\delta _1^1(q^2)}F_\pi ^{}(q^2),`$ (53)
so that
$`\mathrm{tan}\delta _1^1(q^2)={\displaystyle \frac{\mathrm{Im}F_\pi (q^2)}{\mathrm{Re}F_\pi (q^2)}}.`$ (54)
The above is a special case of what is sometimes called the Fermi-Watson-Aidzu phase theorem. In fig. 3 and fig. 4 we plot theoretical curves of the $`l=1,I=1`$ $`\pi \pi `$ phase shift $`\delta _1^1`$ versus $`m_{\pi \pi }`$ and of $`\mathrm{sin}\delta _1^1/p_\pi ^3`$ versus $`m_{\pi \pi }`$ (where $`p_\pi =\frac{1}{2}\sqrt{q^24m_\pi ^2}`$) respectively. We omit the $`\omega `$ contribution from our plots of the phase of $`F_\pi (q^2)`$ for comparing with time-like region pion form factor data . We have also assumed that $`\delta _1^1`$ is purely elastic in the regime shown, i.e., the loop effects of $`\omega \pi `$ are omitted. The curve predicts $`\delta _1^190^{}`$ as $`\sqrt{q^2}774\mathrm{M}\mathrm{e}\mathrm{V}m_\rho `$, and $`\delta _1^1>100^{}`$ for $`\sqrt{q^2}>787`$MeV. These results agree with data very well.
Finally we discuss the near threshold behaviour of the form factor. 1) The chiral perturbative theory predicts the form factor at threshold to be $`[F_\pi (4m_\pi ^2)]_{\mathrm{ChPT}}=1.17\pm 0.01`$, and ours, $`[F_\pi (4m_\pi ^2)]=1.154`$, is close to the ChPT result. 2) The electromagnetic radius of charged pion has been determined to be $`\sqrt{<r^2>_\pi }=0.657\pm 0.027`$fm, whereas the theorectical prediction in this present paper is $`\sqrt{<r^2>_\pi }=0.635`$ fm. 3) The Froggatt-Petersen phase shift function $`\mathrm{sin}\delta _1^1/p_\pi ^3`$ is connected with the vector-isovector $`\pi \pi `$ scattering length $`a_1^1`$ through
$`a_1^1=\underset{q^24m_\pi ^2}{lim}{\displaystyle \frac{\mathrm{sin}\delta _1^1}{p_\pi ^3}}.`$ (55)
Our theoretical prediction is $`a_1^1=0.037`$ in unit of $`m_\pi ^3`$. This value is very close to experimental results from $`K_{e4}`$ data using a Roy equation fit ($`a_1^1=0.038\pm 0.002`$) and ChPT prediction $`a_1^1=0.037\pm 0.01`$ at the two loop order (at $`O(p^4)`$).
In summary, a rigorous, unitary EFT method has been applied to study pion form factor and $`l=1,I=1`$ $`\pi \pi `$ phase shift. The theoretical predictions include all order informations of the chiral perturbative expansion and one-loop effects of pseudoscalar mesons. The Breit-Wigner formula for resonant propagators is derived by the EFT itself instead of an input. The momentum-dependence of $`\mathrm{\Gamma }_\rho (q^2)`$ is predicted by the dynamics. It also has been revealed that the mass parameter in resonant propagators should be different from its physical mass due to momentum-dependent width. This point is confirmed by both of dynamical calculation and phenomenological fit. It also tells us how to understand the mass splitting between $`\rho ^0`$ and $`\omega `$: Although the dynamical calculations show that the mass parameter of $`\rho `$ in effective lagrangian is even larger than one of $`\omega `$, the position of pole localized in real aixs give their right physical mass splitting. The contribution to this mass splitting from $`\rho ^0\omega `$ mixing is very small, the dominant contribution is from one-loops effects of pseudoscalar mesons. The EFT mechanics on $`\rho ^0\omega `$ mass splitting revealed in the present paper is rather subtle. And it is another evident to confirm again that the EFT of QCD proposed in ref. is sound. Actually, to the best of our knowledges, this is the first time to get $`\rho ^0\omega `$ mass splitting through a well-defined quantum field theory calculation.
Due to one-loop effects of pions, the photon-$`\pi \pi `$, photon-vector and vector-$`\pi \pi `$ coupling are all with a nontrivial phase shift instead of purely real in some simple models. In a series of recent papers, we have revealed that the one-loop effects of pseudoscalar mesons play a very important role in low energy hadronic physics. Theoretically, it keeps unitarity of the $`S`$-matrix, and numerically, its contributions are about $`30\%`$. A well-defined EFT must be able to evaluate the high order contributions of the chiral perturbative expansion and $`N_c^1`$ expansion. It is an important criterion to judge a model as a rigorou EFT or a phenomenological model only.
In our study, no parameters need to be fitted by the data of pion form factor and $`l=1,I=1`$ $`\pi \pi `$ phase shift. Thus our results are rigorous theoretical predictions and agree with data very well.
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# 2 and 3-dimensional Hamiltonians with Shape Invariance Symmetry
## I INTRODUCTION
Exactly solvable quantum Hamiltonians (ESQH) have always attracted a lot of interest in theoretical physics and mathematical physics. Hence, construction of exactly solvable models is of great interest . Familiar solvable potentials (particularly one-dimensional one) have the property of shape invariance, where this property has played an important role in calclulation of their determinant by Heat Kernel method . For these potentials, eigenvalues and eigenvectors can be derived using the well known methods of supersymmetric quantum mechanics together with shape invariant factorization. The majority of potentials have also been shown to possess a Lie algebraic symmetry and hence are also solvable by group theoretical techniques. Actually one can establish a connection between ESQH with shape invariance symmetry and ESQH with Lie algebraic symmetry and can show that they are indeed equivalent . One of the authors has introduced some 2 and 3-dimensional shape invariant Hamiltonians . In theses article they have shown that the shape invariance symmetry of these models is due to the existence of some Lie algebraic symmetry. Hence, in this article we construct new 2 and 3-dimensional EQSH with shape invariance symmetry, where $`su(2)`$ and Heisenberg algebra are responsible for the existence of shape invariance symmetry in them.
This paper is organized as follows: In section II, using the left and right invariant vector fields of $`su(2)`$ Lie algebra we first construct its Casimir operator. Then via Fourier transformation over one of the coordinates we construct 2-dimensional Hamiltonian $`H_q(\theta ,\psi )`$ which possess shape invariance symmetry. Using this symmetry we obtain its eigenspectrum analytically. In section III, starting with Hamiltonian of 4-oscillator and Fourier transforming over one of the coordinates, we obtain 3-dimensional Hamiltonian corresponding to motion of a charged particle in presence of an electric field. We show that this 3-dimensional Hamiltonian possess a shape invariance symmetry and using this symmetry we obtain its eigenspectrum. What is so important in both models is that both Hamiltonians factorize shape invariantly into a product of second order differential operators. These second order operators themselves consist of the product of first order differential operators. The paper ends with a brief conclusion.
## II 2-dimensional Hamiltonian obtained from <br>$`SU(2)`$ manifold
### II.1 Left and Right invariant vector field of $`SU(2)`$
Considering the following parametrization of $`su(2)`$ group manifold
$$\mathrm{\Lambda }=\mathrm{exp}(i\stackrel{}{\sigma }.\stackrel{}{n}\psi )=A\left(\begin{array}{cc}\mathrm{exp}(i\psi )& 0\\ 0& \mathrm{exp}(i\psi )\end{array}\right)A^1$$
$$=\left(\begin{array}{cc}\mathrm{cos}(\psi )i\mathrm{cos}(\theta )\mathrm{sin}(\psi )& i\mathrm{sin}(\theta )\mathrm{sin}(\psi )\mathrm{exp}(i\varphi )\\ i\mathrm{sin}(\theta )\mathrm{sin}(\psi )\mathrm{exp}(i\varphi )& \mathrm{cos}(\psi )+i\mathrm{cos}(\theta )\mathrm{sin}(\psi )\end{array}\right),$$
(2.1)
where $`\sigma _i`$, i= 1, 2 and 3 are Pauli matrices and $`\stackrel{}{n}`$ is a unit vector defined as:
$$\stackrel{}{n}=\mathrm{sin}(\theta )\mathrm{cos}(\varphi )\stackrel{}{i}+\mathrm{sin}(\theta )\mathrm{sin}(\varphi )\stackrel{}{j}+\mathrm{cos}(\theta )\stackrel{}{k},$$
and matrix $`A`$ corresponds to the coherent state representation of $`su(2)`$ defined as :
$$A=\left(\begin{array}{cc}1& \tau \\ 0& 1\end{array}\right)\left(\begin{array}{cc}\mathrm{exp}(\frac{\beta }{2})& 0\\ 0& \mathrm{exp}(\frac{\beta }{2})\end{array}\right)\left(\begin{array}{cc}1& 0\\ \tau ^{}& 1\end{array}\right),$$
with $`\tau =\mathrm{tan}(\frac{\theta }{2})\mathrm{exp}(i\varphi )`$ and $`\beta =\mathrm{ln}(1+\tau \tau ^{})`$.
In order to obtain the left and right invariant vector field $`su(2)`$ manifold with the above parametrization, it is convenient first to calculate its left and right invariant one form defined as $`\mathrm{\Lambda }^1d\mathrm{\Lambda }`$ and $`d\mathrm{\Lambda }\mathrm{\Lambda }^1`$ respectively .
As an example, let us write left invariant one form
$$\mathrm{\Lambda }^1d\mathrm{\Lambda }=e_\alpha ^ad\xi ^\alpha \sigma _a,$$
where $`e_\alpha ^a`$ are 3-beins and $`\xi ^\alpha =(\theta ,\varphi ,\psi )`$ are coordinates of $`su(2)`$-manifold. Defining the inverse of 3-bein $`e_a^\alpha =g^{\alpha \beta }\delta _{ab}e_\beta ^b`$ with $`g^{\alpha \beta }`$ as inverse of metric $`g_{\alpha \beta }`$:
$$g_{\alpha ,\beta }=\left(\begin{array}{ccc}1& 0& 0\\ 0& \mathrm{sin}^2(\psi )& 0\\ 0& 0& \mathrm{sin}^2(\psi )\mathrm{sin}^2(\theta )\end{array}\right),$$
then the left invariant vector field is defined as:
$$L_a=e_a^\alpha \frac{}{\xi ^\alpha }.$$
Using the above prescription we obtain the following expression for left and right invariant vector field of $`su(2)`$, respectively,
$$L_+=\frac{i}{2}e^{i\varphi }\left[\mathrm{sin}(\theta )_\psi +(i+\mathrm{cos}(\theta )\mathrm{cot}(\psi ))_\theta +(\mathrm{cot}(\theta )+i\frac{\mathrm{cot}(\psi )}{\mathrm{sin}(\theta )})_\varphi \right],$$
(2.2)
$$L_{}=\frac{i}{2}e^{i\varphi }\left[\mathrm{sin}(\theta )_\psi +(i+\mathrm{cos}(\theta )\mathrm{cot}(\psi ))_\theta +(\mathrm{cot}(\theta )i\frac{\mathrm{cot}(\psi )}{\mathrm{sin}(\theta )})_\varphi \right],$$
(2.3)
$$L_3=\frac{i}{2}\left(\mathrm{cos}(\theta )_\psi +\mathrm{sin}(\theta )\mathrm{cot}(\psi )_\theta _\varphi \right),$$
(2.4)
$$R_+=\frac{i}{2}e^{i\varphi }\left[\mathrm{sin}(\theta )_\psi +(i+\mathrm{cos}(\theta )\mathrm{cot}(\psi ))_\theta +(\mathrm{cot}(\theta )+i\frac{\mathrm{cot}(\psi )}{\mathrm{sin}(\theta )})_\varphi \right],$$
(2.5)
$$R_{}=\frac{i}{2}e^{i\varphi }\left[\mathrm{sin}(\theta )_\psi +(i+\mathrm{cos}(\theta )\mathrm{cot}(\psi ))_\theta +(\mathrm{cot}(\theta )i\frac{\mathrm{cot}(\psi )}{\mathrm{sin}(\theta )})_\varphi \right],$$
(2.6)
$$R_3=\frac{i}{2}\left(\mathrm{cos}(\theta )_\psi +\mathrm{sin}(\theta )\mathrm{cot}(\psi )_\theta +_\varphi \right),$$
(2.7)
where $`L_\pm =L_1\pm iL_2`$ and $`R_\pm =R_1\pm iR_2`$. It is straightforward to show that the left and right invariant vector field fulfill the following $`su(2)`$ Lie algebra:
$$[L_+,L_{}]=2L_3,[L_3,L_\pm ]=\pm L_\pm ,$$
(2.8)
$$[R_+,R_{}]=2R_3,[R_3,R_\pm ]=R_\pm ,$$
(2.9)
also, the left and right invariant generators commute with each other
$$[\stackrel{}{L},\stackrel{}{R}]=0.$$
(2.10)
Considering the Casimir operators of $`su(2)`$ defined as:
$$L^2=\frac{1}{2}(L_+L_{}+L_{}L_+)+L_3^2,$$
and ignoring the scale 1/4, we obtain
$$L^2=\frac{1}{\mathrm{sin}^2(\psi )}_\psi \mathrm{sin}^2(\psi )_\psi \frac{1}{\mathrm{sin}^2(\psi )}\left(\frac{1}{\mathrm{sin}(\theta )}_\theta \mathrm{sin}(\theta )_\theta +\frac{1}{\mathrm{sin}^2(\theta )}_\varphi ^2\right).$$
(2.11)
In obtaining the above formula we have used the left invariant generators. It is straightforward to show that we can obtain the same result with right invariant generators too, that means the Casimir operator of left and right operator are the same.
### II.2 $`H_q(\theta ,\psi )`$ Hamiltonian
Here through dimensional reduction we show that the above Casimir operator reduces to a Hamiltonian of motion of a charged particle in the presence of electric field. Hence, first we make one-dimensional reduction ( eliminate the coordinate $`\varphi `$ ) through the usual Fourier transformation defined as
$$\stackrel{~}{f}(q)=\frac{1}{\sqrt{2\pi }}_0^{2\pi }f(\varphi )\mathrm{exp}(i\varphi q)𝑑\varphi ,$$
(2.12)
over an arbitrary function $`f(\varphi )`$. Obviously the Casimir operator (2.11) reduces to the following operator
$$L_q^2(\theta ,\psi )=\frac{1}{\mathrm{sin}^2(\psi )}_\psi \mathrm{sin}^2(\psi )_\psi \frac{1}{\mathrm{sin}^2(\psi )}\left(\frac{1}{\mathrm{sin}(\theta )}_\theta \mathrm{sin}(\theta )_\theta \frac{q^2}{\mathrm{sin}^2(\theta )}\right),$$
(2.13)
in the Hilbert space of Fourier transformed wavefunctions. In general the non- relativistic Hamiltonian of a charged particle over a 2-dimensional manifold with metric $`g_{\mu \nu }`$ in the presence of magneto static field $`\stackrel{}{B}`$ with vector potential $`\stackrel{}{A}`$ and electro static field $`\stackrel{}{E}`$ with scalar potential $`V`$ can be written as
$$H=\frac{1}{\sqrt{g}}(_\mu iA_\mu )(\sqrt{g}g^{\mu \nu }(_\nu iA_\nu ))+V,$$
(2.14)
where $`g`$ is the determinant of metric $`g_{\mu \nu }`$. After similarity transformation of the Casimir operator (2.13) defined as:
$$\stackrel{~}{L}_q^2(\theta ,\psi )=\mathrm{sin}^{\frac{1}{2}}(\psi )L_q^2(\theta ,\psi )\mathrm{sin}^{\frac{1}{2}}(\psi ),$$
we have
$$\stackrel{~}{L}_q^2(\theta ,\psi )=\frac{1}{\mathrm{sin}(\psi )}_\psi \mathrm{sin}(\psi )_\psi \frac{1}{\mathrm{sin}^2(\psi )}\left(\frac{1}{\mathrm{sin}(\theta )}_\theta \mathrm{sin}(\theta )_\theta \frac{q^2}{\mathrm{sin}^2(\theta )}\right)+\frac{1}{4}\mathrm{cot}^2(\psi )\frac{1}{2}.$$
(2.15)
Comparing the operator (2.15) with the Hamiltonian (2.14) we obtain
$$g_{\psi \psi }=1,g_{\theta \theta }=\mathrm{sin}^2(\psi ),g_{\psi \theta }=g_{\theta \psi }=0$$
and
$$A_\psi =0,A_\theta =\frac{i}{2}\mathrm{cot}(\theta )=d(\frac{i}{2}\mathrm{ln}(\mathrm{sin}(\theta )).$$
(2.16)
It is trivial to see that the vector potential given in (2.16) corresponds to the pure $`u(1)`$ gauge field, hence it can be eliminate through the gauge transform $`AA_\mu +_\mu \chi `$ with gauge function $`\chi =\frac{i}{2}\mathrm{ln}(\mathrm{sin}(\theta ))`$. After the above gauge transformation the $`su(2)`$-Casimir Hamiltonian reduces to
$$H_q(\theta ,\psi )e^\chi \stackrel{~}{L}_q^2(\theta ,\psi )e^\chi =\frac{1}{\mathrm{sin}(\psi )}_\psi \mathrm{sin}(\psi )_\psi \frac{1}{\mathrm{sin}^2(\psi )}_\theta ^2+\frac{q^2\frac{1}{4}}{\mathrm{sin}^2(\psi )\mathrm{sin}^2(\theta )}\frac{3}{4},$$
(2.17)
which can be interpreted as a non-relativistic Hamiltonian of a point particle over
2-dimensional sphere with metric
$$g_{\mu ,\nu }=\left(\begin{array}{cc}1& 0\\ 0& \mathrm{sin}^2(\psi )\end{array}\right)$$
in the presence of electric field with scalar potential
$$V=\frac{q^2\frac{1}{4}}{\mathrm{sin}^2(\psi )\mathrm{sin}^2(\theta )}\frac{3}{4}.$$
Similarly, the left and right invariant vector fields given in (2.2)-(2.7) take the following form after the above given operations, namely, dimensional reduction, similarity transformation and gauge transformation:
$$\stackrel{~}{L}_+^{}(q)=\stackrel{~}{L}_+(q)+g(\theta ,\psi ,q),\stackrel{~}{R}_+^{}(q)=\stackrel{~}{R}_+(q)g^{}(\theta ,\psi ,q),$$
$$\stackrel{~}{L}_{}^{}(q)=\stackrel{~}{L}_{}(q)g^{}(\theta ,\psi ,q),\stackrel{~}{R}_{}^{}(q)=\stackrel{~}{R}_{}(q)+g(\theta ,\psi ,q),$$
$$\stackrel{~}{L}_3^{}(q)=\stackrel{~}{L}_3(q),\stackrel{~}{R}_3^{}(q)=\stackrel{~}{R}_3(q),$$
where
$$\stackrel{~}{L}_+(q)=\frac{i}{2}\left(\mathrm{sin}(\theta )_\psi +(i+\mathrm{cos}(\theta )\mathrm{cot}(\psi ))_\theta +i(q1)(\mathrm{cot}(\theta )+i\frac{\mathrm{cot}(\psi )}{\mathrm{sin}(\theta )})\right)e^\frac{}{q},$$
(2.18)
$$\stackrel{~}{L}_{}(q)=\frac{i}{2}\left(\mathrm{sin}(\theta )_\psi +(i+\mathrm{cos}(\theta )\mathrm{cot}(\psi ))_\theta +i(q+1)(\mathrm{cot}(\theta )i\frac{\mathrm{cot}(\psi )}{\mathrm{sin}(\theta )})\right)e^\frac{}{q},$$
(2.19)
$$\stackrel{~}{L}_3(q)=\frac{i}{2}\left(\mathrm{cos}(\theta )_\psi +\mathrm{sin}(\theta )\mathrm{cot}(\psi )_\theta iq\right),$$
(2.20)
$$\stackrel{~}{R}_+(q)=\frac{i}{2}\left(\mathrm{sin}(\theta )_\psi +(i+\mathrm{cos}(\theta )\mathrm{cot}(\psi ))_\theta +i(q1)(\mathrm{cot}(\theta )+i\frac{\mathrm{cot}(\psi )}{\mathrm{sin}(\theta )})\right)e^\frac{}{q},$$
(2.21)
$$\stackrel{~}{R}_{}(q)=\frac{i}{2}\left(\mathrm{sin}(\theta )_\psi +(i+\mathrm{cos}(\theta )\mathrm{cot}(\psi ))_\theta +i(q+1)(\mathrm{cot}(\theta )i\frac{\mathrm{cot}(\psi )}{\mathrm{sin}(\theta )})\right)e^\frac{}{q},$$
(2.22)
$$\stackrel{~}{R}_3(q)=\frac{i}{2}\left(\mathrm{cos}(\theta )_\psi +\mathrm{sin}(\theta )\mathrm{cot}(\psi )_\theta +iq\right)$$
(2.23)
with $`g(\theta ,\psi ,q)`$ is:
$$g(\theta ,\psi ,q)=\frac{1}{4}\left(\mathrm{cot}(\theta )i\frac{\mathrm{cot}(\psi )}{\mathrm{sin}(\theta )}\right)e^\frac{}{q},$$
where * means the usual complex conjugation. With some calculation one can show that the above algebra, that is, the commutation relations is unchanged under the above mentioned transformation and the Hamiltonian $`H_q(\theta ,\psi )`$ can be written in terms of generators (2.18)-(2.20) in the following form
$$H_q(\theta ,\psi )=\frac{1}{2}\left(\stackrel{~}{L}_+^{}(q)\stackrel{~}{L}_{}^{}(q)+\stackrel{~}{L}_{}^{}(q)\stackrel{~}{L}_+^{}(q)\right)+\stackrel{~}{L}_3^{}(q)^2.$$
Hence, $`H_q(\theta ,\psi )`$ is still Casimir operator $`su(2)`$ Lie algebra with generator given in (2.18)-(2.23).
### II.3 Algebraic Solution of $`H_q(\theta ,\psi )`$ Hamiltonian
In order to obtain eigenspectrum of Hamiltonian (2.17) by algebraic method, first we obtain eigenspectrum of the Casimir operator (2.11). Since this operator commutes with left and right invariant generators given in (2.10), therefore, we can obtain representation of $`su(2)`$ simply by finding simultaneous eigenfunctions of the set of commuting operators, $`(R_3,L_3,L^2)`$. Denoting their simultaneous eigenfunction by $`\chi _{m_L,m_R}^l(\theta ,\psi ,\varphi )`$, we can write
$$L^2\chi _{m_L,m_R}^l(\theta ,\psi ,\varphi )=l(l+1)\chi _{m_L,m_R}^l(\theta ,\psi ,\varphi ),$$
(2.24)
$$L_3\chi _{m_L,m_R}^l(\theta ,\psi ,\varphi )=m_L\chi _{m_L,m_R}^l(\theta ,\psi ,\varphi ),$$
(2.25)
$$R_3\chi _{m_L,m_R}^l(\theta ,\psi ,\varphi )=m_R\chi _{m_L,m_R}^l(\theta ,\psi ,\varphi ).$$
(2.26)
Now solving the difference of the first order differential equations (2.25) and (2.26) we deduce that $`\chi _{m_L,m_R}^l(\theta ,\psi ,\varphi )`$ is proportional to $`e^{i(m_Rm_L)\varphi }`$, hence we have $`\chi _{m_L,m_R}^l(\theta ,\psi ,\varphi )=e^{i(m_Rm_L)\varphi }f(\theta ,\psi )`$, where $`f(\theta ,\psi )`$ can be determined from the solution of the sum of the equations (2.25) and (2.26), that is
$$i\mathrm{cos}(\theta )_\psi f(\theta ,\psi )+i\mathrm{sin}(\theta )\mathrm{cot}(\psi )_\theta f(\theta ,\psi )=(m_L+m_R)f(\theta ,\psi ).$$
(2.27)
Now considering the highest weight defined by $`m_L=m_R=l`$. this happenes if the right hand side of the equation (2.27) vanishes, hence it can be solved by characteristic method which leads to the following results:
$$\chi _{l,l}^l(\theta ,\psi ,\varphi )=\mathrm{exp}(2il\varphi )f^{max}(\mathrm{sin}(\psi )\mathrm{sin}(\theta ))$$
where $`f^{max}`$ is an arbitrary function which can be determined by solving the first order differential equation:
$$R_+\chi _{l,l}^l(\theta ,\psi ,\varphi )=0,L_+\chi _{l,l}^l(\theta ,\psi ,\varphi )=0.$$
Since the highest weight $`\chi _{l,l}^l`$ belongs to the kernel of raising operators $`R_+`$ and $`L_+`$, therefore the sum of the equations (2.2) and (2.5) leads to
$$u\frac{df^{max}(u)}{du}=2lf^{max}(u),$$
where $`u=\mathrm{sin}(\psi )\mathrm{sin}(\theta )`$. Therefore, solving the above equation we obtain $`f^{max}(u)=u^{2l}`$, hence $`\chi _{l,l}^l(\theta ,\psi ,\varphi )`$ has the following form
$$\chi _{l,l}^l(\theta ,\psi ,\varphi )=e^{2il\varphi }(\mathrm{sin}(\psi )\mathrm{sin}(\theta ))^{2l}.$$
(2.28)
The other eigenweights can be obtained through the operation of the lowering operator $`R_{}`$ and $`L_{}`$ over the highest eigenfunction, that is, we have
$$\chi _{m_L,m_R}^l(\theta ,\psi ,\varphi )=(L_{})^{lm_L}(R_{})^{l+m_R}(e^{2il\varphi }(\mathrm{sin}(\psi )\mathrm{sin}(\theta ))^{2l}).$$
(2.29)
In order to eliminate the coordinate $`\varphi `$, first we transfer the function $`e^{2il\varphi }`$ to the left hand side of the operators $`R_{}`$ and $`L_{}`$ in (2.29), then we get:
$$\chi _{m_L,m_R}^l(\theta ,\psi ,\varphi )=e^{i(m_Lm_R)\varphi }L_{}(m_Lm_R+1)L_{}(m_Lm_R+2)\mathrm{}$$
$$\mathrm{}L_{}(lm_R)R_{}(lm_R+1)\mathrm{}R_{}(2l)(\mathrm{sin}(\psi )\mathrm{sin}(\theta ))^{2l},$$
(2.30)
where the operators $`L_{}(m)`$ and $`R_{}(m)`$ are defined as:
$$L_{}(m)=\frac{i}{2}\left(\mathrm{sin}(\theta )_\psi +(i+\mathrm{cos}(\theta )\mathrm{cot}(\psi ))_\theta +im(\mathrm{cot}(\theta )i\frac{\mathrm{cot}(\psi )}{\mathrm{sin}(\theta )})\right),$$
(2.31)
$$R_{}(m)=\frac{i}{2}\left(\mathrm{sin}(\theta )_\psi +(i+\mathrm{cos}(\theta )\mathrm{cot}(\psi ))_\theta +im(\mathrm{cot}(\theta )i\frac{\mathrm{cot}(\psi )}{\mathrm{sin}(\theta )})\right).$$
(2.32)
Finally, the Fourier transformation of (2.30) leads to
$$\chi _{q,m}^l(\theta ,\psi )=L_{}(q+1)L_{}(q+2)\mathrm{}L_{}(l+\frac{qm}{2})R_{}(l+\frac{qm}{2}+1)\mathrm{}$$
$$\mathrm{}R_{}(2l)(\mathrm{sin}(\psi )\mathrm{sin}(\theta ))^{2l},$$
(2.33)
where $`q=m_Lm_R`$ and $`m=m_L+m_R`$.
Since the left and right invariant generators commute with each other, we can exchange these operators in (2.29) before Fourier transformation, whereas after Fourier transformation, we can use only the relation $`L_{}(q)R_{}(q1)=R_{}(q)L_{}(q1)`$. Since the Hamiltonian $`H_q(\theta ,\psi )`$ can be obtained from the relations (2.15) and (2.17) via similarity transformation and gauge transformation, we have:
$$H_q(\theta ,\psi )=\mathrm{exp}(\xi )L_q^2(\theta ,\psi )\mathrm{exp}(\xi ),L_q^2(\theta ,\psi )\chi _{q,m}^l(\theta ,\psi )=l(l+1)\chi _{q,m}^l(\theta ,\psi ),$$
(2.34)
where $`\xi =\frac{1}{2}\mathrm{ln}(\mathrm{sin}(\psi )\mathrm{sin}(\theta )).`$ Hence eigenfunction of Hamiltonian $`H_q(\theta ,\psi )`$ can be written as:
$$\stackrel{~}{\chi }_{q,m}^l(\theta ,\psi )=\mathrm{exp}(\xi )\chi _{q,m}^l(\theta ,\psi ).$$
(2.35)
### II.4 Shape Invariance Symmetry of $`H_q(\theta ,\psi )`$
Here in this section we show that the Hamiltonian $`H_q(\theta ,\psi )`$ possess both degeneracy and shape invariance symmetry . As it is shown in section (II.3), functions $`\stackrel{~}{\chi }_{q,m}^l(\theta ,\psi )=(\mathrm{sin}(\psi )\mathrm{sin}(\theta ))^{\frac{1}{2}}\chi _{q,m}^l(\theta ,\psi )`$ are eigenfunctions of Hamiltonian $`H_q(\theta ,\psi )`$ with the corresponding eigenvalue $`l(l+1)`$. Since $`|m_R|l`$ and $`|m_L|l`$, therefore, $`|q|2l`$ and for a given value of $`q`$ the parameter $`m`$ can take the following values:
$$m=\{\begin{array}{cc}0,\pm 2,\pm 4,\mathrm{},\pm (2l|q|)& for|q|=even,\\ \pm 1,\pm 3,\mathrm{},\pm (2l|q|)& for|q|=odd.\end{array}$$
(2.36)
Since the eigenvalue of Hamiltonian $`H_q(\theta ,\psi )`$ is independent of $`m`$, therefore it has $`(2l+1|q|)`$ degenerate states for a given $`l`$ or given energy $`l(l+1)`$. To see the shape invariance symmetry of Hamiltonian $`H_q(\theta ,\psi )`$, first we consider the Fourier transformed left and right invariant vector fields:
$$\stackrel{~}{L}_+(q)L_+(q1)e^\frac{}{q}$$
$$=\frac{i}{2}\left(\mathrm{sin}(\theta )_\psi +(i+\mathrm{cos}(\theta )\mathrm{cot}(\psi ))_\theta +i(q1)(\mathrm{cot}(\theta )+i\frac{\mathrm{cot}(\psi )}{\mathrm{sin}(\theta )})\right)e^\frac{}{q},$$
(2.37)
$$\stackrel{~}{L}_{}(q)L_{}(q+1)e^\frac{}{q}$$
$$=\frac{i}{2}\left(\mathrm{sin}(\theta )_\psi +(i+\mathrm{cos}(\theta )\mathrm{cot}(\psi ))_\theta +i(q+1)(\mathrm{cot}(\theta )i\frac{\mathrm{cot}(\psi )}{\mathrm{sin}(\theta )})\right)e^\frac{}{q},$$
(2.38)
$$\stackrel{~}{L}_3(q)L_3(q)=\frac{i}{2}(\mathrm{cos}(\theta )_\psi +\mathrm{sin}(\theta )\mathrm{cot}(\psi )_\theta iq)$$
(2.39)
and
$$\stackrel{~}{R}_+(q)R_+(q1)e^\frac{}{q}$$
$$=\frac{i}{2}\left(\mathrm{sin}(\theta )_\psi +(i+\mathrm{cos}(\theta )\mathrm{cot}(\psi ))_\theta +i(q1)(\mathrm{cot}(\theta )+i\frac{\mathrm{cot}(\psi )}{\mathrm{sin}(\theta )})\right)e^\frac{}{q},$$
(2.40)
$$\stackrel{~}{R}_{}(q)R_{}(q+1)e^\frac{}{q}$$
$$=\frac{i}{2}\left(\mathrm{sin}(\theta )_\psi +(i+\mathrm{cos}(\theta )\mathrm{cot}(\psi ))_\theta +i(q+1)(\mathrm{cot}(\theta )i\frac{\mathrm{cot}(\psi )}{\mathrm{sin}(\theta )})\right)e^\frac{}{q},$$
(2.41)
$$\stackrel{~}{R}_3(q)R_3(q)=\frac{i}{2}(\mathrm{cos}(\theta )_\psi +\mathrm{sin}(\theta )\mathrm{cot}(\psi )_\theta +iq).$$
(2.42)
After some tedious algebraic calculation we can derive the following relation between the above operators
$$L_3(q\pm 1)L_\pm (q)L_\pm (q)L_3(q)=\pm L_\pm (q)$$
(2.43)
$$R_3(q\pm 1)R_\pm (q)R_\pm (q)R_3(q)=R_\pm (q).$$
(2.44)
These relations indicate that Hamiltonian $`H_q(\theta ,\psi )`$ possesses shape invariance symmetry. Since by acting the operators $`R_\pm (q)`$ and $`L_\pm (q)`$ on both sides of eigenvalue equations:
$$L_q^2(\theta ,\psi )\chi _{q,m}^l(\theta ,\psi )=l(l+1)\chi _{q,m}^l(\theta ,\psi ),$$
$$R_3(q)\chi _{q,m}^l(\theta ,\psi )=\frac{mq}{2}\chi _{q,m}^l(\theta ,\psi ),$$
$$L_3(q)\chi _{q,m}^l(\theta ,\psi )=\frac{m+q}{2}\chi _{q,m}^l(\theta ,\psi ),$$
we get,
$$R_\pm (q)\chi _{q,m}^l(\theta ,\psi )=A_\pm (q,m)\chi _{q\pm 1,m1}^l(\theta ,\psi ),$$
(2.45)
$$L_\pm (q)\chi _{q,m}^l(\theta ,\psi )=B_\pm (q,m)\chi _{q\pm 1,m\pm 1}^l(\theta ,\psi ),$$
(2.46)
with
$$A_\pm (q,m)=\frac{1}{2}\sqrt{(2l(mq))(2l\pm (mq)+2)},$$
(2.47)
$$B_\pm (q,m)=\frac{1}{2}\sqrt{(2l(m+q))(2l\pm (m+q)+2)}.$$
(2.48)
The above relations imply that the pair of operators $`(L_{},R_+)`$ $`[(L_+,R_{})]`$ map degenerate eigenstates of Hamiltonian $`H_q(\theta ,\psi )`$ for a given value of $`q`$ into each other, that is they decrease \[increase\] the quantum number $`m`$ by 2 units as follows:
$$L_{}(q+1)R_+(q)\chi _{q,m}^l(\theta ,\psi )=A_+(q,m)B_{}(q+1,m1)\chi _{q,m2}^l(\theta ,\psi ),$$
$$L_+(q1)R_{}(q)\chi _{q,m}^l(\theta ,\psi )=A_{}(q,m)B_+(q1,m+1)\chi _{q,m+2}^l(\theta ,\psi ).$$
Now introducing the operator $`Y_+(q):=L_+(q1)R_{}(q)`$ and $`Y_{}(q):=L_{}(q+1)R_+(q)`$ as the raising and lowering operators of degenerates states of Hamiltonian $`H_q(\theta ,\psi )`$, we have the following shape invariance like symmetry between the degenerate states of Hamiltonian $`H_q(\theta ,\psi )`$:
$$Y_{}(q)Y_+(q)\chi _{q,m}^l(\theta ,\psi )=E(q,m)\chi _{q,m}^l(\theta ,\psi )$$
$$Y_+(q)Y_{}(q)\chi _{q,m+2}^l(\theta ,\psi )=E(q,m)\chi _{q,m+2}^l(\theta ,\psi )$$
where
$$E(q,m)=A_{}(q,m)A_+(q,m+2)B_{}(q+1,m+1)B_+(q1,m+1).$$
Thus, for a given value of $`q`$, we can obtain eigenfunction of Hamiltonian $`H_q(\theta ,\psi )`$ with eigenvalue $`l(l+1)`$, simply by acting the pairs of operators $`(L_{},R_+)[(L_+,R_{})]`$ over the highest weight \[lowest weight\], where here we have derived the eigenfunction $`\chi _{q,m}^l(\theta ,\psi )`$ by acting the lowering operator over the highest eigenstate as follows:
$$\chi _{q,m}^l(\theta ,\psi )=k^1(Y_{}(q))^{\frac{2l|q|m}{2}}\chi _{q,(2l|q|)}^l(\theta ,\psi ),$$
(2.49)
where
$$k=B_{}(q+1,m+1)B_{}(q+1,m+3)\mathrm{}$$
$$\times B_{}(q+1,2l|q|1)A_+(q,m+2)A_+(q,m+4)\mathrm{}A_+(q,2l|q|).$$
Using the relation (2.33) we can obtain the highest weight eigenstates for $`q>0`$ and $`q<0`$,
$$\chi _{q,(2l|q|)}^l(\theta ,\psi )=\{\begin{array}{cc}L_{}(q+1)L_{}(q+2)\mathrm{}& \\ \times L_{}(0)R_{}(1)R_{}(2)\mathrm{}R_{}(2l)(\mathrm{sin}(\theta )\mathrm{sin}(\psi ))^{2l}& forq<0\\ R_{}(q+1)R_{}(q+2)\mathrm{}R_{}(2l)(\mathrm{sin}(\theta )\mathrm{sin}(\psi ))^{2l}& forq>0.\end{array}$$
On the other hand pair operator $`(L_+,R_+)[or(L_{},R_{})]`$ leave the eigenvalue $`m`$ and $`l`$ unchanged while they increase \[decrease\] the parameter $`q`$ by 2 units, that is they map eigenfunction of Hamiltonian corresponding to the same energy with different $`q`$ into each other. That is, they map isospectral Hamiltonian into each other, which nothing but shape invariance. In order to show this shape invariance symmetry, we act the related operators over $`\chi _{q,m}^l(\theta ,\psi )`$, we then obtain:
$$L_+(q+1)R_+(q)\chi _{q,m}^l(\theta ,\psi )=A_+(q,m)B_+(q+1,m1)\chi _{q+2,m}^l(\theta ,\psi )$$
$$L_{}(q1)R_{}(q)\chi _{q,m}^l(\theta ,\psi )=A_{}(q,m)B_{}(q1,m+1)\chi _{q2,m}^l(\theta ,\psi ),$$
obviously, the combined action of above operators leave the eigenvalues $`m`$ and $`l`$ unchanged while changing the parameter $`q`$ by 2-units. Hence we define the operator $`X_+(q):=L_+(q+1)R_+(q)`$ and $`X_{}(q):=L_{}(q+1)R_{}(q+2)`$ as raising and lowering operators of parameter $`q`$. Then the shape invariance symmetry means:
$$X_{}(q)X_+(q)\chi _{q,m}^l(\theta ,\psi )=N(q,m)\chi _{q,m}^l(\theta ,\psi )$$
$$X_+(q)X_{}(q)\chi _{q+2,m}^l(\theta ,\psi )=N(q,m)\chi _{q+2,m}^l(\theta ,\psi )$$
where
$$N(q,m)=A_+(q,m)A_{}(q+2,m)B_+(q+1,m1)B_{}(q+1,m+1)$$
or
$$N(q,m)=\frac{1}{16}(2lmq)(2l+m+q+2)$$
$$\times \sqrt{(2lm+q)(2lm+q+4)(2l+mq+2)(2l+mq2)}.$$
For fixed values of energy $`l(l+1)`$ and given values of $`m`$, the parameter $`q`$ can take the following values
$$q=(2l|m|),(2l|m|2),\mathrm{},(2l|m|2),(2l|m|).$$
Hence obtaining the highest eigenstates, by solving the following first order differential equation
$$X_+(2l|m|)\chi _{(2l|m|),m}^l(\theta ,\psi )=0$$
where its integral leads to
$$\chi _{(2l|m|),m}^l(\theta ,\psi )=\{\begin{array}{cc}L_{}(2lm+1)L_{}(2lm+2)\mathrm{}L_{}(2l)(\mathrm{sin}(\theta )\mathrm{sin}(\psi ))^{2l}& form<0,\\ R_{}(2lm+1)R_{}(2lm+2)\mathrm{}R_{}(2l)(\mathrm{sin}(\theta )\mathrm{sin}(\psi ))^{2l}& form>0.\end{array}$$
Therefore using the shape invariance relation, we can obtain the eigenstates of Hamiltonian $`H_q(\theta ,\psi )`$ by consecutive action of $`q`$-lowering operator over $`q`$-highest weight eigenstate,
$$\chi _{q,m}^l(\theta ,\psi )=f^1X_{}(q)X_{}(q+2)\mathrm{}X_{}(2l|m|4)X_{}(2l|m|2)\chi _{(2l|m|),m}^l(\theta ,\psi )$$
$$f=A_(q+2,m)A_(q+4,m)\mathrm{}A_{}(2l|m|,m)$$
$$\times B_{}(q+1,m+1)B_{}(q+3,m+1)\mathrm{}B_{}(2l|m|1,m+1).$$
## III 3-dimensional Hamiltonian obtained from <br>4-Oscillators
Here in this section using the $`su(2)`$-parametrization of previous section, we obtain a special 3-dimensional Hamiltonian from the Hamiltonian of 4-oscillator with the same frequency, where we obtain its spectrum via the corresponding spectrum of 4-oscillator Hamiltonian. We show that such a Hamiltonian possesses shape invariance symmetry. The Hamiltonian of 4-oscillator with same frequency can be written as:
$$H=\frac{1}{2}\mathrm{\Sigma }_{i=0}^4(P_i^2+\frac{1}{2}\omega ^2x_i^2).$$
Now making the following change of variable:
$$\begin{array}{c}x_1=r\mathrm{sin}(\psi )\mathrm{sin}(\theta )\mathrm{sin}(\varphi ),\\ x_2=r\mathrm{sin}(\psi )\mathrm{sin}(\theta )\mathrm{cos}(\varphi ),\\ x_3=r\mathrm{sin}(\psi )\mathrm{cos}(\theta ),\\ x_4=r\mathrm{cos}(\psi ),\end{array}$$
(3.1)
where $`\psi ,\theta ,\varphi `$ are the same coordinates used in the parametrization $`su(2)`$ manifold, the Hamiltonian takes the following form
$$H(r,\theta ,\psi ,\varphi )=\frac{1}{2}[\frac{1}{r^3}_rr^3_r$$
$$+\frac{1}{r}(_\psi ^2+2\mathrm{cot}(\psi )_\psi +\frac{1}{\mathrm{sin}^2(\psi )}(_\theta ^2+\mathrm{cot}(\theta )_\theta +\frac{1}{\mathrm{sin}^2(\theta )}_\varphi ^2))]+\frac{1}{2}\omega ^2r^2.$$
(3.2)
Since angular part of the above Hamiltonian is the same, the one given in (2.11), therefore, its eigenspectrum can be obtained straightforwardly through routine separation variable into radial and angular part which we are not interested in it here in this work. Actually here we are concerned with special Hamiltonian which can be obtained from this 4-oscillator Hamiltonian, that is, those Hamiltonians which possess shape invariance symmetry.
In order to achieve this, we write the above Hamiltonian in terms of raising and lowering operators defined in the usual way:
$$H=\omega (a_1^{}a_1+a_2^{}a_2+a_3^{}a_3+a_4^{}a_4+2),$$
(3.3)
where $`a_i(a_i^{})`$ are defined as:
$$a_i=\sqrt{\frac{\omega }{2}}(x_i+\frac{1}{\omega }\frac{d}{dx_i}),a_i^{}=\sqrt{\frac{\omega }{2}}(x_i\frac{1}{\omega }\frac{d}{dx_i}).$$
These operators have the following form in radial coordinate (3.1)
$$a_1(a_1^{})=\sqrt{\frac{\omega }{2}}[r\mathrm{sin}(\psi )\mathrm{sin}(\theta )\mathrm{sin}(\varphi )+()\frac{1}{\omega }(\mathrm{sin}(\psi )\mathrm{sin}(\theta )\mathrm{sin}(\varphi )_r$$
$$\frac{1}{r}\mathrm{cos}(\psi )\mathrm{sin}(\theta )\mathrm{cos}(\varphi )_\psi \frac{1}{r}\frac{\mathrm{cos}(\theta )\mathrm{sin}(\varphi )}{\mathrm{sin}(\psi )}_\theta \frac{1}{r}\frac{\mathrm{cos}(\varphi )}{\mathrm{sin}(\psi )\mathrm{sin}(\theta )}_\varphi \left)\right],$$
$$a_2(a_2^{})=\sqrt{\frac{\omega }{2}}[r\mathrm{sin}(\psi )\mathrm{sin}(\theta )\mathrm{cos}(\varphi )+()\frac{1}{\omega }(\mathrm{sin}(\psi )\mathrm{sin}(\theta )\mathrm{cos}(\varphi )_r$$
$$+\frac{1}{r}\mathrm{cos}(\psi )\mathrm{sin}(\theta )\mathrm{cos}(\varphi )_\psi +\frac{1}{r}\frac{\mathrm{cos}(\theta )\mathrm{cos}(\varphi )}{\mathrm{sin}(\psi )}_\theta \frac{1}{r}\frac{\mathrm{sin}(\varphi )}{\mathrm{sin}(\psi )\mathrm{sin}(\theta )}\varphi \left)\right],$$
$$a_3(a_3^{})=\sqrt{\frac{\omega }{2}}[r\mathrm{sin}(\psi )\mathrm{cos}(\theta )+()\frac{1}{omega}(\mathrm{sin}(\psi )\mathrm{cos}(\theta )_r$$
$$+\frac{1}{r}\mathrm{cos}(\psi )\mathrm{cos}(\theta )_\psi \frac{1}{r}\frac{\mathrm{sin}(\theta )}{\mathrm{sin}(\psi )}_\theta ))],$$
$$a_4(a_4^{})=\sqrt{\frac{\omega }{2}}\left[r\mathrm{cos}(\psi )+()\frac{1}{\omega }\left(\mathrm{cos}(\psi )_r\frac{1}{r}\mathrm{sin}\psi _\psi \right)\right].$$
Now let us define the set of new operators $`A_i(A_i^{})`$, i= 1, 2 in terms of $`a_i(a_i^{})`$ :
$$A_1=\frac{1}{\sqrt{2}}(a_1+ia_2),A_1^{}=\frac{1}{\sqrt{2}}(a_1^{}ia_2^{}),$$
$$A_2=\frac{1}{\sqrt{2}}(a_1ia_2),A_2^{}=\frac{1}{\sqrt{2}}(a_1^{}+ia_2^{}),$$
where, these new operators have the following differential form in radial coordinates:
$$A_1=\frac{i}{\sqrt{2}}\sqrt{\frac{\omega }{2}}e^{i\varphi }[r\mathrm{sin}(\psi )\mathrm{sin}(\theta )$$
$$+\frac{1}{\omega }(\mathrm{sin}(\psi )\mathrm{sin}(\theta )_r\frac{1}{r}\mathrm{cos}(\psi )\mathrm{sin}(\theta )_\psi +\frac{1}{r}\frac{\mathrm{cos}(\theta )}{\mathrm{sin}(\psi )}_\theta +\frac{1}{r}\frac{i}{\mathrm{sin}(\psi )\mathrm{sin}(\theta )}\varphi ))],$$
(3.4)
$$A_1^{}=\frac{i}{\sqrt{2}}\sqrt{\frac{\omega }{2}}e^{i\varphi }[r\mathrm{sin}(\psi )\mathrm{sin}(\theta )$$
$$+\frac{1}{\omega }(\mathrm{sin}(\psi )\mathrm{sin}(\theta )_r+\frac{1}{r}\mathrm{cos}(\psi )\mathrm{sin}(\theta )_\psi +\frac{1}{r}\frac{\mathrm{cos}(\theta )}{sin(\psi )}_\theta \frac{1}{r}\frac{i}{\mathrm{sin}(\psi )\mathrm{sin}(\theta )}\varphi )],$$
(3.5)
$$A_2=\frac{i}{\sqrt{2}}\sqrt{\frac{\omega }{2}}e^{i\varphi }[r\mathrm{sin}(\psi )\mathrm{sin}(\theta )$$
$$+\frac{1}{\omega }(\mathrm{sin}(\psi )\mathrm{sin}(\theta )_r+\frac{1}{r}\mathrm{cos}(\psi )\mathrm{sin}(\theta )_\psi +\frac{1}{r}\frac{\mathrm{cos}(\theta )}{\mathrm{sin}(\psi )}_\theta \frac{1}{r}\frac{i}{\mathrm{sin}(\psi )\mathrm{sin}(\theta )}\varphi ))],$$
(3.6)
$$A_2^{}=\frac{i}{\sqrt{2}}\sqrt{\frac{\omega }{2}}e^{i\varphi }[r\mathrm{sin}(\psi )\mathrm{sin}(\theta )$$
$$\frac{1}{\omega }(\mathrm{sin}(\psi )\mathrm{sin}(\theta )_r+\frac{1}{r}\mathrm{cos}(\psi )\mathrm{sin}(\theta )_\psi +\frac{1}{r}\frac{\mathrm{cos}(\theta )}{sin(\psi )}_\theta +\frac{1}{r}\frac{i}{\mathrm{sin}(\psi )\mathrm{sin}(\theta )}\varphi )].$$
(3.7)
It is also straightforward to show that they have the following commutator relations:
$$[A_i,A_j^{}]=\delta _{ij},[A_i,A_j]=[A_i^{},A_j^{}]=0,i,j=1,2.$$
The 4-oscillators Hamiltonian (3.3) can be written in terms of the new oscillators as follows:
$$H=\omega (A_1^{}A_1+A_2^{}A_2+a_3^{}a_3+a_4^{}a_4+2).$$
(3.8)
Now its eigenspectrum can be obtained by solving the following eigenvalue equation
$$H\mathrm{\Psi }_{(n_1,n_2,n_3,n_4)}(r,\theta ,\varphi ,\psi )=E_{(n_1,n_2,n_3,n_4)}\mathrm{\Psi }_{(n_1,n_2,n3,n_4)}(r,\theta ,\varphi ,\psi ),$$
(3.9)
by the usual algebraic method. Hence its eigenfunction can be written as:
$$\mathrm{\Psi }_{(n_1,n_2,n_3,n_4)}(r,\theta ,\varphi ,\psi )=N(A_1^{})^{n_1}(A_2^{})^{n_2}(a_3^{})^{n_3}(a_4^{})^{n_4}\mathrm{exp}(\frac{\omega }{2}r^2),$$
(3.10)
with $`N=\frac{\omega }{\pi \sqrt{n_1!n_2!n_3!n_4!}}`$ as the normalization constant, and energy $`E_{(n_1,n_2,n_3,n_4)}=(n_1+n_2+n_3+n_4+2)\omega `$. Using the differential representation of the operator, the wavefunction (3.10) can be written in the following form
$$\mathrm{\Psi }_{(n_1,n_2,n_3,n_4)}(r,\theta ,\varphi ,\psi )=N2^{(1/2)(n_1+n_2)}e^{i(n_2n_1)\varphi }e^{(1/2)r^2}(r\mathrm{sin}(\psi )\mathrm{sin}(\theta ))^{(n_1+n_2)}$$
$$\times _{n_3}(r\mathrm{sin}(\psi )\mathrm{cos}(\theta ))_{n_4}(r\mathrm{cos}(\psi ))\mathrm{\Sigma }_{i=0}^{n_1}(1)^ii!\left(\begin{array}{c}n_1\\ i\end{array}\right)\left(\begin{array}{c}n_2\\ i\end{array}\right)(r\mathrm{sin}(\psi )\mathrm{sin}(\theta ))^{2i},$$
(3.11)
where $``$<sub>n</sub> is the Hermit polynomial of degree $`n`$ and $`\left(\begin{array}{c}n\\ r\end{array}\right)=\frac{n!}{r!(nr)!}`$ . Now with the same prescription used in the previous section, we can eliminate $`\varphi `$, by Fourier transforming over it. Hence, by the Fourier transformation over $`\varphi `$, the 4-oscillator Hamiltonian reduces to the following Hamiltonian:
$$H_m(r,\theta ,\psi )=\frac{1}{2}[\frac{1}{r^3}_rr^3_r$$
$$+\frac{1}{r^2}(_\psi ^2+2\mathrm{cot}(\psi )_\psi +\frac{1}{\mathrm{sin}^2(\psi )}(_\theta ^2+\mathrm{cot}(\theta )_\theta \frac{m^2}{\mathrm{sin}^2(\theta )}))]+\frac{1}{2}\omega ^2r^2,$$
(3.12)
where after similarity transformation through function $`r^{1/2}`$, it reduces to
$$\stackrel{~}{H}_m(r,\theta ,\psi )=r^{1/2}H_m(r,\theta ,\psi )r^{1/2}=\frac{1}{2}[\frac{1}{r^2}_rr^2_r$$
$$+\frac{1}{r^2}(_\psi ^2+2\mathrm{cot}(\psi )_\psi +\frac{1}{\mathrm{sin}^2(\psi )}(_\theta ^2+\mathrm{cot}(\theta )_\theta \frac{m^2}{\mathrm{sin}^2(\theta )})]+\frac{1}{2}\omega ^2r^2+\frac{3}{8r^2}.$$
(3.13)
On the other hand, the Hamiltonian $`H_m(r,\theta ,\psi )`$ given by (3.12) can be written in the following form
$$H_m(r,\theta ,\psi )=\omega (A_1^{}(m+1)A_1(m)+A_2^{}(m1)A_2(m)+a_3^{}a_3+a_4^{}a_4+2),$$
(3.14)
with
$$A_1(m)=\frac{i}{\sqrt{2}}\sqrt{\frac{\omega }{2}}[r\mathrm{sin}(\psi )\mathrm{sin}(\theta )$$
$$+\frac{1}{\omega }(\mathrm{sin}(\psi )\mathrm{sin}(\theta )_r\frac{1}{r}\mathrm{cos}(\psi )\mathrm{sin}(\theta )_\psi +\frac{1}{r}\frac{\mathrm{cos}(\theta )}{\mathrm{sin}(\psi )}_\theta \frac{1}{r}\frac{m}{\mathrm{sin}(\psi )\mathrm{sin}(\theta )})],$$
(3.15)
$$A_1^{}(m)=\frac{i}{\sqrt{2}}\sqrt{\frac{\omega }{2}}[r\mathrm{sin}(\psi )\mathrm{sin}(\theta )$$
$$\frac{1}{\omega }(\mathrm{sin}(\psi )\mathrm{sin}(\theta )_r+\frac{1}{r}\mathrm{cos}(\psi )\mathrm{sin}(\theta )_\psi +\frac{1}{r}\frac{\mathrm{cos}(\theta )}{sin(\psi )}_\theta +\frac{1}{r}\frac{m}{\mathrm{sin}(\psi )\mathrm{sin}(\theta )})],$$
(3.16)
$$A_2(m)=\frac{i}{\sqrt{2}}\sqrt{\frac{\omega }{2}}[r\mathrm{sin}(\psi )\mathrm{sin}(\theta )$$
$$+\frac{1}{\omega }(\mathrm{sin}(\psi )\mathrm{sin}(\theta )_r+\frac{1}{r}\mathrm{cos}(\psi )\mathrm{sin}(\theta )_\psi +\frac{1}{r}\frac{\mathrm{cos}(\theta )}{\mathrm{sin}(\psi )}_\theta +\frac{1}{r}\frac{m}{\mathrm{sin}(\psi )\mathrm{sin}(\theta )})],$$
(3.17)
$$A_2^{}(m)=\frac{i}{\sqrt{2}}\sqrt{\frac{\omega }{2}}[r\mathrm{sin}(\psi )\mathrm{sin}(\theta )\frac{1}{\omega }$$
$$(\mathrm{sin}(\psi )\mathrm{sin}(\theta )_r+\frac{1}{r}\mathrm{cos}(\psi )\mathrm{sin}(\theta )_\psi +\frac{1}{r}\frac{\mathrm{cos}(\theta )}{sin(\psi )}_\theta \frac{1}{r}\frac{m}{\mathrm{sin}(\psi )\mathrm{sin}(\theta )})].$$
(3.18)
It is straightforward to derive the following relation between Hamiltonian (3.12) and operator $`A_i(m)(A_i^{}(m))`$, i= 1, 2:
$$\begin{array}{c}H(m1)A_1^{}(m)A_1^{}(m)H(m)=\omega A_1^{}(m),\\ H(m+1)A_2^{}(m)A_2^{}(m)H(m)=\omega A_2^{}(m),\\ H(m+1)A_1(m)A_1(m)H(m)=\omega A_1(m),\\ H(m1)A_2(m)A_2(m)H(m)=\omega A_2(m),\end{array}$$
(3.19)
where $`H(m):=H_m(r,\theta ,\psi )`$. The above relations indicate that Hamiltonian (3.12) possesses shape invariance symmetry. To see this, we consider the Fourier transformation of eigenvalue equation (3.9):
$$H(m)\mathrm{\Psi }_{(n,m,n_3,n_4)}(r,\theta ,\psi )=E_{(n,n_3,n_4)}\mathrm{\Psi }_{(n,m,n_3,n_4)}(r,\theta ,\psi ),$$
(3.20)
where $`n=n_1+n_2,m=n_2n_1`$ and $`E_{(n,n_3,n_4)}=(n+n_3+n_4+2)\omega `$. Since $`n_1`$ and $`n_2`$ are positive integers, therefore $`n`$ is also a positive integer but $`m`$ is an integer. For a given value of $`m`$, the quantum number $`n`$ can be either even or odd integer, since, quantum numbers $`n_1`$ and $`n_2`$ vary by the same amount, so that $`m`$ remains constant. Actually for some given value of $`m`$, the quantum number $`n`$ can take the following values
$$n=|m|,|m|+2,|m|+4,\mathrm{}.$$
On the other hand, in terms of $`n`$, the quantum number $`m`$ can take the following values
$$m=n,n+2,\mathrm{},n2,n.$$
It is interesting to see that energy of Hamiltonian $`H_m(r,\theta ,\psi )`$ is independent of $`m`$, hence these Hamiltonians are isospectral which is due to the existence of shape invariance symmetry as we show below.
Operating the operator $`A_1^{}(m)`$ on both sides of the eigenvalue relation (3.20) and using the relations (3.19), we get
$$H(m1)(A_1^{}(m)\mathrm{\Psi }_{n,m}(r,\theta ,\psi ))=(E_n+\omega )(A_1^{}(m)\mathrm{\Psi }_{n,m}(r,\theta ,\psi )),$$
therefore, $`A_1^{}(m)\mathrm{\Psi }_{n,m}(r,\theta ,\psi ))`$ corresponds to the eigenfunction of $`H(m1)`$ with corresponding eigenvalue $`E_{n+1}`$, that is
$$A_1^{}(m)\mathrm{\Psi }_{n,m}(r,\theta ,\psi )=\sqrt{\frac{nm}{2}+1}\mathrm{\Psi }_{n+1,m1}(r,\theta ,\psi ),$$
where $`\mathrm{\Psi }_{n,m}(r,\theta ,\psi ):=\mathrm{\Psi }_{(n,m,n_3,n_4)}(r,\theta ,\psi )`$ and $`E_n:=E_{(n,n_3,n_4)}.`$ Similarly, operating $`A_2^{}(m)`$ on both sides of (3.20) and using (3.19) we get:
$$H(m+1)(A_2^{}(m)\mathrm{\Psi }_{n,m}(r,\theta ,\psi ))=(E_n+\omega )(A_2^{}(m)\mathrm{\Psi }_{n,m}(r,\theta ,\psi )),$$
which leads to
$$A_2^{}(m)\mathrm{\Psi }_{n,m}(r,\theta ,\psi )=\sqrt{\frac{n+m}{2}+1}\mathrm{\Psi }_{n+1,m+1}(r,\theta ,\psi ).$$
Also by acting the operators $`A_1(m)`$ and $`A_2(m)`$ on the eigenvalue relation (3.20) and using the relations (3.19) we obtain
$$H(m+1)(A_1(m)\mathrm{\Psi }_{n,m}(r,\theta ,\psi ))=(E_n\omega )(A_1(m)\mathrm{\Psi }_{n,m}(r,\theta ,\psi )),$$
$$H(m1)(A_2(m)\mathrm{\Psi }_{n,m}(r,\theta ,\psi ))=(E_n\omega (A_2(m)\mathrm{\Psi }_{n,m}(r,\theta ,\psi )),$$
which imply the following relations
$$A_1(m)\mathrm{\Psi }_{n,m}(r,\theta ,\psi )=\sqrt{\frac{nm}{2}}\mathrm{\Psi }_{n1,m+1}(r,\theta ,\psi ),$$
$$A_2(m)\mathrm{\Psi }_{n,m}(r,\theta ,\psi )=\sqrt{\frac{n+m}{2}}\mathrm{\Psi }_{n1,m1}(r,\theta ,\psi ).$$
From the above relations we conclude that the pair of operators $`(A_2(m),A_1^{}(m))`$ or
$`(A_2^{}(m),A_1(m))`$ acting at eigenfunction $`\mathrm{\Psi }_{n,m}(r,\theta ,\psi )`$ of Hamiltonian $`H(m)`$, give eigenfunction of Hamiltonian $`H(m\pm 2)`$ with same the energy as follows:
$$A_2(m1)A_1^{}(m)\mathrm{\Psi }_{n,m}(r,\theta ,\psi )=\frac{1}{2}\sqrt{(n+m)(nm+2)}\mathrm{\Psi }_{n,m2}(r,\theta ,\psi ),$$
$$A_2^{}(m+1)A_1(m)\mathrm{\Psi }_{n,m}(r,\theta ,\psi )=\frac{1}{2}\sqrt{(nm)(n+m+2)}\mathrm{\Psi }_{n,m2}(r,\theta ,\psi ).$$
Now introducing the operators $`A_{}(m):=A_2(m1)A_1^{}(m)`$ and $`A_+(m):=A_2^{}(m1)A_1(m2)`$, we have:
$$A_{}(m)A_+(m)\mathrm{\Psi }_{n,m2}(r,\theta ,\psi )=E(n,m)\mathrm{\Psi }_{n,m2}(r,\theta ,\psi ),$$
$$A_+(m)A_{}(m)\mathrm{\Psi }_{n,m}(r,\theta ,\psi )=E(n,m)\mathrm{\Psi }_{n,m}(r,\theta ,\psi ),$$
where
$$E(n,m)=\frac{1}{4}(n+m)(nm+2).$$
The above relations show the existence of shape invariance symmetry between the Hamiltonian $`H(m)`$ and $`H(m2)`$ with same given eigenvalue $`E_n`$. Hence we can obtain the eigenfunction $`\mathrm{\Psi }_{n,m}(r,\theta ,\psi )`$ of Hamiltonian $`H(m)`$ by consecutive action of related raising operators over $`\mathrm{\Psi }_{n,n}(r,\theta ,\psi )`$:
$$\mathrm{\Psi }_{n,m}(r,\theta ,\psi )=c^1A_{}(m+2)A_{}(m+4)\mathrm{}A_{}(n2)A_{}(n)\mathrm{\Psi }_{n,n}(r,\theta ,\psi ),$$
where
$$c=\frac{1}{2^{\frac{nm}{2}}}\sqrt{(nm)!!2n(2n2)\mathrm{}(n+m+4)(n+m+2)},$$
and
$$(nm)!!=(nm)(nm2)\mathrm{}4\times 2,$$
$$\mathrm{\Psi }_{n,n}(r,\theta ,\psi )\mathrm{\Psi }_{(n,n,n_3,n_4)}(r,\theta ,\psi )=(a_3^{})^{n_3}(a_4^{})^{n_4}A_2^{}(n1)A_2^{}(n2)\mathrm{}A_2^{}(1)A_2^{}(0)e^{\frac{1}{2}\omega r^2}.$$
Of course we can obtain the eigenfunction $`\mathrm{\Psi }_{(n,m,n_3,n_4)}(r,\theta ,\psi )`$ by reduction of coordinate $`\varphi `$ via Fourier transformation of (3.11), which has the following form:
$$\mathrm{\Psi }_{(n,m,n_3,n_4)}(r,\theta ,\psi )=N2^{\frac{n}{2}}e^{(1/2)r^2}(r\mathrm{sin}(\psi )\mathrm{sin}(\theta ))^n$$
$$\times _{n_3}(\mathrm{sin}(\psi )\mathrm{sin}(\theta ))_{n_4}(r\mathrm{cos}(\psi ))\mathrm{\Sigma }_{i=0}^{\frac{nm}{2}}(1)^ii!\left(\begin{array}{c}\frac{nm}{2}\\ i\end{array}\right)\left(\begin{array}{c}\frac{n+m}{2}\\ i\end{array}\right)(r\mathrm{sin}(\psi )\mathrm{sin}(\theta ))^{2i}.$$
## IV CONCLUSION
Here in this work having Fourier transformed 3 and 4-dimensional Hamiltonians associated with $`su(2)`$ and Heisenberg Lie algebra we have been able to obtain 2 and 3-dimensional Hamiltonian whit shape invariance symmetry. It would be interesting to obtain many-body Hamiltonian in one dimension or higher, which possesses shape invariance symmetry by appropriate Fourier transformation over some coordinates of Hamiltonian associated with higher ranks semisimple and non semisimple Lie algebra. This is under investigation.
ACKNOWLEDGEMENT
We wish to thank Dr. S.K.A. Seyed Yagoobi for carefully reading the article and for his constructive comments.
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# 1 Introduction
## 1 Introduction
In the last few years many high accuracy calculations have been made of cosmological consequences of Nambu-Goto (NG) cosmic strings . Indeed, such predictions were recently compared to the Boomerang data . There are good reasons why most studies of the cosmological effects of topological defects (see for a summary) have concentrated on NG strings: these are the simplest type of cosmic string, and their equations of motion (at least in Minkowski space) can be solved exactly. On a lattice one can use the highly efficient Smith-Vilenkin algorithm ; and in fact some of the recent predictions are based on Minkowski space codes of NG network evolution .
However NG cosmic strings, of which the simplest example are the strings formed in the Abelian Higgs model, are not likely to be the most realistic type of cosmic string. Cosmic strings that can create significant density perturbations require GUT scale physics. If the Higgs field that forms the string couples to fermions in the GUT theory — as might well be expected of a Higgs field — then these fermions yield zero modes in the core of the string , thereby generating a current (which need not be electromagnetically coupled) along the string.
In this paper we are interested in chiral cosmic strings which arise naturally in supersymmetric (SUSY) theories , where a $`U(1)`$ symmetry is broken with a Fayet-Iliopoulos D-term, resulting in a single fermion zero mode which travels in only one direction along the string — this defines a chiral string. In the cosmological context, chiral strings are automatically formed at the end of inflation in SUSY models with a D-term . Furthermore the zero mode (or chiral nature of the string) survives the subsequent supersymmetry breaking phase transition , and consequently both inflation and chiral cosmic strings contribute to the density and CMBR perturbations in this scenario. Calculations of these observable predictions were carried out recently , and there the $`C_l`$’s were decomposed as
$$C_l=\alpha C_l^{\mathrm{inflation}}+(1\alpha )C_l^{\mathrm{NG}\mathrm{strings}}$$
with $`0\alpha 1`$. However, we believe that the use of the $`C_l`$ from NG strings in the above formula is an oversimplification at least in the inflation plus chiral cosmic string scenario mentioned above. In the case of chiral cosmic strings, the presence of the fermion zero mode is likely to have a significant effect on the dynamics of the string network, which could therefore evolve very differently to a NG network. Indeed the action describing the evolution of chiral cosmic strings is very different from the NG action . The purpose of this paper is to quantify some of the differences between the evolution of these two different types of cosmic string network.
Our starting point is the action for chiral cosmic strings first proposed in . As usual, in order to derive the equations of motion from this action, gauge choices must be made: in Minkowski space with suitable gauge choices the equations of motion reduce to the remarkably familiar form given by
$`{\displaystyle \frac{^2𝐱}{t^2}}{\displaystyle \frac{^2𝐱}{\sigma ^2}}`$ $`=`$ $`0𝐱(t,\sigma )={\displaystyle \frac{1}{2}}[𝐚(t+\sigma )+𝐛(t\sigma )],`$
$`\stackrel{´}{𝐚}^2`$ $`=`$ $`1,`$
$`\stackrel{´}{𝐛}^2`$ $``$ $`1,`$
where for instance $`\stackrel{´}{𝐚}(q)d𝐚(q)/dq.`$ These can be recognized as the usual NG equations of motion with the only difference being the constraint on the derivative of $`𝐛`$ which must now lie within the Kibble-Turok sphere rather than on it. We show in section 2.2 that the physical reason for this stems from the conserved charge on chiral strings: if this charge is zero then $`\stackrel{´}{𝐛}^2=1`$ and one is left with NG strings as required.
If the charge on the strings is maximal then $`\stackrel{´}{𝐛}^2=0`$. As we show in section 3, this latter special case is very interesting since it implies that $`\dot{𝐱}=𝐱^{}`$ with $`|\dot{𝐱}|=1/2`$ so that the strings move along themselves at half the speed of light and never change shape or self-intersect. In the case of loops, this solution corresponds to arbitrary-shape stable vortons. For infinite strings it means that self-intersections (which produce loops) never occur. Since this is an important mechanism of energy loss in the case of NG strings, this result already gives an indication that the evolution of chiral and NG cosmic string networks may be very different. For general chiral string solutions, self-intersection is certainly possible, but we may expect that the probability is lower than in the NG case. To quantify the differences between NG and chiral string evolution, we study in section 3 the self-intersection probability of loops with different numbers of harmonics on them and different conserved charges.
The plan of this paper is the following. For pedagogical reasons, we begin in section 2.1 by briefly reviewing the work of Carter and Peter (CP) . For chiral strings it is not possible to impose exactly the same gauge conditions as are usually chosen for NG strings. In section 2.2 we review the possible gauge choices, in particular those made by Martins and Shellard and by CP. We follow the latter, whose choice leads to very simple equations of motion, which are almost identical to those of the NG string and which, most importantly, are exactly integrable in Minkowski space. In section 3 we use these simple equations to discuss general properties of chiral strings as well as the self-intersection properties of loops with different numbers of harmonics and different charges. Finally, our conclusions and plans for future work are discussed in section 4.
Note: Whilst we were trying to extend the results presented here to FRW universes, a paper by Blanco-Pillado et al. appeared which also obtains the equations of motion above, though from a rather different point of view. Here we follow more closely the work of CP , and extend both of these papers to study some general properties of chiral cosmic strings and the self-intersection of loops.
## 2 Chiral and Nambu-Goto strings
### 2.1 Action and equations of motion
For pedagogical reasons, we here review briefly the work of Carter and Peter .
The action for chiral cosmic strings they proposed involves a dimensionless scalar field $`\varphi `$ which can be interpreted as the phase of the current carriers condensed on the string. The action is
$$S=d^2\sigma \sqrt{\gamma }\left(m^2\frac{1}{2}\psi ^2\gamma ^{ij}\varphi _{,i}\varphi _{,j}\right),$$
(2.1)
where $`\gamma _{ij}`$ ($`i,j=\{0,1\}`$) is the induced metric on the world sheet:
$$\gamma _{ij}=x_{,i}x_{,j}g_{\mu \nu }x_{,i}^\mu x_{,j}^\nu .$$
Here $`g_{\mu \nu }`$ is the background metric, with signature $`(+,,,)`$ (which in the following we shall take to be the Minkowski metric) and $`x^\mu (\sigma ^0,\sigma ^1)`$ denotes the position of the string at world-sheet coordinates $`\sigma ^i`$. The first term in (2.1) is just the NG action for a string with tension $`m^2`$.
The action (2.1) is invariant under reparametrizations, $`\sigma ^i\stackrel{~}{\sigma }^i=\stackrel{~}{\sigma }^i(\sigma ^j)`$, and also under transformations of $`\varphi `$, with a compensating transformation of $`\psi `$:
$$\varphi \stackrel{~}{\varphi }(\varphi ),\mathrm{with}\psi \stackrel{~}{\psi }=\left(\frac{d\stackrel{~}{\varphi }}{d\varphi }\right)^1\psi .$$
(2.2)
These freedoms must be removed by making gauge choices, as discussed below.
The dimensionless Lagrange multiplier $`\psi `$ sets the constraint
$$\gamma ^{ij}\varphi _{,i}\varphi _{,j}=0$$
(2.3)
which ensures that the cosmic strings are indeed chiral. Generally, current-carrying strings are characterized by two currents and their corresponding charges . One of these currents, proportional to $`\gamma ^{ij}\varphi _{,j}`$, is conserved by virtue of the equations of motion; the other, proportional to $`ϵ^{ij}\varphi _{,j}`$, is topologically conserved. However, for chiral strings, because of (2.3) the two currents coincide, and so therefore do the two corresponding charges $`Z`$ and $`N`$.
In Minkowski space, the equations of motion following from (2.1), in addition to (2.3), are
$$_i(\sqrt{\gamma }\psi ^2\gamma ^{ij}\varphi _{,j})=0$$
(2.4)
and finally
$$_i\left[\sqrt{\gamma }\left(\gamma ^{ij}+\frac{\psi ^2}{m^2}\varphi ^{,i}\varphi ^{,j}\right)x_{,j}^\mu \right]=0.$$
(2.5)
In 1+1 dimensions, a scalar field whose gradient is everywhere null is necessarily harmonic, i.e., (2.3) implies
$$0=^j_j\varphi =\frac{1}{\sqrt{\gamma }}_i(\sqrt{\gamma }\gamma ^{ij}\varphi _{,j}),$$
which is consistent with (2.3) only if $`\psi `$ is a function of $`\varphi `$, $`\psi =\psi (\varphi )`$. We shall verify this later in particular coordinate systems.
From (2.4), we may provisionally define the current as $`j^i=\psi ^2\varphi ^{,i}`$, which is conserved and null, satisfying
$$j_ij^i=0.$$
However, this expression is *not* invariant under the $`\varphi `$ transformation (2.2). There is an ambiguity in the definition of $`j^i`$: any current of the form $`j^i=f(\varphi )\varphi ^{,i}`$ is null and conserved; there is an infinity of conservation laws. This is a peculiarity of null currents in 1+1 dimensions. For definiteness, we choose the invariant current
$$j^i=\psi \varphi ^{,i}.$$
(2.6)
### 2.2 Gauge choices
To proceed further, gauge choices must be made. The action (2.1) is reparametrization invariant — we can replace $`\sigma ^0`$ and $`\sigma ^1`$ by any functions of these variables. Moreover, we can replace $`\varphi `$ by any function of itself, provided we change $`\psi `$ to compensate as in (2.2).
For the NG string it is usual to choose the conformal gauge in which $`\gamma _{ij}(\sigma ^k)=\mathrm{\Omega }(\sigma ^k)\eta _{ij}`$, with $`\eta _{ij}=\mathrm{diag}(1,1)`$. Explicitly, if we write $`\tau =\sigma ^0`$, $`\sigma =\sigma ^1`$, and denote derivatives with respect to $`\tau `$ and $`\sigma `$ by a dot and prime respectively, the conformal gauge is specified by two conditions:
$$\dot{x}^2+x^2=0$$
(2.7)
and
$$\dot{x}x^{}=0.$$
(2.8)
This implies that the string’s velocity is perpendicular to its tangent vector. For the NG string, because of the conformal invariance of the action, these two conditions are not in fact independent, and therefore do not fully specify the coordinates. We can in addition impose the temporal gauge condition
$$\tau =tx^0.$$
(2.9)
For chiral strings, however, these three conditions are inconsistent, so different choices are needed.
We first recall the choices made by Martins and Shellard (MS) , who studied a number of properties of chiral cosmic string loops. They opted to maintain the temporal gauge condition (2.9). As noted above, there is no longer freedom to choose the full conformal gauge as well. Instead MS chose the world-sheet metric to be diagonal, maintaining (2.8) but not (2.7). On defining $`ϵ^2=(\sqrt{\gamma }\gamma ^{00})^2=𝐱^2/(1\dot{𝐱}^2)`$, it follows from (2.3) that $`\varphi ^2=ϵ^2\dot{\varphi }^2`$ so that equation (2.4) yields
$$_t[\psi ^2\varphi ^{}]=_\sigma [\psi ^2\dot{\varphi }],$$
confirming the general result that $`\psi =\psi (\varphi )`$. MS chose to fix $`\varphi `$ by setting $`\psi ^2=`$ const $`=1`$. The equations of motion following from (2.5) are then given by
$`[ϵ(1+\mathrm{\Phi })]\dot{}`$ $`=`$ $`\mathrm{\Phi }^{}`$
$`ϵ(1+\mathrm{\Phi })\ddot{𝐱}`$ $`=`$ $`\left[(1\mathrm{\Phi }){\displaystyle \frac{𝐱^{}}{ϵ}}\right]^{}+\dot{\mathrm{\Phi }}𝐱^{}+2\mathrm{\Phi }\dot{𝐱}^{}`$ (2.10)
where $`=d/dt`$ and $`{}_{}{}^{}=d/d\sigma `$ and $`\mathrm{\Phi }=\dot{\varphi }^2/(m^2\gamma _{00})`$. Loop solutions to these equations were studied in .
Since these gauge choices lead to rather complicated equations of motion, we shall opt instead to follow the paper of Carter and Peter (CP) and choose one of the world-sheet coordinates to be proportional to $`\varphi `$. That is, choose $`\eta =m^1\varphi `$ (the factor of $`m^1`$ is introduced for dimensional reasons) to be one world-sheet coordinate and denote the second by $`q`$. By (2.3) this implies
$$\gamma ^{\eta \eta }=0\gamma _{qq}=0.$$
(2.11)
Thus the line element on the world-sheet is
$$ds^2=Ad\eta ^2+2\mathrm{\Omega }dqd\eta $$
where
$$\mathrm{\Omega }\gamma _{\eta q}=\sqrt{\gamma }=x_\eta x_q$$
and
$$A\gamma _{\eta \eta }=x_{,\eta }x_{,\eta }.$$
(2.12)
With this choice of coordinates, we see again that $`\psi =\psi (\varphi )`$ since equation (2.4) gives
$$0=_q[\sqrt{\gamma }\gamma ^{q\eta }m^1\psi ^2]=m^1_q[\psi ^2]\psi =\psi (\varphi ).$$
We also note from (2.6) that
$$j_\eta =m\psi ,j_q=0,j^\eta =0,j^q=\frac{m\psi }{\mathrm{\Omega }}.$$
(2.13)
Now, the equation of motion (2.5) gives
$$2_q_\eta x^\mu +_q[F(_qx^\mu )]=0$$
(2.14)
where
$$F=m^2(\mathrm{\Omega }\gamma ^{qq}+\psi ^2\gamma ^{q\eta })=\frac{m^2}{\mathrm{\Omega }}\left(\psi ^2A\right).$$
It is now clear that we can further simplify the equations (2.14) by choosing the second coordinate $`q`$ in such a way that $`F=0`$. Happily this is a consistent choice because then $`A=x_{,\eta }x_{,\eta }=\psi ^2`$ should be a function of $`\varphi `$ only, independent of $`q`$. But this is indeed the content of the simplified equation of motion, (2.14) with $`F=0`$, namely
$$_q_\eta x^\mu =0.$$
(2.15)
The simplicity of this equation shows the convenience of this gauge choice, in which both $`\eta `$ and $`q`$ are characteristic coordinates. The general solution of the equation of motion is
$$x^\mu (q,\eta )=\frac{1}{2}[a^\mu (q)+b^\mu (\eta )],$$
exactly as for the NG string.
Within the gauge choices so far made, we still have freedom to transform each of the coordinates $`\eta `$ and $`q`$ separately: $`\eta \stackrel{~}{\eta }(\eta )`$ and $`q\stackrel{~}{q}(q)`$. It is convenient to choose them so that $`\eta =a^0`$ and $`q=b^0`$, and hence $`tx^0=\frac{1}{2}(\eta +q)`$. This is essentially a temporal gauge. Finally, let
$$q=t+\sigma ,\eta =t\sigma .$$
Then the equations of motion (2.15) and constraints (2.12) and (2.11) reduce to
$`\ddot{𝐱}𝐱^{\prime \prime }`$ $`=`$ $`0𝐱(t,\sigma )={\displaystyle \frac{1}{2}}[𝐚(q)+𝐛(\eta )],`$
$`\left({\displaystyle \frac{d𝐚}{dq}}\right)^2`$ $``$ $`\stackrel{´}{𝐚}^2=1,`$ (2.16)
$`\left({\displaystyle \frac{d𝐛}{d\eta }}\right)^2`$ $``$ $`\stackrel{´}{𝐛}^21,`$
where $`=d/dt`$ and $`{}_{}{}^{}=d/d\sigma `$. Notice that equations (2.16) resemble very closely the NG equations and constraints in the temporal, conformal gauge: the only difference is that now $`\stackrel{´}{𝐛}`$ is constrained to lie within the Kibble-Turok sphere rather than on it. The physical reason for this will be discussed below.
We note that Blanco-Pillado et al. made essentially the same gauge choice, though without introducing the Lagrange-multiplier variable $`\psi `$.
The meaning of the coordinate $`\sigma `$ can be understood by constructing the stress energy tensor. With the gauge choices made above this is given by
$$T^{\mu \nu }(t,𝐲)=m^2d\sigma (\dot{x}^\mu \dot{x}^\nu x^\mu {}_{}{}^{}x_{}^{\nu }{}_{}{}^{})\delta ^3(𝐲𝐱(t,\sigma )),$$
which is formally identical to the NG stress energy tensor in the conformal-temporal gauge. Since
$$T^{00}(t,𝐲)=m^2𝑑\sigma \delta ^3(𝐲𝐱(t,\sigma ))$$
is conserved in Minkowski space, it follows that $`\sigma `$ again measures the energy or ‘invariant length’ along the string.
Finally, we examine the reason for the inequality $`\stackrel{´}{𝐛}^2<1`$, which follows from the conserved charge on the string. From (2.13) the physical current on the string is given by
$$j^t=j^\sigma =\frac{m\psi }{2\mathrm{\Omega }},$$
so that the conserved charge is
$$N=Z=𝑑\sigma \sqrt{\gamma }j^t=\frac{1}{2}𝑑\sigma m\psi .$$
Now, let
$$\stackrel{´}{𝐛}^2=k^2,$$
so that $`\psi ^2=A=x_{,\eta }x_{,\eta }=\stackrel{´}{b}^2/4=(1k^2)/4`$, from which it clearly follows that $`k^21`$. Note that
$$N=\frac{m}{4}𝑑\sigma \sqrt{1k^2},$$
(2.17)
so the value of $`k`$ determines the charge on the string: this takes its maximum value when $`k=0`$ everywhere (as we will see below this corresponds to interesting vorton solutions), and $`N=0`$ when $`k1`$, which is exactly the NG limit as required.
In the next section we study the self-intersection properties of loops with different values of $`N`$.
## 3 Properties of chiral cosmic strings and loop self-intersections
In this section we first describe some general properties of chiral cosmic strings which follow from equations (2.16). Loop self-intersections are then studied.
### 3.1 General properties
From equations (2.16), the velocity and tangent vectors of the string are given by
$$\dot{𝐱}(t,\sigma )=\frac{1}{2}[\stackrel{´}{𝐚}+\stackrel{´}{𝐛}],𝐱^{}(t,\sigma )=\frac{1}{2}[\stackrel{´}{𝐚}\stackrel{´}{𝐛}].$$
(3.1)
Here $`|\stackrel{´}{𝐚}|=1`$, while $`|\stackrel{´}{𝐛}|=k1`$. We shall generally assume that there is a nonzero current, so that $`k<1`$. It then follows that chiral current-carrying cosmic strings in Minkowski space cannot have zero velocity. (Thus stationary loops do not exist for example.) Similarly the tangent vector of the string never vanishes either so that there are no cusps on these strings.
From equations (3.1) it also follows that
$$\dot{𝐱}𝐱^{}=\frac{1}{4}[1k^2]>0.$$
so that the velocity of a point on the string is not perpendicular to its tangent vector. Notice that this result is due to the gauge conditions we have chosen: with the gauge choice of MS, $`\dot{𝐱}𝐱^{}=0`$.
Observe also that when $`k=0=|\stackrel{´}{𝐛}|`$, equation (3.1) implies that $`\dot{𝐱}=𝐱^{}`$. Since their velocity is always parallel to the tangent vector, these strings do not change their shapes, and thus never self-intersect. Furthermore, given that when $`k=0`$, $`\dot{𝐱}^2+𝐱^{}_{}{}^{}2=\frac{1}{2}`$ and $`\dot{𝐱}𝐱^{}=\frac{1}{4}`$, it follows that for $`k=0`$
$$|\dot{𝐱}|=\frac{1}{2}=|𝐱^{}|.$$
Thus the strings move at half the speed of light. Recall from (2.17) that when $`k=0`$ the strings carry the maximal charge.
### 3.2 Loops
We now turn to the properties of loops which must satisfy the periodicity conditions
$$𝐱(t,\sigma +L)=𝐱(t,\sigma ).$$
From (3.1), this implies, in the centre-of-mass frame,
$$𝐚(q+L)=𝐚(q),𝐛(\eta +L)=𝐛(\eta ).$$
It follows that, as in the NG case, the motion of chiral cosmic string loops is periodic, with period $`L/2`$ (because $`𝐱(t+L/2,\sigma +L/2)=𝐱(t,\sigma )`$). As noted above, for $`k=0`$ these loops do not self-intersect and hence are vorton solutions . The majority of studies of vortons to date have assumed that these are circular loops (see however ). The vortons with $`k=0`$ found here have entirely arbitrary shapes.
For arbitrary $`k`$ it is possible to construct solutions of (2.16) just as in the case of NG strings . For example, a 1-harmonic loop solution with constant $`k`$ is given by
$$𝐚(q)=(\mathrm{cos}q,\mathrm{sin}q,0);𝐛(\eta )=(k\mathrm{cos}\eta ,k\mathrm{sin}\eta ,0).$$
This is a circular string oscillating between maximum and minimum radii of $`(1+k)/2`$ and $`(1k)/2`$. (Such a solution was considered numerically in for arbitrary current carrying loops.) This loop never self-intersects for any value of $`k`$. Higher order harmonic solutions can also be constructed along very similar lines to references .
Since loops with $`k=0`$ never self-intersect, it is interesting to ask how the self-intersection probability of a loop with a given number of harmonics depends on $`k`$ (or equivalently on the conserved charge $`N`$ given in (2.17)).<sup>1</sup><sup>1</sup>1Whilst initially static NG loops always self-intersect , chiral cosmic strings can never be static as observed above. For simplicity, we will study this question for $`k(\varphi )`$ = constant. In that case, $`N`$ is given by
$$N=\frac{m}{4}L\sqrt{1k^2}$$
(3.2)
where $`L`$ is the invariant length of the loop. We now show that self-intersecting loop solutions exist for $`k>0`$ through the construction of an explicitly self-intersecting loop. Then the probabilities of self-intersection will be studied numerically for loops of a fixed invariant length $`L`$ but with different numbers of harmonics on them and different values of $`k`$.
The condition that a loop self-intersects at time $`T`$ is that there exists a solution of
$$𝐚(T+\sigma _1)+𝐛(T\sigma _1)=𝐚(T+\sigma _2)+𝐛(T\sigma _2)$$
(3.3)
for some $`0<\sigma _1\sigma _2<L`$. To show that self-intersection is possible, consider the following solutions for $`𝐚`$ and $`𝐛`$ that satisfy (2.16):
$`𝐚(q)`$ $`=`$ $`{\displaystyle \frac{1}{m}}(\mathrm{cos}mq,\mathrm{sin}mq,0)`$
$`𝐛(\eta )`$ $`=`$ $`{\displaystyle \frac{k}{n}}(\mathrm{cos}n\eta ,\mathrm{cos}\chi \mathrm{sin}n\eta ,\mathrm{sin}\chi \mathrm{sin}n\eta ),`$ (3.4)
where $`n`$ and $`m`$ have no common factors and $`\chi `$ is an arbitrary angle. Now let $`c=(\sigma _1+\sigma _2)/2`$, $`\delta =(\sigma _1\sigma _2)/2`$, $`q=T+c`$ and $`\eta =Tc`$. Then the self intersection condition (3.3) becomes
$$𝐚(q+\delta )𝐚(q\delta )=𝐛(\eta +\delta )𝐛(\eta \delta )$$
for which we must find solutions for $`\eta ,q,\delta `$. On substitution of (3.4), this condition becomes
$`{\displaystyle \frac{1}{m}}(\mathrm{sin}mq\mathrm{sin}m\delta ,\mathrm{cos}mq\mathrm{sin}m\delta ,0)`$
$`=`$ $`{\displaystyle \frac{k}{n}}(\mathrm{sin}n\eta \mathrm{sin}n\delta ,\mathrm{cos}n\eta \mathrm{sin}n\delta \mathrm{cos}\chi ,\mathrm{cos}n\eta \mathrm{sin}n\delta \mathrm{sin}\chi ).`$
Hence the requirement is that
$$\mathrm{cos}mq=\mathrm{cos}n\eta =0\mathrm{sin}n\eta =\pm 1=\pm \mathrm{sin}mq.$$
(3.5)
where $`\delta `$ must satisfy
$$\frac{\mathrm{sin}m\delta }{m}=\pm \frac{k}{n}\mathrm{sin}n\delta .$$
(3.6)
Generically, there are solutions to equations (3.5)–(3.6), and hence self-intersections.
### 3.3 Numerical study of loop self-intersections
More generally one can search for self-intersections numerically and try to determine the self-intersection probability as a function of $`k`$ and the number of harmonics on the loop. To do this, we used a modified version of the code written by Siemens and Kibble to search for self-intersections of NG loops. These authors built on work of Brown and DeLaney who devised a method of generating odd harmonic series satisfying a given constraint in terms of products of rotations. The only difference between the NG and chiral cosmic string loops is that for the former, the constraint is $`|\stackrel{´}{𝐛}|^2=1`$ whilst for the latter the constraint is $`|\stackrel{´}{𝐛}|^2=k^2<1`$. Thus here we carry out a simple extension of the work of Siemens and Kibble to study the self-intersection properties of $`M/P`$ harmonic loops (the notation means that there are $`M`$ harmonics in the solution of $`𝐚`$, and $`P`$ in the solution for $`𝐛`$.)
For more technical details on the code, the reader is referred to . In the results presented in figures 1-3 below, the rotation angles were given a uniform distribution, with the number of points along the string chosen to be $`K=600`$. This gives a resolution of 0.0104712 radians. The cutoff, below which self-intersection was not tested, was taken as 0.084 radians corresponding to 8 step lengths. These are the same parameters as those chosen in which have already been seen to work well. Furthermore, decreasing $`K`$ or increasing the cutoff did not affect our results.
The self-intersection probability was calculated for $`M/M`$ cosmic string loops as a function of $`k`$ (which corresponds to different charges on the loop through (3.2)). Figures 1-3 plot the intersection probability against $`\sqrt{1k^2}N`$ for these $`M/M`$ harmonic loops. (Error bars are one standard deviation.) Notice that for $`k=0`$ the loops do not self-intersect as was already proved above, whereas for a relatively large range of $`k`$ the probability is the same as the NG ($`k=1`$) case. Thus charges on chiral cosmic string loops appear only to have a significant effect on the dynamics of the loops when these charges are large. Indeed, if the loops are formed with large charges, they will scarcely ever intersect and this will lead to a cosmological catastrophe since the loops (vortons) will dominate the energy density of the universe. It therefore remains to understand whether or not these charges are expected to be large or small when the loops form: we leave a discussion of this question to the conclusions.
Finally, we note that while the plots show results for $`M/M`$ harmonic strings, we also ran the code for strings with different numbers of left and right-moving harmonics. This did not substantially change the intersection probability from that of a $`M/M`$ harmonic string if the harmonics were both close to $`M`$.
## 4 Conclusions and discussion
The basis for this paper was the action (2.1) for chiral cosmic strings first proposed in . This is a well defined, unique, action for strings carrying massless zero-modes which travel in one direction along the string at the speed of light. In section 2.2 we reproduced the results of showing how, with suitable gauge choices and treatment of the Lagrange multiplier $`\psi `$, the resulting equations of motion are integrable and reduce to the familiar wave equation with two constraints (2.16). These two constraints are that $`|\stackrel{´}{𝐚}|=1`$ and $`|\stackrel{´}{𝐛}|=k^21`$. We noted that the reason why $`|\stackrel{´}{𝐛}|`$ lies within the Kibble-Turok sphere rather than on it is that the chiral strings carry a conserved charge (associated with the current on them). In the limit of zero charge, the equations of motion and constraints (2.16) reduce to those of NG strings in the conformal-temporal gauge as required.
We placed a certain emphasis on gauge choices in section 2.2 since, as we showed in that section, the same action with less appropriate gauge choices can lead to much more complicated equations of motion which are not readily integrable, as for the equations of motion (2.10) derived in .
In section 3 we showed that chiral current-carrying cosmic strings cannot have cusps on them. Since cosmic rays are predominantly produced at cusps on NG strings , it is likely that a network of chiral cosmic strings will produce fewer cosmic rays.
We also showed that when the charge on the string is maximal (equivalently $`k=0`$), $`\dot{𝐱}=𝐱^{}`$ so that the strings move along themselves at half the speed of light and never self-intersect. In the case of loops these correspond to stationary vorton solutions of arbitrary shape. For infinite strings it means that these can never self-intersect to form loops (at least in Minkowski space). Since this is the main mechanism for removing energy from NG string networks, these results suggest that networks of chiral cosmic strings may evolve very differently from NG cosmic string networks.
In another step to study the evolution of chiral cosmic string networks, we considered the self-intersection probability of loops with $`0<k<1`$ and different numbers of harmonics (section 3). The results show that only when the charge on the loop is relatively large does the self-intersection probability differ significantly from the NG one.
As a result of this work, we are left with a number of important questions to study in the future. Maybe the most significant one of these is to understand what initial value of $`k`$ might be expected for the loops and infinite strings formed at the phase transition (which could be, say, at the end of inflation as discussed in the introduction). If $`k`$ is initially very small (i.e. the charge on the strings is initially close to being maximal) then chiral cosmic strings are already ruled out, as is the mixed scenario of D-term inflation and strings . The reason is that if the chiral cosmic strings effectively never self-intersect they rapidly come to dominate the energy density of the universe. Indeed, since the fermions are traveling in one direction only in the chiral case, the current and corresponding charge are larger than in the non-chiral case. Consequently, we would expect the charge to be close to maximal and hence $`k`$ to be small when the fermion zero modes condense on the string at formation. If however, the zero modes are formed at a subsequent phase transition, then $`k`$ is likely to be closer to unity. Indeed, this is the assumption made in , where theories giving rise to chiral cosmic strings were constrained by the requirement that they should not over produce vortons. We have arguments suggesting that $`k`$ is in fact initially small; these will be presented elsewhere .
Here we have restricted attention to loops in which $`k`$ is constant, but in fact one should also examine the more general case where $`k`$ is a function of $`\varphi `$, though always restricted to the range $`0k1`$.
Another objective would be to try to solve equations (2.16) in a very similar way to the Smith-Vilenkin algorithm which is an exact numerical algorithm for solving the corresponding equations for NG strings in Minkowski space. However, the Smith-Vilenkin algorithm is no longer exact for the chiral string equations: because of the constraint $`|\stackrel{´}{𝐛}|=k<1`$, the vertices will generally not remain on the lattice as the system evolves. Indeed we believe that there is no value of $`k`$ for which the algorithm can be made to work — except perhaps $`k=0`$, a case that is uninteresting in this context as we already know that there the strings are effectively stationary.
Finally, one should also consider to what degree the effects of friction on the evolving chiral cosmic string network are important. Frictional effects on NG and chiral strings are likely to be similar since there are no long range forces in either case; this is unlike the situation for electromagnetically coupled strings . Ultimately the effect of expansion should be incorporated too, though as in the NG case many predictions can be made from Minkowski space results . We are currently studying a number of these questions . Our general conclusion of this paper would be, however, that we have found evidence to suggest that chiral cosmic string networks evolve very differently from NG networks. Hence their cosmological consequences will be very different, and so some caution should be used before simply adding the effects of inflation and NG strings as in , especially when the specific model under consideration actually produces chiral cosmic strings.
## Acknowledgements
We thank B. Carter and P. Peter for useful discussions. This work is supported in part by PPARC, UK and by an ESF network.
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# Dispersion and Symmetry of Bound States in the Shastry-Sutherland Model
\[
## Abstract
Bound states made from two triplet excitations on the Shastry-Sutherland (ShaSu) lattice are investigated. Based on the perturbative unitary transformation by flow equations quantitative properties like dispersions and qualitative properties like symmetries are determined. The high order results (up to $`(J_2/J_1)^{14}`$) permit to fix the parameters of SrCu<sub>2</sub>(BO<sub>3</sub>)<sub>2</sub> precisely: $`J_1=6.16(10)`$meV, $`x:=J_2/J_1=0.603(3)`$, $`J_{}=1.3(2)`$meV. At the border of the magnetic Brillouin zone (MBZ) a general double degeneracy is derived. An unexpected instability in the triplet channel at $`x=0.63`$ indicates a first order transition towards a triplet condensate, related to classical helical order.
Dedicated to Prof. F. Wegner on occasion of his 60<sup>th</sup> birthday.
\]
Quantum antiferromagnets are at the center of research not only because of the high $`T_c`$ superconductors . Of particular interest are systems which do not have an ordered, Néel-type ground state. Their ground state is a spin liquid without long range spin order. Spin liquids are favored by low spin ($`S=\frac{1}{2}`$ mostly), low coordination number ($`Z\{2,3,4\}D\{1,2\}`$), and strong geometric frustration.
Dimer solids are transparent cases of spin liquids. In $`D=1`$, the generic example is the Majumdar-Ghosh model of which Shastry and Sutherland found a $`D=2`$ generalization (ShaSu model) . In both cases frustration is essential. Each spin is coupled to pairs of spins (dimers). If these pairs form singlets the couplings between dimers is without effect and the singlet-on-dimers product state is always an eigen state and for certain parameters the ground state . The systems are gapped. The elementary excited states are dressed $`S=\frac{1}{2}`$ ($`D=1`$) or $`S=1`$ ($`D=2`$) entities. They interact strongly and form bound and antibound states in various spin channels.
Due to its recent realization in SrCu<sub>2</sub>(BO<sub>3</sub>)<sub>2</sub> the ShaSu model (Fig. 1) is presently attracting enormous interest. The Hamiltonian reads
$$H(J_1,J_2)=J_1\underset{𝐢,𝐣\mathrm{dimer}}{}𝐒_𝐢𝐒_𝐣+J_2\underset{𝐢,𝐣\mathrm{square}}{}𝐒_𝐢𝐒_𝐣.$$
(1)
In this Letter we start from the dimer phase . We focus on bound states formed from pairs of the elementary triplets and their symmetries, degeneracies, and dispersion. The perturbative unitary transformation based on flow equations enables us to link smoothly and uniquely $`H(J_1,J_2)`$ at $`x:=J_2/J_10`$ to an effective $`H_{\mathrm{eff}}`$ conserving the number of triplets on dimers $`[H_{\mathrm{eff}},H(J_1,0)]=0`$. This permits a clear distinction between the ground state (without triplets), the 1-triplet sector, the 2-triplet sector etc..
In terms of $`H_{\mathrm{eff}}`$ the dynamics of one triplet is hopping $`t_{h;𝐢}`$ ($`t_{v;𝐢}`$) starting from a horizontal (vertical) dimer by $`i_x`$ dimers right and $`i_y`$ dimers up. Nothing else is possible due to triplet number conservation. The elements $`t`$ are computed in order 15 .
The dynamics of two triplets at large distances is governed by 1-triplet hopping. At smaller distances a 2-particle interaction occurs additionally given by $`W_{h;𝐝;𝐢,𝐝^{}}`$ ($`W_{v;𝐝;𝐢,𝐝^{}}`$) starting with one triplet on a horizontal (vertical) dimer and another at distance $`𝐝`$. The action of $`H_{\mathrm{eff}}`$ is to shift the triplets to $`𝐢`$ and to $`𝐢+𝐝^{}`$. Nothing else is possible due to triplet number conservation. Since the total spin is conserved ($`S\{0,1,2\}`$) the distances are restricted to $`𝐝,𝐝^{}>\mathrm{𝟎}`$, i.e. $`d_x>0`$ or $`d_x=0d_y>0`$, because the exchange parity is fixed.
The action of $`H_{\mathrm{eff}}`$ yields the combined effect of hopping and interaction denoted by $`A_{𝐝;𝐢,𝐝^{}}`$. The true 2-triplet interaction is easily found by subtracting the 1-triplet hopping
$`W_{𝐝;\mathrm{𝟎},𝐝^{}}`$ $`=`$ $`A_{𝐝;\mathrm{𝟎},𝐝^{}}t_{𝐝^{}𝐝}\delta _{𝐝^{},𝐝}t_\mathrm{𝟎}`$ (3)
$`W_{𝐝;𝐝𝐝^{},𝐝^{}}`$ $`=`$ $`A_{𝐝;𝐝𝐝^{},𝐝^{}}t_{𝐝𝐝^{}}\delta _{𝐝^{},𝐝}t_\mathrm{𝟎}`$ (4)
$`W_{𝐝;𝐝^{},𝐝^{}}`$ $`=`$ $`A_{𝐝;𝐝^{},𝐝^{}}t_{𝐝𝐝^{}}`$ (5)
$`W_{𝐝;𝐝,𝐝^{}}`$ $`=`$ $`A_{𝐝;𝐝,𝐝^{}}t_{𝐝+𝐝^{}}`$ (6)
(distinction $`h/v`$ omitted for clarity). Otherwise $`A`$ and $`W`$ are identical. The coefficients $`W`$ for $`S\{0,1,2\}`$ yield the complete 2-particle dynamics. We compute $`W`$ up to $`x^{12}`$, the coefficients for the lowest-lying states even up to $`x^{14}`$.
During the virtual processes the triplet number is changed. Due to the frustration of the ShaSu lattice $`H`$ in (1) changes the number of triplets on the dimers at most by one . An excitation or a de-excitation on a horizontal (vertical) dimer is possible iff at least one of the vertical (horizontal) dimers on the left and right (above and below) are excited. This restriction implies that one triplet hops only in $`x^6`$ (cf. Figs. in ).
Motion of two triplets together is much less restricted (cf. Fig. 2). Matrix elements occur in $`x^2`$ as first observed for total spin $`S=2`$ . But the dispersion of bound states starts only in $`x^3`$ (contrary to $`x^4`$ claimed in Ref. ). Two adjacent triplets interact linearly in $`x`$ ($`x`$ for $`S=0`$, $`x/2`$ for $`S=1`$, $`x/2`$ for $`S=2`$). The energy of the initial and final state in each row in Fig. 2 differ by $`𝒪(x)`$. Hence both rows must be combined making it an $`(x^2)^2/x=x^3`$ process eventually. This applies to the 8 (anti)bound states derived from two triplets on (next) nearest neighbor dimers. The dispersion of any other state sets in at higher order.
We use the following basis for the 2-triplet states
$$|𝐤,𝐝,\sigma :=N^{1/2}\underset{𝐫}{}e^{\left(i(𝐤+\sigma 𝐐)(𝐫+𝐝/2)\right)}|𝐫,𝐫+𝐝,$$
(7)
where $`𝐤`$ is the conserved total momentum in the magnetic Brillouin zone (MBZ) applying due to the two sublattices; $`\sigma \{0,1\}`$, $`𝐐:=(\pi ,\pi )`$, $`N`$ is the number of dimers, $`|𝐫,𝐫+𝐝`$ denotes the state with triplets at $`𝐫`$ and at $`𝐫+𝐝`$. The distance $`𝐝`$ is restricted $`𝐝>\mathrm{𝟎}`$, i.e. $`d_x>0`$ or $`d_x=0d_y>0`$. The matrix elements of $`H_{\mathrm{eff}}`$ in the basis (7) are real due to translation and inversion symmetry.
Before the quantitative analysis of the bound states a qualitative aspect, a general double degeneracy at the border of the MBZ, shall be derived. To see this consider the combination of a shift by the dimer-dimer spacing along a’ (S), a reflection about a (R), and the inversion $`𝐫𝐫`$ (I) (cf. Fig. 1). The combinations SR and I are symmetries of the Hamiltonian. For $`k_x+k_y=\pi `$ (part of the MBZ border) definition (7) implies for the total combination SRI the mapping
$$|𝐤,𝐝,\sigma e^{ik_x+i\pi (d_x+d_y+\sigma +PS)}|𝐤,(\begin{array}{c}d_y\\ d_x\end{array}),1\sigma ,$$
(8)
where $`P\{0,1\}`$ being unity iff $`(d_y,d_x)<\mathrm{𝟎}`$ so that the triplets must be swapped to pass from $`(d_y,d_x)`$ to $`(d_y,d_x)`$. It is crucial that SRI links $`|𝐤,(d_x,d_y),0`$ and $`|𝐤,(d_y,d_x),1`$ like a 2D rotation $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ up to a prefactor. Hence its eigen vectors are complex with linearly independent real and imaginary part and so are the simultaneous eigen vectors of SRI and $`H_{\mathrm{eff}}`$. Because $`H_{\mathrm{eff}}`$ is real the real and the imaginary part constitute in fact linearly independent eigen vectors to the same eigen value. The same double degeneracy is concluded for the other parts of the MBZ border by S and 90 rotation (D). It is also valid in the 1-triplet sector .
The double degeneracy at the MBZ border is interesting for analysing experiment, too. Degeneracy reduces the large number of energetically close states helping to resolve different bound states.
Since 1-triplet hopping is of higher order than interaction an analytic expansion for the energies of the bound states is possible. At finite order in $`x`$ only configurations contribute where the two triplets are not too far away from each other. Of course, higher orders imply larger, but still finite distances. In particular, the energies of the four states which evolve from neighboring triplets can be computed very well since their interaction is linear. Investigating the matrix elements shows that it is sufficient to study the distances $`𝐝\{(0,1),(1,0),(1,\pm 1)\}`$ for order 5. To $`x^{14}`$ only $`𝐝\{(1,\pm 2),(2,\pm 1),(0,2),(2,0),(2,\pm 2)\}`$ must be added. So, for given total momentum only a finite $`8\times 8`$ or $`24\times 24`$ matrix has to be analysed. For illustration consider the elements $`A_{(0,1);𝐢,(2,1)}`$ (the Fourier transform of $`𝐢`$ yields the momentum dependence.) connecting $`(0,1)`$ and $`(2,1)`$ which is $`𝒪(x^4)`$. By second order perturbation one sees that the resulting energy shift is $`(x^4)^2/x=x^7`$ only.
Furthermore, the elements connecting shorter distances to longer distances and the elements among longer distances do not need to be known to very high order. Consider again the process $`(0,1)(2,1)`$. In order $`x^7`$ the element $`A_{(0,1);𝐢,(2,1)}`$ must be known only in $`x^4`$ and $`A_{(2,1);𝐢,(2,1)}`$ only in $`x^1`$; in order $`x^9`$ the element $`A_{(0,1);𝐢,(2,1)}`$ must be known only up to $`x^6`$ and $`A_{(2,1);𝐢,(2,1)}`$ only in $`x^3`$ and so on.
We have analysed the dispersions in $`x^5`$ of the four states bound linearly in $`x`$ in the MBZ. Fukumoto’s results are mostly confirmed . At particular points of high point group symmetry ($`(0,0)`$,$`(0,\pi )`$,$`(\pi /2,\pi /2)`$) the Hamiltonian splits into six blocks corresponding to different representations of the square point group 4mm. At these points the analysis up to $`x^{14}`$ is carried out . The symmetries are classified according to the irreducible representations (four 1D, one 2D) of the point group 4mm $`\mathrm{\Gamma }_1(1),\mathrm{\Gamma }_2(x^2y^2),\mathrm{\Gamma }_3(xy),\mathrm{\Gamma }_4(xy(x^2y^2)),\mathrm{\Gamma }_5(x,y)`$ where simple polynomials are given in brackets to show the transformation behavior.
The extrapolated energies are depicted in Figs. 3 ($`S=0`$) and 4 ($`S=1`$) as functions of $`x`$. For those energies which stay separated from the 2-particle continuum Dlog-Padé approximants are used successfully . The results are stable under changes of the polynomial degrees. The energies close to the continuum (here simply twice the gap $`\mathrm{\Delta }`$ between the ground state and the elementary triplet at $`𝐤=\mathrm{𝟎}`$) are given with less reliability by the truncated series or by a non-defective Dlog-Padé approximant.
In Figs. 3, 4 the modes are sorted in energetically ascending order for small values of $`x`$: solid, long dashed, and short dashed curves. The $`\mathrm{\Gamma }_5`$ modes are naturally degenerate. The double degeneracy for $`𝐤=(0,\pi )`$ does not result from the point group but originates from the complex conjugation as explained above. The dashed-dotted curve at $`(0,\pi )`$ has to be compared to the solid and the long-dashed curve to assess the dispersion of these two modes from $`\mathrm{𝟎}`$ to $`(0,\pi )`$. While for $`S=0`$ this dispersion always has the expected behavior with $`\omega (\mathrm{𝟎})<\omega ((0,\pi ))`$ the energies for $`S=1`$ are reversed for small values of $`x`$ (cf. ). Only above $`x0.55`$ the relation $`\omega (\mathrm{𝟎})<\omega ((0,\pi ))`$ holds for $`S=1`$.
We do not agree with Ref. that the two lowest states are of $`s`$-wave type since this would imply that they are $`\mathrm{\Gamma }_1`$. Instead the $`S=0,1`$ states are odd under reflection about a’ (R1) or about b’ (R2) (cf. Fig. 1). For $`S=0`$, the lowest state is even under SD and the second lowest is odd. For $`S=1`$, it is vice-versa. The $`\mathrm{\Gamma }_5`$ states can be viewed as being of $`p`$-wave type.
For $`S=0`$, the lowest mode vanishes at the same $`x`$ as does the gap $`\mathrm{\Delta }`$. So no additional instability occurs for $`S=0`$. This provides evidence against a competing singlet phase as presumed in Ref. . There is, however, a salient instability for $`S=1`$ (Fig. 4) at $`x=0.63`$. This comes as a surprise since one expects in antiferromagnets binding effects to be largest for $`S=0`$. The singularity at $`x=0.630(5)`$ is very stable occuring in all non-defective Dlog-Padé approximants. We take the vanishing of a bound 2-triplet state with $`S=1`$ as evidence for strong attraction between triplets which are neither parallel ($`S=2`$) nor antiparallel ($`S=0`$). The attraction points towards a first order transition into a condensate of triplets occuring at much lower $`x`$ than thought previously . The angle being neither zero nor $`\pi `$ corroborates a helical order if one adopts a classical view . In the light of the instability for $`S=1`$ we interprete the findings by Koga and Kawakami as indication of the same instability. It is so far not excluded that the bulk triplet condensate (the helical phase) is a singlet or that it can be linked to a singlet .
Next we determine $`x`$ and $`J_1`$ for SrCu<sub>2</sub>(BO<sub>3</sub>)<sub>2</sub>. The steep decrease of the bound $`S=1`$ state enables us to fix $`x`$ very precisely. Based on ESR , FIR as well as INS we assume $`\mathrm{\Delta }=2.98`$meV and $`\omega |_{S=1}=4.7`$meV leading to $`x=0.603(3)`$ and $`J_1=6.16(10)`$meV. The 1-triplet dispersion is in excellent agreement with experiment (cf. inset in Fig. 3 and Ref. ). Raman scattering provides further strong support because the energy of the $`\mathrm{\Gamma }_3`$ singlet matches 30cm<sup>-1</sup> perfectly. The $`\mathrm{\Gamma }_4`$ singlet at 25cm<sup>-1</sup> is forbidden by symmetry since the Raman operator is effectively $`\mathrm{\Gamma }_3`$.
In leading order $`t/U`$ the Raman operator $`R=\gamma _{𝐢,𝐣}𝐒_𝐢𝐒_𝐣`$ couples the same spins as the Hamiltonian. But only the antisymmetric ($`\gamma _1=\gamma _2`$) part of $`R`$ on the dashed bonds (cf. Fig. 1) creates excitations from the ground state. By geometry we have $`\gamma _1^{\prime \prime }=\gamma _2`$ and $`\gamma _2^{\prime \prime }=\gamma _1`$ so that the effective $`R_{\mathrm{eff}}`$ is odd under R1 and R2. But the projection of the vector potential $`𝐀`$ ($`𝐀||𝐄`$) on the bonds under study in the polarisation (ab) implies $`\gamma _1=\gamma _2=\gamma _1^{}=\gamma _2^{}=\gamma \mathrm{cos}(2\alpha )`$ ($`\gamma `$ microscopic constant, $`\alpha `$ angle in Fig. 1), i.e. an even component implying $`R_{\mathrm{eff}}=0`$. On the contrary, polarisation (a’b’) yields $`\gamma _1=\gamma _2=\gamma _1^{}=\gamma _2^{}=\gamma \mathrm{sin}(2\alpha )`$ so that $`R_{\mathrm{eff}}0`$. This finding agrees nicely with experiment where on $`T0`$ the intensities almost vanish for (ab) but grow for (a’b’) . Additionally, $`\gamma _1=\gamma _1^{}`$ and $`\gamma _2=\gamma _2^{}`$ imply odd parity under SD so that $`R_{\mathrm{eff}}`$ is indeed $`\mathrm{\Gamma }_3`$, not $`\mathrm{\Gamma }_4`$. Calculating the next $`\mathrm{\Gamma }_3,S=0`$ bound state (less systematically) yields 45cm<sup>-1</sup> in good agreement with the experimental 46cm<sup>-1</sup> line, too.
We conclude that the 2D model (1) explains the low-lying excitations of SrCu<sub>2</sub>(BO<sub>3</sub>)<sub>2</sub> perfectly. Thermodynamic quantities like the susceptibility $`\chi (T)`$ require the inclusion of the interplain coupling $`J_{}`$ which is fully frustrated not changing the dimer spins . We have employed a Dlog-Padé approximant for the $`1/T`$ expansion of the 2D $`\chi _{2\mathrm{D}}`$ complemented by the condition $`\mathrm{\Delta }=2.98`$meV. This ansatz works fine for $`T>35`$K. The 3D $`\chi _{3\mathrm{D}}`$ is computed from $`\chi _{2\mathrm{D}}`$ on the mean-field level as $`\chi _{3\mathrm{D}}^1=\chi _{2\mathrm{D}}^1+4J_{}`$. The inset in Fig. 4 shows that theory ($`J_{}=1.3(2)`$meV) and experiment agree without flaw above 40K. Our value for $`J_{}`$ is significantly higher than the one in Ref. due to different values of $`x`$ and $`J_1`$.
The above comprehensive analysis of bound states is a fine example of the efficiency of perturbation by unitary transformation based on flow equations. This clear concept allows to distinguish uniquely sectors with different particle numbers and other different quantum numbers like the total spin. Here the concept was put to use to analyse the Shastry-Sutherland lattice as realized in SrCu<sub>2</sub>(BO<sub>3</sub>)<sub>2</sub>. To our knowledge it is the first quantitative description of 2-particle bound states in 2D.
The symmetries of experimentally relevant states were determined. The reliability of the high order results allows to fix the experimental coupling constants very precisely ($`J_1=6.16(10)`$mev, $`J_2/J_1=0.603(3)`$, $`J_{}=1.3(2)`$meV). Thereby, different experiments (ESR, FIR, INS, Raman, $`\chi (T)`$) are explained consistently. We suggest to exploit the double degeneracy derived here to resolve different bound states at the border of the MBZ.
An unexpected instability for the $`S=1`$ 2-triplet bound state is found at $`x0.63`$ indicating a transition to a triplet condensate probably related to the helical phase found previously . We conjecture that this transition is first order occuring at lower $`x`$ than assumed so far. The vicinity of SrCu<sub>2</sub>(BO<sub>3</sub>)<sub>2</sub> to this transition suggests to attempt a closer experimental analysis. Pressure and/or substitution will certainly influence the ratio $`J_2/J_1`$. Thereby one may hope to scan through the transition and to examine the phase beyond.
The authors like to thank H. Kageyama and P. Lemmens for generous provision of data prior to publication and discussion. The work is supported by the DFG in SFB 341 and in SP 1073.
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# Detecting the MSSM Higgs Bosons at Future 𝑒⁺𝑒⁻ Colliders
## I Introduction
Higgs bosons play an imporant role in the Standard Model (SM) ; they are responsible for generating the masses of all the elementary particles (leptons, quarks, and gauge bosons). However, the Higgs-boson sector is the least tested one in the SM. If Higgs bosons are responsible for breaking the symmetry from $`SU(2)_L\times U(1)_Y`$ to $`U(1)_{EM}`$, it is natural to expect that other Higgs bosons are also involved in breaking other symmetries at the grand-unification scale. One of the more attractive extensions of the SM is Supersymmetry (SUSY) , mainly because of its capacity to solve the naturalness and hierarchy problems while maintaining the Higgs bosons elementary.
The minimal supersymmetric extension of the Standard Model (MSSM) doubles the spectrum of particles of the SM and the new free parameters obey simple relations. The scalar sector of the MSSM requires two Higgs doublets, thus the remaining scalar spectrum contains the following physical states: two CP-even Higgs scalar ($`h^0`$ and $`H^0`$) with $`m_{h^0}m_{H^0}`$, one CP-odd Higgs scalar ($`A^0`$) and a charged Higgs pair ($`H^\pm `$), whose detection would be a clear signal of new physics. The Higgs sector is specified at tree-level by fixing two parameters, which can be chosen as the mass of the pseudoscalar $`m_{A^0}`$ and the ratio of vacuum expectation values of the two doublets $`\mathrm{tan}\beta =v_2/v_1`$, then the mass $`m_{h^0}`$, $`m_{H^0}`$ and $`m_{H^\pm }`$ and the mixing angle of the neutral Higgs sector $`\alpha `$ can be fixed. However, since radiative corrections produce substantial effects on the predictions of the model , it is necessary to specify also the squark masses, which are assumed to be degenerated. In this paper, we focus on the phenomenology of the neutral CP-even and CP-odd scalar ($`h^0,H^0,A^0`$) and charged $`(H^\pm )`$.
The search for these scalars has begun at LEP, and current low energy bound on their masses gives $`m_{h^0}`$, $`m_{A^0}`$ $`>`$ 90 $`GeV`$ and $`m_{H^\pm }`$ $`>`$ 120 $`GeV`$ for $`\mathrm{tan}\beta `$ $`>`$ 1 . At $`e^+e^{}`$ colliders the signals for Higgs bosons are relatively clean and the opportunities for discovery and detailed study will be excellent. The most important processes for the production and detection of the neutral and charged Higgs bosons $`h^0`$, $`H^0`$, $`A^0`$ and $`H^\pm `$, are: $`e^+e^{}Z^{}h^0,H^0+Z^0`$, $`e^+e^{}Z^{}h^0,H^0+A^0`$, $`e^+e^{}\nu \overline{\nu }+W^+W^{}\nu \overline{\nu }+h^0,H^0`$ (the later is conventionally referred to as $`WW`$ fusions), and $`e^+e^{}H^+H^{}`$ ; precise cross-section formulas appear in Ref. . The main decay modes of the neutral Higgs particles are in general $`b\overline{b}`$ decays $`(90\%)`$ and $`\tau ^+\tau ^{}`$ decays $`(10\%)`$ which are easy to detect experimentally at $`e^+e^{}`$ colliders . Charged Higgs particles decay predominantly into $`\tau \nu _\tau `$ and $`t\overline{b}`$ pairs.
The $`Z^0h^0`$ production cross-section contains an overall factor $`\mathrm{sin}^2(\beta \alpha )`$ which suppress it in certain parameter regions (with $`m_{A^0}<100`$ $`GeV`$ and $`\mathrm{tan}\beta `$ large); fortunately the $`A^0h^0`$ production cross-section contains the complementary factor $`\mathrm{cos}^2(\beta \alpha )`$. Hence the $`Z^0h^0`$ and $`A^0h^0`$ channels together are well suited to cover all regions in the $`(m_{A^0}\mathrm{tan}\beta )`$ plane, provided that the $`c.m.`$ energy is high enough for $`Z^0h^0`$ to be produced through the whole $`m_{h^0}`$ mass range, and that an adequate event rate can be achieved. These conditions are already shown to be satisfied for $`\sqrt{s}=500`$ $`GeV`$ with assumed luminosity $`10`$ $`fb^1`$, as is expected to be the case of the Next Linear $`e^+e^{}`$ Collider (NLC).
In previous studies, the two-body processes $`e^+e^{}h^0(H^0)+Z^0`$ and $`e^+e^{}h^0(H^0)+A^0`$ have been evaluated extensively. However, the inclusion of three-body process $`e^+e^{}h^0(H^0)+b\overline{b}`$ and $`e^+e^{}A^0+b\overline{b}`$ at future $`e^+e^{}`$ colliders energies is necessary in order to know its impact on the two-body mode processes and also to search for new relations that could have a cleaner signature of the Higgs bosons production.
In the other hand the decay modes of the charged Higgs bosons determine the signatures in the detector. If $`m_{H^\pm }>m_t+m_b`$, the dominant decays modes are $`t\overline{b}`$, $`\overline{t}b`$ and $`\tau ^+\nu _\tau `$, $`\tau ^{}\overline{\nu }_\tau `$. In some part of the parameter space also the decay $`H^+W^+h^0`$ is allowed. If $`m_{H^\pm }<m_t+m_b`$, the charged Higgs boson will decay mainly into $`\tau ^+\nu _\tau `$, $`\tau ^{}\overline{\nu }_\tau `$.
For $`m_{A^0}m_{Z^0}`$, and if 50 events criterion are adecuate, the $`H^+H^{}`$ pair production will be kinematically allowed and easily observable . For $`m_{A^0}>120`$ $`GeV`$, $`e^+e^{}H^+H^{}`$ must be employed for detection of the three heavy Higgs bosons. Assuming that SUSY decays are not dominant, and using the 50 event criterion $`H^+H^{}`$ can be detected up to $`m_{H^\pm }=230`$ $`GeV`$ , assuming $`\sqrt{s}=500`$ $`GeV`$.
The upper limits in the $`H^+H^{}`$ mode are almost entirely a function of the machine energy (assuming an appropriately higher integrated luminosity is available at a higher $`\sqrt{s}`$). Two recent studies show that at $`\sqrt{s}=1`$ $`TeV`$, with an integrated luminosity of 200 $`fb^1`$, $`H^+H^{}`$ detection would extended to $`m_{A^0}m_{H^\pm }450`$ $`GeV`$ even if substantial SUSY decays of these heavier Higgs are present.
In the present paper we study the production of SUSY Higgs bosons at $`e^+e^{}`$ colliders. We are interested in finding regions that could allow the detection of the SUSY Higgs bosons for the set parameter space $`(m_{A^0}\mathrm{tan}\beta )`$. We shall discuss the neutral and charged Higgs bosons production $`b\overline{b}h^0(H^0,A^0)`$, and $`\tau ^{}\overline{\nu }_\tau H^+`$, $`\tau ^+\nu _\tau H^{}`$ in the energy range of a future $`e^+e^{}`$ colliders for large values of the parameter $`\mathrm{tan}\beta `$, where one expects to have a high production. Since the coupling $`h^0b\overline{b}`$ is proportional to $`\mathrm{sin}\alpha /\mathrm{cos}\beta `$, the cross-section will receive a large enhancement factor when $`\mathrm{tan}\beta `$ is large. Similar situation occurs for $`H^0`$, whose coupling with $`b\overline{b}`$ is proportional to $`\mathrm{cos}\alpha /\mathrm{cos}\beta `$. The couplings of $`A^0`$ with $`b\overline{b}`$ and of $`H^\pm `$ with $`\tau ^{}\overline{\nu }_\tau ,\tau ^+\nu _\tau `$ are directly proportional to $`\mathrm{tan}\beta `$, thus the amplitudes will always grow with $`\mathrm{tan}\beta `$. We consider the complete set of Feynman diagramas at tree-level and use the helicity formalism for the evaluation of the amplitudes. The results obtained for the three-body processes are compared with the dominant two-body mode reactions for the plane $`(m_{A^0}\mathrm{tan}\beta )`$. Succintly, our aim in this work is to analyze how much the results of the Bjorken Mechanism \[Fig. 1, (1.4)\] would be enhanced by the contribution from the diagrams depicted in Figs. 1.1-1.3, 1.5 and 1.6 in which the SUSY Higgs bosons are radiated by a $`b(\overline{b})`$ quark. For the case of the charged Higgs bosons the two-body mode \[Figs. 3.1 and 3.4\] would be enhanced by the contribution from the diagrams depicted in Figs. 3.2, 3.3, and 3.5, in which the charged Higgs boson is radiated by a $`\tau ^{}\overline{\nu }_\tau `$ $`(\tau ^+\nu _\tau )`$ lepton.
Recently, it has been shown that for large values of $`\mathrm{tan}\beta `$ the detection of SUSY Higgs bosons is possible at FNAL and LHC . In the papers cited in Ref. the authors calculated the corresponding three-body diagrams for hadron collisions. They pointed out the importance of a large bottom Yukawa coupling at hadron colliders and showed that the Tevatron collider may be a good place for detecting SUSY Higgs bosons. In the case of the hadron colliders the three-body diagrams come from gluon fusion and this fact makes the contribution from these diagrams more important, due to the gluon abundance inside the hadrons. The advantage for the case of $`e^+e^{}`$ colliders is that the signals of the processes are cleaner.
This paper is organized as follows. We present in Sec. II the relevant details of the calculations. Sec. III contains the results for the processes $`e^+e^{}b\overline{b}h^0(H^0,A^0)`$ and $`e^+e^{}\tau ^{}\overline{\nu }_\tau H^+,\tau ^+\nu _\tau H^{}`$ at future $`e^+e^{}`$ colliders. Finally, Sec. IV contains our conclusions.
## II Helicity Amplitude for Higgs Bosons Production
When the number of Feynman diagrams is increased, the calculation of the amplitude is a rather unpleasant task. Some algebraic forms can be used in it to avoid manual calculation, but sometimes the lengthy printed output from the computer is overwhelming, and one can hardly find the required results from it. The CALKUL collaboration suggested the Helicity Amplitude Method (HAM) which can simplify the calculation remarkably and hence make the manual calculation realistic.
In this section we describe in brief the evaluation of the amplitudes at tree-level, for $`e^+e^{}b\overline{b}h^0(H^0,A^0)`$ and $`e^+e^{}\tau ^{}\overline{\nu }_\tau H^+,\tau ^+\nu _\tau H^{}`$ using the HAM . This method is a powerful technique for computing helicity amplitudes for multiparticle processes involving massles spin-1/2 and spin-1 particles. Generalization of this method that incorporates massive spin-1/2 and spin-1 particles, is given in Ref. . This algebra is easy to program and more efficient than computing the Dirac algebra.
A Higgs boson $`h^0,H^0`$, $`A^0`$, and $`H^\pm `$ can be produced in scattering $`e^+e^{}`$ via the following processes:
$`e^+e^{}`$ $``$ $`b\overline{b}h^0,`$ (1)
$`e^+e^{}`$ $``$ $`b\overline{b}H^0,`$ (2)
$`e^+e^{}`$ $``$ $`b\overline{b}A^0,`$ (3)
$`e^+e^{}`$ $``$ $`\tau ^{}\overline{\nu }_\tau H^+,\tau ^+\nu _\tau H^{}.`$ (4)
The diagrams of Feynman, which contribute at tree-level to the different reaction mechanisms, are depicted in Figs. 1-3. Using the Feynman rules given by the Minimal Supersymmetric Standard Model (MSSM), as summarized in Ref. , we can write the amplitudes for these reactions . For the evaluation of the amplitudes we have used the spinor-helicity technique of Xu, Zhang and Chang (denoted henceforth by XZC), which is a modification of the technique developed by the CALKUL collaboration . Following XZC, we introduce a very useful notation for the calculation of the processes (1)-(3) and (4). The complete formulas of the processes (1)-(3) are given in Ref. . For the case of process (4), we present the relevant details of the calculations.
Let us consider the process
$$e^{}(p_1)+e^+(p_2)\{\tau ^{}(k_2)+\overline{\nu }_\tau (k_3)+H^+(k_1),\tau ^+(k_2)+\nu _\tau (k_3)+H^{}(k_1)\},$$
(5)
in which the helicity amplitude is denoted by $`[\lambda (e^{}),\lambda (e^+),\lambda (\tau ^{}),\lambda (\nu _\tau )]`$.
Due to charge invariance, the cross-sections for the production processes $`e^{}e^+\tau ^{}\overline{\nu }_\tau H^+`$ and $`e^{}e^+\tau ^+\nu _\tau H^{}`$ are exactly the same. One should thus calculate the cross-section of one of the two reactions and then multiply by a factor of two to take into account the charge conjugate final state. This will enormously simplify the Feynman diagrams as well as the amplitudes of transition. The Feynman diagrams for this process are shown in Fig. 3. From this figure it follows that the amplitudes that correspond to each graph are
$`_1`$ $`=`$ $`iC_1P_H^{}(k_2+k_3)P_Z(p_1+p_2)T_1,`$ (6)
$`_2`$ $`=`$ $`iC_2P_\tau (k_1+k_3)P_Z(p_1+p_2)T_2,`$ (7)
$`_3`$ $`=`$ $`iC_3P_\nu (k_1+k_2)P_Z(p_1+p_2)T_3,`$ (8)
$`_4`$ $`=`$ $`iC_4P_H^{}(k_2+k_3)P_\gamma (p_1+p_2)T_4,`$ (9)
$`_5`$ $`=`$ $`iC_5P_\tau (k_1+k_3)P_\gamma (p_1+p_2)T_5,`$ (10)
where
$`C_1`$ $`=`$ $`{\displaystyle \frac{g^3}{16\sqrt{2}}}{\displaystyle \frac{m_\tau }{m_W}}\mathrm{tan}\beta {\displaystyle \frac{\mathrm{cos}2\theta _W}{\mathrm{cos}^2\theta _W}},`$ (11)
$`C_2`$ $`=`$ $`{\displaystyle \frac{g^3}{32\sqrt{2}}}{\displaystyle \frac{m_\tau }{m_W}}\mathrm{tan}\beta {\displaystyle \frac{1}{\mathrm{cos}^2\theta _W}},`$ (12)
$`C_3`$ $`=`$ $`C_2,`$ (13)
$`C_4`$ $`=`$ $`{\displaystyle \frac{g^3}{2\sqrt{2}}}{\displaystyle \frac{m_\tau }{m_W}}\mathrm{tan}\beta \mathrm{sin}^2\theta _W,`$ (14)
$`C_5`$ $`=`$ $`C_4,`$ (15)
while the propagators are
$`P_{Z^0}(p_1+p_2)`$ $`=`$ $`{\displaystyle \frac{(sm_{Z^0}^2)+im_{Z^0}\mathrm{\Gamma }_{Z^0}}{(sm_{Z^0}^2)^2+(m_{Z^0}\mathrm{\Gamma }_{Z^0})^2}},`$ (16)
$`P_{H^\pm }(k_2+k_3)`$ $`=`$ $`{\displaystyle \frac{(2k_2k_3m_{H^\pm }^2)+im_H\mathrm{\Gamma }_{H^\pm }}{(2k_2k_3m_{H^\pm }^2)^2+(m_{H^\pm }\mathrm{\Gamma }_{H^\pm })^2}},`$ (17)
$`P_\tau (k_1+k_3)`$ $`=`$ $`{\displaystyle \frac{1}{m_{H^\pm }^2+2k_1k_3}},`$ (18)
$`P_\nu (k_1+k_2)`$ $`=`$ $`{\displaystyle \frac{1}{m_{H^\pm }^2+2k_1k_2}},`$ (19)
$`P_\gamma (p_1+p_2)`$ $`=`$ $`{\displaystyle \frac{1}{s}},`$ (20)
where $`s=(p_1+p_2)^2`$ and the corresponding tensors are
$`T_1^\mu `$ $`=`$ $`\overline{u}(k_2)(1\gamma _5)v(k_3)\overline{v}(p_2)(k\text{/}_1k\text{/}_2k\text{/}_3)(v_e^za_e^z\gamma _5)u(p_1),`$ (21)
$`T_2^\mu `$ $`=`$ $`\overline{u}(k_2)\gamma ^\mu (v_e^za_e^z\gamma _5)(k\text{/}_1+k\text{/}_3)(1\gamma _5)v(k_3)\overline{v}(p_2)\gamma _\mu (v_e^za_e^z\gamma _5)u(p_1),`$ (22)
$`T_3^\mu `$ $`=`$ $`\overline{u}(k_2)(1\gamma _5)(k\text{/}_1+k\text{/}_2)\gamma _\mu (v_\nu ^za_\nu ^z\gamma _5)v(k_3)\overline{v}(p_2)\gamma ^\mu (v_e^za_e^z\gamma _5)u(p_1),`$ (23)
$`T_4^\mu `$ $`=`$ $`\overline{u}(k_2)(1\gamma _5)v(k_3)\overline{v}(p_2)(k\text{/}_1k\text{/}_2k\text{/}_3)u(p_1),`$ (24)
$`T_5^\mu `$ $`=`$ $`\overline{u}(k_2)\gamma _\mu (k\text{/}_1+k\text{/}_3)(1\gamma _5)v(k_3)\overline{v}(p_2)\gamma ^\mu u(p_1).`$ (25)
In fact, we rearrange the tensors $`T^{^{}}`$s in such a way that they become appropriate to a computer program. Then, following the rules from helicity calculus formalism and using identities of the type
$$\{\overline{u}_\lambda (p_1)\gamma ^\mu u_\lambda (p_2)\}\gamma _\mu =2u_\lambda (p_2)\overline{u}_\lambda (p_1)+2u_\lambda (p_1)\overline{u}_\lambda (p_2),$$
(26)
which is in fact the so called Chisholm identity, and
$$p\text{/}=u_\lambda (p)\overline{u}_\lambda (p)+u_\lambda (p)\overline{u}_\lambda (p),$$
(27)
defined as a sum of the two projections $`u_\lambda (p)\overline{u}_\lambda (p)`$ and $`u_\lambda (p)\overline{u}_\lambda (p)`$.
The spinor products are given by
$`s(p_i,p_j)`$ $``$ $`\overline{u}_+(p_i)u_{}(p_j)=s(p_j,p_i),`$ (28)
$`t(p_i,p_j)`$ $``$ $`\overline{u}_{}(p_i)u_+(p_j)=[s(p_j,p_i)]^{}.`$ (29)
Using Eqs. (10)-(12), which are proved in Ref. , we can reduce many amplitudes to expressions involving only spinor products.
Evaluating the tensors of Eq. (9) for each combination of $`(\lambda ,\lambda ^{^{}})`$ with $`\lambda ,\lambda ^{^{}}=\pm 1`$ one obtains the following expressions:
$`_1(+,+)`$ $`=`$ $`F_1f_1^{+,+}s(k_2,k_3)[s(p_2,k_1)t(k_1,p_1)s(p_2,k_2)t(k_2,p_1)s(p_2,k_3)t(k_3,p_1)],`$ (30)
$`_1(,+)`$ $`=`$ $`F_1f_1^{,+}s(k_2,k_3)[t(p_2,k_1)s(k_1,p_1)t(p_2,k_2)s(k_2,p_1)t(p_2,k_3)s(k_3,p_1)],`$ (31)
$`_2(+,+)`$ $`=`$ $`F_2f_2^{+,+}s(k_2,p_2)t(p_1,k_1)s(k_1,k_3),`$ (32)
$`_2(,+)`$ $`=`$ $`F_2f_2^{,+}s(k_2,p_1)t(p_2,k_1)s(k_1,k_3),`$ (33)
$`_3(+,+)`$ $`=`$ $`F_3f_3^{+,+}s(k_2,k_1)t(k_1,p_1)s(p_2,k_3),`$ (34)
$`_3(,+)`$ $`=`$ $`F_3f_3^{,+}s(k_2,k_1)t(k_1,p_2)s(p_1,k_3),`$ (35)
$`_4(+,+)`$ $`=`$ $`F_4s(k_2,k_3)[s(p_2,k_1)t(k_1,p_1)s(p_2,k_2)t(k_2,p_1)s(p_2,k_3)t(k_3,p_1)],`$ (36)
$`_4(,+)`$ $`=`$ $`F_4s(k_2,k_3)[t(p_2,k_1)s(k_1,p_1)t(p_2,k_2)s(k_2,p_1)t(p_2,k_3)s(k_3,p_1)],`$ (37)
$`_5(+,+)`$ $`=`$ $`F_5s(k_2,p_2)t(p_1,k_1)s(k_1,k_3),`$ (38)
$`_5(,+)`$ $`=`$ $`F_5s(k_2,p_1)t(p_2,k_1)s(k_1,k_3),`$ (39)
where
$`F_1`$ $`=`$ $`2iC_1P_{H^\pm }(k_2+k_3)P_{Z^0}(p_1+p_2),`$ (40)
$`F_2`$ $`=`$ $`4iC_2P_\tau (k_1+k_3)P_{Z^0}(p_1+p_2),`$ (41)
$`F_3`$ $`=`$ $`4iC_3P_\nu (k_1+k_2)P_{Z^0}(p_1+p_2),`$ (42)
$`F_4`$ $`=`$ $`2iC_4P_{H^\pm }(k_2+k_3)P_\gamma (p_1+p_2),`$ (43)
$`F_5`$ $`=`$ $`4iC_5P_\tau (k_1+k_3)P_\gamma (p_1+p_2),`$ (44)
and
$`f_1^{+,+}`$ $`=`$ $`(v_e^za_e^z),`$ (45)
$`f_1^{,+}`$ $`=`$ $`(v_e^z+a_e^z),`$ (46)
$`f_2^{+,+}`$ $`=`$ $`(v_e^za_e^z)^2,`$ (47)
$`f_2^{,+}`$ $`=`$ $`((v_e^z)^2(a_e^z)^2),`$ (48)
$`f_3^{+,+}`$ $`=`$ $`(v_\nu ^z+a_\nu ^z)(v_e^za_e^z),`$ (49)
$`f_3^{,+}`$ $`=`$ $`(v_\nu ^z+a_\nu ^z)(v_e^z+a_e^z).`$ (50)
Here, $`v_e^z=1+4\mathrm{sin}^2\theta _W`$, $`a_e^z=1`$, $`v_\nu ^z=1`$ and $`a_\nu ^z=1`$, according to the experimental data .
After the evaluation of the amplitudes of the corresponding diagrams, we obtain the cross-sections of the analyzed processes for each point of the phase space using Eqs. (13)-(17) by a computer program, which makes use of the subroutine RAMBO (Random Momenta Beautifully Organized). The advantages of this procedure in comparison to the traditional “trace technique” are discussed in Refs. .
We use the Breit-Wigner propagators for the $`Z^0`$, $`h^0`$, $`H^0`$, $`A^0`$ and $`H^\pm `$ bosons. The mass of the bottom $`(m_b4.5`$GeV$`)`$ the mass $`(M_{Z^0}=91.2`$GeV$`)`$ and width $`(\mathrm{\Gamma }_{Z^0}=2.4974`$GeV$`)`$ of $`Z^0`$ have been taken as inputs; the widths of $`h^0`$, $`H^0`$, $`A^0`$ and $`H^\pm `$ are calculated from the formulas given in Ref. . In the next sections we present the numerical computation of the processes $`e^+e^{}b\overline{b}h`$, $`h=h^0,H^0,A^0`$ and $`e^+e^{}\tau ^{}\overline{\nu }_\tau H^+,\tau ^+\nu _\tau H^{}`$.
## III Detection of MSSM Higgs Bosons at Future Positro-Electron Colliders Energies
In this paper, we study the detection of neutral and charged MSSM Higgs bosons at $`e^+e^{}`$ colliders, including three-body mode diagrams \[Figs. 1.1-1.3, 1.5, and 1.6; Figs. 2.1-2.3, 2.5 and 2.6; Figs. 3.2, 3.3, and 3.5\] besides the dominant mode diagram \[Fig. 1.4; Fig. 2.4; Figs. 3.1, and 3.4\] consider two stages of a possible Next Linear $`e^+e^{}`$ Collider: the first with $`\sqrt{s}=500`$ $`GeV`$ and design luminosity 50 $`fb^1`$, and the second with $`\sqrt{s}=1`$ $`TeV`$ and luminosity 100-200 $`fb^1`$. We consider the complete set of Feynman diagrams (Figs. 1-3) at tree-level and utilize the helicity formalism for the evaluation of their amplitudes. In the next subsections, we present our results for the case of the different Higgs bosons.
### A Detection of $`h^0`$
In order to illustrate our results on the detection of the $`h^0`$ Higgs boson, we present graphs in the parameters space region $`(m_{A^0}\mathrm{tan}\beta )`$, assuming $`m_t=175`$ $`GeV`$, $`M_\stackrel{}{t}=500`$ $`GeV`$ and $`\mathrm{tan}\beta >1`$ for NLC. Our results are displayed in Fig. 4, for $`e^+e^{}(A^0,Z^0)+h^0`$ dominant mode and for the processes at three-body $`e^+e^{}b\overline{b}h^0`$.
The total cross-section for each contour is $`0.03`$ $`pb`$, and $`0.01`$ $`pb`$, which gives 1500 events, and 500 events to an integrated luminosity of $`=50`$ $`fb^1`$. We can see from this figure, that the effect of the reaction $`b\overline{b}h^0`$ is not more important that $`(A^0,Z^0)+h^0`$, for most of the $`(m_{A^0}\mathrm{tan}\beta )`$ parameter space regions. Nevertheless, there are substantial portions of parameter space in which the discovery of the $`h^0`$ is not possible using either $`(A^0,Z^0)+h^0`$ or $`b\overline{b}h^0`$.
For the case of $`\sqrt{s}=1`$ $`TeV`$, the results of the detection of the $`h^0`$ are shown in Fig. 5. It is clear from this figure that the contribution of the process $`e^+e^{}b\overline{b}h^0`$ becomes dominant, namely $`e^+e^{}(A^0,Z^0)+h^0`$ is small in all parameter space. However, they could provide important information on the Higgs bosons detection. For instance, we give the contours for the total cross-section to, say 0.01 $`pb`$, 0.005, and 0.003 $`pb`$ for both processes. These cross-sections give 1000 events, 500 events, and 300 events in total to a integrated luminosity of $`=100`$ $`fb^1`$. While for $`=200`$ $`fb^1`$ the events number is 2000, 1000, and 600, respectively, then it will be detectable the $`h^0`$ at future $`e^+e^{}`$ colliders.
### B Detection of $`H^0`$
To illustrate our results regarding the detection of the heavy Higgs bosons $`H^0`$, we give the contours for the total cross-section for both processes $`e^+e^{}(A^0,Z^0)+H^0`$, $`e^+e^{}b\overline{b}H^0`$ in Fig. 6 for $`\sqrt{s}=500`$ $`GeV`$ and $`=50`$ $`fb^1`$. The contours for this cross-section are 0.01 $`pb`$, 0.001 $`pb`$ and 0.0001 $`pb`$ for both reactions $`(A^0,Z^0)+H^0`$ and $`b\overline{b}H^0`$. The number of events corresponding to each contour are 500, 50 and 5, respectively.
Our estimate is that if more than 100 total events are obtained for a given process $`(A^0,Z^0)+H^0orb\overline{b}H^0`$ then the Higgs boson $`H^0`$ can be detectable.
For the case of $`\sqrt{s}=1`$ $`TeV`$, the results on the detection of the $`H^0`$ are show in Fig. 7. The events number for each contour is 1000, 100, and 10 for $`=100`$ $`fb^1`$ and 2000, 200, 20 for $`=200`$ $`fb^1`$.
The effect of incorporate $`b\overline{b}H^0`$ in the detection of the Higgs boson $`H^0`$ is more important than the case of two-body mode $`(A^0,Z^0)+H^0`$, because $`b\overline{b}H^0`$ cover a major region in the parameters space $`(m_{A^0}\mathrm{tan}\beta )`$. The most important conclusion from this figure is that detection of all of the neutral Higgs bosons will be possible at Next Linear $`e^+e^{}`$ Collider.
### C Detection of $`A^0`$
For the pseudoscalar $`A^0`$, it is interesting to consider the production mode in $`b\overline{b}A^0`$, since it can have large a cross-section due to the fact that the coupling of $`A^0`$ with $`b\overline{b}`$ is directly proportional to $`\mathrm{tan}\beta `$, thus will always grow with it. In Fig. 8, we present the contours of the cross-sections for the process of our interest $`b\overline{b}A^0`$, at $`\sqrt{s}=500`$ $`GeV`$ and $`=50`$ $`fb^1`$.
We display the contour lines for $`\sigma =0.01,0.003,0.001`$, showing also the regions where the $`A^0`$ can be detected. These cross-sections give a contour of production of 500, 150 and 50 events. It is clear from this figure that very high experimental and analysis efficiencies are necessary for detecting the Higgs boson $`A^0`$.
On the other hand, if we focus the detection of the $`A^0`$ at $`\sqrt{s}=1`$ $`TeV`$ and an integrated luminosity of $`=100`$ $`fb^1`$, the panorama for its detection is more extensive. The Fig. 9 shows the contours lines in the plane $`(m_{A^0}\mathrm{tan}\beta )`$, to the cross-section of $`b\overline{b}A^0`$. The contours for this cross-section correspond to 300, 100 and 10 events. While for $`=200`$ $`fb^1`$ we have 600, 200, and 20 events respectively. It is estimated that if more than 100 total events are obtained for $`b\overline{b}A^0`$, then it is possible to detect the $`A^0`$.
### D Detection of $`H^\pm `$
Our results for the $`H^+H^{}`$ scalars are displayed in Fig. 10, 11 for $`e^+e^{}H^+H^{}`$ dominant mode and for the processes at three-body $`e^+e^{}\tau ^{}\overline{\nu }_\tau H^+,\tau ^+\nu _\tau H^{}`$.
The total cross-section for this reaction with $`\sqrt{s}=500`$ $`GeV`$ and $`=50`$ $`fb^1`$ are shown in Fig. 10 for each contour with 0.01, 0.001, and 0.0001 $`pb`$, which gives 500 events, 50 events, and 5 events, respectively. We can see from this figure that the effect of the reactions $`\tau ^{}\overline{\nu }_\tau H^+`$ and $`\tau ^+\nu _\tau H^{}`$ is slightly more important than $`H^+H^{}`$ for most of the $`(m_{A^0}\mathrm{tan}\beta )`$ parameters space regions. Nevertheless, there are substantial portions of parameters space in which the discovery of the $`H^\pm `$ is not possible using either $`H^+H^{}`$ or $`\tau ^{}\overline{\nu }_\tau H^+`$ and $`\tau ^+\nu _\tau H^{}`$.
In both cases the curves with values of 0.01 $`pb`$, 0.001 $`pb`$, and 0.0001 $`pb`$ give 1000, 100, and 10 events to an integrated luminosity of $`=100`$ $`fb^1`$. Meanwhile, for $`=200`$ $`fb^1`$ we have 2000, 200, and 20 events. These cross-sections are small, however, it is precisely in this curve where the contribution of the processes at three-body is notable. The most important conclusion from this figure is that detection of the charged Higgs bosons will be possible at future $`e^+e^{}`$ colliders.
## IV Conclusions
In this paper, we have calculated the production of the neutral and charged Higgs bosons in association with $`b`$-quarks and with $`\tau \nu _\tau `$ leptons via the processes $`e^+e^{}b\overline{b}h`$, $`h=h^0,H^0,A^0`$ and $`e^+e^{}\tau ^{}\overline{\nu }_\tau H^+,\tau ^+\nu _\tau H^{}`$ using the helicity formalism. We find that these processes could help to detect the possible neutral and charged Higgs bosons at energies of a possible Next Linear $`e^+e^{}`$ Collider when $`\mathrm{tan}\beta `$ is large.
In summary, we conclude that the possibilities of detecting or excluding the neutral and charged Higgs bosons of the Minimal Supersymmetric Standard Model $`(h^0,H^0,A^0,H^\pm )`$ in the processes $`e^+e^{}b\overline{b}h`$, $`h=h^0,H^0,A^0`$ and $`e^+e^{}\tau ^{}\overline{\nu }_\tau H^+,\tau ^+\overline{\nu }_\tau H^{}`$ are important and in some cases are compared favorably with the dominant mode $`e^+e^{}(A^0,Z^0)+h`$, $`h=h^0,H^0,A^0`$ and $`e^+e^{}H^+H^{}`$ in the region of parameters space $`(m_{A^0}\mathrm{tan}\beta )`$ with $`\mathrm{tan}\beta `$ large. The detection of the Higgs boson will require the use of a future high energy machine like the Next Linear $`e^+e^{}`$ Collider.
Acknowledgments
This work was supported in part by Consejo Nacional de Ciencia y Tecnología (CONACyT) (Proyecto I33022-E) and Sistema Nacional de Investigadores (SNI) (México). A.G.R. would like to thank the organizers of the Summer School in Particle Physics 99, Trieste Italy for their hospitality. O. A. S. would like to thank CONICET (Argentina).
FIGURE CAPTIONS
Fig. 1 Feynman Diagrams at tree-level for $`e^+e^{}b\overline{b}h^0`$. For $`e^+e^{}b\overline{b}H^0`$ one has to make only the change $`\mathrm{sin}\alpha /\mathrm{cos}\beta \mathrm{cos}\alpha /\mathrm{cos}\beta `$.
Fig. 2 Feynman Diagrams at tree-level for $`e^+e^{}b\overline{b}A^0`$.
Fig. 3 Feynman Diagrams at tree-level for $`e^+e^{}\tau ^{}\overline{\nu }_\tau H^+,\tau ^+\nu _\tau H^{}`$.
Fig. 4 Total cross-sections contours in $`(m_{A^0}\mathrm{tan}\beta )`$ parameter space for $`e^+e^{}(A^0,Z^0)+h^0`$ and $`e^+e^{}b\overline{b}h^0`$ with $`\sqrt{s}=500`$ $`GeV`$ and an integrated luminosity of $`=50`$ $`fb^1`$. We have taken $`m_t=175`$ $`GeV`$ and $`M_\stackrel{}{t}=500`$ $`GeV`$ and neglected squark mixing.
Fig. 5 Total cross-sections contours for $`\sqrt{s}=1`$ $`TeV`$ and $`=100`$, 200 $`fb^1`$. We have taken $`m_t=175`$ $`GeV`$, $`M_\stackrel{}{t}=500`$ $`GeV`$ and neglected squark mixing. We display contours for $`e^+e^{}(A^0,Z^0)+h^0`$ and $`e^+e^{}b\overline{b}h^0`$, in the parameters space $`(m_{A^0}\mathrm{tan}\beta )`$.
Fig. 6 Same as in Fig. 4, but for $`e^+e^{}(A^0,Z^0)+H^0`$ and $`e^+e^{}b\overline{b}H^0`$.
Fig. 7 Same as in Fig. 5, but for $`e^+e^{}(A^0,Z^0)+H^0`$ and $`e^+e^{}b\overline{b}H^0`$.
Fig. 8 Same as in Fig. 4, but for $`e^+e^{}b\overline{b}A^0`$.
Fig. 9 Same as in Fig. 5, but for $`e^+e^{}b\overline{b}A^0`$.
Fig. 10 Same as in Fig. 4, but for $`e^+e^{}H^+H^{}`$ and $`e^+e^{}\tau ^{}\overline{\nu }_\tau H^+,\tau ^+\nu _\tau H^{}`$.
Fig. 11 Same as in Fig. 5, but for $`e^+e^{}H^+H^{}`$ and $`e^+e^{}\tau ^{}\overline{\nu }_\tau H^+,\tau ^+\nu _\tau H^{}`$.
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# Phase-field modeling of microstructural pattern formation during directional solidification of peritectic alloys without morphological instability
## I Introduction
The spontaneous emergence of complex microstructural patterns during the solidification of alloys is a subject of both fundamental and applied interest . During directional solidification, a sample is pulled in an externally imposed temperature gradient $`G`$ with a fixed pulling speed $`v_p`$. This setup has been used extensively in fundamental studies of solidification patterns because it allows one to study their formation under well-controlled growth conditions. Depending on the type of alloy and the ratio $`G/v_p`$, various patterns are possible. During monophase solidification of a dilute binary alloy, solute redistribution leads to a well-known morphological instability (Mullins-Sekerka instability ) below a critical ratio $`G/v_p`$, and cellular or dendritic patterns are typically formed. For nondilute alloy concentrations close to a eutectic point, two stable solid phases of different compositions can grow from a metastable liquid. In this case, the two phases cooperate and form lamellae or rods parallel to the growth direction (coupled growth). For off-eutectic compositions, coexistence between dendrites and coupled growth structures is also observed.
Much less is known about microstructural pattern formation in peritectic growth , despite the fact that many industrially important metallic alloy systems as well as ceramics such as the high-$`T_c`$ superconductor YBCO are peritectics. A schematic phase diagram of a peritectic AB-alloy (where B will be called the impurity for convenience) is shown in Fig. 1. It contains a peritectic point, analogous to the eutectic point, at which two different solid phases, the parent (primary) and peritectic (secondary) phases, coexist with a liquid of higher composition than either solid phase. Above the peritectic temperature $`T_p`$, the parent phase is stable and the peritectic phase is metastable, whereas below $`T_p`$, the opposite is true. For comparison, in a eutectic, both solid phases are stable below the eutectic temperature, and metastable above, and the impurity concentration in the liquid falls in between the concentrations of the two solid phases. For a sufficiently low $`G/v_p`$ ratio, a dendritic array structure of the parent or the peritectic phase is typically observed, and which of these two phases is selected depends on the alloy composition and $`G/v_p`$ . In contrast, for a high $`G/v_p`$ ratio morphological instability is suppressed. In this case, banded structures made up of alternating layers of primary and peritectic phases perpendicular to the growth direction are formed. These structures have by now been observed in various peritectic systems, including Sn-Cd , Sn-Sb , Zn-Cu , Ag-Zn , and Pb-Bi . It is worth noting that eutecticlike coupled growth structures, which are quite distinct from banded structures, have recently been observed in the Fe-Ni system . Whether stable coupled growth is theoretically possible during peritectic growth has remained an open question for quite some time , and we will address this issue elsewhere. Here, we focus primarily on banded structure formation and phenomena associated with the dynamical spreading of one solid phase onto the other.
Recently, Trivedi has introduced a one-dimensional (1D) model to explain the formation of peritectic banded structures for purely diffusion-controlled growth The conceptual banding cycle assumed in this model is as follows. Consider a melt with homogeneous composition $`C_{\mathrm{}}<C_p`$ being solidified starting from a flat $`\alpha `$-liquid interface in equilibrium. The rejection of impurities B into the liquid during solidification leads to the buildup of a solutal boundary layer. As a result, the interface temperature decreases, following the liquidus curve in the phase diagram. If $`C_{\mathrm{}}`$ is large enough, the interface temperature eventually falls sufficiently below $`T_p`$ for the peritectic phase to nucleate heterogeneously at the solid-liquid interface before the growth of the $`\alpha `$ phase has reached its steady state. The newly nucleated $`\beta `$ phase rejects fewer impurities than the $`\alpha `$ phase. Consequently, the magnitude of the solutal boundary layer decreases and the interface temperature increases, following now the $`\beta `$-liquid coexistence line in the phase diagram. If $`C_{\mathrm{}}`$ is low enough, such that the corresponding interfacial temperature is sufficiently higher than $`T_p`$, the $`\alpha `$ solid may renucleate again before the steady state is reached, and the cycle repeats. Therefore, this model predicts that bands can form only when the composition falls inside a narrow window in the hypoperitectic region ($`C_{p\alpha }<C_{\mathrm{}}<C_{p\beta }`$) whose width depends on the nucleation undercoolings $`\mathrm{\Delta }T_N^\alpha `$ and $`\mathrm{\Delta }T_N^\beta `$.
The first attempts to validate this prediction experimentally yielded contradictory results. Directional solidification experiments with Pb-Bi and Sn-Cd alloys seemed to show that bands also form in the hyperperitectic region ($`C_{p\beta }<C_{\mathrm{}}<C_p`$), in apparent contradiction with this prediction. An attempt was made to resolve this “composition range paradox” by incorporating convection effects , assuming the existence of a fully mixed liquid of uniform composition outside a purely diffusive 1D boundary layer of finite thickness. This model, however, yielded a banding cycle and band spacings that are inconsistent with experimental results, hinting that this boundary-layer approximation (typically valid for strong convection) is inadequate to describe these experiments. Around the same time, careful serial sectioning of solidified Pb-Bi and Sn-Cd alloys revealed that the seemingly banded structures are actually oscillatory treelike structures connected in three dimensions , and not discrete bands, thereby resolving experimentally this composition range paradox. Following this finding, a more accurate model was developed that assumes a planar solidification front, but incorporates a fully two-dimensional convection flow field . This model successfully reproduced the observed oscillatory structures.
Following these studies, experiments were conducted in thin tubes to reduce convection . For tube diameters smaller than 1 mm, truly discrete bands indeed became observable inside a narrow composition range predicted by the 1D diffusive growth model. Surprisingly, however, it was also observed that when the tube diameter was further reduced, “islands” of the $`\beta `$ phase formed inside the matrix of the $`\alpha `$ phase, instead of discrete bands. This observation suggests that there is a microstructural transition from bands to islands if the system size is reduced. It was also observed that islands tend to form more easily for initial compositions closer to $`C_{p\alpha }`$. In addition, some spatially chaotic patterns were observed in some experiments. The formation of these structures is controlled by a subtle interplay between the nucleation process and the competition between the growth of the nuclei and the preexisting phase. In this respect, the one-dimensional model may not always be adequate to describe this competition because it assumes an infinite spreading speed for the newly nucleated phase. Moreover, the 2D convection model assumes a flat interface and is hence not well suited to simulate heterogeneous nucleation and spreading. In order to model accurately the formation of these different structures, a truly 2D model of interface evolution is necessary. The particular difficulty of this problem is that the microstructure formation is controlled by an interplay between nucleation and growth of the different phases. No steady-state growth mode exists, which makes the whole problem explicitly time-dependent.
In this paper, we use a phase-field approach to investigate the formation of this class of banded microstructures in a purely diffusive regime and a 2D geometry. The phase-field method eliminates the need of explicit front tracking and thus greatly simplifies the task of numerically solving the equations of peritectic solidification that involve three-phase junctions. A phase-field model for peritectic growth has recently been proposed . Here, we use an alternative model that is closer to the eutectic model of Wheeler et al. .
We first investigate the spreading of the peritectic phase on the primary phase after a single nucleation event. We characterize in detail the dynamics of the three-phase junction during spreading and find a morphological transition from discrete bands of $`\alpha `$ and $`\beta `$ phases to isolated islands of $`\beta `$ phase when the system size is decreased, in qualitative agreement with experiments. Moreover, our simulations enable us to understand physically the basic mechanism that underlies this transition. We then investigate the effect of multiple nucleations on microstructure formation in large systems by supplementing the phase-field equations with a phenomenological stochastic nucleation law.
The remainder of this article is organized as follows. In Sec. II, we write down the sharp interface and phase-field models. Section III is devoted to the study of the equilibrium properties of the phase-field model and Sec. IV describes the simulation method. Results are presented in Sec. V, followed by a summary and conclusions in Sec. VI.
## II Model
### A Sharp-interface model
The sharp-interface equations are given by
$$_tC=D_L^2C,$$
(1)
$$v_n(C_LC_\nu )=D_L_nC_L,$$
(2)
$$T=T_p+m_\nu (C_LC_p)\mathrm{\Gamma }_\nu K\frac{1}{\mu _\nu }v_n,$$
(3)
where $`C`$ denotes the concentration of impurity B, and the subscript $`\nu `$ labels the solid $`\alpha `$ and $`\beta `$ phases. Equation (1) is the diffusion equation for the solute in the liquid with the solute diffusivity $`D_L`$. We have assumed that diffusion in the solid is negligible (one-sided model). Equation (2) expresses the mass conservation at the moving interface, with $`v_n`$ and $`_n`$ denoting the normal velocity of the interface and the derivative normal to the interface, respectively. Finally, Eq. (3) is the Gibbs-Thomson condition at the solid-liquid interface, with $`K`$, $`m_\nu `$, $`\mu _\nu `$, and $`\mathrm{\Gamma }_\nu `$ being the interface curvature, liquidus slope, kinetic coefficient, and Gibbs-Thomson constant of phase $`\nu `$, respectively. The Gibbs-Thomson constants $`\mathrm{\Gamma }_\nu `$ are defined by
$$\mathrm{\Gamma }_\nu =\frac{\gamma _{\nu L}T_p}{L_\nu },$$
(4)
where $`\gamma _{\nu L}`$ is the surface energy of the $`\nu `$-liquid interface and $`L_\nu `$ is the latent heat of fusion for phase $`\nu `$, both taken at the peritectic temperature. Young’s condition
$$\gamma _{\alpha L}𝐭_{\alpha L}+\gamma _{\beta L}𝐭_{\beta L}+\gamma _{\alpha \beta }𝐭_{\alpha \beta }=0,$$
(5)
must be satisfied at the trijunction points where three phases meet, where $`𝐭_{\mu \nu }`$ is the unit vector parallel to the $`\mu `$-$`\nu `$ interface and pointing away from the trijunction.
### B Phase-field model
To distinguish between the three possible phases (liquid, $`\alpha `$ solid, and $`\beta `$ solid), we follow a similar approach to that of Wheeler et al. for eutectic solidification by introducing two nonconserved order parameters (phase fields) $`\varphi `$ and $`\psi `$. The first distinguishes between solid ($`\varphi =1`$) and liquid ($`\varphi =1`$), the second between the $`\alpha `$ solid ($`\psi =1`$) and the $`\beta `$ solid ($`\psi =1`$). The solid-liquid interface is defined by the level curve $`\varphi =0`$, and the interface between the solid $`\alpha `$ and $`\beta `$ phases is defined by the level curve $`\psi =0`$ when $`\varphi `$ is positive. One important difference from Ref. is that in our model $`\psi `$ takes the well-defined value $`\psi =0`$ in the liquid. This modification is necessary because, in the model of Wheeler et al., the equation of motion for $`\psi `$ becomes a simple diffusion equation in the liquid. This introduces an undesirable new time scale in dynamical simulations that is removed in the present approach.
As a third dynamical variable we need the composition $`C`$, which is a conserved field. We define the scaled composition
$$c(𝐫,t)=[C(𝐫,t)C_{p\beta }]/\mathrm{\Delta }C_\alpha ,$$
(6)
where $`\mathrm{\Delta }C_\nu =(C_pC_{p\nu })`$, $`\nu =\alpha ,\beta `$, is the concentration jump at the $`\nu `$-liquid interface at $`T_p`$.
In terms of these quantities, the equations of motion that govern the dynamics of the system are given by
$`\tau _\varphi {\displaystyle \frac{\varphi }{t}}`$ $`=`$ $`{\displaystyle \frac{\delta F}{\delta \varphi }},`$ (7)
$`\tau _\psi {\displaystyle \frac{\psi }{t}}`$ $`=`$ $`{\displaystyle \frac{\delta F}{\delta \psi }},`$ (8)
$`{\displaystyle \frac{c}{t}}`$ $`=`$ $`\left[M(\varphi ){\displaystyle \frac{\delta F}{\delta c}}\right],`$ (9)
where $`F`$ is the dimensionless free energy of the system (i.e., the Helmholtz free energy, divided by the product of the system size and a typical value of the free energy density that sets the physical energy scale), $`M(\varphi )`$ is the mobility of the impurities, and $`\tau _\varphi `$ and $`\tau _\psi `$ are (fast) relaxation times for the phase fields. These equations are of the standard variational form known from out-of-equilibrium thermodynamics. Note that, since $`\delta F/\delta c`$ is the local chemical potential $`\mu `$, Eq. (9) is simply the continuity equation for the impurity concentration with the mass current $`𝐉`$ given by
$$𝐉=M(\varphi )\mu .$$
(10)
If there are no fluxes across the boundary of the volume where $`F`$ is defined, $`dF/dt0`$ and Eqs. (7)-(9) imply that the dynamics drives the system toward a minimum of free energy.
The free energy functional of the system is assumed to be of the form
$$F=\{\frac{1}{2}W_\varphi ^2|\varphi |^2+\frac{1}{2}W_\psi ^2|\psi |^2+f(\varphi ,\psi ,c)\}d^3𝐫.$$
(11)
Since $`F`$, $`\varphi `$, and $`\psi `$ are dimensionless, the coefficients $`W_\varphi `$ and $`W_\psi `$ have the dimension of length: they determine the width of the diffuse interfaces. The form of the free energy density is chosen such that there are two minima at $`\psi =\pm 1`$ corresponding to the $`\alpha `$ ($`+`$) and $`\beta `$ ($``$) phases for $`\varphi =+1`$. There is a single minimum in the liquid corresponding to $`\varphi =1`$ and $`\psi =0`$, and $`f(\varphi ,\psi ,c)`$ has a single minimum as a function of $`c`$ for fixed values of $`\varphi `$ and $`\psi `$ corresponding to the three equilibrium phases. A convenient way to match these requirements is to construct a free energy density of the form
$`f(\varphi ,\psi ,c)`$ $`=`$ $`{\displaystyle \frac{\lambda }{2}}\left\{c+A_1h(\varphi )+{\displaystyle \frac{1}{2}}A_2[1+h(\varphi )]h(\psi )\right\}^2`$ (14)
$`\lambda \left\{B_1h(\varphi )+{\displaystyle \frac{1}{2}}B_2[1+h(\varphi )]h(\psi )\right\}+g(\varphi )`$
$`+{\displaystyle \frac{1}{2}}\left[1+h(\varphi )\right]g(\psi )+{\displaystyle \frac{1}{2}}\left[1h(\varphi )\right]\psi ^2.`$
Here, $`\lambda `$ is a positive constant, and $`A_1`$, $`A_2`$, $`B_1`$, and $`B_2`$ are functions of temperature. The function $`g`$ is a double-well potential with minima at $`\pm 1`$, and the function $`h`$ must satisfy $`h(\pm 1)=\pm 1`$ and $`h^{}(\pm 1)=0`$ in order to keep the minima of $`f`$ at constant values of $`\varphi `$ and $`\psi `$, independent of the value of $`c`$. We take
$`g(\varphi )`$ $`=`$ $`1/4\varphi ^2/2+\varphi ^4/4,`$ (15)
$`h(\varphi )`$ $`=`$ $`3(\varphi \varphi ^3/3)/2.`$ (16)
The functions $`g(\psi )`$ and $`h(\psi )`$ are similarly defined. It follows trivially from Eq. (14) that the bulk phase free energy densities are given by
$`f_L`$ $``$ $`f(1,0,c)={\displaystyle \frac{\lambda }{2}}(cA_1)^2+\lambda B_1,`$ (17)
$`f_\alpha `$ $``$ $`f(1,1,c)={\displaystyle \frac{\lambda }{2}}(c+A_1+A_2)^2\lambda (B_1+B_2),`$ (18)
$`f_\beta `$ $``$ $`f(1,1,c)={\displaystyle \frac{\lambda }{2}}(c+A_1A_2)^2\lambda (B_1B_2).`$ (19)
For the mobility function $`M(\varphi )`$, we take
$$M(\varphi )=\frac{D_L}{2\lambda }(1\varphi ).$$
(20)
With this choice, the diffusion coefficient of the impurity is a constant equal to $`D_L`$ in the liquid and zero in both solids, which corresponds to the so-called one-sided model. A standard asymptotic analysis of the sharp-interface limit of the present phase-field model shows that Eqs. (7)–(9) reduce as expected to Eqs. (1)–(3). The relation between the parameters in the two sets of equations is given in the next section.
## III Phase diagram and Equilibrium properties
By applying the well-known common tangent construction to the bulk free energy densities given by Eqs. (17)–(19), we can construct the equilibrium phase diagram of the phase-field model. The equilibrium compositions can be expressed in terms of the temperature-dependent functions $`A_1`$, $`A_2`$, $`B_1`$ and $`B_2`$ (see the appendix). However, since there are only four functions, we can at most fit four lines out of six in the phase diagram (i.e., three pairs corresponding to $`\alpha `$-liquid, $`\beta `$-liquid, and $`\alpha `$-$`\beta `$ coexistence). We may choose to construct, say, the two liquidus lines and the two solidus lines and leave the two solid-solid coexistence curves determined by Eq. (A3) and Eq. (A4). Since we are interested only in the behavior of the system at temperatures close to $`T_p`$, we assume for simplicity that the liquidus and solidus lines are straight, and that the concentration jumps at the solid-liquid interface are constant (liquidus and solidus are parallel). We can then choose $`A_1`$ and $`A_2`$ as constants and $`B_1`$ and $`B_2`$ as linear functions of the temperature. The corresponding expressions for the functions $`A_1`$, $`A_2`$, $`B_1`$, and $`B_2`$ are given in the appendix expressed in terms of the dimensionless temperature field
$$\stackrel{~}{T}=\frac{(TT_p)}{|m_\alpha |\mathrm{\Delta }C_\alpha },$$
(21)
which is a measure of the temperature relative to $`T_p`$ normalized by the freezing range of the $`\alpha `$ phase.
In the present model, there exists a temperature-dependent concentration $`c_u`$ such that, in the solid, the solid $`\alpha `$ ($`\beta `$) phase is thermodynamically stable only if $`c<c_u`$ ($`c>c_u`$). By comparing Eq. (18) and Eq. (19), it is easy to show that $`c_u`$ is exactly midway between the two solid-solid coexistence lines. In order to avoid a phase transformation in the solid far behind the solid-liquid interface, we require $`c_u`$ to be independent of temperature. One way to achieve this is to make the solid-solid coexistence lines vertical by choosing suitable parameters. This difference from a real peritectic phase diagram is unlikely to change the qualitative behavior of the system. A phase diagram for the model system used in our simulations is shown in Fig. 2.
The equilibrium interface profiles connecting different phases can be obtained by solving the time-independent one-dimensional version of the equations of motion with suitable boundary conditions. Since the chemical potential must be constant at equilibrium, the relation
$$\mu =\frac{\delta F}{\delta c}=\frac{f}{c}$$
(22)
can be used to eliminate the concentration field from Eqs. (7) and (8). The appropriate value of $`\mu `$ for a certain temperature is obtained from the common tangent construction. The two resulting coupled ordinary differential equations were solved numerically using a Newton-Raphson method on a one-dimensional grid of spacing $`\mathrm{\Delta }x`$. For simplicity, we assumed $`W_\varphi =W_\psi =W`$. Unless otherwise stated, all the results below are obtained for $`\mathrm{\Delta }x/W=0.8`$, which provides a good compromise between computational efficiency and accuracy. The resulting equilibrium profiles, centered at the origin, for the phase fields and the concentration for $`\alpha `$-L equilibrium and $`\beta `$-L equilibrium at $`T_p`$ are shown in Fig. 3. For solid-solid equilibrium, the interface profile of $`\psi `$ can be obtained analytically because $`\varphi =+1`$ is a constant:
$$\psi _0(x)=\mathrm{tanh}(\frac{x}{\sqrt{2}W}).$$
(23)
In all cases, the concentration profiles are given by substituting the equilibrium profiles $`\varphi _0(x)`$ and $`\psi _0(x)`$ obtained previously into Eq. (22).
With the equilibrium interface profiles at hand, we can calculate the surface energies $`\gamma _{\alpha L}`$, $`\gamma _{\beta L}`$, and $`\gamma _{\alpha \beta }`$, defined as the excess Gibbs free energy per unit surface area. They are given by the expressions
$$\gamma _{\mu \nu }=_{\mathrm{}}^{\mathrm{}}[W_\varphi ^2(_x\varphi _0)^2+W_\psi ^2(_x\psi _0)^2]𝑑x,$$
(24)
where $`\varphi _0`$ and $`\psi _0`$ are the equilibrium profiles of the phase fields connecting phases $`\mu `$ and $`\nu `$. The same formula for the surface energies can also be obtained by a matched asymptotic expansion . For the solid-liquid interfaces, the surface energies are obtained by numerical integration. For this purpose, it is more convenient to convert Eq. (24) to a form without the gradients of the fields. Making use of the steady-state one-dimensional version of the equations of motion and the fact that $`\mu `$ is constant in equilibrium, we obtain after some algebra
$$\frac{d}{dx}\left[\frac{1}{2}[W_\varphi ^2(_x\varphi _0)^2+W_\psi ^2(_x\psi _0)^2]\right]=\frac{d}{dx}(f\mu c).$$
(25)
Now one can integrate Eq. (25) from $`\mathrm{}`$ to an arbitrary $`x`$ and make use of the expression for the equilibrium concentration profile and the bulk phase values to show that for the solid-liquid interfaces,
$`{\displaystyle \frac{1}{2}}[W_\varphi ^2(_x\varphi _0)^2+W_\psi ^2(_x\psi _0)^2]`$ (26)
$`=`$ $`g(\varphi _0)+{\displaystyle \frac{1}{2}}[1+h(\varphi _0)]g(\psi _0)+{\displaystyle \frac{1}{2}}[1h(\varphi _0)]\psi _0^2`$ (28)
$`+{\displaystyle \frac{\lambda }{2}}\left({\displaystyle \frac{\overline{B}}{\overline{A}}}B_2\right)[1+h(\varphi _0)][h(\psi _0)1].`$
Here the upper and lower signs are for $`\alpha `$-liquid and $`\beta `$-liquid equilibrium, respectively, and $`\overline{A}`$ and $`\overline{B}`$ are defined in the appendix. For the solid-solid interface, the surface energy $`\gamma _{\alpha \beta }`$ can be calculated exactly and is equal to $`2\sqrt{2}W_\psi /3`$.
Related to the surface energies are the two capillary lengths $`d_0^\alpha `$ and $`d_0^\beta `$ defined as
$$d_0^\nu =\frac{\gamma _{\nu L}}{(\mathrm{\Delta }c_\nu )^2(\mu /c)}$$
(29)
which can also be expressed in terms of the Gibbs-Thomson constants $`\mathrm{\Gamma }_\nu `$ by
$$d_0^\nu =\frac{\mathrm{\Gamma }_\nu }{|m_\nu |\mathrm{\Delta }C_\nu }.$$
(30)
These are two of the physical length scales that are relevant in pattern formation in solidification problems. In real systems, the capillary lengths are microscopic and much smaller than all other physical length scales in the problem. Ideally, one would like to adjust the capillary lengths in the model to match the physical length scale ratios by choosing suitable model parameters. Since $`\mu /c=\lambda `$, it follows that $`d_0^\nu `$ depends on $`\lambda `$ as
$$d_0^\nu \frac{\gamma _{\nu L}}{\lambda }.$$
(31)
Hence, to have small capillary lengths, one would like to increase $`\lambda `$. However, $`\lambda `$ cannot be chosen arbitrarily large for two reasons. First, the surface tensions themselves depend weakly on $`\lambda `$ for $`TT_p`$. As shown in Fig. 4, these variations amount to a few percent over the temperature range of interest when $`\lambda `$ is varied by a factor of $`5`$.
Secondly, the temperature range in which equilibrium interface solutions exist also depends on $`\lambda `$. More precisely, with a fixed value for $`\lambda `$, the $`\alpha `$-liquid equilibrium solution does not exist if $`T`$ is below a certain value, and the $`\beta `$-liquid interface solution ceases to exist if $`T`$ is above another value, because if $`T`$ is too low or too high, the free energy density loses a minimum at $`\psi `$ equal to $`+1`$ or $`1`$, respectively. We have estimated the range of temperatures in which both solutions exist for different $`\lambda `$ by finding equilibrium solutions at different temperatures. The results are shown in Fig. 5.
We can see that this temperature range becomes narrower when $`\lambda `$ increases. From now on, we fix $`\lambda =2.5`$ unless otherwise stated. This is a compromise between having a large $`\lambda `$ and and a sufficient working temperature range in which our two-dimensional simulations can be carried out.
To check Young’s condition, we performed two-dimensional simulations at $`T_p`$ on a square grid with $`\mathrm{\Delta }x/W=0.8`$. The equilibrium angles around a trijunction were measured and found to be consistent with Eq. (5) to within a few degrees.
For a moving interface, there are also nonequilibrium kinetic effects related to the attachment of atoms at the interface and solute trapping. Since we are mostly interested here in qualitative aspects of the growth morphologies, we have not analyzed all these effects in detail. We checked, however, by performing dynamical one-dimensional simulations that nonequilibrium effects only lead to a deviation from local equilibrium that does not exceed the Gibbs-Thomson effect caused by interface curvature in two-dimensional simulations.
## IV Simulations
For our simulations, we cast the equations of motion into a dimensionless form. For simplicity, we take $`W_\varphi =W_\psi =W`$ and $`\tau _\varphi =\tau _\psi =\tau `$. By defining the dimensionless variables
$$\stackrel{~}{𝐫}=\frac{𝐫}{W},\stackrel{~}{t}=\frac{t}{\tau }$$
(32)
and the new variable
$$\stackrel{~}{\mu }=c+A_1h(\varphi )+\frac{1}{2}A_2[1+h(\varphi )]h(\psi ),$$
(33)
the equations of motion can be written in the form
$`{\displaystyle \frac{\varphi }{\stackrel{~}{t}}}`$ $`=`$ $`\stackrel{~}{}^2\varphi {\displaystyle \frac{f}{\varphi }},`$ (34)
$`{\displaystyle \frac{\psi }{\stackrel{~}{t}}}`$ $`=`$ $`\stackrel{~}{}^2\psi {\displaystyle \frac{f}{\psi }},`$ (35)
$`{\displaystyle \frac{c}{\stackrel{~}{t}}}`$ $`=`$ $`\alpha \stackrel{~}{}[\stackrel{~}{D}(\varphi )\stackrel{~}{}\stackrel{~}{\mu }],`$ (36)
where
$$\alpha =\frac{\tau D_L}{W^2}$$
(37)
is the scaled diffusion coefficient of the impurity in the liquid, and
$$\stackrel{~}{D}(\varphi )=(1\varphi )/2.$$
(38)
Instead of the concentration far from the interface, we may also use $`\eta _\beta `$, the volume fraction of $`\beta `$ formed in the solid, to characterize the overall composition of the sample. The two quantities are related by
$$C_{\mathrm{}}=(1\eta _\beta )C_{p\alpha }+\eta _\beta C_{p\beta }.$$
(39)
In a typical directional solidification experiment, the sample is pulled under a temperature gradient $`G`$ with a pulling velocity $`v_p`$. We define the dimensionless temperature gradient and velocity, $`\stackrel{~}{G}`$ and $`\stackrel{~}{v}`$, by
$`\stackrel{~}{G}`$ $`=`$ $`{\displaystyle \frac{G}{|m_\alpha |\mathrm{\Delta }C_\alpha }}W,`$ (40)
$`\stackrel{~}{v}_p`$ $`=`$ $`v_p\tau /W.`$ (41)
Usually, thermal diffusion is orders of magnitude faster than the diffusion of the impurities, and hence we use the “frozen temperature approximation” which assumes that the temperature of the system adjusts instantaneously to the externally imposed temperature gradient. Accordingly, directional growth along the $`x`$ axis is implemented by letting
$$\stackrel{~}{T}=\stackrel{~}{T}_0+\stackrel{~}{G}(\stackrel{~}{x}\stackrel{~}{v}_p\stackrel{~}{t}),$$
(42)
where $`\stackrel{~}{T}_0`$ is some reference temperature.
There are five different physical length scales that control the microstructural pattern formation: the two capillary lengths $`d_0^\alpha `$ and $`d_0^\beta `$ defined by Eq. (29), the two thermal lengths
$$l_T^\nu =\frac{|m_\nu |\mathrm{\Delta }C_\nu }{G}=\frac{|m_\nu |\mathrm{\Delta }C_\nu }{|m_\alpha |\mathrm{\Delta }C_\alpha }\frac{W}{\stackrel{~}{G}},$$
(43)
and the diffusion length
$$l_D=\frac{D_L}{v_p}=\frac{\alpha }{\stackrel{~}{v}_p}W.$$
(44)
Equations. (34)-(36) are integrated numerically on a two-dimensional grid. We use $`\alpha =1`$, $`\mathrm{\Delta }\stackrel{~}{x}=0.8`$, and $`\mathrm{\Delta }\stackrel{~}{t}=0.1`$. Zero-flux boundary conditions are applied to the two sides that are parallel to the growth direction. There are several features in the model that can be exploited in order to speed up the computation. First, the phase fields $`\varphi `$ and $`\psi `$ differ significantly from $`\pm 1`$ only in the interfacial region, and hence we can avoid integrating Eqs. (34) and (35) away from the interface. In addition, Eq. (36) needs to be integrated only in the liquid. Secondly, the concentration field decays exponentially in the growth direction and varies only slowly in space in the liquid region far ahead of the interface. Hence, we can use a coarser and coarser grid as we move away from the interfacial region. Thirdly, in order to simulate a semi-infinite system in the growth direction, we take advantage of the fact that all the fields remain unchanged in the solid in the one-sided model. Whenever the solid-liquid interface has advanced one lattice spacing, we pull the system back by one unit and keep the composition at the end of the liquid side at $`c_{\mathrm{}}`$. With all these implementations, we are able to carry out simulations with typical lengths in the growth direction equal to about ten times the diffusion length. For the results presented in this article, we chose a pulling speed such as to have a diffusion length of $`l_D=200W`$. Other parameters and length scales are listed in Table I.
## V Results
### A Dynamics of spreading
Let us first concentrate on the spreading of the $`\beta `$ phase on the $`\alpha `$ phase, starting from a single nucleus. Similarly to the situation considered in Trivedi’s model, the simulation is started with a homogeneous composition in the liquid and a planar $`\alpha `$-liquid interface. The lateral system size $`L`$ is several times the diffusion length. Nucleations are assumed to occur heterogeneously at the solid-liquid interface when the temperature of the metastable interface reaches a certain undercooling with respect to the stable solid-liquid equilibrium. The nucleation undercoolings $`\mathrm{\Delta }T_N^\alpha `$ ($`\alpha `$ on $`\beta `$) and $`\mathrm{\Delta }T_N^\beta `$ ($`\beta `$ on $`\alpha `$), shown in Fig. 1, are assumed to be constant. Accordingly, in our simulations a circular nucleus of $`\beta `$ phase is put at the solid-liquid interface on one side of the box when the liquid composition at the interface reaches the threshold for nucleation fixed by the nucleation undercooling $`\mathrm{\Delta }T_N^\beta `$. The radius of the nucleus is taken to be $`6W`$, slightly larger than the critical radius for nucleation. Since we are interested here in the deterministic spreading dynamics following a single nucleation event, further nucleation is prohibited. Multiple nucleation events will be treated in Sec. V D.
To characterize the dynamics of spreading, we recorded the position and velocity of the trijunction point. The sideways velocity $`v_y`$ can be regarded as a measure of the spreading speed of the $`\beta `$ phase. Figures. 6(a) and 6(b) show plots of $`v_y/v_p`$ versus time for different nucleation undercoolings and different compositions $`c_{\mathrm{}}`$, respectively. Time is measured in terms of the diffusion time
$$t_D=\frac{l_D}{v_p}=\frac{D_L}{v_p^2}.$$
(45)
Two very different regimes of spreading can be clearly distinguished. Immediately after the nucleation, the spreading velocity is almost independent of the composition, but strongly depends on the nucleation undercooling. The growth of the nucleus is influenced only by its immediate surroundings. On the length scale of the nucleus, which is much smaller than the diffusion length, the impurity concentration can be considered constant and is determined only by the nucleation undercooling. A higher $`\mathrm{\Delta }\stackrel{~}{T}_N^\beta `$ is equivalent to a higher supersaturation, and hence a higher growth speed.
At later times, the modifications of the diffusion field induced by the growing $`\beta `$ phase influence the spreading dynamics, and the spreading velocities for equal undercooling, but different $`c_{\mathrm{}}`$ start to differ \[Fig. 6(b)\]. After a complicated transient, the details of which depend on the choice of parameters, $`v_y`$ becomes a linear function of time, which means a constant lateral acceleration of the trijunction. This acceleration is independent of the nucleation undercooling or the history of the system, but depends on the composition. An explanation of this finding can be deduced from Fig. 7, which shows the interface temperatures on the sides of the box and at the trijunction as functions of time. After the initial transient, the temperature at the trijunction just follows the temperature on the $`\alpha `$ side. This implies that the $`\alpha `$-liquid interface is almost planar up to the trijunction. We can also see that during the whole time of the simulation, the $`\alpha `$-liquid interface is still relaxing toward its steady state below $`T_p`$. Hence, the undercooling that drives the $`\beta `$ phase to spread is increasing.
In the late stages of spreading, the lateral diffusion length $`D/v_y`$ becomes much smaller than the solute boundary layer, and is comparable to or even smaller than the tip of the spreading finger. Therefore, the spreading speed should be a function of local supersaturation only. To check this assumption, we show in Fig. 6(c) the same velocity curves as before, but now plotted against the undercooling of the planar $`\alpha `$-liquid interface with respect to the peritectic temperature, $`\stackrel{~}{T}`$. Since in our phase diagram the liquidus curves are straight lines, this undercooling is simply proportional to the supersaturation. The curves all collapse onto a single master curve after the initial transient, i.e., starting from the time when the interface ahead of the trijunction has become flat. This master curve is not linear, and does not smoothly extrapolate to zero. We did not attempt to calculate it theoretically. We expect that its detailed form should depend on the characteristics of the trijunction, and in particular on the angles between the different interfaces. More theoretical and numerical work would be needed to elucidate in detail the role of the various material parameters. Remarkably, similar observations were very recently reported in experiments on a transparent organic eutectic alloy during spreading of the secondary phase on a planar interface of primary phase. The spreading speed of the secondary phase showed an approximately linear increase with time, and the data could also reasonably well be rescaled onto an analogous master curve.
Since the $`\alpha `$-liquid interface far ahead of the trijunction stays fairly planar before the arrival of the $`\beta `$ phase, the time dependence of the temperature on the $`\alpha `$ side can be well described by the Warren-Langer approximation . The rate of change of the supersaturation is solely determined by the composition $`c_{\mathrm{}}`$, which explains why the final slope of the curves in Figs. 6(a) and 6(b) depends on $`c_{\mathrm{}}`$ but not on the nucleation undercooling.
A completely different behavior is observed when the composition is sufficiently low. As shown in Fig. 6(b) (dashed line), the initial spreading speed is the same as for the other runs. However, at later times, the spreading slows down and the trijunction point turns around such that $`v_y`$ becomes negative. Instead of a band, an isolated island of $`\beta `$ phase is formed. This phenomenon will be addressed in detail in Sec. V B below.
In the final regime of spreading, when the lateral speed becomes much larger than the pulling speed, the lateral diffusion length $`D_L/v_y`$ becomes comparable to the radii of curvature close to the trijunction point. In free growth, such conditions are reached only at very large solidification speeds. Under these circumstances, it is clear that the phase-field model does no longer reflect quantitatively the sharp-interface equations, since it contains corrective terms due to the finite width of the interface. For instance, we observed a violation of Young’s condition at the trijunction point. More precisely, the angles between the interfaces, obtained by taking the tangent vectors to the $`\varphi =0`$ and $`\psi =0`$ level curves at the trijunction, are still consistent with local equilibrium, but the solid-solid interface is highly curved on a length scale comparable to the width $`W`$ of the diffuse interface. This is due to the fact that the diffusivity varies smoothly within the diffuse interface, and hence the part of the solid-solid interface near to the trijunction is still able to move. As a result, the angles between the interfaces, seen on a macroscopic scale, differ from the local equilibrium angles. For the purpose of the present study, where we are mainly interested in the qualitative features of the microstructures, we did not investigate this effect quantitatively. Let us remark, however, that such effects may not be simply an artifact of the phase-field model, but may have a physical significance for high growth speeds if the relaxation of the trijunction toward local equilibrium occurs on a time scale comparable to the time of diffusion through the trijunction region.
### B Morphology transition
The results of the preceding section were obtained for systems with lateral extensions of several times the diffusion length. For some sets of parameters, a surprising event occurs when the system size is reduced while all other parameters are kept constant. After some time, the spreading slows down, and the trijunction point may even turn around, such that $`v_y`$ becomes negative. As a result, the trijunction travels back to the wall where it originated, and an isolated island of $`\beta `$ phase, or partial band, is formed. It hence appears that complete spreading is easier to achieve in larger systems, a quite counterintuitive result. Figsures 8(a) and 8(b) show time series of typical snapshot pictures for the formation of an island and a band, respectively. The scales are the same on both axes and in both figures. Isoconcentration lines in the liquid are also shown. It can be seen that a lateral concentration gradient builds up in the liquid. This concentration gradient plays an important role in the interpretation of the morphological transition from islands to bands and will be discussed below.
To study more systematically the conditions for the formation of islands, we performed simulations with various lateral system sizes $`L`$ and compositions $`c_{\mathrm{}}`$, with the following results1
1. At a fixed $`\mathrm{\Delta }T_N^\beta `$, there exists a critical composition $`c^{}`$ such that if $`c_{\mathrm{}}<c^{}`$ the $`\beta `$ phase always forms islands. This critical composition decreases as the nucleation undercooling increases.
2. At a fixed $`\mathrm{\Delta }T_N^\beta `$ and if $`c_{\mathrm{}}>c^{}`$, there exists a critical lateral system size $`L_c(c_{\mathrm{}})`$ such that if the lateral system size $`L>L_c`$ the $`\beta `$ phase spreads completely and forms bands, whereas for $`L<L_c`$ it forms islands. $`L_c`$ decreases when either $`c_{\mathrm{}}`$ or $`\mathrm{\Delta }T_N^\beta `$ increases.
3. When $`c_{\mathrm{}}<c^{}`$, such that the $`\beta `$ phase always forms islands, the final shape (and also the size) of the islands is independent of $`L`$ when $`L`$ is larger than a certain size.
Figure 9 shows the final morphology of the system for different $`c_{\mathrm{}}`$ (or equivalently different $`\eta _\beta `$) and for different nucleation undercoolings. The dashed lines in Figs. 9(a) and (b) represent an estimate for the critical system size $`L_c(c_{\mathrm{}})`$ for the transition from bands to islands. It can be seen that both $`L_c`$ and $`c^{}`$ are smaller for higher $`\mathrm{\Delta }\stackrel{~}{T}_N^\beta `$.
The existence of the critical size $`L_c`$ can be understood by noticing that both $`\alpha `$ and $`\beta `$ phases have to reject impurities in order to grow, but the concentration jump at the $`\alpha `$-liquid interface is larger than that at the $`\beta `$-liquid interface. Since $`\beta `$ is the stable phase below $`T_p`$, there exists a driving force for the $`\beta `$ phase to spread. On the other hand, as $`\beta `$ rejects fewer impurities than $`\alpha `$, the impurity concentration in front of the $`\beta `$ phase rapidly decreases after the nucleation. This creates a lateral concentration gradient and hence an impurity flow from the $`\alpha `$ to the $`\beta `$ side as can be clearly seen in Fig. 8. This lateral impurity backflow will accelerate the growth of $`\alpha `$ and slow down the growth of $`\beta `$ and hence there is a competition between the two phases.
To be more precise, we can consider the following scaling argument. Let us assume for simplicity a constant spreading speed $`v_s`$ for the $`\beta `$ phase. Then the time required by the $`\beta `$ phase to spread across the system is $`L/v_s`$. On the other hand, impurities diffuse laterally through the system on a time scale of $`L^2/D_L`$. If $`L/v_s<L^2/D_L`$, the $`\beta `$ phase is able to spread over the $`\alpha `$ phase before a significant impurity backflow can occur. If the opposite is true, the impurities have enough time to diffuse and the growth of the $`\beta `$ phase is slowed down. Hence, the critical system size is given by
$$L_c\frac{D_L}{v_s}.$$
(46)
Another way to interpret the above criterion is to note that the “diffusion speed”, which is roughly the speed of the impurities diffusing laterally through the system, is given by $`D_L/L`$. If the diffusion speed is smaller than $`v_s`$, spreading occurs, and Eq. (46) follows immediately.
Clearly, the above argument is only qualitative. We have assumed a constant lateral spreading speed in Eq. (46), although Fig. 6 shows that the spreading speed varies with time, and hence we can give no explicit expression for $`v_s`$ as a function of composition and nucleation undercooling. However, we can see from Fig. 6 that, for any given time, the instantaneous spreading speed increases with increasing nucleation undercooling and increasing volume fraction of $`\beta `$ phase. This observation, together with Eq. (46), allows us to understand qualitatively the shape of the curves $`L_c(c_{\mathrm{}})`$ in Fig. 9. In addition, this criterion allows understanding of the striking finding that spreading is easier in large systems that in small ones.
### C Banding and island formation
So far, we have concentrated on how a single $`\beta `$ nucleus spreads on the $`\alpha `$ phase. It is natural to ask what happens if renucleation is allowed. This is a complicated problem since nucleation is an inherently stochastic phenomenon, which cannot be consistently treated within our deterministic model. However, we can try to gain some insight by incorporating nucleation phenomenologically. We will proceed in two steps. First, we treat repeated nucleation in small samples by deterministic rules to make contact with the recent experiments in the Sn-Cd alloy system . Then, in the next section, we investigate the influence of multiple stochastic nucleations on the pattern formation dynamics in large systems.
For solidification in small systems, it can be assumed that nucleation occurs predominantly at the container walls close to the solid-liquid interface. The density of nuclei and the nucleation rate are very rapidly varying functions of the composition. Therefore, it seems reasonable to assume that a nucleus will form as soon as the concentration in the liquid exceeds the threshold corresponding to the nucleation undercooling. Accordingly, we incorporate repeated nucleation by the following rules (the nucleation of $`\alpha `$ is handled like the nucleation of $`\beta `$ before, by placing a small circular nucleus at the solid-liquid interface).
1. A nucleus of the new phase is placed at one side of the container as soon as the undercooling of the interface exceeds the nucleation undercooling. If both sides of the box reach the threshold at the same time, one side is chosen at random.
2. Once a germ has nucleated, further nucleation of the same phase is prohibited until the germ has either completely spread across the system or completed the formation of an island.
We consider now the two nucleation undercoolings as free parameters, and study different cases. When both $`\mathrm{\Delta }T_N^\alpha `$ and $`\mathrm{\Delta }T_N^\beta `$ are large enough such that the newly nucleated phase spreads completely before the original phase is able to renucleate again, banded structures are obtained. For a smaller $`\mathrm{\Delta }T_N^\beta `$, the $`\beta `$ phase does not spread completely, but forms an island. The $`\alpha `$ phase overtakes the $`\beta `$ phase and continues to grow until the nucleation threshold for $`\beta `$ is reached again. The islands of $`\beta `$ phase form alternately on each side. This is a result of the history of the system: as a result of the formation of the previous island, the concentration of impurities is lower on the side where the last island occurred, and hence nucleation of $`\beta `$ is favored at the other side. Examples of these banded and island structures are shown in Fig. 10(a) to Fig. 10(d).
The scales on both axes in these figures are the same, but Fig. 10(e) has a different scale from Figs. 10(a)–(d). This last picture was obtained by a simulation at much smaller $`\mathrm{\Delta }T_N^\alpha `$, and the lateral size of the system is smaller. In this case, an oscillatory structure is obtained which tends to approach a coupled growth steady state after a complicated transient.
These results are in good qualitative agreement with microstructures obtained in small samples of Sn-Cd alloy . In the experiments, islands tend to form always on the same side of the sample. We believe that this is due to a slight lateral temperature gradient across the sample, which is always present in experiments.
### D Nucleation controlled microstructures in spatially extended systems
Until now, we have mainly focused on the microstructures formed in small samples. It is interesting to ask what kinds of structure are to be expected in large samples in the absence of convection. This situation could be achieved either in quasi-two-dimensional thin samples, or in a microgravity environment. To model the microstructure formation, larger scale computations were carried out. However, in a spatially extended system, multiple nucleations are unavoidable and must be incorporated in a way that is consistent with the predictions of classical nucleation theory. We chose to extend our model by incorporating the effects of multiple nucleation in a phenomenological manner.
For the nucleation of the $`\beta `$ phase on a planar $`\alpha `$ front (the same arguments also apply to the nucleation of $`\alpha `$ on $`\beta `$), classical nucleation theory predicts the nucleation rate
$$I=I_0e^{\mathrm{\Delta }F^{}/k_BT},$$
(47)
where $`I_0`$ is a constant prefactor (with dimension equal to the number of nucleations per unit volume per unit time) and $`\mathrm{\Delta }F^{}`$ is the activation energy for heterogeneous nucleation. Assuming that the critical nucleus is a spherical cap on a planar substrate (the spherical cap model), $`\mathrm{\Delta }F^{}`$ is given, respectively, in two and three dimensions by
$$\mathrm{\Delta }F^{}=\{\begin{array}{cc}\frac{\gamma _{\beta L}^2}{\mathrm{\Delta }F_B}\times \frac{\theta ^2}{\theta (1/2)\mathrm{sin}2\theta },\hfill & 2\mathrm{D}\hfill \\ \frac{\gamma _{\beta L}^3}{\mathrm{\Delta }F_B^2}\times \frac{16\pi (2+\mathrm{cos}\theta )(1\mathrm{cos}\theta )^2}{12},\hfill & 3\mathrm{D},\hfill \end{array}$$
(48)
where $`\mathrm{\Delta }F_B`$ is the difference between the bulk free energy of the $`\beta `$ phase and of the liquid phase. The contact angle $`\theta `$, as shown schematically in Fig. 11, is determined by the balance of surface tensions parallel to the substrate,
$$\gamma _{\alpha L}=\gamma _{\alpha \beta }+\gamma _{\beta L}\mathrm{cos}\theta .$$
(49)
Assuming that the system is locally in thermodynamic equilibrium, it can be shown that $`\mathrm{\Delta }F_B`$ is proportional to $`(TT_p)`$ , such that for a quasi-two-dimensional system the nucleation rate for $`\beta `$ on $`\alpha `$ can be written as
$$I=\{\begin{array}{cc}I_{2D}\mathrm{exp}[A/(TT_p)^2]\hfill & \mathrm{if}T<T_p\hfill \\ 0\hfill & \mathrm{if}TT_p\hfill \end{array}$$
(50)
where $`A`$ is a constant, and $`I_{2D}`$ has now the dimension of number of nucleations per unit time and per unit length of the interface. A similar expression with $`w=0`$ when $`TT_p`$ holds for $`\alpha `$ nucleating on $`\beta `$. The 3D form of $`\mathrm{\Delta }F^{}`$ is used in deriving Eq. (50) since, in practice, the size of a nucleus is still much smaller than the thickness of a thin sample. Equation (50) determines the local nucleation rate and hence the probability per unit time of a nucleus forming as a function of the local temperature at the solid-liquid interface.
Unfortunately, both experimental and theoretical estimates of the free energy barrier and the kinetic prefactor are scarce in the context of heterogeneous nucleation, since the actual values may depend on complicated details of the interfacial structure. Since, in the present study, we focus on morphological aspects of the large scale structure, we decided to treat the two quantities as free parameters. Moreover, we want to compare the stochastic simulations to the deterministic runs of the preceding sections. Consequently, we may eliminate one of those two parameters by the requirement of recovering the rules used previously. That is, in the deterministic simulations a nucleus was put at the solid-liquid interface when it reached the predetermined nucleation undercooling. In the stochastic runs, nucleation should therefore occur with probability $`1`$ for the same interface temperature. This condition will lead to a relation between the prefactor and the energy barrier in the nucleation rate.
To proceed, let us first specify how we treat nucleation in the simulation algorithm. The interface is scanned at a regular time interval $`\mathrm{\Delta }t_N`$, and nucleation is attempted at points regularly spaced by a distance $`\mathrm{\Delta }s_N`$ along the interface. The nucleation rate may be rewritten as
$$I=\frac{w(T)}{\mathrm{\Delta }t_N\mathrm{\Delta }s_N},$$
(51)
where
$$w=\{\begin{array}{cc}w_0\mathrm{exp}[A/(TT_p)^2]\hfill & \mathrm{if}T<T_p\hfill \\ 0\hfill & \mathrm{if}TT_p\hfill \end{array}$$
(52)
is a dimensionless function of the interface temperature. At each test point, a nucleus is generated with probability $`1`$ if $`w>1`$, and with probability $`w`$ otherwise. That is, if $`w<1`$, a random number $`\xi `$ uniformly distributed between $`0`$ and $`1`$ is drawn, and a nucleus is generated if $`\xi <w`$. As before, the nucleus is spherical and has a size of $`6W`$. A possible drawback of the procedure outlined above is that the actual nucleation rate depends on the values chosen for $`\mathrm{\Delta }t_N`$ and $`\mathrm{\Delta }s_N`$. However, it is reasonable to assume that the microstructures should not depend too sensitively on the choice of these parameters as long as their values are much smaller than the time and length scales of the pattern formation process.
Now we can relate the prefactor and the barrier in the nucleation rate. In the preceding sections, nuclei were introduced deterministically when the nucleation undercooling was reached, that is, at a nucleation temperature
$$T_N^\beta =T_p\frac{\mathrm{\Delta }T_N^\beta }{1m_\beta /m_\alpha }.$$
(53)
This implies that in the stochastic algorithm we must choose
$$w(T_N^\beta )=w_0\mathrm{exp}[A/(T_N^\beta T_p)^2]=1.$$
(54)
Given this condition, there is only one remaining free parameter; we choose the dimensionless kinetic prefactor $`w_0`$. This parameter controls the temperature range over which nucleations occur. Indeed, $`w(T)`$ is a function that rapidly increases in a narrow temperature range around $`T_N^\beta `$. Increasing $`w_0`$ and $`A`$ simultaneously while respecting the constraint Eq. (54), the rise of the nucleation rates becomes sharper, as shown schematically in Fig. 12.
Now consider an $`\alpha `$-liquid interface during its transient, when the temperature at the interface is decreasing from above $`T_p`$ to its steady-state temperature. If the temperature range over which the nucleation rate increases significantly is narrow, all nucleation events will occur almost at the same time when the interfacial temperature coincides with $`T_N^\beta `$. As a result, the mean separation between the nuclei will be small. On the other hand, if $`w_0`$ is small, nuclei appear with a broader spread in interfacial temperature, and the mean separation between the nuclei will be larger. Since the mean separation between the nuclei plays a similar role as the system size in a small sample experiment, we might expect a morphology transition from bands to islands as $`w_0`$ is increased.
Figure 13 shows the microstructures obtained in simulations for small $`w_0`$, ranging from $`5\times 10^3`$ to $`5\times 10^5`$. The lateral system size is about twice the diffusion length, and we use periodic boundary conditions in the direction perpendicular to the temperature gradient. The lateral spreading of multiple nuclei leads to a jagged morphology. Each V-shaped site in the figures indicates a nucleation event (either $`\beta `$ on $`\alpha `$ or $`\alpha `$ on $`\beta `$). A transition from irregular banded structures to islands can be observed as $`w_0`$ increases. Note, however, that nucleation events occur in bursts, leading to a spatial periodicity along the growth direction that can be clearly distinguished in Fig. 13(c). This means that there is still a “banding cycle,” now consisting of layers of a two-phase composite structure (particulate structure) and layers of pure $`\alpha `$ matrix.
The values for $`w_0`$ in Fig. 13 are somewhat small. The attempt rate $`w_0`$ should be related to the rate $`I_0`$ in Eq. (47), which is typically about $`10^{30}`$ nuclei/$`\mathrm{cm}^3\mathrm{s}`$ for heterogeneous nucleation in metallic systems . Hence we investigated the microstructures formed with larger values of $`w_0`$, ranging from $`5\times 10^{11}`$ to $`5\times 10^{43}`$. In this range, we always obtain island structures that look qualitatively similar (Fig. 14), and not too different from Fig. 13(c). This implies that the microstructures obtained with multiple nucleation events are not very sensitive to $`w_0`$ when $`w_0`$ is larger than some critical value. Hence we would expect predominantly island structures in spatially extended systems.
The last statement, however, is valid only for the quite restrictive assumptions made in our model. Most importantly, we have assumed that the probability of nucleation depends only on the composition in the liquid, and not on the local geometry of the interface. This neglects the presence of grain boundaries and impurities (bubbles, inclusions) which can considerably enhance nucleation. Such heterogeneities broaden the distribution of the nucleation rate as a function of temperature, and would hence favor bands. The evolution of the grain structure could in principle be modeled by including the local crystalline orientation as an additional order parameter. An interesting perspective is that the interplay between nucleation at grain boundaries and spreading might select a certain grain size, since for small grains a spreading phase can engulf and hence “heal” grain boundaries, whereas for large grains nucleation at the solid-liquid interface (as modeled in our simulations) may occur and lead to the formation of new grains. Such a study, however, is largely beyond the scope of this article.
## VI Summary and Conclusions
We have developed a phase-field model to investigate a class of banded microstructures that form during the directional solidification of peritectic alloys under purely diffusive growth conditions. We focused on a regime of large thermal gradients and low pulling speeds where both phases are morphologically stable and the interface dynamics is controlled by a subtle interplay between the growth and nucleation of two competing solid phases, rather than by the morphological instability of one phase. We restricted our attention to a generic peritectic phase diagram that simplified both the models and the computations, but our approach is in principle flexible enough to be extended to phase diagrams of specific materials.
The two-dimensional simulations of this model have shed light on three main aspects of banding: the transition from islands to bands that has been observed in narrow samples where convection is suppressed , the associated dynamical spreading of one phase onto the other, and the type of structures that one would expect to form in wide samples under purely diffusive growth conditions that are not presently accessible in an earth-based laboratory, at least for the alloys investigated to date.
We have shown that the transition from islands to bands can be understood in terms of a competition between the lateral spreading of $`\beta `$ on $`\alpha `$ and the diffusive backflow of rejected impurities from $`\alpha `$ to $`\beta `$. This competition leads to the surprising result that bands tend to form more easily in wider samples, in qualitative agreement with recent experiments in the Sn-Cd alloy carried out in small samples in order to suppress convection . The critical system size $`L_c`$ at which this transition occurs depends on the nucleation undercooling for the $`\beta `$ phase that influences the spreading rate, and on the alloy composition, with $`L_c`$ becoming infinite when the volume fraction of the $`\beta `$ phase falls below a minimum value necessary for band formation. The influence of other parameters, in particular the form of the phase diagram and the equilibrium angles at the trijunction, has not been investigated in detail here.
When the peritectic phase fully covers the parent phase, its spreading dynamics is characterized by a remarkably uniform acceleration of the moving trijunction that depends on the composition, but not on the nucleation undercooling. This acceleration originates from an increase with time of the local supersaturation (driving force for spreading) associated with the relaxation of the planar parent phase ahead of the trijunction to its steady state below $`T_p`$, together with a direct relationship between the instantaneous speed of the trijunction and this driving force that depends on the material properties, but not on the history or the overall composition of the sample. Moreover, the relative angles between phase boundaries at the trijunction during rapid spreading depart significantly from those prescribed by Young’s condition, indicating a strong departure from local equilibrium. Both predictions might be experimentally testable in transparent organic eutectic systems that exhibit similar spreading transients before coupled growth is established .
The formation of multiple nuclei in wide samples ($`LL_c`$) adds a stochastic element to the interface dynamics that renders the range of possible patterns even richer. One can nonetheless distinguish two basic types of structure that can be understood within the framework of the single island to band transition in narrow samples ($`LL_c`$). The first is a discrete banded structure made up of separate jagged bands that span the whole width of the sample. The second is a particulate banded structure made up of approximate rows of particles (islands) of the peritectic phase embedded in the matrix of the parent phase. The banded (particulate) structure is naturally selected if the mean distance between nuclei is larger (smaller) than the critical sample width $`L_c`$ for the island-band transition, and, moreover, simulations reveal that the particulate structure is preferred if nucleation is assumed to follow a classical nucleation law. Even though we modeled patterns in wide samples with such a law, we expect the transition from a discrete to a particulate banded structure to be generally governed by the mean distance between nuclei even if other nucleation mechanisms (such as wall-induced nucleation and nucleations at grain boundaries) play a dominant role. Both types of structure could conceivably coexist in the same sample if nucleation conditions change during growth.
There are a number of possible extensions of the present study. One is to investigate the patterns that form for a somewhat larger range of pulling speeds where the parent phase is morphologically unstable, but the peritectic phase is still linearly stable. Another is to incorporate the influence of convection in a fully consistent way to make contact with experiments over a wider range of sample sizes, which is now possible within a phase-field context .
###### Acknowledgements.
This research has been supported by the NASA Microgravity Research Program under Grant No. NAG8-1254 and by U.S. DOE grant No. DE-FG02-92ER45471. Computations were made possible by an allocation of time at the Northeastern University Advanced Scientific Computation Center (NUASCC). We thank J. S. Park, R. Trivedi, S. Akamatsu, and G. Faivre for many helpful discussions and for giving us access to their experimental results prior to publication.
## A Results of the common tangent construction
The common tangent construction allows one to determine the equilibrium composition for two-phase equilibrium for given bulk free energies $`f_\nu `$ of the two phases. For two-phase equilibrium, the bulk phases must have equal chemical potentials $`\mu =df_\nu /dc`$ and grand potentials $`\mathrm{\Omega }=f_\nu \mu c`$. Solving the resulting equations for our model bulk free energies, we find for solid-liquid equilibrium
$`c_L^{\alpha ,\beta }`$ $`=`$ $`{\displaystyle \frac{(B_1\pm \frac{1}{2}B_2)}{(A_1\pm \frac{1}{2}A_2)}}+A_1,`$ (A1)
$`c_S^{\alpha ,\beta }`$ $`=`$ $`{\displaystyle \frac{(B_1\pm \frac{1}{2}B_2)}{(A_1\pm \frac{1}{2}A_2)}}A_1A_2.`$ (A2)
The upper (lower) sign is for the $`\alpha `$ ($`\beta `$) phase.
For solid-solid equilibrium, we get
$`c_{SS}^\alpha `$ $`=`$ $`{\displaystyle \frac{B_2}{A_2}}A_1A_2,`$ (A3)
$`c_{SS}^\beta `$ $`=`$ $`{\displaystyle \frac{B_2}{A_2}}A_1+A_2.`$ (A4)
For convenience, we define
$$\overline{A}=\left(A_1\pm \frac{1}{2}A_2\right)$$
(A5)
and
$$\overline{B}=\left(B_1\pm \frac{1}{2}B_2\right)$$
(A6)
where the upper (lower) sign is for the $`\alpha `$ ($`\beta `$) phase.
In order to relate the parameters in our model to a physical system, let us write
$`B_1`$ $`=`$ $`B_{11}+B_{12}\stackrel{~}{T},`$ (A7)
$`B_2`$ $`=`$ $`B_{21}+B_{22}\stackrel{~}{T}.`$ (A8)
Here, $`c`$ and $`\stackrel{~}{T}`$ are the scaled composition and temperature, respectively, defined in the text. Let $`\mathrm{\Delta }C_\nu `$ and $`m_\nu `$ be the concentration jump at the solid-liquid interface and the liquidus slope of phase $`\nu `$, respectively, at $`T_p`$. Let $`r`$ be the ratio $`\mathrm{\Delta }C_\beta /\mathrm{\Delta }C_\alpha `$, then the parameters $`A_1`$, $`A_2`$, $`B_{11}`$, $`B_{12}`$, $`B_{21}`$, and $`B_{22}`$ are related to these quantities in the phase diagram by
$`A_1`$ $`=`$ $`{\displaystyle \frac{1}{4}}(1+r),`$ (A9)
$`A_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1r),`$ (A10)
$`B_{11}`$ $`=`$ $`{\displaystyle \frac{1}{4}}(1+r)\left[r{\displaystyle \frac{1}{4}}(1+r)\right],`$ (A11)
$`B_{21}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1r)\left[r{\displaystyle \frac{1}{4}}(1+r)\right],`$ (A12)
$`B_{12}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(1+r{\displaystyle \frac{m_\alpha }{m_\beta }}\right),`$ (A13)
$`B_{22}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1r{\displaystyle \frac{m_\alpha }{m_\beta }}\right).`$ (A14)
In order to have vertical solid-solid coexistence lines, we have to choose the parameters such as to make $`B_{22}`$ vanish. The parameters for our model peritectic system are listed in Table II.
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# Hipparcos, IUE, and the Stellar Content of the Solar Neighbourhood
## 0.1 Introduction
Understanding how the complex near-UV region of the spectra of stars is shaped within the stellar atmospheres can provide a homogeneous source of information on several of the fundamental stellar parameters, the chemical composition, magnetic activity, rotational velocity, atmospheric velocities, and ages. Many neutron capture elements, whose abundance in metal-deficient stars keeps log of the supernovae rates along the galactic evolution, produce undetectable (or subjected to large measurement errors) features in the optical spectrum, but strong lines in the near-UV (see e.g. Sneden et al. 1998). Interesting absorption lines in the spectrum of light atoms, such as boron and beryllium (see e.g. García López et al. 1998), are only present in the near-UV. Furthermore, F-type main sequence stars dominate this region of the spectrum for intermediate-age stellar populations, and therefore constitute a powerful prospective tool for dating galaxies (Heap et al. 1999).
A lack of proper understanding has prevented full use of the data gathered by the space observatory IUE to study the stellar population in the proximity to the Sun. Indeed all the available determinations of the metallicity distribution in the solar neighbourhood, except for the polemical work by Favata et al. (1997), are based on photometric calibrations (Twarog 1980, Wyse & Gilmore 1995, Rocha-Pinto & Maciel 1998, Flynn & Morell 1997). The star counts do not fit, and the so-called G-dwarf problem, the deficiency in relative numbers of metal-poor stars in the main sequence, seems to extend to other spectral types as well to the more metal-rich end of the metallicity distribution (Rocha-Pinto & Maciel 1998). Ideally, to extract the available information from the near-UV spectrum, the model atmospheres and other spectral synthesis ingredients should be carefully tuned for each of the stars: the abundances and abundance anomalies of all relevant elements, microturbulence, effective temperature, gravity, interstellar extinction, etc. However, as a first step I am neglecting here the abundance anomalies, variations in microturbulence, and the effect of interstellar extinction, to derive only the effective temperatures ($`T_{\mathrm{eff}}`$s) and overall metallicities (\[Fe/H\]s). Selecting a sample limited in volume to 100 pc, makes the availability of Hipparcos (ESA 1997) parallaxes very likely for the stars observed by IUE, and places a safe limit on the role of the interstellar extinction.
## 0.2 Observations and analysis
The Hipparcos Catalogue includes data for 22982 stars within 100 pc from the Sun. Making use of the MAST Cross Correlation Search Tool, I have identified 3421 low-resolution LW ($`18003500`$ Å) spectra of 992 of such stars. The IUE (NEWSIPS) observations have been retrieved from the Villafranca node of the IUE Final Archive in Spain. A newer version of the archive has being released recently (INES; Rodríguez-Pascual et al. 1999). When more than a single spectra was available for a given star, they were combined and cleaned using the IUEDAC IDL Software libraries to produce a single spectrum per star.
I have made use of the flux distributions calculated by Kurucz, and available at CCP7 since 1993. The grid includes models for different gravities (log g), effective temperatures ($`T_{\mathrm{eff}}`$) and metallicities (\[Fe/H\]), while the parameters in the mixing-length treatment of the convection are fixed, as well as it is the microturbulence (2 km/s), and the abundance ratio between different metals (solar proportions). For a given set of ($`T_{\mathrm{eff}}`$, log g, \[Fe/H\]), I obtain the theoretical flux from linear interpolation, therefore using the information of the eight nearest models available in the grid, which is divided in steps of 200 K in $`T_{\mathrm{eff}}`$, 0.5 dex in logg, and 0.5 dex in \[Fe/H\].
Making use of an accurate stellar parallax (p), $`BV`$ photometry and state-of-the-art evolutionary isochrones (Bertelli et al. 1994), one can estimate the stellar radius (R) with an accuracy of roughly 6% (Allende Prieto & Lambert 2000, Lambert & Allende Prieto 2000), and get a small error in the determination of the ratio (pR)<sup>2</sup>, which allows to transform the absolute flux measured at Earth to flux emerging from the stellar atmosphere. It is as well possible to constrain the mass to within 8% (in the range of metallicities we are interested on) and, therefore, the gravity within 0.07 dex. Once the gravity and the dilution factor for the flux are fixed, it is possible to compare the absolute near-UV fluxes measured by IUE with theoretical fluxes, and determine $`T_{\mathrm{eff}}`$ and \[Fe/H\]. As showed by Allende Prieto & Lambert (2000), knowing the absolute visual magnitude and the $`BV`$ color index, comparison with evolutionary isochrones provides an independent estimate of the $`T_{\mathrm{eff}}`$, precise to roughly 2% for stars with $`4500<`$ $`T_{\mathrm{eff}}`$ $`<8500`$ K, which can be used to check for systematic errors (see below).
For each of the analyzed stars I first derive its gravity and the flux dilution factor, then the values of \[Fe/H\] and $`T_{\mathrm{eff}}`$ are obtained by finding the minimum of the square of the difference between the observed and the synthetic spectrum. Previously, the observed spectrum, which has a resolution between 5.2 and 8.0 Å is degraded to that of the synthetic spectra, roughly twice poorer. The search for the optimum ($`T_{\mathrm{eff}}`$,\[Fe/H\]) pair is performed using the Nelder-Mead simplex method, as implemented by Press et al. (1988). Figure 1 shows two examples of the typical goodness-of-fit achieved.
## 0.3 Discussion
The sample was restricted to stars with $`3500<`$ $`T_{\mathrm{eff}}`$ $`<10000`$ K. The employed models are known to be inappropriate close to and below the cooler limit, as the plentiful molecules are not properly taken into account, and the number of nearby stars beyond the upper limit is very small.
A first look at the comparison between the ’evolutionary’ and near-UV $`T_{\mathrm{eff}}`$s reveals systematic differences. The left panel of Figure 2 shows the comparison for the 253 stars whose IUE spectra have perfect quality flags in the considered spectral range (2000-3000 Å). Stars with strong Mg II 2852 Å emission are identified with rhombi, those with a continuum in the region 2000-2400 Å much stronger than the models’ prediction with open circles, and the rest with filled circles. Disregarding the stars for which there is indication of a chromospheric component in the spectrum (open symbols), the spectroscopic $`T_{\mathrm{eff}}`$s of ’normal’ stars exhibit a systematic difference from the evolutionary $`T_{\mathrm{eff}}`$s. Similar – although smaller– effects could be present in other $`T_{\mathrm{eff}}`$ scales that make use of the same type of model atmospheres, such as that derived from the InfraRed Flux Method (IRFM; e.g. Blackwell & Lynas-Gray 1994). di Benedeto (1998) showed indeed that the IRFM $`T_{\mathrm{eff}}`$s of A stars are about 2.3% lower than his empirical scale. The 150 stars whose spectra do not show evidence for a chromosphere are spread out the \[Fe/H\]–$`T_{\mathrm{eff}}`$ plane as shown in the right panel of Figure 2.
It is obviously very difficult to determine the selection effects of the sample, as the stars are required to have been observed by IUE. It is likely that peculiar objects were favoured: multiple systems, stars with abundance anomalies, pulsating/variable stars, metal-poor stars, etc. That explains the excess of stars with $`1.0<`$ \[Fe/H\] $`<0.5`$ find here, in comparison with every other photometric study in the literature. Many of theses stars would indeed be chemically peculiar stars (not being deficient in all metals, but having a low abundance of some of the species relevant to the atmospheric structure and the near-UV opacities) or spatially-unresolved binary (multiple) systems. In both cases the stars would tend to be erroneously shifted towards hotter $`T_{\mathrm{eff}}`$s and lower metallicities, in order to compensate for the excess flux. The same argument could provide an explanation for the lack of stars with $`T_{\mathrm{eff}}`$ $`>7000`$ K at near/super-solar metallicities: the peculiar stars, preferred by IUE observers, would tend to have lower metallicities. Further study is clearly needed to shed light on this issue.
## 0.4 Conclusion
Several issues have been raised in this preliminary analysis, and are fundamental. The first of all is that it is necessary to identify ALL the basic elements (chemical and non-chemical) that affect the shape of the near-UV spectrum. They may be quite different for different types of stars. Once those elements are identified, we should attempt the fit of the spectrum allowing for ALL to vary. The parameters determined should be compared with other scales, of empirical or independent nature, to asses the real possibilities of this kind of analysis and the performance of the model atmospheres, opacities, and other modelling ingredients. That procedure will likely help improving the modelling. It will very likely shed light on another basic problem still far from solved: the effective temperature scale.
At that stage, we should look back to the application of near-UV spectroscopic analysis to other affairs, such as the stellar content of local neighbourhood, and review then the apparent scarcity of solar and super-solar metallicity A-F stars.
###### Acknowledgements.
I am grateful to Benjamín Montesinos and Enrique Solano for their help with the IUE spectra, and to David Lambert for interesting discussions. This research has made use of the cross-correlation tool of the Multimission Archive at the Space Telescope Science Institute (MAST). STScI is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5-26555. Support for MAST for non-HST data is provided by the NASA Office of Space Science via grant NAG5-7584 and by other grants and contracts. I have made extensive use of the IDL software libraries developed at the IUE Data Analysis Center (IUEDAC). All the data used in this work were retrieved from the IUE data server at Villafranca del Castillo satellite tracking station of the ESA. The SIMBAD database, operated at CDS and the NASA ADS were very useful in this work.
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# InAs-AlSb quantum wells in tilted magnetic fields
## I Introduction
The energy spectrum of two-dimensional electron gases (2DEG) in magnetic fields of arbitrary orientation is fairly well understood . Most considerations follow a single-particle approach which is powerful to explain several of the experimentally observed features. For magnetic fields tilted with respect to the sample normal one finds that the Landau splitting, which is proportional to the component of the field perpendicular to the 2DEG can be tuned with respect to the Zeeman splitting which is proportional to the total magnetic field. This is used in the so-called coincidence method where appearance and disappearance of minima in Shubnikov-de Haas oscillations (SdH) as a function of tilt angle is observed in magnetotransport experiments. The analysis in terms of a picture of non-interacting electrons has proven very powerful for the analysis of energy spectra in Si-MOSFETs , InAs-GaSb superlattices , InAs-GaSb quantum wells , GaAs-AlGaAs heterostructures , GaInAs/InP heterostructures and Si/SiGe heterostructures . In this paper we focus on InAs-AlSb quantum wells and extend preliminary studies on this material system . We present several features that are perfectly well explained in the existing single-particle picture, namely 1. the appearance and disappearance of even- and odd-integer SdH minima as a function of tilt angle, 2. a Zeman splitting as large as five times the Landau splitting for tilt angles around 87, and 3. a g-factor for InAs of about 13 in agreement with considerations based on conduction band non-parabolicity . In contrast to straight forward expectations we find 4. non-vanishing SdH minima for even-integer filling factors $`\nu =4,6,8`$ in the range of tilt angles and magnetic fields where these filling factors can be observed and 5. a regime at low magnetic fields where even-integer filling factor SdH minima persist for all tilt angles, while the usual coincidence features occur at higher magnetic fields. These observations are discussed in view of other experiments and theoretical ideas based on exchange enhancement .
## II Level crossing in the single particle regime
All samples contained 15nm wide InAs quantum wells, confined by AlSb or Al<sub>x</sub>Ga<sub>1-x</sub>Sb $`(x0.8)`$ barriers. The sample details are summarized in . Our samples are of very high quality and have mobilities up to 84 m<sup>2</sup>/Vs. In this paper we focus on a sample with a GaSb cap and a carrier density of $`N_s=6.210^{11}`$cm<sup>-2</sup> (UCSB \# 9503-18). The samples were patterned into geometries suitable for transport experiments and equipped with Ohmic contacts to the 2DEG.
The samples were mounted on a revolving stage in several cryostat environments. The angle $`\alpha `$ is measured between the magnetic field orientation and the sample normal. For the data taken at T=1.7 K and magnetic fields up to 8 T the revolving stage was computer controlled. Consequently very dense data sets were obtained. We also measured the samples in a dilution refrigerator at sample temperatures down to 100 mK and magnetic fields up to 15 T as well as in a <sup>3</sup>He system with a base temperature of about 400 mK and magnetic fields up to 28 T. Because of the Landau level broadening the temperature dependence of the SdH oscillations basically levels off below 1.7 K. The difference in experimental resolution of the three setups is mostly determined by the respective measurement electronics.
The results obtained on different samples depend on the carrier concentration. The filling factor is definded by $`\nu =N_sh/eB`$, where $`N_s`$ is the electron density of the 2DEG. For perpendicular fields, i.e. $`\alpha =0`$, all observed features at magnetic fields $`B1.5`$ T, where the Zeeman splitting is not yet resolved, can be analyzed with one single SdH period with very high accuracy . Effects of inversion asymmetry induced zero-field spin-splitting are therefore not considered. From the largest filling factors that we can observe we estimate the Landau level width to about 0.4 meV.
We can follow the disappearance and reappearance of minima at even- and odd-integer filling factors as a function of increasing tilt angle. This interplay between pronounced even- and odd-integer filling factor minima occurs for a series of angles. It also shows up in the respective quantum Hall plateaus .
Figure 1 shows magnetoresistance traces at specific tilt angles where either SdH minima occur only at even-integer filling factors, at even and odd, or only at odd integer filling factors. The amplitude of the SdH oscillations at large tilt angles is magnified with respect to the other traces. The angles are determined by measuring $`\rho _{xy}`$ with high accuracy. The absolute error in the angle becomes larger with increasing tilt angle because of the $`\mathrm{cos}(\alpha )`$-dependence.
We find that the carrier density decreases by up to 5% if parallel magnetic fields larger than 20 T are applied. This also shows up in a non-linear Hall effect for large parallel fields. We attribute this behavior to magnetic freeze-out of carriers due to a redistribution of the electrons from the well into some localized states. The reason for this could be a strong diamagnetic shift of the quantum well state. This effect has no consequences for the results presented in this paper but explains why the SdH minima in Fig. 2 for large tilt angles $`\alpha `$ do not exactly fall onto the dashed lines.
The inset in Fig. 2 describes the various coincidence situations which are characterized by the parameter $`r`$, the ratio of Zeeman and cyclotron energies.
$$r=\frac{g\mu _\mathrm{B}B_{tot}}{\mathrm{}\omega _c}$$
Here $`\omega _c=eB_{}/m^{}`$, $`m^{}`$ is the effective electron mass, $`\mu _\mathrm{B}`$ is the Bohr magneton and $`B_{}=\mathrm{cos}(\alpha )B_{tot}`$. We thus arrive at
$$r\mathrm{cos}(\alpha )=\frac{gm^{}}{2m_e},$$
where $`m_e`$ is the free electron mass. The data in Fig. 1 shows the resistance traces at $`r`$-values always close to the indicated numbers of 1/2, 1, 3/2,…. The larger the tilt angle, the more difficult it is to realize a given coincidence situation accurately since the span of angles at which it takes place decreases with $`\mathrm{cos}\alpha `$. Nevertheless we demonstrate that SdH oscillations can be measured in a situation where the Zeeman splitting is 5 times larger than the Landau splitting.
Figure 2 shows the coincidence situations plotted as $`1/\mathrm{cos}(\alpha )`$ versus $`r`$. The slope of this curve is proportional to the product $`gm^{}`$. We determined the effective mass for this sample by temperature dependent SdH measurements and found a value for the effective mass of $`m^{}=(0.032\pm 0.002)m_e`$ which is in agreement with values reported in the literature. Using this value for $`m^{}`$ we computed $`|g|13`$.
Such experiments have been performed on a series of samples. In first approximation the obtained data can be described by using Landau levels and spin levels behaving and crossing as expected in a single-particle model. Because the $`g`$-factor is so large effects of electron-electron interactions, the so-called exchange enhancement , are expected to be relatively small. Furthermore these effects should increase for decreasing filling factors. In our case the experimental data can best be described with a product $`gm^{}`$ which is constant over the investigated range of magnetic fields and angles.
The effects of non-parabolicity can be estimated using a $`𝐤𝐩`$ formalism which in its simplest case reduces to the two-band model.
$$m^{}(E)=m^{}(E=0)\left(1+2\frac{E}{E_\mathrm{g}}\right)$$
Here $`E_\mathrm{g}=400`$ meV is the band gap of InAs and $`E`$ is the electron energy relative to the conduction band edge. Because of the huge conduction band offset between InAs and AlSb (1.35 eV), we use the model of a quantum well with infintely high walls. The total energy $`E`$ can, to a good approximation, be written as the sum of an approximate Fermi energy $`E_\mathrm{F}=N_s\pi \mathrm{}^2/m^{}`$ and an approximate confinement energy $`E_c=\mathrm{}^2/2m^{}\pi ^2/a^2`$, where $`a`$ is the quantum well width. With this we obtain for the density dependence of the effective mass in the two-band model
$$m^{}(N_s)=\frac{m_0^{}}{2}+\frac{m_0^{}}{2}\sqrt{1+\frac{8}{E_\mathrm{g}}\left(\frac{\mathrm{}^2}{2m_0^{}}\frac{\pi ^2}{a^2}+\frac{\pi \mathrm{}^2}{m_0^{}}N_s\right)}$$
Here $`m_0^{}=m^{}(E=0)`$, i.e. the effective mass at the conduction band edge, which for InAs is $`m_0^{}/m_e=0.023`$. We find $`m^{}(N_s=4.410^{11}`$cm$`{}_{}{}^{2})/m_e=0.032`$ in agreement with our experimentally determined value. The values for the energies are $`E_\mathrm{F}=52`$ meV and $`E_c=51.6`$ meV.
At the same time the $`g`$-factor is reduced in agreement with our experimental findings. For the $`g`$-factor the two-band model results in
$$g(E)=g(E=0)(1\alpha E)$$
The parameter $`\alpha `$ is estimated in Ref. to be $`\alpha =0.0025`$1/meV for a quantum well system very similar to ours. This results in a $`g`$-factor of $`|g|=12`$ very close to our experimental result. Using the expressions for the $`g`$-factor and the effective mass one finds that the total effect of nonparabolicity on the product $`gm^{}`$ almost cancels out.
Several additional aspects should be considered in this discussion. For large tilt angles the in-plane magnetic field component can be as large as 10 T. In this case it is well known that the Fermi surface is no longer a circle but an ellipse. The effective mass thus depends on $`B`$ . We measured the temperature dependence of the SdH oscillations in tilted magnetic fields in order to extract the effective mass as a function of field and tilt angle. Within the experimental accuracy we found that the effective mass is constant to 5% in the investigated parameter regime. On the same footing one also expects that the $`g`$-factor becomes a magnetic field dependent quantity. With these complications in mind one has to take the analysis of the product $`gm^{}`$ from the plot in Fig. 2 with a grain of salt.
## III Level anti-crossings at small filling factors
Figure 3 shows magnetoresistance traces down to even-integer filling factors of $`\nu =6`$. We only present the range of angles where the situation corresponding to $`r=1`$ occurs. The tilt angle is changed in rather small increments which are monitored by the change in the Hall resistance $`\rho _{xy}`$. Similar but less pronounced features alos occur in a sample has a lower carrier density of $`N_s=4.410^{11}`$cm<sup>-2</sup>. The highest perpendicular magnetic fields correspond to total magnetic fields of 28 T. For $`\alpha =73.5`$, minima occur for even- and odd-integer filling factors. As the tilt angle increases, even-integer minima weaken until about $`78.8^{}`$ and then increase again in strength. They never completely disappear even up to filling factors of $`\nu `$=16. This means that there always remains a minimum of the density of states at the Fermi energy when the single-particle model predicts a crossing of spin and Landau levels.
An anti-crossing of single particle levels has been predicted for filling factor $`\nu =2`$ based on the transition from a spin-unpolarized state at small tilt angles to a spin polarized state at large tilt angles. Experimental data obtained on GaInAs/InP heterostructures showed the expected single particle behavior for low-mobility samples while a non-suppression of the SdH minimum at $`\nu =2`$ for high-mobility samples was observed. This was interpreted in the framework of the formation of a spin-polarized ground state induced by the strong parallel magnetic field. In the case of Ref. the SdH minimum corresponding to filling factor $`\nu =4`$ and higher even-integer filling factors were perfectly well suppressed at the same tilt angle as expected in a single particle model. The authors argued that for low mobility samples and higher integer filling factors neighboring levels overlap due to their broadening and the exchange interaction cannot help to further open the gap.
The experimental situation in our case is different in the following aspects. The SdH minima at even-integer filling factors weaken but do not disappear. Furthermore their weaking goes hand in hand with their overall appearance, i.e. the sudden importance of an exchange driven opening of a gap cannot be observed. Unfortunately the carrier density in our samples is too high to observe the behavior of SdH minima corresponding to filling factors $`\nu =2`$ and $`\nu =4`$ at large tilt angles and experimentally accessible magnetic fields.
In order to get an understanding of the energy struture in tilted magnetic fields we calculated the magnetoresistance following Gerhardts . We included a constant background density of states in order to model the broad minima in the magnetoresistance. Based on the single particle energies
$$E_{s,n}=\mathrm{}\omega _c\left(n+\frac{1}{2}\right)+sg\mu _\mathrm{B}B,n=0,1,2,\mathrm{},s=\pm \frac{1}{2}$$
an anti-crossing between neighboring levels of $`\mathrm{\Delta }E=0.29\mathrm{}\omega _c`$ was inserted in the model. At $`B_{}`$=4.2 T ($`\nu =6`$) and $`m^{}=0.032m_e`$, this corresponds to $`\mathrm{\Delta }E=4.4`$ meV. We assumed a Gaussian Landau level broadening $`\mathrm{\Gamma }=\mathrm{}/\tau _q=1.5`$meV with $`\tau _q=0.45`$ ps.
The magnetic field dependence of the anti-crossing was approximated with a smooth parabolic curvature.
Figure 4 shows calculated resistance traces. There is at least qualitative agreement between the calculated (Fig. 4) and experimental (Fig. 3) data sets. From the simulation it is obvious that the situation where even-integer minima in the SdH oscillations are weakened or even suppressed extends over a significantly larger range of angles compared to the experiment. This could arise from our rough modelling but also hints at the importance of interaction effects for the details of SdH oscillations.
What could be the reason for the persistent appearance of even-integer SdH minima in the regime where the underlying single particle energy levels are expected to cross? For small filling factors the effects of exchange enhancement have been demonstrated in various experiments (for a recent example see and references therein). From our experimental data at small filling factors we do not see an indication that eletron-electron interactions in terms of exchange enhancement play a significant role. For the case of perpendicular magnetic fields, $`\alpha =0`$, the energy levels are described by three quantum numbers, namely subband, Landau and spin quantum numbers. This is based on the fact that the Hamiltonian can be separated into a part describing the electron motion in the plane of the 2DEG and another part responsible for the quantization in growth direction. For tilted magnetic fields mixed levels arise whose degenracy is still completely controlled by the perpendicular magnetic field component . For the InAs-AlSb system this approach has to be extended in order to incorporate the strong conduction band non-parabolicity of InAs, as well as the possible strain in the well due to the different lattice constants of barrier, well and GaAs substrate. One can envision that such effects already lead to possible level couplings and anti-crossings as observed in the experiment.
## IV Even-integer SdH minima at low magnetic fields
Figure 5 presents a grey-scale plot composed of magnetoresistance traces taken at very closly spaced tilt angles around the regime of $`r=1`$ and $`r=2`$. Here we focus on the regime of small magnetic fields. For $`\alpha =65^{}`$ SdH minima occur at even-integer filling factors. As the tilt angle is increased, odd-integer minima take over at magnetic fields $`B_{}0.8`$ T and gradually disappear again in favor of even-integer minima. At magnetic fields below 0.8 T minima occur only at even-integer filling factors over the whole range of tilt angles. The inset of Fig. 5 shows a representative resistance trace at an intermediate tilt angle where the SdH oscillations are dominated by even-integer minima at low magnetic fields, a crossover regime and odd integer filling factors at higher magnetic fields. A beating pattern would not display such a phase shift in the pattern of the oscillations.
Starting from the Landau and spin levels in tilted magnetic fields such a behavior can occur in two ways: either the Landau energy is not exactly proportional to the perpendicular component of the magnetic field, or the Zeeman splitting is angle dependent. Both effects have been discussed to some extent before. Non-parabolicity effects are most likely a minor contribution for such small magnetic fields. The large in-plane magnetic field component, which can lead to an anisotropic effective mass dispersion, should become more important for larger magnetic fields. However, the unusual behavior as presented in Fig. 5 occurs in the low-magnetic field regime.
Leadley et al. have shown that there is a critical collapse of the exchange enhanced spin splitting in two-dimensional systems . The authors found that the total spin splitting is a sum of the bare Zeeman splitting proportional to the total magnetic field and a contribution due to exchange enhancement which is proportional to the perpendicular component of the magnetic field.
$$\mathrm{\Delta }_{spin}=g_0\mu _\mathrm{B}B_{tot}+\beta \mathrm{}eB_{}/m^{}$$
For the case of GaAs heterostructures, Leadley at al. found $`\beta =0.2`$ independent of magnetic field. In their case $`g_0`$ is the bare $`g`$-factor because non-parabolicity effects are negligible in GaAs. In our case $`g_0`$ has to be identified with $`g(E)`$ where the non-parabolicity contribution stems from the position of the Fermi energy above the conduction band edge and does not depend on magnetic field in the investigated range of parameters.
In the regime of large magnetic fields discussed before, where spin splitting is well resolved, we found that the exchange enhancement is a minor contribution. However, for small magnetic fields and large tilt angles the exchange contribution could play an important role. If the bare spin splitting is smaller than the Landau level broadening, the exchange enhancement is not expected to play a role. In this case even-integer SdH minima will dominate the magnetoresistance for all tilt angles. Once the bare Zeeman splitting approaches and exceeds the Landau level broadening the exchange enhancement will further increase the spin gap and the usual coincidences between Landau and spin levels will take over.
For any functional dependence of g on B which is smooth one would not expect a sudden crossover from even-integer to odd-integer minima as depcited in the inset of Fig. 5. The sudden change in periodicity over a small magnetic field range requires a mechanism which leads to an abrupt opening of the spin gap similar as it has been observed in Ref. for the cititcal collapse of the exchange enhanced spin-splitting.
## V Summary
We have presented a series of SdH measurements on InAs-AlSb quantum wells in tilted magnetic fields. In a reasonable range of parameters the experimental results can be understood in a straight forward single particle model. The coincidence method is based on independent Landau and spin levels. This way we obtain reasonable numbers for the effective mass and g-factor that agree with results of a two-band model and experimental results of others. For large magnetic fields we find an anti-crossing of neighboring Landau and spin levels. Most likely this is not a consequence of electron-electron interactions. We speculate that this effect arises from the pronounced non-parabolicity of the InAs conduction band as well as from the built-in strain in such samples. For very small magnetic fields SdH minima exist only at even-integer filling factors independent of tilt angle. This is attributed to a critical filling factor necessary for the observation of spin-splitting
We are grateful to R. Warburton and S. Ulloa for helpful discussions and thank ETH Zürich and QUEST for financial support. The hospitality of the High-Magnetic Field Laboratory in Grenoble is gratefully acknowledged.
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# 1 Clifford algebras
## 1 Clifford algebras
Clifford algebras solve an algebraic existence problem. To see this recall that the field of complex numbers arises in two ways. In the first instance it is merely a vector space that helps parametrize the Euclidean plane $`^2`$ but in the second it is an algebra extending the real number field in which square roots exist and which contains an image of the group of rotations. In particular, only by this property we comprehend the law of multiplication of two negative numbers: $`(1)(1)=1`$, since $`1=i^2`$ is the composition of two rotations by $`90`$ degrees. As is well known it took R.W. Hamilton ten years to find out in 1843 that there is no analogue in 3-space. One has to step out of ordinary space to find an algebra which contains $`^3`$ as well its rotations, viz. the skew field of quaternions. What is the appropriate generalization to arbitrary dimension? Starting from a real vector space $`E`$, one has e.g. the exterior algebra $`E`$ introduced by H.G. Grassmann in 1844. It contains $`E`$ and its multiplication $``$ is anti-commutative on basis vectors:
$$e_ie_j+e_je_i=0.$$
But then basis vectors $`e_i`$ are nilpotent, $`e_ie_i=0`$. What we really need is a new multiplication $``$ such that the basic vectors satisfy $`e_ie_i=1`$. How to come to terms with this has first been observed by W.K. Clifford in 1876:
> The system of quaternions differs from this, first in that the squares of the units, instead of being zero, are made equal to $`1`$; and secondly in that the ternary product $`\iota _1\iota _2\iota _3`$ is made equal to $`1`$.…
> I shall now examine the consequence of making, in a system of $`n`$ alternate numbers $`\iota _1,\iota _2,\mathrm{},\iota _n`$, the first of the modifications just named; namely I shall suppose that the square of each of the units is $`1`$.
After the advent of modern abstract algebra the construction of Clifford’s “geometric algebra” runs as follows. We choose an inner product $`,`$ on $`E`$, i.e., we assume a Euclidean vector space $`(E,,)`$, and with respect to this inner product we choose an orthonormal basis $`(e_i)_{1in}`$ that satisfies
$$e_ie_j+e_je_i=2\delta _{ij}=2e_i,e_j.$$
This obtains from assuming
$$vv+v,v=0$$
for any $`vE`$. Just like the exterior algebra the new algebra we are looking for can now be constructed as a quotient of the tensor algebra $`T(E)`$. Here we have to consider the two-sides ideal $`J(E)T(E)`$, which is generated by elements $`vv+v,v1`$, $`vE`$.
###### Definition 1
The $``$-algebra $`𝒞\mathrm{}(E)=T(E)/J(E)`$ (corresponding to a given Euclidean structure) is called the Clifford algebra of $`E`$. In case of $`E=^n`$ with its standard Euclidean structure we write $`𝒞\mathrm{}_n=𝒞\mathrm{}(^n)`$.
The product in $`𝒞\mathrm{}(E)`$ will be denoted by $``$, i.e. for $`u,v𝒞\mathrm{}(E)`$ with $`u=\pi (\stackrel{~}{u})`$, $`v=\pi (\stackrel{~}{v})`$, where $`\pi :T(E)𝒞\mathrm{}(E)`$ denotes the natural projection, let $`uv=\stackrel{~}{u}\stackrel{~}{v}+J(E)`$. We also denote by $`\iota _E:E𝒞\mathrm{}(E)`$ the restriction of $`\pi `$ to $`E`$.
Just like the tensor algebra and the exterior algebra the Clifford algebra solves a universal problem.
###### Theorem 1
Given an associative unital $``$-algebra $`A`$ (with unit 1) and a linear map $`f:EA`$ with $`f(v)f(v)=v,v1`$ for all $`vE`$, there is a unique homomorphism of $``$-algebras, $`\stackrel{~}{f}:𝒞\mathrm{}(E)A`$, such that the following diagram commutes
In particular, the algebra $`𝒞\mathrm{}(E)`$ together with the map $`\iota _E:E𝒞\mathrm{}(E)`$ satisfying $`\iota _E(v)^2=v,v1`$ is uniquely determined by this property up to isomorphism.
Proof: We have
$$\iota _E(v)^2=\pi (v)^2=vv+J(E)=v,v1+J(E)=v,v1\text{ for }vE$$
(here $`1=1+J(E)`$ is the unity of $`𝒞\mathrm{}(E)`$). Since $`T(E)`$ is generated by $`E`$ as an algebra and since $`\pi `$ is surjective, $`𝒞\mathrm{}(E)`$ is generated by $`\iota _E(E)`$. Now given a linear map $`f:EA`$ with $`f(v)^2=v,v1_A`$, $`vE`$, we have an extension to a homomorphism of algebras, $`f:T(E)A`$, given by
$$f(vv+v,v1)=f(v)^2+v,v1_A=0,$$
and hence factorizes to a homomorphism of algebras, $`\stackrel{~}{f}:𝒞\mathrm{}(E)A`$. Then for $`vE`$ we have
$$\stackrel{~}{f}\iota _E(v)=\stackrel{~}{f}\pi (v)=f(v)=f(v)$$
and $`\stackrel{~}{f}`$ is uniquely determined since $`\iota _E(E)`$ generates $`𝒞\mathrm{}(E)`$.
Clifford algebras have entered quite different branches of modern mathematics and physics in the 100 years since their introduction by W.K. Clifford in 1876 \[Cli\] and independently by R. Lipschitz in 1880 \[Lip\]; cf. also his letter from Hades written by his medium A. Weil \[Wei\]. Clifford’s main purpose was to generalize H.G. Grassmann’s exterior algebra and R.W. Hamilton’s quaternions, whereas Lipschitz was looking for a parametrization of orthogonal transformations of $`^n`$. That Clifford algebras indeed meet both purposes turned out in 1935, when R. Brauer and H. Weyl \[BW\] gave a very elegant representation of the spin group.
In 1954 C. Chevalley \[Che\] gave the concise construction presented above. It allows the inner product to be replaced by an arbitrary symmetric bilinear form $`\sigma :E\times E𝕂`$, or, more precisely, by the corresponding quadratic form $`Q`$, and $`𝕂=`$ or $``$ by any field. We preferably consider $`𝕂=`$ or $``$ depending on $`E`$ being a real or a complex vector space. In general, one obtains Clifford algebras $`𝒞\mathrm{}(E,Q)`$, in particular, for $`Q=0`$ the exterior algebra. On $`E=^{r+s}`$ one considers the quadratic forms
$$Q_{r,s}(x)=\underset{i=1}{\overset{r}{}}x_i^2\underset{i=r+1}{\overset{r+s}{}}x_i^2$$
yielding the Clifford algebras $`𝒞\mathrm{}_{r,s}`$. We take a look at some special examples.
Examples 1. For $`E=`$ with inner product $`x,y=xy`$ we have $`𝒞\mathrm{}()=𝒞\mathrm{}_1=`$. For if $`\iota _{}(x)=ix`$, $`x`$, the algebra $``$ is generated by $`\iota _{}()`$ since $`\iota _{}(x)^2=x,x1`$. Given an algebra $`A`$ and $`f:A`$ as above with $`f(x)^2=x,x1_A`$, we get
$$f(x)=xf(1),$$
since $`f`$ is linear, and
$$f(1)^2=1_A.$$
Defining
$$\stackrel{~}{f}(x+iy)=x1_A+yf(1),x,y,$$
we obtain a homomorphism of $``$-algebras and
$$\stackrel{~}{f}\iota _{}(y)=\stackrel{~}{f}(iy)=yf(1)=f(y),y.$$
2. The Clifford algebra $`𝒞\mathrm{}_2`$ is isomorphic with the skew field of quaternions, $``$, which is generated by $`i\iota _^2(e_1)`$ and $`j=\iota _^2(e_2)`$, since $`k=ij`$ and $`i^2=j^2=k^2=1`$, if $`\{e_1,e_2\}`$ denotes the standard basis of $`^2`$.
Remarks 1. By the universal property any isometry $`f:(E,,)(E^{},,^{})`$ induces a homomorphism of algebras $`𝒞\mathrm{}(f):𝒞\mathrm{}(E)𝒞\mathrm{}(E^{})`$: One simply has to lift the map $`\overline{f}=\iota _E^{}f`$ that satisfies
$$\overline{f}(v)^2=\iota _E^{}\left(f(v)\right)^2=f(v),f(v)^{}1=v,v1$$
as in the (commutative) diagram
Given another isometry $`g:(E^{},,^{})(E^{\prime \prime },,^{\prime \prime })`$, one has
$$𝒞\mathrm{}(gf)=𝒞\mathrm{}(g)𝒞\mathrm{}(f).$$
Therefore, $`𝒞\mathrm{}:O(E)Aut𝒞\mathrm{}(E)`$ defines a homomorphism of groups.
2. The involution $`\alpha :EE`$, $`\alpha (v)=v`$, $`vE`$, extends to an involution of $`𝒞\mathrm{}(E)`$ again denoted by $`\alpha `$. Using $`\alpha `$ one defines a $`_2`$-grading
$$𝒞\mathrm{}(E)=𝒞\mathrm{}(E)^0𝒞\mathrm{}(E)^1$$
by $`\alpha |_{𝒞\mathrm{}(E)^j}=(1)^jid`$, $`j=0,1`$; since $`\alpha `$ is a homomorphism, we have
$$𝒞\mathrm{}(E)^i𝒞\mathrm{}(E)^j𝒞\mathrm{}(E)^{i+j\mathrm{mod}2}$$
turning $`𝒞\mathrm{}(E)^0`$ into a subalgebra.
###### Proposition 1
Given $`v,wE`$ with $`v,w=0`$ one has
$$\iota _E(v)\iota _E(w)+\iota _E(w)\iota _E(v)=0.$$
More generally, given $`x=_{\mathrm{}=1}^n\iota _E(v_{\mathrm{}})𝒞\mathrm{}(E)^i`$ and $`y=_{k=1}^m\iota _E(w_k)𝒞\mathrm{}(E)^j`$ with $`v_{\mathrm{}},w_k=0`$ for all $`\mathrm{}`$ and $`k`$, one has
$$xy=(1)^{ij}yx.$$
Proof: We compute $`\iota _E(v+w)^2`$ in two ways. One the one hand
$$\iota _E(v+w)^2=v+w,v+w1=v,v1w,w1$$
and on the other hand
$$\iota _E(v+w)^2=\iota _E(v)^2+\iota _E(w)^2+\iota _E(v)\iota _E(w)+\iota _E(w)\iota _E(v).$$
Equating both sides gives the first assertion. The second one follows by induction since $`\iota _E(E)𝒞\mathrm{}(E)^1`$.
In order to prove the basic structure theorem for Clifford algebras we need the notion of graded tensor product of two graded algebras. Given two unital $``$-algebras $`A`$ and $`B`$ with units $`1_A`$ and $`1_B`$, resp., the tensor product $`AB`$ turns into an $``$-algebra if we put
$$(ab)(a^{}b^{})=aa^{}bb^{}$$
for $`a,a^{}A`$, $`b,b^{}B`$. If $`A`$ and $`B`$ are $`_2`$-graded, i.e. $`A=A^0A^1`$ and $`B=B^0B^1`$ a $`_2`$-grading of $`AB`$ is defined by
$`(AB)^0`$ $`=`$ $`A^0B^0A^1B^1`$
$`(AB)^1`$ $`=`$ $`A^1B^0A^0B^1,`$
where the product is now given by
$$(ab)(a^{}b^{})=(1)^{ij}aa^{}bb^{}$$
for $`a^{}A^i`$, $`bB^j`$. To distinguish the two tensor products, we denote the graded tensor product of $`A`$ and $`B`$ by $`A\widehat{}B`$.
###### Theorem 2
Any orthogonal splitting $`E=E_1E_2`$ gives rise to a canonical isomorphism of algebras $`𝒞\mathrm{}(E)`$ and $`𝒞\mathrm{}(E_1)\widehat{}𝒞\mathrm{}(E_2)`$.
Proof: We start with $`f:E𝒞\mathrm{}(E_1)\widehat{}𝒞\mathrm{}(E_2)`$ defined by
$$f(v_1+v_2)=\iota _{E_1}(v_1)1+1\iota _{E_2}(v_2),v_kE_k.$$
Since $`\iota _{E_k}(v_k)𝒞\mathrm{}(E_k)^1`$, $`1𝒞\mathrm{}(E_k)^0`$, and $`v_1v_2`$, we get
$`f(v_1+v_2)^2`$ $`=`$ $`\iota _{E_1}(v_1)^21+1\iota _{E_2}(v_2)^2`$
$`=`$ $`\left(v_1,v_1v_2,v_2\right)11`$
$`=`$ $`v_1+v_2,v_1+v_211`$
hence a unique homomorphism $`\stackrel{~}{f}:𝒞\mathrm{}(E)𝒞\mathrm{}(E_1)\widehat{}𝒞\mathrm{}(E_2)`$ by Theorem 1. Likewise the isometries $`i_1:E_1E`$ and $`i_2:E_2E`$ induce homomorphisms $`𝒞\mathrm{}(i_k)`$, $`k=1,2`$, and for $`x𝒞\mathrm{}(E_1)^i`$, $`y𝒞\mathrm{}(E_2)^j`$ one has
$$𝒞\mathrm{}(i_1)(x)𝒞\mathrm{}(i_2)(y)=(1)^{ij}𝒞\mathrm{}(i_2)(y)𝒞\mathrm{}(i_1)(x)$$
by the Proposition. Hence $`\stackrel{~}{g}:𝒞\mathrm{}(E_1)\widehat{}𝒞\mathrm{}(E_2)𝒞\mathrm{}(E)`$ defined by
$$\stackrel{~}{g}(xy)=𝒞\mathrm{}(i_1)(x)𝒞\mathrm{}(i_2)(y),x𝒞\mathrm{}(E_1),y𝒞\mathrm{}(E_2),$$
is a homomorphism; and a straightforward computation on generators shows that $`\stackrel{~}{f}`$ and $`\stackrel{~}{g}`$ are mutual inverses.
Remark An analogous result holds in case of a direct composition $`E=E_1E_2`$ into $`𝕂`$-vector spaces with respect to a quadratic form $`Q=Q_1Q_2`$.
Corollary Given an orthonormal basis $`(e_i)_{1in}`$ of $`(E,,)`$ one obtains a basis
$$\{\iota _E(e_{k_1})\mathrm{}\iota _E(e_{k_r})1k_1<\mathrm{}<k_rn,r0\}$$
of $`𝒞\mathrm{}(E)`$. In particular, $`dim𝒞\mathrm{}(E)=2^n`$ and multiplication in $`𝒞\mathrm{}(E)`$ is determined by the relations
$$\iota _E(e_k)\iota _E(e_k)=1,\iota _E(e_k)\iota _E(e_{\mathrm{}})+\iota _E(e_{\mathrm{}})\iota _E(e_k)=0\text{ for }k\mathrm{}.$$
Moreover one has $`𝒞\mathrm{}(E)^i=span\{\iota _E(e_{k_1})\mathrm{}\iota _E(e_{k_r})r=imod2\}`$.
Proof: We decompose $`E`$ orthogonally into
$$E=\underset{k=1}{\overset{n}{}}e_k$$
and apply Theorem 2 repeatedly using Example 1:
$$𝒞\mathrm{}(E)\left(\iota _E(e_1)\right)\widehat{}\mathrm{}\widehat{}\left(\iota _E(e_n)\right).$$
It is clear that the multiplication is determined by the given relations. From
$$\alpha \left(\iota _E(e_{k_1})\mathrm{}\iota _E(e_{k_r})\right)=(1)^r\iota _E(e_{k_1})\mathrm{}\iota _E(e_{k_r})$$
we obtain the final assertion.
From the Corollary we see that $`\iota _E:E𝒞\mathrm{}(E)`$ is injective. Therefore, we can identify $`E`$ with its image $`\iota _E(E)`$ and multiply $`v,wE`$ within $`𝒞\mathrm{}(E)`$, i.e., we write $`vw`$ instead of $`\iota _E(v)\iota _E(w)`$. We also extend the inner product of $`E`$ to the inner product of $`𝒞\mathrm{}(E)`$ that renders the basis of the Corollary an orthonormal basis. Also note that an isomorphism $`𝒞\mathrm{}_{n1}𝒞\mathrm{}_n^0`$ is induces by $`e_ke_ke_n`$, $`k=1,\mathrm{},n1`$, given an orthonormal basis $`\{e_1,\mathrm{},e_n\}`$ of $`^n`$.
Since $`E`$ and $`𝒞\mathrm{}(E)`$ have the same dimensions they are isomorphic as $``$-vector spaces although not as $``$-algebras. A canonical homomorphism $`\varphi :E𝒞\mathrm{}(E)`$ is given by
$$\varphi (v_1\mathrm{}v_k)=\frac{1}{k!}\underset{\sigma S_k}{}(sgn\sigma )v_{\sigma (1)}\mathrm{}v_{\sigma (k)}.$$
It is one-to-one since
$$\varphi (e_{j_1}\mathrm{}e_{j_k})=e_{j_1}\mathrm{}e_{j_k}$$
and actually an isometry if $`E`$ is equipped with the appropriate inner product. The inverse isomorphism $`\sigma :𝒞\mathrm{}(E)E`$ is given by
$$\sigma (x)=c(x)1E,c𝒞\mathrm{}(E),$$
where $`1=^0E`$ and where $`c:𝒞\mathrm{}(E)End(E)`$ denotes the unique extension of the linear map $`c:EEnd(E)`$ defined by
$$c(v)\omega =v\omega v\text{ }\text{ }\omega ,\omega E,vE.$$
We already mentioned the generalized Clifford algebras $`𝒞\mathrm{}_{r,s}`$. It is easily shown that they are generated by multiplying the standard basis elements $`e_1,\mathrm{},e_{r+s}`$ of $`^{r+s}`$ while respecting
$$e_ke_{\mathrm{}}+e_{\mathrm{}}e_k=\{\begin{array}{cc}2,\hfill & k=\mathrm{}r\hfill \\ 2,\hfill & k=\mathrm{}>r\hfill \\ 0,\hfill & \text{else.}\hfill \end{array}()$$
If $`n=r+s`$ is even we put $`\epsilon =e_1\mathrm{}e_n`$. Using $`()`$ we get
$$\epsilon ^2=(1)^{(n1)+(n2)+\mathrm{}+1}e_1^2e_2^2\mathrm{}e_n^2=(1)^{\frac{n(n1)}{2}}(1)^r1$$
and we call $`𝒞\mathrm{}_{r,s}`$ positive or negative if $`\epsilon ^2=+1`$ or $`1`$, respectively. Since the index of a quadratic form does not depend on the chosen basis we can speak of a positive or negative Clifford algebra $`𝒞\mathrm{}(E,Q)`$ in case of any vector space of even dimension and any non-degenerate quadratic form.
###### Theorem 3
Let $`E=E_1E_2`$ and $`Q=Q_1Q_2`$ with dim $`E_1`$ even. Then
$$𝒞\mathrm{}(E,Q)𝒞\mathrm{}(E_1,Q_1)𝒞\mathrm{}(E_2,\pm Q_2),$$
the sign depending on $`𝒞\mathrm{}(E_1,Q_1)`$ being positive or negative, respectively.
Proof: For $`\epsilon =e_1\mathrm{}e_n𝒞\mathrm{}(E_1,Q_1)`$, $`n=dimE_1`$, one has
$$\epsilon e_i=(1)^{n1}e_i\epsilon =e_i\epsilon ,$$
i.e., $`\epsilon v=v\epsilon `$ for any $`vE_1𝒞\mathrm{}(E_1)`$. We define
$$\phi :E=E_1E_2𝒞\mathrm{}(E_1,Q_1)𝒞\mathrm{}(E_2,\pm Q_2)$$
by
$$\phi (v_1,v_2)=v_11+\epsilon v_2,v_iE_i,$$
and obtain
$`\phi (v_1,v_2)^2`$ $`=`$ $`v_1^21+\epsilon ^2v_2^2+v_1\epsilon v_2+\epsilon v_1v_2`$
$`=`$ $`v_1^21\pm 1v_2^2`$
$`=`$ $`\left(Q_1(v_1)+Q_2(v_2)\right)11`$
since $`v_2^2=\left(\pm Q_2(v_2)\right)=Q_2(v_2)`$. Using Theorem 1 (more precisely the corresponding result for an arbitrary quadratic form) we obtain a homomorphism
$$\stackrel{~}{\phi }:𝒞\mathrm{}(E,Q)𝒞\mathrm{}(E_1,Q_1)𝒞\mathrm{}(E_2,\pm Q_2).$$
Since dimensions match we are reduced to verify that $`\stackrel{~}{\phi }`$ is surjective. To this end it suffices to show that $`v_11`$ and $`1v_2`$ belong to the image of $`\stackrel{~}{\phi }`$. But now we have
$$v_11=\stackrel{~}{\phi }\left(\iota _E(v_1)\right)\text{ and }1v_2=\pm \stackrel{~}{\phi }\left(\iota _E(v_2)\epsilon \right),$$
which concludes the proof.
###### Proposition 2
If $`dimE`$ is even and $`𝒞\mathrm{}(E,Q)`$ positive, then
$$𝒞\mathrm{}(E,Q)𝒞\mathrm{}(E,Q).$$
Proof: Employing the canonical maps $`\iota _\pm :E𝒞\mathrm{}(E,\pm Q)`$ we put
$$\epsilon _\pm =\iota _\pm (e_1)\mathrm{}\iota _\pm (e_n).$$
Now $`f:Ev\epsilon _+\iota _+(v)𝒞\mathrm{}(E,Q)`$ satisfies
$$f(v)^2=\epsilon _+^2\iota _+(v)^2=\left(Q(v)\right),$$
hence induces a homomorphism $`\stackrel{~}{f}:𝒞\mathrm{}(E,Q)𝒞\mathrm{}(E,Q)`$. Moreover
$$\stackrel{~}{f}(\epsilon _{}\iota _{}(v))=\left(f(e_1)\mathrm{}f(e_n)\right)f(v)=(1)^{n(n1)/2}\epsilon _+^{n+2}\iota _+(v)=\pm \iota _+(v)$$
whereby $`\stackrel{~}{f}`$ is surjective, hence bijective.
We have already determined $`𝒞\mathrm{}_{1,0}=`$ and $`𝒞\mathrm{}_{2,0}=`$. It is not difficult to see that
$`𝒞\mathrm{}_{0,1}`$ $``$ $`=(1,1)+(1,1),`$
$`𝒞\mathrm{}_{0,2}`$ $``$ $`M_2()\text{with}e_1=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\text{ and }e_2=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$
$`𝒞\mathrm{}_{1,1}`$ $``$ $`M_2()\text{with}e_1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)\text{ and }e_2=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).`$
Combining Theorem 3 with the last Proposition we obtain the following complete classification of Clifford algebras $`𝒞\mathrm{}_{r,s}`$.
###### Theorem 4
The Clifford algebras $`𝒞\mathrm{}_{r+n,s+n}`$ and $`M_{2^n}(𝒞\mathrm{}_{r,s})`$ are isomorphic, in particular
$$𝒞\mathrm{}_{n,n}M_{2^n}(),𝒞\mathrm{}_{r,s}\{\begin{array}{cc}M_{2^s}(𝒞\mathrm{}_{rs,0}),\hfill & r>s\hfill \\ M_{2^r}(𝒞\mathrm{}_{0,sr}),\hfill & r<s\hfill \end{array}$$
and
$$𝒞\mathrm{}_{r+8,s}𝒞\mathrm{}_{r,s+8}M_{16}(𝒞\mathrm{}_{r,s}).$$
Proof: Since $`𝒞\mathrm{}_{1,1}`$ is positive we get
$$𝒞\mathrm{}_{r+1,s+1}𝒞\mathrm{}_{r,s}𝒞\mathrm{}_{1,1}𝒞\mathrm{}_{r,s}M_2()$$
and repeatedly by Theorem 3
$$𝒞\mathrm{}_{r+n,s+n}𝒞\mathrm{}_{r,s}M_2()\underset{n\mathrm{mal}}{\mathrm{}}M_2()𝒞\mathrm{}_{r,s}M_{2^n}()M_{2^n}(𝒞\mathrm{}_{r,s}).$$
Since $`𝒞\mathrm{}_{4,0}`$ is positive again by Theorem 3 and by the Proposition we have
$$𝒞\mathrm{}_{8,0}𝒞\mathrm{}_{4,0}𝒞\mathrm{}_{4,0}𝒞\mathrm{}_{4,0}𝒞\mathrm{}_{0,4}𝒞\mathrm{}_{4,4}M_{16}()$$
and
$$𝒞\mathrm{}_{p+8,q}𝒞\mathrm{}_{p,q}𝒞\mathrm{}_{8,0}𝒞\mathrm{}_{p,q}M_{16}(),$$
because $`𝒞\mathrm{}_{8,0}`$ is positive, too.
We end up with the following table displaying the special Clifford algebras $`𝒞\mathrm{}_n=𝒞\mathrm{}_{n,0}`$
$`n`$ $`\mathrm{\hspace{0.17em}0}`$ $`\mathrm{\hspace{0.17em}1}`$ $`\mathrm{\hspace{0.17em}2}`$ $`\mathrm{\hspace{0.17em}3}`$ $`\mathrm{\hspace{0.17em}4}`$ $`\mathrm{\hspace{0.17em}5}`$ $`\mathrm{\hspace{0.17em}6}`$ $`\mathrm{\hspace{0.17em}7}`$ $`\mathrm{\hspace{0.17em}8}`$ $`𝒞\mathrm{}_n`$ $``$ $``$ $``$ $``$ $`M_2()`$ $`M_4()`$ $`M_8()`$ $`M_8()M_8()`$ $`M_{16}()`$
## 2 Representations of Clifford algebras
We also need representations of abstract Clifford algebras. Recall that a representation $`\rho :𝒞\mathrm{}_nEnd(E)`$ on a real (or complex) finite dimensional vector space $`E`$ is irreducible if for any decomposition $`E=E_1E_2`$ into subspaces invariant under $`\rho `$ one has $`E_1=E`$ or $`E_2=E`$. In the reducible case one has $`\rho =\rho _1\rho _2`$ with $`\rho _j=\rho |_{E_j}`$. Give any non-trivial representation $`\rho `$, one can find an inner product $`,`$ on $`E`$ such that $`\rho (x)`$ acts orthogonally (or unitarily) on $`E`$ for all $`x^n𝒞\mathrm{}_n`$ with $`|x|=1`$. One merely has to average a given inner product $`,^{}`$ over the finite (multiplicative) group $`G_n`$ generated by $`e_1,\mathrm{},e_n𝒞\mathrm{}_n`$, i.e. one puts
$$v,w=\underset{xG_n}{}\rho (x)v,\rho (x)w^{},v,wE.$$
Since $`\rho (x)^2=|x|^2I_E`$, this amounts to
$$\rho (x)v,w=v,\rho (x)w,v,wE,x^n,$$
i.e. $`\rho (x)^{}=\rho (x)`$. If this holds we call $`\rho `$ a skew-adjoint representation.
Now any representation $`\rho `$ can easily be decomposed into a direct sum of irreducible ones: Choosing $`vE`$, $`v0`$, one considers $`E_v=\{\rho (x)vx𝒞\mathrm{}_n\}`$ which is invariant under $`\rho `$. Since $`E_v^{}`$ is also invariant, successively splitting off invariant subspaces (in case also of $`E_v`$) one ends up with $`E=_{j=1}^mE_j`$ and $`\rho =_{j=1}^m\rho _j`$ where $`\rho _j`$ is irreducible. Two representations $`\rho _j:𝒞\mathrm{}_nEnd(E_j)`$ are called equivalent if they are implemented by an isomorphism $`T:E_1E_2`$, i.e.
$$T\rho _1(x)=\rho _2(x)T,\text{for all}x𝒞\mathrm{}_n.$$
In our case we have $`𝒞\mathrm{}_n`$ of the form $`M_m(𝕂)`$ if $`n3`$ and 7 which being a simple algebra does not contain any non-trivial two-sided ideal. To see this consider elementary matrices $`e_{ij}`$ with entries 1 at $`i,j`$ and 0 elsewhere. Now given a two-sided ideal $`VM_m(𝕂)`$ and $`x=_{1i,jm}x_{ij}e_{ij}V\{0\}`$ there is an $`x_{ij}0`$ and therefore $`e_{ij}=x_{ij}^1e_{ii}xe_{jj}V`$. Since $`e_{ij}e_k\mathrm{}=\delta _{kj}e_i\mathrm{}`$ all of the $`e_{ij}`$ belong to $`V`$, i.e. $`V=M_m(𝕂)`$. In particular, we consider the left regular representation
$$\rho _L:M_m(𝕂)End\left(M_m(𝕂)\right)\rho _L(x)y=xy,x,yM_m(𝕂).$$
It decomposed as $`\rho _L=_{j=1}^m\rho _j`$ with irreducible representations $`\rho _j(x)ye_{jj}=xye_{jj}`$ on the left ideals $`V_j=M_m(𝕂)e_{jj}`$. If $`\rho `$ is an arbitrary faithful (i.e. injective) irreducible representation it has to be equivalent to one of the $`\rho _j`$ and hence to each of them. To prove this note that there is a $`vE`$ and an $`xV_1`$ with $`\rho (x)v0`$. Now define $`T:V_1E`$ by $`T(y)=\rho (y)v`$, $`yV_1`$, and observe that
$$T\rho _1(z)y=T(zy)=\rho (zy)v=\rho (z)\rho (y)v=\rho (z)Ty,yV_1,$$
hence by Schur’s Lemma $`T`$ has to be an isomorphism since both representations are irreducible: $`\mathrm{ker}TV_1`$ and $`\mathrm{im}TE`$ are subspaces invariant under $`\rho _1`$ and $`\rho `$, respectively, hence $`\mathrm{im}T=E`$ and $`kerT=\{0\}`$, since $`T0`$. Combining this with Theorem 4 and the table above we obtain:
###### Theorem 5
For $`n3`$ and 7 $`mod(8)`$ the Clifford algebra $`𝒞\mathrm{}_n`$ has up to equivalence exactly one irreducible representation, viz. on $`^{a_n}`$, where
$$a_n=\{\begin{array}{cc}1,\hfill & n=0\text{,}\hfill \\ 2,\hfill & n=1\text{,}\hfill \\ 4,\hfill & n=2,3\text{,}\hfill \\ 8,\hfill & n=4,5,6,7\text{,}\hfill \end{array}$$
and $`a_{n+8k}=2^{4k}a_n`$. In cases $`n3`$ or 7 $`mod(8)`$ there are exactly two non-equivalent irreducible representations on $`^{a_n}`$.
Proof: Noting that $``$ is irreducibly represented in $`M_2()`$ by $`a+ib\left(\begin{array}{cc}a& b\\ b& a\end{array}\right)`$ and $``$ in $`M_2()M_4()`$ by $`z+wj\left(\begin{array}{cc}z& w\\ \overline{w}& \overline{z}\end{array}\right)`$ the first assertion follows from the table. For $`n=3`$ or 7 one has $`𝒞\mathrm{}_nM_n(𝕂)M_n(𝕂)`$ with $`𝕂=`$ or $``$, respectively, and two irreducible representations on $`𝕂^n^{a_n}`$ are given by $`\rho _1(x,y)=\rho (x)`$ and $`\rho _2(x,y)=\rho (y)`$. They are not equivalent since $`\rho _1(I_n,I_n)=I_n`$ and $`\rho _2(I_n,I_n)=I_n`$.
Writing $`n=(2\mathrm{}+1)16^\alpha 2^\beta `$ with $`\beta =0,1,2`$, or $`3`$ and $`\rho (n)=8\alpha +2^\beta `$ the highest power of 2 dividing $`n`$ being just $`a_{\rho (n)1}`$ we obtain:
Corollary The Clifford algebra $`𝒞\mathrm{}_{\rho (n)1}`$ has a non-trivial representation on $`^n`$. In particular, there are matrices $`A_1,\mathrm{},A_{\rho (n)1}O(n)`$ with $`A_i^2=I_n`$ and $`A_iA_j=A_jA_i`$, $`ij`$, $`i,j=1,\mathrm{},\rho (n)1`$.
Proof: Let $`n=pa_{\rho (n)1}`$, $`p`$ odd, and $`\delta :𝒞\mathrm{}_{\rho (n)1}End(^{a_{\rho (n)1}})`$ the previous representation. Then
$$\overline{\delta }=\underset{k=1}{\overset{p}{}}\delta :𝒞\mathrm{}_{\rho (n)1}End\left(\underset{k=1}{\overset{p}{}}^{a_{\rho (n)1}}\right)$$
is the one we are looking for, since, as seen before, we can choose an inner product that renders $`A_i`$ orthogonal with respect to a suitable orthonormal basis.
The matrices $`A_j`$ and the numbers $`a_{\rho (n)1}`$ which are guaranteed by the Corollary are often called Hurwitz-Radon matrices and Radon numbers, respectively, after A. Hurwitz \[Hur\] and J. Radon \[Rad\] who around 1920 independently constructed such matrices in order to factorize quadratic forms. They also solved the linear vector field problem: There are exactly $`a_{\rho (2n)1}`$ linear vector fields, given by $`X_j(x)=A_jx`$, $`xS^{2n1}^{2n}`$, that are linearly independent at each point; cf. \[Eck\].
We also consider complex Clifford algebras $`𝒞\mathrm{}_n^{}=𝒞\mathrm{}_n_{}`$ and their irreducible representations on complex vector spaces. Complexifying immediately entails
$$𝒞\mathrm{}_n^{}\{\begin{array}{cc}M_{2^k}(),\hfill & \text{if }n=2k\text{,}\hfill \\ M_{2^k}()M_{2^k}(),\hfill & \text{if }n=2k+1\text{.}\hfill \end{array}$$
This also shows that (up to equivalence) $`𝒞\mathrm{}_n^{}`$ has exactly one irreducible representation if $`n=2k`$ and exactly two if $`n=2k+1`$. The isomorphism with $`M_{2^k}()`$ can be made explicit using the Pauli matrices $`\sigma _j`$. The basis elements $`e_j`$, $`1j2k`$, are represented (up to a choice of sign) by the following skew-hermitian unitary matrices:
$`A_{2\mathrm{}1}=`$ $`iA_{2\mathrm{}1}^{}`$ $`=\sigma _3\underset{\mathrm{}1\text{-}times}{\mathrm{}}\sigma _3i\sigma _1I_2\underset{n\mathrm{}\text{-}times}{\mathrm{}}I_2,1\mathrm{}k,`$
$`A_2\mathrm{}=`$ $`iA_2\mathrm{}^{}`$ $`=\sigma _3\underset{\mathrm{}1\text{-}times}{\mathrm{}}\sigma _3i\sigma _2I_2\underset{n\mathrm{}\text{-}times}{\mathrm{}}I_2,1\mathrm{}k.`$
This is a simple consequence of the construction in Theorem 3.
These matrices allow to classify complex Clifford algebras and their irreducible representations directly. If $`n=2k`$, i.e., $`dimM_{2^k}()=2^{2k}=dim𝒞\mathrm{}_n`$ one only has to show that the representation $`\rho (e_j)=A_jM_{2^k}()`$, $`j=1,\mathrm{},2k=n`$, is faithful, i.e. that the matrices $`A_I=A_{i_1}\mathrm{}A_i_{\mathrm{}}`$, $`1i_1<\mathrm{}<i_{\mathrm{}}2k`$, are linearly independent. To this end one uses the trace which defines an inner product on $`M_{2^k}()`$ by $`A,B=tr(A^{}B)`$. Now for $`\mathrm{}`$ even one has
$$tr(A_I)=tr(A_i_{\mathrm{}}A_{i_1}\mathrm{}A_{i_\mathrm{}1}=(1)^\mathrm{}1tr(A_I),$$
hence $`tr(A_I)=0`$, and for $`\mathrm{}<2k`$ odd and $`i_{\mathrm{}+1}I`$ one has
$$tr(A_I)=tr(A_IA_{i_{\mathrm{}+1}}A_{i_{\mathrm{}+1}})=tr(A_{i_{\mathrm{}+1}}A_IA_{i_{\mathrm{}+1}})=(1)^{\mathrm{}}tr(A_I),$$
hence again $`tr(A_I)=0`$. Given a linear combination $`a_IA_I=0`$ this implies
$$0=tr\left(a_IA_IA_J\right)=\pm a_J2^k.$$
Note that the argument does not use the special shape of the matrices $`A_j`$.
If $`n=2k+1`$ there is another matrix
$$A_{2k+1}=iA_{2k+1}^{}=i\sigma _3\mathrm{}\sigma _3.$$
However, the extended representation $`\rho :𝒞\mathrm{}_{2k+1}M_{2^k}()`$ by $`\rho (e_{2k+1})=A_{2k+1}`$ is no longer faithful, since
$$\omega =i^{[(n+1)/2]}e_1\mathrm{}e_n=i^{k+1}e_1\mathrm{}e_{2k+1}$$
is represented by $`\rho (\omega )=I_{2^k}`$.
A non-equivalent representation $`\rho ^{}`$ will be defined by $`\rho ^{}(e_j)=A_j`$, $`1j2k+1`$, since $`\rho ^{}(\omega )=I_{2^k}`$. To obtain a faithful (reducible) representation one takes the direct sum $`\rho \rho ^{}:𝒞\mathrm{}_{2k+1}^{}M_{2^k}()M_{2^k}()M_{2^{k+1}}()`$.
###### Definition 2
If $`\rho :𝒞\mathrm{}_n^{}End(E)`$ is an irreducible faithful representation, then the vector space $`E^{2^k}`$ is called a space of spinors; usually, it will be denoted by $`S_0`$.
Remarks 1. Different realizations of $`S_0`$ will be given in the following examples.
2. In $`𝒞\mathrm{}_{2k}^{}`$ the element $`\omega =i^ke_1\mathrm{}e_{2k}`$ satisfies $`\omega ^2=1`$ and $`\omega e_j=e_j\omega `$, $`j=1,\mathrm{},2k`$, hence defines a $`_2`$-grading on $`S_0`$, i.e. $`S_0=S_0^0S_0^1`$ where $`S_0^j=\frac{1}{2}\left(1+(1)^j\omega \right)S_0`$, $`j=0,1`$, are the so-called spaces of half-spinors.
The uniqueness of irreducible representations by complex $`2^k\times 2^k`$-matrices that contain and generalize Pauli’s matrices \[Pau\] has first been proved by P. Jordan and E. Wigner \[JW\] using group theoretical arguments (in connection with the quantum theory of many electron systems in 1927). The shortest proof without any theory of real Clifford algebras can be found in H. Weyl’s “Group Theory and Quantum mechanics” of 1931. He explicitly gives the matrices $`A_j^{}`$ and expresses by them all of the elementary matrices that generate the simple algebra $`M_{2^k}()`$; cf. also \[BW\] and \[Wey\]. We give his construction in the following example.
Examples 3. The reducible representation $`\rho _L:𝒞\mathrm{}_n^{}End(𝒞\mathrm{}_n^{})`$ be decomposed into a sum of irreducible ones if $`n=2k`$. In the first case one needs a minimal left ideal $`V`$ to act on. Starting from an orthonormal basis $`\{e_1,\mathrm{},e_{2k}\}`$ of $`^n`$ one can construct $`V`$ as follows: Put
$$f_{\mathrm{}}=\frac{1}{\sqrt{2}}(e_{2\mathrm{}1}+ie_2\mathrm{})\text{and}g_{\mathrm{}}=\frac{1}{\sqrt{2}}(e_{2\mathrm{}1}ie_2\mathrm{})$$
as well as
$$p_{\mathrm{}}^\pm =\frac{1}{2}(1\pm ie_{2\mathrm{}1}e_2\mathrm{})\text{for}1\mathrm{}k,$$
hence $`p_{\mathrm{}}^+=\frac{1}{2}f_{\mathrm{}}g_{\mathrm{}}`$ and $`p_{\mathrm{}}^{}=\frac{1}{2}g_{\mathrm{}}f_{\mathrm{}}`$. The idempotents $`p_{\mathrm{}}^\pm `$ mutually commute, and for any $`n`$-tuple $`\epsilon =(\epsilon _1,\mathrm{},\epsilon _k)`$ with $`\epsilon _j=\pm `$ they define a projection $`p^\epsilon =p_1^{\epsilon _1}\mathrm{}p_k^{\epsilon _k}`$. Then $`V^\epsilon =𝒞\mathrm{}_{2k}^{}p^\epsilon `$ is a minimal left ideal and $`𝒞\mathrm{}_{2k}^{}End(V^\epsilon )`$.
Note that each projection $`p^\epsilon `$ is associated with an elementary matrix $`e_{jj}`$ in $`M_{2^k}()`$, e.g. $`e_{11}`$ with $`p=p^\epsilon `$ where $`\epsilon =(1,\mathrm{},1)`$. Thus $`V`$ is isomorphic with the vector space of matrices that have non-trivial entries only in its $`j^{th}`$ column. Therefore, one has $`1=_\epsilon p^\epsilon `$.
With regards to this example B.L. van der Waerden writes in 1966 \[vdW2\]:
> If you want to determine the structure of an algebra or of a group defined by generating elements and relations and to find a representation of the algebra or group by linear transformations or by permutations, construct the regular representation.
4. To decompose the reducible representation $`c:𝒞\mathrm{}_n^{}End(^n)`$ which in fact is equivalent to the previous one one starts with the orthogonal decomposition $`^n=W\overline{W}`$, where $`W`$ or $`\overline{W}`$ denote the subspaces spanned by $`g_{\mathrm{}}`$ or $`f_{\mathrm{}}`$, respectively. From the relations
$`f_jf_{\mathrm{}}+f_{\mathrm{}}f_j=0,`$
$`g_jg_{\mathrm{}}+g_{\mathrm{}}g_j=0,`$
$`f_jg_{\mathrm{}}+g_{\mathrm{}}f_j=2\delta _j\mathrm{}.`$
and since $`f_jg_Ip=0`$ if $`jI`$ and $`f_jg_Ip=(1)^{\mathrm{}}2g_I^{}p`$ if $`I=I^{}\{j=i_{\mathrm{}}\}`$ one obtains that the subspace
$$V=𝒞\mathrm{}_{2k}^{}p=Wp$$
is a left ideal and isomorphic with $`W`$ as a vector space. Modifying $`c`$ on $`W^{2k}`$ by taking
$$\stackrel{~}{c}(w)=\sqrt{2}\left(\epsilon (v)i(\overline{v})\right)End(W)$$
for $`w=v+\overline{v}W\overline{W}`$ one obtains the appropriate irreducible representation. Indeed, from the previous relations one easily verifies for
$$w=\underset{j=1}{\overset{k}{}}(x_je_{2j1}+y_je_{2j})=\frac{1}{\sqrt{2}}\underset{j=1}{\overset{k}{}}(\overline{z_j}f_j+z_jg_j)$$
with $`z_j=x_j+iy_j`$ the relation $`\stackrel{~}{c}(w)^2=_{j=1}^k|z_j|^2I=|w|^2I`$.
The main problem with the space of spinors is that there is no canonical way to decompose a given representation, even a natural one as in the previous examples, into irreducible ones. Therefore, the spin structure to be defined in the next section and whose construction rests on a proper choice of irreducible representations will in its last analysis always be superficial.
We have shown that up to equivalence any representation $`\rho :𝒞\mathrm{}_{2k}End(E)`$ can be written as $`\rho _0I:𝒞\mathrm{}_{2k}End(S_0W)`$ with $`ES_0W`$ and where
$$\rho (v)(ew)=\rho _0(v)ew,ewS_0W.$$
Now given $`S_0`$, at least, $`W`$ is canonically defined. To see this we have to digress and recall some general results about tensor products.
Given two real (or complex) vector spaces $`E`$ and $`F`$, which moreover are right respectively left modules for some real (or complex) algebra $`A`$ the tensor product $`E_AF`$ is defined as the quotient space of $`EF`$ by the subspace generated by $`vawvaw`$, $`vE`$, $`wF`$, $`aA`$. It is the unique vector space with the following universal property. If $`H`$ is a vector space and $`f:E\times FH`$ is a bilinear $`A`$-balanced map, i.e. $`f(va,w)=f(v,aw)`$ for $`vE`$, $`wF`$, $`aA`$, then there is a unique linear map $`f_A`$ such that the following diagram commutes:
Let $`B`$ denote another algebra and let $`G`$ be a left-$`B`$-module. Then the following results hold:
(a) If $`F`$ is also a right-$`B`$-module (hence an $`(A,B)`$-bimodule), then $`E_AF`$ is a right-$`B`$-module, $`F_BG`$ an left-$`A`$-module, and
$$(E_AF)_BGE_A(F_BG)$$
is a natural isomorphism.
(b) If $`E`$ is a $`(B,A)`$-bimodule, then $`Hom_B(E,G)`$, the space of $`B`$-module homomorphisms consisting of linear maps $`fHom(E,G)`$, which satisfy $`f(bv)=bf(v)`$ for $`vE`$ and $`bB`$, is an left-$`A`$-module by $`af(v)=f(va)`$, $`aA`$, $`vE`$. One has the natural isomorphism
$$Hom_A((F,Hom_B(E,G))Hom_B(E_AF,G),$$
induced by $`f\stackrel{~}{f}`$ mit $`\stackrel{~}{f}(uv)=f(v)(u)`$ for $`fHom_A(F,Hom_B(E,G))`$, $`uE`$, and $`vF`$.
(c) Moreover, by $`(fv)(w)=f(w)v`$ for $`fHom_B(G,E)`$, $`vF`$, and $`wG`$, one obtains a natural homomorphism
$$Hom_B(G,E)_AFHom_B(G,E_AF),$$
which is one-to-one and onto if $`F`$ is a finitely generated projective module, i.e. a direct summand of $`A^n`$ for some $`n`$.
(d) If, however, $`E`$ is an right-$`A`$-module, $`G`$ a left-$`B`$-module, and $`F`$ a $`(B,A)`$-bimodule, then $`Hom_A(E,F)`$ is a left-$`B`$-module by $`(bf)(v)=bf(v)`$, $`bB`$, $`vE`$, and one has a natural isomorphism
$$E_AHom_B(F,G)Hom_B(Hom_A(E,F),G),$$
induced by $`uf\stackrel{~}{h}`$ with $`h(g)=fg(u)`$ for $`uE`$, $`fHom_B(F,G)`$, and $`gHom_A(E,F)`$.
We only need these results in the special case $`A=`$ and leave its proofs to the reader; cf. \[AF\].
As a simple consequence of the last one we obtain that the module $`W`$ in the decomposition $`E=S_0W`$ can be chosen as $`W=Hom_{𝒞\mathrm{}_n^{}}(S_0,E)`$: Since $`𝒞\mathrm{}_n^{}=End(S_0)`$, there are isomorphisms
$$SHom_{𝒞\mathrm{}_n^{}}(S_0,E)Hom_{𝒞\mathrm{}_n^{}}(Hom(S_0,S_0),E)Hom_{𝒞\mathrm{}_n^{}}(𝒞\mathrm{}_n^{},E)E.$$
We conclude this section and the algebraic part of the paper with a classical result of representation that will be essential in the proof of the main theorems of the next section. It is a special case of the Theorem of Skolem-Noether.
LemmaLet $`A_1`$ and $`A_2`$ be two isomorphic simple subalgebras of $`M_n()`$, say both isomorphic to $`M_k()`$. Then each isomorphism $`\mathrm{\Phi }:A_1A_2`$ is an inner automorphism $`\mathrm{Ad}(U)`$ of $`M_n()`$, i.e., there is a $`UGL_n()`$ with
$$\mathrm{\Phi }(a)=\mathrm{Ad}(U)a=UaU^1,aA_2.$$
In particular, each automorphism of $`M_k()`$ is inner and each derivation $`D`$ of $`M_k()`$ is an inner derivation, i.e. given by
$$Da=\mathrm{ad}(v)a=[v,a]=vaav,aM_k(),$$
for some $`vM_k()`$.
Proof: The simple algebra $`M_k()`$ is represented by $`A_1`$ and $`A_2`$ in $`M_n()`$, respectively. There are decompositions $`^n=_{j=1}^{\mathrm{}}E_j`$ and $`^n=_{j=1}^rF_j`$ which reduce $`A_1`$ and $`A_2`$, respectively. $`A_1`$ and $`A_2`$ being isomorphic, one has $`\mathrm{}=r`$, and since the restricted irreducible representations have to equivalent, one has
$$E_j^kF_j\text{for all}j.$$
Now $`U`$ is given as the direct sum of such isomorphisms. To prove the second assertion, one simply has to take $`k=n`$ and to choose $`A_2`$ as the image of $`A_1=M_k()`$ under a given automorphism. For the last assertion take $`n=2k`$, $`A_1=\{\left(\begin{array}{cc}a& 0\\ 0& a\end{array}\right)aM_k()\},`$ and $`A_2=\{\left(\begin{array}{cc}a& D(a)\\ 0& a\end{array}\right)aM_k()\}`$ for a given derivation $`D`$. Then there is a $`U=\left(\begin{array}{cc}u& v\\ w& z\end{array}\right)`$ with
$$\left(\begin{array}{cc}u& v\\ w& z\end{array}\right)\left(\begin{array}{cc}x& 0\\ 0& x\end{array}\right)=\left(\begin{array}{cc}x& D(x)\\ 0& x\end{array}\right)\left(\begin{array}{cc}u& v\\ w& z\end{array}\right)$$
for all $`xM_k()`$. This entails $`wx=xw`$ and $`zx=xz`$ for all $`xM_k()`$, hence, by Schur’s Lemma, $`w`$ and $`z`$ are multiples of the identity. If say $`z0`$, the further condition $`xv+D(x)z=vx`$ on $`D`$ leads to $`D=\mathrm{ad}(z^1v)`$.
Remark The maps $`\nu :GL_k()Aut(M_k())`$, $`\tau (u)=\mathrm{Ad}(u)`$ and $`\mu :M_k()\mathrm{Der}(M_k())`$, $`\mu (v)=\mathrm{ad}(v)`$ (into the space of derivations) are both onto but in general not one-to-one, since $`\mathrm{ker}\nu ^{}=\{0\}`$ and $`\mathrm{ker}\mu `$. In the case of $`\mu `$ one can, however, consider its restriction $`\mu _0`$ to the subspace $`M_k^0()`$ of matrices with vanishing trace and obtains an isomorphism.
## 3 Spinor bundles and Dirac operators
We now want to globalize the results of the previous section, i.e. to perform the constructions on vector bundles over smooth manifold.
###### Definition 3
Let $`E`$ be a Euclidean vector bundle of rank $`k`$ over $`M`$. The vector bundle $`𝒞\mathrm{}(E)=_{pM}𝒞\mathrm{}(E_p)`$ will be called the Clifford bundle of $`E`$. If $`M`$ is endowed with a Riemannian structure one particularly has $`𝒞\mathrm{}M=𝒞\mathrm{}(TM)`$, the Clifford bundle of $`M`$.
Starting from a local orthonormal frame $`(\text{e}_i)_{1ik}`$ of $`E`$ over $`U`$ one obtains a local trivialization
$$E|_Uv_q=\underset{i=1}{\overset{k}{}}a_i\text{e}_i(q)\phi (v_q)=(q,(a_i)_{1ik})U\times ^k$$
and $`\phi |_{E_q}:E_q\{q\}\times ^k=^k`$ is isometric for any $`qU`$. Choosing an atlas $`A=\{(U_\alpha ,\phi _\alpha )\alpha A\}`$ in this way one obtains a cocycle of transition maps with $`g_{\alpha \beta }:U_\alpha U_\beta O(k)`$. To trivialize $`𝒞\mathrm{}(E|_U)`$ we choose
$$𝒞\mathrm{}(\phi _\alpha )(x_q)=𝒞\mathrm{}(\phi _\alpha |_{E_q})(x_q)\{q\}\times 𝒞\mathrm{}_k,x_q𝒞\mathrm{}(E_q)$$
according to Remark 1 after Theorem 1. The corresponding transition maps are given by $`f_{\alpha \beta }:U_\alpha U_\beta Aut(𝒞\mathrm{}_k)`$ with
$$f_{\alpha \beta }(p)=𝒞\mathrm{}\left(g_{\alpha \beta }(p)\right),pU_\alpha U_\beta .$$
They possess the cocycle property and are differentiable, since the group homomorphism $`O(k)f𝒞\mathrm{}(f)Aut(𝒞\mathrm{}_k)`$ is a polynomial in the coefficients with respect to a fixed orthonormal basis of $`^k`$ and the induced basis of $`𝒞\mathrm{}_k`$. This makes $`𝒞\mathrm{}(E)`$ a smooth vector bundle. In particular,
$$\{\text{e}_{i_1}\mathrm{}\text{e}_{i_k}1i_1<\mathrm{}<i_kr,k=0,\mathrm{},n\}$$
provides an orthonormal frame of $`𝒞\mathrm{}(E)|_U`$. The Clifford bundle $`𝒞\mathrm{}(E)`$ depends on the Euclidean structure and is itself a Euclidean vector bundle. On the other hand, the $`C^{\mathrm{}}`$-structure and the Riemannian structure of $`𝒞\mathrm{}(E)`$ do not depend on the choice of the frame.
Each fiber of $`𝒞\mathrm{}(E)`$ comes with an algebra structure and fiber-wise multiplication makes $`C^{\mathrm{}}\left(𝒞\mathrm{}(E)\right)`$, the space of smooth sections, into an algebra, too. Suitably modifying the definition of a vector bundle one obtains the notion of an algebra bundle $`(A,\pi ,M)`$:
Each fiber $`\pi ^1(p)`$ is a finite dimensional topological algebra with respect to the topology induced by $`A`$, and at each point $`pM`$ there exists a chart $`\phi :\pi ^1(U)U\times A`$ with a fixed given algebra $`A_0`$, such that
$$\phi |_{\pi ^1(q)}:\pi ^1(q)\{q\}\times A_0$$
is an algebra isomorphism for any $`qU`$. We only consider the special case of unital algebras $`A_0`$ and $`A_p`$ with unit elements $`e_0`$ and $`e_p`$, respectively. Then we have a global section e in $`A`$.
Alternatively, we may assume a bundle morphism $`(\mu ,id_M)`$ with $`\mu :AAA`$ and $`\mu (e_pa_p)=a_p=\mu (a_pe_p)`$ for any $`a_pA_p`$.
Moreover a vector bundle $`F`$ will be called a (left-)$`A`$-bundle if there is a bundle morphism $`\tau :AFF`$ with
$$\tau \left(a_p\tau (b_pv_p)\right)=\tau \left(\mu (a_pb_p)v_p\right)\text{for}a_p,b_pA_p,v_pF_p.$$
In other words, $`F_p`$ is a left-$`A_p`$-module for any $`pM`$ and
$$\sigma \text{s}(p)=\sigma (p)\text{s}(p)=\tau \left(\sigma (p)\text{s}(p)\right),pM$$
defines a smooth section, $`\sigma \text{s}C^{\mathrm{}}(F)`$, for $`\sigma C^{\mathrm{}}(A)`$ and $`\text{s}C^{\mathrm{}}(F)`$. An $`A`$-bundle morphism is a bundle morphism $`(\widehat{f},id_M)`$ between two $`A`$-modules $`F`$ and $`G`$, that is an $`A_p`$-linear map from $`F_p`$ to $`G_p`$ for any $`pM`$. The space of $`A`$-bundle morphisms, $`HOM_A(F,G)`$, is a $`C^{\mathrm{}}(A)`$-module and can be identified with $`C^{\mathrm{}}\left(Hom_A(F,G)\right)`$. Here $`Hom_A(F,G)`$ is the sub-bundle of $`Hom(F,G)`$, whose fibers are $`Hom_{A_p}(F_p,G_p)`$, $`pM`$. In particular, we can extend the natural isomorphisms at the end of the previous section to the setting of $`A`$-bundles.
We are now going to define geometric differential operators that are closely connected with the topological or geometrical structure of an oriented Riemannian manifold $`M`$.
###### Definition 4
A smooth vector bundle $`E`$ over $`M`$ is called a spinor bundle over $`M`$ if it is a left-$`𝒞\mathrm{}M`$-bundle.
If the module structure is given by the morphism $`\tau :𝒞\mathrm{}MEE`$ we also consider the bundle morphism $`c_E:TMEnd(E)`$, $`c_E(v_p)(e_p)=\tau (ve)`$, $`vT_pM`$, $`e_pE_p`$, induced by $`\tau `$ and its extension $`c_E:𝒞\mathrm{}MEnd(E)`$ to a morphism of algebra bundles. To emphasize the underlying Clifford multiplication we sometimes denote a spinor bundle by $`(E,c_E)`$.
Examples 5. The Clifford bundle $`𝒞\mathrm{}M`$ itself is a spinor bundle if
$$c_{𝒞\mathrm{}M}:𝒞\mathrm{}MEnd(𝒞\mathrm{}M)$$
is in each fiber given by the left regular representation $`\rho _L`$.
6. Likewise the Grassmann bundle $`^{}M`$, the exterior bundle of the cotangent bundle $`T^{}M`$, is turned into a spinor bundle using the isomorphism of $`𝒞\mathrm{}M`$ with $`^{}M`$. Here and in the following we use the “musical isomorphisms” $`{}_{}{}^{\mathrm{}}:^{}MM`$ and its inverse $`{}_{}{}^{\mathrm{}}:M^{}M`$ that extend the pairing between tangent vectors and cotangent vectors provided by the Riemannian metric g of $`M`$.
7. Given a spinor bundle $`E`$ and a smooth vector bundle $`F`$ we can turn $`EF`$ into a spinor bundle, $`E`$ twisted by $`F`$. Here $`𝒞\mathrm{}M`$ operates on $`EF`$ by $`v(ef)=(ve)f`$ for $`efEF`$.
Given a spinor bundle $`E`$ over $`M`$ via the isomorphism $`{}_{}{}^{\mathrm{}}:T^{}MTM𝒞\mathrm{}M`$ the bundle morphism $`\tau `$ induces a linear map $`T:C^{\mathrm{}}(T^{}ME)C^{\mathrm{}}(E)`$ given by
$$T(\omega \text{s})(p)=\tau \left(\omega (p)^{\mathrm{}}\text{s}(p)\right),pM.$$
It is easy to see that $`T`$ is a differential operator of order zero. To define more sophisticated differential operators on $`C^{\mathrm{}}(E)`$ we need a (Koszul) connection $``$ on $`E`$, i.e. a linear differential operator $`:C^{\mathrm{}}(E)C^{\mathrm{}}(T^{}ME)`$ satisfying
$$(f\text{s})=df\text{s}+f\text{s},fC^{\mathrm{}}(M),\text{s}C^{\mathrm{}}(E).()$$
For a vector field $`XC^{\mathrm{}}(TM)`$ this gives rise to a covariant derivative $`_X`$ that satisfies
$$_X(f\text{s})=X(f)\text{s}+f_X(\text{s}).$$
The dual connection $`^{}`$ on $`C^{\mathrm{}}(E^{})`$ can be defined by its covariant derivatives
$$_X^{}\text{s}^{}(\text{s})=X\left(\text{s}^{}(\text{s})\right)\text{s}^{}(_X\text{s})$$
for $`\text{s}^{}C^{\mathrm{}}(E^{})`$, $`\text{s}C^{\mathrm{}}(E)`$, and $`XC^{\mathrm{}}(TM)`$. We also note the following elementary constructions that can be performed with connections $`^E`$ and $`^F`$ for vector bundles $`E`$ and $`F`$, respectively. By
$$^{EF}(\text{s}\text{t})=^E\text{s}^F\text{t},$$
and by
$$^{EF}(\text{s}\text{t})=(^E\text{s})\text{t}+\mathrm{\Psi }\left(\text{s}(^F\text{t})\right),\text{s}C^{\mathrm{}}(E),\text{t}C^{\mathrm{}}(F)$$
one defines connections $`^{EF}`$ for $`EF`$ and $`^{EF}`$ for $`EF`$. Here $`\mathrm{\Psi }`$ is induced by the isomorphism of vector bundles, $`\psi :ET^{}MFT^{}MEF`$. In particular, one obtains a connection $`^{End(E)}`$ on $`End(E)E^{}E`$.
###### Definition 5
Let $`E`$ be a spinor bundle over $`M`$, and $``$ a connection for $`E`$. Then
$$D=T:C^{\mathrm{}}(E)C^{\mathrm{}}(E)$$
defines a first order differential operator, the Dirac operator associated with $`(E,)`$.
###### Proposition 3
Given a local orthonormal frame $`(E_i)_{1im}`$ of $`TM`$ over $`U`$ one has
$$D\text{s}=\underset{k=1}{\overset{m}{}}E_k_{E_k}\text{s}$$
for $`\text{s}C^{\mathrm{}}(E|_U)`$.
Proof: Since $`X=_{k=1}^mE_x,XE_k`$ for $`XC^{\mathrm{}}(TM)`$ one has
$$_X\text{s}=\underset{k=1}{\overset{m}{}}E_k,X_{E_k}\text{s},$$
hence
$$\text{s}=\underset{k=1}{\overset{m}{}}E_k^{\mathrm{}}_{E_k}\text{s}$$
and so the representation of $`D`$ as stated.
With respect to a local frame $`(s_j)_{1jr}`$ of $`E`$ a connection is given by
$$\left(\underset{j=1}{\overset{r}{}}f_j\text{s}_j\right)=\underset{j=1}{\overset{r}{}}\left(df_j\text{s}_j+f_j\underset{k=1}{\overset{r}{}}\omega _{jk}\text{s}_k\right),$$
where the local connection form $`\omega =(\omega _{jk})_{1j,kr}`$ defined on say $`U`$ uniquely determines $``$ on $`U`$ and vice versa.
Recall that the tangent bundle of a Riemannian manifold $`M`$ itself comes with a unique torsion-free Riemannian connection, the Levi-Civita connection which we denote by $`\overline{}`$. Here torsion-free means that
$$\overline{}_XY\overline{}_YX=[X,Y]$$
and Riemannian that
$$\overline{}_XY,Z+Y,\overline{}_XZ=XY,Z$$
for any vector fields $`X`$, $`Y`$ and $`Z`$. Moreover, the Levi-Civita connection $`\overline{}`$ extends to $`T^{}M`$ and to the tensor bundle by the previously mention constructions and also to the exterior bundle $`^{}M`$ and to the Clifford bundle if we assume the product formula
$$\overline{}_X(\omega _1\omega _2)=(\overline{}_X\omega _1)\omega _2+\omega _1(\overline{}_X\omega _2)$$
for forms $`\omega _1`$, $`\omega _2\mathrm{\Omega }(\mathrm{M})`$ respectively
$$\overline{}_X(\sigma _1\sigma _2)=(\overline{}_X\sigma _1)\sigma _2+\sigma _1(\overline{}_X\sigma _2)$$
for sections $`\sigma _1`$, $`\sigma _2^{\mathrm{}}(𝒞\mathrm{}M)`$. Combined with the action of the Clifford bundle we obtain Dirac operators that are defined on any oriented Riemannian manifold. The Dirac operator on $`\mathrm{\Omega }(\mathrm{M})`$ has been introduced by E. Kähler in 1961 \[Kae\] and so is sometimes called Dirac-Kähler. The extension $`\overline{}`$ to $`^{}M`$ also satisfies
$$\overline{}(\sigma \omega )=(\overline{}\sigma )\omega +\sigma (\overline{}\omega )$$
in the sense that
$$\overline{}_X(\sigma \omega )=(\overline{}_X\sigma )\omega +\sigma (\overline{}_X\omega )$$
for $`XC^{\mathrm{}}(TM)`$, $`\omega \mathrm{\Omega }(M)`$, and $`\sigma C^{\mathrm{}}(𝒞\mathrm{}M)`$, hence in both cases $`\overline{}`$ and Clifford multiplication are compatible. Also recall that Clifford multiplication by unit tangent vectors $`X_pT_pM`$ is orthogonal on the spinor bundles $`𝒞\mathrm{}M`$ and $`^{}M`$ equipped with the Riemannian metric induced by g. This suggests the following definition.
###### Definition 6
Let $`E`$ be complex vector bundle with a Hermitian metric $`,`$, a connection $``$ and a left $`𝒞\mathrm{}M_{}`$-module structure $`c_E`$. We call the triple $`(E,,,)`$ a Dirac triple and, for short, $`E`$ a Dirac bundle if the given data are compatible, i.e. if
(1) $`c_E`$ is a skew-adjoint representation in each fiber,
(2) $``$ is a compatible connection, i.e.
$$(\sigma \text{s})=(\overline{}\sigma )\text{s}+\sigma \text{s}\sigma C^{\mathrm{}}(𝒞\mathrm{}M),\text{s}C^{\mathrm{}}(E),$$
(3) $``$ is a Riemannian connection, i.e.
$$_X\text{s}_1,\text{s}_2+\text{s}_1,_X\text{s}_2=X\left(\text{s}_1,\text{s}_2\right),XC^{\mathrm{}}(TM),\text{s}_1,\text{s}_2C^{\mathrm{}}(E).$$
Remarks 1. By definition of the Levi-Civita connection on $`𝒞\mathrm{}M`$ to ensure (2) it suffices that
$$(X\text{s})=(\overline{}X)\text{s}+X\text{s}$$
for $`XC^{\mathrm{}}(TM)C^{\mathrm{}}(𝒞\mathrm{}M)`$ and $`\text{s}C^{\mathrm{}}(E)`$.
2. If $`(E,^E)`$ is a Dirac bundle and $`F`$ is a Riemannian vector bundle with Riemannian connection $`^F`$, then $`(EF,^E^F)`$ with Clifford multiplication as in Example 3 is again a Dirac bundle, since for $`\text{s}_1C^{\mathrm{}}(E)`$, $`\text{s}_2C^{\mathrm{}}(F)`$, and $`\sigma C^{\mathrm{}}(𝒞\mathrm{}M)`$ one has
$`^E^F\left(\sigma (\text{s}_1\text{s}_2)\right)`$ $`=`$ $`^E(\sigma \text{s}_1)\text{s}_2+(\sigma \text{s}_1)^F\text{s}_2`$
$`=`$ $`\left((\sigma )\text{s}_1\right)\text{s}_2+(\sigma ^E\text{s}_1)\text{s}_2+(\sigma \text{s}_1)^F\text{s}_2`$
$`=`$ $`(\sigma )(\text{s}_1\text{s}_2)+\sigma (^E\text{s}_1\text{s}_2)+\sigma (\text{s}_1^F\text{s}_2)`$
$`=`$ $`(\sigma )(\text{s}_1\text{s}_2)+\sigma ^E^F(\text{s}_1\text{s}_2).`$
Condition (1) also holds, since for $`XC^{\mathrm{}}(TM)`$
$`X(\text{s}_1\text{t}_1),\text{s}_2\text{t}_2`$ $`=`$ $`(X\text{s}_1)\text{t}_1,\text{s}_2\text{t}_2`$
$`=`$ $`X\text{s}_1,\text{s}_2\text{t}_1,\text{t}_2=\text{s}_1,X\text{s}_2\text{t}_1,\text{t}_2`$
$`=`$ $`\text{s}_1\text{t}_1,(X\text{s}_2)\text{t}_2`$
$`=`$ $`\text{s}_1\text{t}_1,X(\text{s}_2\text{t}_2).`$
In this way we obtain a Dirac operator with coefficients in the bundle $`F`$ or a Dirac operator by twisting the Dirac operator $`D^E`$ on $`E`$ with the connection $`^F`$. It will be denoted by $`D^E^F`$ or simply by $`D^EI_F`$.
It is well known that any complex vector bundle can be equipped with a Hermitian structure and with a Riemannian connection. Recall that one defines inner products and connections locally and in a second step uses partitions of unity to paste the local data to obtain global ones. So, in general, there is a lot of freedom to do this. In case of a complex spinor bundle one can ask whether these data can be chosen to satisfy (1) to (3). We shall prove that this can indeed be achieved. But before doing so we address the question of uniqueness, i.e. the impact that irreducibility has on the choice of these data.
###### Proposition 4
Let $`S`$ be an irreducible complex spinor bundle with a Hermitian metric $`,`$ and a connection $``$ satisfying properties (1) to (3). Then the following results hold:
(a) Any Hermitian metric $`,^{}`$ with property (1) is of the form
$$,^{}=\lambda ,$$
for some positive real-valued function $`\lambda C^{\mathrm{}}(M)`$.
(b) Any connection $`^{}`$ with property (2) is of the form
$$^{}=+\omega $$
for some complex-valued one-form $`\omega \mathrm{\Omega }^1(M,)`$.
(c) If moreover $`^{}`$ is a Riemannian connection with respect to the given metric, the one-form $`\omega `$ is purely imaginary, i.e. $`^{}=+i\eta `$ for some real-valued one-form $`\eta \mathrm{\Omega }^1(M,)`$.
Proof: (a) For $`pM`$ let $`TEnd(S_p)`$ be a hermitian endomorphism, such that
$$s_1,s_2^{}=Ts_1,s_2$$
for all $`s_1`$, $`s_2S_p`$. Then
$$TX_ps_1,s_2=X_ps_1,s_2^{}=s_1,Xs_2^{}=Ts_1,X_ps_2=X_pTs_1,s_2$$
for all $`X_pT_pM`$. Since the $`X_p`$ generate $`End(S_p)`$, $`T`$ commutes with each element of $`End(S_p)`$, hence by Schur’s Lemma $`T=\lambda I`$ with $`\lambda `$. Since $`T`$ is hermitian and positive, we have $`\lambda `$.
(b) Analogously we conclude that the section $`\varphi =_X^{}_X`$ into $`End(S)`$, which satisfies
$$\varphi (\sigma \text{s})=\sigma \varphi (\text{s})$$
for all $`\sigma C^{\mathrm{}}(𝒞\mathrm{}M_{})`$ and $`\text{s}C^{\mathrm{}}(S)`$ because of the derivation property that $`\varphi =\omega (X)I`$ with $`\omega (X)`$.
(c) This is immediate, since $`\omega =^{}`$ has to be skew-hermitian, i.e. $`\overline{\omega }=\omega `$.
###### Theorem 6
Let $`E`$ be a complex spinor bundle over the Riemannian manifold $`M`$ (of dimension $`m=2n`$). Then there are a Hermitian structure and a Riemannian connection for $`E`$ compatible with Clifford multiplication which possess properties (1) and (2).
Proof: It suffices to prove this locally. Using a partition of unity local metrics as well as local connections can be pasted to global ones ensuing properties (1) to (3). Let $`(U,\phi )`$ be a chart of $`M`$ at $`pM`$ trivializing $`E|_U`$. We shall show that on a possibly smaller $`U`$ there are complex vector bundles $`S`$ and $`W`$ with $`E|_U=SW`$ and the $`𝒞\mathrm{}M`$-action irreducible on $`S`$ and trivial on $`W`$. By the previous remarks it suffices to consider only $`S`$ and to equip $`W`$ with an arbitrary Hermitian structure and an arbitrary Riemannian connection. Starting from a local orthonormal frame $`\{E_1,\mathrm{},E_m\}`$ of $`TM|_U`$ we obtain sections $`\text{p}^\epsilon C^{\mathrm{}}(𝒞\mathrm{}M_{}|_U)`$ consisting of orthogonal projections. If $`\text{s}_1C^{\mathrm{}}(E|_U)`$ is a non-vanishing section one has $`\text{p}^\epsilon (q)\text{s}_1(q)0`$ in the possibly smaller open set $`U`$ for some $`\epsilon `$. Then
$$f(q,(a_\sigma )_{\sigma G_m})=\underset{\sigma G_m}{}a_\sigma \sigma (q)\text{p}^\epsilon (q)\text{s}_1(q),qU,(a_\sigma )_{\sigma G_m}^{|G|},$$
defines a vector bundle morphism $`f:U\times ^{|G_m|}E|_U`$ of constant rank $`\mathrm{rk}f(p,)=N`$ hence $`F_1=Imf`$ is a subbundle of $`E`$ whose fibers are irreducible $`𝒞\mathrm{}_m^{}`$-modules. We have $`E=F_1F_1^{}`$ and proceeding likewise with a second non-vanishing section $`\text{s}_2C^{\mathrm{}}(F_1^{}|_U)`$ etc. we eventually obtain that $`E|_USW`$ as a $`𝒞\mathrm{}U_{}`$-bundle where $`S=F_1`$ and $`W=\epsilon _U^{\mathrm{}}`$. Now the products of sections $`E_j`$ in $`C^{\mathrm{}}(𝒞\mathrm{}M|_U)`$ generate a finite group $`G_m`$. Given an arbitrary Hermitian structure $`,,^{}`$ on $`E|_U`$ we may define a new one by putting
$$v,w=\underset{\sigma G_m}{}\sigma (q)v,\sigma (q)w^{},v,wE_q.$$
Now given the irreducible spinor bundle $`S`$ we have an isomorphism of algebra bundles $`\mathrm{\Phi }:𝒞\mathrm{}M_{}|_UEnd(S)`$. Extending the Levi-Civita connection $`\overline{}`$ to $`𝒞\mathrm{}M|_U`$ and then to $`𝒞\mathrm{}M_{}|_U`$, by $`\mathrm{\Phi }^1`$ we induce a connection $`=\mathrm{\Phi }\overline{}\mathrm{\Phi }^1`$ on $`End(S)`$. We only have to show that $`=^{End(S)}`$, i.e. induced by a Riemannian connection $`^S`$ on $`S`$. This one will automatically possess property (2), since
$`^S(\sigma \text{s})`$ $`=`$ $`^S\left(\mathrm{\Phi }(\sigma )(\text{s})\right)=\left(\mathrm{\Phi }(\sigma )\right)(\text{s})+\mathrm{\Phi }(\sigma )(^S\text{s})`$
$`=`$ $`\mathrm{\Phi }(\overline{}\sigma )(\text{s})+\mathrm{\Phi }(\sigma )(^S\text{s})`$
$`=`$ $`\overline{}\sigma \text{s}+\sigma ^S\text{s}.`$
Note that from the Remark concluding section 2 we have sub-bundles $`End_0(S)`$ of fiber-wise endomorphisms with trace 0 and $`\mathrm{Der}(S)`$ of fiber-wise derivations of $`End(S)`$, as well as a bundle isomorphism $`\mu _0:End_0(S)\mathrm{Der}(S)`$.
If $`_0`$ is an arbitrary Riemannian connection on $`S`$ and $`\stackrel{~}{}_0`$ the connection induced on $`End(S)`$, then $`\eta =\stackrel{~}{}_0`$ is a section in $`T^{}MEnd\left(End(S)\right)`$ and from the derivation property even a section in $`T^{}M\mathrm{Der}(S)`$. For $`\gamma =\mu _0^1\eta `$ we then have
$$_X\text{t}\stackrel{~}{}_{0X}\text{t}=\gamma (X)\text{t}\text{t}\gamma (X),XC^{\mathrm{}}(TM),\text{t}C^{\mathrm{}}\left(End(S)\right).$$
And putting
$$_{0X}^S\text{s}=_{0X}\text{s}+\gamma (X)(\text{s}),\text{s}C^{\mathrm{}}(S),$$
we obtain, for the induced connection on $`End(S)`$,
$$\stackrel{~}{}_{0X}^S\text{t}=\stackrel{~}{}_{0X}\text{t}+\gamma (X)\text{t}\text{t}\gamma (X),$$
hence $`\stackrel{~}{}_0^S=`$. Although we started from a Riemannian connection $`_0`$ the construction does not guarantee that $`_0^S`$ is also a Riemannian connection. Now putting
$$\text{s}_1,\text{s}_2^{}=X\text{s}_1,\text{s}_2_{0X}^S\text{s}_1,\text{s}_2\text{s}_1,_{0X}^S\text{s}_2$$
we get a sesquilinear form on $`S`$, hence $`,^{}=\omega (X),`$ by (a) of the Proposition. It is easily seen, that $`\omega `$ is a (real-valued) one-form. We can finally put
$$^S=_0^S+\frac{1}{2}\omega $$
which by (b) of the Proposition satisfies property (2) and by a simple computation is seen to be a Riemannian connection with respect to $`,`$.
For $`E`$ we have found, at least locally, a decomposition $`E=SW`$ with $`𝒞\mathrm{}M`$ acting irreducibly on $`S`$. However, there are topological obstructions for a global such decomposition to hold. We come back to this point later on. However, if $`S`$ is given globally, $`W`$ is naturally determined by $`W=Hom_{𝒞\mathrm{}M_{}}(S,E)`$. It is easy to show that this is indeed a sub-bundle of the bundle of $`Hom(S,E)`$. This gives rise to the following definition.
###### Definition 7
An oriented Riemannian manifold $`M`$ of dimension $`m=2n`$ is said to be spin<sup>c</sup> if there is a complex spinor bundle $`S`$ over $`M`$ with $`𝒞\mathrm{}MEnd(S)`$.
If $`M`$ is spin<sup>c</sup> any spinor bundle $`E`$ can be written as $`E=SW`$ with some complex vector bundle $`W`$. In particular, for any further irreducible spinor bundle $`S^{}`$ there exists a complex line bundle $`L`$ with $`S^{}=SL`$, viz. $`L=Hom_{𝒞\mathrm{}M_{}}(S,S^{})`$. Given $`S`$ we can now make it a Dirac bundle by properly choosing a Hermitian structure and a Riemannian connection $``$. However, this connection is is only determined up to an additional purely imaginary one-form. Most desirable would be a unique connection on $`S`$ induced by the Levi-Civita connection of $`M`$. Then the connection on any further Dirac bundle $`S^{}=SL`$ could be chosen as the product connection only depending on the connection on the line bundle $`L`$. To ensure this we need a spin structure for $`M`$ given by an additional structure on $`S`$.
We start with the algebraic setting and consider the complex vector space $`S_0`$ of spinors bearing an operation of the real Clifford algebra $`𝒞\mathrm{}_m`$. This is not irreducible but depending on the dimension $`m=2n=8k+2\mathrm{}`$ one can find an irreducible real subspace of $`S_0`$. More precisely, there exist an antilinear map $`\theta _0:S_0S_0`$ with $`\theta _0^2=I_{S_0}`$ for $`\mathrm{}=0`$ or $`3`$ and $`\theta _0^2=I_{S_0}`$ for $`\mathrm{}=1`$ or $`2`$, a so-called structural map. In the first case $`S_0`$ carries a real structure, in the second case a quaternionic structure. This is obvious if $`\mathrm{}=0`$ or $`3`$, since then $`𝒞\mathrm{}_m=M_{2^n}()`$ is acting irreducibly on $`^{2^n}`$ and the complex Clifford algebra and the spinor space $`S_0=^{2^n}`$ are obtained therefrom by complexification. Here $`\theta `$ can be chosen the complex conjugation $`c:^{2^n}^{2^n}`$ taken component-wise. In the other two cases we consider the explicit representation of $`S_0=^{2^n}`$, and by periodicity may restrict to $`k=0`$. With $`c`$ as before and $`\tau =i\sigma _2`$ we now put $`\theta _0=\tau c`$ if $`\mathrm{}=1`$ and $`\theta _0=(\tau \sigma _3)c`$, if $`\mathrm{}=2`$. Then $`\theta _0`$ is the required structural map, and in all cases it commutes with the representation of $`𝒞\mathrm{}_m`$. The antilinear map $`\theta _0:S_0S_0`$ can also be seen as a linear map $`\theta _0:S_0\overline{S}_0`$, where by $`\overline{S}_0`$ we denote the complex vector space $`S_0`$ with scalar multiplication changed to $`\lambda v=\overline{\lambda }v`$, $`\lambda `$, $`vS_0`$. The representation of $`𝒞\mathrm{}_m`$ on $`S_0`$ induces a representation of $`𝒞\mathrm{}_m`$ on $`\overline{S}_0`$, and extending both representations to $`𝒞\mathrm{}_m^{}`$ we obtain an element $`\theta _0`$ of $`Hom_{𝒞\mathrm{}_m^{}}(S_0,\overline{S}_0)`$ with $`\theta _0^2=\pm I_{S_0}`$.
Since $`c`$ and $`\tau `$ both depend on a basis of $`S_0`$ it is in general not possible to extend this local construction to a global one on the spinor bundle $`S`$. So at first we will assume a global structural map and afterwards will establish sufficient conditions for it existence.
###### Definition 8
Let $`M`$ be an oriented Riemannian manifold of dimension $`m=8n+2\mathrm{}`$. We say that $`M`$ carries a spin structure or that $`M`$ is spin, if $`M`$ is spin<sup>c</sup> and if the irreducible complex spinor bundle $`S`$ allows a structural map $`\theta C^{\mathrm{}}\left(Hom_{𝒞\mathrm{}M_{}}(S,\overline{S})\right)`$ with $`\theta ^2=I_S`$ or $`\theta ^2=I_S`$ inducing respectively a real ($`\mathrm{}=0`$ or $`3`$) or quaternionic ($`\mathrm{}=1`$ oder $`2`$) structure on $`S`$, that is compatible with the complex conjugation of $`𝒞\mathrm{}M_{}=𝒞\mathrm{}M`$.
Remark Equivalently, we may require the existence of a real spinor bundle on which the real Clifford bundle $`𝒞\mathrm{}M`$ acts irreducibly on each fiber. If $`\mathrm{}=3`$ or $`4`$ one can choose the fixed-point bundle of $`\theta `$, and conversely the complexified real spinor bundle will define a spin structure.
Of course, any spin manifold is spin<sup>c</sup> but the converse does not hold in general. We address this question in the next section. Here we only prove the following general characterization.
###### Theorem 7
Let $`M`$ be an oriented Riemannian manifold $`M`$ with spin<sup>c</sup> structure given by the irreducible complex spinor bundle $`S`$. Then $`S`$ defines a spin structure, i.e., allows a global structural map $`\theta `$ if and only if the vector bundle $`Hom_{𝒞\mathrm{}M_{}}(S,\overline{S})`$ is trivial.
Proof: We already know that $`Hom_{𝒞\mathrm{}M_{}}(S,\overline{S})`$ is a complex line bundle: Each fiber contains a $`𝒞\mathrm{}(T_pM_{})`$-linear isomorphism $`\theta _p:S_p\overline{S}_p`$ and by irreducibility of $`S_p`$ and $`\overline{S}_p`$ and Schur’s Lemma any $`𝒞\mathrm{}(T_pM_{})`$-linear map $`\theta _p^{}:\overline{S}_pS_p`$ satisfies $`\theta _p^{}\theta _p=\lambda I_{S_p}`$ for some $`\lambda `$. In case of a spin structure $`\theta `$ defines a non-vanishing section in $`Hom_{𝒞\mathrm{}M_{}}(S,\overline{S})`$, hence $`Hom_{𝒞\mathrm{}M_{}}(S,\overline{S})`$ is trivial. Conversely, if this bundle is trivial and if $`\stackrel{~}{\theta }^{}`$ is a non-vanishing section, then $`\stackrel{~}{\theta }^2=\lambda I_S`$ for some non-vanishing map $`\lambda C^{\mathrm{}}(M,)`$. But
$$\lambda (p)\stackrel{~}{\theta }_p(v)=\stackrel{~}{\theta }_p\stackrel{~}{\theta }_p\stackrel{~}{\theta }_p(v)=\stackrel{~}{\theta }(\lambda (p)v)=\overline{\lambda (p)}\stackrel{~}{\theta }(v),$$
i.e., $`\lambda C^{\mathrm{}}(M)`$ is real-valued, and replacing $`\stackrel{~}{\theta }`$ by $`\theta =|\lambda |^{1/2}\stackrel{~}{\theta }`$ we obtain a structural map.
Now given an irreducible complex spinor bundle $`S`$ and a structural map $`\theta )`$ we can choose a Riemannian structure compatible with Clifford multiplication and such that $`\theta `$ is an isometry. Moreover, we can choose a Riemannian connection $`^S`$ with properties (1) and (2) uniquely determined up to a purely imaginary one-form. If we also require that $`^S`$ is compatible with $`\theta `$, i.e.
$$^S\text{s}=(\mathrm{id}_{T^{}M}\theta )^{\overline{S}}(\theta \text{s}),\text{s}C^{\mathrm{}}(S),$$
then such a connection $`^S`$ is uniquely determined:
###### Theorem 8
If $`M`$ is a spin manifold of dimension $`m=2n`$ with corresponding spinor bundle $`S`$ and structural map $`\theta `$, then:
(a) On $`S`$ there exists a Riemannian structure compatible with Clifford multiplication and with $`\theta `$, i.e. $`\theta (\text{s}_1),\theta (\text{s}_2)=\overline{\text{s}_1,\text{s}_2}`$ for $`\text{s}_1,\text{s}_2C^{\mathrm{}}(S)`$.
(b) There is a unique Riemannian connection $`^S`$ with properties (1) and (2) and compatible with $`\theta `$.
Proof: (a) We consider $`\theta `$ as an antilinear map on $`S`$ and change a given Riemannian metric $`,^{}`$ with property (1) to
$$\text{s}_1,\text{s}_2=\frac{1}{2}\left(\text{s}_1,\text{s}_2^{}+\overline{\theta (\text{s}_1),\theta (\text{s}_2)^{}}\right).$$
Then the new metric will also be compatible with Clifford multiplication. Moreover, one has
$$\theta (\text{s}_1),\theta (\text{s}_2)=\frac{1}{2}\left(\theta (\text{s}_1),\theta (\text{s}_2)^{}+\overline{\theta ^2(\text{s}_1),\theta ^2(\text{s}_2)^{}}\right)=\overline{\text{s}_1,\text{s}_2}.$$
In particular, $`\theta (\text{s}_1),\text{s}_2=\pm \theta (\text{s}_2),\text{s}_1`$, hence $`\theta (\text{s}_1),\text{s}_1=0`$ in the quaternionic case.
(b) It suffices to prove uniqueness. We choose a local orthonormal frame $`\text{s}_j`$ of $`S`$ (which is a local orthonormal frame of $`\overline{S}`$ simultaneously) and the corresponding local connection form $`\omega `$. In the real case $`S`$ is a complexified real spinor bundle, and we can choose the frame such that $`\theta (\text{s}_j)=\text{s}_j`$, $`j=1,\mathrm{},2^n`$. In the quaternionic case we can choose the frame such that $`\text{s}_{2^{n1}+j}=\theta (\text{s}_j)`$, $`j=1,\mathrm{},2^{n1}`$. By compatibility of $`^S`$ and $`\theta `$ in the real case we obtain
$$^{\overline{S}}\theta (\text{s}_j)=\underset{i=1}{\overset{2^n}{}}\omega _{ji}\text{s}_i=\underset{i=1}{\overset{2^n}{}}\overline{\omega }_{ji}\text{s}_i=\theta (^S\text{s}_j),$$
i.e. $`\omega =\overline{\omega }`$. In the quaternionic case we obtain
$`^{\overline{S}}\theta (\text{s}_j)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{2^n}{}}}\omega _{j+2^{n1},i}\text{s}_i={\displaystyle \underset{i=1}{\overset{2^{n1}}{}}}\overline{\omega }_{ji}\theta (\text{s}_i)+{\displaystyle \underset{i=1}{\overset{2^{n1}}{}}}\overline{\omega }_{j,i+2^{n1}}\theta (\text{s}_{i+2^{n1}})`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{2^{n1}}{}}}\overline{\omega }_{ji}\text{s}_{i+2^{n1}}{\displaystyle \underset{i=1}{\overset{2^{n1}}{}}}\overline{\omega }_{j,i+2^{n1}}\text{s}_i=\theta (^S\text{s}_j),`$
for $`j=1,\mathrm{},2^{n1}`$, i.e.
$$\omega _{j+2^{n1},i}=\{\begin{array}{cc}\overline{\omega }_{j,i+2^{n1}},\hfill & i=1,\mathrm{},2^{n1}\text{,}\hfill \\ \overline{\omega }_{j,i2^{n1}},\hfill & i=2^{n1}+1,\mathrm{},2^n,\hfill \end{array}$$
and, in particular, $`\omega _{jj}=\overline{\omega }_{j+2^{n1},j+2^{n1}}`$ for $`j=1,\mathrm{},2^{n1}`$.
Thus, in both cases addition of a purely imaginary one-form is prohibited.
Examples 8. Any oriented complex manifold (or, more generally, an almost-complex manifold) is spin<sup>c</sup>: Since the complex cotangent bundle $`T^{}M_{}`$) splits orthogonally $`T^{}M_{}=(T^cM)^{}(\overline{T^cM})^{}=^{1,0}M^{0,1}M`$, we can choose $`S=^{}\overline{T^cM}=^{0,}M`$.
9. However, in general a complex manifold is not spin, e.g. it can be proved that $`P^n`$ is spin if and only if $`n`$ is odd.
10. Any oriented compact hyper surface $`M^{2n+1}`$ (that is the boundary of a compact $`2n+1`$-dimensional submanifold $`N`$ with boundary) is spin:
Using the matrices $`A_jM(^{2^n})`$, $`j=1,\mathrm{},2n+1`$ the Clifford multiplication $`E_jv=A_jv`$ for $`v^{2^n}=S_0`$ and the standard orthonormal frame $`E_1,\mathrm{},E_{2n+1}`$ of $`^{2n+1}`$ makes $`^{2n+1}\times ^{2^n}`$ a (trivial) complex spinor bundle over $`^{2n+1}`$. If we restrict to $`M`$ and consider $`TM`$ as a subbundle of $`TN|_M`$ the Clifford modules $`\{p\}\times ^{2^n}`$, $`pM`$, are irreducible $`𝒞\mathrm{}(T_pM_{})`$ modules, since we can generate $`𝒞\mathrm{}_{2n+1}^{}`$ by an orthonormal basis $`E_1^{}(p),\mathrm{},E_{2n}^{}(p)`$ of $`T_pM`$ and the exterior normal vector $`E_{2n+1}^{}(p)=X_N(p)`$. Therefore, $`H=M\times S_0`$ defines a spinor bundle for $`M`$. Since, moreover, the $`E_j^{}`$ as real linear combinations of the $`E_j`$ also commute with the structural map $`\theta `$ of $`S_0`$, we even have a spin structure. The grading operator on $`H`$ is defined by $`ϵ=iX_N`$ with the exterior normal vector field $`X_N`$ at $`M`$. If $`M=S^{2n}`$, we have $`X_N(x)=_{k=1}^{2n+1}x_kE_k`$ and $`H^{n(mod2)}`$, the bundles of half-spinors are non-trivial smooth vector bundles.
Special examples are oriented compact surfaces $`T_g`$ in $`^3`$ or spheres $`S^{2n}`$ in $`^{2n+1}`$. The former allow $`2^{2g}`$ different spin structures whereas there is only one spin structure on $`S^{2n}`$. To see this we need the following result.
###### Theorem 9
Let $`M`$ be a connected oriented Riemannian manifold. If $`M`$ carries a spin<sup>c</sup> structure, then all of the non-equivalent spin<sup>c</sup> structures are parametrized by $`H^2(M,)`$. If moreover $`M`$ is spin, then all of the different spin structures are parametrized by $`H^1(M,_2)Hom(\pi _1(M),_2)`$. In particular, $`M`$ allows at most one spin structure if $`M`$ is simply connected.
Proof: Starting from a irreducible complex spinor bundle $`S`$, any further irreducible complex spinor bundle on $`M`$ is of the form $`S^{}=SL`$ where $`L=Hom_{𝒞\mathrm{}M_{}}(S,S^{})`$. If $`S^{}`$ and $`S^{\prime \prime }=S(E)L^{}`$ are isomorphic as spinor bundles, i.e. determine equivalent spin<sup>c</sup> structures, there is a $`\mathrm{\Phi }Iso_{𝒞\mathrm{}M_{}}(S^{},S^{\prime \prime })`$, and so $`LL^{}`$. This shows that $`H^2(M,)`$ acts transitively on the set of different spin<sup>c</sup> structures. Now for $`Hom_{𝒞\mathrm{}M_{}}(S^{},\overline{S}^{})`$ we obtain
$`Hom_{𝒞\mathrm{}M_{}}(S^{},\overline{S}^{})`$ $``$ $`Hom_{𝒞\mathrm{}M_{}}(S_{}L,\overline{S}_{}\overline{L})Hom(L,Hom_{𝒞\mathrm{}M_{}}(S,\overline{S}_{}\overline{L}))`$
$``$ $`L^{}_{}Hom_{𝒞\mathrm{}M_{}}(S,\overline{S}_{}\overline{L})L^{}_{}Hom_{𝒞\mathrm{}M_{}}(S,\overline{S}_{}\overline{L}`$
$``$ $`L^{}_{}\overline{L}Hom(L,\overline{L}),`$
if $`Hom_{𝒞\mathrm{}M_{}}(S,\overline{S})`$ is trivial. Therefore, there is a structural map on $`S^{}`$ if and only if $`L^{}\overline{L}\overline{L}^2`$ is trivial. If $`H^2(M,)`$ has no 2-torsion, $`\overline{L}`$ has to be trivial, too, and likewise $`L`$. In any case different spin structures are classified by isomorphy classes of real line bundles, i.e., by $`H^1(M,_2)`$; cf. \[Kar\].
Remarks 1. We always started with the Clifford bundle of the tangent bundle. Only with literate changes we can start with a real Riemannian vector bundle $`E`$ of even rank. A spin<sup>c</sup> structure is then given by a complex spinor bundle $`S(E)`$ with $`𝒞\mathrm{}^{}(E)`$ acting irreducibly on the fibers, and a spin structure by an additional structural map compatible with Clifford multiplication. If $`E`$ comes with a Riemannian connection $`^E`$ there is unique connection $`^{𝒞\mathrm{}(E)}`$ on $`𝒞\mathrm{}(E)`$ and in the spin case a unique Riemannian connection $`^{S(E)}`$ on $`S(E)`$ that satisfy properties (1) and (2) and
$$^{S(E)}(\sigma \text{s})=^{𝒞\mathrm{}(E)}(\sigma )\text{s}+\sigma \left(^{S(E)}\text{s}\right),$$
for $`\sigma C^{\mathrm{}}(𝒞\mathrm{}(E))`$, $`\text{s}C^{\mathrm{}}(S(E))`$.
2. On an oriented Riemannian vector bundle $`E`$ of odd rank $`m=2n+1`$ (in particular, on an odd-dimensional Riemannian manifold) spin<sup>c</sup> or spin structures can be defined, too. Here a spin<sup>c</sup> structure is given by a complex spinor bundle $`S(E)`$, on which $`𝒞\mathrm{}^{}(E)`$ acts irreducibly, and where for each oriented orthonormal frame $`e_1(p),\mathrm{},e_m(p)`$ of $`E_p`$ the element $`i^{n+1}e_1(p)\mathrm{}e_m(p)`$ acts as $`I_{E_p}`$.
## 4 Spin groups and principal bundles
There are topological obstructions for a spin<sup>c</sup> or a spin structure to exist on a manifold $`M`$. We know that if $`M`$ is spin<sup>c</sup> and $`S`$ an irreducible complex spinor bundle structure then $`M`$ is spin if and only if $`L=Hom_{𝒞\mathrm{}M_{}}(S,\overline{S})`$ is trivial. Now if $`M`$ is simply connected this can be decided by computing a topological invariant. It is well known (cf. \[Sdr2\]) that $`L`$ is trivial if and only if the first Chern class $`c_1(L)`$ vanishes. But this does not apply in general if $`M`$ is not simply connected. Then the obstructions are better expressed in terms of the so-called second Stiefel-Whitney class $`w_2(TM)`$, an element of $`H^2(M,_2)`$ (cf. \[Hae\]). This is a cohomology class with coefficients in $`_2=\{\pm 1\}`$, and can be represented by lifts of cocycles of $`SO(n)`$-valued transition maps to the covering group $`Spin(n)`$.
At this point we have to digress and take a closer look at the covering group $`Spin(n)`$ of $`SO(n)`$. Here again Clifford algebras are the appropriate tool to generalize classical constructions. We first inspect how Clifford algebras help represent orthogonal transformations. It is well known that $`S^3`$ is the two-fold simply connected covering of the Lie group $`SO(3)`$. Identifying $`^3`$ with $`\mathrm{Im}=\{is+jt+kus,t,u\}`$ an element $`xS^3=\{y|y|^2=\overline{y}y=1\}`$ acts on $`^3`$ by
$$Ad_x(v)=xvx^1=xv\overline{x},v^3.$$
Note that $`x`$ and $`x`$ define the same element of $`SO(3)`$. More generally one could use any $`x^{}=\{0\}`$ since $`Ad_x=Ad_{x/|x|}`$.
To find the covering group of $`SO(n)`$ for $`n4`$ or of $`SO(E)`$ for a Euclidean vector space $`E`$ we start from the regular group $`G𝒞\mathrm{}(E)`$ of invertible elements of the algebra $`𝒞\mathrm{}(E)`$. For $`xE\{0\}G𝒞\mathrm{}(E)`$ and $`vE𝒞\mathrm{}(E)`$ we have $`xv+vx=2x,v1`$, hence
$$Ad_x(v)=v2\frac{x,v}{x,x}x.$$
From a geometric point of view this is the reflection at the hyperplane perpendicular to $`x`$. Using the involution $`\alpha `$ (that induces the grading $`𝒞\mathrm{}(E)=𝒞\mathrm{}(E)^0𝒞\mathrm{}(E)^1`$) we pass over to the “twisted” adjoint representation on $`E`$ given by
$$\stackrel{~}{Ad}_x(v)=\alpha (x)vx^1$$
which is is naturally defined on the Clifford group
$$\mathrm{\Gamma }(E)=\{xG𝒞\mathrm{}(E)\alpha (x)vx^1E\text{ for all }vE\}.$$
###### Proposition 5
The twisted adjoint representation $`\stackrel{~}{Ad}:\mathrm{\Gamma }(E)Aut(E)`$ is a homomorphism of groups and induces an exact sequence
$$1^{}\mathrm{\Gamma }(E)\stackrel{\stackrel{~}{Ad}}{}O(E)1.$$
Any $`x\mathrm{\Gamma }(E)`$ can be written as $`x=v_1\mathrm{}v_k`$, $`v_iE`$, $`v_i0`$, $`i=1,\mathrm{},k`$.
Proof: Obviously, $`\stackrel{~}{Ad}`$ is a homomorphism. Next we show that $`x^{}=\{0\}\mathrm{\Gamma }(E)`$ if $`\alpha (x)v=vx`$ for all $`vE`$ or equivalently if this holds elements $`v`$ of an orthonormal basis $`(e_i)_{1in}`$ of $`E`$. To this end we write $`x=x^0+x^1𝒞\mathrm{}(E)^0𝒞\mathrm{}(E)^1`$ with $`x^0=a_i^0+e_ib_i^1`$ and $`x^1=a_i^1+e_ib_i^0`$, where $`a_i^j`$ and $`b_i^j`$ are of degree $`j`$ (mod 2) and both do not contain $`e_i`$. Then we get
$$\alpha (x)e_i=(x^0x^1)e_i=e_i(a_i^0+a_i^1)+b_i^1+b_i^0$$
and
$$e_ix=e_i(x^0+x^1)=e_i(a_i^0+a_i^1)b_i^1b_i^0,$$
which entails $`b_i^0=b_i^1=0`$, i.e. $`x^{}`$.
Since $`O(E)`$ is generated by reflections it is at least contained in the image of $`\stackrel{~}{Ad}`$.
It remains to show $`\stackrel{~}{Ad}\left(\mathrm{\Gamma }(E)\right)O(E)`$, i.e. $`|\stackrel{~}{Ad}_x(v)|=|v|`$ for $`vE`$. To prove this we consider the anti-automorphism of $`𝒞\mathrm{}(E)`$ induced by
$$x=v_1\mathrm{}v_kx^t=v_k\mathrm{}v_1$$
and the anti-automorphism
$$𝒞\mathrm{}(E)x\overline{x}=\alpha (x^t)=\left(\alpha (x)\right)^t𝒞\mathrm{}(E)$$
which allows to extend the quadratic form $`Evv\overline{v}=v,v1=|v|^21𝒞\mathrm{}(E)`$ to the so-called spinor norm
$$𝒞\mathrm{}(E)xN(x)=x\overline{x}𝒞\mathrm{}(E)$$
of $`𝒞\mathrm{}(E)`$. Since the anti-automorphisms leave $`\mathrm{\Gamma }(E)`$ invariant, we have $`N\left(\mathrm{\Gamma }(E)\right)\mathrm{\Gamma }(E)`$. Actually $`N\left(\mathrm{\Gamma }(E)\right)^{}`$, because
$`\stackrel{~}{Ad}_{N(\overline{x})}(v)`$ $`=`$ $`\alpha \left(\alpha (x^t)x\right)v\left(\alpha (x^t)x\right)^1=x^t\alpha (x)vx^1\alpha (x^1)^t`$
$`=`$ $`\left(\alpha (x^1)\alpha (x)vx^1x\right)^t=v.`$
Now $`N|_{\mathrm{\Gamma }(E)}`$ is a homomorphism of groups, since
$$N(xy)=xy\overline{xy}=xy\alpha (y^t)\alpha (x^t)=xN(y)\alpha (x^t)=N(x)N(y),$$
as $`N(\mathrm{\Gamma }(E))^{}`$. In particular,
$$N\left(\alpha (x)vx^1\right)=N\left(\alpha (x)\right)N(v)N(x)^1=N(v)N\left(\alpha (x)\right)N(x)^1=N(v),$$
since $`N\left(\alpha (x)\right)=\alpha (x)x^t=\alpha \left(N(x)\right)=N(x)^{}`$, and we conclude
$$|\stackrel{~}{Ad}_x(v)|^2=|\alpha (x)vx^1|^2=|v|^2,$$
i.e. $`\stackrel{~}{Ad}_xO(E)`$.
###### Definition 9
We put $`Pin(E)=N^1(1)\mathrm{\Gamma }(E)`$ and define the spin group of the Euclidean vector space $`E`$ by $`Spin(E)=Pin(E)𝒞\mathrm{}(E)^0`$. In the case $`E=^n`$ with its standard inner product we write $`Spin(n)`$ instead of $`Spin(^n)`$.
Remarks 1. The group $`Spin(E)`$ is compact, in fact a Lie group as a closed subgroup of the group of invertibles of the algebra $`𝒞\mathrm{}^0(E)`$.
2. Of course, $`Pin(E)`$ and $`Spin(E)`$ both depend on the Euclidean structure. More generally, one can also define $`Spin(E,Q)`$ for a real vector space $`E`$ and a non-degenerate quadratic form $`Q`$.
3. One has $`\mathrm{Pin}(E)=\{v_1\mathrm{}v_k𝒞\mathrm{}(E)v_iE,v_i,v_i=1,i=1,\mathrm{},k\}`$ and $`Spin(E)=\{v_1\mathrm{}v_{2k}𝒞\mathrm{}(E)v_iE,v_i,v_i=1,i=1,\mathrm{},2k\}`$.
Corollary The groups $`Pin(E)`$ and $`Spin(E)`$ fit into the following exact sequences
$$1_2Pin(E)O(E)1$$
$$1_2Spin(E)SO(E)1.$$
In particular,
$$1_2Spin(n)SO(n)1$$
is exact, i.e., $`Spin(n)`$ is a non-trivial two-sheeted covering of $`SO(n)`$. For $`n3`$ it is simply connected, i.e. the universal covering group of $`SO(n)`$.
Proof: Given $`x\mathrm{\Gamma }(E)`$ and $`\lambda =1/\sqrt{N(x)}`$ one has $`\lambda xPin(E)`$ hence
$$\stackrel{~}{Ad}|_{Pin(E)}:Pin(E)O(E)$$
is onto and
$$ker\stackrel{~}{Ad}|_{Pin(E)}=\{\lambda ^{}N(\lambda )=\lambda ^2=1\}_2.$$
Any element of $`SO(E)`$ may be written as $`\stackrel{~}{Ad}_{v_1}\mathrm{}\stackrel{~}{Ad}_{v_{2k}}`$ hence
$$\rho =\stackrel{~}{Ad}|_{\mathrm{\Gamma }(E)𝒞\mathrm{}(E)^0}:\mathrm{\Gamma }(E)𝒞\mathrm{}(E)^0SO(E)$$
is onto with $`ker\rho =^{}`$. Now the restriction to $`Spin(E)`$ yields the analogous exact sequence. To prove the last assertion we only have to find a continuous path connecting $`+1`$ and $`1`$ in $`Spin(n)`$. To this end we choose $`e_1,e_2^n`$ with $`e_1e_2,|e_i|=1`$, and
$`c(t)`$ $`=`$ $`\mathrm{exp}(2\pi te_1e_2)=\mathrm{cos}2\pi t+e_1e_2\mathrm{sin}2\pi t`$
$`=`$ $`(e_1\mathrm{cos}\pi t+e_2\mathrm{sin}\pi t)(e_1\mathrm{cos}\pi t+e_2\mathrm{sin}\pi t),`$
for $`0t\frac{1}{2}`$. Thus the covering is non-trivial and $`Spin(n)`$ is connected. For $`n3`$ it is also simply connected by the classic topological result $`\pi _1\left(SO(n)\right)=_2`$, $`n3`$.
If $`E_{}`$ is the complexification of $`E`$ with $``$-linear extension $`Q_{}`$ of $`Q`$, then $`𝒞\mathrm{}(E_{},Q_{})`$ and $`𝒞\mathrm{}(E,Q)`$ are isomorphic. We put $`\alpha (xz)=\alpha (x)z`$ and $`(xz)^t=x^t\overline{z}`$ and with $`\overline{}`$ and $`N`$ as before we also define $`Pin^c(E)`$ and the group $`Spin^c(E)𝒞\mathrm{}^0(E,Q)`$. The latter is isomorphic with $`Spin(E)\times S^1/_2`$ where $`_2=\{(1,1),(1,1)\}`$. If $`E=^n`$ we simply denote it by $`Spin^c(n)`$. The group $`Spin^c(E)`$ is also compact and fits into the exact sequences
$$1S^1Spin^c(E)\stackrel{\rho _0}{}SO(E)1$$
$$1Spin(E)Spin^c(E)\stackrel{\rho _1}{}S^11,$$
where the left hand homomorphisms are canonical inclusions and the right hand ones are defined by $`\rho _0([(x,z)])=\rho (x)`$ and $`\rho _1([(x,z)])=z^2`$, $`(x,z)Spin(E)\times S^1`$, respectively.
Usually, spin and spin<sup>c</sup> structures are defined with the help of corresponding principal bundles; cf. \[BH\] and \[Mil\]. One starts with the orthonormal frame bundle $`P_{SO(m)}`$ of the tangent bundle of an $`m`$-dimensional oriented Riemannian manifold $`M`$ (or of an oriented Riemannian vector bundle of rank $`m`$). A spin structure for $`M`$ consists of a principal bundle $`P_{Spin(m)}`$ with structure group $`Spin(m)`$ and a two-sheeted covering
$$\xi :P_{Spin(m)}P_{SO(m)}\text{ with }\xi (pg)=\xi (p)\rho _0(g),pP_{Spin(m)},gSpin(m),$$
where $`\rho _0:Spin(m)SO(m)`$ is the standard covering. A spin<sup>c</sup> structure is given by a principal bundle $`P_{Spin^c(m)}`$ and a map
$$\xi :P_{Spin^c(m)}P_{SO(m)}\text{ with }\xi (pg)=\xi (p)\rho _0(g),pP_{Spin^c(m)},gSpin^c(m)$$
where $`\rho _0:Spin^c(m)SO(m)`$ is again the standard map.
To show that this approach is equivalent with the one presented so far one has to go two ways. A spinor bundle can be obtained as an associated bundle: If $`F`$ is a real or a complex vector space, which is also a $`𝒞\mathrm{}_m`$-module or a $`𝒞\mathrm{}_m^{}`$-module with compatible inner product, representations $`\rho :Spin(m)SO(F)`$ or $`\rho :Spin^c(m)U(F)`$ will be induced by left-multiplication with elements of $`Spin(m)𝒞\mathrm{}_m^0`$ or $`Spin^c(m)𝒞\mathrm{}_m^0`$, respectively. Then $`S=P_{Spin(n)}\times _\rho F`$ is a real or a complex spinor bundle, which moreover is irreducible if $`F=S_0`$ the space of spinors.
If on the other hand a spin structure is given by an irreducible complex spinor bundle $`S`$ the corresponding principal bundles can be recovered as follows. First recall that $`P_{SO(m)}`$ can be considered as the subset of $`Hom(M\times ^m,TM)`$ that consists of all orientation preserving isometries $`f_p:^mT_pM`$, $`pM`$. Then we define $`P_{Spin(m)}`$ and $`P_{Spin^c(m)}`$ to be appropriate subsets of $`Hom(M\times S_0,S)`$. In the second case it consists of all isometries $`\varphi _p:S_0S_p`$ that respect the decompositions $`S_0^0S_0^1`$ and $`S_p^0S_p^1`$ and satisfy $`\varphi _p(v\varphi _p^1)T_pM𝒞\mathrm{}(T_pM_{})`$ for all $`v^m𝒞\mathrm{}_m^{}=End(S_0)`$. In the first case we additionally require that these isometries respect the real or quaternionic structure. The map $`\xi :P_{Spin^c(m)}P_{SO(m)}`$ is now defined by $`\xi (\mathrm{\Phi }_p)=Ad(\mathrm{\Phi }_p)`$. Then one has
$$\xi (\mathrm{\Phi }_pg)=Ad(\mathrm{\Phi }_pg)=Ad(\mathrm{\Phi }_p)Ad(g)=\xi (\mathrm{\Phi }_p)\rho _0(g).$$
This action from the right is transitive, since for $`\mathrm{\Phi }_p,\mathrm{\Phi }_p^{}P_{Spin^c(m)}`$ one has $`\mathrm{\Phi }_p=\mathrm{\Phi }_p^{}((\mathrm{\Phi }_p^{})^1\mathrm{\Phi }_p)`$ and by definition $`x=(\mathrm{\Phi }_p^{})^1\mathrm{\Phi }_p𝒞\mathrm{}_m^0`$ as well as $`N(x)=1`$, hence $`xSpin^c(m)`$. Here $`N(x)=1`$ does hold, since $`x`$ is unitary and since $`\overline{x}=x^{}`$ for $`x𝒞\mathrm{}_m=End(S_0)`$ as $`\alpha (v)=v=v^{}`$ for $`v^m`$. In the spin case $`\xi `$ is defined likewise and obviously such an element $`x`$ belongs to $`Spin(m)`$.
Finally, we can come back to the topological obstructions that decide upon spin<sup>c</sup> or spin structures. A spin<sup>c</sup> structure can be supplied if and only if $`w_2(TM)`$ is the $`mod2`$-reduction of some integral cohomology class (or, what amounts to the same, if the integral Stiefel-Whitney class $`W_3(TM)`$ vanishes). A spin structure exists if and only if $`w_2(TM)=1`$. We refer to \[Kar2\], where the first assertion is proved explicitly and the second one implicitly – in the case of a spin structure in the notation of \[Kar2\] one has to replace $`^{}`$ by $`^{}`$, which makes $`l_1`$ automatically an isomorphism. These conditions can be checked combinatorically. We refer to \[Gil\] and \[LM\] for some specific computations. In particular, it is proved that $`w_2(TP^n)`$ and $`w_2(TP^{2n+1})`$ only vanish if $`n`$ is odd. Using deeper results of algebraic topology one can show that any compact oriented 3-manifold is spin (since according to E. Stiefel it is parallizable) and that any compact oriented 4-manifold is spin<sup>c</sup> (according to a theorem of Whitney; cf. \[HH\]).
## 5 The geometric Dirac operators
Now we want to look more closely at some Dirac operators. First we consider the special case $`M=^m`$ with its standard metric and the global orthonormal frame $`E_j=\frac{}{x_j}`$, $`j=1,\mathrm{},m`$, of $`T^m`$. If $`V`$ is an $`n`$-dimensional $`𝒞\mathrm{}_m`$-module defined by an algebra homomorphism $`\rho :𝒞\mathrm{}_mEnd(V)`$, say $`\rho (e_j)=A_j`$, $`\rho (1)=I_V`$, and $`(v_i)_{1in}`$ is a basis of $`V`$ a global frame on $`E=^m\times V`$ is given by $`\text{s}_i(p)=(p,v_i)`$, $`p^m`$, $`i=1,\mathrm{},n`$. Let $``$ denote a flat connection on $`E`$, i.e. $`\omega 0`$ with respect to the frame $`\text{s}_i`$, hence $`f\text{s}_i=df\text{s}_i`$ for $`fC^{\mathrm{}}(^m)`$. Then Clifford multiplication $`E_j\text{s}_i`$ is given by
$$(E_j\text{s}_i)(p)=(p,A_j(v_i)),$$
hence
$$D\left(\underset{i=1}{\overset{n}{}}f_i\text{s}_i\right)(p)=\underset{j=1}{\overset{m}{}}\underset{i=1}{\overset{n}{}}E_j_{E_j}(f_i\text{s}_i)(p)=(p,\underset{j=1}{\overset{m}{}}\underset{i=1}{\overset{n}{}}\frac{}{x_j}f_i(p)A_j(v_i)).$$
Let $`A_j`$ also denote the matrix with respect to the basis $`(v_i)`$. Then a (local) representation of $`D`$ is given by
$$Df=\underset{j=1}{\overset{n}{}}A_j\frac{}{x_j}f,$$
where $`f=(f_1,\mathrm{},f_n)^T`$. In particular
$$D^2f=\underset{j,k=1}{\overset{m}{}}A_jA_k\frac{}{x_j}\frac{}{x_k}f=\mathrm{\Delta }I_nf=\underset{j=1}{\overset{m}{}}\frac{^2}{x_j^2}f,$$
since
$$A_jA_k+A_kA_j=\rho (e_je_k+e_ke_j)=\rho (2\delta _{jk})=2\delta _{ik}I_V.$$
Thus $`D`$ is a square-root of $`\mathrm{\Delta }`$. In cases $`m=1,2`$ we have the following classical operators.
Examples 11. If $`m=1`$, i.e. $`𝒞\mathrm{}_1=`$, we choose $`V=^2`$ with $`\rho (e_1)=i`$ and obtain $`D=i\frac{}{x}`$.
12. If $`m=2`$, i.e. $`𝒞\mathrm{}_2=`$, we choose $`V=`$ with $`\rho (e_1)=i`$, $`\rho (e_2)=j`$ and get a grading $`𝒞\mathrm{}_2=𝒞\mathrm{}_2^0𝒞\mathrm{}_2^1`$ by
$`𝒞\mathrm{}_2^0u+ve_2e_1u+iv`$
$`𝒞\mathrm{}_2^1ue_1+ve_2u+iv.`$
Identifying $`E=^2\times V`$ and $`\times ()`$ the Dirac operator $`D=i\frac{}{x_1}+j\frac{}{x_2}`$ becomes
$$\frac{1}{2}D(fg)=\frac{}{z}g\frac{}{\overline{z}}f,f,gC^{\mathrm{}}().$$
If we write $`D=D^0D^1`$ with $`D^j:C^{\mathrm{}}(^2\times 𝒞\mathrm{}_2^j)C^{\mathrm{}}\left(^2\times 𝒞\mathrm{}_2^{j+1(mod2)}\right)`$ then $`\frac{1}{2}D^0`$ is just the Cauchy-Riemann operator $`\overline{}=\frac{}{\overline{z}}`$ which is studied in the theory of complex functions.
If $`M`$ is a (compact) oriented Riemannian manifold there are several Dirac operators related to additional geometric structures. We cannot go into the analytic properties of these Dirac operators; cf. \[Gil\], \[LM\], or \[Sdr1,2\]. We only note that they are symmetric (elliptic) differential operators.
If $`M`$ is of even dimension $`m=2k`$ we have a global section $`\omega C^{\mathrm{}}(𝒞\mathrm{}M_{})`$ which is locally given by
$$\omega =i^kE_1\mathrm{}E_m$$
with respect to an oriented orthonormal frame $`(E_i)_{1im}`$ of $`TM`$. Obviously, one has $`\omega ^2=1`$ and $`\omega X=X\omega `$ for $`XC^{\mathrm{}}(TM)`$. Since $`\omega `$ does not depend on the local frame we may assume $`\overline{}_{E_i}E_j(p)=0`$ at a fixed point $`pM`$ and conclude that
$$\overline{}_{E_i}\omega (p)=i^{\mathrm{}}\underset{j=1}{\overset{m}{}}E_1\mathrm{}_{E_i}E_j\mathrm{}E_m(p)=0,$$
hence $`\overline{}\omega =0`$. Using $`\omega `$ any Dirac bundle $`E`$ on $`M`$ will be graded by $`E^0=(1+\omega )E`$ and $`E^1=(1\omega )E`$. For $`\text{s}C^{\mathrm{}}(E^j)`$ and $`XC^{\mathrm{}}(TM)`$ we then obtain
$$_X\text{s}=(1)^j_X(\omega \text{s})=(1)^j\left((\underset{X}{\overline{}}\omega )\text{s}+\omega _X\text{s}\right)=(1)^j\omega _X\text{s},$$
i.e. $`_X\text{s}=\frac{1}{2}\left(1+(1)^j\omega \right)_X\text{s}C^{\mathrm{}}(E^j)`$, and
$$X\text{s}=(1)^jX\omega \text{s}=(1)^{j+1}\omega X\text{s},$$
hence $`X\text{s}=\frac{1}{2}\left(1+(1)^{j+1}\omega \right)X\text{s}C^{\mathrm{}}(E^{j+1mod2})`$. Since $`\omega \text{s},\text{t}=\text{s},\omega \text{t}`$ for $`\text{s},\text{t}C^{\mathrm{}}(E)`$, the decomposition $`E=E^0E^1`$ is orthogonal. This gives rise to the following definition:
###### Definition 10
Let $`E`$ be a Dirac bundle on $`M`$. An orthogonal decomposition $`E=E^0E^1`$ is called admissible if
(1) $`𝒞\mathrm{}^i(M)E^jE^{i+jmod2}`$,
(2) $`_X\text{s}C^{\mathrm{}}(E^i)`$ for $`\text{s}C^{\mathrm{}}(E^i)`$, $`XC^{\mathrm{}}(TM)`$.
Now given an admissible Dirac bundle $`E=E^0E^1`$ the corresponding Dirac operator induces first order differential operators
$$D^i:C_0^{\mathrm{}}(E^i)C_0^{\mathrm{}}(E^{i+1mod2}),i=0,1.$$
If $`M`$ is compact $`D`$ and $`D^j`$ are elliptic and extend to bounded operators on appropriate Sobolev space. These extensions are Fredholm operators, i.e. have an index
$$indD=dim\mathrm{ker}DdimCokerD.$$
Of course, $`indD=0`$ but
$$indD^0=dim\mathrm{ker}D^0dim\mathrm{ker}D^1$$
turns out to be an interesting geometric invariant. We already met the Kähler-Dirac operator $`D=d+\delta `$. Since the Dirac bundle $`(^{}M=^{\mathrm{ev}}M^{\mathrm{odd}}M,\overline{})`$ is admissible, we obtain $`D=D^0D^1`$ and $`indD^0=\chi (M)`$, the Euler characteristic of $`M`$. There is a different admissible decomposition of $`^{}M`$ given by $`\omega `$. If $`M`$ is of dimension $`m=4k`$ the index of the corresponding Dirac operator is just the signature of $`M`$; cf. \[Gil\], \[LM\], or \[Sdr2\]. We can now, finally, define the Spin-Dirac or Atiyah-Singer operator.
###### Definition 11
Let $`M`$ be a compact spin manifold of dimension $`m=2n`$ with spinor bundle $`S`$, and $`^S`$ the unique connection on $`S`$ that is induced by the Levi-Civita connection. The Dirac operator associated with the Dirac bundle $`S`$ is called the Spin-Dirac operator or Atiyah-Singer operator and will be denoted by $`D_{AS}`$.
The Dirac bundle $`S=S^0S^1`$ with decomposition induced by $`\omega `$ is admissible, i.e. $`D_{AS}=D_{AS}^0D_{AS}^1`$. The operator $`\begin{array}{c}/\hfill \\ D\hfill \end{array}=D_{AS}^0`$ is also often called the Spin-Dirac operator.
Its index $`ind\begin{array}{c}/\hfill \\ D\hfill \end{array}=\widehat{A}(M)`$ is called the $`\widehat{A}`$-genus of $`M`$. It has topological significance which is expressed by the famous Atiyah-Singer index theorem:
$$ind\begin{array}{c}/\hfill \\ D\hfill \end{array}=_m\widehat{A}(TM)$$
where $`\widehat{A}(TM)H^m(M,)`$ is the cohomology class first indroduced by F. Hirzebruch in 1954; cf. \[Gil\] for a detailed history of the subject matter.
We finally take a closer look on the Dirac-Laplace operator $`D^2`$ and on its relation to the curvature tensor $`\text{curv }()`$ which is defined by
$$\text{curv }()(X,Y)=[_X,_Y]_{[X,Y]}$$
for vector fields $`X`$ and $`Y`$. Here the brackets denote the respective commutators. Note that the curvature tensor $`\text{curv }()(X,Y)`$ of a Riemannian connection is skew-adjoint,
$$\text{curv }(X,Y)\text{s},\text{t}=\text{s},\text{curv }(X,Y)\text{t}(R)$$
and that
$$\text{curv }()(X,Y)(\sigma \text{s})=\text{curv }(\overline{})(X,Y)(\sigma )\text{s}+\sigma \text{curv }()(X,Y)\text{s},(D)$$
if $``$ satisfies property (2) of a Dirac triple.
Recall that the second covariant derivative
$$_{X,Y}^2:C^{\mathrm{}}(E)C^{\mathrm{}}(E),$$
is defined for $`X,YC^{\mathrm{}}(TM)`$ by
$$_{X,Y}^2\text{s}=_X_Y\text{s}_{\overline{}_XY}\text{s},\text{s}C^{\mathrm{}}(E),$$
and the curvature tensor of $``$ is given by
$$\text{curv }()(X,Y)=_{X,Y}^2_{Y,X}^2.$$
Since $`_{X,Y}^2\text{s}(p)`$ only depends on $`X_p`$ and $`Y_p`$ it makes $`_,^2`$ and $`\text{curv }()`$ tensors with values in $`E_p`$. The Bochner-Laplace operator of the connection $``$ is defined as
$$^{}\text{s}=tr(_,^2\text{s}),\text{s}C^{\mathrm{}}(E),$$
i.e. as
$$^{}\text{s}=\underset{j=1}{\overset{m}{}}_{E_j,E_j}^2\text{s}$$
when computed in some local orthonormal frame $`(E_j)_{1jm}`$. This definition does not depend on the chosen frame. Moreover, $`^{}:C^{\mathrm{}}(E)C^{\mathrm{}}(E)`$ is a second order (elliptic) differential operator.
Using the smooth section $`C^{\mathrm{}}(Hom(E,E))`$ given by
$$(\text{s})=\frac{1}{2}\underset{j,k=1}{\overset{m}{}}E_jE_k\text{curv }()(E_j,E_k)(\text{s})$$
we obtain the following fundamental result.
###### Theorem 10
(Bochner-Weitzenböck) Let $`E`$ be a Dirac bundle over $`M`$ with associated Dirac operator $`D`$. Then the Dirac-Laplace operator satisfies $`D^2`$
$$D^2=^{}+.$$
Proof: With the frame $`(E_j)_{1jm}`$ at $`p`$ as above we have
$`D^2`$ $`=`$ $`{\displaystyle \underset{j,k=1}{\overset{m}{}}}E_j_{E_j}(E_k_{E_k})={\displaystyle \underset{j,k=1}{\overset{m}{}}}E_jE_k_{E_j}_{E_k}`$
$`=`$ $`{\displaystyle \underset{j,k=1}{\overset{m}{}}}E_jE_k_{E_j,E_k}^2`$
$`=`$ $`{\displaystyle \underset{j=1}{\overset{m}{}}}_{E_j,E_j}^2+{\displaystyle \underset{1j<km}{}}E_jE_k(_{E_j,E_k}^2_{E_k,E_j}^2)`$
$`=`$ $`^{}+.`$
A simple but important application is the following vanishing theorem.
Corollary If $`M`$ is compact and connected and $`(p)`$ positive semi-definite for all $`pM`$ and positive definite for at least one $`p`$, then the differential equation $`D^2\text{s}=0`$ has only the trivial solution $`\text{s}=0`$, i.e. there are no non-trivial harmonic sections.
Proof: For a fixed point $`pM`$ one can choose $`(E_j)_{1jm}`$ such that $`\overline{}_{E_i}E_j(p)=0`$, and given sections $`\text{s},\text{t}C^{\mathrm{}}(E)`$ there is a vector field $`X`$ with
$$X,Y=_Y\text{s},\text{t}_E,YC^{\mathrm{}}(TM).$$
These data help to prove that
$`^{}\text{s},\text{t}(p)`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{m}{}}}_{E_j}_{E_j}\text{s},\text{t}(p)`$
$`=`$ $`{\displaystyle \underset{j=1}{\overset{m}{}}}\left(E_j_{E_j}\text{s},\text{t}_{E_j}\text{s},_{E_j}\text{t}\right)(p)`$
$`=`$ $`\mathrm{div}X(p)+\text{s},\text{t}(p).`$
When integrated over $`M`$, by Gauß’ theorem, the divergence term does not occur and we obtain
$$0_M(\text{s}),\text{s}=_M^{}\text{s},\text{s}=_M\text{s},\text{s}0,$$
if $`D^2\text{s}=0`$. Therefore, $`\text{s}=0`$ and $`\text{s}`$ is constant, since $``$ is Riemannian. Assuming $`s(p)0`$ and $`(p)`$ positive definite gives $`_M(\text{s}),\text{s}_E>0`$, which cannot hold.
There are a lot of special cases of the Bochner-Weitzenböck formula. The Bochner-Weitzenböck formula for the Laplace operator can already be found in Weitzenböck’s monograph “Invariantentheorie” of 1923. It has been rediscovered and applied in 1946 by S. Bochner \[Boc\]. Here we only consider one special case and deduce a special vanishing result.
###### Theorem 11
If $`M`$ is a spin manifold with spinor bundle $`S`$ and connection $`^S`$, then the Spin-Dirac-Laplace operator $`D_{AS}^2`$ and the Bochner-Laplace operator $`^S^S`$ are related by
$$D_{AS}^2=^S^S+\frac{1}{4}\tau .$$
Here $`\tau `$ denotes the scalar curvature of the Riemannian manifold $`M`$.
Proof: We only have to prove that $`=\frac{1}{4}\tau `$. It suffices to show that with respect to a local orthonormal frame $`\{E_1,\mathrm{},E_m\}`$ of $`TM`$ the curvature $`\text{curv }(^S)`$ is given by
$$\text{curv }(^S)(X,Y)=\frac{1}{4}\underset{k,\mathrm{}=1}{\overset{m}{}}R(X,Y)E_k,E_{\mathrm{}}E_kE_{\mathrm{}},X,YC^{\mathrm{}}(TM_U)()$$
since then we obtain
$``$ $`={\displaystyle \frac{1}{2}}{\displaystyle \underset{i,j=1}{\overset{m}{}}}E_iE_j\text{curv }(^S)(E_i,E_j)`$
$`={\displaystyle \frac{1}{8}}{\displaystyle \underset{i,j,k,\mathrm{}=1}{\overset{m}{}}}R(E_i,E_j)E_k,E_{\mathrm{}}E_iE_jE_kE_{\mathrm{}}`$
$`={\displaystyle \frac{1}{8}}{\displaystyle \underset{\mathrm{}=1}{\overset{m}{}}}({\displaystyle \underset{ijk\mathrm{}}{}}R(E_i,E_j)E_k+R(E_k,E_i)E_j+R(E_j,E_k)E_i,E_{\mathrm{}}E_iE_jE_k`$
$`+{\displaystyle \underset{i,j}{}}R(E_i,E_j)E_i,E_{\mathrm{}}E_iE_jE_i+{\displaystyle \underset{i,j}{}}R(E_i,E_j)E_j,E_{\mathrm{}}E_iE_jE_j)E_{\mathrm{}}`$
$`={\displaystyle \frac{1}{4}}{\displaystyle \underset{i,j,\mathrm{}=1}{\overset{m}{}}}R(E_i,E_j)E_j,E_{\mathrm{}}E_iE_{\mathrm{}}={\displaystyle \frac{1}{4}}{\displaystyle \underset{i,j=1}{\overset{m}{}}}R(E_i,E_j)E_j,E_i={\displaystyle \frac{1}{4}}\tau `$
by the symmetries of the Riemann curvature tensor $`R`$ and by the definition of $`\tau `$.
Now it is a straight-forward computation to show that for fixed vector fields $`X`$ and $`Y`$ the right-hand side of $`()`$ which we denote by $`R(X,Y)`$ shares the same properties (R) and (D) as the left hand-side and so does their difference $`T=\text{curv }()(X,Y)R(X,Y)`$. In particular, by (D) it commutes with the left-action of $`𝒞\mathrm{}M`$ and so acts as multiplication by an element $`\gamma C^{\mathrm{}}(M,)`$ which by (R) is skew-adjoint, i.e. $`\gamma =i\eta `$ with $`\eta C^{\mathrm{}}(M,)`$. Actually, $`\eta `$ has to vanish, since $`T`$ also respects the real structure on $`S`$, i.e. commutes with the structural map $`\theta `$.
This Bochner-Weitzenböck formula for the Spin-Dirac operator is used by A. Lichnerowicz \[Lic\] to prove the following vanishing theorem. The relation of $`D_{AS}^2`$ and the scalar curvature had however already been noted by E. Schrödinger in 1932 \[Sch\].
Corollary (Lichnerowicz) Let $`M`$ be a compact spin manifold with positive scalar curvature. Then there are no harmonic spinors on $`M`$. If $`dimM=4k`$, then $`\widehat{A}(M)=0`$.
Proof: The first assertion is immediate while the second one is a consequence of the Atiyah-Singer index theorem.
We also study the twisted Dirac operator $`\begin{array}{c}/\hfill \\ D\hfill \end{array}I_E`$, where $`E`$ is a Hermitian vector bundle $`E`$ with connection $`^E`$, i.e. the Dirac operator of the Dirac bundle $`(SE,^{SE})`$. Let $`^E:C^{\mathrm{}}(SE)C^{\mathrm{}}(SE)`$ denote the zero order differential operator, which for sections $`\sigma \text{s}`$ and the frame $`(E_i)_{1im}`$ is defined by
$$^E(\sigma \text{s})=\frac{1}{2}\underset{j,k=1}{\overset{m}{}}E_jE_k\sigma \text{curv }(^E)(E_j,E_k)\text{s}.$$
###### Theorem 12
Let $`M`$ be spin and $`S`$ and $`E`$ as before. Then the Spin-Dirac operator $`D_{AS}I_E`$ and the Bochner-Laplace operator $`^{}`$ of the tensor bundle $`SE`$ are related by
$$(D_{AS}I_E)^2=^{}+\frac{1}{4}\tau +^E.$$
Here $`\tau `$ is again the scalar curvature of $`M`$ .
Proof: For $`\sigma C^{\mathrm{}}(S)`$ and $`\text{s}C^{\mathrm{}}(E)`$ we have
$$^{SE}(\sigma \text{s})=(^S\sigma )\text{s}+\sigma (^E\text{s}).$$
This entails
$$\text{curv }(^{SE})(\sigma \text{s})=\text{curv }(^S)(\sigma )\text{s}+\sigma \text{curv }(^E)(\text{s})$$
and
$`(\sigma \text{s})`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{j,k=1}{\overset{m}{}}}E_jE_k\text{curv }(^{SE})(E_j,E_k)(\sigma \text{s})`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{j,k=1}{\overset{m}{}}}E_jE_k\left(\text{curv }(^S)(E_j,E_k)(\sigma )\right)\text{s}`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{j,k=1}{\overset{m}{}}}E_jE_k\sigma \text{curv }(^E)(E_j,E_k)(\text{s})`$
$`=`$ $`{\displaystyle \frac{1}{4}}\tau (\sigma \text{s})+^E(\sigma \text{s}).`$
Remark For the Spin<sup>c</sup>-Dirac operator $`D_{S^c}`$ there is a Bochner-Weitzenböck formula, too. If $`M`$ is spin a spin<sup>c</sup> structure is given by $`S^c=SL`$ for some complex line bundle $`L`$. Choosing the product connection on $`S^c`$ with some Hermitian connection $`^L`$ on $`L`$ the square of the corresponding Spin<sup>c</sup>-Dirac operator $`D_{S^c}`$ satisfies
$$D_{S^c}^2=^{}+\frac{1}{4}\tau +^L.$$
Because of
$`^L(\sigma \text{s})`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{j,k=1}{\overset{m}{}}}E_jE_k\sigma \text{curv }(^L)(E_j,E_k)(\text{s})`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{j,k=1}{\overset{m}{}}}E_jE_k\sigma \mathrm{\Omega }^L(E_j,E_k)(\text{s})`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{j,k=1}{\overset{m}{}}}\mathrm{\Omega }^L(E_j,E_k)E_jE_k\sigma \text{s}=(\mathrm{\Omega }^L\sigma )\text{s}`$
(with $`\mathrm{\Omega }^L`$ denoting the curvature form of $`^L`$) we obtain
$$D_{S^c}^2=^{}+\frac{1}{4}\tau +\mathrm{\Omega }^L.$$
Since this computation is local, we can also apply it in the general non-spin case. Although $`S^c`$ is a product $`SL`$ only locally the line bundle $`L_{S^c}=Hom_{𝒞\mathrm{}M_{}}(\overline{S^c},S^c)=LL`$ is nevertheless globally defined. Choosing a Hermitian connection, the corresponding curvature form $`\mathrm{\Omega }`$ satisfies $`\mathrm{\Omega }=2\mathrm{\Omega }^L`$. Thus we obtain the Weitzenböck formula
$$D_{S^c}^2=^{}+\frac{1}{4}\tau +\frac{1}{2}\mathrm{\Omega }.$$
The index formula for $`D_{S^c}^0`$ can also be established by a local computation (cf. \[Sdr2\]). With $`c=c_1(L_{S^c})`$ one obtains
$$indD_{S^c}^0=_Me^{c/2}\widehat{A}(TM).$$
The non-vanishing of the $`\widehat{A}`$-genus is the simplest obstruction for a Riemannian metric with positive scalar curvature. N. Hitchin \[Hit\] has introduced an invariant $`\alpha (M)`$, which can be defined for spin manifolds $`M`$ of any dimension and which coincides with $`\widehat{A}(M)`$ if $`m=4k`$. It again vanishes in case of positive scalar curvature. For simply connected manifolds $`\alpha (M)=0`$ is even sufficient for such a metric to exist as S. Stolz \[Sto\] proved in 1989; cf. \[RS\] for a survey of the current state.
The Bochner-Weitzenböck formula for the Spin<sup>c</sup>-Dirac operator on oriented compact 4-manifolds is the footing of the so-called Seiberg-Witten theory in which the theoretical physicists N. Seiberg and E. Witten initiated new differential topological invariants in 1994. These lead to new essential contributions for the classification of 4-manifolds \[Mor\].
References
\[AT\] Anderson, F.W., Fuller, K.R.: Rings and Categories of Modules, Springer, New York - Heidelberg - Berlin, 1974
\[ABS\] Atiyah, M.F., Bott, R., Shapiro, A.: Clifford modules, Topology 3 (1964), Suppl. 1, 3-38
\[AS\] Atiyah, M.F., Singer, I.M.: The index of elliptic operators on compact manifolds, Bull. Amer. Math. Soc. 69 (1963) 422-432
\[BD\] Baum, P., Douglas, R.G.: Index theory, bordism, and $`K`$-homology, in: Operator Algebras and $`K`$-Theory, Contemp. Math. 10, 1-31, Amer. Math. Soc., Providence, RI, 1982
\[Boc\] Bochner, S.: Vector fields and Ricci curvature, Bull. Amer. Math. Soc. 52 (1946) 776-797
\[BBW\] Booss-Bavnbek, B., Wojciechowski, K.P.: Elliptic Boundary Problems for Dirac Operators, Birkhäuser, Basel-Boston, 1993
\[BH\] Borel A., Hirzebruch, F.: Characteristic classes on homogeneous soaces II, Amer. J. Math. 81 (1959) 315-382
\[BW\] Brauer, R., Weyl, H.: Spinors in $`n`$ dimensions, Amer. J. Math. 57 (1935) 425-449
\[Car\] Cartan, E.: Sur les groupes projectifs qui laissent invariante aucune multiplicité plane, Bull. Soc. Math. France 41 (1913) 53-96
\[Che\] Chevalley, C.: The Algebraic Theory of Spinors, Columbia Univ. Press, New York, 1954, together with The Construction and Study of Certain Important Algebras, Math. Soc. Japan, Tokyo, 1955, reprinted in: Collected Works Vol. 2, Springer, New York - Heidelberg - Berlin, 1998
\[Cli\] Clifford, W.K.: Applications of Grassmann’s extensive algebra, Amer. J. Math. 1 (1878) 350-358 or Math. Papers, 266-276; cf. also: On the classification of geometric algebras , ibid. 397-401
\[Con1\] Connes, A.: Noncommutative Geometry, Academic Press, New York, 1994
\[Con2\] Connes, A.: Noncommutative geometry and reality, J. Math. Phys. 36 (11) (1995) 6194-6231
\[Con3\] Connes, A.: Gravity coupled with matter and foundation of noncommutative geometry, Commun. Math. Phys. 182 (1996) 155-176
\[Dar\] Darwin, C.G.: The wave equation of the electron, Proc. Roy. Soc. London (A) 118 (1928) 654-680
\[Dir\] Dirac, P.A.M.: The quantum theory of the electron, Proc. Roy. Soc. London (A) 117 (1927) 610-624
\[Eck\] Eckmann, B.: Hurwitz-Radon matrices revisited: from effective solution of the Hurwitz matrix equations to Bott periodicity, CRM Proc. & Lecture Notes 6, 23-35, Amer. Math. Soc., 1994
\[Fri\] Friedrich, T.: Dirac Operatoren in der Riemannschen Geometrie, Vieweg, Braunschweig, 1997
\[Gil\] Gilkey, P.B.: Invariance Theory, the Heat Equation, and the Atiyah-Singer Index Theorem, (2nd ed.), CRC, Baton Rouge, 1995
\[GH\] Greub, W.H., Halperin, S.: An intrinsic definition of the Dirac operator, Collect. Math. 26 (1975) 19-37
\[Hae\] Haefliger, A.: Sur l’extension du groupe structural d’un espace fibré, C. R. Acad. Sci. Paris 243 (1956) 558-560
\[Hit\] Hitchin, N.: Harmonic spinors, Adv. Math. 14 (1974) 1-55
\[HH\] Hirzebruch, F., Hopf, H.: Felder von Flächenelementen in 4-dimensionalen Mannigfaltigkeiten, Math. Ann. 136 (1958) 156-172
\[Hur\] Hurwitz, A.: Über die Komposition der quadratischen Formen, Math. Ann. 88 (1923) 1-25 (posthumous)
\[JW\] Jordan, P., Wigner, E.: Über das Paulische Äquivalenzverbot, Zeit. f. Physik 47 (1927) 631-651
\[Kae\] Kaehler, E.: Der innere Differentialkalkül, Rend. Mat. 21 (1962) 425-523
\[Krb\] Karoubi, M.: $`K`$-Theory, Springer, New York - Heidelberg - Berlin, 1974
\[Kar1\] Karrer, G.: Einführung von Spinoren auf Riemannschen Mannigfaltigkeiten, Ann. Acad. Scient. Fennicae Ser. A. I. Mathematica 336/5 (1963)
\[Kar2\] Karrer, G.: Darstellung von Cliffordbündeln, Ann. Acad. Scient. Fennicae Ser. A. I. Mathematica 521 (1973)
\[LM\] Lawson, H.B., Michelsohn, M.L.: Spin Geometry, Princeton Univ. Press, Princeton, 1989
\[Lic\] Lichnerowicz, A.: Spineurs harmoniques, C. R. Acad. Sci. Paris A 257 (1963) 7-9
\[Lip\] Lipschitz, R.: Untersuchungen ueber die Summen von Quadraten, Max Cohen & S., Bonn, 1886 (French extract in Bull. Sci. Math. 2 Sér. 10 (1886) 163-183)
\[Mil\] Milnor, J.: Spin structures on manifolds, l’Ens. Math. 9 (1963) 198-203
\[Mor\] Morgan, J.: The Seiberg-Witten Equations and Applictions to the Topology of Smooth 4-Manifolds, Princeton Univ. Press, Princeton, 1996
\[Pau\] Pauli, W.: Zur Quantenmechanik des magnetischen Elektrons, Zeit. f. Physik 43 (1927) 601-623
\[Ply\] Plymen, R.: Strong Morita equivalence, spinors and symplectic spinors, J. Operator Theory 16 (1986) 305-324
\[Rad\] Radon J.: Lineare Scharen orthogonaler Matrizen, Abh. Math. Seminar Hamburg 1 (1922) 1-14
\[Ren\] Rennie, A.: Commutative geometries are spin manifolds, preprint, Univ. Adelaide, http://xxx.lanl.gov/math-ph/9903021
\[RS\] Rosenberg, J.M., Stolz, S.: Manifolds of positive scalar curvature, Algebraic Topology and its Applications, 241-267, MSRI Publ. 27, Springer, New York, 1994
\[Sdr1\] Schröder, H.: Funktionalanalysis, Verlag Harri Deutsch, Frankfurt a.M., 2000
\[Sdr2\] Schröder, H.: Globale Analysis, Textbook (manuscript), Univ. Dortmund, 2000
\[Sch\] Schrödinger, E.: Diracsches Elektron im Schwerefeld I., Sitz.-ber. Preuss. Akad. Wiss. Berlin, Phys.-Math. Kl. XI (1932) 105-128
\[Sto\] Stolz, S.: Simply connected manifolds of positive scalar curvature, Bull. Amer. Math. Soc. 23 (1990) 427-432
\[Var\] Várilly, J.C.: An Introduction to Noncommutative Geometry, Lecture Notes EMS Summer School on Noncommutative Geom. and Appl., Monsaraz and Lisabon, 1997, http://xxx.lanl.gov/physics/9709045
\[vdW1\] Waerden, B. L. van der: Exclusion principle and spin, in: Theoretical Physics in the Twentieth Century, (eds. M. Fierz, V.F. Weisskopf), Interscience Publ., New York, 1960, pp. 199-244
\[vdW2\] Waerden, B. L. van der: On Clifford algebras, Indag. Math. 28 (1966) 78-83
\[Wei\] Weil, A.: Correspondence, Ann. of Math. 69 (1959) 247-251
\[Wey1\] Weyl, H.: The Theory of Groups and Quantum Mechanics, Dover Publ., New York, (transl. by H.P. Robertson from Gruppentheorie und Quantenmechanik, Hirzel, Leipzig, 1931)
\[Wey2\] Weyl, H.: The classical groups, Princeton Univ. Press, Princeton, 1949<sup>2</sup>
Address:
Herbert Schröder
Fachbereich Mathematik
Universität Dortmund
Postfach 50 05 00
D-44221 Dortmund
e-mail: schroed@math.uni-dortmund.de
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# Polyelectrolyte Titration: Theory and Experiment
## I Introduction
In recent years the search for environment-friendly materials has promoted the development of numerous water-soluble polymer applications. Polymers can be made water soluble by making them compatible with the strong polar environment of the aqueous media, e.g., by introducing either charges or strong dipoles on the chains .
Most charged polymers (polyelectrolytes) have a hydrophobic backbone. This hydrophobicity induces an effective attraction between monomers which competes with the Coulomb repulsion between charges. As a result, hydrophobic polyelectrolytes exhibit complex behavior including conformation changes, macro- and meso-phase separation, self-association and aggregation .
In spite of considerable theoretical and experimental effort, many questions remain open. In particular, the role of the coupling between hydrophobic attractions and long range Coulomb interactions on the physico-chemical properties of polyelectrolyte solutions is not fully understood. In this paper, we consider the role of such coupling in solutions of weak polyacids.
Weak acid monomers (denoted HA) can undergo dissociation of the type
$$\mathrm{HA}\mathrm{H}^++\mathrm{A}^{}$$
where A<sup>-</sup> is the charged monomer attached to the backbone, while the dissociated H<sup>+</sup> charge dissolves into the solution. The dissociation/association is an equilibrium process satisfying detailed balance. The ionization degree determines the effective amount of charge on the polyelectrolyte chain, and it depends on the pH of the solution. At low pH the polymer is weakly charged while at high pH a larger fraction of monomers is dissociated and the polymer charge saturates to its maximal value. The most visible consequence is the solubility in water: hydrophobic polyacids can become water soluble at high enough pH, where the polymer charge is strong enough to overcome the chain hydrophobicity. In contrast to low molecular weight acids, the charged groups of polyacids are correlated as they are linked together along the chain. Indeed, the dissociation of one acid group is correlated in a complex way to the position and number of other charged groups on the chain. As a result, when the amount of charge on the chain varies, the chain conformation is affected, and in turn, influences the dissociation of other groups . Studies of the interplay of these competing effects has attracted a large amount of experimental and theoretical interest. Additional motivation for these efforts is related to the use of water-soluble polyelectrolytes and in particular alkali-swellable polymers in many industrial applications, such as coatings, food and cosmetic industries .
An elementary and wide spread experimental tool to characterize polyacids is to perform titration experiments. In these experiments, a strong base like NaOH is added to a solution of weakly charged polyelectrolytes. The pH of the solution and the equilibrium dissociation depends not only the amount of added base (like for ordinary acids ), but also on the polyelectrolyte concentration and presence of salt . The overall shape of the titration curve depends on the nature of the polyelectrolyte which is titrated . In the case of poly(acrylic) acid, the pH increases steadily with the dissociation degree. In contrast, the titration curve of poly(methacrylic) acid shows a maximum at low dissociation. Such a non-monotonous dependence has been related to a conformational transition. The case of hydrophobic polyelectrolytes is still more intriguing: the titration curve shows a large plateau where the pH is almost constant before the neutralization point is obtained. It has been argued in the literature that this plateau could be associated with a first-order phase transition between collapsed and swollen states of the polyacid chains .
In this paper we show that this peculiar shape of the titration curves can be explained without the need to rely on conformational phase transition of the polyacid chains. By including the effect of correlations of charges along the chains in the free energy, titration curves are calculated and show a behavior similar to the curves obtained in experiments. In particular, we explain how the pH depends on the polymer concentration and the amount of added salt. This dependence is special to polyacid solutions and is much weaker in monomeric acid solutions.
The paper is structured in the following way: in the next section the mean-field free energy, its one-loop correction and the resulting titration equations and curves are presented. Then, the experimental measurements of the titration of methacrylic acid / ethyl-acrylate copolymers (MAA-EA) are discussed in Sec. III, and the comparison between theory and experiment is presented in Sec. IV. The full formalism relying on a field theoretical approach will be presented in a forthcoming publication.
## II Theory
We present first the free energy leading to titration curves of weak polyacids in presence of added salt. Our system consists of four dissociating species: water, methacrylic acid monomers (denoted HA), NaOH titrating base and NaCl salt. The dissociation reactions are written as:
$`\mathrm{H}_2\mathrm{O}\mathrm{H}^++\mathrm{OH}^{},`$ (1)
$`\mathrm{HA}\mathrm{H}^++\mathrm{A}^{},`$ (2)
$`\mathrm{NaOH}\mathrm{Na}^++\mathrm{OH}^{},`$ (3)
$`\mathrm{NaCl}\mathrm{Na}^++\mathrm{Cl}^{}.`$ (4)
The partial dissociation of water and acid monomers is accompanied by an energy cost of breaking the molecular bond. They are denoted as $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_2`$, respectively (in units of the thermal energy $`k_\mathrm{B}T`$), and they are related to the mass action law as will be detailed below. Since NaOH is a strong base it is fully dissociated. Similarly, the salt is fully dissociated.
The polymer used in the experiment is a statistical copolymer composed randomly of two monomers: methacrylic acid and ethyl-acrylate. The total polymerization index is denoted as $`N`$. A fraction $`f=1/3`$ of the monomers are composed of the methacrylic acid, i.e. they are ionizable. Since the charged monomers are distributed uniformly along the chain, we assign a partial charge $`fe`$ to each monomer (“smearing” the charges along the chain). Note that $`fe`$ is the nominal charge of each monomer, while the actual charge is related to the partial dissociation of the acid monomers.
The concentration of monomers in the solution is denoted $`c_\mathrm{m}`$ while the concentrations of the small ions are denoted $`c_\mathrm{H}`$, $`c_{\mathrm{OH}}`$, $`c_{\mathrm{Na}}`$, $`c_{\mathrm{Cl}}`$.
The base concentration added to the solution is denoted $`c_\mathrm{B}`$ and its ratio to the MAA monomer concentration is defined as the degree of neutralization $`\gamma `$
$$\gamma \frac{c_\mathrm{B}}{fc_\mathrm{m}}$$
(5)
where $`\gamma `$ varies between zero (no added base) and infinity (large base excess). It is also useful to define the degree of ionization $`\alpha `$ of the acid
$$\alpha =\frac{[\mathrm{A}^{}]}{[\mathrm{HA}]+[\mathrm{A}^{}]}=\frac{[\mathrm{A}^{}]}{fc_\mathrm{m}}$$
(6)
where $`[\mathrm{A}^{}]`$ and $`[\mathrm{HA}]`$ are the concentrations of dissociated and non-dissociated monomers, respectively. This degree of ionization is related to the pH=$`\mathrm{log}_{10}c_\mathrm{H}`$ of the solution via
$$\mathrm{pH}=\mathrm{pK}_\mathrm{A}+\mathrm{log}_{10}\frac{\alpha }{1\alpha }$$
(7)
In addition, due to charge neutrality, $`\alpha `$ and $`\gamma `$ are related via
$$\alpha =\gamma +\frac{[\mathrm{H}^+][\mathrm{OH}^{}]}{fc_\mathrm{m}}$$
(8)
The conservation of mass implies the following relation between the ion concentrations:
$$c_{\mathrm{Na}}=c_\mathrm{B}+c_{\mathrm{Cl}}$$
Note that the Cl<sup>-</sup> ions come only from the salt, while the Na<sup>+</sup> come both from the base and the salt. Similarly, the OH<sup>-</sup> ions come from the dissociation of water and base, whereas the H<sup>+</sup> ions come from the dissociation of water and monomeric acid.
### A Free energy
We denote the full free energy as $`F=F_0+\mathrm{\Delta }F`$, where $`F_0`$ is the mean field term and $`\mathrm{\Delta }F`$ is the one-loop correction. Below we discuss the two terms of the free energy separately.
#### 1 Mean-field free energy
In a separate work, the mean-field free energy (per unit volume and in units of $`k_\mathrm{B}T`$) is shown to be
$`\beta F_0={\displaystyle \frac{c_\mathrm{m}}{N}}\left(\mathrm{log}c_\mathrm{m}w_\mathrm{m}1\right)+{\displaystyle \underset{j=\mathrm{H},\mathrm{OH},\mathrm{Na},\mathrm{Cl}}{}}c_j\left(\mathrm{log}c_jw_j1\right){\displaystyle \frac{v}{2}}c_\mathrm{m}^2+{\displaystyle \frac{w}{6}}c_\mathrm{m}^3\lambda _0c_\mathrm{m}`$ (9)
(10)
$`\left(c_{\mathrm{OH}}c_\mathrm{B}\right)\mathrm{\Delta }_1\left(c_\mathrm{H}+c_\mathrm{B}c_{\mathrm{OH}}\right)\mathrm{\Delta }_2+\left(c_\mathrm{H}+c_\mathrm{B}c_{\mathrm{OH}}\right)\beta e\phi _0fc_\mathrm{m}\mathrm{log}(1+\mathrm{e}^{\beta e\phi _0})`$ (11)
It turns out to be identical to the usual Flory-Huggins theory of polymer-solvent mixtures. The factor $`1/N`$ in the first term accounts for the reduction of translational entropy of chains of $`N`$ monomers. The second term represents the entropy of mixing of the small ions. The parameter $`w_j`$ is the molar volume of the $`j`$ species. The next two terms represent the short range monomer-monomer interactions. The hydrophobic effect is modeled by an attractive second virial coefficient $`v`$ (having dimension of volume), and a repulsive third virial coefficient $`w`$ (having dimension equal to volume squared). It is introduced to avoid collapse of the chain. The next term represents the chain conformational entropy, where $`\lambda _0`$ is the conformational entropy per monomer of a free Brownian chain.
The following two terms represent the energy cost of dissociation of water molecules and acid monomers. We define by $`\mathrm{\Delta }_1`$ the energy loss for each dissociation of a water molecule
$$\mathrm{H}_2\mathrm{O}\mathrm{H}^++\mathrm{OH}^{}$$
and by $`\mathrm{\Delta }_2`$ the energy loss for each dissociation of an acid molecule
$$\mathrm{HA}\mathrm{H}^++\mathrm{A}^{}$$
The number of OH<sup>-</sup> ions coming from dissociation of water molecules is equal to the difference between the total number of OH<sup>-</sup> ions and those coming from the NaOH base. Similarly, the number of dissociated acid groups is equal to the difference between the total number of H<sup>+</sup> ions and the number of H<sup>+</sup> ions coming from the water.
Finally, the last two terms account for the electrostatic energy, $`\phi _0`$ being the electrostatic potential in the solution. The first one is the electrostatic energy of the small ions, while the second is the electrostatic free energy of partially dissociated monomers .
#### 2 One-loop correction to the free energy
The one-loop correction to the mean-field free energy $`\mathrm{\Delta }F`$ will be presented in detail in a forthcoming publication. It is obtained by integrating over the quadratic fluctuations of the concentration fields. The correction term to the free energy is given (up to an additive constant) by
$`\beta \mathrm{\Delta }F`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{d}^3𝐪}{(2\pi )^3}\mathrm{log}\mathrm{\Sigma }(q)}`$ (12)
$`\mathrm{\Sigma }(q)`$ $`=`$ $`\left[1+(wc_\mathrm{m}v)c_\mathrm{m}ND(\eta )\right]\left(q^2+\kappa ^2\right)+4\pi l_\mathrm{B}f^2A^2ND(\eta )c_\mathrm{m}`$ (13)
$`\eta `$ $`=`$ $`{\displaystyle \frac{1}{6}}a^2q^2N`$ (14)
where the Debye-Hückel screening length $`\kappa ^1`$ depends on the total concentration of small ions $`c_\mathrm{I}`$ and the effective charge of the chain $`fA(1A)c_\mathrm{m}`$
$`\kappa ^2`$ $`=`$ $`4\pi l_\mathrm{B}\left(c_\mathrm{I}+fA(1A)c_\mathrm{m}\right)4\pi l_\mathrm{B}c_{\mathrm{eff}}`$ (15)
$`c_\mathrm{I}`$ $`=`$ $`c_\mathrm{H}+c_{\mathrm{OH}}+c_\mathrm{B}+2c_{\mathrm{Cl}}`$ (16)
$`l_\mathrm{B}=e^2/\epsilon k_\mathrm{B}T7`$ A is the Bjerrum length. The dissociation fraction of the acid monomers is given by
$`A`$ $`=`$ $`{\displaystyle \frac{\mathrm{e}^{\beta e\phi _0}}{1+\mathrm{e}^{\beta e\phi _0}}}`$ (17)
The Debye function $`D(\eta )`$ entering eq. (13) is given by
$$D(\eta )=\frac{2}{\eta }(1+\frac{\mathrm{e}^\eta 1}{\eta })$$
(18)
In addition to the small ion contribution to the electrostatic screening, eq. (15) includes a term proportional to $`A(1A)`$, where $`A`$ is defined in eq. (17) above. This term accounts for changes in the dissociation degree of monomers depending on the local electrostatic potential. At small $`A`$ values, $`A(1A)A`$, and the polymer contribution is the same as that of disconnected monomers. However, as $`A1`$ there is a substantial reduction (by the factor $`1A`$) of the effective polymer charge that contributes to screening . This reduction can be understood in terms of charge correlation along the chain. Dissociated monomers prevent further dissociation of other monomers.
In our model, the concentrations $`c_{\mathrm{Na}}`$, $`c_{\mathrm{Cl}}`$ and $`c_\mathrm{m}`$ are fixed by the amount of NaCl salt, NaOH base and polymer in the solution while the concentrations of dissociated H<sup>+</sup> and OH<sup>-</sup> ions depend on the degree of dissociation of water molecules and acid monomers (through the electrostatic potential $`\phi _0`$). Thus $`c_\mathrm{H}`$, $`c_{\mathrm{OH}}`$ and $`\phi _0`$ are variational parameters, determined by the requirement that the free energy $`F=F_0+\mathrm{\Delta }F`$ is an extremum. In titration experiments $`c_\mathrm{H}`$ can be monitored directly through the pH of the solution.
#### 3 Simplified free energy
In order to better understand the main contributions to the correction $`\mathrm{\Delta }F`$ let us return to the expression for $`\mathrm{\Sigma }(q)`$, eq. (13). First, we note that the Debye function is bound between 0 and 1 (see also eq. (34)). For typical values of the physical parameters (see also Sec. IV),
$`v25\text{A}^3`$ $`w100\text{A}^6`$ (19)
$`N10^4`$ $`c_\mathrm{m}1\text{mM}`$ (20)
it is easily seen that for long chains ($`N1`$), the Debye function can be approximated by $`ND(\eta )12/a^2q^2`$, and
$$|wc_\mathrm{m}v|c_\mathrm{m}ND(\eta )1$$
(21)
Thus, the effect of the solvent is small and will be neglected in this section
$$\mathrm{\Delta }F\frac{1}{2}\frac{\mathrm{d}^3𝐪}{(2\pi )^3}\mathrm{log}\left[q^2+\kappa ^2+z\right]$$
(22)
where
$$z=4\pi l_\mathrm{B}f^2A^2ND(\eta )c_\mathrm{m}$$
(23)
For the same range of parameters $`\kappa ^2z`$ and $`\mathrm{\Delta }F`$ can be expanded to first order in $`z`$, $`\mathrm{\Delta }F=\mathrm{\Delta }F_0+\mathrm{\Delta }F_1`$:
$$\mathrm{\Delta }F_0=\frac{1}{2}\frac{\mathrm{d}^3𝐪}{(2\pi )^3}\mathrm{log}\left[q^2+\kappa ^2\right]$$
(24)
The integral can be calculated by including a cut-off at large $`q`$, which cancels out the Coulomb self-energy, yielding
$$\mathrm{\Delta }F_0=\frac{\kappa ^3}{12\pi }$$
(25)
This result is the well known Debye-Hückel correlation energy of an electrolyte, because the polymer contribution is negligible and does not appear in this leading term. The correction $`\mathrm{\Delta }F_1`$ is given by
$$\mathrm{\Delta }F_1\frac{1}{2}\frac{\mathrm{d}^3𝐪}{(2\pi )^3}\frac{z}{q^2+\kappa ^2}6f^2A^2\frac{l_\mathrm{B}\kappa ^1}{a^2}c_\mathrm{m}$$
(26)
This correction is due to correlations along the chain between dissociated monomers as can be seen from the following simple argument. Consider the screened electrostatic interaction between charged monomers on a single infinite chain. For a specific chain configuration the Coulomb energy per monomer is
$$U_{\mathrm{el}}=\underset{j0}{}\frac{f^2A^2e^2}{\epsilon |𝐫_j|}\mathrm{e}^{\kappa r_j}$$
(27)
where $`r_j`$ is the spatial distance between the $`j=0`$ monomer and another $`j=\pm 1,\pm 2,\mathrm{}`$ monomer. For a Gaussian random walk, the typical (most probable) distance between monomers is given by
$$r_j=\left(\frac{2j}{3}\right)^{1/2}a$$
(28)
Assuming that $`r_j`$ can be approximated by $`r_j`$ in eq. (27), and $`\kappa a1`$, the sum over $`j`$ can be replaced by a continuous integral leading to eq. (26).
### B Titration equations
The free energy depends on the species concentrations: $`c_{\mathrm{OH}},c_{\mathrm{Na}},c_\mathrm{H},c_\mathrm{m},c_{\mathrm{Cl}}`$. However, the concentration of monomers, Na<sup>+</sup> and Cl<sup>-</sup>, is fixed by the amount of polymer, base and salt added to the solution. On the other hand, the concentration of H<sup>+</sup> and OH<sup>-</sup> ions is determined self-consistently by minimizing the full one-loop free energy $`F=F_0+\mathrm{\Delta }F`$, eqs. (9), (12), with respect to $`c_\mathrm{H},c_{\mathrm{OH}}`$ and $`\phi _0`$. The resulting equations of state determine the dependence of the pH of the solution on the other system parameters.
$`\mathrm{log}(c_\mathrm{H}\omega _\mathrm{H})+\beta e\phi _0\beta \mathrm{\Delta }_2+2\pi l_\mathrm{B}I_1=0`$ (29)
$`\mathrm{log}(c_{\mathrm{OH}}\omega _{\mathrm{OH}})\beta e\phi _0\beta \mathrm{\Delta }_1+\beta \mathrm{\Delta }_2+2\pi l_\mathrm{B}I_1=0`$ (30)
$`c_\mathrm{H}+c_\mathrm{B}c_{\mathrm{OH}}fc_\mathrm{m}A\left(12\pi l_\mathrm{B}(1A)\left((12A)I_1+2ANfI_2\right)\right)=0`$ (31)
where the quantities $`I_1`$ and $`I_2`$ are defined as:
$`I_1`$ $`=`$ $`{\displaystyle _0^\mathrm{\Lambda }}{\displaystyle \frac{k^2dk}{2\pi ^2}}{\displaystyle \frac{1+Nc_\mathrm{m}(wc_\mathrm{m}v)D(\eta )}{\mathrm{\Sigma }(k)}}`$ (32)
$`I_2`$ $`=`$ $`{\displaystyle _0^\mathrm{\Lambda }}{\displaystyle \frac{k^2dk}{2\pi ^2}}{\displaystyle \frac{D(\eta )}{\mathrm{\Sigma }(k)}}`$ (33)
and $`\mathrm{\Lambda }`$ is a short distance (large $`q`$) cut-off.
With the full expression (18) of the Debye function $`D(\eta )`$, the integrals in (32), (33) cannot be performed analytically. However, for the values of the parameters used in the experiments, namely polymerization index $`N10^410^5`$ and monomer length $`a5`$A, it is possible to use a simplified form for the Debye function:
$$D(\eta )\frac{1}{1+\eta /2}$$
(34)
As can easily be seen by comparing eq. (34) with the exact expression (18), the above form has the right behavior both at small and large values of $`\eta `$. Within this approximation, the integrals are given by:
$`I_1`$ $`=`$ $`{\displaystyle \frac{1}{2\pi R^2}}{\displaystyle \frac{R^2\kappa ^2\kappa _+\kappa _{}/2+(\lambda +1)\kappa ^2+u}{\kappa _+\kappa _{}(\kappa _++\kappa _{})}}`$ (35)
$`I_2`$ $`=`$ $`{\displaystyle \frac{1}{2\pi R^2(\kappa _++\kappa _{})}}`$ (36)
where we have used the notation:
$`\lambda `$ $`=`$ $`Nc_\mathrm{m}(wc_\mathrm{m}v)`$ (37)
$`u`$ $`=`$ $`4\pi l_\mathrm{B}A^2Nf^2c_\mathrm{m}`$ (38)
$`R^2`$ $`=`$ $`Na^2/6`$ (39)
and
$$\kappa _\pm ^2=\frac{1}{R^2}\left(\lambda +1+\frac{R^2}{2}\kappa ^2\pm \sqrt{\left(\lambda +1R^2\kappa ^2/2\right)^22uR^2}\right)$$
(40)
Note that in eq. (35) we have omitted a term, equal to $`\mathrm{\Lambda }/2\pi ^2`$, which exactly cancels the Coulomb self-energy.
With these definitions, eqs. (29) and (30) can be recast in the form:
$$c_\mathrm{H}c_{\mathrm{OH}}=10^{14}\mathrm{e}^{4\pi l_\mathrm{B}I_1}$$
(41)
which shows the change induced by fluctuations on the water dissociation constant and
$$c_\mathrm{H}=\frac{1A}{A}10^{\mathrm{pK}_\mathrm{A}}\mathrm{e}^{2\pi l_\mathrm{B}I_1}$$
(42)
These two equations, together with eq. (31), which expresses charge neutrality at the one-loop level, are solved iteratively. The numerical solution is obtained by using the mean field values as a starting point for the iterations, and convergence is usually achieved after a few iterations.
### C Structure Function — $`S(q)`$
In scattering experiments the structure function $`S(q)`$ is readily obtained. It is related to the Fourier transform of the various density-density correlations in the sample. For example, we can regard the monomer-monomer correlations:
$$S(q)=\delta c_\mathrm{m}(q)\delta c_\mathrm{m}(q)/c_\mathrm{m}$$
(43)
Since the one-loop expansion takes into account Gaussian fluctuations, it can also be used to calculate the structure function. The result is
$$S^1(q)=\frac{1}{ND(\eta )}+wc_\mathrm{m}^2vc_\mathrm{m}+\frac{f^2A^24\pi l_\mathrm{B}c_\mathrm{m}}{q^2+\kappa ^2}$$
(44)
The inverse structure function $`S^1(q)`$ is the energy penalty associated with density fluctuations at a wavenumber $`q`$. A minimum in $`S^1(q)`$ corresponds to the strongest fluctuating wavenumber $`q^{}`$. As a result, incoming radiation at this wavenumber interacts most strongly with the sample and a peak appears in the structure function and consequently in the scattering intensity.
An instability appears when $`S^1(q)`$ vanishes, corresponding to a divergence of the peak in the structure function. When the instability appears at a finite $`q=q^{}`$ the system undergoes a meso-phase separation and becomes spatially modulated. If, however, the instability is at $`q=0`$ the system undergoes a macrophase separation. In our case this occurs when
$$vc_\mathrm{m}=wc_\mathrm{m}^2+\frac{1}{N}+\frac{f^2A^2c_\mathrm{m}}{c_{\mathrm{eff}}}$$
(45)
Recall that $`c_{\mathrm{eff}}`$ and $`A`$ depend on the pH and the degree of neutralization.
## III Experiment
### A Preparation of samples
The copolymer polyelectrolyte chains used in this study are prepared by standard emulsion polymerization techniques using neutral ethyl-acrylate (EA) monomers and methacrylic acid (MAA) monomers. The methacrylic acid is a weak acid with pK<sub>A</sub>=4.5. The weight fraction of MAA in the copolymer is equal to 0.35 taken to be $`f=1/3`$ in the theoretical section. The emulsion polymerization is performed under starved monomer conditions. Under these conditions, the MAA and EA monomers are evenly distributed along the polymer chain. From size exclusion chromatography, we have estimated the molecular weight of the chains to be of the order of 106 daltons, corresponding to a polymerization index $`N10^4`$. The MAA-EA copolymers precipitate from the aqueous solution, as they are insoluble in water. The precipitate is carefully washed by ultrafiltration in order to remove surfactants, unreacted monomers and initiators. This cleaning procedure is stopped when the resistivity of the water flushed through the separation membrane is that of pure water (18.2 M$`\mathrm{\Omega }`$/m). The solid content of the stock solution is then determined accurately by drying and weighing.
### B Titration measurements
Polymer solutions at different weight concentrations are prepared by mixing weighted amounts of the stock solutions with de-ionized water. The dissolution of CO<sub>2</sub> is prevented by carefully de-gassing the solution with nitrogen. Each of the polyelectrolyte solutions is neutralized by an NaOH solution with a molar concentration ranging from 0.1M to 2M depending on the polymer concentration. Upon neutralization, the methacrylic acid is neutralized and repulsive forces due to the negative charges cause the chain to expand, resulting in the progressive solubilization of the polymer chains. As a result the solution becomes transparent and its viscosity increases. In the following, we shall characterize the neutralization of the polymer chains by the degree of neutralization, $`\gamma =c_\mathrm{B}/(fc_\mathrm{m})`$, which is the ratio of the amount of added base to the amount of available acids groups. Titration experiments are performed using a pH-meter (Metrohm 691) with a combined glass electrode. The measurements are made at 20<sup>o</sup>C with constant stirring under a nitrogen atmosphere. In parallel, the conductivity of the solution is measured (Metrohm 712).
The measurements reveal that all small ions present in solution are free and contribute to the conductivity. This was checked at different $`\gamma `$ values by changing the polymer concentration. The measured conductivity is the exact sum of the conductivities coming from the Na<sup>+</sup> ions (whose concentration is known from the amount of added NaOH base), and H<sup>+</sup>, OH<sup>-</sup> (known from the pH), while the small contribution of the polymer is negligible. These conductivity measurements demonstrate that the distribution of small ions is homogeneous in the solution, and no evidence for the Donnan effect and counter-ion condensation is observed.
## IV Results: comparison of experiments with theory
In this section we present the experimental results for the titration curves of MAA-EA copolymers and compare them with the theory of Sec. II. Fig. 1(a) shows the titration curves measured for different polymer concentrations. They differ substantially from the titration curves of monomeric methacrylic acid (MAA). Methacrylic acid in aqueous solution has the typical behavior of a weak monomeric acid. The pH increases monotonously with $`\gamma `$, takes the value pK$`{}_{\mathrm{A}}{}^{}=4.5`$ for $`\gamma `$=0.5 and then jumps near $`\gamma `$=1. By contrast, the titration curves of MAA-EA copolymers exhibit the following behavior: in the range $`0\gamma 0.2`$, the pH increases sharply to a plateau value which remains nearly constant up to $`\gamma 1`$ where the amount of added base equals that of the monomer (neutralization point). The jump of the pH for $`\gamma 0.2`$ is associated with a swelling phase transition of the latex particles i.e. a change of conformation of the chains. A similar phenomenon has been observed during the titration of pure poly(methacrylic acid) .
Let us mention that it is possible to calculate the swelling transition of the polymer as function of monomer concentration, $`c\mathrm{m}`$. This can be done by minimizing the free energy with respect to $`c_\mathrm{m}`$ (in addition to the other annealed degrees of freedom discussed above). Any non-convexity of the free energy signals the existence of a polymer precipitate in excess water.
On Fig. 1 several titration curves are plotted for different concentrations of the same MAA-EA copolymer. As the concentration $`c_\mathrm{m}`$ increases, the jump in the pH at $`\gamma 1`$ increases. Note that the deviation between the pH of the MAA monomeric acid and the MAA-EA copolymers ($`c_\mathrm{m}=0.1`$ M in Fig. 1), is large in the plateau region ($`\gamma <1`$). On the other hand, the deviation is quite small for $`\gamma >1`$ when the acid is almost completely dissociated. The value of the pH at the plateau depends strongly on the polymer concentration: the larger the $`c_m`$, the lower the plateau. It is surprising that even though MAA-EA copolymers contain carboxylic groups, the plateau value of the pH below $`\gamma =1`$ may be neutral or even greater than 7.
This difference in behavior between the MAA monomers and MAA-EA copolymers might be associated with a complex structure of the MAA-EA copolymer for $`\gamma 1`$. Collapsed microdomains may exist on the chains due to the competition between the hydrophobic attraction and Coulomb repulsion. The plateau region suggests the existence of such collapsed microdomains along the chains. For $`\gamma 1`$ the chains are in a swollen state and their behavior resembles that of monomeric weak acid.
In Fig. 1(b) we plot for comparison the titration curves as calculated from the theory. The titration equations presented in Sec. II.B are solved numerically by an iteration procedure starting from the mean field values as the first iteration and including the corrections of the second iteration. Since there are a couple of unknown physical parameters (like the second and third virial coefficients, $`v,w`$ and the monomer size $`a`$), we do not try to fit the experimental titration curve. Rather, we note that the corresponding theoretical curves look very similar to the experimental ones, demonstrating the same type of plateau for low $`\gamma `$ and the same trend with the monomer concentration.
Figure 2 shows the effect of adding a monovalent salt (NaCl). The overall shape of the titration curves remain unchanged. However, the value of the pH at the plateau increases upon the addition of salt, while the pH for $`\gamma 1`$ is essentially independent on the salt concentration. This indicates that neutralization of the carboxylic groups carried by the chains becomes easier as the ionic strength is increased. As the Coulomb interaction is screened by the salt, there is less correlation between charged groups along the chain when their distance is larger than the Debye length. Therefore, monomers separated by distances larger than the Debye length can dissociate independently from each other. For $`\gamma 1`$, almost all the charged groups on the chains are dissociated; thus, the salt has no effect on the pH. The comparison between experiment and theory is presented in part (a) and (b), where a good agreement can be observed. Note the larger deviation from the experimental results for higher salt concentration ($`c_{\mathrm{NaCl}}=0.1`$M).
In Fig. 3, the plateau value of the pH is plotted on a semi-log plot as function of the ionic strength $`c_\mathrm{I}=c_{\mathrm{Na}}+c_\mathrm{H}+c_{\mathrm{Cl}}+c_{\mathrm{OH}}`$. The plateau value of the pH is taken at $`\gamma =0.5`$. It is interesting to note that the different points taken at various salt and base concentrations collapse on a single curve. This result indicates that the pH is dominated by electrostatic effects and not by conformation changes of the chains. When the ionic strength increases, the fixed charges carried by the polymer are screened and the energy associated with the dissociation of a carboxylic group decreases.
This behavior can be characterized semi-quantitatively by noting that when $`\gamma 1/2`$ the pH is much higher than the pK<sub>A</sub> and, therefore, $`\gamma ^{(0)}1`$ and $`A1`$. Consequently, at the plateau $`\mathrm{\Delta }\gamma 1/2`$. To first order in $`10^{\mathrm{pK}_\mathrm{A}\mathrm{pH}}`$ the plateau value can be obtained from the simplified free energy, Sec. II.A.3 and reads:
$`\mathrm{pH}_{\mathrm{plateau}}\mathrm{pK}_\mathrm{A}+\mathrm{log}_{10}\left[\kappa l_\mathrm{B}+24\pi f^2c_m{\displaystyle \frac{l_\mathrm{B}^2\kappa ^3}{a^2}}+24f{\displaystyle \frac{l_\mathrm{B}\kappa ^1}{a^2}}\right]`$ (46)
Recall that $`\kappa `$ depends on the total amount of ions in the system through eq. (15). The above expression is used for the plot of the solid curve of Fig. 3b, where the pH of the plateau is shown as a function of the effective ionic strength $`c_{\mathrm{eff}}`$ defined in eq. 15. In this regime, the difference between $`c_{\mathrm{eff}}`$ and $`c_\mathrm{I}`$ is a small correction of order $`10^{\mathrm{pK}_\mathrm{A}\mathrm{pH}}`$.
Equation 46 exhibits three different regimes depending on the relative importance of the various contributions to the logarithmic term. In each regime, $`c_\mathrm{H}`$ has a different power law dependence on the total ionic strength $`c_{\mathrm{eff}}=c_\mathrm{I}+fA(1A)c_m`$ leading to different slopes in Fig. 3. At low ionic strength the second term dominates and $`c_\mathrm{H}c_{\mathrm{eff}}^{3/2}`$ (long dashed line in Fig. 3b). At intermediate ionic strength the last term dominates and $`c_\mathrm{H}c_{\mathrm{eff}}^{1/2}`$ (short dashed line). Finally, at high ionic strength (beyond the experimentally accessible values) $`c_\mathrm{H}c_{\mathrm{eff}}^{1/2}`$.
In Figure 4 the effect of polymer concentration and ionic strength on $`S(q)`$ is shown. The structure function is calculated at the plateau regime were $`A1`$. The dependence on the controlled system parameters: the pH, the salt concentration and the degree of neutralization $`\gamma `$ is taken implicitly into account in $`c_{\mathrm{eff}}`$. At low values of $`c_{\mathrm{eff}}`$, the structure function exhibits a peak at finite $`q`$ corresponding to the most favorable density–density fluctuations. As $`c_{\mathrm{eff}}`$ increases, $`S(q0)`$ increases until the peak disappears. This increase is accompanied by an increase in the osmotic compressibility of the solution.
The effect of varying the polymer concentration $`c_\mathrm{m}`$ at a fixed $`c_{\mathrm{eff}}`$ is depicted in the inset. The peak is stronger at lower polymer concentrations indicating that fluctuations become considerably stronger at lower concentrations. This effect agrees well with the titration curves shown in Fig. 1. Indeed, at low polymer concentrations the shift in the titration curve is stronger than at high concentrations. The peak position shifts to higher wavenumbers (smaller length scale) at higher concentration, similar to the correlation length of polymer solutions that decreases with increasing polymer concentration.
## V Conclusions
We have presented an experimental study of the titration of weak polyacids by a strong base. Our experimental findings performed on solutions of MAA-EA copolymers, are supported by theoretical calculations, which show similar trends of the pH variation with the concentration of added base, salt and polyacid.
Titration experiments are one of the simplest and most useful experimental tools to probe the degree of neutralization of monomeric and polymeric acids. Since titration curves of weak acids are universal, any deviation from this behavior, as observed here, can be the signature of a complex behavior coupling Coulombic and hydrophobic interactions.
The most striking feature is the existence of a plateau of the pH for low degree of neutralization. The plateau in the pH is at higher value than the corresponding pH of the monomeric acid indicating that the charges on the chain inhibit the dissociation of other charged groups. The pH is not affected strongly by further addition of the base till the neutralization point $`\gamma =1`$. This is probably due to the existence of collapsed microdomains along the chains for $`\gamma 1`$. The pH at the plateau decreases as function of the ionic strength. This demonstrates that non-specific electrostatic interactions are responsible for the existence of the plateau, since the pH depends mostly on the amount of small ions and not on their type (for a fixed degree of neutralization).
The theory presented above includes one-loop corrections to the mean-field free energy. On the mean-field level, there is no difference between the polymeric and monomeric titration curves except for the translational entropy of the chains. The one-loop correction couples the chain connectivity with the electrostatics and induces the large deviations in the titration curves. We modeled the polyacid as flexible chains using the standard Debye function. Our formalism can also be applied to different models of chain elasticity. In particular, for the case of semi-flexible chains the same titration equations are obtained with a modified Debye function taking into account the persistence length of the chains.
Finally, it will be interesting to complement this study by scattering experiments where it might be possible to resolve the chain microstructures and relate them to the degree of ionization of the chains.
Acknowledgments:
IB gratefully acknowledges the support of the Chateaubriand postdoctoral fellowship and the hospitality of the Elf-Atochem research center at Levallois-Perret. DA acknowledges partial support from the US-Israel Binational Foundation (BSF) under grant No. 98-00429, and the Israel Science Foundation founded by the Israel Academy of Sciences and Humanities — centers of Excellence Program.
Borukhov et al: Fig. 1
Borukhov et al: Fig. 2
Borukhov et al: Fig. 3
Borukhov et al: Fig. 4
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# ASCA Observations of NLS1s
## 1 Introduction
ASCA observations have proved to be instrumental in our understanding of Narrow-line Seyfert 1 galaxies (NLS1s). One of the most compelling first indications that NLS1s may be characterized by a high accretion rate came from the ASCA spectrum of RE 1034+39, reported by Pounds, Done & Osborne (1995). This spectrum revealed a strong soft excess component and a steep hard X-ray power law that seemed to be similar to the spectra of black hole candidates in the high state. While it is now quite clear that not all NLS1s have such spectra and indeed a range of strengths of the soft excess are seen (e.g. Leighly 1999b), and that perhaps comparison with the very high state seen in some Galactic X-ray emitting objects may be more appropriate, the result has been extremely important in the development of our understanding of these objects.
The results from the ASCA observations of NLS1s have now been reported in several places (e.g. Leighly 1999ab and references therein, see also Vaughn et al. these proceedings). Therefore, only part of this review will be devoted to a few results from the variability analysis of NLS1s and a very brief discussion of models and implications. In the second half, I present a preview of our work on the HST spectra from NLS1s (Leighly & Halpern 2000).
## 2 X-ray Variability of NLS1s
The X-ray variability in NLS1 ASCA observations is discussed in detail in Leighly 1999a. Because of the gaps in the light curves, the excess variance, also known as the fractional amplitude of variability, was used to quantify the variability. The results are shown in Fig. 1 (left). This figure shows that for a given X-ray luminosity, the amplitude of variability is consistently larger for NLS1s than for Seyfert 1 galaxies with broad optical lines (BLS1s). The uniform sampling and other properties of the ASCA data mean that the excess variance should be proportional to the inverse of a time scale of variability $`T_i`$ as $`T_i^{1\alpha }`$, where $`\alpha `$ is the slope of the variability power spectrum (see also Lawrence & Papadakis 1993). Thus, the simplest interpretation of Fig. 1 is that, if the hard X-ray luminosity is characteristic of the absolute mass accretion rate (i.e. the efficiency of conversion of accretion energy to radiation is the same in all objects), then NLS1s have a shorter variability time scale than do Seyfert 1s with broad optical lines (BLS1s). This would be equivalent to compressing a BLS1 light curve into an interval approximately 10 times shorter in NLS1s. The variability time scale may be characteristic of the black hole mass, simply from light-crossing time-scale arguments. Thus, the simplest explanation of this result is that NLS1s have a smaller black hole mass than BLS1s, and since they are no fainter than BLS1s, they must be accreting at a relatively larger accretion rate.
Closer examination indicates that the situation may not be as simple as the argument above would imply. Specifically, Fig. 1 shows a large spread in the excess variance values for NLS1s, much larger than the measurement error, but also much larger than the model-dependent systematic error in $`\mathrm{log}`$(excess variance) of $``$0.31–0.47 arising from the weakly nonstationary nature of a $`1/f^\alpha `$ power spectrum (see Leighly 1999a for details). It is possible that the large spread is a consequence of the fact that the nature of the variability may not be homogeneous among NLS1s. Leighly 1999b report a correlation between the strength of the soft excess in the ASCA spectrum and the amplitude of the variability, and suggests that strong soft excess objects are characterized by flaring, high amplitude variability. In Fig. 1 (right), I separate the 6 objects identified as having strong soft excesses in Leighly 1999b and plot the excess variance as a function of 0.5–2.0 keV luminosity<sup>1</sup><sup>1</sup>1Note that this is the inferred intrinsic luminosity and therefore there may be some model dependence associated with the modeling of the absorption.. When these six objects are excluded, the slope of the regression at $`0.3`$ is no longer biased by the scatter. The slope is still somewhat flatter than that expected if the time scale is proportional to the luminosity assuming that the slope of the variability power spectrum is $`\alpha =1.5`$ (dashed lines)<sup>2</sup><sup>2</sup>2Lawrence & Papadakis (1993) found that the slopes of the power spectra on 1-day time scales in a sample of objects are consistent with a constant value of 1.55..
What is the origin of the enhanced variance observed in NLS1s? As stated above, it may be simply a consequence of a smaller emission region characteristic of a smaller black hole. However, the light curves in some strong soft excess NLS1s are characterized by very high amplitude flares, and they could not be simply time-compressed versions of BLS1 light curves. A few other emission mechanisms have been suggested that instead would be equivalent to stretching the amplitude of flares in BLS1 lightcurves to produce the enhanced excess variance observed in NLS1s. Naturally, because of the scale invariant nature of the variability power spectrum (at least in the range of frequencies that the ASCA data probe) these general scenarios cannot be distinguished using the power spectrum or the excess variance.
* Noting that the very large amplitude variability on short time scales implies a very high efficiency of conversion of accretion energy to radiation in the NLS1 PHL 1092, it has been proposed that beaming plays an important role in amplifying flares (Brandt et al. 1999). A specific scenario has been proposed: the emission regions are found on the very inner edge of the accretion disk, and the accretion disk is inclined at a very high angle with respect to the viewer so that Doppler effects cause amplification of flares (Boller et al. 1997). The results presented in the next section of this contribution cast significant doubt on an edge-on orientation.
* It has recently been suggested that if NLS1s are accreting at a higher rate then the magnetic field energy, assumed to be in equipartion with gravitational potential energy, should be proportionally larger (Mineshige et al. 2000). Large amplitude flares may result from reconnection of this more powerful magnetic field.
* It has been suggested that occultations by large optically-thick clouds may produce high-amplitude X-ray variability (Brandt et al. these proceedings). An observed steep $`\alpha _{ox}`$ in some objects seems to support this idea; however, there is no physical reason that luminous NLS1s should not have intrinsically steep $`\alpha _{ox}`$, for example, if the X-ray emitting corona is weak or not present.
What are the implications of the high amplitude flaring? It has been previously suggested that high amplitude flaring is evidence that the variability is nonlinear (Green, McHardy & Done 1999; Boller et al. 1997). This conclusion is based on the assumption that the AGN variability is Gaussian, since in that case a very high variance compared with the mean can only be produced if the variability is nonlinear. In fact, non-Gaussian variability makes more sense for AGNs because our usual physical picture supposes that the light curve should be built up from the superposition of flares, and flares are inherently non-Gaussian. Linear, non-Gaussian variability can easily produce high values of excess variance. Non-Gaussian variability can be detected using a parameter related to the skew of the flux distribution, and evidence for non-Gaussianity was found in several of the NLS1s with the highest amplitude of variability. Demonstration that a light curve is nonlinear is a much harder problem; it is a very important one, however, because of the potentially strong constraints on emission processes and geometry. See Leighly 1999a for a detailed discussion. We find some evidence for nonlinear variability in the broad-line radio galaxy 3C 390.3 (Leighly & O’Brien 1997) in that quiescent periods occur before and after large flares, and this behavior cannot be reproduced by a linear non-Gaussian model. Evidence for nonlinear variability has been recently reported in Cyg X-1 (Timmer et al. 2000).
## 3 HST Observations of NLS1s IRAS 13224$``$3809 and 1H 0707$``$495
In 1997, we reported the detection of an absorption feature near 1 keV in the ASCA spectra from three NLS1s, IRAS 13224$``$3809, 1H 0707$``$495, and PG 1404+226 (Leighly et al. 1997). Absorption by ionized oxygen is common in the X-ray spectra from Seyfert 1 galaxies; however, if we are to interpret the 1 keV features in this way, the absorbing material must have highly relativistic velocities (0.2–0.6$`c`$, depending on whether the features are interpreted as absorption lines or edges). Acceleration of ionized gas to these high velocities should be difficult; however, we noted that NLS1s exhibit suggestive similarities to a subclass of Broad-Absorption Line Quasars (BALQSOs), objects characterized by the signature of absorption by outflowing gas in the UV: both types of objects exhibit strong or extreme Fe II and weak \[O III\] emission, they often have red continua and strong infrared emission, and they are predominately radio-quiet (see Leighly et al. 1997 for details). An alternative explanation as absorption by highly ionized neon and iron has been presented by Nicastro et al. (1999). With luck, the origin of these features will be resolved by our upcoming Chandra observation of 1H 0707$``$495.
Absorption in X-rays is often accompanied by absorption in the UV. To test our hypothesis that the X-ray absorption features are produced by relativistically outflowing gas, we applied for and were awarded HST STIS UV spectroscopic observations of IRAS 13224$``$3809 and 1H 0707$``$495 (Leighly & Halpern 2000) and the results of the observations were presented for the first time at this meeting. The broad-band continuum spectra, corrected for Galactic reddening, are presented in Fig. 2 along with the average QSO spectrum compiled by Paul Francis. This figure shows that these NLS1s have continua as blue as the average QSO<sup>3</sup><sup>3</sup>3IRAS 13224$``$3809 has been previously reported as having a red spectrum (Mas-Hesse et al. 1994). We find that our optical spectrum obtained from a 1 hour exposure on the CTIO 4 meter telescope is indeed rather red and does not join smoothly to the UV spectrum. We infer, from a difference in the spatial profile of the lines and continuum, that there is a strong, red galaxy component in the spectrum. In contrast, our 1H 0707$``$495 optical spectrum is quite blue and joins smoothly onto the UV spectrum.. There is also no evidence for resonance-line absorption intrinsic to the object; all of the narrow absorption lines in the spectra originate in our Galaxy.
The emission line profiles are the most interesting feature of these spectra. The high-ionization lines, including Ly$`\alpha `$$`\lambda 1216`$, N V$`\lambda 1240`$, Si IV$`\lambda 1397`$ and C IV$`\lambda 1549`$ have a different profile than the low-ionization lines, including Mg II$`\lambda 2800`$ and H$`\beta `$$`\lambda 4861`$. The difference is illustrated in Fig. 3 which overlays the rescaled high-ionization line C IV and the low-ionization line Mg II profiles as a function of velocity<sup>4</sup><sup>4</sup>4Note that the profile of the doublet Mg II$`\lambda `$2796, $`\lambda `$2804 is consistent with that of H$`\beta `$.. As is commonly found in AGN (e.g. Marziani et al. 1997), the average quasar high-ionization lines are broader and blueshifted relative to the low-ionization lines. The same trend is true in NLS1s, except the difference is more profound. The C IV FWHM is $`4\times `$ that of H$`\beta `$, but it is significant that all of the extra width is on the blue side, whereas the red side lines up well with the low-ionization lines. Similar line profiles were reported from the NLS1 I Zw 1 (Laor et al. 1997).
We also find evidence for emission from gas with very high densities. In Fig. 4 we plot the Al III–Si III\]–C III\] line region of our spectra, as well as archived spectra from several other NLS1s. We find a very high Si III\] to C III\] ratio in our spectra from IRAS 13224$``$3809 and 1H 0707$``$495. The ions responsible for these lines are found under the same physical conditions; however, Si III\] has a higher critical density than C III\] and thus the ratio becomes large when the C III\] emission has saturated. A correlation between this ratio and the width of H$`\beta `$ has been previously reported by Wills et al. 1999 (see also Wills, these proceedings). We note that these emission lines are relatively narrow and symmetric about their rest wavelengths and thus their profiles are more similar to those of the low-ionization lines.
Fig. 4 compares line profiles from 6 NLS1s. Three of the objects (I Zw 1, IRAS 13224$``$3809 and 1H 0707$``$495) display blueshifted high-ionization lines and high Si III\] to C III\] ratios. RX J0134$``$42 also appears to have a blueshifted C IV line. However, the spectra from two lower luminosity NLS1s, Akn 564 and RE 1034+39, lack both of these attributes. We suspect that this reflects a dependence of the line properties on luminosity.
### 3.1 Disk-wind Models of AGN
We propose that the characteristic spectra we observe in the higher luminosity NLS1s can be naturally explained by a disk-wind model, as shown schematically in Fig. 5. It has long been known that it is not possible to produce all broad-line emission in gas characterized by one set of physical properties. Several multicomponent models, including disk-wind models, have been proposed; unfortunately, space limitations prohibit a through review here. What is new here is that our observations provide evidence that a disk-wind model is strongly favored over all other models. The fact that the red side of the high- and low-ionization lines line up suggests that the low-ionization lines and low-velocity cores of the high-ionization lines are produced in the same relatively low-velocity material (i.e. the accretion disk) while the high-velocity blue wings of the lines are produced in a wind. The disk is optically thick, and therefore, we see only emission from the wind accelerated toward us.
We are currently constructing a simple disk-wind model to test our assertion and Fig. 5 shows a schematic of the assumed geometry. We are following Murray & Chiang (1998) in general: we use a $`\beta `$ velocity law, often applied in star and CV winds, and we are using Cloudy to estimate the ionization structure. However, since we are interested in the emission produced in the wind in particular, we treat the radiative transfer numerically using the Sobolev approximation. A complementary and more sophisticated model being constructed by Proga, Stone and Kallman (2000; hereafter PSK) aims to determine the detailed dynamics of the wind. A definitive model has not yet been constructed; however, it is clear that there are potentially complex relationships between the radiative transfer, dynamics and ionization, and the results will quite possibly not conform to naive intuition when everything has been taken into account. Examples of two complications:
* The wind is quite likely to be optically thick to scattering of the resonance lines, and thus radiative transfer will be strongly enhanced along directions in which the velocity gradient is high. This effect can strongly alter the observed line profiles for the high-ionization lines emerging from the wind, and prompts the use of the Sobolev approximation. The He II recombination lines are immune to this effect, however.
* The acceleration of the wind is produced by resonance-line driving and thus the velocities attained depend intimately on the ionization structure of the gas. Specifically, there must be a mechanism for shielding the wind gas from the photoionizing X-ray source. This problem seems to have a natural solution: PSK find that there can be adequate “failed” wind material falling back onto the object to do this.
Our preliminary results and also the results presented by PSK are promising enough to conjecture that many of the features of the HST NLS1 spectra will be explainable using a disk-wind system. For example:
* A disk-wind model coupled with a high accretion rate may naturally explain NLS1 spectra. PSK find that a primary condition on formation of the wind is that the wind-driving (UV) flux should be large compared with the wind-photoionizing (X-ray) flux. A high accretion rate predicts a spectrum strongly peaked in the extreme UV and therefore the UV to X-ray ratio will be large enough that a strong wind can be driven. If the accretion rate is high, then densities in the disk will be high, resulting in a large Si III\] to C III\] ratio. More material will be available to be blown from the nucleus, and since the black hole mass is relatively smaller, the wind is more likely to reach escape velocity. The wind may partially shield the accretion disk from the intense radiation field, thus permitting excitation but not causing ionization of the Fe<sup>+</sup> ion, so that strong optical Fe II emission may be produced. This could explain the observed correlation between high-ionization blue asymmetry and optical Fe II (e.g. Marziani et al. 1996).
* The large UV to X-ray ratio requirement found by PSK may explain why lower luminosity NLS1s do not show blue wings on their high-ionization lines. Quasars are known to have typically relatively larger UV to X-ray ratios than Seyferts (Wilkes et al. 1994), and thus gas in Seyferts may become too highly ionized to be accelerated to high velocities.
* The observed line profiles will be affected by the observer’s viewing angle, and radiative transfer effects, as noted above, potentially affect the profiles as well in a non-obvious way. We interpret the fact that the red sides of the high- and low-ionizations lines line up as evidence that the low-ionization lines and the low-velocity core of the high-ionization lines are produced in the disk. However, for there to be only blue-side wind emission, either the observer’s viewing angle must be nearly face on, or the angle of the wind stream lines with respect to the disk must be large. It is quite possible that in reality the wind stream line angle and the radius at which the wind is launched are not independent of the accretion rate, and thus results from a self-consistent model would be potentially very interesting.
### 3.2 Summary and Further Implications
We have described how the particular UV spectra from higher luminosity NLS1s strongly suggests that the emission lines are produced in a disk-wind system. Further, we conjecture that, through a potentially complicated dependence of model parameters such as the amount of material and velocity and angle of stream lines on the physical parameters such as the black hole mass and accretion rate, it may be possible to explain Eigenvector 1 self-consistently. There are a few more important issues:
* Broad emission lines in AGN are often assumed to be produced in clouds that have virialized bulk motions (e.g. Peterson & Wandel 1999). In contrast, the strong winds that we propose are responsible for the blue wings in our spectra and thus for a significant fraction of the high-ionization line emission in luminous NLS1s will not have virialized bulk motion.
* It has been proposed that NLS1s bear some similarity to high redshift quasars (Mathur 2000): in particular, very high redshift quasars ($`z4`$) appear to have rather narrow lines compared with local quasars (Shields et al. 1997). The spectra presented here show that this conjecture may have serious weaknesses. At high redshift, only the rest-frame UV high-ionization lines are observed in the optical spectrum. We have discovered that the high-ionization lines in higher-luminosity NLS1s are typically broad; furthermore, Wills et al. (these proceedings) find no correlation between the high-ionization and H$`\beta `$ line widths in the PG-quasar subsample. These results imply that the observed optical spectra of high-redshift NLS1s will show broad emission lines, and therefore, such objects may be difficult to identify. Finally, if there is a wind, it is possible that some of the N V emission is a result of excitation of high velocity N<sup>+4</sup> ions by absorption of Ly$`\alpha `$ photons, and abundance enhancements may be less necessary (however, see also Hamann & Korista 1996).
* Will the wind be homogeneous? Resonance-line driven winds are a characteristic of Wolf-Rayet stars and there is very strong evidence for density enhancements in these winds. Furthermore, PSK find density inhomogeneities in their simulations. Clumping may be necessary to match the intensity of observed emission lines. If present, clumping could produce interesting time-dependent and ionization-dependent behavior.
Acknowledgements: I would like to thank the Wilhelm and Else Heraeus Foundation for travel support, Jules Halpern for helpful discussions, and Daniel Proga for an early look at his submitted paper. I gratefully acknowledge support through NAG5-7971 (NASA LTSA).
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# On two continuous models for the dynamics of sandpile surfaces
## I Introduction
Recent much interest to the physics of granular media was, in particular, stimulated by two salient features of the granular state: multiplicity of metastable pile shapes and occurrence of avalanches upon pile surfaces. It has been realized that, to account for metastability, the model of pile surface dynamics should not be written as an evolutionary equation for the pile surface alone. An additional unknown characterizing the flow of grains down the pile surface is useful because such flows are not uniquely determined by the external source and local free surface topography.
A large spatio-temporal scale pile growth model involving two coupled dependent variables and able to account for metastability has been proposed in . This model neglects avalanches as small fluctuations of the pile surface and describes the evolving mean surface of a pile that grows on an arbitrary support under a given distributed source of bulk material. The model permits an equivalent formulation as an evolutionary variational or quasivariational inequality; such a formulation simplifies significantly both the mathematical study of the problem and its numerical solution . As it has been shown in , the shapes of real piles on flat open platforms are described by the analytical solutions of this inequality. A modification of the model, able to account for avalanches as almost instantaneous slides, is also discussed in ; according to observations made in the same work, such a slide may, indeed, be a possible avalanche scenario (see also ).
Independently and using different arguments, the same pile growth model in the form of a variational inequality has been derived by Aronsson, Evans, and Wu . In , Evans et al. studied its discontinuous solutions corresponding to avalanches; in Evans and Rezakhanlou showed that the cellular automata models of sandpiles, presented as intuitively attractive examples in almost all works on self-organized criticality , converge in a continuous limit to a similar variational inequality with an anisotropy inherited from the cellular structure of these crude models.
A different continuous model, also involving two coupled dependent variables and describing the granular surface flow and pile surface dynamics, has been proposed by Bouchaud, Cates, Ravi Prakash, and Edwards (the BCRE model) . Although the choice of the basic variables in this model is equivalent to that in , the model is written for the free surfaces only slightly deviating from the critical slope and employs different phenomenological constitutive relations. The emphasis is put onto the simulation of fast processes, like amplification and distinction of rolling grains population during an avalanche. The BCRE model has been simplified by de Gennes , applied to various one-dimensional surface flow problems (see, e.g., ), and modified for thick surface granular flows . Further exact solutions to simplified BCRE equations can be constructed by the methods proposed in . Using the BCRE model, Bouchaud and Cates explained another type of avalanches (in a thin granular layer on an inclined plane, see ).
Our aim here is to investigate a relation between the two models mentioned above. After reminding briefly the variational and BCRE models, we propose a full-dimensional generalization of the latter, originally formulated by BCRE in the one-dimensional case. To do this, we modify and extend the constitutive relations determining the surface flow velocity and the rolling-to-immobilized-state transition rate: BCRE’s assumption that the slope is everywhere almost critical is too restrictive for our purpose. Rescaling the variables, we show that the modified BCRE model contains a small parameter, the ratio of a characteristic rolling grains layer thickness to the pile size, and hence may often be simplified by employing a quasistationary equation for the rolling grains layer. The issue of scaling turns out to be very important in description of pile growth: another dimensionless parameter in the model thus obtained is the ratio of a typical rolling grain path length to the pile size. For large piles, this coefficient is also small and we show that in the long-scale limit the modified BCRE model tends to the variational model . For small piles, the corresponding term can be significant. These results make clear why differ the shapes of small and large piles and, correspondingly, why different models should be used to simulate, say, formation of large sand dunes and small Aeolian ripples.
## II Variational model of pile growth
Let a cohesionless granular material having an angle of repose $`\alpha _r`$ be tipped out onto a given rough rigid surface $`y=h_0(x)`$, where $`x=(x_1,x_2)`$<sup>2</sup>. We want to find the shape of a pile thus generated.
The real process of pile growth is often intermittent: discharged granular material not only flows continuously over the pile slopes but is also able to build up and then to pour suddenly down the slope in an avalanche. However, the avalanches usually involve only a small amount of particles in a pile and cause small fluctuations of the pile free surface. The model neglects these fluctuations and is a model for the mean surface evolution. Whether the pile evolution is governed by a continuous surface flow or results from many small avalanches, the surface flow is typically confined to a thin boundary layer which is distinctly separated from the motionless bulk .
Let us assume for simplicity that the support surface has no steep slopes, i.e.,
$`|h_0|k,`$
where $`k=\mathrm{tan}\alpha _r`$ (see for the general case). Assuming the bulk density of material in a pile is constant we can write the conservation law as
$`_th+𝒒=w,`$
where $`h(x,t)`$ is the free surface, $`𝒒(x,t)`$ is the horizontal projection of the flux of rolling particles, and $`w(x,t)`$ – the source intensity. We neglect the inertia and suppose that surface flow is directed towards the steepest descent,
$`𝒒=mh,`$
where
$$m(x,t)0$$
(1)
is an unknown scalar function. The conservation law takes now the form
$$_th(mh)=w.$$
(2)
It is assumed in this model that the surface slope angle cannot exceed the angle of repose,
$$|h|k,$$
(3)
and that no pouring occurs over the parts of the pile surface which are inclined less:
$$|h(x,t)|<km(x,t)=0.$$
(4)
To complete the model we have to specify the initial,
$$h|_{t=0}=h_0,$$
(5)
and a boundary condition. Let the granular material be allowed to leave the system freely through part $`\mathrm{\Gamma }_1`$ of the boundary of domain $`\mathrm{\Omega }`$<sup>2</sup>, and the other part of the boundary, $`\mathrm{\Gamma }_2`$, presents an impermeable wall. The boundary conditions are then, respectively,
$$h|_{\mathrm{\Gamma }_1}=h_0|_{\mathrm{\Gamma }_1},m_nh|_{\mathrm{\Gamma }_2}=0.$$
(6)
The model (1)-(6) contains two coupled unknowns, the free surface $`h`$ and an auxiliary function $`m`$ determining the rolling grains flux magnitude. Conditions (1), (3), and (4) define $`m`$ as a multivalued function of $`|h|`$, see Fig. 1.
The problem (1)-(6) may be considered an anomalous diffusion problem and solved by approximating this highly nonlinear multivalued relation. However, a better way to solve this problem is based on its following reformulation in the form of an evolutionary variational inequality (see and for variational inequalities in mechanics and physics and their numerical solution, respectively).
Let us define the set $`K`$ of possible surfaces as
$`K=\{\phi (x)||\phi |k,\phi |_{\mathrm{\Gamma }_1}=h_0|_{\mathrm{\Gamma }_1}\}`$
and the scalar product of two functions as $`(\varphi ,\psi )=_\mathrm{\Omega }\varphi \psi 𝑑x`$. We can now consider the following problem (variational inequality):
$$\{\begin{array}{c}Find\text{ }h(x,t)\text{ }such\text{ }that\text{ }hK\text{ }for\text{ }all\text{ }t>0,\\ (_thw,\phi h)0\text{for all}\phi K,\\ \text{and}h|_{t=0}=h_0.\end{array}$$
(7)
Theorem. Function $`h(x,t)`$ is a solution of the variational inequality (7) if and only if there exists $`m(x,t)`$ such that the pair $`\{h,m\}`$ is a solution to (1)-(6).
The outline of the proof is given in (see for mathematical details and a proof of existence of a unique solution to the variational inequality (7)). It has been also shown that the surface flux magnitude $`m(x,t)`$ is, in this model, a Lagrange multiplier related to the point-wise constraint (3). The values of such multipliers are not uniquely determined by the local conditions, which is the ”mathematical explanation” of long-range interactions typical of extended dissipative systems in a critical (marginally stable) state, see .
The model (1)-(6) or, equivalently, (7), has simple analytical solutions , such as the conical pile growing under a point-like source or the piles on flat open platforms described in . Numerical solutions also demonstrate simple geometrical structures that agree with one’s sandbox memories. Although this model is much simplified in many respects, it allows for the multiplicity of possible pile shapes. The avalanches may be introduced into the model as solution discontinuities (in time) triggered by sudden changes of the admissible set $`K`$ and are instantaneous events. On the time scale of a slow pile growth the life of an avalanche is, indeed, very short.
## III Modified BCRE equations
The BCRE equations involve two coupled variables: the pile height, $`y=h(x,t)`$, and the effective thickness (density) of the rolling grains layer, $`R(x,t)`$ ($`R(x,t)d\mathrm{\Omega }`$ is the volume that the material, currently rolling above the area $`d\mathrm{\Omega }`$, would occupy in the pile). The model has been formulated for a two-dimensional pile ($`x`$<sup>1</sup>); free surface slope deviations from the critical angle were assumed small. Original BCRE equations included diffusion terms to account for a non-locality of grains dislodgement and for fluctuations of rolling grains velocity. Although diffusion plays a crucial role in BCRE’s scenario of avalanches , these terms were regularly omitted by other researchers who either assumed that in their problems diffusion is insignificant and simplified the model, or proposed a different avalanche scenario (see, e.g., ). Below, we also omit the diffusion terms at first but introduce small diffusion at a later stage as a means for model regularization in transition to a large-scale limit.
Simplified BCRE equations may be written as follows:
$`_th=\mathrm{\Gamma }[h,R],_tR+_x(vR)=w\mathrm{\Gamma }[h,R].`$
Here the term $`\mathrm{\Gamma }[h,R]`$ accounts for the conversion of rolling grains into immobilized grains and vice versa, $`v`$ is the horizontal projection of rolling grains velocity, and $`w(x,t)`$ is the source intensity (we assume that the tipped grains do not stick to the pile surface but join the rolling grains first).
Limiting their consideration to the slopes that are close to critical, BCRE assumed constant downslope drift velocity $`v`$. The surface flux magnitude, $`q=vR`$, is thus determined solely by the rolling layer thickness $`R`$. Since in the previous model $`q=mk`$ for the critical slopes, $`m`$ and $`R`$ play similar roles and the two choices of basic variables, $`\{h,m\}`$ and $`\{h,R\}`$, are essentially equivalent. The exchange term $`\mathrm{\Gamma }`$ in BCRE model is linearized in a vicinity of the critical angle $`\alpha _r`$ and is proportional (for thin surface flows) to $`R`$: $`\mathrm{\Gamma }[h,R]=\gamma R(\alpha _r\theta ),`$ where $`\theta (x,t)`$ is the surface slope angle and $`\gamma `$ is a coefficient.
For a three-dimensional pile ($`x`$<sup>2</sup>), the model equations are similar,
$`_th=\mathrm{\Gamma }[h,R],_tR+(𝒗R)=w\mathrm{\Gamma }[h,R],`$ (8)
but the constitutive relations determining $`𝒗`$ and, probably, $`\mathrm{\Gamma }[h,R]`$ should be modified; here we will follow (see also ). We assume that the rolling particles drift towards the steepest descent of the free surface with a mean velocity $`v`$ depending on the slope angle (the steeper the slope, the higher is the velocity). On their way downslope, these particles may be trapped and absorbed into the motionless bulk (the steeper the slope, the lower is the trapping rate $`\mathrm{\Gamma }`$). If the surface is horizontal, the mean flow velocity is zero and the trapping rate is maximal, for $`\theta =\alpha _r`$ the rolling particles follow without trapping. Below, we will not consider the overcritical slopes and assume also that the trapping rate is proportional to the amount of rolling grains $`R`$.
At least partially, this simplified picture can be justified by recent experimental, theoretical, and numerical studies on the motion of a spherical particle on a rough inclined plane . For the relevant region of slope angles $`\theta `$, the energy dissipation due to the multiple shocks experienced by a moving particle is equivalent to the action of a viscous friction force . Because of that such particles reach a constant mean velocity proportional to $`\mathrm{sin}\theta `$. Sometimes, however, the particles are suddenly trapped in a well and completely lose their momentum in the direction of motion . Of course, conditions in the collective flow of grains over the pile surface are somewhat different. In particular, the flow velocity may depend on the thickness of rolling grains layer and the exchange rate is not exactly proportional to $`R`$ . Various improved dependencies can be incorporated into the model. The limiting behavior of the modified BCRE model is, however, robust and does not depend on details. For clarity of presentation we will consider the long-scale limit of a thin-flow model with the simplest phenomenological relations determining the flow velocity and rolling-to-immobilized-state transition.
Since the mean velocity of surface flow is proportional to $`\mathrm{sin}\theta `$ , its horizontal projection $`v`$ is proportional to $`\mathrm{sin}\theta \mathrm{cos}\theta =\mathrm{tan}\theta /(1+\mathrm{tan}^2\theta )`$. Postulating that the flow is in the steepest descent direction, we obtain $`𝒗=\mu h/(1+|h|^2),`$ where $`\mu `$ is a coefficient. Simplifying this relation we assume
$$𝒗=\mu h.$$
(9)
The exchange rate $`\mathrm{\Gamma }`$ should not depend on the slope orientation and we assume it to be a smooth decreasing function of $`|h|^2`$ that becomes zero for critical slopes. Assuming $`\mathrm{\Gamma }`$ is proportional to $`R`$ (thin flows) we arrive at
$$\mathrm{\Gamma }[h,R]=\gamma R\left(1\frac{|h|^2}{k^2}\right)$$
(10)
as the simplest constitutive relation . We will now derive a dimensionless formulation for the modified BCRE model (8)-(10).
The parameters in this model have the following dimensions: $`[\gamma ]=𝒯^1`$ and $`[\mu ]=𝒯^1`$. Let us denote by $`\overline{w}`$ the characteristic intensity of the external source; $`[\overline{w}]=𝒯^1`$. The three length scales characterizing the pile surface dynamics and surface granular flow may be defined as follows:
* typical thickness of the rolling grains layer, $`L_R=\overline{w}/\gamma `$;
* mean path of a rolling particle before it is trapped strongly depends on the slope steepness but, for a fixed subcritical slope, is proportional to the ratio $`L_P=\mu /\gamma `$ characterizing the competition between rolling and trapping;
* the pile size $`L`$.
The time $`T=L/\overline{w}`$ needed for a source with given intensity $`\overline{w}`$ to produce a pile of size $`L`$ may be used as a long time scale. Rescaling the variables,
$`x^{}={\displaystyle \frac{1}{L}}x,h^{}={\displaystyle \frac{1}{L}}h,R^{}={\displaystyle \frac{1}{L_R}}R,w^{}={\displaystyle \frac{1}{\overline{w}}}w,t^{}={\displaystyle \frac{1}{T}}t,`$
we arrive at the following dimensionless formulation:
$$_th=\mathrm{\Gamma }[h,R],$$
(11)
$$\frac{L_R}{L}_tR\frac{L_P}{L}(Rh)=w\mathrm{\Gamma }[h,R],$$
(12)
$$\mathrm{\Gamma }[h,R]=R\left(1\frac{|h|^2}{k^2}\right).$$
(13)
Typically, $`L_RL_P<L.`$ The first coefficient in (12) is very small, so it may often be possible to omit the corresponding term and use a quasistationary equation for the rolling layer. Such an approach has already been employed in simulation of the dynamics of sand ripples, see . The second coefficient, $`L_P/L`$, may be significant for small piles, like sand ripples, but becomes small too for large piles. Further simplification of the model is then appropriate.
## IV The long-scale limit of BCRE model
Let us denote $`\nu =L_P/L`$ and study the $`\nu 0`$ behavior of the model (11)-(13). This limit corresponds to the case of large piles ($`LL_P`$). We want to show that in this limit the pile shape evolution is described by the variational inequality (7) which remains invariant under the rescaling employed.
Physically, the situation is clear: although the model (11)-(13) permits grains to roll down upon any inclined slope, the rolling particles are quickly stopped and their paths are short comparing to the pile size for all except the almost critical slopes. This is essentially what is assumed in the model which permits rolling upon the critical slopes only. Mathematically, the situation is somewhat more complicated.
Since $`L_RL_P`$, we assume $`L_R/L`$ is $`o(\nu )`$ and set $`L_R/L=\nu \lambda (\nu ),`$ where $`\lambda `$ tends to zero as $`\nu 0`$. Let us introduce a new variable, $`m=\nu R`$, define $`\psi (u)=1u^2/k^2`$, and rewrite the model (11)-(13) as
$$_th=\frac{m\psi (|h|)}{\nu },\lambda _tm(mh)=w\frac{m\psi (|h|)}{\nu }.$$
For any $`\nu >0`$ this system consists of two coupled hyperbolic equations. The second equation, which can be regarded as an equation for $`m`$, contains in its main part the coefficient $`h`$ which may be discontinuous. The theory for such equations is complicated and not well developed. To circumvent the difficulty, we add small diffusion to both equations and consider the regularized model
$$_th=\frac{m\psi (|h|)}{\nu }+\epsilon _h\mathrm{\Delta }h,$$
(14)
$$\lambda _tm(mh)=w\frac{m\psi (|h|)}{\nu }+\epsilon _m\mathrm{\Delta }m,$$
(15)
where the positive coefficients $`\epsilon _h(\nu )`$ and $`\epsilon _m(\nu )`$ vanish as $`\nu `$ tends to zero. It should be noted that, although small diffusion may be physically meaningful and has been included into the original BCRE formulation , here we introduce it merely as a parabolic regularization of hyperbolic equations convenient for analyzing the model’s behavior at $`\nu 0`$.
We assume the same initial and boundary conditions, (5) and (6) correspondingly, for the function $`h`$. The non-negative values of $`m`$ both in $`\mathrm{\Omega }`$ at $`t=0`$ and on the boundary of this domain for $`t>0`$ may be chosen arbitrary: these initial and boundary conditions result only in the appearance of boundary layers in the solution for any finite $`\nu >0`$ and are lost in the $`\nu 0`$ limit. Rigorously, convergence of the problem (14)-(15) to the variational inequality (7) is proved elsewhere . Here we present a simplified scheme of the proof and avoid technicalities.
The main step is, as usual, to obtain uniform in $`\nu >0`$ a priori estimates on the solutions of the equations (14)-(15). First, taking the gradient and multiplying by $`h`$, we derive from (14) a parabolic partial differential equation for $`|h|^2`$. Since $`\psi (k)=0`$ and $`|h_0|k`$, we are able, using the maximum principle for this equation, to show that for $`\nu >0`$
$$|h|k\text{for all}(x,t);h\text{is uniformly bounded.}$$
(16)
Second, using the non-negativeness of the source function $`w(x,t)`$ and applying the maximum principle to the equation (15), we deduce that for each $`\nu >0`$
$$m0\text{ for all }(x,t).$$
(17)
Applying the estimates (16), (17) to the equation (14) we obtain
$$m\psi (|h|)=O(\nu ).$$
(18)
Sending $`\nu `$ to zero in (16), (17), (18) we establish the fulfillment in this limit of the conditions (1), (3), and (4). Finally, adding the equations (14) and (15) we obtain
$`\lambda _tm+_th(mh)=w+\epsilon _h\mathrm{\Delta }h+\epsilon _m\mathrm{\Delta }m.`$
Since $`\lambda `$, $`\epsilon _h`$, and $`\epsilon _m`$ vanish as $`\nu 0`$, we can show that the corresponding limits of $`h`$ and $`m`$ satisfy also the balance equation (2) in some weak (integral) sense. This completes the proof, because the model (1)-(6) is equivalent to the variational inequality (7).
To illustrate this result we will now compare solutions of the BCRE-type model (14)-(15), solutions of the variational inequality (7), and real shapes of small and large piles. Let us consider the simplest situation: a pile growing under a point source on an infinite horizontal support $`h_0=0`$. Although in this case the piles are known to be almost perfect cones, sometimes one can notice curved tails near the bottom of a small pile (Fig. 2a). As the pile becomes larger, the tail remains of only, say, tens of grain diameters long, so the tail of a large pile is difficult to see (Fig. 2b).
The modified BCRE model (14)-(15) describes this situation quite satisfactory, see Fig. 3. Although the tails of small piles ($`\nu =0.2`$) are clearly seen, tails of the larger piles ($`\nu =0.01`$) is difficult to detect. We see also that piles, predicted by the BCRE model with small $`\nu `$ and $`\lambda `$, are very close to the growing cone, exact analytical solution of the variational inequality (7). It may be noted that for small values of $`\nu `$ and $`\lambda `$ the model equations are stiff and their numerical solution becomes difficult. Thus, even using an implicit finite-difference approximation of (14)-(15), we needed $`10^5`$ time steps in the latter example.
## V Conclusion
We considered two different continuous models for the pile surface dynamics: the BCRE model and the variational model . Both models are written for two coupled dependent variables and are able, in principle, to account for multiplicity of metastable pile shapes and surface avalanches. It has been found that the models are related and describe the pile surface dynamics on different spatio-temporal scales.
BCRE-type models may be used to simulate the fast processes, such as the initiation, spreading, and settling down of an avalanche. To describe the much slower dynamics of the mean shape of a pile, the model may often be simplified by employing a quasistationary equation for the rolling grains layer. Such a model is able to predict some peculiarities of small pile shapes and was recently employed for simulating the nonlinear dynamics of sand ripples .
Unlike the BCRE models, the variational model of pile growth does not permit the discharged grains to roll upon subcritical slopes and is therefore unable to account for such features of sand surface as sand-ripple instability or surface slope deviation from the critical angle near the bottom of a conical pile. Indeed, these effects are determined by rolling of particles upon the subcritical slopes and are exhibited on the length scale comparable to the mean path of a particle prior to its being trapped.
On the other hand, sand ripples on the dune surface or tiny tails at the bottom of a pile are seen only from a short distance. These small details are difficult to distinguish watching from a larger distance allowing one to follow the evolution of a big dune or formation of a large pile. In such situations the BCRE model contains another small parameter. This complicates simulations and makes them inefficient. As has been shown in our work, in the long-scale limit the BCRE-type models converge to the variational model of pile growth. The latter model is more appropriate for simulating the pile surface dynamics on a large spatio-temporal scale.
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# RELATIVISTIC CONIC BEAMS AND SPATIAL DISTRIBUTION OF GAMMA-RAY BURSTS
## 1 INTRODUCTION
The BATSE experiment on the Compton Gamma Ray Observatory and the study of the afterglows (e.g., Piran 1999 and references therein). have established that the gamma-ray bursts (GRBs) are cosmological (Mao and Paczy$`\stackrel{´}{\mathrm{n}}`$ski, 1992; Meegan et al., 1992; Piran, 1992). Even though the distance scale seems settled (Metzger et al., 1997), it appears that uncertainties remain in the total energy and the burst rate of GRBs (Kumar, 1999; Kumar and Piran, 1999). These two important issues depend on the level of beaming in GRB emission. That is, the issues critically depend on whether the geometry of the gamma-ray emitting ejecta is spherical or jet-like (Harrison et al., 1999; Kulkarni et al., 1999; $`\mathrm{M}\stackrel{´}{\mathrm{e}}\mathrm{sz}\stackrel{´}{\mathrm{a}}\mathrm{ros}`$ and Rees, 1999; Sari et al., 1999). A number of authors have studied energetics and geometry of the ejecta (Mao and Yi, 1994; Rhoads, 1997; Panaitescu and $`\mathrm{M}\stackrel{´}{\mathrm{e}}\mathrm{sz}\stackrel{´}{\mathrm{a}}\mathrm{ros}`$, 1998; Rhoads, 1999; Moderski et al., 2000). It is also important to put constraints on the width of the luminosity function by comparing the observed intensity distribution with those predicted by a physical model (Mao and Yi, 1994; Yi, 1994). In essence, the rate, the energy, and the luminosity function of GRBs are all closely related to whether or not the geometry of the ejecta is spherical.
Two most frequently quoted statistics in GRB observations are $`<V/V_{max}>`$ and $`\mathrm{log}N(>F)\mathrm{log}F`$, where $`F`$ refers to the peak flux (or peak count rate) and $`N`$ denotes the number of GRBs with fluxes higher than $`F`$ (e.g., Yi 1994). These two quantities contain information on the lumonosity function of GRBs and the spatial number density of the sources. A value of $`<V/V_{max}>`$ consistent with that of an observed sample is a necessary condition but not a sufficient condition for a luminosity function $`\mathrm{\Phi }(L)`$ which is neither directly observable nor theoretically well undertood. The luminosity function of GRBs can be obtained for an assumed source distribution $`n(z)`$ such that the calculated $`\mathrm{log}N(>F)\mathrm{log}F`$ fits the observed distribution, and vice versa. The density $`n(z)`$ refers to the rate of GRBs per unit time per unit comoving cosmological volume. However, due to the very nature of $`N(>F)`$, which is the convolution of $`n(z)`$ and $`\mathrm{\Phi }(L)`$, one almost always obtains $`n(z)`$ for a given $`\mathrm{\Phi }(L)`$ such that the theoretical $`\mathrm{log}N(>F)\mathrm{log}F`$ curve fits the observed intensity distribution. Therefore, in order to extract information concerning $`n(z)`$ or $`\mathrm{\Phi }(L)`$, one has to assume one of these two functions or to develop a techinique to separate the effects of these two unknown functions (Horack and Emslie 1994; Horack et al. 1994; Ulmer et al. 1995; Ulmer and Wijers 1995). It is therefore of great interest to construct $`\mathrm{\Phi }(L)`$ on the basis of a physical model, which is one of our major goals in this Letter. Since there remain uncertainties in GRB engine models, we focus on the consequences of the conical beaming without specifying how a beam is formed in a physical engine model.
Using the first BATSE catalog of gamma-ray bursts (Fishman et al., 1994), Mao and Yi (1994) studied the effects of the relativistic bulk motion in a conical beam on the statistics of gamma-ray bursts. They found that the luminosity function is naturally introduced by the random distribution of the space orientation of the cone axis and that the case of the standard candle is not easily distinguished from that of the beaming-induced luminosity function with a sharp peak. This is especially the case for large beam opening angle and the large Lorentz factor $`\gamma `$, as one may expect. Different Lorentz factors and opening angles however result in non-trivial changes for the distances to GRBs and especially the highest redshift of or the maximum distance to the most distant GRB for a given sample. For instance, the maximum redshift $`z_{max}`$ increases as the ratio of the opening angle to $`\gamma ^1`$ decreases. We modify the conical beam model by allowing a spatial variation of $`\gamma `$ and the density profile of gamma-ray emitting electrons on the photon-emitting surface of the cone.
From numerical simulations of relativistic jets (Marti et al., 1997; Renaud and Henri, 1998; Rosen et al., 1999) and observations of the astrophysical jets (Zensus, 1997; Spruit, 2000), it is clear that jets do have some significant structure in them and the bulk Lorentz factors evolve as the jets propagate. Therefore, it is plausible to extend the simplest jet model such as that of Mao and Yi (1994). In a more realistic jet model, the bulk Lorentz factor has a spatial profile at the surface where the observed gamma-ray emission occurs and the spatial electron density distribution is significantly inhomogeneous.
In $`\mathrm{\S }`$ 2 we begin with a brief presention of data we use, which are parts of the BATSE 4B catalog (Paciesas et al., 1999), and we describe our conical jet geometry, following Mao and Yi (1994), in $`\mathrm{\S }`$ 3 we present results. Finally, we conclude with summary of our results and discussions in $`\mathrm{\S }`$ 4.
## 2 OBSERVATIONAL DATA AND MODEL
The BATSE 4B catalog (Paciesas et al., 1999) provides 1637 triggered GRBs detected from 1991 April through 1996 August. We use the bursts which are detected on the 1024 ms trigger time scale. We choose the bursts of which peak count rates are above 0.4 $`\mathrm{photons}\mathrm{cm}^2\mathrm{s}^1`$ in order to avoid the threshold effects (cf. Mao and Yi 1994). Of those bursts, we select the GRBs whose $`\mathrm{C}_{\mathrm{max}}/\mathrm{C}_{\mathrm{min}}`$ is greater than $`1.0`$, which gives a sample of 651 bursts. The energy range in which the peak flux is measured is $`50300\mathrm{keV}`$. The peak fluxes of bursts vary by about 2 orders of magnitude in this data set.
Throughout this Letter we adopt the simplest cosmological model; the universe is flat with the density parameter $`\mathrm{\Omega }_0=1`$, the Hubble constant is 65 $`\mathrm{kms}^1\mathrm{Mpc}^1`$, the cosmological constant is absent, all the GRBs are ’standard bursts’ with an identical power-law spectrum, and the rate of bursts per unit comoving volume per unit comoving time is constant (cf. Mao and Paczy$`\stackrel{´}{\mathrm{n}}`$ski 1992; Yi 1994).
In our beaming model, the ejecta is flowing outward relativistically in a cone with the geometrical opening angle $`\mathrm{\Delta }\theta `$. The observed gamma-ray emission is produced at radius $`R`$ from the central engine where the radiation first becomes optically thin. In the minimal beaming model of Mao and Yi (1994), the ejected material has the same bulk Lorentz factor $`\gamma `$ at this distance and the photon-emitting electrons’ density at the surface of the cone is uniform. In this Letter those simplifications are relaxed as described below.
In the cylindrical symmetry one may consider the colatitude alone, which is defined by the angle between the line of sight and the symmetry axis of the cone. Besides the spherical coordinate system centered on the central engine, we introduce an auxiliary spherical coordinate system with the $`z^{^{}}`$-axis along the symmetry axis of the cone. It can be shown that the angle $`\mathrm{\Theta }`$ between a position of a direction within the cone and the line of sight is given by
$$\mathrm{cos}\mathrm{\Theta }=\mathrm{cos}\theta ^{^{}}\mathrm{cos}\theta \mathrm{sin}\theta ^{^{}}\mathrm{sin}\varphi ^{^{}}\mathrm{sin}\theta .$$
(1)
The monochromatic flux received by a local observer at a distance $`D`$ from the source, taking the cosmological redshift effects into account, reads
$`F(\nu ,\theta )`$ $`=`$ $`{\displaystyle \frac{(1+z)R^2}{D_L^2(z)}}\times {\displaystyle _0^{2\pi }}𝑑\varphi ^{^{}}`$ (2)
$`\times {\displaystyle _0^{\mathrm{\Delta }\theta }}\mathrm{sin}\theta ^{^{}}d\theta ^{^{}}\mathrm{cos}\mathrm{\Theta }\mathrm{\Gamma }^3I_0[\nu (1+z)\mathrm{\Gamma }^1],`$
where $`D_L(z)=2c/H_0[1+z(1+z)^{1/2}]`$, $`c`$ and $`H_0`$ being the speed of light and the Hubble constant respectively, and $`\mathrm{\Gamma }=[\gamma (1\beta \mathrm{cos}\mathrm{\Theta })]^1`$, $`\beta =(1\gamma ^2)^{1/2}`$. And we have used the relation $`\nu =\mathrm{\Gamma }\nu _0`$; here $`\nu _0`$ is the corresponding frequency in the comoving frame. In this model we have ignored the structure of light curves due to the relative time delay of radiation from different parts of the cone, and other complicated cosmological effects.
We define the terms, excluding terms for the cosmological information in the above equation (2), as the local peak count rate,
$`P_{loc}(\theta )`$ $`=`$ $`{\displaystyle \nu ^{1\alpha }𝑑\nu \times R^2_0^{2\pi }𝑑\varphi ^{^{}}}`$ (3)
$`\times {\displaystyle _0^{\mathrm{\Delta }\theta }}\mathrm{sin}\theta ^{}\mathrm{cos}\mathrm{\Theta }\mathrm{\Gamma }^{2+\alpha }d\theta ^{}.`$
The maximum local peak count rate is achieved at $`\theta =0`$ and can be obtained analytically :
$`P_{loc,max}(\theta =0)`$ $`=`$ $`{\displaystyle \nu ^{1\alpha }𝑑\nu \frac{2\pi R^2}{\beta ^2\gamma ^{2+\alpha }}}`$ (4)
$`\times [{\displaystyle \frac{1}{\alpha }}(x_u^\alpha x_l^\alpha ){\displaystyle \frac{1}{1+\alpha }}(x_u^{1+\alpha }x_l^{1+\alpha })],`$
where $`x_u^1=1\beta \mathrm{cos}\mathrm{\Delta }\theta `$, and $`x_l^1=1\beta `$.
Based on the randomness of the direction of the cone axis with respect to the line of sight, we obtain the probability function which directly reflects the angle between the line of sight and the direction of the cone, $`\theta `$. Assuming that the cone is uniformly distributed in space, the probability function for a one-sided cone is given by
$$p(\theta )d\theta =\frac{1}{2}\mathrm{sin}\theta d\theta ,$$
(5)
where $`\theta `$ is between $`0`$ and $`\pi `$. Since the local peak count rate is a function of the orientation of the cone, one may translate the probability function of the cone’s angle into the probability function of the local peak count rate.
In order to obtain the luminosity function and compare with observations, we calculate the local peak count rate in a fixed frequency range. We allow variations of $`\gamma `$ and the electron density at the surface of the cone by introducing a window function for $`\gamma `$ and electron density, $`N_e`$. The window function is axisymmetric with respect to the symmetry axis of the cone. The window function we adopt is the Gaussian function centered at the center of the cone:
$`W(\theta ^{})=\mathrm{exp}[A({\displaystyle \frac{\theta ^{^{}}}{\mathrm{\Delta }\theta }})^2],`$ (6)
where $`A`$ is a constant, $`\mathrm{\Delta }\theta `$ is a given opening angle, and $`\theta ^{^{}}`$ varies from $`0`$ to $`\mathrm{\Delta }\theta `$.
## 3 RESULTS
The probability function of the local peak count rate is shown in Figure 1 for three different ratios of $`\mathrm{\Delta }\theta `$ to $`\gamma ^1`$, in the case where $`\gamma `$ is uniform over the photon-emitting surface of the cone. For a fixed $`\gamma `$, as the opening angle increases the peak of the probability of the local peak count rate becomes higher and narrower, as one may expect. A model with a large value of Lorentz factor ($`\mathrm{\Delta }\theta \gamma ^1`$) essentially provides a luminosity function indistinguishable from that of the standard candle model. For a given value of the ratio, the whole distribution function moves vertically when the absolute value of the Lorentz factor changes. This is because the local peak count rate is decreasing faster with the angle for larger $`\gamma `$. In the extreme case where $`\mathrm{\Delta }\theta \gamma ^1`$ the distribution is a power law with the index of $`1/3`$, as expected from the analytical result (Yi, 1993).
In Figure 2, we show the probability distribution functions for both the varying bulk Lorentz factor and the inhomogeneous electron density. We assume the axisymmetry of the Lorentz factor profile around the cone axis, $`\gamma =\gamma (\theta ^{^{}})`$, and hence the Lorentz factor profile could mimic a simplified model for the jet-environment drag. At the center of the cone, $`\gamma `$ has the maximum value and decreases with $`\theta ^{}`$. As $`\gamma `$ decreases more steeply with $`\theta ^{}`$, the $`\mathrm{log}p_L(\mathrm{log}P_{loc})`$ shows a smaller peak at $`\mathrm{log}(P_{loc}/P_{loc,max})=0`$, and a higher level of the ’tail’ of the $`\mathrm{log}p_L(\mathrm{log}P_{loc})`$. It is because this type of $`\gamma `$ effectively reduces the ’average’ value of $`\gamma `$ over the cone and the ’effective opening angle’ simultaneously. For a very narrow window function it essentially reduces to the pencil-beam case. The photon-emitting electrons are supposed to be distributed according to the $`\gamma `$ distribution such that the local electron number density is inversely proportional to the square of the bulk Lorentz factor, $`\gamma `$. That is, the electron number density is a function of the angular position in the cone. In this profile, the electron number density increases outwards from the center of the cone, since $`\gamma `$ is a decreasing function. In effect, the cone is reminiscent of the hollow cone.
Once the probability distribution is obtained, we are in a position to calculate $`\mathrm{log}N(>F)\mathrm{log}F`$ with an assumed spatial distribution of GRBs, $`n(z)`$. In Figure 3, we compare the cumulative intensity distribution curve produced by our luminosity function with observational data. Since the goal of this study is looking at effects of the conic beam with varying distributions of $`\gamma `$ and the electron density over the surface of the cone, we assume there is no source evolution, that is, $`n(z)=n_0`$. We adopt the power law index $`\alpha `$ of unity instead of 2 (cf. Mao and Yi, 1994). This is a simplification in a sense that observed GRBs show various power-law indices. This value, however, represents the averaged power law indices of observed GRBs (Band et al., 1993; Mallozzi et al., 1996; Preece et al., 2000). All the theoretical cumulative probability distribution with the luminosity functions shown in Figure 2 are plotted in Figure 3 along with the observed distribution. The curves shown in Figure 3 are best-fit functions determined by the Kolmogorov-Smirnov (K-S) test for each parameter set. Parameters in our model are not sensitive enough to constrain the luminosity function.
For a given set of model paramters, $`z_{max}`$ plays a role of a parameter in the K-S test in that the luminosity function studied here optimizes $`\mathrm{log}N(>F)\mathrm{log}F`$ to fit observational data. The obtained $`z_{max}`$’s in this procedure are shown in the Table 1. These $`z_{max}`$’s obtained with the 4B BATSE catalog are greater than those for the 1B BATSE catalog by about $`0.1`$ on average, which indicates that we are seeing fainter GRBs in the 4B BATSE than in the 1B BATSE catalog and therefore more distant GRBs.
It shows that the redshift of the most distant GRBs becomes larger as $`\gamma `$ falls more steeply from the cone center. As shown in the Table 1, our toy models with the narrow Gaussian profiles could easily reproduce $`z_{max}`$ values as high as the highest reported $`z_{max}=3.42`$ (quoted from Bulik (1999), see references therein). This simple case indicates that the effects of the luminosity function have significant implications on the cosmological spatial distribution of GRBs. It is interesting to see that such a high redshift would be hard to explain in the homogeneous beam model with a constant $`n(z)`$. The effects of the Gaussian window on the values of $`z_{max}`$ are substantial enough for further comments. The model for a narrow window with $`A=4`$ gives $`z_{max}=3.7(\mathrm{average}\mathrm{redshift}<z>=1.03)`$ while the uniform beam model (i.e. without any variation of $`\gamma `$) gives $`z_{max}=1.6(<z>=0.51)`$. The average value of redshifts, $`<z>`$, is taken over the cumulative redshift distribution function shown in Figure 4. The cumulative redshift distribution is defined by
$`𝒩(z^{^{}})={\displaystyle \frac{N(z^{^{}}<z<z_{max})}{N(0<z<z_{max})}}`$ (7)
where
$`N(z^{^{}}<z<z_{max})={\displaystyle _z^{^{}}^{z_{max}}}{\displaystyle \frac{4\pi }{1+z}}n(z)r^2(z)𝑑r(z).`$ (8)
The cumulative redshift distributions shown in Figure 4 also indicate that the luminosity distribution induced by the beaming has a significant implication on the GRBs distances in the observed samples. When the beam is sharply peaked at the beam center (solid curves in the upper panel in Figure 4) the fraction of the high redshift GRBs (e.g. $`z>3`$) could be as high as $`10\%`$ while the broad beam case (e.g. $`A=1/8`$) essentially rules out any high redshift GRBs. The cumulative distribution slowly decreases with the redshift when the $`\gamma `$ in the conic beam is decreasing rapidly in a sense that GRBs are spead out in a broader region beyond the averaged $`z`$ for the highly concentrated beam. That is, the ratio of $`z_{max}`$ to $`<z>`$ is 3.59 and 3.14 when $`A=4`$ and $`A=1/8`$, respectively. The beams’ electron density structure also has a similar effect on the redshifts of GRBs. The effects of the luminosity function have to be explicitly considered when observed high redshift values are interpreted (Krumholz et al., 1998). For instance, in the standard candle case, a significant probability for high redshift GRBs could directly imply a substantial source evolution effect. However, the beaming induced luminosity function could make this simple interpretation much uncertain (Blain and Natarajan, 2000).
## 4 DISCUSSION
The theoretical models for GRBs are abundant. Despite remarkable progresses in understanding physical mechanisms involved in these models, the GRB prompt emission mechanisms and engine models have so far been unable to constrain the extent of beaming and the luminosity distribution of GRBs. This in turn has been a major uncertainty in interpreting the observed flux data in terms of the cosmological spatial distribution of the bursts. In this regard, the present work has shown that the simple beaming models and their resulting apparent luminosity functions have significant effects in interpreting the observed data. If the GRBs are indeed standard candles with a single well-defined luminosity, the spatial distribution of GRBs in connection with the cosmic star formation rate could be translated into the cosmological source evolution. However, the luminosity distributions we have considered affect the maximum redshift and the average redshift significantly. It is therefore important to derive a theoretical luminosity function for a given GRB model.
The jet models we adopted are obviously over-simplified. Despite this major drawback, the models capture the essential ingredients of the beamed relativistic jets concerning the apparent luminosity function. One of the major uncertainties is that the jets and GRB sources differ greatly and GRB luminosities and jets’ physical conditions are intrinsically different in each source. Given the wide range of burst durations and the diverse burst types (Fishman and Meegan 1995, and references therein), such a possibility cannot be ruled out. If this is indeed the case, our standard source approach is not applicable.
We thank C. Kim, H. Kim, and K. Kwak for useful discussions and especially C. Kim for her help with Figure 3. IY is supported in part by the KRF grant No. 1998-001-D00365.
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# Dynamics of the Galactic Bulge using Planetary Nebulæ
## 1 Introduction
A spiral galaxy consists of a relatively flattened stellar disk in nearly circular rotation and, in most systems, a central bulge. It is estimated that about 30 % of these galaxies also show a central bar in the visible; however the real fraction of barred galaxies is probably significantly higher because some apparently normal spirals show a bar feature in the near-IR that was not visible in their optical images (e.g. Sellwood & Wilkinson, 1993). In addition, barred galaxies often show a lens and/or ring around the bar. The flattened disks contain objects of all ages, from the interstellar gas and very young stars to the old disk stars which in our Galaxy are almost as old as the globular clusters. The bulges appear to be made up mainly of old stars.
The disks of most disk galaxies are relatively thin, with the ratio of their radial to vertical scale heights mostly in the range 5 to 15. In the later-type barred galaxies, the central bar may be no thicker than the host disk. Kormendy (1993) has argued that many of the features identified as bulges from the surface photometry of more face-on galaxies may also be as thin as the disks. However, many edge-on galaxies show bulges which clearly do extend beyond the disk.
The bulges of spiral galaxies show a wide range of shapes, from spheroidal through boxy or peanut shaped bulges. The boxy versus spheroidal structure of bulges is roughly understood in terms of their orbital properties but not in terms of origin. Many possibilities have been suggested for the origin of boxiness in bulges, including the formation and dissolution of bars, dissipative processes during the collapse of a rapidly rotating inner region, or later accretion events (see Sellwood 1993; Rowley 1986; Whitmore & Bell 1988; Combes et al. 1990; Pfenniger et al. 1991).
Kormendy and Illingworth (1982) pointed out that the boxy bulges are frequently cylindrical rotators, unlike the more spheroidal bulges. This led to a burst of observational and theoretical studies of these systems (e.g. Binney & Petrou 1985; Rowley 1986; Shaw 1993), with the growing indication that these boxy or peanut-shaped edge-on systems may be associated with bars (Combes et al. 1990; Sellwood & Wilkinson 1993).
The Milky Way has an excellent example of a box-shaped bulge. This feature was seen in the early $`2.4\mu `$m balloon scans (Matsumoto et al. 1982), and spectacularly confirmed by the $`2.2\mu `$m image of the Galaxy from COBE<sup>1</sup><sup>1</sup>1Cosmic Background Explorer/DIRBE<sup>2</sup><sup>2</sup>2Diffuse Infrared Background Experiment<sup>,</sup><sup>3</sup><sup>3</sup>3The COBE datasets were developed by NASA Goddard Space Flight Center under the guidance of the COBE Science Working Group and were provided by the NSSDC. (Weiland et al. 1994; Arendt et al. 1994) as seen from the contours plot of the COBE/DIRBE 2.2$`\mu `$m image (Figure 1). See Binney et al. (1997) for a dust-corrected non-parametric recovery of the light distribution in the inner few kpc of the Milky Way from the COBE/DIRBE surface brightness map.
The bulge of the Milky Way provides a unique opportunity to investigate the detailed pattern of rotation and velocity dispersion in a boxy Galactic bulge. We can study the structure of the bulge to see if this boxy bulge is really a stellar bar, and we can also see how the bulge and disk are related dynamically.
This paper is outlined as follows. We start with an overview of recent studies of the bar/bulge problem through axisymmetric and N-body models (Section 2). In Section 3, we discuss the wide range of tracers available to study the kinematics of the Galactic bulge, and the data obtained for this study. A preliminary visual assessment of the data is presented in Section 4 using the mean velocity and velocity dispersion versus the Galactic longitude and latitude. Section 5 compares our planetary nebulæ (PNe) distribution with the distribution of light in the COBE/DIRBE images in the 1.25, 2.2 and $`3.5\mu `$m wavelength regions. In Section 6, we compare our PNe data with four Galactic bar-bulge models: three are N-body models and one is a relaxed Schwarzschild realization of the COBE light distribution. A summary and conclusions are given in Section 7.
## 2 Dynamics of the Bulge
### 2.1 The Bar/Bulge of the Milky Way
Evidence is accumulating that the boxy peanut-shaped bulges seen in edge-on disk galaxies are associated with bar structures (Combes et al. 1990; Jenkins & Binney 1994; Blitz & Spergel 1991; Kuijken & Merrifield 1995; Bureau & Freeman 1999). For the bulge of our Galaxy, the $`2.4\mu `$m balloon scans and the near-IR COBE/DIRBE images show such a boxy peanut shape. An unambiguous direct identification of a bar at the center of the Galaxy is difficult because the Sun is located in the plane of the Galaxy and our view of the Galactic center is obscured by the dust. The patchy extinction in the plane of the Galaxy is obvious from the optical image of the Galactic bulge taken at ESO (Madsen & Laustsen, 1986). It is clear that the southern part of the bulge is much less affected by extinction than the northern part. The southern part includes two famous regions of relatively low extinction, Sgr I ($`l=1.4^{}`$, $`b=2.6^{}`$) and Baade’s Window (BW) ($`l=1.0^{}`$, $`b=3.9^{}`$) which are widely used for studies of the stellar population and dynamics of the inner bulge. The distribution of extinction over the bulge is also nicely shown from the work of the COBE/DIRBE group (Arendt et al. 1994, figure 3b, plate L7).
Nevertheless, much observational evidence is now pointing to the existence of such a bar. Here, we list only a few: see Gerhard (1999) for a more detailed review.
* de Vaucouleurs (1964) was the first to point out that a central bar is probably responsible for the non-circular motions of the HI in the inner part of the Milky Way. He had already noted that similar non-circular motions were present in the inner parts of barred spiral galaxies.
* The asymmetry in longitude of the distribution of the $`2.4\mu `$m emission derived from the balloon scans indicated that the stars in the central kpc lie in a bar with its near side at positive Galactic longitude and suggested that the bar is tilted relative to the Galactic plane (Blitz & Spergel 1991).
* The COBE/DIRBE images (Weiland et al. 1994) confirmed the asymmetry in the surface brightness distribution of the bulge in the near-IR, but show no evidence for an out-of-plane tilt of the bar.
* Nikolaev & Weinberg (1997) reported that the distribution of variables in the IRAS Point Source Catalogue (PSC) is consistent with a bar with semi-major axis of 3.3 kpc and position angle of 24$`{}_{}{}^{}\pm 2^{}`$ (where position angle is the angle between the major axis of the bar and the Sun-center line and is taken as positive for a bar pointing into the positive Galactic longitude quadrant).
* Rohlfs and Kampmann (1993) showed that the HI terminal velocities indicate the presence of a bar with a semi-major axis of $`23`$ kpc and a position angle of about 45.
* Binney et al. (1991) used CO kinematics in the inner parts of the Galaxy to show the presence of a bar with a pattern speed of 63 $`\mathrm{kms}^1\mathrm{kpc}^1`$, a corotation radius of 2.4 kpc and a position angle of $`16^{}\pm 2^{}`$. More recent gas dynamical studies (Englmaier & Gerhard 1999; Weiner & Sellwood 1999; Fux 1999) all support for a substantially larger corotation radius.
From some of these studies, and others on the brightnesses of tracer objects like Mira variables (e.g. Whitelock 1993) and clump giants in the bulge (e.g. Stanek et al. 1994), it seems fairly clear that the bulge objects at positive Galactic longitude are brighter than those at negative longitude. This is generally interpreted as evidence that we are viewing the bar/bulge at an angle from its major axis and that the closer end of the bar is at positive longitude. There is still disagreement on the parameters of the bar, i.e. its length, strength, pattern speed and position angle. But, if we were to view the Galaxy edge-on from outside, it would probably look much like NGC 891, with probably more bulge than NGC 891 but less than NGC 4565 (see the Hubble Atlas).
### 2.2 Axisymmetric models
Kent (1992) used infrared ($`2.4\mu `$m) surface photometry from the Spacelab infrared telescope to make an axisymmetric model for the luminosity density distribution in the inner galaxy. For the disk, he modelled the luminosity density $`L`$ as a double exponential in $`R`$ and $`z`$, and for the bulge he adopted
$$L(R,z)=3.53K_{}(s/667)L_{}\mathrm{pc}^3\mathrm{for}\mathrm{s}>938$$
and
$$L(R,z)=1.04\times 10^6(s/0.482)^{1.85}L_{}\mathrm{pc}^3\mathrm{for}\mathrm{s}<938.$$
where $`K_{}`$ is a modified Bessel function. Here $`s^4=R^4+(z/0.61)^4`$ and the units of s in the equation above are parsecs. This form of the $`L(R,z)`$ distribution for the bulge leads to box-shaped isophotes.
Kuijken (1995) used a quadratic programming technique on a bilinear tessellation in the energy, angular momentum $`(E,L)`$ plane to construct a two-integral distribution function $`f(E,L)`$ for a slightly modified version of Kent’s axisymmetric model for the inner Galaxy. The distribution function is forced to give an isotropic velocity dispersion. With Kent’s values for the mass to light ratios for the disk and bulge, the predicted line-of-sight velocity distribution in Baade’s window is in excellent agreement with the distribution observed for the M giants by Sharples et al. (1990). However, the agreement is not so good for the velocity distribution of the K giants in Minniti’s (1992) field at $`l=8^{},b=7^{}`$: the discrepancy between the data and the prediction from the distribution function $`f(E,L)`$ is seen in the mean velocity and in the shape of the velocity distribution in this region. Kuijken suggests that the discrepancy might be associated with the triaxiality of the bulge, and points out how remarkable it is that his oblate, isotropic and axisymmetric model gives such a good fit to the velocity distribution in Baade’s Window.
Durand et al. (1996) used a two-integral axisymmetric model with a Kuzmin-Kutuzov Stäckel potential (with a halo-disk structure) to study the dynamics of a sample of 673 PNe taken from the Acker et al. (1992) catalogue. The method fits the kinematics to the projected moments of a distribution function by means of Quadratic Programming. They conclude that their two-integral model does not adequately characterize the dynamical state of their sample of PNe.
Our particular interest here is in investigating the triaxial structure of the bar-bulge further, so we will not pursue the axisymmetric models in this paper. The question is about the origin of central bar-bulges: do they arise from instabilities of the disk of galaxies or from other processes like the accretion of satellites or as part of the dissipative collapse of the galaxies ? The quantitative study of the formation of bars through disk instabilities is now well advanced through N-body models, which we now discuss briefly. In §6.3, we will compare the kinematical properties of our PNe and the models. For the N-body models, this comparison will show whether the instability picture gives a plausible description of the observed bulge kinematics.
### 2.3 N-body models
In the last few years, the growing evidence for a bar at the center of our Galaxy initiated much interest in developing detailed dynamical models of the Galactic bar-bulge. Different kinds of models are now available, but the observational constraints on their stellar dynamics are not yet well advanced.
N-body models of the bar-forming instabilities of disks provide theoretical predictions of the dynamics of the resulting bar-bulges which can be tested against dynamical data from the Galactic bulge and other bulges. For example, Fux (1996), Sellwood (1993) and Kalnajs (1996) have all modelled the central bar-bulge through the instabilities of self-gravitating stellar disks. As tests of the relevance of these models to the dynamics of the Galactic bulge, the detailed kinematics of their models can be compared with the observed kinematics of tracer bulge objects like the PNe which are the subject of this paper.
Another kind of numerical model for the Galactic bulge comes from the work of Zhao (1996) who constructed a rotating Schwarzschild model for the COBE light distribution. Although this model does not provide direct insight into the formation of the bulge, in the way that the studies of disk instabilities can do, the Schwarzschild model is of much interest for evaluating the present dynamical state of the bulge. For this purpose, we can compare the kinematics of N-body realizations of this model with observational data, as above.
It would be most desirable if we could obtain an unbiased spatial distribution and the radial velocities of a subset of bulge objects. Such a database would allow us to distinguish between the various proposed models, and no doubt suggest others. Unfortunately most of the stellar objects have to contend with the high and patchy absorption near the Galactic plane \[OH-IR stars are a clear exception\].
## 3 Planetary Nebulæ as tracers
To study the kinematics of the Galactic bulge, we have access to a wide range of tracers: OH/IR stars (Habing 1993; Sevenster et al. 1997a,b), Miras (Whitelock 1993), M giant stars, K giant stars (both individually and through the integrated bulge light) (Walker et al. 1990; Minniti et al. 1992; Minniti 1996a,b; Terndrup 1993; Ibata & Gilmore 1995a,b), carbon stars (Whitelock 1993), SiO Maser sources (Deguchi 1997), RR Lyræ stars (Walker & Terndrup 1991) and PNe (Kinman et al. 1988; Durand et al. 1998). The highly evolved OH/IR stars, Miras and M giants stars are probably biased towards the metal-rich population: the radial distribution of these objects is significantly steeper than the distribution of integrated light in the bulge (e.g. de Zeeuw 1993), and the kinematics of these objects reflects the kinematics of the metal-rich component of the bulge (Sevenster 1997). The carbon stars are rare and are also an indication of an intermediate age metal-rich population. The K giant stars are found at all metallicities and would be the ideal tracers to use since all bulge stars are likely to go through a K giant phase, but they are relatively faint. The K giants have already provided important dynamical information (e.g. Terndrup et al. 1995; Ibata & Gilmore 1995a,b), and much more will appear in the future from the large fiber surveys in progress (e.g. Harding & Morrison 1993). The RR Lyræ stars are also useful bulge tracers but they are biased toward the metal-poor population and are fainter than the K giants. The PNe are not biased towards the metal-rich population (e.g. Hui et al. 1993): recall the presence of PNe in the very metal-poor globular cluster M15 (Pease 1928). Their spatial distribution and their high velocity dispersion indicate that most of the bulge PNe are old objects. Their strong H$`\alpha `$ and \[OIII\] emission lines make their velocities easy to measure. We have thus decided to use the PNe as probes to study the kinematics and dynamics of the Galactic bulge.
The distances of PNe are still poorly known. Using the optical diameter as a distance criterion is not adequate because PNe have a wide range of absolute diameters. Nevertheless, using the angular diameters, spatial distribution and radial velocities of a sample of PNe, Gathier et al. (1983) estimated that probably 80% of the small (diameter $`<20^{\prime \prime }`$) PNe within $`10^{}`$ of the Galactic center belong to the bulge. While it is clear that most of the PNe towards the bulge are associated with the bulge, it is also evident that their apparent spatial distribution at low Galactic latitudes is affected by the interstellar absorption.
### 3.1 The Data
In 1994 and 1995, we conducted an H$`\alpha `$ imaging survey of the Galactic bulge in order to detect new PNe (Beaulieu et al. 1999). The survey yield 56 new and 45 already catalogued PNe. We obtained radial velocities for each new PNe plus a sample of 317 catalogued PNe (i.e. 272 catalogued and the 45 rediscovered PNe) taken from the Strasbourg-ESO Catalogue of Galactic Planetary Nebulæ (Acker et al. 1992). Although we intended to observe only the southern part of the bulge (less affected by extinction), we have obtained a few fields in the northern part as well. Our data have already been used in a study of Galactic kinematics by Durand et al. (1998).
Our database of PNe contains two samples. The first sample comprises the 97 PNe (new and rediscovered) found in the southern bulge from our uniform survey with the 1.0m telescope. The region covered by this survey is $`20^{}<l<20^{}`$ and $`5^{}>b>10^{}`$. We will refer to this uniformly selected sample as the Survey fields only sample. A note is needed here: this sample, in fact, contains 98 PNe but we are using 97 PNe for the analysis. The reason for this is that we accepted one PN as “probable” after we have completed the Survey fields only sample analysis. This PN is SB15 : $`\mathrm{PNG009}.306.5`$.
The second sample is less homogeneous, with the 98 PNe Survey fields only sample (including, this time, SB15), the 3 PNe which we discovered in the northern bulge, and the 272 PNe from the Acker et al. (1992) catalogue for which we have measured new radial velocities. This larger sample contains 373 PNe and covers the more extended region $`30^{}<l<30^{}`$ and $`3.3^{}<|b|<15^{}`$. We will refer to this sample as the Survey fields \+ Catalogue sample. Figure 2 shows the ($`l,b`$) distributions for the two samples.
In the absence of information on distances for our PNe, we made no attempt to separate disk and bulge PNe in our two samples. Therefore, disk contamination is likely. We note, however, that some of the dynamical models used in this study (see §6) include a disk.
## 4 Analysis
In the first part of this section, we present several plots showing the kinematics of these two samples for preliminary visual assessment. We then go on to compare the properties of the PNe samples with the properties of several recent dynamical models. This comparison will be first presented visually in the form of plots of individual velocities, mean velocities and velocity dispersions against $`l`$ and $`b`$. Then we will use a statistical technique by Saha (1998) to make a more quantitative comparison of the data with the models, and to estimate the Galactic scaling parameters and orientations which best match the models to our data.
The typical radial velocity error for our PNe is 11 km s<sup>-1</sup> (Beaulieu et al. 1999). For the Galactic bulge, the velocity dispersion ranges from about 60 km s<sup>-1</sup> to 125 km s<sup>-1</sup> (Fig. 13), so this radial velocity error is negligible.
In the presentation of the kinematics of our samples, in order to illustrate the systemic rotational properties of the bulge PNe more clearly, we will show the velocities of the PNe corrected for the solar reflex motion. We adopted the circular velocity of the Local Standard of Rest (LSR) at the Sun as 220 km s<sup>-1</sup> (Kerr et al. 1986). For the Sun’s peculiar velocity relative to the LSR we use 16.5 km s<sup>-1</sup> towards $`l=53^{}`$, $`b=25^{}`$ (e.g. Mihalas & Binney 1981). The corrected line-of-sight ($`V_{los,GC}`$) velocity (i.e. the line-of-sight velocity in km s<sup>-1</sup> that would be observed by a stationary observer at the location of the Sun) is then given by
$$V_{los,GC}=V_{obs}+220\mathrm{sin}l\mathrm{cos}b+16.5[\mathrm{sin}b\mathrm{sin}25+\mathrm{cos}b\mathrm{cos}25\mathrm{cos}(l53)]$$
where $`V_{obs}`$ is the heliocentric observed line-of-sight velocity in km s<sup>-1</sup>.
Figure 3 shows the longitude versus velocity diagram for the Survey fields only (top panel) and Survey fields \+ Catalogue (lower panel) (corrected for the solar reflex motion).
In the figures that follow, we note that there must be some level of distance bias in our PNe samples. The longitude distribution of the PNe shows some evidence for depletion at $`l<0`$ (the more distant side of the bar) relative to $`l>0`$ (Fig. 14), although this depletion is only marginally significant (Fig. 16). In the comparisons of the PNe distribution and kinematics with the various models (§6), we will ignore this distance bias.
### 4.1 Survey fields only
Figure 4 shows the longitude versus mean velocity (top panel) and the longitude versus velocity dispersion (lower panel) using 8 bins in longitude, with approximately equal numbers of PNe in each bin (12 to 13 PNe). The rotation of the bulge is clearly seen, with an amplitude of about $`\pm 100`$ km s<sup>-1</sup>. The velocity dispersion of the bulge is approximately constant with longitude, except for the apparent drop in $`\sigma `$ for $`l>+12^{}`$. This drop is seen again in the larger sample described in §4.2 but on both sides of the Galactic center: see Figure 7. It is probably due to the contribution of the inner disk at these longitudes (see Lewis & Freeman 1989).
Figure 5 shows the latitude versus mean velocity (top panel) and the latitude versus velocity dispersion (lower panel) for 2 bins with equal number of PNe in latitude. Each bin in latitude contains 48 to 49 PNe. We see that the total velocity dispersion about the mean velocity does not appear to change significantly with latitude. \[Note that this total velocity dispersion in the plots against latitude includes the systemic rotation and random velocities of the stars.\]
Tables 1 and 2 summarize the binned data shown in Figure 4 and 5 respectively. Column 1: the mean latitude and longitude, Column 2: the mean velocity (km s<sup>-1</sup>), Column 3: the velocity dispersion, Column 4: the error (standard deviation) in the mean velocity, and Column 5: the error (standard deviation) in the velocity dispersion.
We have also divided the $`lV`$ diagram of Figure 3 (top panel) into two bins in Galactic latitude (with 48 to 49 PNe in each bin) (Figure 6) in order to see if contamination from disk PNe is affecting our data. Disk contamination is potentially more serious at higher latitudes because of the steeper density gradient of the bulge. Therefore, if contamination were present, we would expect the lower latitude bin ($`b=04.9^{}`$ to $`06.5^{}`$) to be significantly hotter (i.e. have higher velocity dispersion) than the higher latitude bin ($`b=06.6^{}`$ to $`10.2^{}`$). We see no evidence in Figure 6 for serious disk contamination in our sample, except possibly for $`l>+12^{}`$.
### 4.2 Survey fields and Catalogue objects
Now we present the data for the larger and more extended but less homogeneous Survey fields \+ Catalogue sample of 373 PNe. (see Figure 3 (lower panel)). Figure 7 shows the longitude versus mean velocity (top panel) and the longitude versus velocity dispersion (lower panel) using 12 bins in longitude with approximately equal numbers (31 to 32) of PNe in each bin. Again, the rotation of the bulge is clearly seen. For $`|l|>12^{}`$, the mean rotational velocity continues to rise as the data become dominated by PNe of the inner disk. In this larger sample, beyond $`|l|>12^{}`$, we see again an apparent drop in the velocity dispersion, due presumably to the contribution of the inner disk PNe at these longitudes.
Figure 8 shows the latitude versus mean velocity (top panel) and the latitude versus total velocity dispersion (lower panel) using 6 bins in latitude with approximately equal numbers (62 to 63) of PNe in each bin; 2 bins are in the northern bulge and 4 bins in the southern bulge. \[Note again that the total velocity dispersion in the latitude plots includes the systemic rotation and random velocities of the stars.\]
Tables 3 and 4 summarize the binned data shown in Figure 7 and 8 respectively. Column 1: the mean latitude and longitude, Column 2: the mean velocity (km s<sup>-1</sup>), Column 3: the velocity dispersion, Column 4: the error (standard deviation) in the mean velocity and Column 5: the error (standard deviation) in the velocity dispersion.
In Figure 9, we are looking again at the disk contamination using the same latitude bins as for the Survey fields only sample. In this larger and more extended sample, contamination from the disk PNe becomes evident outside the longitude region $`|l|=12^{}`$ where the PNe velocity distribution becomes significantly colder.
We also present a series of longitude-velocity diagrams for 6 bins in latitude. Figure 10: $`b=+03.3^{}`$ to $`+05.2^{}`$ (top panel) and $`b=+05.2^{}`$ to $`+15.1^{}`$ (lower panel). Figure 11: $`b=03.3^{}`$ to $`04.4^{}`$ (top panel) and $`b=04.5^{}`$ to $`05.8^{}`$ (lower panel). Figure 12: $`b=05.8^{}`$ to $`07.4^{}`$ (top panel) and $`b=07.4^{}`$ to $`14.9^{}`$ (lower panel). For this less homogeneous (and generally brighter) sample of PNe, the disk contamination really starts to show in the two high latitude bins: the PNe velocity dispersion becomes much colder at all longitudes, as we would expect to see if the disk contamination is significant at higher latitudes.
### 4.3 Comparison with other studies
In recent years, there have been some important studies of the kinematics of K and M giants in the Galactic bulge. Although the regions observed are mostly not as extended as our survey, we should now compare the kinematics derived from these studies with the results from the PNe.
In Figure 13, we show again the mean velocity and velocity dispersion against longitude for our extended sample, and have overplotted data from kinematic studies of giants, which fall in our Survey fields \+ Catalogue sample region. Minniti (1996a) presented data for three bulge fields. He gives kinematical data for the more metal-rich (\[Fe/H\] $`>1`$: filled symbols) and metal-poor stars (\[Fe/H\] $`<1`$: open symbols) separately. Data for one field come from Harding and Morrison (1993), and again we show the data points for the more metal-rich and more metal-poor stars separately. For Baade’s Window (K giants: Terndrup et al. 1995), we show the only available data point, the velocity dispersion value, for his stars with V $``$ 16.0: these fainter stars are likely to be a relatively uncontaminated sample of bulge stars. Sharples et al. (1990) find an almost identical dispersion for their M giants in Baade’s Window. Finally, we present three data points (higher latitude \[$`b=12^{}`$\]) from Ibata & Gilmore (1995a,b). We derived equivalent $`<V_{los,GC}>`$ values for their three negative longitude fields from the gradients $`\mathrm{\Omega }_G`$ that they estimated assuming an isotropic velocity dispersion. We used the formalism of Morrison et al. (1990), assuming that the stars in each field lie where the line-of-sight passes closest to the center of the bulge. Ibata and Gilmore give kinematical solutions for several assumptions about the shape of the bulge velocity ellipsoid $`\sigma `$. The derived $`<V_{los,GC}>`$ values depend very weakly on the assumptions about the shape of $`\sigma `$, so we have only plotted the isotropic solution (asterisks) in Figure 13 (upper panel). Their velocity dispersions are more sensitive to the shape of $`\sigma `$. We show their velocity dispersions for an isotropic bulge (asterisks) in Figure 13 (lower panel). The isotropic solution for $`\sigma `$ appears to give better agreement of the Ibata and Gilmore data with the other bulge samples. For comparison we also show (line) the slope of the linear rotation curve found for 279 bulge PNe by Durand et al. (1998). (Part of our data is included in their analysis.) The slope of this line is 9.9 $`\mathrm{km}\mathrm{s}^1\mathrm{degree}^1`$. Table 5 summarizes the symbols associated with each study. Column 1: the study, Column 2: the field’s ($`l,b`$) coordinates, and Column 3: the symbol used on the plot.
For Minniti’s three fields, the data for the metal-rich giants clearly matches our PNe data better than do the metal-poor giants. For the Harding-Morrison field, although we see the same match of the metal-rich giants with our PNe in the velocity dispersion, it is in fact the opposite that is seen for the mean velocity. This disagreement was also observed by de Zeeuw (1993) when he compared the Minniti and Harding-Morrison samples with Kent’s model (Kent 1992). Nevertheless, the otherwise good agreement seen so far identifies the bulge PNe with the more metal-rich giant (\[Fe/H\] $`>1`$) of the bulge, as we would expect. For Baade’s Window, it is interesting to see that the velocity dispersion is perhaps somewhat higher than the mean of the velocity dispersion values for our PNe at lower $`|l|`$, but we note that our PNe are mostly more distant from the Galactic plane than Baade’s Window (cf Figure 8). (The velocity dispersion along the minor axis of the Galactic bulge is known to decrease with increasing $`|b|2^{}`$: e.g. Rich 1996.)
We note that the $`<V_{los,GC}>`$ values shown in Figure 13 for the Ibata and Gilmore sample pertain to their more metal-poor stars with \[Fe/H\] $`<0.5`$. The shallower slope of the $`<V_{los,GC}>l`$ relation for their stars is consistent with the metallicity trends seen in the Minniti and Harding-Morrison samples.
## 5 Comparison with COBE images
The COBE/DIRBE images in the 1.25, 2.2 and $`3.5\mu `$m wavelength regions allow us to compare the distribution of the PNe with the integrated near-infrared emission from the Galactic bulge. In this region of the spectrum, the light distribution comes from various stellar populations but is dominated by the more metal-rich K and M-giants which have kinematics similar to those of the PNe. The $`1.25\mu `$m map also gives an indication of the distribution of the dust.
### 5.1 Histogram of the Longitude Distributions
Figure 14 shows a histogram of the longitude distribution of the COBE light and the PNe in our southern surveyed fields ($`5^{}>b>10^{}`$). The COBE histograms were constructed from the COBE light distribution within the individual 30 arcmin fields used for the PNe survey (see Beaulieu et al. 1999), so the distributions are directly comparable. The dashed lines represents the three bands (1.25, 2.2 and $`3.5\mu `$m) of the COBE light distribution and the solid line represents the PNe distribution. We see immediately that the three COBE distributions agree very well and that the PNe distribution follows the COBE light distribution. The fact that the three COBE light distributions agree so well is an indication that extinction, in our surveyed fields, is not severe and that its distribution is fairly uniform.
We also compare the three COBE light distributions with their cumulative distributions, in preparation for the next section. Figure 15 shows that the cumulative distributions for the three COBE bands are very similar.
### 5.2 K-S Test
The Kolmogorov-Smirnov (K-S) test estimates the probability that a set of observed values can be excluded as coming from a given specified distribution.
We performed a one-sample, two-tailed K-S test in our surveyed fields (Galactic longitude $`l=+20^{}`$ to $`20^{}`$ and Galactic latitude $`b=5^{}`$ to $`10^{}`$) using the well-determined COBE light distribution in longitude as the specified distribution and our sample of PNe as the observed distribution. The test uses the largest value $`D`$ of the deviation $`|F_0(X)S_N(X)|`$ where $`|F_0(X)`$ and $`S_N(X)|`$ are the cumulative distributions of the specified distribution (the COBE light distribution) and the set of observed values (the longitude distribution of our PNe counts). We have seen in Figure 15 that the cumulative distributions of the three COBE colors agree very well and the results for the maximum deviation will be similar in all three colors.
Figure 16 shows the two cumulative distributions for our PNe sample and the $`2.2\mu `$m COBE light. The ordinate, N, has been normalized to 1.0 for both distributions. We have used table E of Siegel (1956) to estimate the probabilities. Table 6 gives the results obtained for the maximum deviation $`D`$ and the associated probability that the deviation $`D`$ could occur by chance from the same parent distribution. Column 1: the COBE band, Column 2: the maximum deviation value $`D`$ and and Column 3: the associated probability of occurrence. This probability is between 0.23 and 0.30, and we conclude that there is no significant difference between the longitude distribution of the PNe and the COBE light in the zone of our deep survey.
## 6 Comparison with models
The evidence for a bar in our Galaxy initiated much interest in developing detailed dynamical models (N-body and Schwarzschild) of the Galactic bar-bulge. Several different kinds of models are now available, but the observational constraints on their stellar dynamics are still weak. Our kinematical data for the PNe of the Galactic bar-bulge provide further constraints on the models.
In this section, we present the data of our survey with velocities relative to the LSR, using the parameters for the sun’s peculiar motion as given in the equation in §4. The motivation for doing so is that most observational studies are presented in that manner and it would therefore be easier for future comparison. Also, we will use our data to estimate the best value of the tangential velocity of the LSR for each model.
### 6.1 Presentation of the models
At the time of conducting this study, there were four triaxial numerical models available to study the dynamics of the Galactic bulge. They offer interesting and different approaches to studying the formation and structure of the bar-bulge. There are three N-body models (Sellwood 1993; Fux 1996; Kalnajs 1996) and one Schwarzschild model with an N-body realization (Zhao 1996). (Very recently, a more elaborate Schwarzschild model has appeared (Häfner et al. 1999), constrained by a subset of the data in Figure 13 plus some proper motions.)
#### 6.1.1 Sellwood’s model
Sellwood’s model is one of the earlier N-body dynamical models. It is a purely stellar N-body system with $`5\times 10^4`$ particles. It starts from a Q = 1.2 axisymmetric Kuz’min-Toomre disk which contains 70% of the total mass. The remaining 30% is in a rigid Plummer sphere which has half the scale length of the disk. The bar-bulge forms through the instability of the disk. The resulting model shows a peanut-shaped bulge. At a viewing angle of 30 to the major axis and a finite distance from the center, the model shows an asymmetry in longitude between the positive and negative sides, which is consistent with the one seen in the COBE/DIRBE image (Weiland et al. 1994).
Figure 17 presents the face-on view (XY) and the edge-on view (YZ) as seen from infinity, with the Sun-center line at an angle of 30 from the major axis.
#### 6.1.2 Fux’s model
Fux’s model is an N-body system of stars. It has four components: an exponential stellar disk of constant thickness ($`15\times 10^5`$ particles), a composite power-law stellar nucleus-spheroid ($`5\times 10^4`$ particles), a dark halo ($`2\times 10^5`$ particles), and a dissipative gas component (a smoothly truncated Mestel disk with $`2\times 10^4`$ particles). The system starts in equilibrium and the rotating bar forms through instabilities. The model provided to us by Fux is a gas-free version which has evolved for 5 Gyr: we note that Fux (1997) has built more elaborate models of the Milky Way including gas, which we have not considered here.
Figure 18 presents the face-on view (XY) and the edge-on view (YZ) as seen from infinity, with the Sun-center line at an angle of 30 from the major axis.
#### 6.1.3 Kalnajs’ model
Kalnajs has been conducting numerical experiments on thin self-gravitating disks which turn into triaxial rotating objects because of buckling instabilities. The projected shapes of these objects, when viewed from the right distance and orientation, resemble the light distribution of the Galactic bulge, and the line-of-sight velocities can be scaled to match observed motions of planetary nebulae in the bulge. The experiments use only 8000 particles, but since the triaxial objects appear to be stationary in a rotating frame, one can add the distributions at different times and obtain models containing effectively $`10^5`$ particles.
Figure 19 presents the face-on view (XY) and edge-on view (YZ) as seen from infinity, with the Sun-center line at an angle of 45 from the major axis.
#### 6.1.4 Zhao’s model
The last is a model of the COBE bar, constructed from 10K orbits (direct, retrograde and chaotic) in the rotating bar potential plus a rigid Miyamoto-Nagai disk potential, using the non-negative least square fitting technique pioneered by Schwarzschild.
The model provided by Zhao for our comparison is the system allowed to evolve as an N-body system after 10 rotations and it contains 32634 particles.
Figure 20 presents the face-on view (XY) and the edge-on view (YZ) as seen from infinity, with the Sun-center line at an angle of 20 from the major axis..
Table 7 summarizes the parameters suggested by the authors of each model. Column 1: the model, Column 2: the total number of particles in the model, Column 3: the solar galactocentric radius $`R_{}`$ (in model units), Column 4: the viewing angle $`\varphi `$ (in degrees) of the bar. $`\varphi `$ is the angle between the major axis of the bar and the Sun-center line, and is taken as positive for a bar pointing into the first quadrant of l, Column 5: the velocity scale $`V_{scale}`$ (km s<sup>-1</sup> per model unit) of the model, and Column 6: the solar tangential velocity $`V_{,T}`$ (km s<sup>-1</sup>).
In the next section, we use a statistical technique to estimate these scaling parameters for each model from our data.
### 6.2 Search for best parameters
The authors of each model have suggested values for the Sun’s galactocentric distance (in model units), the viewing angle of the bar and a velocity scale (Table 7). However, by varying these parameters, we may hope to obtain somewhat better fits to the present data. There are four parameters one can vary: (i) the overall spatial scale of the model, or equivalently $`R_0`$ in model units; (ii) the overall velocity scale; (iii) the viewing angle of the bar; and (iv) the tangential velocity of the LSR. Saha (1998) has developed a method for searching the space of these four parameters for values which are most likely to have given rise to the observed data. We used his code, which gives a median fit for the four parameters and error bar estimates under the assumption that the models and the data are drawn from the same underlying distribution function.
We are going to compare the positions and the radial velocities of the 97 PNe from the Survey fields only sample with those of the four models. We choose to restrict the comparison to the 97 PNe in our survey region, because they were selected in a homogeneous manner. In making our comparison we must only use that part of the model which would fall into our surveyed window. Since our window lies several scale lengths below the Galactic plane, only a small fraction of the model particles are used in the comparison. The number of model particles is held fixed as the observer’s position changes: the respective numbers for Sellwood, Fux, Kalnajs and Zhao were 400, 6000, 9000 and 1700.
We use Saha’s procedure to make a quantitative comparison of the $`(l,b,V_{los})`$ distributions for samples of observed objects and N-body models. Saha’s statistic is
$$W=\underset{i=1}{\overset{B}{}}\frac{(m_i+s_i)!}{m_i!s_i!}$$
where the $`(l,b,V_{los})`$ space has been partitioned into a total of $`B`$ cells, $`m_i`$ and $`s_i`$ are the numbers of model and sample objects in the $`i`$-th cell; $`W`$ is proportional to the probability that both the observed sample and the model come from the same underlying (but unknown) distribution, so $`W`$ can be used to compare the goodness of fit of various models. As described above, the $`W`$ statistic also serves to estimate the scaling parameters for each model from the observed sample. (see Sevenster et al. 1999 for a previous application of this statistic.)
For choosing the number $`B`$ of cells, our guideline is that the average number of model particles per cell should be 5 or more, and the spatial cells should not be smaller than important features in the distribution function, such as the scale height (see Saha 1998 for more discussion). After some experimentations, we used a total of 260 cells in $`(l,b,V_{los})`$: 13 in $`l`$, 2 in $`b`$ and 10 in $`V_{los}`$. Table 8 presents the results: Column 1: the model, Column 2: the total number of particles in our window, Columns 3: the four parameters (i) the solar orbit radius $`R_{}`$ in model units, (ii) the orientation angle $`\varphi `$ in degrees, (iii) the velocity scale (in km s<sup>-1</sup> per model unit), and (iv) the solar tangential velocity (in km s<sup>-1</sup>), Column 4: give the median and the 90 % confidence limits for these parameters.
These results are produced by the program after searching through the region of parameter space given by $`7<R_{}<9`$, $`0^{}<\varphi <90^{}`$, $`200<V_{,T}<240`$ and $`0<V_{scale}/V_{scale,model}<2`$, where $`V_{scale,model}`$ is the suggested velocity scale value from each model (see table 7). This choice of search region was partly guided by the likely values of the corresponding galactic parameters and appears to be satisfactory: for every parameter and every model, the median estimate of the parameter lies away from the boundary of the search region by at least the 90% confidence limit. Of these parameters, $`V_{scale}`$ is the least constrained by the data, and $`V_{,T}`$ the best constrained. For all of the models, the $`W`$ statistic indicates that the probability that the models and data come from the same underlying distribution exceeds 98%.
### 6.3 Models versus Data
With the estimated parameters given in Table 8, we now present some visual comparisons of the kinematics of the models and the data. The figures are similar to those shown earlier for our data alone, except for the fact that the PNe velocities and the velocity data for the models are heliocentric.
Figures 21 to 24 present the longitude-velocity diagrams for Sellwood (400 particles), Fux (6000 particles), Kalnajs (9000 particles) and Zhao (1700 particles) respectively.
Figures 25 to 28 show $`<V_{los}>`$ (top panel) and $`\sigma `$ (lower panel) against the longitude for the data and models, with the model represented by thick lines.
The main features of the ($`<V_{los}>`$,$`\sigma `$) versus longitude relations are that all models give a fair representation of the observed $`<V_{los}>l`$ distribution, but the Sellwood and Zhao models have a velocity dispersion that is relatively low. The $`V/\sigma `$ values for the Sellwood and Zhao models appear to be somewhat higher than for the bulge of the Galaxy, at least in the region of our survey. But, as indicated by the Saha procedure, all of the models are good representations of the PNe data.
### 6.4 Models versus Models
We attempted to use the program to discriminate between the models by intercomparing the maximum $`W`$ value from Saha’s procedure for samples of similar total numbers of particles. For example, Sellwood’s model has 400 particles in our survey region, so we estimated values of $`W`$ for Sellwood’s model and random samples of 400 particles drawn from the larger simulations (Zhao, Fux, Kalnajs) within our region. Similarly, Zhao’s model has 1700 particles within our region, so we compared W values for Zhao’s model and random samples of 1700 particles from the larger simulations (Fux, Kalnajs). Table 9 presents values of $`\mathrm{ln}W`$ for each set of comparisons.
The total number of particles is shown for each comparison. The sampling standard deviation of $`\mathrm{ln}W`$ is derived empirically by the program. Table 9 shows: Column 1: the models being compared and Column 2: the value of $`\mathrm{ln}W`$ for each set of comparisons, i.e. 400, 1700 and 6000 particles. The last line of Table 9 shows the sampling standard deviation of $`\mathrm{ln}W`$ of each run. We recall that for all of the models, the probability that the models and the data come from the same underlying distribution is more than 98 %.
We see that the values of $`\mathrm{ln}W`$ for each model do not differ by more than about $`1.9\sigma `$, indicating again that there is no significant difference between the ability of the various N-body models to represent our data. Table 9 shows that Sellwood’s model comes out best in the N=400 comparison of the four models, despite the apparently large deviations in the velocity dispersion (Figure 25). We recall that the $`W`$-statistic involves comparison of data and model over cells in velocity and $`(l,b)`$. The quality of the velocity comparisons is seen in Figures 25-28. Figure 29 shows the cumulative distributions over $`l`$ of the four models (all with N=400) and the survey fields only PNe sample (over the same interval in $`b`$). Sellwood’s model lies closest to the data in Figure 29, followed by Zhao’s model. This help to understand the ordering of the $`\mathrm{ln}W`$ values for the models as given in Table 9.
## 7 Summary and Conclusions
Planetary Nebulæ are good tracers for a dynamical study of the Galactic bulge because they are less affected by metallicity bias than most other tracers and they are strong emitters in H$`\alpha `$ \- this make their velocities easy to measure. We chose to survey the southern Galactic bulge in the region $`l=\pm 20^{}`$ and $`b=5^{}`$ to $`10^{}`$ because of its lower extinction relative to the northern bulge.
We compared the longitude distribution of PNe in our surveyed fields with the COBE light distribution at 1.25, 2.2 and $`3.5\mu `$m. We conclude that (i) the light distributions in the three COBE bands agree very well, indicating that the extinction in our surveyed fields is not severe and that its distribution is fairly uniform and, (ii) there is no significant difference between the longitude distribution of the PNe and the COBE light in the zone of our deep survey.
Recent studies of stellar kinematics in a few clear windows in the Galactic bulge have provided mean velocities and velocity dispersions which can be compared with our data. We thus compared data from Minniti (1996a), Harding and Morrison (1993), Terndrup et al. (1995) and Ibata and Gilmore (1995a,b), and found that the metal-rich stars in Minniti’s three fields agree very well with our data. Harding and Morrison’s metal-rich stars agree well with our velocity dispersion data for the PNe, but not so well with our mean velocity. We also found that the velocity dispersion in Baade’s Window (Terndrup et al. 1995) is somewhat higher than ours near $`l=0`$ but note that Baade’s window is closer to the Galactic plane than most of our PNe. For the Ibata and Gilmore data, the velocity gradient over $`l`$ is shallower than for the PNe and other samples of giants; this is presumably due to their restriction to more metal-weak giants (\[Fe/H\]$`<0.5`$). Their velocity dispersion estimates for an isotropic bulge agree better with the PNe values.
To assist in the comparison of the four N-body models with our sample of data, we used a procedure proposed by Saha (1998) to make a quantitative comparison of the $`(l,b,V_{los})`$ distributions for samples of observed objects and N-body models. The main conclusion from this comparison is that all four models show a fairly good fit to our data.
Sellwood, Fux and Kalnajs’ models are all bar-forming systems via the instabilities of a disk and, after scaling, are kinematically more or less similar. Zhao’s model is constructed to fit the COBE light: in this sense, it is a step up from Kent’s (1992) axisymmetric model for the Spacelab near-IR photometry. Kent’s predicted velocity dispersion, as quoted by de Zeeuw (1993), was already a fairly good fit to the existing data; therefore, it is not surprising that Zhao’s model should also fit well.
Using the estimated parameters obtained from Saha’s procedure, we made some visual comparisons of the kinematics of the models and the data. The Kalnajs and Fux models give a good visual representation of the mean velocity and velocity dispersion of the bulge in our survey region; the Sellwood and Zhao models represent the mean velocity well but their velocity dispersion is marginally low relative to the PNe observations.
It will be interesting to use our PNe sample as a more detailed kinematical test of Kuijken’s axisymmetric isotropic two-integral model. One important goal of this comparison would be to look for kinematical disagreements between the data and the axisymmetric model that might be kinematical signatures of triaxiality. In the same spirit, it would be interesting to compare Kuijken’s model in detail with the numerical triaxial systems discussed in §6.
We saw earlier that Minniti’s data (Minniti et al. 1992) is apparently not consistent with Kuijken’s model. As a preliminary comparison with Kuijken’s model, we examined the distribution of LSR velocities for our Survey fields \+ Catalogue sample. Figure 30 shows a histogram of LSR velocities for the PNe with $`5^{}<l<10^{}`$. We can compare the velocity distribution in the region $`5^{}<l<10^{}`$ with the distribution measured by Minniti et al (1992) for the giants towards $`l=8^{},b=7^{}`$ and discussed by Kuijken (1995). The PNe in our $`5^{}<l<10^{}`$ region cover a larger region of sky than the Minniti sample; however, the mean value of $`|b|`$ for the PNe is about $`6^{}`$, so we might expect the velocity distributions of the PNe and the giants to be at least qualitatively similar. We see from Figure 30 that the velocity distribution in the region $`5^{}<l<10^{}`$ is asymmetric; the asymmetry is in the opposite sense to that found by Minniti et al. (1992) but closer to that seen for the more metal-rich stars in Minniti’s later (1996b) study for this field. The mean LSR velocity for our sample is $`36\pm 11`$ km s<sup>-1</sup>, compared with $`5\pm 10`$ km s<sup>-1</sup> for the Minniti et al. (1992) sample and the predicted value of 32 km s<sup>-1</sup> for the Kuijken model. There seems to be better agreement between the PNe and the two-integral model in this region than was found between the model and the giants.
It may be that we are seeing an effect of metallicity in the Minniti et al. (1992) sample. There was no information about metallicity at that time so the sample could be suffering from pollution by the more slowly rotating metal-poor stars. We recall here that the metal-rich stars in Minniti’s three fields (Minniti 1996a) are in good agreement with our Survey fields \+ Catalogue sample (cf Figure 13).
So far, only a few clear Galactic bulge windows have been extensively studied. Although these studies provide important information on the kinematics in the bulge, their small region do not give us the entire picture of the bulge kinematics. Two major studies of tracers in the Galactic bulge, the K giants (Harding & Morrison 1993) and the OH/IR stars (Sevenster et al. 1997a,b) and a new PNe H$`\alpha `$ survey of the Southern Galactic Plane (Parker & Phillips 1998) are presently under way. A comparison of the PNe surveys with the results coming from the OH/IR and K giants large-scale surveys should clearly indicate any dynamical differences between the populations from which these different tracers come.
Finally, we conclude that the existing studies give a more or less consistent picture of the kinematics of the Galactic bulge, as summarized in Figure 13, at least for the metal-rich bulge tracers. We find it interesting that the N-body models, in which the bar/bulge grows from the disk via bar-forming instabilities, give a good representation of the detailed stellar kinematics of the bulge.
SFB wish to acknowledge funding from The Australian Government through an Australian National University Scholarship and an Overseas Postgraduate Research Scholarship. We are most grateful to Roger Fux and Jerry Sellwood for allowing us to use their models. Our thanks go to Gerry Gilmore for critical comments on the manuscript.
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# 1 Introduction
## 1 Introduction
The production of scalar mesons in radiative decays is a valuable source of information on hadron spectroscopy. It was argued that the branching ratio for $`\varphi \gamma f_0(980)`$ can be used to make a unique choice among different models of $`f_0(980)`$: a conventional quark-antiquark state, an exotic $`qq\overline{q}\overline{q}`$ state , and a $`K\overline{K}`$ molecule . Such a possibility to resolve the long debated problem of the $`f_0(980)`$ structure using a single partial width looks very attractive. However, it was argued in that the dependence of the theoretical predictions for the $`\varphi \gamma f_0(980)`$ decay on the $`f_0(980)`$ structure is partly due to differences in modeling. Recent measurements of the radiative decays $`\varphi \gamma \pi ^0\pi ^0`$ by SND and $`\varphi \gamma \pi ^+\pi ^{},\gamma \pi ^0\pi ^0`$ by CMD-2 in Novosibirsk have made it possible to confront alternative models of the light scalar mesons with experimental data.
The goal of this paper is to reanalyze the calculation of the decay width for $`\varphi \gamma f_0(980)`$ in coupled channel models where the $`f_0(980)`$ state arises as a dynamical state (a $`K\overline{K}`$ molecule which may also have a substantial admixture of a quark-antiquark component). Since the $`\varphi `$ meson is nearly a pure $`s\overline{s}`$ state, the decay $`\varphi \gamma \pi \pi `$ is an OZI–rule violating process which is expected to proceed via a two–step mechanism with intermediate $`K\overline{K}`$ states. Therefore this decay is well suited for probing the $`K\overline{K}`$ content of the scalar mesons. The $`K\overline{K}`$ molecular state was originally proposed in the potential quark model . A dynamical state close to $`K\overline{K}`$ threshold is also introduced in the coupled channel models and in the meson exchange interaction models . A state strongly coupled to the $`s\overline{s}`$ and $`K\overline{K}`$ channels near the $`K\overline{K}`$ threshold was as well found in the unitarized quark model . The coupled-channel model derived from the lowest order chiral Lagrangian produces a scalar state dominated by the $`K\overline{K}`$ channel. General discussions of the nature of the $`f_0(980)`$ can be found in and references therein.
Our approach is based on a coupled channel model (CCM) for the $`\pi \pi `$ and $`K\overline{K}`$ systems which is similar to the one studied in . The calculation of the decay $`\varphi \gamma f_0(980)`$ for point-like particles is summarized in Sec. 2. The details of the coupled channel model are given in Sec. 3, and the model parameters are determined from a fit to the $`\pi \pi `$ scattering data. The reaction $`\varphi \gamma \pi \pi `$ in a CCM framework is studied in Sec. 4. The analytic structure of the $`\pi \pi `$ and $`K\overline{K}`$ scattering amplitudes is investigated in Sec. 5. The mixing between the two–meson and $`q\overline{q}`$ channels is discussed in Sec. 6. In Sec. 7 the physical properties of the scalar mesons in the model proposed are discussed and compared to other approaches in the literature. The details of the formalism are collected in the Appendices.
## 2 The Decay $`\varphi \gamma f_0`$
For the benefit of the reader, we begin with a brief summary of the formulas describing radiative transitions between vector and scalar states. The amplitude of the radiative $`\varphi `$ decay into the scalar meson $`f_0`$ has the following structure which is imposed by gauge invariance:
$`M(\varphi \gamma f_0)`$ $`=`$ $`ϵ_\varphi ^\mu ϵ_\gamma ^\nu (p_\nu q_\mu g_{\nu \mu }(pq))H(p^2,(pq)^2)`$ (1)
where ($`ϵ_\varphi `$, $`p`$) and ($`ϵ_\gamma `$, $`q`$) are the polarizations and four-momenta of the $`\varphi `$ and $`\gamma `$, correspondingly, and $`H(p^2,(pq)^2)`$ is the scalar invariant amplitude which depends on the invariant masses of the initial and final mesons. The polarization vectors satisfy the constraints $`ϵ_\varphi p=0`$ and $`ϵ_\gamma q=0`$. Using the three-dimensional gauge $`ϵ_\gamma =(0,𝜺_\gamma )`$ in the center-of-mass system (CMS) one gets the amplitude
$`M(\varphi \gamma f_0)`$ $`=`$ $`(𝜺_\varphi 𝜺_\gamma )m_\varphi \omega H(p^2,(pq)^2)`$ (2)
where $`m_\varphi `$ is the $`\varphi `$ mass, $`\omega `$ is the photon energy in the CMS, and $`𝜺_\varphi `$ and $`𝜺_\gamma `$ are the $`\varphi `$ and $`\gamma `$ three–dimensional polarization vectors in the CMS, respectively.
For point-like particles, the amplitude $`M(\varphi \gamma f_0)`$ is given by the diagrams in Fig.1 (see and references therein). While the diagram Fig.1(a) corresponding to loop radiation is overall logarithmically divergent, its contribution to the $`p_\nu q_\mu `$ term in Eq.(1) is finite. Gauge invariance enforces the appearance of the seagull diagram Fig.1(b) which contributes only to the $`g_{\nu \mu }(pq)`$ term in Eq.(1). Since the sum of the loop–radiation and seagull terms is gauge invariant, one obtains the total amplitude $`M(\varphi \gamma S)`$ by calculating only the $`p_\nu q_\mu `$ term of the loop radiation. The result for the scalar invariant amplitude is
$`H(p^2,(pq)^2)`$ $`=`$ $`{\displaystyle \frac{eg_\varphi g_{f_0KK}}{2\pi ^2m_K^2}}I(a,b)`$ (3)
$`a={\displaystyle \frac{m_\varphi ^2}{m_K^2}}`$ , $`b={\displaystyle \frac{m_{f_0}^2}{m_K^2}}`$ (4)
where $`g_\varphi `$ and $`g_{f_0KK}`$ are the $`\varphi K^+K^{}`$ and $`f_0K^+K^{}`$ coupling constants and the function $`I(a,b)`$ is defined in Appendix A. The radiative decay width is
$`\mathrm{\Gamma }(\varphi \gamma f_0)`$ $`=`$ $`{\displaystyle \frac{\omega ^3|H(p^2,(pq)^2)|}{12\pi }}=`$ (5)
$`=`$ $`{\displaystyle \frac{\alpha g_\varphi ^2g_{f_0KK}^2}{3(2\pi )^4}}{\displaystyle \frac{\omega }{m_\varphi ^2}}|(ab)I(a,b)|^2.`$ (6)
The generalization of Eqs. (1,3,6) to the case $`\varphi \gamma \pi \pi `$ where the $`\pi \pi `$ system has total angular momentum $`J=0`$ and isospin $`I=0`$ is straightforward :
$`M(\varphi \gamma \pi \pi )`$ $`=`$ $`ϵ_\varphi ^\mu ϵ_\gamma ^\nu (p_\nu q_\mu g_{\nu \mu }(pq))H_{\pi \pi }(p^2,(pq)^2).`$ (7)
Here the scalar invariant amplitude $`H_{\pi \pi }(p^2,(pq)^2)`$ is given by (compare with Eq.(3))
$`H_{\pi \pi }(p^2,(pq)^2)`$ $`=`$ $`{\displaystyle \frac{eg_\varphi }{2\pi ^2m_K^2}}I(a,b)t_{K^+K^{}\pi \pi }`$ (8)
and $`t_{K^+K^{}\pi \pi }`$ is the $`J=0`$ part of the $`T`$-matrix for the $`K^+K^{}\pi \pi `$ scattering. The $`\pi \pi `$ invariant mass distribution has the form
$`{\displaystyle \frac{d\mathrm{\Gamma }}{dM_{\pi ^+\pi ^{}}}}=2{\displaystyle \frac{d\mathrm{\Gamma }}{dM_{\pi ^0\pi ^0}}}`$ $`=`$ $`{\displaystyle \frac{\alpha g_\varphi ^2\omega }{18(2\pi )^6m_\varphi ^2}}|(ab)I(a,b)|^2|t_{K\overline{K}\pi \pi }^0|^2k_{\pi \pi }`$ (9)
where $`t_{K\overline{K}\pi \pi }^0=\sqrt{6}t_{K^+K^{}\pi ^+\pi ^{}}`$ is the isoscalar $`K\overline{K}\pi \pi `$ amplitude and $`k_{\pi \pi }`$ is the relative momentum of the pions in the final state:
$`k_{\pi \pi }`$ $`=`$ $`\sqrt{M_{\pi \pi }^2/4m_\pi ^2},M_{\pi \pi }^2=(pq)^2.`$ (10)
Equation (9) leads to Eq.(6) in the Breit–Wigner (BW) approximation for the $`K^+K^{}\pi \pi `$ scattering amplitude
$`t_{K^+K^{}\pi \pi }`$ $`=`$ $`{\displaystyle \frac{g_{f_0KK}g_{f_0\pi \pi }}{M_{\pi \pi }^2(M_{f_0}i\mathrm{\Gamma }_{f_0}/2)^2}}`$ (11)
under the assumption that the integral over the $`\pi \pi `$ mass spectrum is saturated by the narrow resonance with mass $`M_{f_0}`$ and width $`\mathrm{\Gamma }_{f_0}`$ (see Appendix A).
According to Eq.(8), the hadronic part of the amplitude, which contains the information about the scalar mesons, is factored out in the form of the $`T`$-matrix for the $`K^+K^{}\pi \pi `$ scattering, with both the physical region ($`M_{\pi \pi }2m_K`$) and the unphysical region ($`M_{\pi \pi }<2m_K`$) being relevant to the $`\varphi \gamma \pi \pi `$ decay. It is known from the studies of scalar mesons, see and references therein, that the analytical structure of the scalar–isoscalar amplitudes near the $`K\overline{K}`$ threshold is far from being a trivial BW resonance. Therefore a coupled channel model of the $`K^+K^{}\pi \pi `$ scattering is required to describe the decay $`\varphi \gamma \pi \pi `$ beyond the BW approximation.
## 3 The $`\pi \pi K\overline{K}`$ Coupled Channel Model
To describe the interaction in the $`\pi \pi K\overline{K}`$ system with total angular momentum $`J=0^{++}`$ and isospin $`I^G=0^+`$ we exploit a coupled channel model similar to that of . The two scattering channels, 1 and 2, correspond to the $`\pi \pi `$ and $`K\overline{K}`$ systems and channel 3 contains a single $`q\overline{q}`$ bound state. The $`T`$-matrix, as a function of the invariant mass squared $`s`$, is defined by the Lippmann-Schwinger equation
$`𝑻(s)`$ $`=`$ $`𝑽+𝑽𝑮^0(s)𝑻`$ (12)
where $`𝑮^0(s)`$ is the free Green function.
The interaction potentials are taken in separable form:
$`𝑽`$ $`=`$ $`\left(\begin{array}{ccc}v_{11}(s)|11|& v_{12}(s)|12|& g_{13}|1q\overline{q}|\\ v_{12}(s)|21|& v_{22}(s)|22|& g_{23}|2q\overline{q}|\\ g_{13}|q\overline{q}1|& g_{23}|q\overline{q}2|& 0\end{array}\right)`$ (13)
where the form factors in channel 1 and 2 depend on the corresponding relative three–momentum $`k`$:
$`k|1`$ $`=`$ $`\xi _1(k)={\displaystyle \frac{\lambda _1^2}{k^2+\lambda _1^2}}`$ (14)
$`k|2`$ $`=`$ $`\xi _2(k)={\displaystyle \frac{\lambda _2^2}{k^2+\lambda _2^2}}`$ (15)
In , the potentials $`v_{ij}(s)`$ were assumed to be energy independent, and the chiral symmetry constraints on the scattering amplitude were strictly imposed only in the $`\pi \pi `$ channel by adjusting the strength of $`v_{11}`$ so that the Adler zero is at the correct position. In the present case, it is essential to ensure a correct behaviour of the $`K\overline{K}\pi \pi `$ scattering amplitude not only in the physical scattering region but also down to the $`\pi \pi `$ threshold. We found it easier to impose the chiral symmetry constraints by using energy dependent potentials<sup>1</sup><sup>1</sup>1The unitarization of the lowest order in chiral perturbation theory goes along a similar way, see and references therein.. The energy dependence is taken in the form
$`v_{11}(s)`$ $`=`$ $`bg_{11}s`$ (16)
$`v_{12}(s)`$ $`=`$ $`g_{12}s`$ (17)
$`v_{22}(s)`$ $`=`$ $`g_{22}s.`$ (18)
With our choice of interaction (13-18) the analytical solution for the $`T`$-matrix can be easily obtained. Further details of the model are given in Appendix B. As in we assume that the diagonal interaction in the $`K\overline{K}`$ channel produces a weakly bound state in the absence of coupling to the other channels, thus simulating a “molecular” origin of the $`f_0(980)`$ resonance. The state $`|q\overline{q}`$ in channel 3 has a bare mass $`M_r>2m_K`$.
The model parameters have been determined from the fit of the $`\pi \pi `$ scattering phase $`\delta _0^0`$ and the inelasticity parameter $`\eta _0^0`$ in the mass range $`M_{\pi \pi }1.5`$GeV, see Table 1 and Figs.2(a,b). Since the fit was found to be only weakly sensitive to the form–factor parameter $`\lambda _2`$, its value was fixed, and only the coupling constants $`g_{ij}`$, the bare mass of $`q\overline{q}`$, and $`\lambda _1`$ were treated as free parameters. The parameter $`b`$ in the diagonal $`\pi \pi `$ potential Eq.(16) was used for fine tuning of the $`\pi \pi `$ scattering length which was fixed at $`a_0^0=0.22m_\pi ^1`$. The best fit of only the $`\pi \pi `$ scattering data (fit 1) does not lead automatically to a very good description of the $`\pi \pi `$ invariant mass distribution in the decay $`\varphi \gamma \pi \pi `$. However, the shape of the $`\varphi \gamma \pi \pi `$ improves if the $`d\mathrm{\Gamma }/dM_{\pi \pi }`$ data are added to the fit as shown in Fig.3 (the details are discussed in Sec.4). As a result, our model provides a good description of the whole data set.
The energy dependence of the $`T`$-matrix element $`t_{K^+K^{}\pi \pi }`$ calculated in our coupled channel model is shown in Fig.4. In addition to a narrow peak due to the $`f_0(980)`$ state, this matrix element has a significant contribution from the lower mass region corresponding to the $`\sigma `$ meson.
## 4 The Decay $`\varphi \gamma \pi \pi `$ in the Coupled Channel Model
The formulas given in Sec. 2 are valid for point-like particles. When a form factor in the $`K\overline{K}\pi \pi `$ vertex is included, the $`\varphi \gamma \pi \pi `$ amplitude contains an additional term arising from the minimal substitution $`k_\mu (k_\mu eA_\mu )`$ in the momentum dependence of the form factors (see for details). With an appropriate choice of the form factors, the sum of three diagrams shown in Fig. 5 becomes explicitly finite. In order to study the influence of the form factors on the results for the $`\varphi \gamma \pi \pi `$ we use a nonrelativistic approximation for the $`K^+`$ and $`K^{}`$, which is justified by the fact that the most interesting region corresponding to the $`f_0(980)`$ resonance is very close to the $`K\overline{K}`$ threshold. The electric dipole matrix element is factorized into two parts describing the $`K^+K^{}`$ loop radiation with gauge invariant complement and the final state rescattering $`K^+K^{}\pi \pi `$, correspondingly. The total scalar invariant amplitude has the form
$`H_{\pi \pi }^{(\lambda )}(p^2,q^2)`$ $`=`$ $`{\displaystyle \frac{eg_\varphi }{2\pi ^2m_K^2}}J_\lambda (M_{\pi \pi })t_{K^+K^{}\pi \pi }(M_{\pi \pi })`$ (19)
which is similar to the relativistic point-like case defined by Eq.(8) where the function $`I(a,b)`$ is replaced by $`J_\lambda (M_{\pi \pi })`$. For the definition of $`J_\lambda (M_{\pi \pi })`$ and further technical details we refer to Appendix C. The parameter $`\lambda `$ refers to the form–factor dependence on the relative $`K^+K^{}`$ momentum in the $`K^+K^{}\pi \pi `$ vertex given by Eq.(15). Figure 6 shows the dependence of the electric dipole matrix element on the $`\pi \pi `$ invariant mass for different values of the form–factor parameter $`\lambda `$ in comparison with the relativistic point-like case.
With our choice of the form factor (15) there is no substantial suppression of the nonrelativistic result<sup>2</sup><sup>2</sup>2The stronger dependence of the total amplitude on the form factor found in results from the use of the dipole form factor which falls off faster than the monopole form factor used in our case. in comparison to the point-like case for the relevant range of $`\lambda =0.60.8`$GeV (which corresponds to the data fit in our coupled channel model). For larger values of $`\lambda `$, the full relativistic treatment is needed as the real part of the nonrelativistic result becomes sensitive to the short distance contribution (see the discussion in ). The dependence of the imaginary part of $`M(\varphi \gamma \pi \pi )`$ on the form factor in the mass region close to the $`K\overline{K}`$ threshold is rather weak. As a result, we can neglect the form–factor dependence in the loop calculations and just use the relativistic point–like result given by Eq.(8). It is interesting to note that the imaginary part can be obtained from the divergent triangle diagram Fig.5(a) alone by using the Siegert theorem which ensures a correct behaviour of the electric dipole matrix element.
## 5 The Poles of the $`S`$-Matrix
The poles of the $`S`$-matrix in the complex $`s`$-plane corresponding to the fits in Sect.3 are shown in Table 2 and Fig.7. There are five poles related to the resonances in our model. The pole $`M_A`$ located on the sheet II ($`\mathrm{I}mk_1<0`$, $`\mathrm{I}mk_2>0`$) is very close to the $`K\overline{K}`$ threshold. The pair of poles, $`M_B`$ on sheet II and $`M_D`$ on the sheet III ($`\mathrm{I}mk_1<0`$, $`\mathrm{I}mk_2<0`$), corresponds to a broad structure associated with the $`\sigma `$ meson. The other pair of poles, $`M_C`$ on sheet II and $`M_E`$ on sheet III, corresponds to a broad resonance above the $`K\overline{K}`$ threshold, which can be associated with the $`f_0(1370)`$ state (for a more realistic description of the $`S`$-matrix above 1.3 GeV additional poles are needed which are not included in the present model). The position of the resonance poles depends on the coupling constants, and this distinguishes them from the fixed poles originating from the singularities of the form factors (15). The latter are located at $`k_1=\pm i\lambda _1`$ and $`k_2=\pm i\lambda _2`$, their distance to the physical region being determined by the range of the interaction. In our model these fixed poles approximate the potential singularities which correspond to the left hand cut in a more general case.
The origin and the nature of the resonance poles found in our model can be elucidated by studying how these poles move in the complex $`s`$-plane when the model parameters are varied between the physical case determined by the fit and the limit of vanishing couplings in the $`\pi \pi `$ channel and between the $`\pi \pi `$ and the other channels ($`K\overline{K}`$ and $`q\overline{q}`$):
$`v_{11}(s)`$ $``$ $`xv_{11}(s),0x1`$ (20)
$`v_{12}(s)`$ $``$ $`xv_{12}(s)`$ (21)
$`g_{13}`$ $``$ $`x^{1/2}g_{13}`$ (22)
$`g_{23}`$ $``$ $`x^{1/2}g_{13}`$ (23)
The diagonal interaction in the $`K\overline{K}`$ channel with the physical strength of the coupling $`g_{22}`$ produces a bound state close to the $`K\overline{K}`$ threshold with mass $`m_{K\overline{K}}=0.97`$GeV.
Our coupled channel model has only one pole $`M_A=0.975i0.017`$GeV near the $`K\overline{K}`$ threshold, which is sufficient for a good description of the $`\pi \pi `$ scattering data. This pole is directly related to a molecular $`K\overline{K}`$ state in the absence of coupling to the $`\pi \pi `$ channel. The number of the $`S`$-matrix poles near the $`K\overline{K}`$ threshold has been discussed in the literature (see and references therein) for a long time. While the relation of the $`f_0(980)`$ meson to at least one $`S`$-matrix pole close to $`K\overline{K}`$ threshold is well established, the exact location and even the number of the relevant $`S`$-matrix poles is model dependent as demonstrated in Fig.8. The models based on dynamical input (coupled channel, potential, unitarized chiral perturbation theory) usually produce only one pole which can be traced to a weakly bound $`K\overline{K}`$ state in an appropriate limit of the channels coupling. One exception to this observation is a coupled channel model where the second pole near the $`K\overline{K}`$ threshold results from the interplay of a nearby $`q\overline{q}`$ pole with dynamical singularities. We were not able to find a good fit with the same feature in our model; the main reason appears to be due to the using of the energy dependent potentials in Eq.(13) contrary to the case in . However, it is not excluded that two–pole solutions can be found with some other parametrization of interactions in CCM. It remains to be investigated whether two–pole solutions are compatible with the data on the $`\varphi \gamma \pi \pi `$ decay.
The $`K`$-matrix parameterizations routinely produce the second pole on sheet III. There is, however, a much larger spread in this pole location than on sheet II. As shown above, the reaction $`\varphi \gamma \pi \pi `$ allows one to probe the scattering $`K\overline{K}\pi \pi `$ both above and below the $`K\overline{K}`$ threshold, and the $`\pi \pi `$ mass distribution in $`\varphi \gamma \pi \pi `$ is very sensitive to the $`S`$-matrix poles related to the $`f_0(980)`$. Therefore more detailed experimental data on $`\varphi \gamma \pi \pi `$ would be very useful for reducing the present uncertainty about the analytical structure of the $`S`$-matrix in the $`f_0(980)`$ region.
## 6 The Mixing between the $`q\overline{q}`$ and Mesonic Channels
The mixing of the quark-antiquark states with the open meson channels can be studied in the CCM by using the probability sum rule for a resonance embedded into a continuum which is described in Appendix B. The spectral function $`\rho (s)`$ which determines the probability density for the $`q\overline{q}`$ component in the scattering states, as defined by Eq.(56), is shown in Fig.9(a). In the case of weak coupling, the probability density would be well localized near the position of the bare $`q\overline{q}`$ state at $`M_r=1.1`$GeV. For the physical case, we find a broad peak centered above the $`K\overline{K}`$ threshold in the region corresponding to the $`f_0(1370)`$ resonance (the poles $`M_E`$ and $`M_C`$). There is also a sizable contribution to the $`q\overline{q}`$ spectral density from the low–mass region of the $`\sigma `$ meson. The $`\rho (s)`$ distribution in the $`f_0(980)`$ resonance has a characteristic dip-bump structure resulting from an interplay of the resonance pole $`M_A`$ with a nearby zero of the mass operator $`\mathrm{\Pi }(s)`$. The position and the width of the $`\rho (s)`$ distribution indicates that an essential contribution to the saturation of the sum rule (56) comes from the pole $`M_B`$ related to the $`\sigma `$ meson, while the narrow structure associated with the pole $`M_A`$ alone plays a minor role. The fact that $`q\overline{q}`$ coupling with the $`\pi \pi `$ channel significantly enhances the spectral density $`\rho (s)`$ in the region of the $`\sigma `$ meson is related to the interplay of the $`S`$-matrix poles demonstrated in Fig.7c: the pole $`M_B`$ corresponding to the $`\sigma `$ meson is pushed towards the $`\pi \pi `$ threshold by the pole $`M_C`$ originating from the $`q\overline{q}`$ state.
Using the spectral density $`\rho (s)`$ we can calculate the contribution of the $`q\overline{q}`$ scalar mesons to the QCD sum rule related to the scalar quark condensate . Figure 9b shows the Laplace transform of the spectral density which is used in the sum rule analysis:
$`I_n(M^2)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}s^ne^{s/M^2}\rho (s)𝑑s`$ (24)
The advantage of CCM with respect to the previous studies (see and references therein) where the contribution of the $`q\overline{q}`$ mesons was usually approximated by one narrow resonance is a more realistic shape of the $`q\overline{q}`$ spectral density in the low-mass region which is important in the Laplace sum rules. As a result, the momenta $`I_0(M^2)`$ and $`I_1(M^2)`$ have quite different slopes in their $`M^2`$–dependence as shown in Fig.9(b).
## 7 Discussion
The theoretical calculations of the decays $`\varphi \gamma \pi \pi `$ and $`\varphi \gamma f_0(980)`$ and the corresponding experimental data are summarized in Table 3. The predicted branching ratios for the $`\gamma \pi ^0\pi ^0`$ and $`\gamma \pi ^+\pi ^+`$ channels should be compared with the experimental data with some caution. While these branching ratios can be easily defined theoretically ($`\mathrm{\Gamma }_{\varphi \gamma \pi ^+\pi ^{}}/\mathrm{\Gamma }_{\varphi \gamma \pi ^0\pi ^0}=2`$ for the isoscalar $`\pi \pi `$ states, e.g.), the data analysis relies on the modeling of competing reaction mechanisms. In particular, the interpretation of the experimental results for the $`\gamma \pi ^+\pi ^+`$ channel involves the consideration of the interference between the $`e^+e^{}\varphi \gamma \pi \pi `$ and $`e^+e^{}\gamma \rho \gamma \pi \pi `$ mechanisms. In this respect, we wish to emphasize again the importance of using a realistic $`K^+K^{}\pi \pi `$ amplitude instead of a simple superposition of BW resonances. The channel $`\gamma \pi ^0\pi ^0`$ which is free from the $`\rho `$ meson contribution in the $`\pi \pi `$ channel appears to be better suited for the study of the $`\pi \pi `$ invariant mass distribution, although this case has some background due to the mechanism $`\varphi \pi ^0\rho ^0\pi ^0\pi ^0\gamma `$.
Our results are in good agreement with recent calculations within the chiral unitary approach both for the total width and for the $`\pi \pi `$ mass distribution. This is not surprising because in both cases the $`f_0(980)`$ resonance is produced mainly by the attractive interaction in the $`K\overline{K}`$ channel. Our result for the total width $`\mathrm{\Gamma }_{\varphi \gamma \pi \pi }`$ is close to the earlier calculations of the two–step mechanism with the intermediate $`K^+K^{}`$ state which used the BW approximation (6) with the coupling constant $`g_{f_0K\overline{K}}^2/4\pi =0.6\text{GeV}^2`$. In our model we can define an effective coupling constant $`g_{f_0K\overline{K}}`$ by approximating the $`K\overline{K}\pi \pi `$ amplitude (52) by a BW resonance which leads to slightly higher value of $`g_{f_0K\overline{K}}^2/4\pi =1.1\text{GeV}^2`$.
The effect of the form factor in the $`K\overline{K}\pi \pi `$ vertex was studied earlier in where a suppression by a factor of about 5 was given as an estimate. This significant suppression resulted from the use of a very soft dipole form factor with the characteristic parameter $`\mu =0.14`$GeV (Eq.(4.14) in ) which was suggested by the study of the $`\varphi \gamma \gamma `$ decay . However, the $`2\gamma `$ decay is related to the short distance behavior of the $`K\overline{K}`$ wave function, while we are interested in the $`\varphi K\overline{K}`$ vertex at moderate relative momenta (in general, there is no unique relation between these two properties, see the discussion in ). It is difficult to justify such a low value of the form–factor parameter in the CCM: a good fit of the scattering data requires much harder form factors as discussed in Sec. 3. A weak effect of the form factor found in our case is consistent with the results in if the form–factor parameter $`\mu 0.5`$GeV is used there.
In view of this overall agreement between different calculations of the two–step mechanisms with the $`K^+K^{}`$ intermediate state we find it difficult to agree with the statements that the $`K\overline{K}`$ model of $`f_0(980)`$ can be excluded on the ground of its alleged conflict with the $`\varphi \gamma f_0(980)`$ data. Since the hadronic part of the amplitude $`\varphi \gamma \pi \pi `$ is factored out as the matrix element $`t_{K\overline{K}\pi \pi }`$ or, in the simplified case, as the coupling constant $`g_{f_0K\overline{K}}`$, the self-consistency of the $`K\overline{K}`$ model of $`f_0(980)`$ can be examined by analyzing the value of $`g_{f_0K\overline{K}}`$. For this purpose we use a well known result from nonrelativistic scattering theory: the residue of the scattering matrix at a pole corresponding to a bound state is uniquely related to the asymptotic normalization constant of the bound state wave function which for a weakly bound state depends, in leading order, only on the binding energy. For a weakly bound $`K\overline{K}`$ state, the relation between the $`g_{f_0K\overline{K}}`$ and the binding energy $`E_b=2m_Km_{f_0}`$ has the form
$`{\displaystyle \frac{g_{f_0K\overline{K}}^2}{4\pi }}`$ $`=`$ $`32m_K\kappa (1+O(\kappa /\lambda ))`$ (25)
$`\kappa `$ $`=`$ $`\sqrt{m_KE_b}`$ (26)
where $`\lambda `$ is the range of the $`K\overline{K}`$ interaction. Taking $`E_b=4`$MeV and neglecting small corrections of the order $`\kappa /\lambda `$ one gets $`\frac{g_{f_0K\overline{K}}^2}{4\pi }=0.7\text{GeV}^2`$ which is very close to the values discussed above. Therefore a large $`K\overline{K}`$ component is naturally expected for the $`f_0(980)`$ state regardless of further details of particular models.
## 8 Conclusion
The decay $`\varphi \gamma \pi \pi `$ has been studied in an exactly solvable coupled channel model containing the $`\pi \pi `$, $`K\overline{K}`$, and $`q\overline{q}`$ channels using separable potentials. The $`f_0(980)`$ resonance corresponds to one $`S`$-matrix pole close to the $`K\overline{K}`$ threshold; this pole has a dynamical origin and represents the molecular-like $`K\overline{K}`$ state. The molecular picture of the $`f_0(980)`$ meson is found to be in a fair agreement with the experimental data. We confirm the assessment of that the earlier conclusions about suppression of the $`\varphi \gamma f_0(980)`$ branching ratio in the molecular $`K\overline{K}`$ model were partly related to differences in modeling.
The lightest scalar meson, $`\sigma `$, has a dynamical origin resulting from the attractive character of the effective $`\pi \pi `$ interaction, with a partial contribution from the coupling via the intermediate scalar $`q\overline{q}`$ states. The distinction between genuine $`q\overline{q}`$ states and dynamical resonances, $`\sigma `$ and $`f_0(980)`$, can be illuminated by considering the limit $`N_c\mathrm{}`$ where the $`q\overline{q}`$ states turn into infinitely narrow resonances while the dynamical states disappear altogether.
The structure of the $`q\overline{q}`$ state embedded into the mesonic continuum has been analyzed using the calculated $`q\overline{q}`$ spectral density. The gross structure of the quark–antiquark spectral density $`\rho (s)`$ is related to the $`f_0(1370)`$ resonance. There is also a significant contribution to $`\rho (s)`$ in the low mass region ($`\sigma `$ meson) which is related to the strong coupling between the $`\pi \pi `$ and $`q\overline{q}`$ channels. The same approach can also be used for the QCD sum rules related to the gluon condensate by extending the coupled channel model to include the mixing with the scalar glueballs. The consideration of this topic is beyond the scope of this paper.
## Acknowledgments
The author thanks M.P. Locher for a fruitful collaboration which lead to this paper, S.I. Eidelman for bringing attention to the problem of the radiative $`\varphi `$ decays and a discussion of the experimental data, D. Bugg and B.S. Zou for a discussion of the nature of the $`f_0(980)`$.
## Appendix Appendix A The function $`I(a,b)`$ and the decay widths
The function $`I(a,b)`$ is given by (see e.g. and references therein)
$`I(a,b)`$ $`=`$ $`{\displaystyle \frac{1}{2(ab)}}{\displaystyle \frac{2}{(ab)^2}}\left(f({\displaystyle \frac{1}{b}})f({\displaystyle \frac{1}{a}})\right)+{\displaystyle \frac{a}{(ab)^2}}\left(g({\displaystyle \frac{1}{b}})g({\displaystyle \frac{1}{a}})\right)`$ (27)
$`f(x)`$ $`=`$ $`\{\begin{array}{cc}(\mathrm{arcsin}\frac{1}{2\sqrt{x}})^2,\hfill & x>\frac{1}{4}\hfill \\ \frac{1}{4}\left(\mathrm{ln}\frac{1+\sqrt{14x^2}}{1\sqrt{14x^2}}i\pi \right)^2,\hfill & x\frac{1}{4}\hfill \end{array}`$ (30)
$`g(x)`$ $`=`$ $`\{\begin{array}{cc}\sqrt{4x^21}\mathrm{arcsin}\frac{1}{2\sqrt{x}},\hfill & x>\frac{1}{4}\hfill \\ \frac{1}{2}\sqrt{14x^2}\left(\mathrm{ln}\frac{1+\sqrt{14x^2}}{1\sqrt{14x^2}}i\pi \right),\hfill & x\frac{1}{4}\hfill \end{array}`$ (33)
The $`\varphi K^+K^{}`$ coupling constant $`g_\varphi `$ is related to the decay width by
$`\mathrm{\Gamma }(\varphi K^+K^{})`$ $`=`$ $`{\displaystyle \frac{g_\varphi ^2}{48\pi m_\varphi ^2}}(m_\varphi ^24m_K^2)^{3/2}.`$ (34)
The relation between the coupling constants and decay widths for the scalar meson $`f_0`$ has the form
$`\mathrm{\Gamma }(f_0\pi \pi )`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{g_{f_0\pi \pi }^2}{16\pi m_{f_0}^2}}(m_{f_0}^24m_\pi ^2)^{1/2}`$ (35)
where the extra factor $`\frac{1}{2}`$ accounts for the identity of the two pions in the final state.
## Appendix Appendix B The Coupled Channel Model
The free Green function $`G^0(s)`$ is a diagonal matrix:
$`𝑮^0(s)`$ $`=`$ $`\left(\begin{array}{ccc}𝑮_1^0(s)& 0& 0\\ 0& 𝑮_2^0(s)& 0\\ 0& 0& 𝑮_3^0(s)\end{array}\right)`$ (36)
where the single–channel Green functions have the form
$`𝑮_1^0(s)`$ $`=`$ $`{\displaystyle \frac{2}{\pi }}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{|k_1k_1|}{s/4(m_\pi ^2+k_1^2)}}k_1^2𝑑k_1`$ (37)
$`𝑮_2^0(s)`$ $`=`$ $`{\displaystyle \frac{2}{\pi }}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{|k_2k_2|}{s/4(m_K^2+k_2^2)}}k_2^2𝑑k_2`$ (38)
$`𝑮_3^0(s)`$ $`=`$ $`G_3^0(s)|q\overline{q}q\overline{q}|,G_3^0(s)={\displaystyle \frac{1}{sM_r^2}}.`$ (39)
Here $`|k_1`$ and $`|k_2`$ denote the free $`\pi \pi `$ and $`K\overline{K}`$ states with relative momenta $`k_1`$ and $`k_2`$, respectively. The state $`|q\overline{q}`$ in channel 3 has a bare mass $`M_r`$. With the form factors given by Eq.(15) the matrix elements of the Green functions are
$`G_n^0(s)=n|𝑮_n^0(s)|n`$ $`=`$ $`{\displaystyle \frac{\lambda _n^3}{2(k_n(s)+i\lambda _n)^2}},n=1,2`$ (40)
where $`k_n(s)`$ is the relative momentum in the channel $`n`$:
$`k_1(s)`$ $`=`$ $`\sqrt{s/4m_\pi ^2}`$ (41)
$`k_2(s)`$ $`=`$ $`\sqrt{s/4m_K^2}.`$ (42)
The $`\pi \pi `$ elastic scattering amplitude $`f_{\pi \pi }(s)`$ and the $`\pi \pi K\overline{K}`$ amplitude $`f_{\pi \pi K\overline{K}}(s)`$ have the form:
$`f_{\pi \pi }(s)`$ $`=`$ $`k_1|T(s)|k_1=\xi (k_1)^2{\displaystyle \frac{N_{11}(s)}{D(s)}}`$ (43)
$`f_{\pi \pi K\overline{K}}(s)`$ $`=`$ $`k_1|T(s)|k_1=\xi (k_1)\xi (k_2){\displaystyle \frac{N_{12}(s)}{D(s)}}`$ (44)
where
$`D(s)`$ $`=`$ $`1u_{11}(s)G_1^0(s)u_{22}(s)G_2^0(s)`$ (45)
$`+(u_{11}(s)u_{22}(s)u_{12}^2(s))G_1^0(s)G_2^0(s)`$
$`N_{11}(s)`$ $`=`$ $`u_{11}(s)(u_{11}(s)u_{22}(s)u_{12}^2(s))G_3^0(s)`$ (46)
$`N_{12}(s)`$ $`=`$ $`u_{12}`$ (47)
$`u_{11}(s)`$ $`=`$ $`v_{11}(s)+{\displaystyle \frac{g_{13}^2}{sM_r^2}}`$ (48)
$`u_{12}(s)`$ $`=`$ $`v_{12}(s)+{\displaystyle \frac{g_{13}g_{23}}{sM_r^2}}`$ (49)
$`u_{22}(s)`$ $`=`$ $`v_{22}(s)+{\displaystyle \frac{g_{23}^2}{sM_r^2}}.`$ (50)
The connection between the partial wave $`S`$-matrix and the scattering amplitude $`f_{\pi \pi }`$ is given by
$`S_{J=0}^{I=0}(s)=\eta _0^0(s)e^{2i\delta _0^0(s)}=1+2ik_1f_{\pi \pi }(s)`$ (51)
where $`\delta _0^0(s)`$ is the scattering phase and $`\eta _0^0(s)`$ is the inelasticity parameter. The $`T`$-matrix element is related to the amplitude (44) by
$`t_{K\overline{K}\pi \pi }`$ $`=`$ $`8\pi \sqrt{s}f_{K\overline{K}\pi \pi }.`$ (52)
The spectral density of the $`q\overline{q}`$ state has the form
$`\rho (s)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi i}}(G_3(siϵ)G_3(s+iϵ)),`$ (53)
where $`G_3(s)`$ is the exact Green function in the $`|q\overline{q}`$ subspace:
$`G_3(s)`$ $`=`$ $`q\overline{q}|G(s)|q\overline{q}={\displaystyle \frac{1}{sM_r^2\mathrm{\Pi }(s)}}`$ (54)
and $`\mathrm{\Pi }(s)`$ is the mass operator of the $`q\overline{q}`$ state:
$$\mathrm{\Pi }(s)=\frac{4(g_{13}^2G_1^0(s)+g_{23}^2G_2^0(s)+(2v_{12}(s)g_{13}g_{23}g_{13}^2v_{22}(s)g_{23}^2v_{11}(s))G_1^0(s)G_2^0(s))}{1v_{11}(s)G_1^0(s)v_{22}(s)G_2^0(s)+(v_{11}(s)v_{22}(s)v_{12}^2(s))G_1^0(s)G_2^0(s)}.$$
(55)
The spectral density $`\rho (s)`$ satisfies the normalization condition
$`{\displaystyle _{4m_\pi ^2}^{\mathrm{}}}\rho (s)𝑑s`$ $`=`$ $`1.`$ (56)
Equations (53,56) represent the completeness relation projected onto the $`q\overline{q}`$ channel and the normalization $`q\overline{q}|q\overline{q}=1`$.
## Appendix Appendix C The Nonrelativistic Approximation
The amplitudes corresponding to the diagrams in Fig. 5 have the form:
$`M^{NR}`$ $`=`$ $`{\displaystyle \frac{eg_\varphi t_{K^+K^{}\pi \pi }}{(2\pi )^3}}(I_a+I_b+I_c)={\displaystyle \frac{eg_\varphi t_{K^+K^{}\pi \pi }m_\varphi \omega }{2\pi ^2m_K^2}}J_\lambda (M_{\pi \pi })`$ (57)
where
$`I_a`$ $`=`$ $`8{\displaystyle \frac{d^3𝐤}{2E_k}\frac{(𝜺_\varphi 𝐤)(𝜺_\gamma (𝐤+𝐪/2))F(|𝐤+𝐪/2|)}{(k^2m_K^2+iϵ)((kq)^2m_K^2+iϵ)}}=`$ (58)
$`=`$ $`{\displaystyle \frac{4m_K}{m_\varphi M_{\pi \pi }}}{\displaystyle d^3𝐤\frac{(𝜺_\varphi 𝐤)(𝜺_\gamma 𝐤)F(|𝐤+𝐪/2|)}{(\mathrm{\Delta }_\varphi 𝐤^2+iϵ)(\mathrm{\Delta }_{\pi \pi }(𝐤+𝐪/2)^2+iϵ)}}`$ (59)
$`I_b`$ $`=`$ $`2(𝜺_\varphi 𝜺_\gamma ){\displaystyle \frac{d^3𝐤}{2E_k}\frac{F(|𝐤+𝐪/2|)}{((kq)^2m_K^2+iϵ)}}=`$ (60)
$`=`$ $`{\displaystyle \frac{(𝜺_\varphi 𝜺_\gamma )}{M_{\pi \pi }}}{\displaystyle d^3𝐤\frac{F(|𝐤|)}{(\mathrm{\Delta }_{\pi \pi }𝐤^2+iϵ)}}`$ (61)
$`I_c`$ $`=`$ $`2{\displaystyle \frac{d^3𝐤}{2E_k}\frac{(𝜺_\varphi 𝐤)(𝜺_\gamma 𝐤/k)\frac{d}{d|𝐤|}F(|𝐤|)}{(k^2m_K^2+iϵ)}}=`$ (62)
$`=`$ $`I_{c1}+I_{c2}`$ (63)
$`I_{c1}`$ $`=`$ $`{\displaystyle \frac{(𝜺_\varphi 𝜺_\gamma )}{M_\varphi }}{\displaystyle d^3𝐤\frac{F(|𝐤|)}{(\mathrm{\Delta }_\varphi 𝐤^2+iϵ)}}`$ (64)
$`I_{c2}`$ $`=`$ $`{\displaystyle \frac{2(𝜺_\varphi 𝜺_\gamma )}{3M_\varphi }}{\displaystyle d^3𝐤\frac{𝐤^2F(|𝐤|)}{(\mathrm{\Delta }_\varphi 𝐤^2+iϵ)^2}}`$ (65)
Here $`E_k=\sqrt{k^2+m_K^2}`$, $`\mathrm{\Delta }_\varphi =(m_\varphi 2m_K)m_K`$, $`\mathrm{\Delta }_{\pi \pi }=(M_{\pi \pi }2m_K)m_K`$, the form factor $`F(|𝐤|)=\xi _2(|𝐤|)`$ according to Eq.(15), the subscript $`\lambda `$ in $`J_\lambda (M_{\pi \pi })`$ refers to the form–factor parameter. In deriving the nonrelativistic approximation (59-63) we keep only the positive-energy parts of the kaon propagators and substitute $`d^3k/E_kd^3k/m_K`$ in the final integral. The total amplitude $`M^{NR}`$ vanishes at $`\omega =0`$ as expected for the electric dipole transition. The diagram Fig.5(c) contains two terms according to Eq.(63) which can be combined with the loop radiation term (59) and the contact term (61) in a way that the combinations $`(I_a+I_{c2})`$ and $`(I_b+I_{c1})`$ are explicitly finite and both vanish at $`\omega =|𝐪|0`$. The integrals (59-65) with our choice of the form factor can be straightforwardly calculated in analytical form.
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# 1 INTRODUCTION
## 1 INTRODUCTION
The changes in the emission and absorption spectra of a gas placed in a strong electromagnetic field are the result of three effects. One consists of the formation of a nonequilibrium velocity distribution (Bennett’e ”holes” and “peaks”<sup></sup>). This factor significantly influences the spectral characteristics of lasers and was studied in detail by many authors. The second effect stems from the splitting of atomic levels; it was directly observed in the optical portion of the spectrum only very recently<sup>\[2,3,\]</sup> in the case of potassium atoms placed in the tremendous fields of a ruby laser. In gas lasers the fields are weaker, level splitting is much smaller than the Doppler line width, and the observability of the effect is not a simple matter. For example, according to Feld and Javan<sup></sup>, splitting is not possible at all in this case. This conclusion however is the consequence of an error in their calculations (see discussion of (3.4) below). Finally, the third effect of a strong external field consists in the fact that the probability of absorption or emission of photons turns out to depend not only on level populations but also on the polarization induced by the external field, i.e., on the nonlinear interference effect (NIE)<sup>\[5-7\]</sup>. This effect is the subject of the present paper.
The interest in NIE is due to several causes. First, it is this effect that is responsible for causing the spectral densities of Einstein coefficients, of absorption or emission to be different frequency functions leading to characteristic changes in the pure emission or absorption lines \[7-9\]. The NIE contribution should depend significantly oh the relaxation characteristics<sup></sup>, providing new opportunities to study collisions. For gas systems with large Doppler broadening the theory predicts an angular anisotropy of spectral characteristics and a possibility of obtaining an extremely sharp structure<sup>\[4-6,10\]</sup>. Although the early experiments with spontaneous<sup></sup> and stimulated emission<sup></sup> have so far failed to provide a quantitative verification of the theory, they have undoubtedly established the existence of the anisotropy effect.
The present work investigates NIE in gaseous systems and considers the problem under what conditions the plays a major role. It is shown that under certain conditions the velocity (distribution of atoms In a strong field does not change at all while the interference effects remain.
## 2 GENERAL EXPRESSIONS
We consider the photon emission of two monochromatic fields interacting with an atom whose term system is shown in Fig. 1. One of the two fields is regarded as strong and it resonates with the $`mn`$ transition, the matrix element of interaction (traveling wave) is
$`V_{mn}\mathrm{exp}\{i\omega _{mn}t\}=G\mathrm{exp}\{i(\mathrm{\Omega }tkr)\}.`$ (2.1)
$`G=d_{mn}E/2\mathrm{},\mathrm{\Omega }=\omega _\mu \omega _{mn}.`$ (2.2)
We are interested in emission or absorption of photons of a field resonating with one of the four transitions. $`nj`$, $`ml`$, $`fm`$, and $`gn`$ (Fig.1). For example in the case of $`nj`$
$`V_{nj}\mathrm{exp}\{i\omega _{nj}t\}=G_\mu \mathrm{exp}\{i(\mathrm{\Omega }_\mu tk_\mu r)\},`$ (2.3)
$`G_\mu =d_{nj}E_\mu /2\mathrm{},\mathrm{\Omega }_\mu =\omega _\mu \omega _{nj}`$ (2.4)
The system of equations for the density matrix has the form
$`L_{jj}\rho _{jj}=V_{nj}\rho _{nn}+q_j,`$ (2.5)
$`L_{jn}\rho _{jn}=iV_{mn}\mathrm{exp}\{i\omega _{mn}t\}\rho _{jm}=`$ (2.6)
$`=iV_{nj}^{}\mathrm{exp}\{i\omega _{nj}t\}(\rho _{nn}\rho _{jj}),`$ (2.7)
$`L_{jm}\rho _{jm}=iV_{mn}^{}\mathrm{exp}\{i\omega _{mn}t\}\rho _{jn}=`$ (2.8)
$`=iV_{nj}^{}\mathrm{exp}\{i\omega _{nj}t\}\rho _{nm};`$ (2.9)
$`L_{mm}\rho _{mm}=+2Re[iV_{mn}\mathrm{exp}\{i\omega _{mn}t\}\rho _{nm}]=q_m,`$ (2.10)
$`L_{nn}\rho _{nn}=2Re[iV_{mn}\mathrm{exp}\{i\omega _{mn}t\}\rho _{nm}]=q_n+\gamma _{mn}\rho _{mm},`$ (2.11)
$`L_{nm}\rho _{nm}=iV_{nm}\mathrm{exp}\{i\omega _{mn}t\}(\rho _{nn}\rho _{mm})=q_m,`$ (2.12)
$`L_{ik}=/t+v+\mathrm{\Gamma }_{ik},\mathrm{\Gamma }_{ll}\mathrm{\Gamma }_l,`$ (2.13)
$`\mathrm{\Gamma }_{ik}`$ are transition widths and $`q_i`$ is the rate of excitation of atoms to the state $`i`$, $`\mathrm{v}`$.
According to (2.9) and (2.13) the field $`V_{jn}`$ does not affect the population (”weak field”). Therefore the entire system of equations was found to be split up; eqs. (2.13) include only $`\rho _{mm}`$, $`\rho _{nn}`$, and $`\rho _{nm}`$, and the solution of the system serves as a ”source” for the computation of $`\rho _{jm}`$, $`\rho _{jn}`$ and $`\rho _{jj}`$ from (2.9). In the case of (2.1) and (2.3) the system (2.9) (2.13) reduces to equations whose solution has the form
$`\rho _{jj}=n_j+{\displaystyle \frac{\gamma _{nj}}{\mathrm{\Gamma }_j}}\rho _{nn},`$ (2.14)
$`\rho _{nn}=n_n+{\displaystyle \frac{2\pi G^2}{\mathrm{\Gamma }_n\sqrt{1+\ae }}}\left(1{\displaystyle \frac{\gamma _{mn}}{\mathrm{\Gamma }_n}}\right)(n_mn_n)W_B(v),`$ (2.15)
$`\rho _{mm}=n_m{\displaystyle \frac{2\pi G^2}{\mathrm{\Gamma }_m\sqrt{1+\ae }}}(n_mn_n)W_B(v),`$ (2.16)
$`\rho _{nm}=r_{nm}\mathrm{exp}\{i(\mathrm{\Omega }tkr)\},`$ (2.17)
$`r_{nm}=iG(\rho _{mm}\rho _{nn})/(\mathrm{\Gamma }+i\mathrm{\Omega }^{})`$ (2.18)
where
$`W_B(v)=\mathrm{\Gamma }_B/\pi [\mathrm{\Gamma }_B^2+(\mathrm{\Omega }kv)^2],\mathrm{\Gamma }_B=\mathrm{\Gamma }\sqrt{1+\ae },`$ (2.19)
$`\mathrm{\Gamma }\mathrm{\Gamma }_{nm},\mathrm{\Omega }^{}=\mathrm{\Omega }kv,`$ (2.20)
$`\mathrm{\Omega }_{\mu }^{}{}_{}{}^{}=\mathrm{\Omega }_\mu k_\mu v,\ae =\tau ^2G^2={\displaystyle \frac{2(\mathrm{\Gamma }_m+\mathrm{\Gamma }_n\gamma _{mn})}{\mathrm{\Gamma }_m\mathrm{\Gamma }_n\mathrm{\Gamma }}},`$ (2.21)
$`n_i={\displaystyle \frac{q_i(v)}{\mathrm{\Gamma }_i}}+{\displaystyle \frac{\gamma _{ki}}{\mathrm{\Gamma }_i}}{\displaystyle \frac{q_k(v)}{\mathrm{\Gamma }_k}}`$ (2.22)
The quantities $`n_i(v)`$ represent velocity distributions of atoms in the absence of a strong field $`(G=0)`$ determined by excitation processes $`q_i(v)`$.
The emission (absorption) power is determined by the general formula
$$w_{nj}=2\mathrm{}\omega _{nj}ReiV_{nj}\mathrm{exp}\{i\omega _{nj}t\}\rho _{jn},$$
(2.23)
where the angle brackets designate averaged velocities $`v`$ of atoms. Using the system (2.9) we can express $`\rho _{jn}`$ in terms of (2.18) and obtain an expression for power (2.23) in the form
$$w_{nj}=2\mathrm{}\omega _{nj}|G_\mu |^2Re\frac{[\mathrm{\Gamma }_{jm}+i(\mathrm{\Omega }_{\mu }^{}{}_{}{}^{}+\mathrm{\Omega }^{})](\rho _{nn}\rho _{jj})iGr_{nm}}{[\mathrm{\Gamma }_{jm}+i(\mathrm{\Omega }_{\mu }^{}{}_{}{}^{}+\mathrm{\Omega }^{})][\mathrm{\Gamma }_{jn}+i\mathrm{\Omega }_{\mu }^{}{}_{}{}^{}]+G^2}.$$
(2.24)
Equation (2.24) clearly reflects the classification of effects due to the external field. The denominator contains squares $`\mathrm{\Omega }_\mu `$ terms, i.e., it contains resonances at two frequencies. This can be interpreted as a splitting of the atom levels in the external field The numerator in (2.24) contains two terms with significantly different properties. The first term is proportional to the population difference $`\rho _{nn}\rho _{jj}`$ containing Bennett’s ”holes,” as reflected in the factor $`W_B(v)`$ (henceforth called the Bennett distribution). The second term proportional to $`r_{nm}`$ varies only the the shape but not its integral intensity, since
$$\underset{\mathrm{}}{\overset{+\mathrm{}}{}}w_{nj}𝑑\mathrm{\Omega }_\mu =2\pi \mathrm{}\omega _{nj}|G_\mu |^2<\rho _{nn}\rho _{jj}>.$$
The fact that this term appeared and its property are not at all specific to the special case under consideration. According to (2.9) the ”sources” that ”excite” $`\rho _{jm}`$ and $`\rho _{jn}`$ are both the population difference $`\rho _{nn}\rho _{jj}`$ and the non-diagonal element $`\rho _{nm}`$ stimulated by the strong field for any spectral composition of the strong field. Therefore $`w_{nj}`$ contains $`\rho _{nm}`$ also in the general case, and not only in a monochromatic field. We can say that this term reflects the ”coherence” that is contributed to the atomic state by the strong field, so that a weak field ”mixes” the m and j states as well as the n and j elates. The last circumstance causes oscillations at the frequency $`\omega +\omega _\mu `$. The above properties of the term with $`r_{nm}`$ allow us to call the associated phenomena nonlinear interference effects.
We can regard (2.24) as the difference between the number of acts of emission and absorption of the $`\mathrm{}\omega _\mu `$ photon. All the terms of $`w_{nj}`$ except $`\rho _{jj}`$ determine emission processes. Conversely terms associated with $`\rho _{jj}`$ control the weak field energy absorption rate. According to (2.24) only the level splitting effect stands out in the absorption probability<sup></sup>. This is due to the fact that absorption corresponds to the transition from the unexcited level $`j`$ to excited level $`n`$. NIE is due to the reverse transition from an excited to unexcited state, i.e., in the case when $`nj`$ are contained only in the emission. Therefore the line shapes of pure emission and absorption turn out to be different due to NIE. The sign of their difference, i.e., of $`w_{nj}`$, is determined not only by the sign of population difference $`\rho _{nn}\rho _{jj}`$; in particular the sign of $`w_{nj}`$ can change with the change of $`\mathrm{\Omega }_\mu `$ <sup>\[7-9\]</sup>.
Equation (2.24) makes it possible to analyze also spontaneous emission. For this purpose it is merely necessary to drop the term $`\rho _{jj}`$ from (2.24) and replace $`|G_\mu |^2`$ by a quantity corresponding to the atomic interaction with zero oscillations of the field<sup></sup>: $`\gamma _{nj}(8\pi ^2)^1\mathrm{\Delta }\mathrm{\Omega }_\mu \mathrm{\Delta }O`$. Equations for other transitions are of the same type and can be obtained from (2.24) by a simple substitution of indices and signs. For example, $`w_{ml}`$ is obtained from the substitutions $`mn`$, $`jl`$, and $`\mathrm{\Omega }^{}\mathrm{\Omega }^{}`$.
## 3 EMISSION AND ABSORPTION LINE SHAPE IN TRAVELING MONOCHROMATIC WAVE FIELD
We analyze the role of nonequilibrium velocity distribution and nonlinear interference effects. We consider first two directions of $`k_\mu `$ in detail: along and against k. The value of $`w_{nj}`$ averaged over $`\mathrm{v}`$ for these two directions is
$`w_{nj}^\pm =2\mathrm{}\omega _{nj}|G_\mu |^2{\displaystyle \frac{\sqrt{\pi }}{k\overline{v}}}\mathrm{exp}\{{\displaystyle \frac{\mathrm{\Omega }_\mu ^2}{(k_\mu \overline{v})^2}}\}\times `$ (3.1)
$`\times \{N_nN_j+(N_mN_n)Re[F_\pm (\mathrm{\Omega }_\mu )+f_\pm (\mathrm{\Omega }_\mu )]\},`$ (3.2)
$`F_\pm +f_\pm ={\displaystyle \frac{k_\mu }{k}}{\displaystyle \frac{2G^2}{\sqrt{1+\ae }}}\times `$ (3.3)
$`\times {\displaystyle \frac{\mathrm{\Gamma }_n^1(1\gamma _{mn}/\mathrm{\Gamma }_m)[\mathrm{\Gamma }_\pm +iz]+[1\pm \sqrt{1+\ae }]/2}{[\mathrm{\Gamma }_0+iz][\mathrm{\Gamma }_\pm +iz]+G^2}},`$ (3.4)
$`z=\mathrm{\Omega }_\mu \mathrm{\Omega }k_\mu /k,\mathrm{\Gamma }_0=\mathrm{\Gamma }_{jn}+\mathrm{\Gamma }_Bk_\mu /k,`$ (3.5)
$`\mathrm{\Gamma }_\pm =\mathrm{\Gamma }_{jm}+\mathrm{\Gamma }_B(k_\mu /k\pm 1),\mathrm{\Gamma }_B=\mathrm{\Gamma }\sqrt{1+\ae }.`$ (3.6)
The signs + and - in (3.3) correspond to $`k_\mu `$ directed along and against k; $`f_\pm `$ and $`F_\pm `$ represent the interference term and a term due to the nonequilibrium addition to the velocity distribution, respectively. Equation (3.3) is not applicable if $`k_\mu <k`$ and $`𝐤_\mu 𝐤<0`$. Velocity averaging can be performed also in this case. However the obtained expression can be used to some extent in the analysis only if $`\ae `$ is small Then (3.3) is valid if $`\mathrm{\Gamma }_{}`$ is replaced by $`\mathrm{\Gamma }_{jm}k_\mu /k+(1k_\mu /k)\mathrm{\Gamma }_{jn}`$, $`G=0`$ and $`\ae =0`$ everywhere (except for the common factor $`G^2`$), and $`[1+\sqrt{1+\ae }]/2`$ is replaced by $`k_\mu /k`$.
A comparison of (3.3) with (2.24) shows that $`w_{nj}`$ has the same formal structure as the corresponding expression for the fixed atom whose resonant frequency is converted with respect to the Bennet distribution maximum and which has the widths $`\mathrm{\Gamma }_\pm `$ and $`\mathrm{\Gamma }_0`$ instead of $`\mathrm{\Gamma }_{jm}`$ and $`\mathrm{\Gamma }_{jn}`$ respectively. The physical meaning of $`\mathrm{\Gamma }_0`$ and $`\mathrm{\Gamma }_\pm `$ is as follows. The perturbation theory distinguishes between step-wise and two-photon processes whose line shape is determined by the factors $`<[\mathrm{\Gamma }_{jn}+i(\mathrm{\Omega }_\mu 𝐤_\mu 𝐯)]^1>`$ and $`<\{\mathrm{\Gamma }_{jm}+i[(\mathrm{\Omega }_\mu +\mathrm{\Omega })(𝐤_\mu +𝐤)𝐯]\}^1>`$. In our case the averaging is carried out essentially with the Bennett distribution (since $`\mathrm{\Gamma }_Bk_\mu \overline{v}`$) and the result of the averaging is $`[\mathrm{\Gamma }_0+iz]^1`$ and $`[\mathrm{\Gamma }_\pm +iz]^1`$<sup></sup>. Consequently $`\mathrm{\Gamma }_0`$ is the line width of a step-wise transition that is the sum of the width $`\mathrm{\Gamma }_Bk_\mu /k`$ of the velocity distribution converted with respect lo Doppler shifts in the $`\omega _{nj}`$ region and the natural width $`\mathrm{\Gamma }_{jn}`$ of the $`nj`$ transition. Correspondingly $`\mathrm{\Gamma }_\pm `$ is the line width of two-photon transition consisting of the natural part $`\mathrm{\Gamma }_{jm}`$ and the Doppler part $`\mathrm{\Gamma }_B(k_\mu /k\pm 1)`$. Thus the physical meaning of the analogy between (3.3) and the line shape of an ”effective atom” is quite clear. The ”effective atom” represents the group of atoms that interact with a strong field. The ”effective atom” has the same system of terms as in Fig.1 except that the widths are changed in accordance with the Bennett distribution and frequency-correlated properties of the step-wise and two-photon processes<sup></sup>.
Just as in the case of an individual atom, the step- wise and two-photon processes in the ”effective atom” cannot be considered independently if $`G`$ is sufficiently large<sup></sup>. In fact the numerator in (3.3) contains $`G^2`$ and its expansion in terms of simple fractions
$`{\displaystyle \frac{1}{[\mathrm{\Gamma }_0+iz][\mathrm{\Gamma }_\pm +iz]+G^2}}`$ (3.7)
$`={\displaystyle \frac{1}{(z_1+iz)(z_2+iz)}}`$ (3.8)
$`={\displaystyle \frac{1}{z_1z_2}}\left[{\displaystyle \frac{1}{z_2+iz}}{\displaystyle \frac{1}{z_1+iz}}\right].`$ (3.9)
$`z_{1,2}=1/2\{\mathrm{\Gamma }_0+\mathrm{\Gamma }_\pm \pm \sqrt{(\mathrm{\Gamma }_0\mathrm{\Gamma }_\pm )^2}4G^2\}`$ (3.10)
yields resonant numerators with $`z_1,z_2`$ rather than with $`\mathrm{\Gamma }_0,\mathrm{\Gamma }_\pm `$. Under certain conditions the radical in (3.3) can turn out to be imaginary, which would correspond to the splitting of the levels of an effective atom.
Equation (3.3) shows that when $`\gamma _{mn}=\mathrm{\Gamma }_m`$ the effect of velocity distribution variation is completely eliminated and only the NIE remains. The physical meaning of this is quite clear. The external field transfers some atoms from the upper level to the lower; at the same time however the relaxation transition is reduced by the same quantity since there are no other channels of decay from the upper level. On the other hand the polarization stimulated by the field at the transition $`mn`$ does not turn to zero (see (2.24), expression for $`r_{nm}`$ and NIE remains unchanged. The transition $`6p^1P_2^07s^3S_1`$ of mercury, $`\lambda =1.529`$, at which generation was observed<sup></sup> can serve as an example of a case in which the condition $`\gamma _{mn}=\mathrm{\Gamma }_m`$ is valid.
The interference effect. We examine the interference term $`f_\pm (\mathrm{\Omega }_\mu )`$ in greater detail. Based on (3.3) and (3.10) we have
$$f_\pm (\mathrm{\Omega }_\mu )=\frac{k_\mu }{k}\frac{G^2}{\sqrt{1+\ae }}\frac{1\sqrt{1+\ae }}{z_1z_2}\left[\frac{1}{z_2+iz}\frac{1}{z_1+iz}\right].$$
(3.11)
The line contour of $`Re[f_\pm (\mathrm{\Omega }_\mu )]`$ has the simplest shape when $`z_{1,2}`$ are real. In this case it follows from (3.11) that the function $`Ref_\pm `$ changes sign in going from the center of the line to the wings. The sign of $`Ref_+`$ at the point $`z=0`$ is determined by the factor $`1+\sqrt{1+\ae }`$ and depends therefore on the relative direction $`𝐤_\mu `$ and $`𝐤`$. When $`𝐤_\mu 𝐤>0`$ the value in the center is negative and in the opposite direction it is positive. When the values of the external field are small ($`\ae 1`$) we have $`Ref_+\ae ^2`$ and $`Ref_{}\ae `$.
The function
$$f_\pm (z)=\left[\frac{k_\mu }{k}\frac{G^2}{\sqrt{1+\ae }}\frac{1\sqrt{1+\ae }}{z_1^2}\right]^1Ref_\pm (z)$$
is illustrated in Fig.2 for $`z_1/z_2=1;2.5;5`$. According to Fig. 2 the graphs have an approximately similar shape (the positive maximum in the center and broad negative wings) for any values of $`z_1/z_2`$. However the larger $`z_1/z_2`$ the narrower and more intense the maximum. When $`z_2z_1`$ its width is approximately equal to $`z_2`$ and its intensity in the center is proportional to $`z_2^1`$. This case seems to be the most interesting from the practical point of view.
We consider the conditions for which the relation $`z_2z_1`$ is valid. For the ”interference” direction $`𝐤_\mu 𝐤<0`$, in which the effect is sharper, the expressions for $`z_{1,2}`$ can be represented in the form
$`z_{1,2}={\displaystyle \frac{1}{2}}\{\mathrm{\Gamma }_{jn}+\mathrm{\Gamma }_{jm}+\mathrm{\Gamma }_B({\displaystyle \frac{2k_\mu }{k}}1)`$ (3.12)
$`\pm \sqrt{(\mathrm{\Gamma }_B+\mathrm{\Gamma }_{jn}\mathrm{\Gamma }_{jm})^24G^2}\}.`$ (3.13)
According to this formula the absence of splitting and the considerable difference between $`z_1`$ and $`z_2`$, are due to the conditions
$$\mathrm{\Gamma }+\mathrm{\Gamma }_{jn}\mathrm{\Gamma }_{jm},k_\mu k,\mathrm{\Gamma }^2\ae /G^2=(\mathrm{\Gamma }\tau )^21.$$
(3.14)
Here the radical in (3.12) can be expanded into a series:
$`z_1=\mathrm{\Gamma }_{jn}+\mathrm{\Gamma }{\displaystyle \frac{k_\mu }{k}}\sqrt{1+\ae }{\displaystyle \frac{\ae /\tau ^2}{\mathrm{\Gamma }\sqrt{1+\ae }+\mathrm{\Gamma }_{jn}\mathrm{\Gamma }_{jm}}},`$ (3.15)
$`z_2=\mathrm{\Gamma }_{jm}+\mathrm{\Gamma }\left({\displaystyle \frac{k_\mu }{k}}1\right)\sqrt{1+\ae }+{\displaystyle \frac{\ae /\tau ^2}{\mathrm{\Gamma }\sqrt{1+\ae }+\mathrm{\Gamma }_{jn}\mathrm{\Gamma }_{jm}}}.`$ (3.16)
We see from (3.15) that the minimum value of $`z_2`$ equals the line width of the forbidden transition $`\mathrm{\Gamma }_{jm}`$. In many cases we can expect that $`\mathrm{\Gamma }_{jm}\mathrm{\Gamma }jn`$. Consequently the emission spectrum at the transition $`jn`$ can contain a structure with a considerably smaller width than is typical of the given transition. The value of $`z_2`$ increases with the field but much slower than $`z_1`$ when $`(k_\mu k)/k1`$.
The amplitude of the interference term
$$f_{}(0)=\frac{k_\mu }{k}\frac{1+\sqrt{1+\ae }}{\sqrt{1+\ae }}\frac{G^2}{z_1z_2}=\frac{k_\mu }{k}\frac{1+\sqrt{1+\ae }}{\sqrt{1+\ae }}\frac{G^2}{\mathrm{\Gamma }_0\mathrm{\Gamma }_{}+G^2}$$
(3.17)
as a function of $`G^2`$ is a curve with saturation where one half of the maximum value is reached approximately for $`G^2=\mathrm{\Gamma }_0\mathrm{\Gamma }_{}`$. Therefore the ratio $`G^2/\mathrm{\Gamma }_0\mathrm{\Gamma }_{}\ae _{}`$ can be interpreted as the saturation parameter of the effective atom. If $`(k_\mu k)/k1`$ and $`\mathrm{\Gamma }_0\mathrm{\Gamma }_{}`$, the width $`z_2\mathrm{\Gamma }_{jm}[1+\ae _{}]`$ is also determined by the quantity $`\ae _{}`$. We note that $`\ae _{}<\ae `$. In tact, according to (3.14) and (2.22)
$`{\displaystyle \frac{\ae }{\ae _{}}}=\mathrm{\Gamma }_0\mathrm{\Gamma }_{}\tau ^2=2[\mathrm{\Gamma }_{jm}+({\displaystyle \frac{k_\mu }{k}}1)\mathrm{\Gamma }\sqrt{1+\ae }]\times `$ (3.18)
$`\left[\mathrm{\Gamma }_{jn}+{\displaystyle \frac{k_\mu }{k}}\mathrm{\Gamma }\sqrt{1+\ae }\right]{\displaystyle \frac{\mathrm{\Gamma }_m+\mathrm{\Gamma }_n\gamma _{mn}}{\mathrm{\Gamma }_m\mathrm{\Gamma }\mathrm{\Gamma }_n}}.`$ (3.19)
By virtue of the obvious inequalities $`2\mathrm{\Gamma }_{}>\mathrm{\Gamma }_m`$, $`\mathrm{\Gamma }_0>\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }_m+\mathrm{\Gamma }_n\gamma _{mn}>\mathrm{\Gamma }_n`$, the right-hand side in (3.19) is larger than unity. Therefore as $`G^2`$ increases the population difference in the center of the Bennett distribution is equalized first since it is proportional to $`\ae /(1+\ae )`$. The amplitude of the interference term is determined by the ratio $`\ae _{}/(1+\ae _{})`$, retains its linear dependence up to large values of $`G^2`$, and becomes saturated at $`\ae _{}1`$. At the same time the width of the central maximum increases, becoming twice as large at $`\ae _{}=1`$ at the same value of the field.
We now consider the behavior of the interference term when $`𝐤_\mu `$ is parallel to $`𝐤`$. We first show that $`z_1`$ and $`z_2`$ cannot differ significantly in this case. In fact, it follows from (3.10) that $`z_1`$ and $`z_2`$ differ sharply if $`\mathrm{\Gamma }_0+\mathrm{\Gamma }_+`$$`\mathrm{\Gamma }_0\mathrm{\Gamma }_+`$ or $`\mathrm{\Gamma }_0+\mathrm{\Gamma }_+`$$`\mathrm{\Gamma }_+\mathrm{\Gamma }_0`$. These conditions in turn are equivalent to the inequality systems (see (3.5)) $`\mathrm{\Gamma }_{jm}\mathrm{\Gamma }`$, $`\mathrm{\Gamma }_{jm}\mathrm{\Gamma }_{jn}`$ or $`\mathrm{\Gamma }_{jn}\mathrm{\Gamma }`$, $`\mathrm{\Gamma }_{jn}\mathrm{\Gamma }_{jm}`$ which can be readily shown to be invalid in spontaneous relaxation and in impact broadening of lines. Consequently the roots $`z_1`$ and $`z_2`$ are of the same order of magnitude in the direction $`𝐤_\mu 𝐤>0`$ and the structure is relatively not sharp. According to (3.11) the amplitude $`f_+(0)`$ is
$$f_+(0)=\frac{k_\mu }{k}\frac{\sqrt{1+\ae }1}{\sqrt{1+\ae }}\frac{G^2}{\mathrm{\Gamma }_0\mathrm{\Gamma }_++G^2}.$$
(3.20)
Comparing (3.20) and (3.17) we see that $`|f_+(0)|<f_{}(0)`$, i.e., the amplitude of the structure in the direction $`𝐤_\mu 𝐤>0`$ is always smaller than for $`𝐤_\mu 𝐤<0`$.
So far we considered $`z_1,z_2`$ to be real. Now let
$`z_{1,2}=z_0+i\zeta ,z_0=(\mathrm{\Gamma }_0+\mathrm{\Gamma }_\pm )/2,`$ (3.21)
$`\zeta =\sqrt{G^2(\mathrm{\Gamma }_0\mathrm{\Gamma }_\pm )^2/4},`$ (3.22)
$`Ref_\pm (z)={\displaystyle \frac{k_\mu }{k}}{\displaystyle \frac{1\sqrt{1+\ae }}{\sqrt{1+\ae }}}{\displaystyle \frac{G^2}{2\zeta }}\times `$ (3.23)
$`\times \left[{\displaystyle \frac{z+\zeta }{z_0^2+(z+\zeta )^2}}{\displaystyle \frac{z\zeta }{z_0^2+(z\zeta )^2}}\right].`$ (3.24)
The general shape of the graph $`Ref_\pm `$ depends on the ratio $`\zeta /z_0`$, as is apparent from Fig.5. When $`\zeta /z_0`$ is small the contours are qualitatively indistinguishable from the case of real, but similar, $`z_1,z_2`$ (see curves 1 and 2 in Fig.5). It is of interest therefore to determine the maximum possible values for the ratio $`\zeta /z_0`$S. We can show using (3.21) and (3.3) that under the most favorable conditions $`\zeta \sqrt{3}z_0`$. The curve in Fig.5 corresponding to $`\zeta =\sqrt{3}z_0`$ indicates the maximum effect of line splitting. The ”fuzzy” splitting of the interference term has a physical meaning: the increasing $`G^2`$ is accompanied by a rise in the atomic level splitting occurring together, however, with an increase in the line widths of effective atom, $`\mathrm{\Gamma }_0`$, and $`\mathrm{\Gamma }_\pm `$ due to the broadening of Bennett distribution (see (2.22)). Nevertheless we can observe level splitting even with a large Doppler broadening since the shape of curve 3 in Fig.5 is still significantly different from the others.
Nonequilibrium addition to the velocity distribution. We turn to the term $`F_\pm (\mathrm{\Omega }_\mu )`$ in (3.3):
$`F_\pm (\mathrm{\Omega }_\mu )={\displaystyle \frac{k_\mu }{k}}\mathrm{\Gamma }_n^1(1{\displaystyle \frac{\gamma _{mn}}{\mathrm{\Gamma }_m}}){\displaystyle \frac{G^2}{\mathrm{\Gamma }_n\sqrt{1+\ae }}}{\displaystyle \frac{1}{z_1z_2}}\times `$ (3.25)
$`\times \left[{\displaystyle \frac{z_1\mathrm{\Gamma }_\pm }{z_1+iz}}{\displaystyle \frac{z_2\mathrm{\Gamma }_\pm }{z_2+iz}}\right].`$ (3.26)
In the case of real $`z_{1,2}`$ the sign of $`z_1\mathrm{\Gamma }_\pm `$ and $`z_2\mathrm{\Gamma }_\pm `$ is the same but depends on the sign of $`\mathrm{\Gamma }_0\mathrm{\Gamma }_\pm `$. If $`\mathrm{\Gamma }_0>\mathrm{\Gamma }_\pm `$ then $`z_{1,2}\mathrm{\Gamma }_\pm >0`$; on the other hand, if $`\mathrm{\Gamma }_0<\mathrm{\Gamma }_\pm `$ then $`z_{1,2}\mathrm{\Gamma }_\pm <0`$ (see (3.10)). According to Fig.5 of particular interest is the case of strongly different $`z_1`$ and $`z_2`$ when $`Re[F_\pm (z)]`$ has the form of a broad dispersive contour (the width $`z_1`$) with a sharp notch (or spike) in the center (the width of $`z_2z_1`$). The conditions that allow for $`z_1z_2`$ were analyzed above. We note that $`z_2z_1`$ can be realized when $`𝐤_\mu 𝐤<0`$.
If $`z_{1,2}`$ are complex, $`Re[F_\pm (\mathrm{\Omega }_\mu )]`$ has the form
$`Re[F_\pm (\mathrm{\Omega }_\mu )]={\displaystyle \frac{k_\mu }{k}}{\displaystyle \frac{z_0}{\mathrm{\Gamma }_n}}(1{\displaystyle \frac{\gamma _{mn}}{\mathrm{\Gamma }_m}}){\displaystyle \frac{G^2}{\sqrt{1+\ae }}}\{{\displaystyle \frac{1}{z_0^2+(z+\zeta )^2}}`$ (3.27)
$`+{\displaystyle \frac{1}{z_0^2+(z\zeta )^2}}{\displaystyle \frac{\mathrm{\Gamma }_0\mathrm{\Gamma }_\pm }{\mathrm{\Gamma }_0+\mathrm{\Gamma }_\pm }}{\displaystyle \frac{1}{\zeta }}[{\displaystyle \frac{z+\zeta }{z_0^2+(z+\zeta )^2}}{\displaystyle \frac{z\zeta }{z_0^2+(z\zeta )^2}}]\}.`$ (3.28)
In contrast to (3.23) the possibility to observe splitting is determined now not only by the ratio $`\zeta /z_0`$ but also by the magnitude and sign of the factor $`(\mathrm{\Gamma }_0\mathrm{\Gamma }_\pm )/(\mathrm{\Gamma }_0+\mathrm{\Gamma }_\pm )`$. From (3.5) for $`\mathrm{\Gamma }_0,\mathrm{\Gamma }_\pm `$ we can see that $`1<s(\mathrm{\Gamma }_0\mathrm{\Gamma }_\pm )/(\mathrm{\Gamma }_0+\mathrm{\Gamma }_\pm )<1`$. Figure 5 shows plots of
$$F_\pm =\left[\frac{k_\mu }{k}\left(1\frac{\gamma _{mn}}{\mathrm{\Gamma }_m}\right)\frac{G^2}{z_0\mathrm{\Gamma }_n\sqrt{1+\ae }}\right]ReF_\pm $$
for the limiting values of the factor $`s`$ and for $`\mathrm{\Gamma }_0=\mathrm{\Gamma }_\pm `$. According to Fig.5, a sharply defined splitting effect can occur even with $`\zeta =z_0`$ which is less than the possible limit of $`\zeta z_0\sqrt{3}`$. Particularly significant is curve 3 in Fig.5 according to which the intensity is much lower in the center than in the side maxima. Using (3.28) we can obtain for $`\zeta =z_0`$, $`𝐤_\mu 𝐤<0`$ and $`k_\mu =k`$:
$$\frac{Re[F_{}(0)]}{Re[F_{}(\zeta )]}=\frac{5}{2}\frac{\mathrm{\Gamma }_{}}{\mathrm{\Gamma }_0+2\mathrm{\Gamma }_{}}\frac{5}{2}\frac{\mathrm{\Gamma }_{jm}}{\mathrm{\Gamma }_{jn}+2\mathrm{\Gamma }_{jm}+\mathrm{\Gamma }\sqrt{1+\ae }}.$$
(3.29)
Consequently if $`\mathrm{\Gamma }_{jn}+\mathrm{\Gamma }_B\mathrm{\Gamma }_{jm}`$, the ratio (3.29) is much smaller than unity. The condition $`\mathrm{\Gamma }_0\mathrm{\Gamma }_{}`$ corresponds to the value $`s=1`$ and it can be satisfied for $`\mathrm{\Gamma }\sqrt{1+\ae }\mathrm{\Gamma }_{jm}`$.
Comparison of $`F_\pm (\mathrm{\Omega }_\mu )`$ and $`f_\pm (\mathrm{\Omega }_\mu )`$. It is clear from the preceding discussion that the frequency dependences of $`F_\pm `$ and $`f_\pm `$ are similar in general and in some cases one term can emphasize or, conversely, concentrate the effects contributed by the other.
We now consider the properties of the sum $`F_\pm `$ and $`f_\pm `$ and determine the weight of each of the two terms. We begin with the case of real roots $`z_{1,2}`$. In this case the curves $`Re[F_\pm (z)]`$ and $`Re[f_\pm (z)]`$ are of the same type throughout and we may limit the analysis to a single point $`z=0`$ (maximum or minimum). From (3.5) and (3.10) we find
$`Re[F_\pm (0)+f_\pm (0)]={\displaystyle \frac{k_\mu }{k}}{\displaystyle \frac{G^2\mathrm{\Gamma }_\pm }{z_1z_2\sqrt{1+\ae }}}[{\displaystyle \frac{2}{\mathrm{\Gamma }_n}}(1{\displaystyle \frac{\gamma _{mn}}{\mathrm{\Gamma }_m}})+`$ (3.30)
$`+{\displaystyle \frac{1}{\mathrm{\Gamma }_\pm }}(1\sqrt{1+\ae })].`$ (3.31)
The first term in the brackets is associated with $`f_\pm `$ and the second with $`f_\pm `$. The appearance of the factors $`1/\mathrm{\Gamma }_n`$ and $`1/\mathrm{\Gamma }_\pm `$ is understandable: $`1/\mathrm{\Gamma }_n`$ determines the time of interaction of an atom at the n level with the field. An analog of such an ”accumulation time” for the interference term is the quantity $`1/\mathrm{\Gamma }_\pm `$.
In addition to the factor $`1\gamma _{mn}/\mathrm{\Gamma }_m`$, whose role was discussed above, the relation between $`F_\pm (0)`$ and $`f_\pm (0)`$ depends on the relaxation constants, field amplitude, direction of observation, and the ratio $`k_\mu /k`$. To observe NIE even with $`\gamma _{mn}\mathrm{\Gamma }_m`$ the most convenient conditions obtain when $`k_\mu =k`$ and $`\mathrm{\Gamma }_{jm}\mathrm{\Gamma }_n`$; furthermore its role increases with the rise in field intensity. Conversely when $`𝐤`$ and $`𝐤_\mu `$ are parallel we can expect an almost complete elimination of NIE because the inequality $`\mathrm{\Gamma }_+\mathrm{\Gamma }_n[\sqrt{1+\ae }1]/2`$ can be assured by $`\mathrm{\Gamma }_{jm}\mathrm{\Gamma }_n`$, $`\mathrm{\Gamma }\mathrm{\Gamma }_n`$, $`\ae 1`$ and $`k_\mu >k`$. Therefore $`Re[F_\pm ]`$ as well as $`Re[f_\pm ]`$ can be predominant depending on the values of the numerous variable parameters.
If $`z_{1,2}`$ are complex the expression for $`Re[F_\pm +f_\pm ]`$ differs from (3.28) only by the substitution of factor $`s`$
$$c=\frac{\mathrm{\Gamma }_0\mathrm{\Gamma }_\pm }{\mathrm{\Gamma }_0+\mathrm{\Gamma }_\pm }\frac{\mathrm{\Gamma }_n}{\mathrm{\Gamma }_0+\mathrm{\Gamma }_\pm }\left(1\frac{\gamma _{mn}}{\mathrm{\Gamma }_m}\right)^1[1\sqrt{1+\ae }],$$
(3.32)
where the second term reflects the role of $`Re[f_\pm ]`$. We can show that the value of $`c`$ varies between +1 and -1 Therefore the total contour can be deformed within the same limits as $`Re[F_\pm ]`$ (see Fig.5).
We now consider $`w_{nj}`$ for the intermediate values of the angle $`\theta `$ between $`𝐤`$ and $`𝐤_\mu `$. We denote the velocity component perpendicular to $`𝐤`$ by $`𝐮`$:
$$\mathrm{\Omega }_\mu ^{}=\mathrm{\Omega }_\mu \mathrm{𝐤𝐮}\mathrm{sin}\theta 𝐤_\mu 𝐯\mathrm{cos}\theta ,\mathrm{\Omega }^{}=\mathrm{\Omega }\mathrm{𝐤𝐯}.$$
(3.33)
According to (3.33) the averaging with respect to $`v`$ leads as before to (3.5), except that $`𝐤_\mu `$ must be replaced by $`k_\mu \mathrm{cos}\theta `$ (apart from the common factor in $`F_\pm `$ and $`f_\pm `$) and $`\mathrm{\Omega }_\mu `$ by $`\mathrm{\Omega }_\mu ku\mathrm{sin}\theta `$. The subsequent averaging with respect to $`u`$ can be carried out although only its result is given here When the angles are small, $`\theta \mathrm{\Gamma }_+/k\overline{v}`$, $`\mathrm{\Gamma }_0/k\overline{v}`$, there is practically no variation of $`w_{nj}`$.
The same consideration applies to the angles $`|\pi \theta |\mathrm{\Gamma }_{}/k\overline{v}`$, $`\mathrm{\Gamma }_0/k\overline{v}`$. When $`|\theta |`$ (or $`|\pi \theta |`$) increases above the indicated values the spectral width of the functions $`F_\pm ,f_\pm `$ increases approximately as $`k\overline{v}\mathrm{sin}\theta `$ and reaches the full Doppler width when $`\theta \pi /2`$. Since the integrated intensity of the correction to $`w_{nj}`$ due to strong field does not depend on $`\theta `$, the amplitude of this correction is $`k\overline{v}/\mathrm{\Gamma }_0`$ times lower than in the above cases. All these phenomena are due to the fact that the strong field represents a plane monochromatic wave and causes changes in the distribution of only one velocity component Therefore the case of $`\theta =0`$ and the adjacent directions of $`k_\mu `$ is the most interesting one.
Our analysis deals with the case where both fields represent plane traveling waves. The experimenter may find it convenient to use a strong field within the resonator of a suitable gas laser<sup></sup>. The strong field then has the form of a standing wave and the pattern of events is somewhat different. When the departure from resonance in the strong field is greater than the width of Bennett distribution ($`|\mathrm{\Omega }|>\mathrm{\Gamma }_0,\mathrm{\Gamma }_\pm `$), one can regard the two traveling waves as fully independent because they interact with different groups of atoms. Therefore the expression for $`w_{nj}`$ now contains, instead of $`F_+(\mathrm{\Omega }_\mu )+f_+(\mathrm{\Omega }_\mu )`$ or $`F_{}(\mathrm{\Omega }_\mu )+f_{}(\mathrm{\Omega }_\mu )`$, the sum of these terms
$$F_+(\mathrm{\Omega }_\mu )+f_+(\mathrm{\Omega }_\mu )+F_{}(\mathrm{\Omega }_\mu )+f_{}(\mathrm{\Omega }_\mu ).$$
(3.34)
All the singularities of the terms with indices + or - are now at the distance $`\pm \mathrm{\Omega }k_\mu /k`$ from the line center (see definition of $`z`$ in (3.5)) and they overlap. Thus all that we said for the case of a strong field in the form of a traveling wave remains valid for that of a standing wave. At the same time different frequencies should produce effects corresponding to ”interference” and ”non-interference” directions.
On the other hand if the condition $`|\mathrm{\Omega }|>\mathrm{\Gamma }`$ does not hold, the Bennett distributions stemming from two opposed waves overlap and we have a different situation. We can say that the additive property of nonlinear effects due to opposed waves appears a priori in the first approximation (with respect to $`G^2`$), i.e., (3.34) is valid if $`G^2`$ is left in the expression for $`F_\pm +f_\pm `$ only in the form of a common factor. The invariance of (3.3) in successive approximations with respect to $`G^2`$’ is due to the fact that large fields generate a spatial inhomogeneity of the medium (with a period of $`\lambda /2`$)<sup></sup>. Consequently the atomic probability amplitudes are subject to a form of phase modulation and the atomic levels are split into a number of sublevels larger than the two sublevels typical of the traveling wave. The above modulation was Investigated in <sup></sup> in the case of resonance fluorescence and it was found that the emission spectrum changed significantly.
## 4 GENERATION IN THE PRESENCE OF EXTERNAL FIELD
In Secs. 2 and 3 the fields that resonated with transitions $`nj`$, $`gn`$, etc., were considered weak (Fig.1). Experiments<sup></sup> showed that generation at these transitions was a convenient method of studying NIE. Therefore we now consider generation at the $`gn`$ transition (since it was studied in<sup></sup>) The unsaturated (with respect to $`G_\mu `$) gain at the $`gn`$ transition changes in an external field $`G`$ that is resonant with $`mn`$ (see Sec. 3). To compute the generation power at $`gn`$ we must know the saturation function of the $`gn`$ transition We can show that once the conditions
$$|N_mN_n|\frac{G^2}{\mathrm{\Gamma }^2}|N_gN_n|,|N_mN_n|\frac{G^4}{\mathrm{\Gamma }^4}|N_gN_n|\frac{G_\mu ^2}{\mathrm{\Gamma }^2}$$
(4.1)
are satisfied, saturation at the $`gn`$ transition is the same as in the case of $`G=0`$. Therefore the generation power is determined by the standard formula
$`{\displaystyle \frac{\mathrm{\Gamma }_n+\mathrm{\Gamma }_g\gamma _{ng}}{\mathrm{\Gamma }_n\mathrm{\Gamma }_g\mathrm{\Gamma }_{ng}}}G_\mu ^2=[1{\displaystyle \frac{\mathrm{\Delta }N\mathrm{exp}\{\mathrm{\Omega }_\mu ^2/(k_\mu \overline{v})^2\}+\alpha }{N_gN_n}}]\times `$ (4.2)
$`\times \left[1+{\displaystyle \frac{\mathrm{\Gamma }_{ng}^2}{\mathrm{\Gamma }_{ng}^2+\mathrm{\Omega }_\mu ^2}}\right]^1;`$ (4.3)
$`\alpha ={\displaystyle \frac{k_\mu }{k}}(N_mN_n)G^2\{{\displaystyle \frac{1\gamma _{mn}/\mathrm{\Gamma }_m}{\mathrm{\Gamma }_n\mathrm{\Gamma }_0}}[{\displaystyle \frac{\mathrm{\Gamma }_0^2}{\mathrm{\Gamma }_0^2+(\mathrm{\Omega }_\mu +k_\mu \mathrm{\Omega }/k)^2}}`$ (4.4)
$`+{\displaystyle \frac{\mathrm{\Gamma }_0^2}{\mathrm{\Gamma }_0^2+(\mathrm{\Omega }_\mu k_\mu \mathrm{\Omega }/k)^2}}]+{\displaystyle \frac{1}{\mathrm{\Gamma }+\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm}}}\times `$ (4.5)
$`\times [{\displaystyle \frac{\mathrm{\Gamma }_+}{\mathrm{\Gamma }_+^2+(\mathrm{\Omega }_\mu k_\mu \mathrm{\Omega }/k)^2}}{\displaystyle \frac{\mathrm{\Gamma }_0}{\mathrm{\Gamma }_0^2+(\mathrm{\Omega }_\mu k_\mu \mathrm{\Omega }/k)^2}}]\},`$ (4.6)
$`\mathrm{\Gamma }_0=\mathrm{\Gamma }_{gn}+{\displaystyle \frac{k_\mu }{k}}\mathrm{\Gamma },\mathrm{\Gamma }_+=\{\begin{array}{cc}\mathrm{\Gamma }_{gm}k_\mu /k+(1k_\mu /k)\mathrm{\Gamma }_{gn},\hfill & k_\mu <k\hfill \\ \mathrm{\Gamma }_{gm}+(k_\mu /k1)\mathrm{\Gamma },\hfill & k_\mu >k\hfill \end{array}`$ (4.7)
where $`\mathrm{\Delta }N`$ is the threshold population difference for $`G=0`$ and $`\mathrm{\Omega }_\mu =0`$. In the absence of the external field (4.2) determines the usual dependence of power on $`\mathrm{\Omega }_\mu `$ with the ”Lamb dip”. The term $`\alpha `$ introduces an additional spectral structure.
We consider the case when the role of atomic collisions is small, so that $`\mathrm{\Gamma }+\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm}=\mathrm{\Gamma }_n`$. A ”spike” or a ”dip” (depending on the sign of $`N_mN_n`$) then appears at the frequency $`\mathrm{\Omega }_\mu =\mathrm{\Omega }k_\mu /k`$
$$I_{}=\frac{N_mN_n}{N_nN_g}\frac{k_\mu }{k}\frac{|G|^2}{\mathrm{\Gamma }_n\mathrm{\Gamma }_0}\left(1\frac{\gamma _{mn}}{\mathrm{\Gamma }_m}\right)\frac{\mathrm{\Gamma }_0^2}{\mathrm{\Gamma }_0^2+(\mathrm{\Omega }_\mu +k_\mu \mathrm{\Omega }/k)^2}.$$
(4.8)
Another ”spike” or ”dip” appears at $`\mathrm{\Omega }_\mu =k_\mu \mathrm{\Omega }/k`$ (Fig.6).
$`I_+={\displaystyle \frac{N_mN_n}{N_nN_g}}{\displaystyle \frac{k_\mu }{k}}{\displaystyle \frac{|G|^2}{\mathrm{\Gamma }_n\mathrm{\Gamma }_0}}[{\displaystyle \frac{\mathrm{\Gamma }_0}{\mathrm{\Gamma }_+}}{\displaystyle \frac{\mathrm{\Gamma }_+^2}{\mathrm{\Gamma }_+^2+(\mathrm{\Omega }_\mu k_\mu \mathrm{\Omega }/k)^2}}`$ (4.9)
$`{\displaystyle \frac{\gamma _{mn}}{\mathrm{\Gamma }_m}}{\displaystyle \frac{\mathrm{\Gamma }_0^2}{\mathrm{\Gamma }_0^2+(\mathrm{\Omega }_\mu k_\mu \mathrm{\Omega }/k)^2}}];`$ (4.10)
if $`\mathrm{\Gamma }_m,\mathrm{\Gamma }_g\mathrm{\Gamma }_n`$ and $`|1k_\mu /k|1`$ then $`\mathrm{\Gamma }_+\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }_+\mathrm{\Gamma }_{gn}`$ (see (refe4.4)). Consequently we see from (4.8) and (4.9) that in this case the ”spikes” $`I_{}`$ and $`I_+`$ differ sharply from each other in width and height. The second term in (4.9) contributes significantly only to the wings of the $`I_+`$, contour so that the width of this ”spike” is much smaller than the natural width at the $`gn`$ transition. When $`\gamma _{mn}=\mathrm{\Gamma }_m`$ the ”spike” $`I_{}`$ vanishes and only the interference ”spike” $`I_+`$, remains with singularities in the wings (a ”spike” in a ”trough”). In the other limiting case of $`\mathrm{\Gamma }_m\mathrm{\Gamma }_n,\mathrm{\Gamma }_g`$; $`\mathrm{\Gamma }_+\mathrm{\Gamma }_0`$ both spikes have the same width and vanish when $`\gamma _{mn}/\mathrm{\Gamma }_m1`$. When $`\mathrm{\Omega }=0`$ and $`\mathrm{\Gamma }_+\mathrm{\Gamma }_{gn}`$, the above singularities occur in the floor of the Lamb ”dip” as shown schematically in Fig.6.
Two generation peaks differing in width were observed in<sup></sup>. A strong frequency dependence of generation in the region $`I_+`$ can be utilized for effective output power stabilization of generation frequency.
We consider the dependence of generating emission frequency on the natural resonator frequency The generation frequency is determined by the requirement that the field phase shift in a double pass of the resonator be a multiple of $`2\pi `$. The value of the refraction index necessary to compute the phase can be found from
$$n_0=1+2\pi NRe\{r_{ng}d_{ng}\}(E_\mu /4)^1,$$
where $`E_\mu `$ is intensity of the field resonating with the $`ng`$ transition. If $`|\mathrm{\Omega }_\mu |k_\mu \overline{v}`$ the generation frequency is determined from the equation
$`\mathrm{\Omega }_r\omega _r\omega _{gn}=\mathrm{\Omega }_\mu +{\displaystyle \frac{l}{l_r}}{\displaystyle \frac{\mathrm{\Delta }\omega _r}{2}}\times \{{\displaystyle \frac{2}{\sqrt{\pi }}}{\displaystyle \frac{N_gN_n}{\mathrm{\Delta }N}}{\displaystyle \frac{\mathrm{\Omega }_\mu }{k\overline{v}}}`$ (4.11)
$`[{\displaystyle \frac{N_gN_n}{\mathrm{\Delta }N}}1]{\displaystyle \frac{\mathrm{\Omega }_\mu \mathrm{\Gamma }_{ng}}{2\mathrm{\Gamma }_{ng}^2\mathrm{\Omega }_\mu ^2}}{\displaystyle \frac{k_\mu }{k}}|G|^2{\displaystyle \frac{N_mN_n}{\mathrm{\Delta }N}}\mathrm{\Phi }(\mathrm{\Omega }_\mu )\},`$ (4.12)
where $`\omega _r`$ is the natural frequency of the resonator and
$`\mathrm{\Phi }(\mathrm{\Omega }_\mu )=(1{\displaystyle \frac{\gamma _{mn}}{\mathrm{\Gamma }_m}}){\displaystyle \frac{1}{\mathrm{\Gamma }_n}}[(\mathrm{\Omega }_\mu +{\displaystyle \frac{k_\mu }{k}}\mathrm{\Omega }{\displaystyle \frac{\mathrm{\Gamma }_0\mathrm{\Gamma }_{ng}\mathrm{\Omega }_\mu }{2\mathrm{\Gamma }_{ng}^2+\mathrm{\Omega }_\mu ^2}})\times `$
$`\times {\displaystyle \frac{1}{\mathrm{\Gamma }_0^2+(\mathrm{\Omega }_\mu +k_\mu \mathrm{\Omega }/k)^2}}+`$
$`+(\mathrm{\Omega }_\mu {\displaystyle \frac{k_\mu }{k}}\mathrm{\Omega }{\displaystyle \frac{\mathrm{\Gamma }_0\mathrm{\Gamma }_{ng}\mathrm{\Omega }_\mu }{2\mathrm{\Gamma }_{ng}^2+\mathrm{\Omega }_\mu ^2}}){\displaystyle \frac{1}{\mathrm{\Gamma }_0^2+(\mathrm{\Omega }_\mu k_\mu \mathrm{\Omega }/k)^2}}]`$
$`+{\displaystyle \frac{1}{\mathrm{\Gamma }+\mathrm{\Gamma }_{gn}\mathrm{\Gamma }_{gm}}}[(\mathrm{\Omega }_\mu {\displaystyle \frac{k_\mu }{k}}\mathrm{\Omega }{\displaystyle \frac{\mathrm{\Gamma }_+\mathrm{\Gamma }_{gn}\mathrm{\Omega }_\mu }{2\mathrm{\Gamma }_{ng}^2+\mathrm{\Omega }_\mu ^2}})\times `$
$`\times {\displaystyle \frac{1}{\mathrm{\Gamma }_+^2+(\mathrm{\Omega }_\mu k_\mu \mathrm{\Omega }/k)^2}}`$
$`(\mathrm{\Omega }_\mu {\displaystyle \frac{k_\mu }{k}}\mathrm{\Omega }{\displaystyle \frac{\mathrm{\Gamma }_0\mathrm{\Gamma }_{gn}\mathrm{\Omega }_\mu }{2\mathrm{\Gamma }_{gn}^2+\mathrm{\Omega }_\mu ^2}}){\displaystyle \frac{1}{\mathrm{\Gamma }_0^2+(\mathrm{\Omega }_\mu k_\mu \mathrm{\Omega }/k)^2}}].`$
The first term in the curved brackets of (4.11) describes the known phenomenon of ”pulling” the generation frequency by the natural resonator frequency towards the center of the atomic line. The second describes a ”repulsion” of the generation frequency from the transition frequency towards the resonator frequency proportional to the quantity $`(N_gN_n)/\mathrm{\Delta }N1`$. On the curve of $`\mathrm{\Omega }_\mu `$ as a function of $`\mathrm{\Omega }_r`$ (Fig.7) the first effect corresponds to the deviation of the $`\mathrm{\Omega }_\mu `$ asymptote from the straight line $`\omega _r\omega _{gn}=\mathrm{\Omega }_\mu `$ by an angle of the order of $`\mathrm{\Delta }\omega _r/k_\mu \overline{v}`$, and the second effect corresponds to the singularity of the order of $`\sqrt{2}\mathrm{\Gamma }_{gn}`$ near $`\mathrm{\Omega }_\mu =0`$.
We consider singularities occurring in the curve $`\mathrm{\Omega }_\mu `$ in the region of frequencies $`|\mathrm{\Omega }_\mu \pm \mathrm{\Omega }k_\mu /k|\mathrm{\Gamma }_{+,0}`$ if $`\mathrm{\Gamma }_{ng}\mathrm{\Omega }k_\mu \overline{v}`$. For a purely spontaneous relaxation and $`\gamma _{mn}\mathrm{\Gamma }_m`$ we obtain from (4.11)
$$\mathrm{\Omega }_{r^\pm }=\mathrm{\Omega }_\mu \frac{l_r}{l}\frac{\mathrm{\Delta }\omega _r}{2}\frac{N_mN_n}{\mathrm{\Delta }N}\frac{k_\mu }{k}\frac{|G|^2}{\mathrm{\Gamma }_n\mathrm{\Gamma }_{+,0}}\frac{\mathrm{\Gamma }_{+,0}(\mathrm{\Omega }_\mu \mathrm{\Omega }k_\mu /k)}{\mathrm{\Gamma }_{+,0}^2+(\mathrm{\Omega }_\mu \mathrm{\Omega }k_\mu /k)^2}.$$
(4.13)
The term proportional to $`\mathrm{\Delta }\omega _r/k_\mu \overline{v}`$ has been dropped. It appears from (4.13) that in the presence of an external field when $`\mathrm{\Omega }_\mu =\pm \mathrm{\Omega }k_\mu /k`$ the dependence of generation frequency on the natural resonator frequency increases when $`N_mN_n>0`$ and decreases when $`N_mN_n<0`$:
$$\left(\frac{d\mathrm{\Omega }_\mu }{d\mathrm{\Omega }_\mu ^\pm }\right)_{\mathrm{\Omega }_\mu =\pm k_\mu \mathrm{\Omega }/k}=\left[1\frac{\mathrm{\Delta }\omega _r}{2\mathrm{\Gamma }_\pm }\frac{l}{l_r}\frac{N_mN_n}{\mathrm{\Delta }N}\frac{k_\mu }{k}\frac{G^2}{\mathrm{\Gamma }_n\mathrm{\Gamma }_{+,0}}\right].$$
In the latter case this phenomenon can be used for passive stabilization of the generation frequency. The lower the resonator $`Q`$ the greater this effect. If $`\gamma _{mn}=\mathrm{\Gamma }_m`$ the singularity at $`\mathrm{\Omega }_\mu =\mathrm{\Omega }k_\mu /k`$ vanishes. At $`\mathrm{\Omega }=0`$ all the singularities in $`\mathrm{\Omega }_\mu `$ as a function of $`\mathrm{\Omega }_r`$ appear only when $`|\mathrm{\Omega }_\mu |\mathrm{max}\{\mathrm{\Gamma }_{ng},\mathrm{\Gamma }_0,\mathrm{\Gamma }_+\}`$. The dependence of $`\mathrm{\Omega }_\mu `$ on $`\mathrm{\Omega }_r`$ can be cumbersome in this case. However if $`\mathrm{\Gamma }_+\mathrm{\Gamma }_{ng},\mathrm{\Gamma }_0`$ the most pronounced is only the contribution from $`\mathrm{\Gamma }_+`$.
Translated by S. Kassel 97
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# A Catalogue of Galaxies in the HDF-South: Photometry and Structural ParametersThe full catalogue is available in electronic form at the CDS via anonymous ftp to cdsarc.u-strasbg.fr (130.79.128.5) or via http://cdsweb.u-strsbg.fr/Abstract.html
## 1 Introduction
The Hubble Deep Field South (HDF-S) was observed in October 1998 by the Hubble Space Telescope (HST). It is the southern counterpart of the Hubble Deep Field North (HDF-N) and shares its characteristics of depth and spatial resolution. The HDF-S is a four arcmin<sup>2</sup> survey located at RA 22$`h`$ 32$`m`$ 56$`s`$, DEC -60 33’ 02”, observed during 150 orbits with the Wide Field Planetary Camera 2 (WFPC2). The WFPC2 detector is composed by 4 chips: the Planetary Camera (PC) with higher resolution (0.05 arcsec/pixel) and smaller field of view (35$`\times `$35 arcsec<sup>2</sup>) and 3 Wide Field Cameras (WF), with a spatial resolution of 0.10 arcsec/pixel and a field of view of 77$`\times `$77 arcsec<sup>2</sup>.We will refer to the area of sky observed by the Planetary Camera as PC field, and to the one observed by the three Wide Field Cameras as WF field.
HDF-S images cover a wavelength range from the ultraviolet to near-infrared (4 broad-band filters roughly corresponding to standard UBVI). Images were taken in four filters: F300W, F450W, F606W, and F814W, by a dithering technique. Released images are a combination of all individual exposures, weighted for the background signal, and resampled to a pixel scale of 0.0398 arcsec/pixel using the ”Drizzle” package (Fruchter, Hook et al. 1997). Details about observations and data reduction may be found in the Hubble Deep Field South web page (http:// www.stsci.edu/ftp/science/hdfsouth/hdfs.html).
The main goal of the two HDF campaigns is to study the characteristics of faint galaxies and provide constraints on models explaining the formation and evolution of galaxies, in particular the excess of faint blue sources. It is therefore crucial to have a reliable catalogue, both for sources detection and photometry. In this paper we present a galaxy catalogue derived from the public version 1 images in fits format of the WFPC2 data, released on 23 November 1999, along with documentation about data reduction and absolute calibration in the different bands.
This catalogue gives for each galaxy, besides photometric information in all the optical bands, metric size, mean surface brightness, asymmetry index and light concentration indexes, all fundamental for different cosmological and evolutionary tests.
In Section 2 we describe the applied detection procedure and in Section 3 the technique used to compute magnitudes, while in Section 4 we deal with the selection criteria adopted to extract the sample. In Section 5 we summarize the estimation of the sample completeness, while we discuss the possible oversampling of sub-galactic structures, such as HII regions, in high resolution imaging in Section 6. Finally in Section 7 and 8 we describe the technique used to recover galaxy colours and structural parameters respectively.
## 2 Sources Detection
Object detection was performed across the area observed in optical bands (WFPC2 field) using the software for automatic sources detection SExtractor (Bertin and Arnouts 1996). The software firsts smoothes the images and then applies a threshold to identify peaks. As a smoothing function we have chosen a Gaussian function having FWHM equal to the one measured on the images (i.e. 0.16 arcsec) and a detection threshold of 1 sigma per pixel, with a minimum detection area equal to the seeing disk, have been adopted to pick up objects. The software allows to examine every identified group of contiguous pixels to deblend overlapping sources: two sources belonging to the same group are considered different objects if they differ of 5 magnitudes or less (*deblend-mincont*=0.01, *deblend-ntresh*=32).
This choice of low thresholds aimed at obtaining a raw catalogue not biased against very faint or small source, in fact the first sample consisted of 6093, 4747, 9850, 5229 detections in the F300W, F450W, F606W and F814W band respectively. From this sample we then extracted a catalogue optimizing completeness while minimizing the number of spurious detections by selection criteria based on the signal-to-noise ratio of sources and simulations (as explained in Section 4).
## 3 Photometry
Galaxies in this survey span a wide range both in redshift and in physical size. A fixed aperture would measure different fractions of flux for galaxies dissimilar in shape and size and at different redshifts, while isophotal magnitudes suffer from the $`(1+z)^5`$ cosmological dimming in surface brightness.
“Pseudo-total” magnitudes were therefore estimated using the method of Djorgovski et al. (1995) and Smail et al. (1995): they assigned the isophotal magnitude to sources with isophotal diameter larger than $`\theta _123`$ FWHM, while smaller sources are assigned an aperture corrected magnitude, that is the magnitude within an aperture $`\theta _1`$ corrected to the magnitude within $`\theta _2`$, larger than $`\theta _1`$: $`m=m(\theta 1)+\mathrm{\Delta }m`$, $`\mathrm{\Delta }m=<m(\theta 2)m(\theta 1)>`$. In literature the choice of $`\theta _1`$ and $`\theta _2`$ is based on multiples of the FWHM of images. In the HDF-S the excellent seeing needs a different approach, we therefore estimated magnitudes for our sample on the basis of the following steps:
* For large sources, i.e. those sources having an isophotal diameter $`D_{iso}>\theta _1`$, and for “blended” sources, as flagged by SExtractor, we selected the SExtractor isophotal corrected magnitude;
* to small sources, having $`D_{iso}<\theta _1`$, we assigned an aperture corrected magnitude (estimated within $`\theta _1`$ and then corrected by $`\mathrm{\Delta }m`$ to $`\theta _2`$, being $`\theta _1<\theta _2`$)
$`\theta _1`$ has been defined as the minimum apparent diameter of a galaxy having an effective diameter $`r_e=10`$ Kpc. Hereafter we use a $`\mathrm{\Lambda }=0`$ cosmology, with $`q_0=0.5`$ and $`H_0=50`$ kms<sup>-1</sup>Mpc<sup>-1</sup> unless differently specified. With this choice $`\theta _1=1.2`$ arcsec. Since the correction $`\mathrm{\Delta }m`$ is measured on a subsample of relatively bright galaxies, we defined an area for each band, $`A_{90}`$, such that 90$`\%`$ of galaxies belonging to the *djorg* subsample has isophotal area smaller than $`A_{90}`$. $`\theta _2`$ is defined as the diameter corresponding to a circle of area $`A_{90}`$.
At bright magnitudes Kron’s technique, based on an adaptive aperture, $`r_1=\frac{{\scriptscriptstyle rI(r)}}{{\scriptscriptstyle I(r)}}`$, gives very good results. Kron (1980) and Bertin & Arnouts (1996) demonstrated that a photometry within an adaptive aperture ($`2.5r_1`$) is expected to measure a fraction of the total flux between 0.9 and 0.94. We then chose to use Kron’s magnitude (m<sub>kr</sub>) as a reference in order to test our method: if the flux fraction measured by Kron’s technique is 0.94, the fraction estimated by our “pseudo-total” magnitude (m<sub>dj</sub>) is $`x=0.9410^{0.4(m_{kr}m_{dj})}`$. Our procedure is intended to correct the systematic underestimate (6-10$`\%`$) of total flux typical of Kron’s technique, which may be important for very faint sources.
Table 1 and Figure 1 show clearly that statistically $`\mathrm{\Delta }m`$ corrects for the flux underestimate typical of Kron’s magnitudes. Moreover our estimate of total magnitudes has a narrower distribution at low S/N than Kron magnitude, see Figure 3, and in a plot magnitude-isophotal area (Figure 2) it is not evident any discontinuity in the passage between large and small sources (i.e. sources with $`\theta >\theta _1`$ or vice versa). This test confirms the validity of our choice of $`\theta _1`$.
If not differently specified, magnitudes are expressed in the AB system, that is a system based on a spectrum which is flat in $`f_\nu `$: m=$`2.5\mathrm{log}f_\nu 48.60`$ (Oke 1974) .
## 4 Selection of the Catalogue
Our low detection threshold led to a large number of detections (more than 4500 on less than 5 arcmin<sup>2</sup>): it is therefore necessary to evaluate the number of spurious sources. The depth and coverage for each filter is not homogeneous in the field of view, due to the variety of pointings that were combined together. The image depth fades towards the edges of the area covered, as well as in a cross-shaped area between detectors which received much lower coverage than the central region of each chip. As a consequence of the decreased image quality, the outer regions of each image are less reliable: in fact sky RMS is higher than the average one and some sources near the edges revealed at inspection to be spurious. We selected reliable sources by means of criteria based on S/N ratio and comparisons with simulations.
The ”Drizzle” algorithm (Variable-Pixel Linear Reconstruction) used to combine the various pointings preserves photometry and resolution and removes the effects of geometric distortion, but it causes adjacent pixels to be correlated. The pixel-to-pixel noise ($`\sigma _{sky}`$)therefore underestimates the true noise of a larger area by a factor 1.9. The noise measured on PC field is greater by a factor $`2`$ than that measured on the WF area. The S/N is then computed by a semi-empirical model (Pozzetti et al. 1998, Williams et al. 1996): S/N=R/$`\sigma _{tot}`$, where R are net counts and $`\sigma _{tot}^2=R/(\mathrm{\Gamma }t_{exp})+21.9^2\sigma _{sky}^2A_{obj}`$ for WF sources, $`\sigma _{tot}^2=R/(\mathrm{\Gamma }t_{exp})+241.9\sigma _{sky}^2A_{obj}`$ for PC sources. In the above formulas $`\sigma _{sky}`$ is the pixel-to-pixel sky RMS, t<sub>exp</sub> is the exposure time, $`\mathrm{\Gamma }`$ is the gain expressed in electrons per ADU, A<sub>obj</sub> and A<sub>sky</sub> are respectively the object and the sky isophotal areas (in pixel) used to estimate the local background.
The former term in the sum represents the Poissonian noise due to the source, the latter estimates statistical fluctuations in the mean value of sky, in the Poissonian approximation. The factor 2 is linked to uncertainties in the determination of local background: the correct term would be $`1.9^2\mathrm{\Gamma }^2\sigma _{sky}^2A_{obj}^2/A_{sky}`$, but since A<sub>sky</sub> differs less than 30$`\%`$ by the mean value of A<sub>obj</sub>, we considered A$`{}_{sky}{}^{}`$A<sub>obj</sub>.
Our detection threshold corresponds to a minimum signal-to-noise ratio of $`S/N_{WF}=1.34`$ and $`S/N_{PC}=0.67`$ for the faintest sources detectable on the WF area and on the PC area respectively.
In Table 2 we report for each filter the zeropoint (AB magnitude, Oke 1974), the sky RMS estimated by SExtractor and the corresponding 5$`\sigma `$ magnitude limit for a point source.
We treated this problem statistically, in the hypothesis that noise is symmetrical with respect to the mean sky value. Operationally we have first created for each filter a noise frame by reversing the original images, in order to reveal the negative fluctuations and to make negative (i.e. undetectable) real sources (Saracco et al. 1999). Then we run SExtractor with the same detection parameter set used to search for sources in the original images detecting, by definition, only spurious sources. Applying a S/N=5 cut off, after removing the edges of the images, we were able to reduce the spurious contamination to a negligible fraction (4$`\%`$) on the WF area, while such a cut off is not able to reduce spurious detections to a reasonable level on the PC area being them more than 35$`\%`$. In Figure 4 the magnitude distribution of spurious sources obtained on the WF area and the PC area in the F606W band are shown. It is clear that the influence of spurious sources on the PC field is still remarkable after applying selection criteria, while the contamination is suppressed in the WF field.
Thus, to avoid introducing such a large number of spurious by the PC data, we restricted the selection of sources to the central WF area only corresponding to 4.38 arcmin<sup>2</sup>. On this area 450, 1153, 1694 and 1416 sources have been selected accordingly to the above criteria in the F300W, F450W, F606W and F814W band respectively, while the raw catalogues had 6093,4747, 9850, 5229 detections in the same bands.
In every magnitude bin we compared the number of sources in our final catalogue with the number of spurious detections in order to get the contamination of false detection, shown in Table 3.
We then removed stars from the sample by using the SExtractor morphological classifier. We defined as stars those sources brighter than I<sub>814</sub>=22 and having a value of the “stellarity” index larger than 0.9. This choice tends to underestimate stars both at faint magnitudes where no classification is considered, and at bright magnitudes where some fuzzy stars could be misclassified as galaxy. On the other hand this will ensure that our galaxy sample is not biased against compact galaxies. The star “cleaning” procedure has classified and removed 14 stars at I$`{}_{814}{}^{}<22`$ in agreement with the number of stars found in the HDF-N by Mendez et al. (1998) to this depth and in excess by a factor of two with respect to the prediction of the galaxy model of Bahcall & Soneira (1981).
## 5 Completeness Correction
Completeness correction for faint undetected sources strongly depends on the source apparent spatial structure besides to their magnitude. To estimate completeness via simulations we must account for sub-galactic structures. These could be quite common at high redshift and detectable by HST, and for the morphology which is that of the co-moving UV and B pass-bands and hence is strongly affected by star formation episodes. These features imply that “typical” profiles of galaxies are not able to well describe the shapes of a lot of galaxies in the HDF-S. Thus, in order to reproduce the manifold of shapes which characterizes sources in the HDF-S we generated a set of simulated frames by directly dimming the original frames themselves by various factors while keeping constant the RMS (Saracco et al., 2000). This procedure has allowed us to avoid any assumption on the source profile while providing an artificial fair dimmed sample in a real background noise. We thus define the correction factor $`\overline{c}`$ as the mean number of dimmed galaxies which should enter the fainter magnitude bin over the mean number of detected ones. It represents the inverse of the fraction of galaxies undetected in each bin. If $`n_i`$ is the number of galaxies detected in the i-th bin of the original catalogue and $`m_{i+1}`$ is the number of sources in the i+1-th bin of the simulated catalogue, the correction factor $`c_{i+1}`$ corresponding to the ratio between the expected number of galaxies ($`n_i`$) and the number of galaxies recovered ( $`m_{i+1}`$). The “true” number of galaxies in the i+1-th bin of the original catalogue is then $`N_{i+1}`$= $`n_{i+1}\overline{c}_{i+1}`$, where $`\overline{c}`$ is the mean over different simulations for the same frame.
When dimming fluxes by a factor F=$`10^{0.4y}`$, magnitudes are $`y`$ mag fainter, but also noise is lowered: if $`\sigma `$ is the sky RMS on original images, the dimmed frames have a RMS of $`\sigma /F`$. We then added a frame of pure Poissonian noise, with a sky RMS
$$\sigma _{noise}^2=\sigma ^2(\sigma /F)^2=\sigma ^2(1\frac{1}{F^2}).$$
The final images have, by construction, the correct RMS. By choosing $`y`$=0.5 mag, the artificial noise added is $`\sigma _{noise}=0.6\sigma `$, i.e. about one third of the final image is due to pure Poissonian noise. As showed in Figures 5-6 and in Table 4, this choice allowed us to correct up to V<sub>606</sub>=29, with brighter limits in the other bands.
We tried also $`y`$=1 mag, but the estimated incompleteness was catastrophic, corresponding to 96$`\%`$. This result seemed caused by the high simulated noise ( $`\sigma _{noise}`$=0.84 $`\sigma `$, i.e. almost one half of noise is artificial), which caused also brighter bins to be incomplete.
The completeness correction allowed the computation of differential number counts up to fainter magnitudes: we thus determined the slopes of the number counts relation. Our best fit gives $`\gamma _U=0.47\pm 0.05`$, $`\gamma _B0.35\pm 0.02`$, $`\gamma _V0.28\pm 0.01`$ and $`\gamma _I0.28\pm 0.01`$ (see Volonteri et al. 2000, for a detailed discussion).
## 6 Galaxies or HII Regions?
Colley et al. (1996) suggested that galaxies in the HDF-N may suffer from a wrong selection. High redshift galaxies on optical images have a clumpy appearance: first the redshift moves the ultraviolet rest-frame light into the optical, so galaxies are observed in UV rest-frame where star-forming regions are more prominent, second the fraction of irregular galaxies is higher than locally (van den Bergh et al. 1996, Abraham et al. 1996) and a large number of galaxies display asymmetry and multiple structure.
Colley et al. (1996) stressed that compact high-redshift objects may appear more prominently than diffuse objects if their angular size is smaller than the point-spread function (PSF). The cosmological dimming in surface brightness for these sources is less significant, leading to an enhancement of compact sources over diffuse, resolved objects. Since HST images have an excellent seeing, these clumps are not smoothed, so they may confuse detection algorithms (SExtractor, DAOFIND, FOCAS, etc), and be counted as several distinct faint sources.
As suggested by Colley et al. (1996) in this case strong correlations between sources on scales $`<10`$ kpc should be found. A good test is therefore the two-point angular correlation function.
In order to check our sample against this effect, we computed this function on small angular scales: HII regions physical sizes are about 0.5 kpc, so a wrong selection of the catalogue should bring a positive peak in the two-point angular correlation function at $`0.251`$ arcsec. This scale corresponds to sizes less than 10 kpc for a wide range in redshift ($`0.8<z<3.5`$).
We computed the correlation function by comparing the number of data pairs at given angular separation to the number of data-random simulated pairs at the same separation (Davis & Peebles, 1983). At fixed $`\theta `$ the sum of data pairs is
$$DD(\theta )=\underset{i}{}\underset{j}{}\delta _i\delta _j$$
where $`\delta _i`$ e $`\delta _j`$ are delta-functions on i-th and j-th galaxies positions. The sum over $`i`$ is over all sources in the sample, and the sum over $`j`$ includes only objects within a distance $`\theta `$ from particle $`i`$.
We then created a random sample and computed the cross count sum
$$DR(\theta )=\underset{i}{}\underset{j}{}\delta _i\delta _j^R$$
where $`\delta _i`$ is as before, and $`\delta _j^R`$ is the delta-function for positions of objects in the random sample within a distance $`\theta `$ from particle $`i`$.
The resulting two-point angular correlation function is given by:
$$w(\theta )=\frac{n_R}{n_D}\frac{DD(\theta )}{DR(\theta )}1$$
where $`\frac{n_R}{n_D}`$ is the ratio of the mean density of random and data samples respectively.
We analyzed our sample in I<sub>814</sub> and B<sub>450</sub> bands. The I<sub>814</sub>-band catalogue should suffer less from the effects described above, being selected in the reddest filter, the vice versa is true for the B<sub>450</sub>-band catalogue.
The two-point angular correlation function is compatible with zero in each band if computed on the whole catalogue. Following Colley et al. (1996) we selected “small galaxies” as defined by $`𝒟=\sqrt{ab}<0.2`$ arcsec, where $`a`$, $`b`$ are the intensity-weighted second moments. At $`\theta =0.8`$ arcsec it is evident a peak in the function, $`w(\theta )=0.67\pm 0.59`$ (Figure 7). This feature may not be significant due to large errors.
However about 20-30$`\%`$ of sources in the B<sub>450</sub>-band catalogue have separation$`<1`$ arcsec. We therefore analyzed these sources, by cross-correlating I<sub>814</sub> and B<sub>450</sub> catalogues. In the B<sub>450</sub>-band catalogue we selected pairs with a separation $`<1`$ arcsec which were not included in the I<sub>814</sub>-band catalogue. These objects were single sources splitted in the B<sub>450</sub>-band (with B$`{}_{450}{}^{}`$ 27-29), corresponding to a single detection in the I<sub>814</sub>-band. We then used SExtractor on the B<sub>450</sub> frame, after choosing a higher *deblend-mincont*=0.1. 87 sources, with $`21<B_{450}<26`$, corresponding to about 7$`\%`$ of the whole sample, were then considered as single galaxies.
This simple analysis shows that a remarkable fraction of sources in the HDF-S has a neighbour at small angular distance. It should be pointed out that the HDF-S is a pencil-beam survey, and projection effects may be sizeable.
## 7 Optical Colours
The estimate of galaxy colours could be biased towards the selection band and affected by an aperture effect. In order to test for this, we first determined colours simply by using SExtractor in the so-called *double image mode* on the I<sub>814</sub>-selected catalogue, we measured B<sub>450</sub> magnitudes within the 1$`\sigma `$ isophotal area determined in the I<sub>814</sub> band (B<sub>col</sub>). The (B-I)<sub>I</sub> colour is then the difference between B<sub>col</sub> and the 1$`\sigma `$ isophotal magnitude in the I<sub>814</sub> band. We then measured the (B-I)<sub>B</sub> colour for the B-selected catalogue using the same method, but with the B<sub>450</sub> band 1$`\sigma `$ isophotal area as a reference for both the I<sub>814</sub> (I<sub>col</sub>) and B<sub>450</sub> magnitudes.
For every source the B-I colours measured on the basis of the B<sub>450</sub> or I<sub>814</sub> area are different. In Figure 8 we show the residuals (B-I)<sub>I</sub>-(B-I)<sub>B</sub> vs I<sub>814</sub>. (B-I)<sub>B</sub> is slightly redder than (B-I)<sub>I</sub>), probably because the B<sub>450</sub>-band isophote underestimates the I<sub>814</sub> band flux (broadly speaking I<sub>814</sub> isophotal areas are larger than B<sub>450</sub> isophotal areas).
In order to measure colours unbiased by the selection band we chose a different approach, that is we created a *meta-image* UBVI by summing all four frames, after normalizing each to have the same rms sky noise. We then run SExtractor in the *double image mode*: detection and isophote boundaries were measured on the combined image, while isophotal magnitudes were measured on U<sub>300</sub>, B<sub>450</sub>, V<sub>606</sub>, I<sub>814</sub> images individually. Using this procedure (Moustakas et al. 1997) both object detection and isophote determination are based on the summed image, and isophotes are not biased towards any of the bands.
We finally cross-correlated the catalogue obtained from the combined image with the I<sub>814</sub> sample selected according to our criteria (see *Image Analysis*). We assigned a lower limit in magnitude to sources undetected in any of other bands (that is missing in our final U<sub>300</sub>, B<sub>450</sub>, V<sub>606</sub> catalogues). The limiting magnitude is the 5$`\sigma `$ isophotal magnitude within the isophote measured in the combined image.
## 8 Structural Parameters
For a subsample (I<sub>814</sub>$`<26`$) we computed size, surface brightness, light concentration index, asymmetry index.
The measure of the size of distant galaxies needs the use of a metric radius, that is a size independent on the redshift and light profile. We estimated a metric radius based on the Petrosian function (Petrosian, 1976), defined (Kron, 1995) as
$$\eta (\theta )=\frac{1}{2}\frac{d\mathrm{ln}l(\theta )}{d\mathrm{ln}(\theta )}$$
where $`l(\theta )`$ is the light growth curve. This function has the property:
$$\eta (\theta )=\frac{I(\theta )}{<I>_\theta }$$
where $`I(\theta )`$ is the surface brightness at the radius $`\theta `$ and $`<I>_\theta `$ is the mean surface brightness within $`\theta `$. We defined for each galaxy the angular size as the value of $`\theta _{0.5}`$ such that $`\eta (\theta )=`$0.5 (Bershady et al., 1998, Saracco et al, 1999). In order to determine the function $`\eta (\theta )`$ we obtained the intensity profile, through multi-aperture photometry at equispaced (0.04 pixel) diameters, and subsequently interpolated it by a spline fit. There are sources whose $`\eta (\theta )`$ is always larger than 0.5 (considering as a border the aperture $`ith`$ such that mag$`{}_{i+1}{}^{}>`$ mag<sub>i</sub>); in these cases we defined $`\theta _{0.5}=1`$, and also all quantities linked to $`\theta _{0.5}`$ were arbitrarily set $`1`$. We also measured the effective radius (half-light radius, $`r_{eff}`$); for I<sub>814</sub>$`<26`$ the relation $`\theta `$ vs $`r_{eff}`$ is well fitted by $`\theta _{0.5}=1.2r_{eff}`$. After determining $`\theta _{0.5}`$, we computed for each galaxy the mean surface brightness within $`\theta _{0.5}`$.
Abraham et al. (1994) and Abraham et al. (1996) showed that two indexes, namely an asymmetry index ($`A`$) and a central concentration index ($`C_{abr}`$), are very useful in order to estimate a quantitative galaxy morphology. The former is determined by rotating the galaxy by 180 and subtracting the resulting image from the original one. The asymmetry index is given by the sum of absolute values of the pixels in the residual image, normalized by the sum of the absolute value of the pixels in the original image and corrected for the intrinsic asymmetry of the background. The concentration index is given by the ratio of fluxes in two isophotes, based on the analysis of light profiles. The measure of these indexes is independent of colour, though they correlate well with optical colours.
We therefore computed both the asymmetry and the concentration indexes following Abraham et al. (1994) and Abraham et al. (1996), but also computed a different light concentration parameters.
That is we computed
$$C_\eta =\frac{F(<\theta _{0.5})}{F(<1.5\theta _{0.5})}$$
(Saracco et al., 1999), i.e. the ratio between the flux within radius $`\theta _{0.5}`$ and the flux within $`1.5\theta _{0.5}`$.
As discussed by Saracco et al. (1999), $`C_\eta `$ is independent from the redshift of the source, since it is related to a metric size and is independent of the asymptotic profile. Brinchmann et al. (1998) on the contrary pointed out that the central concentration $`C_{abr}`$ defined by Abraham et al. (1994, 1996) is redshift-dependent. $`C_{abr}`$ has been computed in order to classify our galaxies and compare our results with HDF-North.
We tested the various parameters against apparent magnitude and checked any correlation with colours (Conselice et al. 2000). We selected two subsamples, composed respectively by galaxies with both B<sub>450</sub>-V<sub>606</sub> and V<sub>606</sub>-I<sub>814</sub> redder or bluer than a local elliptical or a local irregular galaxy. We used these subsamples as tracers of colours.
The asymmetry index (Figure 9) seems not to be biased: the faintest sources are on the whole more symmetric than brighter ones, but the presence of asymmetric objects also in the last bins suggest this feature to be linked to the nature of these galaxies. The trend towards high symmetry may be due to the influence of noise, which makes the profile smoother.
The central concentration index $`C_{abr}`$ defined by Abraham et al. (1994, 1996) seems to be biased against compact sources at faint magnitudes (I$`{}_{814}{}^{}>24.5`$), while $`C_\eta `$ does not correlate with apparent magnitude (Figures 10-11)
## 9 Description of the Catalogue
We presented a catalogue of galaxies in the HDF-S, created using the public version 1 images of the WFPC2 data. We created a catalogue for each pass-band (I<sub>814</sub>, V<sub>606</sub>, B<sub>450</sub><sub>450</sub>, U<sub>300</sub>), which is available upon request. In V<sub>606</sub>, B<sub>450</sub>, U<sub>300</sub> for each galaxy the catalogue reads:
* isophotal flux and magnitude, with their errors (SExtractor)
* isophotal corrected flux and magnitude, with their errors (SExtractor)
* fixed circular aperture (0.6 arcsec radius) flux and magnitude, with their errors (SExtractor)
* Kron’s technique flux and magnitude, with their errors (SExtractor)
* SExtractor best choice for photometry (Kron’s or isophotal corrected)
* isophotal area (pix<sup>2</sup>)(SExtractor)
* coordinates in pixels of the barycentre, related to version1 frames (SExtractor)
* right ascension and declination (SExtractor)
* second order moments of the light distribution (SExtractor)
* semi-major and semi-minor axis lenghts in pixels (SExtractor)
* flags produced by detection and measurement processes (SExtractor)
* stellar classification (SExtractor)
* signal-to-noise ratio (Section 4)
* pseudo-total magnitude and error (Section 3)
For the I<sub>814</sub>-selected catalogue besides the above entries, we estimated also the colours (Section 7) and for I$`{}_{814}{}^{}<26`$, the petrosian radius (arcsec), the mean surface brightness within the petrosian radius (mag arcsec<sup>-2</sup>), light concentration indexes, that is $`C_\eta `$ and $`C_{abr}`$, and the asymmetry index as computed by Abraham software (Section 8). In Table 5 a small part of the catalogue is shown.
## Appendix A The Catalogue
Example of a part of the catalogue: we show only some photometric entries for sake of concision. The real catalogue reads also the isophotal corrected flux and magnitude, with their errors, a fixed circular aperture (0.6 arcsec radius) flux and magnitude, with their errors, Kron’s technique flux and magnitude, with their errors, and SExtractor best choice for photometry (Kron’s or isophotal corrected). The full catalogue is available at http://cdsweb.u-strsbg.fr/Abstract.html and via
http://www.merate.mi.astro.it/ saracco/science.html.
###### Acknowledgements.
Thanks to M. Bolzonella for useful discussions and support and R. Abraham for making available the software for computing asymmetry and concentration index. We would like to thank the referee, G. Paturel, for his helpful comments and suggestion. MV acknowledges financial support by Fondazione Cariplo.
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# Spectroscopic confirmation of clusters from the ESO imaging survey Based on observations collected at the European Southern Observatory (La Silla, Chile), Proposal ID: 62.O-0601
## 1 Introduction
Clusters of galaxies are the largest virialized structures observed in the Universe. Since they arise from exceptionally high peaks of the primordial fluctuation density field, their properties are highly sensitive to the nature of such cosmic fluctuations. Therefore, the mass function of both local (e.g. White et al. whi93 (1993); Girardi et al. gir98 (1998)) and distant clusters (e.g. Oukbir & Blanchard ouk92 (1992); Carlberg et al. car97 (1997); Eke et al. eke98 (1998); Borgani et al. bor99 (1999)) is a powerful tool to constrain cosmological models for the formation and evolution of cosmic structures. Moreover, clusters are useful laboratories for testing models of galaxy evolution. While early-type galaxies only show evidence for passive evolution (e.g. Stanford et al. sta98 (1998)), the fraction of blue galaxies increases significantly with redshift (Butcher & Oemler but78 (1978)), at least up to $`z0.5`$, and the fraction of S0’s decreases (Dressler et al. dre99 (1999)). It is therefore essential to have reliable cluster catalogues over the largest possible redshift range.
Most distant clusters, at $`z0.5`$, have so far been identified through optical follow-ups of X-ray selected clusters (see, e.g. Gioia et al. gio90 (1990) and Rosati et al. ros00 (2000) for a recent review), or by looking at the environment of high-redshift radio galaxies (e.g. Smail & Dickinson sma95 (1995); Deltorn et al. del97 (1997)).
In the optical, clusters at $`z0.5`$ and beyond started to be classified in the 80’s (Gunn et al. gun86 (1986)). In the 90’s a large catalogue of objectively selected distant clusters, identified in the optical, became available (Postman et al. pos96 (1996)). These last clusters are identified using a matched-filter algorithm using both positional and photometric data. In brief, this algorithm filters a galaxy catalogue to remove fluctuations in the projected distribution of galaxies that are not likely to be galaxy clusters. For this purpose, the filter is built around parametrizations of the spatial distribution and luminosity function of cluster galaxies. This algorithm also provides an estimate of the redshift for each candidate cluster (hereafter we refer to the matched-filter estimated redshift as $`z_{mf}`$). Currently, $`30`$ PDCS clusters have been confirmed spectroscopically, most of them at $`z<0.5`$ (Holden et al. 1999a , 1999b ; Oke et al. oke98 (1998)).
Recently, Olsen et al. (1999a , 1999b ) and Scodeggio et al. (sco99 (1999)) have presented a catalogue of 302 cluster candidates from the $`I`$-band images of the ESO Imaging Survey (EIS, see Renzini & da Costa ren97 (1997)). Clusters are identified in two dimensions (hereafter, 2-d) using the matched filter algorithm of Postman et al. (pos96 (1996); see Olsen et al. 1999a ). The estimated redshifts for EIS clusters span the range $`0.2z_{mf}1.3`$, with a median redshift $`z_{mf}=0.5`$.
Several EIS cluster candidates have been confirmed so far, most at $`z<0.5`$, either from the existence of the red sequence of cluster ellipticals/S0’s in colour-magnitude diagrams (Olsen et al. 1999b ), or from a combination of photometric and spectroscopic data (da Costa et al. dac99 (1999)).
The EIS cluster catalogue is the largest optically selected cluster sample currently available in the Southern Hemisphere to this depth. This catalogue constitutes an obvious reference for follow-up observations at the ESO VLT aimed at determining the structure and dynamics of distant clusters, as well as the spectroscopic properties of their member galaxies. Unfortunately, little is currently known on the performance of the matched filter algorithm in detecting real clusters at $`z0.5`$. As we already pointed out, most confirmed PDCS and EIS clusters have redshifts $`z<0.5`$. Therefore, to point blindly at EIS cluster candidates would make for an inefficient use of VLT time, because we expect several of these candidate clusters not to be real, in particular at $`z_{mf}0.5`$.
The aim of our investigation is twofold: we want to confirm as many EIS clusters as possible, in order to build a reliable sample of distant clusters with well determined redshift, and, at the same time, evaluate the performance of the matched filter algorithm in the detection of high-redshift clusters. In order to achieve this purpose, we use two independent methods: (1) multi-object spectroscopic observations of EIS cluster candidates in the redshift range $`0.5z_{mf}0.7`$, and (2) the detection of the colour-magnitude sequences traced by early-type galaxies through multi-colour optical and near-IR photometry of the most distant EIS cluster candidates (Scodeggio et al., in preparation).
In this paper we report the first results of the spectroscopic investigations of 6 EIS clusters. We are able to confirm the existence of significant concentrations in redshift space in correspondence of four of the six EIS fields targeted. For two of these confirmed clusters, the spectroscopic mean redshift agrees with the matched-filter estimate to within $`\mathrm{\Delta }z=\pm 0.1`$.
In Sect. 2 we describe our spectroscopic observations, data reduction, and give the new galaxy redshifts. In Sect. 3 we analyse the data, and define sets of galaxies in redshift space. We also discuss the concordance of the mean redshifts of these sets with the matched-filter estimates of the cluster mean redshifts. We then make a likelihood analysis of the reality of the galaxy sets, and flag four of them as reliable at $`>95`$ % confidence level (Sect. 4). Finally, we discuss our results and give our conclusions in Sect. 5.
We use H$`{}_{0}{}^{}=`$ h<sub>75</sub> 75 km s<sup>-1</sup> Mpc<sup>-1</sup>, $`\mathrm{\Omega }_0=0.2`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ throughout this paper, unless otherwise stated.
## 2 Observations and data reduction
We select our targets among the candidate EIS clusters in patches C and D, with estimated redshift $`0.5z_{mf}0.7`$ (Scodeggio et al. sco99 (1999)). We do not apply any additional criterion for the selection of our cluster candidates. The size of our sample is one tenth of the total of 36$`+`$28 EIS clusters in patches C and D within the above-mentioned redshift range.
Our targets are listed in Table 1. In Column (1) we list the cluster candidate identification name, in Column (2) and (3) the right ascension and declination (J2000), in Column (4) the cluster richness (see Olsen et al. 1999a ), and in Column (5) the matched-filter redshift estimate, $`z_{mf}`$ (see Scodeggio et al. sco99 (1999)). In Column (6) we list the number of galaxies targeted for multi-slit spectroscopy in each cluster field, and in Column (7) the number of successful redshift estimates.
The observations were carried out at the 3.6 m ESO telescope at La Silla, Chile, during two nights in February 1999. The weather conditions were good, with seeing slightly above 1”, in partial moonlight. With a total useful observing time of 16 hours over the two nights, we were able to obtain $`3\times 45`$ min exposures for each cluster.
We observed with EFOSC2 in Multi-Object Spectroscopy (MOS) mode. EFOSC2 was equipped with a Loral CCD of 2048$`\times `$2048 with 15 $`\mu `$m pixels, allowing for an unvignetted field-of-view of 3.8’$`\times `$5.5’. We used Grism # 1, giving a spectral range 3185–10940 Å, and a dispersion of 6.3 Å/pixel. On the MOS masks our slits were 1.2” wide.
We obtained spectra for 102 objects in the magnitude range $`17.0m_I21.3`$, where $`m_I`$ is the apparent magnitude in the $`I_c`$ band (Nonino et al. non99 (1999)). In Fig. 16 we show $`I_c`$-band images of the six EIS candidate clusters. Small circles mark galaxies with redshift, large circles mark galaxies belonging to significant overdensities in redshift space (see Sect. 4).
We reduce the data with standard IRAF<sup>1</sup><sup>1</sup>1IRAF is distributed by the National Optical Astronomy Observatories, which is operated by AURA Inc. under contract with NSF packages. We determine redshifts using the task XCSAO that implements the cross-correlation technique of Tonry & Davis (ton79 (1979)). We use several real and synthetic templates for the cross-correlation. We use emission lines, where present in the spectrum of the object, to determine the redshift with the task EMSAO. We examine visually all spectra, by overplotting the positions of the major spectral features redshifted at the redshift(s) determined by the automatic techniques described above. We employ particular care in flagging those features that could be contaminated by night-sky lines.
In total, we determine 67 galaxy redshifts, from a minimum of $`z=0.09`$ to a maximum of $`z=0.79`$, with an average $`\overline{z}=0.380`$. One of our objects turns out to be a QSO at $`z=3.2`$. We do not consider this object in our analysis. An internal estimate of the typical redshift uncertainty is $`\delta z0.001`$. The success-rate is magnitude-dependent, as can be seen in Fig. 7: it is 85 % for $`m_I19.5`$ and decreases to 57 % for fainter galaxies.
We list in Table 2 the galaxies with measured redshift. In Column (1) we list the name of the EIS candidate cluster field, in Column (2) a galaxy identification number, in Column (3) and (4) the (J2000) right ascension and the declination of the galaxy, in Column (5) the $`I_c`$ magnitude, in Column (6) the redshift, and in Column (7) the galaxy set to which the galaxy is assigned. These galaxy sets are defined in Sect. 3 and listed in Table 3.
## 3 The definition of the galaxy systems in redshift space
Since the six cluster candidates are drawn from the EIS cluster catalogue, they obviously correspond to significant density enhancements in projection. Here we search for systems of galaxy redshifts that could be associated to the 2-d over-densities. In this way we assign a spectroscopic redshift to 4 of our cluster candidates. We also evaluate the probability that these systems correspond to a genuine three-dimensional cluster.
Olsen et al. (1999a ) search for clusters in projection assuming a 2-d radial density profile with a cutoff radius $`r_{co}=1.33`$ h$`{}_{}{}^{1}{}_{75}{}^{}`$ Mpc. This size is well fitted to the EFOSC2 field-of-view, corresponding to $`1.9\times 1.3`$ h$`{}_{}{}^{2}{}_{75}{}^{}`$ Mpc<sup>2</sup>, at $`z0.6`$ (the average estimated redshift of our candidate clusters). Therefore, we search for cluster members in redshift-space within the whole EFOSC2 field.
Several refined algorithms for the definition of systems of galaxies in redshift space can be found in the literature (e.g. Katgert et al. kat96 (1996); Pisani pis93 (1993)). However, with only a dozen galaxy redshifts per field, these sophisticated algorithms can not be applied. We choose to identify galaxy systems using a physical criterion based on the well established properties of nearby clusters of galaxies.
Within each EFOSC2 field, we identify any set of two or more galaxies contained within a given redshift range, $`\mathrm{\Delta }z`$. In order to define $`\mathrm{\Delta }z`$, we note that Abell-like clusters of galaxies have mean velocity dispersions $`\sigma _v750`$ km s<sup>-1</sup> (Girardi et al. gir93 (1993)). Since the line-of-sight velocity distributions of clusters are approximately gaussian in shape (Girardi et al. gir93 (1993)), $`>99`$ % of the cluster members have a velocity within $`\pm 3\sigma _v`$. As a consequence, galaxies in a given cluster should be located within a redshift range $`\mathrm{\Delta }z0.015\times (1+z)`$, taking into account the cosmological factor (Danese et al. dan80 (1980)).
We list in Table 3 the sets of galaxies we identify in redshift space. In Column (1) we list the EIS cluster field identification, in Column (2) the galaxy set identification, in Column (3) the number of galaxies in the set, in Column (4) the median redshift of the set, in Column (5) the total redshift range covered by the galaxies within the galaxy set. In Column (6) we list the probability of the galaxy set to correspond to a significant overdensity in redshift space, as estimated from resamplings of the Canada-France Redshift Survey data-base (CFRS; Lilly et al. lil95 (1995); Le Fèvre et al. lef95 (1995); Hammer et al. ham95 (1995); Crampton et al. cra95 (1995)) – see Sect. 4. The real systems (probability $`0.95`$) are flagged in Column (7).
In total we identify 19 galaxy sets along the line-of-sight of six EIS candidate clusters. It is interesting to detail the comparison of the estimated mean redshifts, $`z_{mf}`$’s, of these clusters (see Table 1), to the spectroscopic redshifts of the 19 sets (see Table 3). In this comparison, we take into account the uncertainties in the mean redshifts of the galaxy sets, and note that the matched-filter redshift estimates are at most accurate to within $`\mathrm{\Delta }z=\pm 0.05`$.
In the case of the candidate clusters EIS0533-2353, EIS0950-2154, and EIS0951-2047, we do not find any galaxy set close to the estimated cluster redshifts. According to the matched-filter algorithm, these three clusters are located at a higher redshift than any of the galaxy sets found in their fields. In each of the fields of EIS0955-2113 and of EIS0956-2009, there is one set of galaxies with mean redshift close to the estimated redshift ($`\overline{z}=0.673`$ vs. $`z_{mf}=0.6`$, and $`\overline{z}=0.445`$ vs. $`z_{mf}=0.5`$, respectively). Finally, in the field of EIS0540-2418, there are two sets of galaxies with mean redshifts close to the cluster estimated redshift, $`z_{mf}=0.6`$.
The detection of galaxy concentrations close to the estimated redshifts of three EIS clusters supports the reliability of the matched-filter redshift estimates.
As far as the failed detections are concerned, EIS0533-2353 may have escaped detection because of its high (estimated) redshift, $`z_{mf}=0.7`$ – we may simply have not been observing deep enough. Moreover, the line-of-sight to a single cluster can intercept several different galaxy sets. Katgert et al. (kat96 (1996)) estimate that 10 % of nearby Abell clusters result from the superposition of two almost equally rich systems (and this fraction is probably higher for more distant clusters). The nearest of these systems has the highest chance of being detected. In this context, we note that the low-$`z`$ set of galaxies, 3a, detected in the field of EIS0950-2154 is somewhat off-centered with respect to the nominal EIS cluster center (see Fig. 3). The same is not true for the sets with $`zz_{mf}`$ in the fields of EIS0540-2418, EIS0955-2113, and EIS0956-2009. This fact suggests that the set 3a does not correspond to the EIS cluster, but is a foreground group.
Da Costa et al. (dac00 (2000)) suggest that all the six EIS candidate clusters could be real clusters at redshift $`>0.5`$. Da Costa et al. base their suggestion on the detection of red-sequences of early-type galaxies in the colour-magnitude diagrams. They note that under-sampling of the redshift distribution of galaxies in the cluster fields may explain the lack of spectroscopic detections of some clusters.
## 4 The reality of the galaxy sets
In this Sect. we assign likelihoods to the 19 sets of galaxies which have been identified in the previous Sect. If $`s(m)`$ is the incompleteness of our redshift sample in the magnitude interval $`[m_1,m_2]`$ (see Fig. 8), then the redshift distribution reads
$$N(z)dz=\frac{dV(z)}{dz}dz_{m_1}^{m_2}f[L(m,z)]s(m)𝑑m,$$
(1)
where $`V(z)`$ is the volume element at the redshift $`z`$ and $`f[L(m,z)]`$ is the luminosity function (LF hereafter). We take the LF in the $`I_4`$ band given by Postman et al. (pos96 (1996)) and convert it to our $`I_c`$ band, using the Postman et al. transformation procedure. Our reference LF has therefore a Schechter (sch76 (1976)) form with parameters $`\alpha =1.1`$ and $`M^{}=22.15+5\mathrm{log}h_{75}`$. We also assume negligible evolution of the LF out to $`z0.8`$ (see Lilly et al. lil95 (1995) and Lin et al. lin99 (1999)). Our LF then depends on $`z`$ only through the evolution of the stellar populations and the K-correction, that we take from Poggianti (pog97 (1997)).
In Fig. 9, we plot $`N(z)`$, computed according to Eq. 1, along with the observed $`z`$-distribution. Based on the estimated $`N(z)`$, we find that 16 out of 19 sets correspond to significant overdensities in redshift space, with a probability $`95`$ %. Since Eq. 1 provides $`N(z)`$ for a uniform galaxy distribution, these probabilities do not include the effect of large scale clustering (see, e.g., Zaritsky et al. zar97 (1997)). In order to account for the large scale clustering, and following the approach by Holden et al. (1999a ), we modulate $`N(z)`$ with the redshift distribution derived from the CFRS assuming our selection function, $`s(m)`$.
We first convert the CFRS $`I_{AB}`$ magnitudes to the $`I_c`$ system (see Lilly et al. lil95 (1995)). We then extract 50000 galaxies from the CFRS, with a bootstrap sampling, adopting the magnitude distribution of our sample (see Fig. 7). A Kolmogorov-Smirnov test shows that the bootstrapped CFRS and our data-set have similar redshift distributions. This is expected since magnitude selection is the main process leading to the inclusion of galaxies in the sample.
We extract random subsamples of 11 galaxies from the bootstrapped CFRS reference sample (11 is the mean number of galaxies with redshift in our EIS cluster fields). Using the same procedure described in Sect. 3, we identify sets of galaxies within these subsamples, and compute their probabilities relative to the uniform redshift distribution $`N(z)`$ given by Eq. 1. In this way we construct a distribution of probabilities to detect a real system within a galaxy sample which includes large-scale clustering, but not galaxy clusters. We finally obtain the likelihoods of the 19 observed sets, by comparing their original $`N(z)`$-based probabilities to the distribution of probabilities for the random sets. These 19 likelihoods are listed in Table 3.
We find that four of our 19 sets have a likelihood $`>95`$ %; all of them have at least three galaxy members. The four sets are flagged in the last column of Table 3. We refer to these four sets as the ’real systems’ hereafter. As expected, many of the sets with a significant overdensity with respect to the uniform redshift distribution are no longer significant when compared to a redshift distribution which includes the large scale clustering.
Our results are robust against modifications of the adopted LF (we change $`\alpha `$ by $`\pm 0.2`$, and $`M^{}`$ by $`\pm 0.5`$ mag), and of the galaxy-type for which the evolutionary- and K-corrections are computed. Furthermore, we verify that varying cosmological parameters (h<sub>75</sub>, $`\mathrm{\Omega }_0`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }`$) within conservative ranges, only induces marginal changes in the system likelihoods. Finally, we also checked that narrowing the $`\mathrm{\Delta }z`$ range used to define the redshift-sets, from 0.015 to 0.010, hardly modifies the membership and likelihoods of the sets.
## 5 Discussion and conclusions
We obtain 67 new redshifts for galaxies in six EIS candidate cluster fields. Based on these data, we establish the existence of real systems in redshift space in the direction of four of these candidate clusters. The reality of the systems is established at $`>95`$ % confidence level, and in two cases, at $`>99`$ %. The redshift overdensities, coupled with the 2-d overdensities detected by the use of the matched-filter algorithm, strongly supports the reality of four of the six examined EIS clusters. These 4 clusters add to the other two spectroscopically confirmed EIS clusters (da Costa et al. dac99 (1999)).
Two of the four $`z`$-systems have a median redshift in good agreement with the matched-filter estimate for the EIS cluster redshift ($``$ median$`(z)z_{mf}<0.1`$). The other two have significantly lower redshifts (median$`(z)=0.129`$, 0.236 vs. $`z_{mf}=0.5`$).
Taken at face value, these results suggest that, in several cases, the matched-filter algorithm over-estimates the mean cluster redshift by a large amount. However, it is quite possible that in some cases we have not detected the EIS cluster, but a foreground galaxy system projected along the same line-of-sight of the cluster. Similarly, it is difficult to conclude about the reality of the EIS clusters where we do not detect any real redshift system.
In particular, we note that in the field of EIS0540-2418 we have a marginal detection (92 % probability) of the galaxy set 2d at median$`(z)=0.698`$ (see Table 3), in fair agreement with the matched-filter estimate of the cluster redshift, $`z_{mf}=0.6`$. We also note that the cluster EIS0533-2353 has $`z_{mf}=0.7`$, larger than any of our galaxy sets. This suggests that it could have escaped detection because our observations were not deep enough. In fact, da Costa et al. (dac00 (2000)) suggest that all our six EIS cluster candidates could be real, based on the analysis of the colour-magnitude diagrams for galaxies in the cluster fields.
We conclude that our spectroscopic confirmation rate must be considered as a lower limit. If at least one third of the EIS clusters in the redshift range sampled by our observations are real, there are more than 25 EIS clusters with $`z_{mf}`$ in the range 0.5–0.7. This sample is large enough for the derivation of the properties of clusters at intermediate to high redshifts.
Optical selection of clusters of galaxies at high redshifts is a necessary complementary approach to X-ray selection. While X-ray selection tends to detect only rich Abell-like clusters, optically selected cluster samples contain a large number of poor clusters. In fact, the space density of PDCS clusters is five times higher than that of rich Abell clusters (Holden et al. 1999a ), and very few PDCS clusters are X-ray bright (Holden et al. hol97 (1997)). Consistently, the velocity dispersions of our two systems with $`5`$ galaxy redshifts (system 4c, at $`\overline{z}=0.236`$ and system 6c, at $`\overline{z}=0.445`$, see Table 3) are $`600`$ km s<sup>-1</sup>, typical of low-richness clusters ($`R1`$, see Girardi et al. gir93 (1993)).
With the current and near-future ground-based facilities for wide-field optical and near-infrared imaging, we can expect a rapid increase in the samples of optically selected clusters. Currently, our spectroscopic sample only comprises $`10`$ % of all the clusters in the two patches C and D, and in the (estimated) redshift range 0.5–0.7. We plan to extend our sample in forthcoming observing runs. Confirmed EIS clusters at high redshift will be the natural targets of VLT observations aimed at determining their dynamical properties.
###### Acknowledgements.
We thank Nando Patat for imaging the six EIS fields in preparation for our spectroscopic run.
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# Dense gas in nearby galaxies Based on observations with the Heinrich-Hertz-Telescope (HHT) and the IRAM 30-m telescope. The HHT is operated by the Submillimeter Telescope Observatory on behalf of Steward Observatory and the Max-Planck-Institut für Radioastronomie
## 1 Introduction
Low lying rotational transitions of CO are widely used as tracers of molecular hydrogen and are essential to determine dynamical properties and total molecular masses of galaxies. The widespread use of CO $`J`$ = 1–0 and 2–1 spectroscopy is however not sufficiently complemented by systematic surveys in higher rotational CO transitions to confine the excitation conditions of the dense interstellar medium (ISM). While the $`J`$ = 1 and 2 states of CO are only 5.5 and 17 K above the ground level, the $`J`$ = 3 to 7 states are at 33, 55, 83, 116, and 155 K and trace a component of higher excitation. ‘Critical densities’, at which collisional deexcitation matches spontaneous decay in the optically thin limit, are $``$ 10<sup>5-6</sup>$`\text{cm}^3`$ for CO $`J`$ = 3–2 to 7–6 in contrast to 10<sup>3.5</sup> and 10<sup>4.3</sup>$`\text{cm}^3`$ for the ground rotational CO transitions.
Starburst galaxies are known to contain large amounts of molecular gas that may be heated to $`T_{\mathrm{kin}}`$ $``$ 100 K by young massive stars, cosmic rays or turbulent heating. Therefore highly excited CO transitions, observed at submm-wavelengths, are the appropriate tool to study this interstellar gas component. Among the three nearest ($`D`$ $``$ 3 Mpc) nuclear starburst galaxies, NGC 253, NGC 4945, and M 82 (NGC 3034) M 82 is most readily accessible from telescopes of the northern hemisphere. Containing one of the brightest IRAS point sources beyond the Magellanic Clouds ($`S_{100\mu \mathrm{m}}`$ $``$ 1000 Jy), M 82 has been observed at a variety of wavelengths, ranging from the radio to the $`\gamma `$-ray domain of the electromagnetic spectrum. The starburst in M 82 is likely triggered by a tidal interaction with M 81, causing a high infrared luminosity ($`L_{\mathrm{FIR}}`$ $``$ 4 10<sup>10</sup>L), a high density of supernova remnants, and copious amounts of dense gas with strong OH and $`\mathrm{H}_2\mathrm{O}`$ masers and a large number of molecular high density tracers (for CO maps, see Sutton et al. 1983; Olofsson & Rydbeck 1984; Young & Scoville 1984; Nakai et al. 1986, 1987; Lo et al. 1987; Loiseau et al. 1988, 1990; Phillips & Mampaso 1989; Turner et al. 1991; Tilanus et al. 1991; Sofue et al. 1992; White et al. 1994; Shen & Lo 1995; Kikumoto et al. 1998; Neininger et al. 1998).
So far, few CO 4–3 maps of external galaxies were published (for M 51, M 82, M 83, and NGC6946 see White et al. 1994; Petitpas & Wilson 1998; Nieten et al. 1999). Among these M 82 is the only true starburst galaxy but its CO 4–3 map (White et al. 1994) is confined to the very central region. With respect to higher rotational CO transitions, only a few CO 6–5 spectra were presented from nearby galaxies (Harris et al. 1991; Wild et al. 1992).
We have used the Heinrich-Hertz-Telescope (HHT) on Mt. Graham (Baars & Martin 1996) to map M 82 in the CO $`J`$= 7–6, 4–3, and <sup>13</sup>CO 3–2 transitions. These data are complemented by new $`J`$ = 2–1 and 1–0 spectra taken with the IRAM 30-m telescope.
## 2 Observations
### 2.1 Observations with the Heinrich-Hertz-Telescope
<sup>13</sup>CO 3–2 (331 GHz $``$ 907$`\mu `$m), <sup>12</sup>CO 4–3 (461 GHz $``$ 650$`\mu `$m), and <sup>12</sup>CO 7–6 (807 GHz $``$ 372$`\mu `$m) line emission was observed at the HHT during Feb. 1999 with beamwidths of $``$ 22<sup>′′</sup>, 18<sup>′′</sup>, and 13<sup>′′</sup>, respectively. For the CO 3–2 and 4–3 transitions, SIS receivers were employed; the CO 7–6 line was observed with a Hot Electron Bolometer (HEB) kindly provided by the Center for Astrophysics (Kawamura et al. 1999). The backends consisted of two acousto optical spectrometers, each with 2048 channels (channel spacing $``$ 480 kHz, frequency resolution $``$ 930 kHz) and a total bandwidth $``$ 1 GHz.
Spectra were taken using a wobbling (2 Hz) secondary mirror with a beam throw of $`\pm `$120 to $`\pm `$240<sup>′′</sup> in azimuth. Scans obtained with reference positions on either azimuth were coadded to ensure flat baselines. Receiver temperatures were at the order of 170 K at 331 GHz, 150 K at 461 GHz, and 1000 K at 807 GHz; system temperatures were $``$ 900, 3500, and 8000 K on a $`T_\mathrm{A}^{}`$ scale, respectively.
The receivers were sensitive to both sidebands. Any imbalance in the gains of the lower and upper sideband would thus lead to calibration errors. To account for this, we have observed Orion-KL (at 807 GHz) and IRC+10216 (at 807 and 331 GHz) prior to our M 82 measurements with the same receiver tuning setup. Peak temperatures were 70, 20, and 2.3 K on a $`T_\mathrm{A}^{}`$ scale, respectively (cf. Howe et al. 1993; Groesbeck et al. 1994). At 461 GHz, Orion-KL ($`T_\mathrm{A}^{}`$ $``$ 70 K; cf. Schulz et al. 1995) was mapped but the tuning was later changed by $`\mathrm{\Delta }V`$ = 200 $`\mathrm{km}\mathrm{s}^1`$ for M 82 (see Sect. 3.2).
All results displayed are given on a main beam brightness temperature scale ($`T_{\mathrm{mb}}`$). This is related to $`T_\mathrm{A}^{}`$ via $`T_{\mathrm{mb}}`$ = $`T_\mathrm{A}^{}`$ ($`F_{\mathrm{eff}}`$/$`B_{\mathrm{eff}}`$) (see Downes 1989). Main beam efficiencies, $`B_{\mathrm{eff}}`$, were 0.5, 0.38, and 0.36 at 330, 461, and 806 GHz, as obtained from measurements of Saturn. Forward hemisphere efficiencies are 0.9, 0.75, and 0.70, respectively (D. Muders, priv. comm.). With an rms surface accuracy of $``$ 20$`\mu `$m ($`\lambda `$/18 at 806 GHz), the HHT is quite accurate. This reduces the effect of the source coupling efficiency on the measured source size.
At the center of M 82, CO lineshapes depend sensitively on the position observed, so that in each of our maps the dynamical center could be identified with an accuracy better than 5<sup>′′</sup>. While relative pointing errors should be small when compared to the spacing (10<sup>′′</sup>) of our <sup>13</sup>CO $`J`$ = 3–2 and CO 4–3 maps, relative pointing is less reliable in the case of the CO $`J`$ = 7–6 map; here deviations may reach $`\pm `$5<sup>′′</sup> for a few positions.
### 2.2 Observations with the IRAM 30-m telescope
<sup>12</sup>CO and <sup>13</sup>CO $`J`$ = 1–0 and 2–1 observations were made with SIS receivers of high image sideband rejection ($``$25 db for 1–0 and $``$13 db for 2–1 line data) of the inner 50<sup>′′</sup> $`\times `$ 50<sup>′′</sup> (<sup>12</sup>CO; 5<sup>′′</sup> spacing for the central 20<sup>′′</sup>, otherwise 10<sup>′′</sup>) and 20<sup>′′</sup> $`\times `$ 20<sup>′′</sup> (<sup>13</sup>CO; 5<sup>′′</sup> spacing) of M 82 in June 1999. The measurements were made in a position switching mode with the off-position displaced by 15 in right ascension. The two <sup>12</sup>CO lines (as well as those of <sup>13</sup>CO) were measured simultaneously. Beamwidths at 115 ($`J`$ = 1–0) and 230 GHz ($`J`$ = 2–1) were 21 and 13<sup>′′</sup>; forward hemisphere and beam efficiencies were 0.92 and 0.72 for the $`J`$ = 1–0 and 0.89 and 0.45 for the $`J`$ = 2–1 data, respectively. Calibration was checked by observing IRC+10216. Measured line intensities were $`T_{\mathrm{mb}}`$ = 17 and 43 K for <sup>12</sup>CO $`J`$ = 1–0 and 2–1 (channel spacings were 3.2 and 1.6 $`\mathrm{km}\mathrm{s}^1`$, respectively) and 2.3 and 6.0 K for <sup>13</sup>CO $`J`$ = 1–0 and 2–1 (channel spacings: 2.7 and 1.4 $`\mathrm{km}\mathrm{s}^1`$; cf. Mauersberger et al. 1989).
## 3 Results
### 3.1 Overall morphology of CO submillimeter line emission
Maps of integrated <sup>13</sup>CO $`J`$ = 3–2, <sup>12</sup>CO 4–3, and <sup>12</sup>CO 7–6 intensity are presented in Fig. 1. Detectable emission is strongly confined to the central part of the galaxy. In each map a particularly wide spectral feature, with slightly smaller peak intensity than the most intense lines in the SW and NE, could be identified with the dynamical center (see Fig. 2 which also contains IRAM spectra). Two main peaks of emission are detected, being displaced by almost 10<sup>′′</sup> from the center (this corresponds to a projected distance of 150 pc); the south-western hot spot is most prominent, while evidence for a third peak in the NE (at $``$ (17<sup>′′</sup>,17<sup>′′</sup>)), seen in the CO $`J`$ = 4–3 and 7–6 maps, is not conclusive. Velocity channel maps of the <sup>13</sup>CO $`J`$ = 3–2 and <sup>12</sup>CO $`J`$ = 4–3 line emission are displayed in Fig. 3. These outline the extent of the emission at various velocities, show the dominant rotation pattern with the red-shifted lobe in the north-east and the blue-shifted lobe in the south-west, and indicate the rapid change in radial velocity near the dynamical center (cf. Neininger et al. 1998).
### 3.2 Consistency of CO line temperatures
Calibration at sub-millimeter wavelengths is critical because of rapidly changing weather conditions, high atmospheric opacities, imbalances in the receiver gains between the sidebands, and because of uncertainties in beam and forward hemisphere efficiencies. Calibration uncertainties introduced by these effects may rise up to $`\pm `$30% and a comparison with data published elsewhere is needed. Intensities of our <sup>13</sup>CO $`J`$ = 3–2 spectra from the center, the south-western, and north-eastern lobes are smaller by $``$30% than those given by Tilanus et al. (1991; their Fig. 3). Our spectrum from IRC+10216 is however 30% stronger (on a $`T_{\mathrm{mb}}`$ scale) than that given by Groesbeck et al. (1994), so that our scaling is intermediate between those of Tilanus et al. and Groesbeck et al. CO $`J`$ = 4–3 line intensities (Fig. 2) are $``$30% larger than those given by White et al. (1994). Since their data were obtained with higher angular resolution (JCMT beamwidth: 11<sup>′′</sup>), this difference is significant. A comparison with the three 4–3 spectra from the 10-m CSO shown by Güsten et al. (1993) shows good agreement for the lobe positions. Differences in 4–3 peak temperatures between central and lobe positions are however less pronounced than reported by Güsten et al. (we find peak line temperature ratios of $``$1.3 instead of $``$2.0 between the lobes and the center). Our CO 7–6 line intensities from Orion-KL agree well with those obtained by Howe et al. (1993) with the 10-m CSO antenna.
Previous high angular resolution CO data from the 30-m telescope (compare Fig. 9a of Loiseau et al. 1990 with our Table 1) were affected by uncertainties in the image sideband ratios. Our new data are not significantly affected by this problem, should be reliably calibrated within $`\pm `$10%, and agree well with recent spectra obtained independently in an ‘on-the-fly’ observing mode (Weiß, in preparation).
### 3.3 Line intensity ratios
Claims that the large scale ($``$20<sup>′′</sup>) integrated CO $`J`$ = 2–1/1–0 line intensity ratio were much larger than unity (e.g. Knapp et al. 1980; Olofsson & Rydbeck 1984; Loiseau et al. 1990) can be firmly rejected. The more recently measured ratios of 1.0 (Wild et al. 1992), 1.1 (Mauersberger et al. 1999) and 1.0–1.4 (Table 1) imply that <sup>12</sup>CO line intensities of the three lowest rotational CO transitions must be similar over the central 22<sup>′′</sup>. Our <sup>13</sup>CO $`J`$ = 3–2 line temperatures are a factor of $``$10 smaller than those of the $`J`$ = 3–2 <sup>12</sup>CO transition observed with the same telescope (R. Wielebinski, priv. comm.). The situation with respect to the $`J`$ = 6–5 line remains unclear. 6–5 spectra from the western lobe (Harris et al. 1991) and the center (Wild et al. 1992; their Fig. 2) do not allow a beam convolution to the angular resolution of the lower frequency data and the beam pattern appears to be complex. Our $`J`$ = 7–6 data show peak line temperatures of order 2–3 K. This is smaller than the 4 K measured in the $`J`$ = 4–3 and lower $`J`$ transitions with larger beam sizes. Beyond $`J`$ = 4–3 we thus find clear evidence for a weakening of CO emission with increasing rotational quantum number $`J`$.
## 4 Spatial distributions: Are there differences?
### 4.1 CO
A comparison of the CO $`J`$ = 4–3 data presented in Fig. 1 with those of White et al. (1994) shows a strong discrepancy in the overall spatial distribution: Our CO $`J`$ = 4–3 map contains (at least) two maxima of emission, while the higher angular resolution data of White et al. (1994) have only one peak. Data with only one peak (that are based on spectra with sufficient resolution to separate the lobes) are also presented by Wild et al. (1992; their Fig. 10 and Table 5) for <sup>13</sup>CO and C<sup>18</sup>O $`J`$ =2–1. The original <sup>13</sup>CO $`J`$ = 2–1 spectra displayed by Loiseau et al. (1988; their Fig. 1), however, clearly show a double-lobed distribution. Maps in the low-$`J`$ <sup>12</sup>CO and <sup>13</sup>CO transitions (e.g. Lo et al. 1987; Nakai et al. 1987; Loiseau et al. 1988, 1990; Tilanus et al. 1991; Shen & Lo 1995; Neininger et al. 1998) as well as our $`J`$ = 7–6 (Fig. 3), 2–1 and 1–0 data all show a double-lobed structure. We thus conclude that, in spite of previous evidence to the contrary, the overall spatial distribution of emission from highly excited CO shows two main centers of emission.
Are the two main hotspots observed in the $`J`$ = 1–0 and 2–1 transitions identical with those observed in higher excited CO lines? The position angles (east of north) of the lines connecting the hotspots are slightly smaller in our submillimeter data (Fig. 1) than in interferometric maps (cf. Lo et al. 1987; Shen & Lo 1995; Kikumoto et al. 1987; Neininger et al. 1998). This is likely an error in our data caused by measurements at varying hour angles with pointing offsets along the azimuth and elevation axes.
More significant is a difference in angular separation: While the two main peaks of line emission observed by us are separated by 15<sup>′′</sup>$`\pm `$2<sup>′′</sup>, the interferometric maps (Lo et al. 1987; Shen & Lo 1995; Neininger et al. 1998) show a separation of 27<sup>′′</sup>$`\pm `$2<sup>′′</sup>. This difference is larger than the positional uncertainties. A larger separation in the low-$`J`$ transitions is also supported by filled-aperture measurements of Nakai et al. (1986, 1987) and Wild et al. (1992) for CO $`J`$ = 1–0, by Loiseau et al. (1990) for CO $`J`$ = 2–1, by Loiseau et al. (1988) for <sup>13</sup>CO $`J`$ = 2–1, and by us for the <sup>12</sup>CO $`J`$ = 1–0 and 2–1 lines. The CO $`J`$ = 3–2 distribution (Tilanus et al. 1991) shows an intermediate lobe separation.
Which peaks are detected in the CO submillimeter lines? The two main lobes of CO emission are located almost symmetrically with respect to the kinematical center, both in lines of low (e.g. Neininger et al. 1998) and high (see Fig. 3) excitation. We can therefore exclude that we see the north-eastern lobe and the ‘compact central core’ (see Shen & Lo 1995; Neininger et al. 1998) that are separated by $``$15<sup>′′</sup>. Instead, the sudden drop of angular separation from 27<sup>′′</sup> to 15<sup>′′</sup> must reveal inhomogeneities in the molecular ring that are characterized by changes in density and temperature. In addition to the ‘high CO excitation component’ there exists a ‘low CO excitation component’, mainly emitting in the CO $`J`$ = 1–0 and 2–1 lines. The transition in lobe separation occurs at the $`J`$ = 3–2 line: In CO $`J`$ = 3–2 (Tilanus et al. 1991) the separation is still $``$20<sup>′′</sup>. <sup>13</sup>CO $`J`$ = 3–2 emission with smaller optical depths and less photon trapping requires, however, higher excited gas so that the lobe separation becomes smaller.
Since lobe separations are often comparable to the angular resolution of the particular observation, integrated intensity maps may be misleading as they easily exhibit structure dominated by the superposition of components with identical lines of sight but different velocities. To clarify the situation, we thus present in Fig. 4 position-velocity maps of the CO $`J`$ = 4–3 (our HHT map with highest signal-to-noise ratio and best relative pointing) and CO 2–1 (Weiß, in preparation) line emission. Beam widths are 18<sup>′′</sup> and 13<sup>′′</sup>, respectively. Surprisingly, the p-v diagrams do not reproduce the strikingly different lobe separations seen in the integrated intensity maps. This also holds when comparing CO $`J`$ = 2–1 with 7–6 emission. This hints at wider line profiles for the higher excited CO transitions at the inner edges of the CO $`J`$ = 1–0 and 2–1 lobes. This is corroborated by Fig. 4 that shows, for the CO $`J`$ = 4–3 line, a slightly smaller lobe separation and a significantly larger full-width-to-half-power line width towards the inner edge of the south-western lobe.
### 4.2 Other tracers of the interstellar medium
Table 2 displays lobe separations determined by various tracers of atomic line, molecular line, or dust continuum emission. There are three preferred angular distances: $``$26<sup>′′</sup>, 15<sup>′′</sup>, and 10<sup>′′</sup>. Spectral lines of low excitation show two maxima at an angular distance of $``$26<sup>′′</sup>. Highly excited lines and the far-infrared and mm-wave continuum show a lobe separation of $``$15<sup>′′</sup>. The continuum in the near and mid infrared and the Br $`\gamma `$ line show an angular spacing of $``$10<sup>′′</sup>.
Since the dynamical center of the galaxy is located not too far from the mid point of the line connecting the lobes, these data show evidence for an even more complex structure than indicated by the CO data alone. The hot dust and Br $`\gamma `$ emission from the inner parts of the ring may trace the most recent star formation activity as it propagates outwards into the molecular lobes located at larger galactocentric distances (Alton et al. 1999).
Since the far infrared and submillimeter continuum are associated with the central portion of the ring, this is the location where column densities of the cool dense gas must be largest.
## 5 Excitation analysis
The $`J`$ = 7–6 transition is the CO line with highest rotational quantum number so far observed in M 82. Measurements of this line widen the range of excitation accessible by CO data considerably. To analyse these excitation conditions, we have performed radiative transfer calculations using a Large Velocity Gradient (LVG) model describing a cloud of spherical geometry (see Appendix A).
### 5.1 Physical parameters: Our data
The three submillimeter CO transitions measured by us arise from the high excitation component, show a similar spatial distribution, and can therefore be used simultaneously in a radiative transfer analysis. To obtain true line intensity ratios, the CO 4–3 and 7–6 spectra were smoothed to the angular resolution of the <sup>13</sup>CO 3–2 data. Table 1 displays these ratios for a beamwidth of 22<sup>′′</sup>, assuming gaussian beamshapes. As indicated by the spatial distributions (see Fig. 1), line ratios are similar toward the lobes and the central beam. CO 7–6/4–3 (integrated) line intensity ratios are at the order of 0.3, while <sup>12</sup>CO 4–3/<sup>13</sup>CO 3–2 ratios are 7 – 10 (see Table 1).
In Appendix A (Fig. 10) we plot our LVG line intensity ratios as a function of density, kinetic temperature, <sup>12</sup>CO/<sup>13</sup>CO ratio, and CO ‘abundance’ (i.e., $`X`$(CO)/grad $`V`$ with $`X`$(CO) = \[CO\]/\[H<sub>2</sub>\]). For the line parameters given in Table 1, reasonable solutions can be found for small CO abundances ($`X`$(CO)/grad $`V`$ $``$ 10<sup>-6</sup> pc/$`\mathrm{km}\mathrm{s}^1`$) and high <sup>12</sup>CO/<sup>13</sup>CO ratios ($``$ 75). High CO abundances ($`X`$(CO)/grad $`V`$ $``$ 10<sup>-3</sup> pc/$`\mathrm{km}\mathrm{s}^1`$) and small <sup>12</sup>CO/<sup>13</sup>CO ratios are however not providing a realistic solution, because H<sub>2</sub> densities would approach 10<sup>2</sup>$`\text{cm}^3`$ and would thus become prohibitively small. Calculations assuming a plane-parallel instead of a spherical cloud geometry would yield even smaller densities.
A comparison of linewidths and velocity drifts with the size of the region suggests grad $`V`$ $``$ 1 $`\mathrm{km}\mathrm{s}^1`$/pc. With this value and $`X`$(CO) $``$ 10<sup>-5</sup> – 10<sup>-4</sup> (e.g. Blake et al. 1987; Farquhar et al. 1994) we can exclude a <sup>12</sup>CO/<sup>13</sup>CO abundance ratio as small as 25 and are guided to the four diagrams of Fig. 10 with log $`X`$(CO) $``$ log ($`X`$(CO)/\[grad $`V`$/$`\mathrm{km}\mathrm{s}^1`$ pc<sup>-1</sup>\]) = $`5`$ or $`4`$ and <sup>12</sup>CO/<sup>13</sup>CO = 50 or 75. Our parameters are summarized in Table 3. Kinetic temperatures are $`T_{\mathrm{kin}}`$ $``$ 60 – 130 K, densities $`n(\mathrm{H}_2`$) $``$ 10<sup>3.3-3.9</sup>$`\text{cm}^3`$, and area filling factors $`f_{\mathrm{a},22^{\prime \prime }}`$ $``$ 0.04 – 0.07 and $`f_{\mathrm{a},15^{\prime \prime }}`$ $``$ 0.07 – 0.11. While the density is well constrained, the kinetic temperature is less well determined. Solutions with $`T_{\mathrm{kin}}`$ $``$ 150 K are however unlikely because they would require H<sub>2</sub> densities $``$10<sup>3</sup>$`\text{cm}^3`$. Even excluding such extreme solutions, the density of the CO emitting gas is small when compared with that of the Orion Hot Core (e.g. Schulz et al. 1995; van Dishoeck & Blake 1998), but column densities are large: With $`N`$(CO) = 3.08 10<sup>18</sup> \[$`n_{\mathrm{CO}}`$/$`\text{cm}^3`$\] \[($`\mathrm{\Delta }V`$/grad $`V`$)/pc\] $`\text{cm}^2`$, we obtain the 22<sup>′′</sup> beam averaged and cloud averaged column densities displayed in Table 3. For a line-of-sight source size of 350 pc and $`N`$(H<sub>2</sub>)$`_{22^{\prime \prime }}`$ = 10<sup>23</sup>$`\text{cm}^2`$ (see Table 3), the beam averaged mean molecular density is $`<`$$`n(\mathrm{H}_2`$)$`>_{22^{\prime \prime }}`$ $``$ 200 $`\text{cm}^3`$ and the volume filling factor becomes $`f_{\mathrm{v},22^{\prime \prime }}`$ = $`<`$$`n(\mathrm{H}_2`$)$`>`$/$`n(\mathrm{H}_2`$) $``$ 0.05. Since the ratio $`f_{\mathrm{v},22^{\prime \prime }}`$/$`f_{\mathrm{a},22^{\prime \prime }}`$ denotes the line of sight dimension of the clouds in units of the beam size, we obtain with $`r_{\mathrm{cloud}}`$ = 0.5 tg 22<sup>′′</sup> $`D_{\mathrm{pc}}`$ \[$`f_{\mathrm{v},22^{\prime \prime }}`$/$`f_{\mathrm{a},22^{\prime \prime }}`$\] $``$ 150 pc a characteristic cloud radius (that will be discussed and revised in Sects. 6.2 and 6.3). The total molecular mass is $`M_{\mathrm{mol},22^{\prime \prime }}`$ $``$ 1–7 10<sup>8</sup>M.
### 5.2 Physical parameters: All CO data
So far, we have only analysed the ‘warm CO component’ of M 82, exclusively seen in <sup>13</sup>CO $`J`$ = 3–2 and higher excited rotational transitions. To combine these results with data from lower $`J`$ rotational transitions and to further elucidate the physical state of the gas, Figs. 57 show (integrated) line intensities as a function of quantum number $`J`$. Calibration errors are at the order of $`\pm `$10% in the $`J`$ = 1–0 and 2–1 lines and $`\pm `$20% in the higher excited lines.
Figs. 5 and 6 show (integrated) line temperatures for a 22<sup>′′</sup> beam toward the dynamical center of the galaxy. Fig. 7 displays line temperatures for a 15<sup>′′</sup> beam toward the lobes and demonstrates that CO excitation is similar toward the south-western and north-eastern hotspot. Results from radiative transfer calculations (see Sect. 5.1 and Appendix A) are also given. Input parameters correspond to the four boxes in Fig. 10 that provide a promising fit to the data (log ($`X`$(CO)/\[grad $`V`$/$`\mathrm{km}\mathrm{s}^1`$ pc<sup>-1</sup>\] = $`5`$ or $`4`$ and <sup>12</sup>CO/<sup>13</sup>CO = 50 or 75).
Apparently there is a problem with the CO $`J`$ = 4–3 line: The integrated intensity (Fig. 6) does not allow a reasonable fit, while the peak intensity (Fig. 5) is ‘appropriate’. This is caused by the narrow lineshape of our 4–3 spectrum (Fig. 2). Compared with other lines, the integrated line intensity is too small, but the peak line temperature is almost ‘normal’. The CO $`J`$ = 4–3 profile, shown by Güsten et al. (1993) for the central position, is also weak, both with respect to its peak and integrated line intensity. The emission from the lobes, however, fits into the general trend (Fig. 7). An interpretation in terms of a diminished lobe separation for CO $`J`$ $``$ 4–3 (see Sect. 4.1) is not conclusive: The 22<sup>′′</sup> beam centered on the dynamical core of M 82 is mainly confined to the inner parts of the lobes and should not be greatly affected by emission from further out. A CO $`J`$ = 4–3 deficiency is neither seen in Fig. 7 nor in a corresponding plot showing integrated intensities for a 15<sup>′′</sup> beam.
To summarize: Data from the $`J`$ = 4–3 line are contradictory so that convincing evidence for a true anomaly is missing. In spite of differences in lobe separation (Sect. 4.1), the CO data can be reproduced, within observational errors, with densities, temperatures, filling factors, \[CO\]/\[H<sub>2</sub>\] abundances, and <sup>12</sup>CO/<sup>13</sup>CO isotope ratios deduced from the three submillimeter transitions mapped by us with the HHT (Fig. 1).
There are few constraints for the low excitation component that is mainly seen in the CO $`J`$ = 1–0 and 2–1 lines. The column density must be smaller than for the high excitation component because the latter is more closely related to the far infrared and submillimeter continuum from the dust (Sect. 4.2). In view of the remarkable number of ‘super’-star clusters (O’Connell et al. 1995) and supernovae (Kronberg et al. 1985) near the outer portions of the ring, cloud temperatures for the low excitation component should also be $`T_{\mathrm{kin}}`$ $``$ 10 K. For $`T_{\mathrm{kin}}`$ $``$ 50 K, LVG densities are at the order of $`n(\mathrm{H}_2`$) $``$ 10<sup>3</sup>$`\text{cm}^3`$ or less.
## 6 Discussion
### 6.1 A comparison with other LVG simulations
The most detailed models of CO emission from M 82 were so far provided by Güsten et al. (1993). In order to fit the <sup>12</sup>CO/<sup>13</sup>CO $`J`$ = 1–0 and 2–1 line intensity ratios then available, they rejected a one component LVG scenario and introduced two gas components, one of low ($`n(\mathrm{H}_2`$) $``$ 10<sup>3</sup>$`\text{cm}^3`$) and one of high ($``$10<sup>5</sup>$`\text{cm}^3`$) density. The low density component is similar to that proposed by us for the gas mainly emitting in the CO $`J`$ = 1–0 and 2–1 lines (see Sect. 4.1). Our parameters for the high CO excitation component agree to a large extent with those of the one component scenario of Güsten et al. (1993).
For the highly excited gas we thus also find optically thick <sup>12</sup>CO low-$`J`$ emission and a high \[<sup>12</sup>CO\]/\[<sup>13</sup>CO\] abundance ratio ($``$50; see Table 3), that is further supported by an independent chain of arguments involving CN and <sup>13</sup>CN data (see Henkel et al. 1998). While the LVG model result is shown to be inconclusive in Sect. 6.3.2, the CN data suggest that the \[<sup>12</sup>CO\]/\[<sup>13</sup>CO\] ratio is larger than that observed in the galactic center region. Likely explanations are radial infall of <sup>13</sup>CO deficient gas from the outer parts of the galaxy or a <sup>12</sup>C excess in the ejecta from massive stars (e.g. Henkel & Mauersberger 1993). In agreement with Güsten et al. (1993) we also find that excitation in the south-western and north-eastern lobe is similar. Since our new data allow us to constrain kinetic temperatures to the high end of those predicted by the one component scenario of Güsten et al. ($`T_{\mathrm{kin}}`$ $``$ 30–70 K), our densities are at the low end ($`n(\mathrm{H}_2`$) $``$ 10<sup>4</sup>$`\text{cm}^3`$, since the product $`T_{\mathrm{kin}}`$ $`n(\mathrm{H}_2`$)<sup>1/2</sup> is approximately conserved among models simulating optically thick subthermal CO emission). There is a remarkable agreement with respect to CO column density, H<sub>2</sub> density, and molecular gas mass between various studies (cf. Tilanus et al. 1991; Wild et al. 1992; Güsten et al. 1993).
### 6.2 Intrinsic inconsistencies of the model
Our CO data could be reproduced by assuming the presence of two gas components. Selecting such a ‘best’ model, we could discriminate between previously proposed scenarios and could constrain cloud conditions giving rise to highly excited CO emission to slightly higher temperatures and smaller densities than previously suggested. On a deeper level, however, there remain problems. While calculated CO column densities (Sect. 5.1) appear to be correct (see e.g. Fig. 2 of Lo et al. 1987 and Fig. 3 of Smith et al. 1991), an obvious puzzle is the large volume filling factor ($`f_{\mathrm{v},22^{\prime \prime }}`$ $``$ 0.05) that is comparable to the area filling factor (Sect. 5.1). This forced us to postulate a characteristic cloud radius ($`r_{\mathrm{cloud}}`$ $``$ 150 pc) that encompasses a large part of the studied volume. Such a large cloud radius is inconsistent with the expectation of similar scale lengths along the three dimensions ($`f_{\mathrm{v},22^{\prime \prime }}`$ $``$ $`f_{\mathrm{a},22^{\prime \prime }}^{3/2}`$) and with the spatial fine structure deduced from CO (e.g. Shen & Lo 1995; Neininger et al. 1998), high density tracers like HCN and HCO<sup>+</sup> (e.g. Brouillet & Schilke 1993; Paglione et al. 1997; Seaquist et al. 1998), and infrared fine structure lines (e.g. Lugten et al. 1986; Lord et al. 1996).
Another problem is the density of the gas observed. A density of $`n(\mathrm{H}_2`$) $``$ 10<sup>4</sup>$`\text{cm}^3`$ is small when compared to the prototypical ‘hot core’ associated with the Orion nebula, violating a sometimes noted similarity between these tiny, highly obscured galactic star forming regions and the more extended starbursts studied in external galaxies (e.g. Lo et al. 1987; Wolfire et al. 1990). More seriously, the molecular density determined by us for the starburst in M 82 is smaller than most theoretical studies and observational data permit: Assuming ‘reasonable’ density stucture ($`n(\mathrm{H}_2`$) $``$ $`r^{1\mathrm{}2}`$) and accounting for the intense UV field, Brouillet & Schilke (1993) find that molecular clouds with densities less than a few times 10<sup>4</sup>$`\text{cm}^3`$ should not exist in the central region of M 82. For the transition region between the atomic and molecular gas they propose a density of 10<sup>4-5</sup>$`\text{cm}^3`$ which disagrees with the range of densities deduced from our LVG analysis in Sect. 5.1 (for cloud stability against tidal stress, see Appendix B).
Studies of high density tracers (e.g. CS, HCN, and HCO<sup>+</sup>) commonly reveal densities $`n(\mathrm{H}_2`$) $``$ 10<sup>4</sup>$`\text{cm}^3`$ (e.g. Mauersberger & Henkel 1989; Brouillet & Schilke 1993; Paglione et al. 1997; Seaquist et al. 1998), even for the low excitation component (compare Mauersberger & Henkel 1989 with Baan et al. 1990). In view of $`T_{\mathrm{dust}}`$ = 48 K (Hughes et al. 1994; Colbert et al. 1999), radiative excitation by infrared photons should not significantly alter the density estimates (see Carrol & Goldsmith 1981 and Appendix C). Millimeter wave recombination lines indicate the presence of an ionized component with low filling factor and high electron density ($`n_\mathrm{e}`$ $`>`$ 10<sup>4.5</sup>$`\text{cm}^3`$; Seaquist et al. 1996). This gas may be associated with a population of (ultra)compact HII regions or with shock ionized dense molecular material. Far infrared fine-structure lines, if tracing the interface between the molecular and the ionized gas, indicate densities of 10<sup>3.3-4.0</sup>$`\text{cm}^3`$ (e.g. Lugten et al. 1986; Wolfire et al. 1990; Lord et al. 1996; Stutzki et al. 1997; Colbert et al. 1999) that may agree too well with those derived by us for CO (Sect. 5.1).
An interesting aspect is also provided by kinetic temperature estimates: With $`T_{\mathrm{kin}}`$ $``$ 20–60 K (Seaquist et al. 1998; this temperature is consistent with the estimated cosmic ray flux; see Völk et al. 1989, Suchkov et al. 1993, and Fig. 2 of Farquhar et al. 1994) the temperature of the dense molecular gas appears to agree fairly well with that of the dust ($`T_{\mathrm{dust}}`$ = 48 K; Hughes et al. 1994; Colbert 1999). This is further supported by an apparent lack of CH<sub>3</sub>OH and SiO emission, two tracers of high temperature gas that are easily seen in other nearby starburst galaxies (Henkel et al. 1991; Mauersberger & Henkel 1993). The surrounding neutral and ionized layers have larger temperatures, at the order of 50–100 to 200 K (e.g. Lugten et al. 1986; Wolfire et al. 1990; Lord et al. 1996; Colbert et al. 1999). An analysis of \[C i\] 492 and 809 GHz emission from the south-western lobe of M 82 (Stutzki et al. 1997) shows a particularly striking similarity in density, temperature, and area (but not volume) filling factor with our LVG parameters derived for CO (Sect. 5.1; the \[C i\]/CO abundance ratio then becomes $``$0.4). It seems as if CO were an integral part of the dense atomic gas layers of M 82.
### 6.3 A PDR model for the CO emission from M 82
#### 6.3.1 General aspects
To resolve inconsistencies related to CO fine scale structure, H<sub>2</sub> density, and kinetic temperature (Sect. 6.2), we note that a significant fraction of the molecular gas in the Milky Way lies in photon dominated regions (PDRs; see Hollenbach & Tielens 1997). A PDR model is also successfully applied to simultaneously explain CO and \[Cii\] line intensities in the central region of the late-type spiral IC 342 that is not a starburst galaxy but that is believed to be a face-on ‘mirror image’ of the galactic center region (Schulz et al. in preparation).
In the central region of M 82 with its high UV flux, the bulk of the interstellar line radiation likely arises from PDRs. As we have seen (Sect. 6.2), the coolest gas component is observed toward the dense cloud cores, a situation that is consistent with the inside-out temperature gradient expected in the case of PDRs. While CO excitation by infrared radiation from the dust can be neglected (see Appendix C), the impact of UV photons may strongly affect rotational level populations and measured CO line intensities in regions with high UV flux. Estimating the far-UV flux (6.0 – 13.6 eV) in the starburst region, Wolfire et al. (1990), Stacey et al. (1991), Lord et al. (1996), and Colbert et al. (1999) find $`\chi `$ $``$ 10<sup>2.8-3.9</sup> ($`\chi `$: incident far-UV flux in units of the local galactic flux, 1.6 10<sup>-3</sup> erg s<sup>-1</sup>$`\text{cm}^2`$). Such a high value leads to outer cloud layers that are predominantly heated by collisions with electrons photoejected from dust grains or through collisional deexcitation of vibrationally excited, UV-pumped H<sub>2</sub>.
Exploring effects of finite cloud size for clumps with plane-parallel geometry, Köster et al. (1994) presented comprehensive computations of CO rotational line intensities from PDRs. In contrast to the constant temperature LVG treatment presented in Sect. 5, PDR models account for strong kinetic temperature gradients. The low-$`J`$ lines typically arise from deeper inside the clouds than the mid-$`J`$ lines observed by us (Sect. 3.1). Reproducing observed line ratios of optically thick <sup>12</sup>CO emission, PDR model densities are larger than those derived from a one temperature, one density LVG model. PDR line intensity ratios only surpass unity if the higher-$`J`$ line has a critical density (for the values, see Sect. 1) that is smaller than the actual gas density. Then, both lines are approximately thermalized but the higher $`J`$-line is emitted from warmer layers further out. For densities at the order of 5 10<sup>3</sup>$`\text{cm}^3`$ (the density deduced from our one component LVG model in Sect. 5.1), <sup>12</sup>CO PDR line temperatures rapidly decrease with rotational quantum number $`J`$. For $`\chi `$ $``$ 10<sup>3</sup> and $`N`$(H<sub>2</sub>) $``$ 10<sup>22</sup>$`\text{cm}^2`$/$`\mathrm{km}\mathrm{s}^1`$ (Sect. 5.1), densities are then at the order of 10<sup>4</sup> and 10<sup>5</sup>$`\text{cm}^3`$ for the low and high excitation components, respectively.
For the gas with low CO excitation, a density of $`n(\mathrm{H}_2`$) $``$ 10<sup>4</sup>$`\text{cm}^3`$ is sufficiently high to explain the detection of CS $`J`$ = 2–1 emission (e.g. Mauersberger & Henkel 1989; Baan et al. 1990). A density of $`n(\mathrm{H}_2`$) $``$ 10<sup>5</sup>$`\text{cm}^3`$ in the high excitation region would fulfill all theoretical (Brouillet & Schilke 1993) and observational density requirements outlined in Sect. 6.2. Furthermore, the volume filling factor $`f_{\mathrm{v},22^{\prime \prime }}`$ = $`<`$$`n(\mathrm{H}_2`$)$`>_{22^{\prime \prime }}`$/$`n(\mathrm{H}_2`$) would drop by one to two orders of magnitude below the value estimated in Sect. 5.1. If the area filling factor is not drastically altered, this yields reasonable molecular cloud radii $``$10 pc (see Sect. 5.1 and Lugten et al. 1986; Wolfire et al. 1990; Brouillet & Schilke 1993; Shen & Lo 1995; Seaquist et al. 1996; Stutzki et al. 1997). The apparently intermediate temperature of the CO emitting gas between those of the dust and the dense atomic medium is naturally explained by mid-$`J`$ CO emission predominantly arising in the heated surface layers of dense molecular clouds.
To summarize: PDR simulations of <sup>12</sup>CO emission remove inconsistencies related to spatial fine scale structure, density, and kinetic temperature. So far published PDR results fail however when <sup>13</sup>CO is also considered. Calculated <sup>12</sup>CO/<sup>13</sup>CO line intensity ratios are much smaller than the ratios observed (see Table 1 and Köster et al. 1994) and a more detailed numerical analysis is therefore needed.
#### 6.3.2 PDR model calculations
For a numerical approach we note that the beam averaged column density ($`N`$(H<sub>2</sub>)$`_{22^{\prime \prime }}`$ $``$ 10<sup>23</sup>$`\text{cm}^2`$) and the beam averaged density ($`<`$$`n(\mathrm{H}_2`$)$`>_{22^{\prime \prime }}`$ $``$ 200 $`\text{cm}^3`$) are observationally determined (Sects. 5.1 and 6.2) and do not depend on the choice of the excitation model if a significant fraction of the dust is associated with molecular clouds. Complementing these boundary conditions, we obtain (cf. Sect. 5.1)
$$\frac{r_{\mathrm{cloud}}}{[\mathrm{pc}]}175\frac{200\mathrm{cm}^3/n(\mathrm{H}_2)_{\mathrm{PDR}}}{T_{\mathrm{CO},\mathrm{observed}}/T_{\mathrm{CO},\mathrm{PDR}}}$$
(1)
and
$$\frac{N(\mathrm{H}_2)_{\mathrm{PDR}}}{[\mathrm{cm}^2\mathrm{km}\mathrm{s}^1]}\frac{10^{23}/300}{T_{\mathrm{CO},\mathrm{observed}}/T_{\mathrm{CO},\mathrm{PDR}}}$$
(2)
with 200 $`\text{cm}^3`$/$`n(\mathrm{H}_2`$)<sub>PDR</sub> denoting the volume filling factor, $`T_{\mathrm{CO},\mathrm{observed}}/T_{\mathrm{CO},\mathrm{PDR}}`$ being the area filling factor, $`n(\mathrm{H}_2`$)<sub>PDR</sub> giving the average density, and $`N`$(H<sub>2</sub>)<sub>PDR</sub> representing the average column density of an individual clump per $`\mathrm{km}\mathrm{s}^1`$ (300 $`\mathrm{km}\mathrm{s}^1`$ is the total CO linewidth of the nuclear region of M 82; $`N`$(H<sub>2</sub>)<sub>PDR</sub> = (4/3) $``$$`r_{\mathrm{cloud}}`$$``$$`n(\mathrm{H}_2`$)<sub>PDR</sub> and $`n(\mathrm{H}_2`$)<sub>PDR-cloud-surface</sub> = 0.5 $``$$`n(\mathrm{H}_2`$)<sub>PDR</sub> for the assumed cloud geometry and density structure, see below).
Since embedded clumps of dense gas may more often be spherical than plane-parallel, spherical clouds with a power law density distribution ($`n`$($`r`$) $``$ $`r^{1.5}`$) were modeled with the dust being heated by the external UV radiation and intrinsic infrared emission (cf. Störzer et al. 1996; 2000). While we can reproduce the observed relative intensities of the various rotational CO transitions with a density at the order of $`n(\mathrm{H}_2`$) $``$ 10<sup>5</sup>$`\text{cm}^3`$ (see Fig. 8 and Sect. 6.3.1), $`r_{\mathrm{cloud}}`$ and $`N`$(H<sub>2</sub>) are inconsistent with Eqs. 1 and 2. Furthermore, computed <sup>12</sup>CO/<sup>13</sup>CO line intensity ratios show little dependence on cloud structure and remain with characteristic values of 2–4 (see also Gierens et al. 1992 for <sup>12</sup>C/<sup>13</sup>C = 40; Köster et al. 1994 and Störzer et al. 2000 for <sup>12</sup>C/<sup>13</sup>C = 67) much smaller than the observed ratios of 10–15 (see Table 1). Varying $`\chi `$ over the permitted range (10<sup>3.3±0.4</sup>) does not change these line intensity ratios significantly. Drastically increasing the carbon isotope ratio leads to <sup>12</sup>C/<sup>13</sup>C $`>`$ 100 which is inconsistent with data from the solar system, the Milky Way, and the Magellanic Clouds (e.g. Wilson & Rood 1994; Chin et al. 1999).
Other parameters that can be varied are the density and column density of the far-UV irradiated cloud. The <sup>12</sup>CO lines are formed in warm layers near the surface whereas <sup>13</sup>CO lines are predominantly emitted from cooler regions deeper inside. With the column density being fixed, a higher density leads to a moderate increase in the <sup>12</sup>CO/<sup>13</sup>CO line intensity ratio (see Fig. 8), because <sup>12</sup>CO is then emitted from CO layers closer to the heated cloud surface. To reproduce observed <sup>12</sup>CO/<sup>13</sup>CO line ratios in this way requires, however, densities $``$10<sup>8</sup>$`\text{cm}^3`$.
More important are variations in column density: A spherical UV illuminated cloud is effectively smaller in <sup>13</sup>CO than in <sup>12</sup>CO and this difference becomes more pronounced when $`N`$(H<sub>2</sub>) decreases (see Störzer et al. 2000). In the extreme case of totally photodissociated <sup>13</sup>CO, some <sup>12</sup>CO may still exist in the clump so that, in principle, <sup>12</sup>CO/<sup>13</sup>CO line intensity ratios up to very large values can be reproduced. Once the <sup>13</sup>CO abundance gets very small, <sup>12</sup>CO also tends to become optically thin. In Fig. 9 we show expected line temperatures in this optically thin regime, with $`n(\mathrm{H}_2`$) = 5 10<sup>3</sup>$`\text{cm}^3`$ successfully reproducing measured <sup>12</sup>CO line intensity ratios for the high CO excitation component. <sup>12</sup>CO/<sup>13</sup>CO = 19, 10, 8.5, 9.4, 13, 20, and 30 for $`J`$ = 1–0 … 7–6, respectively. This is compatible with observed ratios of 15, 10, and 9 for the three lowest rotational transitions. For the low excitation component, $`n(\mathrm{H}_2`$) $``$ 10<sup>3</sup>$`\text{cm}^3`$; Fig. 9 shows a fit with <sup>12</sup>CO/<sup>13</sup>CO = 17, 10, and 11 for $`J`$ = 1–0 to 3–2.
It is remarkable that PDR model densities are similar to those derived with our LVG simulation, contradicting theoretical predictions of minimum densities in excess of 10<sup>4</sup>$`\text{cm}^3`$ (Brouillet & Schilke 1993). Column densities ($`N`$(H<sub>2</sub>) $``$ 5 10<sup>20</sup>$`\text{cm}^2`$/ $`\mathrm{km}\mathrm{s}^1`$) and line temperatures of individual cloudlets are, however, drastically different. An average cloudlet shows optically thin, not optically thick CO line emission. Line temperatures of individual clumps are much smaller than the observed average over the inner 300 pc, so that the PDR area filling factor $`f_\mathrm{a}`$ is not $``$1 as in the LVG approximation but $``$1. Thus not a higher density (as anticipated in Sect. 6.3.1) but a higher $`f_\mathrm{a}`$ leads to plausible $`r_{\mathrm{cloud}}`$ values in the sub-parsec range (see Eq. 1). Both $`f_\mathrm{a}`$ $``$ 1 and $`r_{\mathrm{cloud}}`$ $`<`$ 1 pc are consistent with early CO studies, favoring CO lines of low optical depth. These conclusions are supported by observations of fine-structure lines in the infrared (e.g. Olofsson & Rydbeck 1984; Lugten et al. 1986; Wolfire et al. 1990; Brouillet & Schilke 1993; Schilke et al. 1993; Lord et al. 1996). In view of our PDR simulations, the high density ($`>`$10<sup>4.5</sup>$`\text{cm}^3`$) ionized gas found by Seaquist et al. (1996) is likely related to (ultra)compact HII regions and not to evaporating dense molecular clouds.
Obviously, our study does not imply that the core of M 82 does not contain regions of high density and column density. The bulk of the CO emission, however, must arise from a warm low density interclump medium. The gas may be barely dense enough to avoid tidal disruption (see Appendix B). For typical cloud cores in the galactic disk, Eq. 5 of Larson (1981) predicts a column density of $`N`$(H<sub>2</sub>) $``$ 10<sup>22</sup>$`\text{cm}^2`$. With $`N`$(H<sub>2</sub>) $``$ 5 10<sup>20</sup>$`\text{cm}^2`$/$`\mathrm{km}\mathrm{s}^1`$ (Table 4) this infers a linewidth of 20 $`\mathrm{km}\mathrm{s}^1`$. The relations defined for approximately virialized clouds by Larson (1981) then yield cloud sizes well in excess of 100 pc, in strong contradiction with our PDR model. The bulk of the CO emission therefore arises from gas that may not be virialized. We conclude that CO is tracing a different gas component than molecular high density tracers like CS or HCN. The dominance of CO emission from a diffuse medium with cloud fragments of low column density per $`\mathrm{km}\mathrm{s}^1`$, $`X`$ = $`N`$(H<sub>2</sub>)/$`I_{\mathrm{CO}}`$ conversion ratio, and CO intensity explains why HCN is a better tracer of star formation and infrared luminosity (Solomon et al. 1992a). It also explains why \[C i\]/CO abundance ratios are higher than those observed in spatially confined star forming regions of the galactic disk (see also Sect. 6.2, Schilke et al. 1993, and White et al. 1994).
An argument against our scenario could be that Eqs. 1 and 2 only hold approximately: The denominator $`T_{\mathrm{CO},\mathrm{observed}}/T_{\mathrm{CO},\mathrm{PDR}}`$ is too small by an order of magnitude to yield the proper $`r_{\mathrm{cloud}}`$ value assumed in the PDR approximation, while it is too large by the same amount to provide the proper $`N`$(H<sub>2</sub>)<sub>PDR</sub> value. In order to check whether our scenario is plausible, we therefore extend our analysis to other galaxies.
### 6.4 Other galaxies
#### 6.4.1 The Milky Way and other ‘quiescent’ galaxies
Studying the physics of the molecular gas in the galactic center region, Dahmen et al. (1998) compared <sup>12</sup>CO and C<sup>18</sup>O $`J`$ = 1–0 data obtained with a linear resolution of $``$22 pc (8). Lineshapes, spatial distributions, and line ratios indicate the presence of an extended diffuse gas component that is not apparent in studies focusing on individual cloud cores (see also Oka et al. 1998a). A significant fraction of the gas is not virialized. Higher resolution data (e.g. Hüttemeister et al. 1998) reveal that the dense molecular gas ($`n(\mathrm{H}_2`$) $``$ 10<sup>4</sup>$`\text{cm}^3`$) is relatively cool ($`T_{\mathrm{kin}}`$ $``$ 25 K) in comparison to that at lower density ($``$ 100 K at a few 10<sup>3</sup>$`\text{cm}^3`$). Apparently, a warm diffuse molecular medium is not only ubiquitous in the starburst galaxy M 82 but also dominates the CO emission in the more quiescent central region of the Milky Way (for extended CO maps of a galactic disk star forming region, see Wilson et al. 1999). This also holds for the nuclear regions of IC 342 (Schulz et al. in preparation) and NGC 7331 (Israel & Baas 1999), another two galaxies with $`L_{\mathrm{FIR}}`$ $``$ 10<sup>10</sup>L. Like M 82, NGC 7331 shows a warmer inner and cooler outer dust ring and a molecular ring that appears to be related to its cool dust component.
Shocks and cloud-cloud collisions induced by the presence of bars (e.g. Achtermann & Lacey 1995; Morris & Serabyn 1996; Fux 1997, 1999; Hüttemeister et al. 1998), tidal disruption of clouds near the center (Güsten 1989), a high gas pressure (e.g. Helfer & Blitz 1997) that can help to keep the gas molecular, and a high stellar density that can affect molecular cloud dynamics (e.g. Mauersberger et al. 1996b; Oka et al. 1998b) may all contribute to the disintegration of molecular clouds and to the synthesis of an extended warm molecular spray consisting of low column density cloud fragments.
#### 6.4.2 Nearby starburst galaxies
So far, M 82 is the only starburst galaxy where we could show that a CO-LVG excitation analysis (inferring $`f_\mathrm{v}`$ $``$ $`f_\mathrm{a}`$; see Sect. 5.1) does not lead to results that are free of inconsistencies. The lack of CO $`J`$ = 7–6 data does not constrain LVG parameters sufficiently to search for a similar discrepancy in the other two nearby starburst galaxies NGC 253 and NGC 4945. This implies that previous studies did not provide sufficient motivation to replace LVG excitation analyses by PDR scenarios.
Towards NGC 253, M 82, and NGC 4945 <sup>12</sup>CO and <sup>12</sup>CO/<sup>13</sup>CO line intensity ratios are similar. This even holds for absolute integrated line intensities (within 30%) for a 22–24<sup>′′</sup> beam size (compare Table 1 with Mauersberger et al. 1996a,b; Harrison et al. 1999). Not surprisingly, published LVG densities and area filling factors for NGC 253 and NGC 4945 match those for M 82 (Henkel et al. 1994; Mauersberger et al. 1996a,b; Harrison et al. 1999). Following the procedure outlined in Sect. 5.1, average densities along the line-of-sight ($`<`$$`n(\mathrm{H}_2`$)$`>`$ $``$ 100 $`\text{cm}^3`$) yield volume filling factors at the order of $`f_{\mathrm{v},22^{\prime \prime }}`$ = $`<`$$`n(\mathrm{H}_2`$)$`>`$/$`n(\mathrm{H}_2`$) $``$ 0.02. Since volume filling factors seem to be slightly smaller than area filling factors, average cloud sizes become $``$30 pc (for the Milky Way, see Oka et al. 1998b), not enough to request the use of PDR models.
Because of the similarity of CO line strengths and line ratios and since high \[C i\]/CO ratios are observed toward both M 82 and NGC 253 (Schilke et al. 1993; White et al. 1994; Harrison et al. 1995; Israel et al. 1995; Stutzki et al. 1997) we feel nevertheless that PDR simulations relating the bulk of the CO emission to a warm diffuse molecular medium are relevant not only to M 82 but to NGC 253 and NGC 4945 as well.
#### 6.4.3 Mergers
Active disk galaxies exhibit <sup>12</sup>CO/<sup>13</sup>CO ratios $``$ 10 (this is independent of inclination, Hubble Type, and metallicity; Sage & Isbell 1991). Perturbed gas-rich mergers with infrared luminosities $``$ 10<sup>11</sup>L tend to show larger values, up to 40–60 (Henkel & Mauersberger 1993). This may be caused by excitation effects (see Aalto et al. 1999) or by a deficiency of <sup>13</sup>CO emission since $`L_{\mathrm{FIR}}`$ is better correlated with <sup>12</sup>CO (Taniguchi & Ohyama 1998). Inflow of gas from the outer disks may provide <sup>13</sup>CO deficient gas and nucleosynthesis in shortlived massive stars may further enhance <sup>12</sup>C (Casoli et al. 1992). While this may help to enhance <sup>12</sup>CO/<sup>13</sup>CO line intensity ratios, this is not sufficient to explain measured <sup>12</sup>CO/<sup>13</sup>CO values within the context of PDRs. Instead we suggest that the direct interaction of two galactic nuclei and their associated disks is even more efficient than the presence of a bar to trigger cloud-cloud collisions and to create a warm diffuse molecular debris containing a large number of small cloud fragments. If such fragments are smaller than in individual non-merging galaxies, higher <sup>12</sup>CO/<sup>13</sup>CO line ratios might result. While the detailed spatial CO fine-structure may only be revealed by next generation mm-wave telescopes (e.g. Downes 1999), this would imply that in mergers with <sup>12</sup>CO/<sup>13</sup>CO $``$ 10, \[C i\]/CO line intensity ratios should be as large as or even larger than in M 82 and NGC 253. For a first detection of \[C i\] in a merger, see Gerin & Philips (1998).
#### 6.4.4 Galaxies at high redshifts
Is it a general property of the integrated spectrum of a starburst that, beginning with the CO $`J`$ = 4–3 transition, line intensities (accounting for beam dilution) gradually decrease with increasing rotational quantum number $`J`$? Ground based measurements of nearby galaxies in mid and high-$`J`$ CO transitions require exceptional weather conditions. In galaxies of high redshift, however, many otherwise inaccessible transitions are shifted into the observable mm- and submm-wavelength bands. While linear resolutions remain poor, the bulk of the CO emission from higher excited states should arise from the nuclear starburst environment. Suitable examples of distant sources with a number of detected CO transitions are IRAS F 1024+4724 ($`z`$ = 2.286), the Cloverleaf quasar ($`z`$ = 2.558), and BR 1202–0725 ($`z`$ = 4.692). Toward IRAS F 1024+4724, the $`J`$ = 3–2 line is stronger than the 4–3 and 6–5 lines; line intensity ratios w.r.t. 3–2 are $``$0.75 and 0.6 (Solomon et al. 1992b). Toward the Cloverleaf quasar, CO lines are characterized by a constant or increasing brightness temperature from 3–2 to 4–3, followed by constant or (more likely) decreasing line intensities in the higher $`J`$ transitions (Barvainis et al. 1997). Toward BR 1202–0725, the $`J`$ = 7–6 to 5–4 line intensity ratio is $``$0.65 (Omont et al. 1996). Although the sample is too small for a reliable statistical analysis, it seems that CO line intensity ratios are fairly uniform in starbursts and depend little on redshift (age) and temperature ($`T_{\mathrm{cmb}}`$/\[K\] = 2.73 (1+$`z`$)) of the microwave background. Unfortunately, <sup>12</sup>CO/<sup>13</sup>CO line intensity ratios are not yet known. If they are $``$10, densities at the order of a few 10<sup>3</sup>$`\text{cm}^3`$ and small column densities ($``$10<sup>21</sup>$`\text{cm}^2`$/$`\mathrm{km}\mathrm{s}^1`$) should be a characteristic feature for the bulk of the CO emitting cloudlets in all starburst galaxies.
## 7 Conclusions
We have studied millimeter and submillimeter CO line emission up to the $`J`$ = 7–6 rotational transition toward the central region of the starburst galaxy M 82 and obtain the following main results:
1. The spatial structure of the millimeter and submillimeter CO emission is distinct. While integrated intensity maps suggest that the lobe separation of the low-$`J`$ transitions is $``$26<sup>′′</sup>, it is $``$15<sup>′′</sup> for the mid-$`J`$ transitions. Major-axis position-velocity maps in the CO $`J`$ = 2–1 and 4–3 lines show however agreement in the lobe positions. This indicates that, at the inner edges of the low-$`J`$ CO lobes, line profiles are wider in the higher excited CO transitions. We thus distinguish between a ‘low’ and a ‘high’ CO excitation component, the latter coinciding with the main source of millimeter and submillimeter dust emission.
2. An LVG excitation analysis of CO submillimeter lines leads to internal inconsistencies. While measured line intensities are reproduced with $`T_{\mathrm{kin}}`$ $``$ 60 – 130 K, $`n(\mathrm{H}_2`$) $``$ 10<sup>3.3-3.9</sup>$`\text{cm}^3`$, cloud averaged column densities $`N`$(CO)<sub>cloud</sub> $``$ 10<sup>20</sup> and $`N`$(H<sub>2</sub>)<sub>cloud</sub> $``$ 10<sup>24-25</sup>$`\text{cm}^2`$, \[<sup>12</sup>CO\]/\[<sup>13</sup>CO\] abundance ratios $``$50, and a total molecular mass of a few 10<sup>8</sup>M, area filling factors ($`f_\mathrm{a}`$ $``$ 0.05–0.10) and volume filling factors ($`f_\mathrm{v}`$ $``$ 0.05) are similar. This results in cloud sizes that do not match their angular scale. On the other hand, the resulting H<sub>2</sub> column density is consistent with that derived from the dust continuum at millimeter and submillimeter wavelengths. For the low excitation component, densities are $``$ 10<sup>3</sup>$`\text{cm}^3`$.
3. An application of PDR models resolves the inconsistencies of the LVG calculations. LVG densities, column densities, and total mass are confirmed. The bulk of the CO emission arises, however, from a diffuse, low column density ($`N`$(H<sub>2</sub>) $``$ 5 10<sup>20</sup>$`\text{cm}^2`$/$`\mathrm{km}\mathrm{s}^1`$) interclump medium with small $`X`$ = $`N`$(H<sub>2</sub>)/$`I_{\mathrm{CO}}`$ conversion factors, area filling factors $``$1, and sub-parsec cloud sizes. The relations defined by Larson (1981) are not fulfilled and the gas may not be virialized. Such a scenario explains why CS or HCN are better tracers of global star formation rate and infrared luminosity than CO. Our scenario also explains observed high \[C i\]/CO line intensity ratios, while relative abundances of <sup>12</sup>CO versus <sup>13</sup>CO cannot be accurately determined. Higher column density clouds, even accounting for variations in far-UV flux and <sup>12</sup>C/<sup>13</sup>C isotope ratios, do not reproduce observed <sup>12</sup>CO/<sup>13</sup>CO line intensity ratios $``$10. Densities are close to the minimum values required for tidal stability in the absence of magnetic fields.
4. In regard to <sup>12</sup>CO line intensity ratios, the central region of M 82 appears to be representative for the entire family of starburst galaxies, both at small and at high redshifts. A comparison of the starburst regions in M 82 and NGC 253 demonstrates that this similarity extends to <sup>12</sup>CO/<sup>13</sup>CO and \[C i\]/CO line intensity ratios. The large <sup>12</sup>CO/<sup>13</sup>CO line intensity ratios ($``$10) observed toward ‘nearby’ mergers prove, however, that differences exist at least w.r.t. rare CO isotopomers. Galaxy pairs with such high <sup>12</sup>CO/<sup>13</sup>CO line ratios require the presence of a particularly diffuse highly fragmented low column density ISM.
Apparently, dropping the assumption of constant temperature in the CO excitation model is a necessary step to provide a self-consistent approach to the physical properties of molecular clouds in the nuclear starburst region of M 82. While the use of PDR models is crucial for a better understanding of the molecular gas phase in a starburst environment, important information is still missing. Interferometic observations of high density tracers (e.g. CN, CS, HCN, HNC, N<sub>2</sub>H<sup>+</sup>), coupled with PDR model calculations including chemical aspects, are needed to fully understand the spatial morphology, density distribution, and molecular excitation of this archetypical starburst complex. An interesting aspect is provided by HCO<sup>+</sup> $`J`$ = 1–0 line emission (Table 2). The bulk of this emission might arise from regions intermediate between those of the low and high CO excitation component. Since this molecule (as well as N<sub>2</sub>H<sup>+</sup>) is a sensitive tracer of ionization conditions in the dense gas, a detailed knowledge of its spatial distribution would be crucial for a better understanding of structure and excitation.
So far, models were calculated for the cosmic ray flux of the solar neighbourhood. A flux enhancement by two to three orders of magnitude (with all necessary chemical implications) has still to be incorporated into PDR codes (but see Schilke et al. 1993). Another important quantity is the spatial distribution of the UV flux. For M 82, we do not know the variation of the UV flux as a function of galactocentric radius.
## Appendix A The radiative transfer model
Systematic motions and microturbulence are frequently used as simplifying assumptions to facilitate treatment of line formation in molecular clouds. For the large amplitude systematic motions assumed in the Large Velocity Gradient (LVG) approximation, the source functions are locally defined. For the homogeneous velocity field assumed in the microturbulent approximation, the source functions are generally coupled throughout the cloud by scattered radiation. The crucial difference between the two approximations is that between local and non-local excitation and both can be viewed as limiting cases for the treatment of molecular line formation. In the case of CO, line intensities determined with a standard LVG model (e.g. Castor 1970; Scoville & Solomon 1974) do not differ by more than a factor of three from models using the microturbulent approach (White 1977; see also Ossenkopf 1997). This is within the uncertainties that can be attributed to cloud geometry. Furthermore, Wild et al. (1992) found no significant difference between results from their ‘clumpy cloud’ and standard LVG models. Given the inhomogeneity of the ISM, the assumption of uniform physical conditions is crude and only yields average gas properties. In view of the limited quality of the available data and the lack of information on source morphology and its fine-scale structure, however, an LVG treatment of the radiative transfer is an appropriate first step to analyse M 82.
Applying our LVG model to simulate a cloud of spherical geometry, CO collisional cross sections were taken from Green & Chapman (1978). For $`T_{\mathrm{kin}}`$ $`>`$ 100 K, we assumed collision rates $`C_{\mathrm{ij}}`$ $``$ $`T_{\mathrm{kin}}^{1/2}`$. Using instead the collision rates recommended by De Jong et al. (1975; these include higher temperatures) leads to similar results. A comparison of $`J`$=2–1/$`J`$=1–0 line intensity ratios with those of Castets et al. (1990), that were computed with yet another set of collision rates, also shows consistency within 20%.
Fig. 10 displays our calculated CO 7–6/4–3 and CO 4–3/<sup>13</sup>CO 3–2 line intensity ratios as a function of density (10<sup>2-7</sup>$`\text{cm}^3`$), kinetic temperature (5–150 K), <sup>12</sup>CO/<sup>13</sup>CO abundance ratio (25–75), and CO fractional abundance in terms of $`X`$(CO)/grad $`V`$ (10<sup>-6…-3</sup> pc/$`\mathrm{km}\mathrm{s}^1`$; $`X`$(CO) is the fractional abundance parameter, grad $`V`$ denotes the velocity gradient in $`\mathrm{km}\mathrm{s}^1`$/pc). The choice of $`X`$(CO)/grad $`V`$ is motivated by a source size of $``$ 300 pc and a velocity range of $``$300 $`\mathrm{km}\mathrm{s}^1`$ (this leads to grad $`V`$ $``$ 1 $`\mathrm{km}\mathrm{s}^1`$ pc<sup>-1</sup>), a solar system \[C\]/\[H\] abundance ratio of 3.5 10<sup>-4</sup> (Grevesse et al. 1994), and the assumption that a significant fraction of the available carbon ($``$10%) is forming CO molecules. Since molecular clouds with significant CO 7–6 emission must be warm, strong CO fractionation in favor of enhanced <sup>13</sup>C abundances (cf. Watson et al. 1976) can be excluded. Fig. 10 thus presents calculations for <sup>12</sup>CO/<sup>13</sup>CO = 25 (the carbon isotope ratio in the galactic center region; e.g. Wilson & Rood 1994), <sup>12</sup>CO/<sup>13</sup>CO = 50 (the <sup>12</sup>C/<sup>13</sup>C ratio in the inner galactic disk; e.g. Henkel et al. 1985), and <sup>12</sup>CO/<sup>13</sup>CO = 75 (close to the carbon isotope ratio of the local ISM and the solar system; e.g. Stahl et al. 1989; Stahl & Wilson 1992).
## Appendix B Cloud stability
In order for a cloud to be gravitationally bound, it must be sufficiently dense to withstand the tidal stresses caused by the gravitational potential of the galaxy. Neglecting rotation, turbulence, and magnetic fields, and compensating tidal forces by a cloud’s own gravity, we can derive a minimum density that is required for survival in a hostile medium. With $`V_{\mathrm{rot}}`$ $``$ 140 $`\mathrm{km}\mathrm{s}^1`$ at a galactocentric radius of $`R`$ $``$ 75 pc from the dynamical center of M 82 (Neininger et al. 1998) and applying Eq. 5 of Güsten and Downes (1980), we find a limiting cloud density of $`n_{\mathrm{min}}`$ $``$ 10<sup>3.2-3.5</sup> \[$`R`$/120 pc\]<sup>-2</sup>$`\text{cm}^3`$. $`R`$ = 120 pc refers to the molecular lobes of M 82, located at offsets $`\pm `$7$`.^{\prime \prime }`$5 from the dynamical center of the galaxy (see Table 2). The minimum density $`n_{\mathrm{min}}`$ is near the lower limit of the densities deduced in Sect. 5.1 with an LVG model and just matches densities derived in Sect. 6.3.2 with a PDR model.
## Appendix C Radiative pumping
Because the kinetic temperature is constrained by the CO $`J`$=7–6/$`J`$=4–3 line temperature ratios (Table 1), collisional excitation to excited vibrational or electronic levels, $``$ 3000 K above the ground state, is not effective. Densities and temperatures obtained with Fig. 10 might however be affected by radiative excitation, that is neglected in our LVG treatment (apart of the 2.7 K microwave background).
To operate efficiently, radiative pumping must be faster than the collisional rates for rotational excitation in the $`v`$ = 0 vibrational ground state. With $`<`$$`\sigma v`$$`>`$ $``$ 2 10<sup>-11</sup> cm<sup>3</sup> s<sup>-1</sup> (Green & Chapman 1978) and a density of 5 10<sup>3</sup>$`\text{cm}^3`$, collision rates are at the order of 10<sup>-7</sup> s<sup>-1</sup>. Direct rotational excitation by $`\lambda `$ $``$ 1 mm photons from the dust is inefficient, because $`\tau _{\mathrm{dust}}`$ $``$ 1. For vibrational exitation by 4.7$`\mu `$m photons we obtain with the Einstein coefficients $`A_{\mathrm{rot}}`$ = 6.1 10<sup>-6</sup> and 3.4 10<sup>-5</sup> s<sup>-1</sup> for the $`J`$ = 4–3 and 7–6 rotational transitions, $`A_{\mathrm{v}=10}`$ $``$ 20 s<sup>-1</sup> (Kirby-Docken & Liu 1978), and the condition $`T_{\mathrm{dust}}`$ $``$ 3070 K/ln \[$`A_{\mathrm{v}=10}`$/$`A_{\mathrm{rot}}`$\] (see Eq. 6 of Carroll & Goldsmith 1981; $`\nu _{\mathrm{v}=10}`$ $``$ 6.4 10<sup>13</sup> Hz) a minimum dust temperature of $`T_{\mathrm{dust}}`$ $``$ 200 K. This is much larger than $`T_{\mathrm{dust}}`$ = 48 K, estimated by Hughes et al. (1994) and Colbert et al. (1999). The condition for efficient pumping by infrared photons is therefore, if at all, only fulfilled for a small part of the molecular complex in the central region of M 82 (cf. McLeod et al. 1993). Adopting in spite of this $`T_{\mathrm{dust}}`$ $``$ 200 K, the radiative pumping rate becomes $`B_{\mathrm{v}=01}`$ $`u_{\mathrm{v}=10}`$ $``$ 10<sup>-15</sup> s<sup>-1</sup> (or less, if beam dilution plays a role). This is negligible in comparison with the collisional pumping rate.
###### Acknowledgements.
It is a pleasure to thank the HHT staff, in particular H. Butner, B. Hayward, D. Muders, F. Patt, and B. Stupak, for their enthusiastic support of the project and for their flexibility in changing schedules according to variable weather conditions. For the permission to use the HEB, we also thank the Harvard-Smithonian Center for Astrophysics (CfA). We acknowledge useful discussions with S. Hüttemeister, E. Ros, P. Schilke, C.M. Walmsley, A. Weiß, and the useful comments of an anonymous referee. R.Q.M. acknowledges support by the exchange program between the Chinese Academy of Sciences and the Max-Planck-Gesellschaft; C.H. acknowledges support from NATO grant SA.5-2-05 (CRG. 960086) 318/96.
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# Windings of the 2D free Rouse chain
## 1 Introduction
In order to study the dynamics of dilute polymer solutions, P.E. Rouse proposed in 1953 his famous model of harmonically bound Brownian particles (Rouse chain) . Since that time, this model has become very popular in the field of polymer science. It appears that, despite its drawbacks and limitations (in particular, absence of excluded volume and hydrodynamic interactions), it is conceptually important and useful to study the dynamics of polymers in melts . In this paper, we will consider the free planar motion of such a chain of $`n`$ particles (monomers) and especially address its long time ($`t\mathrm{}`$) properties from the Brownian motion viewpoint.
A configuration of this chain being represented by a complex $`n`$-vector $`z`$ (the components $`z_i,`$ $`i=1,\mathrm{},n`$ are the complex coordinates of the particles), we will study closed trajectories of length $`t`$, i.e. $`z(t)=z(0)`$ or open ones ($`z(0)`$ fixed, $`z(t)`$ left unspecified, i.e. integrated over).
More precisely, if we consider some given bounded domain $`S`$ of area $`𝒮`$ and define the occupation time $`T_j`$ as the time spent inside $`S`$ by the $`j^{\mathrm{th}}`$ particle, our goal is to compute the joint probability distribution $`P(T_1,T_2,\mathrm{},T_n)(P(\{T_j\}))`$. Similarly, $`A_j`$ and $`\theta _j`$ being respectively the area enclosed by the trajectory of the $`j^{\mathrm{th}}`$ particle and its winding angle around O, we will be interested in the joint laws $`P(\{A_j\})`$ and $`P(\{\theta _j\})`$.
On general grounds, we expect that the various properties of the chain will be strongly influenced by the free Brownian motion of the center of mass (c.o.m.) of the chain. But will they strictly satisfy the same laws ? With the same scaling variables ? And what about the correlations among the different variables ? Before answering those questions, we first recall some standard results concerning a planar Brownian particle with a diffusion constant $`D`$ \[4-8\].
Results i) and ii) concern open trajectories when $`t\mathrm{}`$ while iii) concerns closed trajectories and is valid for all times $`t`$:
1. Kallianpur-Robbins’ law for the probability distribution of the occupation time $`T`$ of a bounded domain of area $`𝒮`$:
$$P\left(T^{}=\frac{4\pi DT}{𝒮\mathrm{ln}t}\right)=\theta (T^{})e^T^{}$$
(1)
2. Spitzer’s law for the angle $`\theta `$ wound around a given point:
$$P\left(\theta ^{}=\frac{2\theta }{\mathrm{ln}t}\right)=\frac{1}{\pi }\frac{1}{1+(\theta ^{})^2}$$
(2)
with the characteristic function:
$$e^{i\lambda \theta ^{}}=e^{|\lambda |}$$
(3)
3. Lévy’s law for the area $`A`$ enclosed by the closed trajectory of the particle:
$$P\left(A^{}=\frac{A}{2Dt}\right)=\frac{\pi }{2}\frac{1}{\mathrm{cosh}^2(\pi A^{})}$$
(4)
$$e^{iBA^{}}=\frac{\left(\frac{B}{2}\right)}{\mathrm{sinh}\left(\frac{B}{2}\right)}$$
(5)
The distributions i) and iii) have moments of all orders in contrast with ii) that has none.
Those laws were discovered more than 40 years ago and since that time, many refinements have been made. For instance, in , the authors found the asymptotic ($`t\mathrm{}`$) joint law of the small ($`\theta _{}`$) and big ($`\theta _+`$) windings. $`\theta _{}`$( resp. $`\theta _+`$) are the angles wound around O and only counted when $`r`$ is smaller (resp. greater) than some fixed $`r_0`$ ($`r`$ is the distance separating the particle from O). With the rescaled angles $`\theta _\pm ^{}=\frac{2\theta _\pm }{\mathrm{ln}t}`$, the characteristic function writes :
$$e^{i(\lambda _{}\theta _{}^{}+\lambda _+\theta _+^{})}=\frac{1}{\mathrm{cosh}(\lambda _+)+\frac{|\lambda _{}|}{\lambda _+}\mathrm{sinh}(\lambda _+)}$$
(6)
($`\lambda _+=\lambda _{}=\lambda `$ gives back Spitzer’s law (2)).
Remark that (2) and (6) don’t depend on the diffusion constant. This is quite different from the Brownian motion on a bounded domain surrounding O. In that case, we have :
$$e^{i\lambda \theta }=e^{cD|\lambda |t}$$
(7)
where $`c`$ is a constant depending on the geometry and the boundary conditions. Here, $`D`$, appears as a multiplicator of $`|\lambda |`$. We will use this remark at the end of the paper.
## 2 The free Rouse chain
Starting our study, we consider the following set of coupled Langevin equations:
$`\dot{z}_1`$ $`=`$ $`k(z_2z_1)+\eta _1`$
$`\dot{z}_l`$ $`=`$ $`k(z_{l+1}+z_{l1}2z_l)+\eta _l,2ln1`$ (8)
$`\dot{z}_n`$ $`=`$ $`k(z_{n1}z_n)+\eta _n`$
where $`k`$ is the spring constant and $`\eta _m`$ ($`\eta _{mx}+i\eta _{my}`$) a gaussian white noise:
$`\eta _m(t)`$ $`=`$ $`0`$
$`\eta _m(t)\eta _m^{}(t^{})`$ $`=`$ $`2\delta _{mm^{}}\delta (tt^{})`$ (9)
(This noise would correspond to a $`D=1/2`$ diffusion constant if particles were free).
For the chain c.o.m., we get $`\dot{z}_G=\frac{1}{n}\left(_{i=1}^n\eta _i\right)\eta _G`$ with $`\eta _G(t)\eta _G(t^{})=\frac{2}{n}\delta (tt^{})`$. The c.o.m. motion is free with $`D=1/(2n)`$.
Introducing the complex $`n`$-vector $`\eta `$, eq.(8) can be written in a matrix form:
$$\dot{z}=k𝐌z+\eta $$
(10)
where $`𝐌`$ is the tridiagonal $`(n\times n)`$ matrix:
$$𝐌=\left(\begin{array}{ccccc}1& 1& 0& \mathrm{}& 0\\ 1& 2& 1& \mathrm{}& 0\\ 0& 1& 2& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& \mathrm{}& 1\end{array}\right)$$
with eigenvalues:
$$\omega _j=\mathrm{\hspace{0.17em}2}\left(1\mathrm{cos}\frac{\pi (j1)}{n}\right),1jn$$
(11)
($`\omega _1=0`$; $`det^{}𝐌_{j=2}^n\omega _j=n`$)
With the matrix $`\omega =\mathrm{diag}(\omega _i)`$, we can write:
$`\omega `$ $`=`$ $`𝐑^1𝐌𝐑`$ (12)
$`z`$ $`=`$ $`𝐑Z`$ (13)
where $`𝐑`$ is an orthogonal matrix and the components of $`Z`$ are the normal coordinates, that we will widely use in the sequel. From $`𝐑_{j1}=\frac{1}{\sqrt{n}}`$, $`j=1,\mathrm{},n`$, we deduce that $`Z_1(=_{i=1}^nz_i/\sqrt{n})`$ is essentially the c.o.m. coordinate. Remark also that $`_{i=2}^n\omega _i|Z_i|^2=_{i=2}^n|z_iz_{i1}|^2=^t\overline{z}𝐌z`$.
Let us call $`𝒫(z,z^{(0)},t)`$ the probability for the chain to go from configuration $`z^{(0)}`$ at $`t=0`$ to $`z`$ at time $`t`$. $`𝒫`$ satisfies a Fokker-Planck equation :
$$_t𝒫=\left({}_{}{}^{t}_{z}^{}k𝐌z+^t_{\overline{z}}k𝐌\overline{z}+2^t_{\overline{z}}_z\right)𝒫$$
(14)
where $`_z`$ (resp. $`_{\overline{z}}`$) is a $`n`$-vector of components $`_{z_i}`$ (resp. $`_{\overline{z}_i}`$) and $`{}_{}{}^{t}_{z}^{}`$ (resp. $`{}_{}{}^{t}_{\overline{z}}^{}`$) is the transpose of $`_z`$ (resp. $`_{\overline{z}}`$). The solution can be written in terms of a path integral ($`𝒟z𝒟\overline{z}=_{i=1}^n𝒟z_i𝒟\overline{z_i}`$):
$`𝒫(z,z^{(0)},t)`$ $`=`$ $`det\left(e^{tk𝐌}\right){\displaystyle _{z^{(0)}}^z}𝒟z𝒟\overline{z}\mathrm{exp}({\displaystyle \frac{1}{2}}{\displaystyle _0^t}{}_{}{}^{t}(\dot{\overline{z}}+k𝐌\overline{z})(\dot{z}+k𝐌z)\mathrm{d}\tau )`$
$``$ $`F(z,z^{(0)}).G_0(z,z^{(0)},t)`$
with
$`F(z,z^{(0)})`$ $`=`$ $`e^{\frac{k}{2}\left({}_{}{}^{t}\overline{z}𝐌z^t\overline{z}^{(0)}𝐌z^{(0)}\right)}=`$
$`=`$ $`e^{\frac{k}{2}_{i=2}^n\left(|z_iz_{i1}|^2|z_i^{(0)}z_{i1}^{(0)}|^2\right)}=e^{\frac{k}{2}_{i=2}^n\omega _i\left(|Z_i|^2|Z_i^{(0)}|^2\right)}`$
$`G_0(z,z^{(0)},t)`$ $`=`$ $`{\displaystyle _{z^{(0)}}^z}𝒟z𝒟\overline{z}\mathrm{exp}\left({\displaystyle \frac{1}{2}}{\displaystyle _0^t}\left({}_{}{}^{t}\dot{\overline{z}}\dot{z}+k^2{}_{}{}^{t}\overline{z}𝐌^2z2k\mathrm{Tr}𝐌\right)d\tau \right)=`$ (16)
$`=`$ $`z\left|e^{tH_0}\right|z^{(0)}`$
$`H_0`$ $`=`$ $`2^t_{\overline{z}}_z+{\displaystyle \frac{1}{2}}k^2{}_{}{}^{t}\overline{z}𝐌^2zk\mathrm{Tr}𝐌`$ (17)
In fact, $`𝒫`$, eq.(2), can be easily deduced from the gaussian distribution of $`\eta `$ (use (10); det($`e^{tk𝐌}`$) is the functional Jacobian for the change of variable $`\eta z`$ ).
$`G_0(z,z^{(0)},t)`$ is most conveniently written in terms of the normal coordinates $`Z_i`$ and $`Z_i^{(0)}`$, clearly exhibiting the free motion of the c.o.m. :
$`G_0(z,z^{(0)},t)={\displaystyle \frac{1}{2\pi t}}e^{\frac{1}{2t}|Z_1Z_1^{(0)}|^2}\times `$
$`\times {\displaystyle \underset{i=2}{\overset{n}{}}}\left({\displaystyle \frac{s_ie^{k\omega _it}}{2\pi }}\mathrm{exp}\{{\displaystyle \frac{1}{2}}(\overline{Z}_ic_iZ_i+\overline{Z}_i^{(0)}c_iZ_i^{(0)}\overline{Z}_i^{(0)}s_iZ_i\overline{Z}_is_iZ_i^{(0)})\}\right)`$ (18)
$`s_i={\displaystyle \frac{k\omega _i}{\mathrm{sinh}(k\omega _it)}};c_i=k\omega _i\mathrm{coth}(k\omega _it)`$
When $`kt1`$, we get, for $`G_0`$, the limiting expression:
$`G_0^{\mathrm{}}(z,z^{(0)},t)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi t}}e^{\frac{1}{2t}|Z_1Z_1^{(0)}|^2}{\displaystyle \underset{i=2}{\overset{n}{}}}\left({\displaystyle \frac{k\omega _i}{\pi }}e^{\frac{k\omega _i}{2}(|Z_i|^2+|Z_i^{(0)}|^2)}\right)`$
$``$ $`𝒢_0(z,z^{(0)},t).g_0(z,z^{(0)})`$
where $`𝒢_0`$ is the c.o.m. propagator and $`g_0`$ can be simply written in terms of the $`z_i`$:
$$g_0(z,z^{(0)})=n\left(\frac{k}{\pi }\right)^{n1}e^{\frac{k}{2}_{i=2}^n\left(|z_iz_{i1}|^2+|z_i^{(0)}z_{i1}^{(0)}|^2\right)}$$
(20)
Furthermore, as can be easily checked, $`𝒫`$ is properly normalized:
$`dzd\overline{z}𝒫(z,z^{(0)},t)=1`$ ($`G_0`$ given by (18) or (2)).
Now, we turn to the computation of the joint law $`P(\{T_j\})`$.
## 3 Occupation times distribution
Recall that $`T_j`$ is the time spent by particle $`j`$ inside a bounded domain $`S`$ of area $`𝒮`$. We consider trajectories starting at $`t=0`$ from some given configuration $`z^{(0)}`$ and reaching at time $`t`$ the final configuration $`z`$. Leaving $`z`$ unspecified, we have, with positive $`p_i`$ ’s:
$`e^{_{i=1}^np_iT_i}=det\left(e^{tk𝐌}\right)\times `$
$`\times {\displaystyle }\mathrm{d}z\mathrm{d}\overline{z}{\displaystyle _{z^{(0)}}^z}𝒟z𝒟\overline{z}\mathrm{exp}({\displaystyle _0^t}({\displaystyle \frac{1}{2}}^t(\dot{\overline{z}}+k𝐌\overline{z})(\dot{z}+k𝐌z)+V_P(z))\mathrm{d}\tau )`$ (21)
$`={\displaystyle dzd\overline{z}F(z,z^{(0)})G_P(z,z^{(0)},t)}`$ (22)
with
$$G_P(z,z^{(0)},t)=z\left|e^{t(H_0+V_P)}\right|z^{(0)}$$
(23)
$$V_P(z)=\underset{i=1}{\overset{n}{}}p_i\mathrm{𝟏}_S(z_i)$$
(24)
$`\mathrm{𝟏}_S(z_i)`$ is the indicatrix function of the domain $`S`$. Symbolically, we write:
$$G_P=\underset{m=0}{\overset{\mathrm{}}{}}(1)^mG_0(V_PG_0)^m$$
(25)
with
$`G_0(V_PG_0)^m={\displaystyle _0^t}\mathrm{d}t_m{\displaystyle _0^{t_m}}\mathrm{d}t_{m1}\mathrm{}{\displaystyle _0^{t_2}}\mathrm{d}t_1{\displaystyle }\left({\displaystyle \underset{j=1}{\overset{m}{}}}\mathrm{d}\overline{z}^{(j)}\mathrm{d}z^{(j)}\right)G_0(z,z^{(m)},tt_m)\times `$
$`\times V_P(z^{(m)})G_0(z^{(m)},z^{(m1)},t_mt_{m1})V_P(z^{(m1)})\mathrm{}V_P(z^{(1)})G_0(z^{(1)},z^{(0)},t_1)`$ (26)
($`z^{(j)}`$ is the chain configuration at time $`t_j`$; $`\mathrm{d}\overline{z}^{(j)}\mathrm{d}z^{(j)}=_{i=1}^n\mathrm{d}\overline{z}_i^{(j)}\mathrm{d}z_i^{(j)}`$).
Let us compute the contribution $`N_m(t)`$ of this generic term to (22). Integrating over $`z`$, we have:
$`N_m(t)`$ $`=`$ $`(1)^m{\displaystyle _0^t}\mathrm{d}t_m{\displaystyle _0^{t_m}}\mathrm{d}t_{m1}\mathrm{}{\displaystyle _0^{t_2}}\mathrm{d}t_1{\displaystyle }\left({\displaystyle \underset{j=1}{\overset{m}{}}}\mathrm{d}\overline{z}^{(j)}\mathrm{d}z^{(j)}\right)F(z^{(m)},z^{(0)})\times `$ (27)
$`\times V_P(z^{(m)})G_0(z^{(m)},z^{(m1)},t_mt_{m1})\mathrm{}V_P(z^{(1)})G_0(z^{(1)},z^{(0)},t_1)`$
$``$ $`(1)^m{\displaystyle _0^t}dt_m{\displaystyle \left(\underset{j=1}{\overset{m}{}}\mathrm{d}\overline{z}^{(j)}\mathrm{d}z^{(j)}V_P(z^{(j)})\right)F(z^{(m)},z^{(0)})\mathrm{\Phi }(t_m,\{z^{(l)}\})}`$ (28)
$`\mathrm{\Phi }`$ is a time convolution product of the free propagators $`G_0`$. Disregarding for the moment the spatial integrations, the above expression is well-suited, in the limit $`t\mathrm{}`$, for applying Tauberian theorems . Introducing the Laplace Transform $`\widehat{\mathrm{\Phi }}`$:
$$\widehat{\mathrm{\Phi }}(u,\{z^{(l)}\})_0^{\mathrm{}}e^{ut^{}}\mathrm{\Phi }(t^{},\{z^{(l)}\})dt^{}=\underset{k=1}{\overset{m}{}}\widehat{G_0}(z^{(k)},z^{(k1)},u)$$
(29)
we notice that, when $`u0^+`$:
$$\widehat{G_0}(z^{(k)},z^{(k1)},u)_0^{\mathrm{}}e^{ut^{}}G_0(z^{(k)},z^{(k1)},t^{})dt^{}_a^{\mathrm{}}e^{ut^{}}G_0(z^{(k)},z^{(k1)},t^{})dt^{}$$
(30)
for some large $`a`$. So, we can use the asymptotic form $`G_0^{\mathrm{}}`$ in the computation of $`\widehat{G_0}`$ and get:
$$\widehat{G_0}(z^{(k)},z^{(k1)},u)_{u0^+}\mathrm{ln}\left(\frac{1}{u}\right)\frac{1}{2\pi }g_0(z^{(k)},z^{(k1)})$$
(31)
A weak Tauberian theorem gives for the time integration in (28):
$$_0^tdt_m\mathrm{\Phi }(t_m,\{z^{(l)}\})_t\mathrm{}\left(\frac{\mathrm{ln}t}{2\pi }\right)^m\underset{k=1}{\overset{m}{}}g_0(z^{(k)},z^{(k1)})$$
(32)
Finally, the result for $`N_m(t)`$ is:
$$N_m(t)_t\mathrm{}(1)^m\left(\frac{nL_P\mathrm{ln}t}{2\pi }\right)^m$$
(33)
$$L_P=\left(\underset{i=1}{\overset{n}{}}\mathrm{d}z_i\mathrm{d}\overline{z_i}\right)V_P(z)\left(\frac{k}{\pi }\right)^{n1}e^{k_{i=2}^n|z_iz_{i1}|^2}$$
(34)
$`L_P`$ is computed with $`V_P`$, eq.(24): $`L_P=\left(_{i=1}^np_i\right)𝒮`$
Rescaling the occupation times $`T_i`$:
$$T_i^{}=\frac{2\pi T_i}{n𝒮\mathrm{ln}t}$$
(35)
we get:
$$e^{_{i=1}^np_iT_i^{}}=\underset{m=0}{\overset{\mathrm{}}{}}(1)^m\left(\underset{i=1}{\overset{n}{}}p_i\right)^m=\frac{1}{1+\left(_{i=1}^np_i\right)}$$
(36)
This relationship is actually valid, by analytic continuation, for all the positive $`p_i`$ ’s (and not only when $`p_i<1`$). This is because the distribution $`P(\{T_j\})`$ has moments of all orders and consequently $`e^{_{i=1}^np_iT_i}`$ is holomorphic when $`\mathrm{Re}(p_i)0`$.
(36) leads to the probability distribution:
$$P(\{T_i^{}\})=\theta (T_1^{})e^{T_1^{}}\underset{i=2}{\overset{n}{}}\delta (T_i^{}T_{i1}^{})$$
(37)
($`n=1`$ gives back the Kallianpur-Robbins’ law).
So, in the large time limit, the $`T_i`$ ’s are strongly correlated leading to identical $`(T_i^{})`$ ’s. Moreover, we remark that $`T_i`$ scales like $`n`$ and, also, that the law for the c.o.m. would be the same as for one monomer (compare (35) to – (1) with $`D=1/(2n)`$): the c.o.m. free motion dominates this process.
We also got similar exponential distributions for the rescaled variables $`T^{}`$ in the following cases:
1. $`T`$ is the time spent when the whole chain is inside $`S`$. $`L_P`$, eq.(34), is now computed with $`V_P(z)=p\left(_{i=1}^n\mathrm{𝟏}_S(z_i)\right)`$. Introducing
$$w(𝒮)=\left(\underset{i=1}{\overset{n}{}}\mathrm{d}z_i\mathrm{d}\overline{z_i}\mathrm{𝟏}_S(z_i)\right)e^{k_{i=2}^n|z_iz_{i1}|^2}$$
(38)
the rescaled variable writes:
$$T^{}=\left(\frac{\pi }{k}\right)^{n1}\frac{2\pi T}{nw(𝒮)\mathrm{ln}t}$$
(39)
In contrast with (35), $`k`$ is now present in the asymptotic law. For instance, if $`S`$ is a small disk of radius $`r_0`$ ($`kr_0^21`$), then $`w(𝒮)𝒮^n`$ and $`T`$ scales like $`k^{n1}`$: when $`k`$ grows, the chain collapses and it is easier to confine it inside a given domain.
2. $`T`$ is the time spent by the chain when at least one of its particles is inside $`S`$. $`T^{}`$ is similar to (39) except that $`w(𝒮)`$ must be changed: $`L_P`$ is now computed with $`V_P(z)=p\left(1_{i=1}^n(1\mathrm{𝟏}_S(z_i))\right)`$.
To conclude this section, let us draw two lessons:
1. The scaling variables take the same general form as for the free Brownian particle. In the sequel, we will show that it is still true for the other quantities we study.
2. For the computation of the perturbation theory, when $`t\mathrm{}`$, we can systematically use the asymptotic form $`G_0^{\mathrm{}}`$ of the unperturbed propagator. Obviously, for this consideration to hold, we must be sure that the perturbation series is well behaved. In those conditions, we will make a wide use of this remark.
## 4 Areas distribution
Now, we compute the areas distribution $`P(\{A_j\})`$ for closed trajectories of length $`t`$ starting and ending at some fixed $`z^{(0)}`$. To do so, we insert the constraint:
$$\underset{j=1}{\overset{n}{}}\delta \left(A_j\frac{1}{4i}_0^t(z_j\dot{\overline{z}_j}\overline{z}_j\dot{z_j})d\tau \right)$$
(40)
in the measure (2) and use the relationship $`\delta (x)=\frac{1}{2\pi }e^{iBx}dB`$. It is easy to show that this manipulation amounts to add $`n`$ different magnetic fields $`B_j`$ to the initial system. Those fields are uniform, orthogonal to the motion plane and such that particle $`j`$ is submitted to $`B_j`$.
With the $`(n\times n)`$ diagonal matrix $`𝐁`$ ($`𝐁_{ij}=B_i\delta _{ij}`$), we get
$`P(\{A_i\})`$ $`=`$ $`{\displaystyle \left(\underset{j=1}{\overset{n}{}}\frac{\mathrm{d}B_j}{2\pi }e^{iB_jA_j}\right)\left(\frac{G_𝐁(z^{(0)},z^{(0)},t)}{G_0(z^{(0)},z^{(0)},t)}\right)}`$ (41)
$`\mathrm{with}G_𝐁(z^{(0)},z^{(0)},t)`$ $`=`$ $`z^{(0)}\left|e^{tH_𝐁}\right|z^{(0)}`$ (42)
$`H_𝐁`$ $`=`$ $`H_0+V_𝐁`$ (43)
$`V_𝐁(z)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(^tz𝐁_z+^t\overline{z}𝐁_{\overline{z}}\right)+{\displaystyle \frac{1}{8}}^t\overline{z}𝐁^2z`$ (44)
$`=`$ $`{\displaystyle \frac{1}{2}}\left(^tZ𝐁^{}_Z+^t\overline{Z}𝐁^{}_{\overline{Z}}\right)+{\displaystyle \frac{1}{8}}^t\overline{Z}𝐁^{\prime \prime }Z`$ (45)
($`𝐁^{}=𝐑^1\mathrm{𝐁𝐑},𝐁^{\prime \prime }=𝐑^1𝐁^2𝐑`$). In principle, $`P(\{A_j\})`$ depends on $`z^{(0)}`$ but we will show that, actually, this is not the case when $`t\mathrm{}`$ (remark that, for all $`t`$, $`P(\{A_j\})`$ doesn’t depend on the c.o.m. of $`z^{(0)}`$: this is due to translation invariance and this is the reason why we consider the propagator and not the partition function that would diverge like the area of the plane, leading to serious problems in the perturbation theory).
Now, let us sketch the perturbative computation of $`G_𝐁`$. Following our previous remarks, we will use $`G_0^{\mathrm{}}`$ for the unperturbed propagator. The generic term writes:
$`(1)^m{\displaystyle _0^t}dt_m{\displaystyle _0^{t_m}}dt_{m1}\mathrm{}{\displaystyle _0^{t_2}}dt_1{\displaystyle \left(\underset{j=1}{\overset{m}{}}\mathrm{d}\overline{z}^{(j)}\mathrm{d}z^{(j)}\right)\mathrm{}}`$
$`\mathrm{}G_0^{\mathrm{}}(z^{(j+1)},z^{(j)},t_{j+1}t_j)V_𝐁(z^{(j)})G_0^{\mathrm{}}(z^{(j)},z^{(j1)},t_jt_{j1})\mathrm{}`$ (46)
With normal coordinates, $`V_𝐁(z^{(j)})`$ takes the form:
$$V_𝐁(z^{(j)})=\underset{i,l=1}{\overset{n}{}}\left(\frac{1}{2}\left(Z_i^{(j)}𝐁_{il}^{}_{Z_l^{(j)}}+\overline{Z}_i^{(j)}𝐁_{il}^{}_{\overline{Z}_l^{(j)}}\right)+\frac{1}{8}\overline{Z}_i^{(j)}𝐁_{il}^{\prime \prime }Z_l^{(j)}\right)$$
(47)
We now proceed by inspection:
1. For the terms $`Z_i^{(j)}𝐁_{il}^{}_{Z_l^{(j)}}`$, only will survive, after integration, those contributions with $`i=j`$. The same holds for $`\overline{Z}_i^{(j)}𝐁_{il}^{}_{\overline{Z}_l^{(j)}}`$. Moreover, the diagonal contributions exactly cancel except when $`i=j=1`$.
2. We reach the same conclusion for the terms $`Z_i^{(j)}𝐁_{il}^{\prime \prime }Z_l^{(j)}`$ (non diagonal terms vanish after integration. Diagonal contributions, $`i=j>1`$, are subleading – compared to $`i=j=1`$ – when $`kt1`$: we recover the fact that the process is dominated by the c.o.m. motion).
We are finally left with an effective perturbation $`V_𝐁^{\mathrm{eff}}`$:
$$V_𝐁^{\mathrm{eff}}=\frac{1}{2}\left(Z_1𝐁_{11}^{}_{Z_1}+\overline{Z}_1𝐁_{11}^{}_{\overline{Z}_1}\right)+\frac{1}{8}𝐁_{11}^{\prime \prime }|Z_1|^2$$
(48)
Only the first mode is affected by the magnetic fields and we can disregard the other modes that will cancel when taking the ratio $`G_𝐁/G_0`$.
Remark that
$$𝐁_{11}^{}=\frac{1}{n}\left(\underset{i=1}{\overset{n}{}}B_i\right)\mathrm{and}𝐁_{11}^{\prime \prime }=\frac{1}{n}\left(\underset{i=1}{\overset{n}{}}B_i^2\right)$$
(49)
The effective hamiltonian for the remaining mode writes:
$$H_𝐁^{\mathrm{eff}}=2_{Z_1}_{\overline{Z}_1}+\frac{1}{2}\left(Z_1𝐁_{11}^{}_{Z_1}+\overline{Z}_1𝐁_{11}^{}_{\overline{Z}_1}+\frac{1}{4}(𝐁_{11}^{})^2|Z_1|^2\right)+\frac{1}{8}\left(𝐁_{11}^{\prime \prime }(𝐁_{11}^{})^2\right)|Z_1|^2$$
(50)
It describes the behavior of a charged particle submitted to an uniform magnetic field $`𝐁_{11}^{}`$ and an harmonic oscillator of frequency $`\omega =\frac{1}{2}\sqrt{𝐁_{11}^{\prime \prime }(𝐁_{11}^{})^2}`$. Using known results about this problem , we immediately get:
$`{\displaystyle \frac{G_𝐁(z^{(0)},z^{(0)},t)}{G_0(z^{(0)},z^{(0)},t)}}`$ $`=`$ $`{\displaystyle \frac{t\sqrt{𝐁_{11}^{\prime \prime }}}{2\mathrm{sinh}\left(\frac{t}{2}\sqrt{𝐁_{11}^{\prime \prime }}\right)}}\times `$ (51)
$`\times \mathrm{exp}\left({\displaystyle \frac{\sqrt{𝐁_{11}^{\prime \prime }}\left(\mathrm{cosh}\left(\frac{t}{2}\sqrt{𝐁_{11}^{\prime \prime }}\right)\mathrm{cosh}\left(\frac{t}{2}𝐁_{11}^{}\right)\right)|Z_1^{(0)}|^2}{2\mathrm{sinh}\left(\frac{t}{2}\sqrt{𝐁_{11}^{\prime \prime }}\right)}}\right)`$
However, for our computation to be consistent, we must consider this expression in the large time limit. This is readily done in rescaling the areas $`A_i^{}=A_i/t`$ and doing $`t\mathrm{}`$. The final expression for the characteristic function of $`P(\{A_i^{}\})`$ is quite simple:
$$e^{i_{j=1}^nB_jA_j^{}}=\frac{\sqrt{𝐁_{11}^{\prime \prime }}}{2\mathrm{sinh}\left(\frac{\sqrt{𝐁_{11}^{\prime \prime }}}{2}\right)}$$
(52)
($`𝐁_{11}^{\prime \prime }=\frac{1}{n}_{i=1}^nB_i^2`$ ; when $`n=1`$, (52) gives back Lévy’s result, eq.(5)).
Owing to the form of $`𝐁_{11}^{\prime \prime }`$, $`P(\{A_i^{}\})`$ will only be a function of the variable
$`\sqrt{_{i=1}^n(A_i^{})^2}(A^{})`$, showing clearly that the ($`A_i^{}`$) ’s are correlated. Its determination is reduced to the computation of the following integral :
$$P(\{A_i^{}\})P(A^{})=\left(\frac{2n}{\pi }\right)^{n/2}\frac{1}{(A^{})^{n/21}}_0^{\mathrm{}}J_{n/21}(A^{}r)\frac{r^{n/2+1}}{\mathrm{sinh}(r)}dr$$
(53)
where $`J_\nu `$ is a Bessel function. Closed form expressions can be given for odd $`n`$ values. For instance:
$`n=3P(A^{})`$ $`=`$ $`{\displaystyle \frac{3\pi }{2A^{}}}{\displaystyle \frac{\mathrm{tanh}(\pi \sqrt{3}A^{})}{\mathrm{cosh}^2(\pi \sqrt{3}A^{})}}`$ (54)
$`n=5P(A^{})`$ $`=`$ $`{\displaystyle \frac{5}{4A^3}}{\displaystyle \frac{\mathrm{tanh}(\pi \sqrt{5}A^{})(\pi \sqrt{5}A^{})(13\mathrm{tanh}^2(\pi \sqrt{5}A^{}))}{\mathrm{cosh}^2(\pi \sqrt{5}A^{})}}`$ (55)
Now, if we consider the distribution of the sum of the areas $`𝒜=_{i=1}^nA_i`$, it is obtained by setting $`B_j=B,j`$. (52) leads to ($`𝒜^{}=𝒜/t`$):
$$e^{iB𝒜^{}}=\frac{B}{2\mathrm{sinh}\left(\frac{B}{2}\right)}$$
(56)
With (5), we see that the sum of areas has, asymptotically, exactly the same distribution as the area enclosed by a single Brownian particle. In fact, we can compute $`e^{iB𝒜^{}}`$ for all $`t`$ values (and not only when $`t\mathrm{}`$). This is because, that time, the matrix $`𝐁`$ ($`𝐁_{ij}=B\delta _{ij})`$ commutes with $`𝐌`$. So, we are left with an $`\{`$ harmonic oscillator + uniform magnetic field $`\}`$ problem for each normal coordinate (except for $`Z_1`$, that only feels a pure magnetic field). We get the result :
$`e^{iB𝒜^{}}`$ $`=`$ $`{\displaystyle \frac{B}{2\mathrm{sinh}\left(\frac{B}{2}\right)}}{\displaystyle \underset{i=2}{\overset{n}{}}}{\displaystyle \frac{F_i(B)}{F_i(0)}}`$ (57)
$`F_i(B)`$ $`=`$ $`{\displaystyle \frac{\omega _i^{}}{2\pi \mathrm{sinh}(t\omega _i^{})}}\mathrm{exp}\left({\displaystyle \frac{\omega _i^{}}{2\pi \mathrm{sinh}(t\omega _i^{})}}\left(\mathrm{cosh}(t\omega _i^{})\mathrm{cosh}(B/2)\right)|Z_i^{(0)}|^2\right)`$ (58)
$`\omega _i^{}`$ $`=`$ $`\sqrt{\omega _i^2+\left({\displaystyle \frac{B}{2t}}\right)^2}`$ (59)
We recover (56) in the limit $`t\mathrm{}`$ ($`_{i=2}^nF_i(B)/F_i(0)1`$ when $`t\mathrm{}`$).
To close this section, it is interesting to consider the asymptotic law for the area $`A_j^{}`$ ($`=A_j/t`$) enclosed by a given monomer $`j`$. (52) gives:
$$e^{iB_jA_j^{}}=\frac{B_j}{2\sqrt{n}\mathrm{sinh}\left(\frac{B_j}{2\sqrt{n}}\right)}$$
(60)
It follows that $`A_j`$ satisfies Lévy’s law (5) and scales like $`\frac{t}{\sqrt{n}}`$. Remark that the area swept by the chain c.o.m., G, should scale like $`\frac{t}{n}`$. On the other hand, for the same gaussian noise, we would get $`A_jt`$ if particle $`j`$ was free (i.e. $`k=0`$). The actual scaling of $`A_j`$ is intermediate: this particle moves more freely than G but it is embedded in the chain, thus not completely free!
Those considerations allow us to give a more precise sense to the statement: “the process is dominated by the c.o.m. motion”. This one is true as long as we look at occupation times. However, when we study finer quantities like areas, this sentence must be corrected. Similar (even more dramatic) deviations will occur when we look at winding angles.
To end up with areas, let us remark that the case of open trajectories can be treated exactly along the same lines as the one developed here, without additionnal difficulties (in particular, (48) still holds). We will not address this problem in the present work.
## 5 Winding angles distribution
The last part of this paper will be devoted to the distribution $`P(\{\theta _j\})`$ ($`\theta _j`$ is the angle wound around O by particle $`j`$ during a time $`t`$). We consider the same conditions as for Spitzer’s law ($`z^{(0)}`$, initial configuration, fixed, with $`z_j^{(0)}0,j`$ ; $`z`$, final configuration, unspecified ; $`t\mathrm{}`$).
We want to proceed as before and insert the constraint:
$$\underset{j=1}{\overset{n}{}}\delta \left(\theta _j\frac{1}{2i}_0^t\left(\frac{z_j\dot{\overline{z}_j}\overline{z}_j\dot{z_j}}{z_j\overline{z}_j}\right)d\tau \right)$$
(61)
in the Wiener measure (2). We are now faced with the problem of $`n`$ harmonically bound particles submitted to the magnetic fields of $`n`$ different point-like vortices located at the origin. The corresponding hamiltonian is:
$`H_\lambda `$ $`=`$ $`H_0+V_\lambda `$ (62)
$`V_\lambda `$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}\lambda _i\left({\displaystyle \frac{1}{z_i}}_{\overline{z}_i}{\displaystyle \frac{1}{\overline{z}_i}}_{z_i}\right)+{\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \frac{\lambda _i^2}{2z_i\overline{z}_i}}`$ (63)
and the distribution $`P(\{\theta _i\})`$ is given by:
$$P(\{\theta _i\})=\left(\underset{j=1}{\overset{n}{}}\frac{\mathrm{d}\lambda _j}{2\pi }e^{i\lambda _j\theta _j}\right)dzd\overline{z}F(z,z^{(0)})G_\lambda (z,z^{(0)},t)$$
(64)
$$G_\lambda (z,z^{(0)},t)=z\left|e^{tH_\lambda }\right|z^{(0)}$$
(65)
Studying the limit $`t\mathrm{}`$, we cannot develop directly as before a perturbation theory with $`V_\lambda `$: this is because the last term in $`V_\lambda `$ gives a divergent contribution . Due to this term, all the eigenfunctions of $`H_\lambda `$ must vanish in O at least as $`_{i=1}^n|z_i|^{|\lambda _i|}`$ ($`U(z)`$). So, we redefine those eigenfunctions :
$$\mathrm{\Psi }=U\stackrel{~}{\mathrm{\Psi }}$$
(66)
The new hamiltonian acting on $`\stackrel{~}{\mathrm{\Psi }}`$ is
$`\stackrel{~}{H_\lambda }`$ $`=`$ $`H_0+\stackrel{~}{V_\lambda }`$ (67)
$`\stackrel{~}{V_\lambda }(z)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}\left((\lambda _i|\lambda _i|){\displaystyle \frac{1}{z_i}}_{\overline{z}_i}(\lambda _i+|\lambda _i|){\displaystyle \frac{1}{\overline{z}_i}}_{z_i}\right)`$ (68)
with a propagator $`\stackrel{~}{G_\lambda }`$
$$\stackrel{~}{G_\lambda }(z,z^{(0)},t)=z\left|e^{t\stackrel{~}{H_\lambda }}\right|z^{(0)}=\frac{U(z^{(0)})}{U(z)}G_\lambda (z,z^{(0)},t)$$
(69)
($`\stackrel{~}{G_0}=G_0`$). That time, the perturbation theory is properly defined and we can compute the characteristic function:
$$C(\{\lambda _j\})e^{i_{j=1}^n\lambda _j\theta _j}=dzd\overline{z}\left(\underset{j=1}{\overset{n}{}}\frac{|z_j|^{|\lambda _j|}}{|z_j^{(0)}|^{|\lambda _j|}}\right)F(z,z^{(0)})\stackrel{~}{G_\lambda }(z,z^{(0)},t)$$
(70)
with, symbolically,
$$\stackrel{~}{G_\lambda }=\underset{m=0}{\overset{\mathrm{}}{}}(1)^mG_0^{\mathrm{}}(\stackrel{~}{V_\lambda }G_0^{\mathrm{}})^m$$
(71)
Using integration by parts and also the relationship $`_{z_i}\left(\frac{1}{\overline{z}_i}\right)=\pi \delta (z_i)`$, we first calculated $`C(\{\lambda _i\})`$ up to $`4^{\mathrm{th}}`$ order in $`\stackrel{~}{V_\lambda }`$, with the result:
$`C(\{\lambda _j\})e^{X/2}D(X)`$ (72)
$`D(X)=1+`$
$`+\left({\displaystyle \frac{n+1}{2}}\right)\left({\displaystyle \frac{X}{1!}}+{\displaystyle \frac{X^2}{2!}}n{\displaystyle \frac{X^3}{3!}}\left({\displaystyle \frac{3n^21}{2}}\right)+{\displaystyle \frac{X^4}{4!}}(3n^32n)\mathrm{}\right)`$ (73)
$`X=\left({\displaystyle \underset{i=1}{\overset{n}{}}}|\lambda _i|\right)\mathrm{ln}t`$ (74)
The prefactor $`e^{X/2}`$ comes out from $`U(z)`$ in (70) when integrated over the final configuration: it will be present at all orders of the computation. Morover, (72) suggests that $`C(\{\lambda _i\})`$ is only a function of $`X`$: this is actually the case, as will be shown in the sequel.
Let us consider the $`m^{\mathrm{th}}`$ order term in (70,71) and suppose that we integrate, first, over $`z,z^{(m)},z^{(m1)},\mathrm{},z^{(k+1)}`$. Following the computation step by step, it is not difficult to convince oneself that the integration over $`z^{(k)}`$ involves expressions of the form:
$$d\overline{z}^{(k)}dz^{(k)}\varphi (z^{(k)},T)\stackrel{~}{V_\lambda }(z^{(k)})G_0^{\mathrm{}}(z^{(k)},z^{(k1)},t_kt_{k1})$$
(75)
$$\mathrm{where}\varphi (z^{(k)},T)=e^{\frac{|Z_1^{(k)}|^2}{2T}}e^{\frac{1}{2}_{i=2}^nk\omega _i|Z_i^{(k)}|^2}$$
(76)
and $`T=t_lt_k,k+1lm`$. Let us call $`J_k`$ the result of (75). In the limit of long times, it reads:
$`J_k=({\displaystyle \underset{i=1}{\overset{n}{}}}|\lambda _i|)\times `$ (77)
$`\times `$ $`\left({\displaystyle \frac{n+1}{2(t_kt_{k1})}}\varphi (z^{(k1)},t_kt_{k1})+{\displaystyle \frac{1}{T+t_kt_{k1}}}\varphi (z^{(k1)},T+t_kt_{k1})\right)`$
The $`m`$ successive spatial integrations produce the factor $`\left(|\lambda _i|\right)^m`$ and, at the end, we are left with time integrals of the form:
$$I_m(i_{m1},\mathrm{},i_0)(t)=_0^tdt_m_0^{t_m}dt_{m1}\mathrm{}_0^{t_2}dt_1\frac{e^{_{i=1}^m\frac{\alpha _i}{t_i}_{i=1}^{m1}\frac{\beta _i}{t_{i+1}t_i}}}{(t_{i_{m1}}t_{m1})\mathrm{}(t_{i_1}t_1)t_{i_0}}$$
(78)
with $`\alpha _i,\beta _i>0`$ and
$$i_{m1}=i_{m2}=\mathrm{}=i_k=m;i_{k1}=i_{k2}=\mathrm{}=i_l=k;i_{l1}=i_{l2}=\mathrm{}=i_j=l;\mathrm{}$$
(79)
We have proved, step by step, that;
$$I_m(i_{m1},\mathrm{},i_0)(t)_t\mathrm{}\frac{\left(\mathrm{ln}t\right)^m}{_{l=0}^{m1}(i_ll)}$$
(80)
Those considerations show that, actually, $`C(\{\lambda _i\})`$ is only a function of $`X`$. So, we can write $`D(X)=_{m=0}^{\mathrm{}}a_mX^m`$, with $`a_0=1`$ (see eq.(73)).
Moreover, with the help of the above equation (80), and also looking at the tree structure exhibited in eq.(77), the following recursion relation can be shown:
$`a_m`$ $`=`$ $`y{\displaystyle \underset{k=0}{\overset{m1}{}}}{\displaystyle \frac{a_k}{(mk)!}}`$ (81)
$`y`$ $`=`$ $`{\displaystyle \frac{n+1}{2}}`$ (82)
It allows to write a closed form formula for $`D(X)`$:
$$D(X)=\frac{1}{1y(e^X1)}=\frac{e^{X/2}}{\mathrm{cosh}(X/2)+n\mathrm{sinh}(X/2)}$$
(83)
With (72) and, also, a rescaling of the angles $`\left(\theta _i^{}=\frac{2\theta _i}{\mathrm{ln}t}\right)`$, we get the desired characteristic function:
$`e^{i_{j=1}^n\lambda _j\theta _j^{}}`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{cosh}(u)+n\mathrm{sinh}(u)}}`$ (84)
$`u`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}|\lambda _i|`$ (85)
(with $`n=1`$, we recover eq.(3)).
We consider (84) as the main result of this paper.
Finally, Fourier transformation shows that $`P(\{\theta _j^{}\})`$ is an “infinite sum of products of Spitzer’s laws” (!) with highly correlated variables:
$`P(\{\theta _j^{}\})={\displaystyle \frac{2}{n+1}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}\left\{\left({\displaystyle \frac{n1}{n+1}}\right)^k\left({\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \frac{1}{\pi (2k+1)}}{\displaystyle \frac{1}{1+\left(\frac{\theta _j^{}}{2k+1}\right)^2}}\right)\right\}`$ (86)
All the moments of this distribution are infinite (unless they trivially vanish).
For a given particle $`j`$ of the chain, we have:
$$e^{i\lambda _j\theta _j^{}}=\frac{1}{\mathrm{cosh}(\lambda _j)+n\mathrm{sinh}(|\lambda _j|)}$$
(87)
that leads to:
$$P(\theta _j^{})=\frac{2}{n+1}\underset{k=0}{\overset{\mathrm{}}{}}\left\{\left(\frac{n1}{n+1}\right)^k\frac{1}{\pi (2k+1)}\frac{1}{1+\left(\frac{\theta _j^{}}{2k+1}\right)^2}\right\}$$
(88)
The difference with Spitzer’s law is due to the presence of $`n`$ in the denominator of (87).
To shed some light on this problem, let us go back to the joint law (6) of small and big windings for the chain c.o.m.. What could we expect for the corresponding windings of particle $`j`$? With little effort, we can say that:
1. The big windings will be roughly the same for both (when the chain is far from O, particle $`j`$ follows the c.o.m. and winds around O in the same way). So, we keep $`\lambda _+`$ unchanged in (6).
2. The small windings will be quite different. This is because particle $`j`$ is artificially maintained in the vicinity of O: despite its higher mobility, it spends the same time as the c.o.m. in a given domain surrounding O. As a consequence, its small windings law will be broadened. Assuming that the remark following (7) holds, we get this broadening by changing $`|\lambda _{}|`$ into $`n|\lambda _{}|`$ in (6) ($`n`$ is the ratio of the diffusion constants; of course, we don’t say at all that (7) is the law of small windings!).
Thus, our guess for particle $`j`$ is:
$$e^{i(\lambda _+\theta _{j+}^{}+\lambda _{}\theta _j^{})}=\frac{1}{\mathrm{cosh}(\lambda _+)+n\frac{|\lambda _{}|}{\lambda _+}\mathrm{sinh}(\lambda _+)}$$
(89)
Setting $`\lambda _+=\lambda _{}=\lambda _j`$, we recover (87). We are aware that this argument is strictly heuristic and that (89) remains to be proved. Nevertheless, we think that it allows to explain correctly the presence of $`n`$ in (87).
## 6 Conclusion
Let us briefly summarize this work. We have computed explicitly the asymptotic joint laws of the occupation times, areas and winding angles of a chain of harmonically bound Brownian particles.
For all these properties, we have shown that the scaling variables take the same general form as for the standard Brownian motion. However, a detailed study reveals important specific features that reflect a subtle interplay between the free c.o.m. motion – that strongly influences the whole chain properties – and the relative freedom of a given particle of the chain. For occupation times distributions, it appears that the c.o.m. satisfies the same law as a given monomer; now, for the areas, the scaling becomes slightly different and, finally, for the winding angles, the law itself is changed. Remark also that correlations are systematically present.
Moreover, we observe that the scaling variables and the laws are very different from those met in our study of the attached Rouse chain (in the latter case, $`\theta t,A\sqrt{t}`$, and the winding angles are uncorrelated). These differences are not so surprising since, in that case, we had no translation invariance.
One of us (O.B.) acknowledges Dr. G. Oshanin for drawing his attention to this problem.
e-mail:
benichou@lptl.jussieu.fr
desbois@ipno.in2p3.fr
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# The lines of the Kontsevich integral and Rozansky’s rationality conjecture
## 1. Introduction
In some lectures at the Joseph Fourier Institute in June of 1999, Lev Rozansky formulated an important conecture concerning the structure of the Kontsevich integral \[Roz\], and mentioned something of a related program for a “finite-type theory of knots’ complements”.
By means of some notation let us now describe what of this will be proved in this paper.
A generating diagram is a diagram with oriented trivalent vertices (which is to say that the incident edges at a trivalent vertex are cyclically ordered) and edges decorated with oriented bivalent vertices labelled by elements of $`[[k]]`$, the ring of formal power series in a variable $`k`$. A generating diagram represents a series of elements of $`B(k)`$ by expanding these power series into series of diagrams, as follows. If $`f(k)[[k]]`$ is
$$f=f_0+f_1k+f_2k^2+f_3k^3\mathrm{},$$
where $`f_i`$, then an edge labelled with $`f(k)`$ is to be expanded as follows.
$$\text{}=f_0\text{}+f_1\text{}+f_2\text{}+f_3\text{}\mathrm{}$$
###### Remark 1.0.1.
The incoming edges at the label (which Definition 3.0.2 will introduce as a “winding coupon”) are ordered, which determines the orientation of the introduced trivalent vertices, as shown. The opposite ordering with the label $`f(k)`$ gives the same series.
###### Definition 1.0.2.
Define the gd-degree of a generating diagram to be half the number of trivalent vertices of the (original) diagram.
###### Definition 1.0.3.
Take a knot $`K`$ in an integral homology three-sphere $`M`$. Let $`A_{(M,K)}(t)`$ denote the Alexander polynomial of the pair $`(M,K)`$ fixed by the requirement that it satisfy:
1. $`A_{(M,K)}(t)=A_{(M,K)}(\frac{1}{t}),`$
2. $`A_{(M,K)}(1)=1.`$
Let $`^1(t)`$ be the ring of rational functions in $`t`$ that are non-singular at 1. Denote the inclusion
$$\iota :^1(t)[[k]],$$
defined by subtituting $`e^k`$ into $`t`$.
###### Definition 1.0.4.
Let $`L_{(M,K)}`$ be the $``$-vector subspace of $`[[k]]`$ that is the image under $`\iota `$ of rational functions of the form
$$\frac{P(t)}{A_{(M,K)}(t)}$$
where $`P(t)[t,t^1]`$. This notation can be read as the space of labels.
The next two definitions require the “wheel with $`2n`$ spokes”:
$$\omega _2=\text{},\omega _4=\text{},\omega _6=\text{},\mathrm{}$$
###### Definition 1.0.5.
If the rational numbers $`b_{2n}`$ are determined by the equality:
$$b_{2n}x^{2n}=\frac{1}{2}\text{log}\left(\frac{\text{sinh}(\frac{x}{2})}{\frac{x}{2}}\right),$$
then the series $`\nu (k)B(k)`$ is defined to be
$$\text{exp}_{}(b_{2n}\omega _{2n}).$$
###### Remark 1.0.6.
This has recently been shown to be $`\widehat{Z}(U)`$, the Kontsevich integral of the unknot, by Bar-Natan, Le and Thurston \[TW\].
###### Definition 1.0.7.
Let $`Wh(M,K)`$ be defined by
$$Wh(M,K)=\text{exp}_{}\left(\left[\frac{1}{2}\text{log}\left(A_{(M,K)}(e^h)\right)\right]|_{h^{2n}\omega _{2n}}\right)\nu (k),$$
where the operation indicated is to expand the term inside the square brackets into a power series in $`u`$, and then to replace terms like $`ch^{2n}`$ by $`c\omega _{2n}`$.
The LMO invariant was introduced by Thang Le, Jun Murakami and Tomotada Ohtsuki \[LMO\] (following an earlier investigation also with Hitoshi Murakami \[LMMO\]). Our (perhaps non-standard) normalisation of the non-surgered component specialises to the three-sphere as follows:
(1.0.1)
$$Z^{LMO}(S^3,K)=\widehat{Z}(K).$$
This brings us to the rationality conjecture. In his lectures Rozansky conjectured the $`MS^3`$ case of the following (see also the new paper by Garoufalidis and Rozansky \[GR\]).
###### Theorem 1.0.8.
Let $`K`$ be a zero-framed knot in an integral homology three-sphere $`M`$. The LMO invariant of this pair may be represented
$$Z^{LMO}(M,K)=Wh(M,K)\text{exp}_{}(r)B(k),$$
where $`r=_{m=1}^{\mathrm{}}r^{(m)}`$ with $`r^{(m)}`$ a finite $``$-linear combination of connected generating diagrams of $`gd`$-degree $`m`$ whose edges are labelled from $`L_{(M,K)}`$.
This conjecture is motivated by an analogous property of the coloured Jones function, as shown by Rozansky (see, for example, \[Roz2\]).
The structures we describe which lead to the proof of this depend on a delicate assembly of results from the literature. A significant debt is to the papers of the Aarhus group, who are Bar-Natan, Garoufalidis, Rozansky and Thurston \[A1, A2, A3\].
This theory was described at the workshop “Art of Low-Dimensional Topology VII”, January 8th, 2000, for which invitation the author thanks Toshitake Kohno.
A sequel to this paper \[KS\] will develop some technical issues raised within this work.
###### Note 1.0.9.
Between versions of this paper, a new paper by Garoufalidis and Rozansky appeared \[GR\]. We (the author) suggest reading these papers in tandem, as their concerns are somewhat complementary.
###### Acknowledgements 1.0.10.
The author is supported by a Japan Society for the Promotion of Science Postdoctoral Fellowship. Thanks to Tomotada Ohtsuki and the Department of Mathematical and Computing Sciences at the Tokyo Institute of Technology for their support; to Kazuo Habiro and Dylan Thurston for many helpful comments regarding this work; and also to Louis Funar, Stavros Garoufalidis and Hitoshi Murakami.
## 2. The outline
### 2.1. Special surgery presentations
The strategy of the calculation to be presently described is to apply the LMO surgery formula to a special surgery presentation which exists for any knot in a $`HS^3`$. The following is well-known.
###### Lemma 2.1.1.
A zero-framed knot in a $`HS^3`$ may be obtained from the zero-framed unkot $`U`$ in $`S^3`$ by performing surgery on some framed link which has the property that every component of it has linking number 0 with $`U`$.
It may be useful to keep the following example in mind. The figure of 8 knot is obtained by performing surgery on the (blackboard-framed) component marked with a $``$ below:
Thus a knot in a $`HS^3`$ can be presented by a framed link in a solid torus (fixed in $`S^3`$) such that every component has linking zero with the core of the torus.
It will prove technically advantageous to work with a slightly different object: a framed string link in the solid torus. For this definition, realise the solid torus $`ST`$ as the complement in the cube $`\{(x,y,z)^3;0x1,0y1,0z1\}`$ of the hole $`\{(x,y,z)^3;\frac{1}{4}<x<\frac{3}{4},\frac{1}{4}<y<\frac{3}{4},0z1\}`$.
###### Remark 2.1.2.
We thus use the definite article “the”, as in “the solid torus”, to remind that we are referring to a particular solid torus embedded in $`S^3`$.
###### Definition 2.1.3.
A $`\mu `$-string string link in the solid torus is a proper embedding $`[0,1]\mathrm{}[0,1]ST`$ such that the $`i`$th $`\{0\}`$ is mapped to $`(\frac{i}{\mu +1},0,\frac{1}{2})`$, such that the $`i`$th $`\{1\}`$ is mapped to $`(\frac{i}{\mu +1},1,\frac{1}{2})`$, and with a framing in the familiar sense of a framed tangle. These are identified up to framed isotopies in the solid torus.
###### Remark 2.1.4.
We will refer to the $`y=0`$ plane as the base, and the $`y=1`$ plane as the top. In this work string links in the solid torus will always be oriented from the base to the top.
###### Definition 2.1.5.
We will call the meridional disc the disc $`\{x[\frac{3}{4},1],y=\frac{1}{2},z[0,1]\}`$ in the solid torus.
To draw a diagram of a string link in a solid torus, we will take a projection in general position onto the $`xy`$ plane, in the familiar sense. It is convenient to represent the “hole” as as a fixed dashed loop, or to represent the “meridional disc” as a short dashed line segment. A diagram will be called in general position with respect to the meridional disc if it intersects that dashed line transversally. For example:
or
###### Definition 2.1.6.
1. Let a marked framed tangle (resp. link) be a framed tangle (resp. link), possibly with some distinguished components.
2. For a string link in the solid torus $`T`$, let Thr$`(T)`$ denote the marked framed tangle obtained by marking all components, threading the hole with an unmarked zero-framed unknot (to fix this let us say we thread on the $`x>1`$ side), and then forgetting the hole.
3. For a marked framed tangle $`T^{}`$, let Clos$`(T^{})`$ denote the marked framed link obtained by closing the tangle.
4. For a marked framed link $`L`$, let KII$`(L)`$ denote the class of $`L`$ modulo Kirby move IIs, where markings indicate to-be-surgered components (slides of unmarked components over marked components are also allowed).
We now restrict to the presentations guaranteed by Lemma 2.1.1. In the following, the adjective special will just mean one of examples in question. In particular:
###### Definition 2.1.7.
Let a special string link in the solid torus be a string link in the solid torus $`T`$ such that:
1. Every component has zero algebraic intersection number with the meridional disc, for diagrams in general position with respect to the meridional disc.
2. The determinant of the linking matrix of the marked components of
Clos(Thr($`T`$)) is $`\pm 1`$.
###### Definition 2.1.8.
Let a special tangle be a marked, framed tangle with one closed, unmarked component; forgetting that component leaves the tangle a string link whose linking matrix has determinant $`\pm 1`$.
###### Definition 2.1.9.
Let a special link be a marked, framed link with one unmarked component; the determinant of the linking matrix of marked components is $`\pm 1`$.
### 2.2. The master diagram
Here is the plan. The basic idea is that formulae for the Kontsevich integral of a knot can be obtained by applying the LMO invariant to the surgery presentations just considered. We will see that the resulting factorisation of the calculation through a certain invariant of string links in the solid torus has important implications for the result.
The following diagram records this factorisation, as we will now explain.
###### Diagram 2.2.1.
The invariant
$$\sigma \stackrel{ˇ}{Z}^{ST}:\left\{\begin{array}{c}Specialstringlinks\\ inthesolidtorus.\end{array}\right\}B^{ST}(X)^{Int}$$
is introduced in Section 3. This is an enhancement of the usual Kontsevich integral, taking values in a space of winding diagrams. A winding diagram is, in an appropriate sense, a uni-trivalent diagram decorated by winding coupons.
Intuitively speaking, this decoration (modulo some relations) describes a homotopy class of proper mappings of that diagram into the solid torus. In this intuitive picture winding coupons record intersections of the edges of some representative in general position with respect to some fixed meridional disc, with that disc. This is illustrated in the next figure.
This invariant is more or less pre-existent in the literature, though our approach and the structures we describe depart from existent works in certain ways. Our formal presentation, with labelled edges, is closest to that of Goryunov \[G\], and the intuitive picture is closest to that of Andersen-Mattes-Reshetikhin \[AMR\] (see also Suetsugue \[S\]).
The destination indicated, $`B^{ST}(X)`$, is a space of symmetrised winding diagrams. This plays the part of Bar-Natan’s algebra $`B`$: legs on skeletons are to be symmetrised. This space is introduced in Section 3.4. The map $`\sigma `$ is the appropriate version of Bar-Natan’s “formal Poincare-Birkhoff-Witt” map.
The utility of this invariant is expressed by the top face of the cube. Namely, the Kontsevich integral of one of the tangles of interest, $`Thr(T)`$, factors through $`\stackrel{ˇ}{Z}^{ST}`$. The map which completes the square is expressed in terms of a map $`Thr^D`$, threading diagrams. This is introduced in Section 5. This map,
$$Thr^D:B^{ST}(X)B(X,\underset{¯}{k})$$
is the operation of replacing winding coupons (“intersections with the meridional disc”) with exponentials of legs. The commutativity of this face depends crucially on a recent calculation due to Bar-Natan, Le and Thurston, as is indicated in that section.
$$Thr^D\left(\text{}\right)=\text{}+\text{}+\frac{1}{2!}\text{}+\frac{1}{3!}\text{}\mathrm{}$$
The subspace $`B^{ST}(X)^{Int}B^{ST}(X)`$ is the subspace of “integrable” elements (adapting a concept of the Aarhus papers to the present context), as is defined in Section 4. In this case it refers to the subspace of elements of the form
$$S=exp_{}\left(\frac{1}{2}\underset{i,j}{}\text{}\right)R,$$
where $`W(t)M_\mu ([t,t^1])`$ is a Hermitian matrix of Laurent polynomials, such that $`det(W(1))=\pm 1`$, and $`R`$ is a series of diagrams without chords. In the case at hand ($`T`$ a string link in the solid torus), $`\sigma (\stackrel{ˇ}{Z}^{ST}(T))`$ is of this form with matrix $`W(T,t)`$, the winding matrix of $`T`$, introduced in Section 3.5. This is a generalisation of the notion of linking matrix which incorporates winding information of the link around the hole of the solid torus.
The mapping $`^{FGinST}𝑑X`$ (again, an adaption of a concept from the Aarhus papers) is defined in the following way. Take an element $`SB^{ST}(X)^{Int}`$, with decomposition as above.
$$^{FGinST}𝑑XS=exp_{}\left(\frac{1}{2}\underset{i,j}{}\text{}\right),R.$$
This takes values in the space $`B^{QST}(\varphi )`$ of rational winding diagrams. This space is in some sense an extension of $`B^{ST}(\varphi )`$ which admits rational functions as labels on winding coupons. This mapping and space are introduced in Section 3.5.
The front face of the master diagram is detailed in Section 7. The commutativity of the front face indicates that this formula calculates the LMO invariant (in the event, an extension of a theorem due to the Aarhus group \[A3\]). Note that $`\sigma _+`$ and det are just some normalisation factors that need to be carried along for the diagram to make sense.
The theorem to take home is the following.
###### Theorem 2.2.2 (Surgery formula).
Let a pair of a zero-framed knot $`K`$ in an integral homology three-sphere $`M`$ be presented by $`T`$, some special string link in the solid torus. Then $`Z_n^{LMO}(M,K)`$ is equal to
$$\frac{Wh(M,K)Thr^D\left(^{FGinST}𝑑X\sigma (\stackrel{ˇ}{Z}^{ST}(T))\right)}{\left((1)^n^{(n)}𝑑U\sigma (\stackrel{ˇ}{Z}(U_+))\right)^{\sigma _+(W(T,1))}\left(^{(n)}𝑑U\sigma (\stackrel{ˇ}{Z}(U_{}))\right)^{\sigma _{}(W(T,1))}}B_n(\underset{¯}{k}).$$
Observe that in this setting ($`HS^3`$s), $`Z_n^{LMO}`$ is the degree less or equal to $`n`$ truncation of the full LMO invariant $`Z^{LMO}`$. We may alternatively present this formla as follows.
###### Theorem 2.2.3 (Surgery formula, $`\beta `$-version.).
Let a pair of a zero-framed knot $`K`$ in an integral homology three-sphere $`M`$ be presented by $`T`$, some special string link in the solid torus. Then $`Z^{LMO}(M,K)`$ is equal to
$$\frac{Wh(M,K)Thr^D\left(^{FGinST}𝑑X\sigma (\stackrel{ˇ}{Z}^{ST}(T))\right)}{\left(^{FG}𝑑U\sigma (\stackrel{ˇ}{Z}(U_+))\right)^{\sigma _+(W(T,1))}\left(^{FG}𝑑U\sigma (\stackrel{ˇ}{Z}(U_{}))\right)^{\sigma _{}(W(T,1))}}B(\underset{¯}{k}).$$
### 2.3. Conjecture - a winding diagram valued invariant of knots
Clearly we are only seeing half of a cube in Diagram 2.2.1. It will be interesting to describe the other vertex and faces. More immediately, the dashed line
$$\left\{\frac{Speciallinks}{KirbymoveIIs}\right\}B^{QST}(\varphi )\times \times ^1[t,t^1]$$
would be a consequence of the following conjecture.
###### Conjecture 2.3.1.
$`Thr^D:B^{QST}(\varphi )B(\underset{¯}{k})`$ is injective.
Actually, this seems clear; we defer a careful explanation of this to the sequel, which will also discuss the relation with some normalisation and other technical issues (see Section 3.7). Can this corollary be proved without reference to the Kontsevich integral?
###### Corollary 2.3.2 (to the conjecture.).
Take a pair $`(M,K)`$. Choose $`T`$, a special string link in the solid torus presenting $`(M,K)`$. Then
$$\frac{^{FGinST}𝑑X\sigma (\stackrel{ˇ}{Z}^{ST}(T))}{\left(^{FG}𝑑U\sigma (\stackrel{ˇ}{Z}(U_+))\right)^{\sigma _+(W(T,1))}\left(^{FG}𝑑U\sigma (\stackrel{ˇ}{Z}(U_{}))\right)^{\sigma _{}(W(T,1))}}B^{QST}(\varphi )$$
is an invariant of the pair $`(M,K)`$.
Strictly speaking, this does not increase our pool of knot invariants. It may, however, be a presentation more appropriate for topological applications.
### 2.4. Notation and conventions
Spaces of diagrams
We use standard definitions for spaces of uni-trivalent diagrams \[BN\]. There is one point that may be unfamiliar to some readers. To introduce this, we note that in generality a space may be denoted:
$$A_n(SK;x_1,\mathrm{},x_p,\underset{¯}{w}_1,\mathrm{},\underset{¯}{w}_q)$$
which indicates that univalent vertices may:
1. be located, up to orientation-preserving diffeomorphisms, on a skeleton $`SK`$,
2. or be labelled from $`\{x_1,\mathrm{},x_p,\underset{¯}{w}_1,\mathrm{},\underset{¯}{w}_q\}`$.
The underlining of a variable indicates that link relations for that variable are to be included in the definition. These were identified in Section 5.2 of \[A2\]. The point is that link relations are what must be included to obtain an isomorphism:
$$A_n(SK;x_1,\mathrm{},x_p,\underset{¯}{w}_1,\mathrm{},\underset{¯}{w}_q)A_n(SK\underset{p}{\underset{}{\mathrm{}}}\text{}).$$
On occasions when there is no skeleton (that is to say that all univalent vertices are to be labelled), we may use a $`B`$ in place of the $`A`$, following the conventions of \[BN\]. When there is only one label, then $`B(\underset{¯}{x})B(x)`$.
In this work, we usually work with a set of labels $`X=\{x_1,\mathrm{},x_\mu \}`$ corresponding to string link components (ultimately providing surgery components), and/or a label $`k`$ corresponding to the (closed) knot component. Such spaces will be denoted, for example, $`B(X,\underset{¯}{k})`$.
An element of this space may be indicated $`S(\overline{x},k)`$, where $`\overline{x}`$ is thought of as a vector of variables. The logic of this notation should be clear after reading Section 7.1.
The invariants
The notation $`\widehat{Z}`$ denotes precisely the functorial representation of the category of framed $`q`$-tangles, according to the definition of \[LM\] (see also \[BN2\]). We use, unless otherwise stated, the associator with rational coefficients.
Our definition of $`\stackrel{ˇ}{Z}`$ differs slightly from existent usage. In general we will be considering tangles which have a closed component, which, when forgotten, leaves the tangle a string link. Take such a tangle $`T`$, which has $`n`$ such string link components. Our usage of $`\stackrel{ˇ}{Z}(T)`$ is:
$$\stackrel{ˇ}{Z}(T)=(\underset{n}{\underset{}{\nu \mathrm{}\nu }})\mathrm{\Delta }^{n1}(\nu )\widehat{Z}(T).$$
The LMO invariant
The LMO group are Thang Le, Jun Murakami and Tomotada Ohtsuki. In Theorem 6.2 of \[LMO\] there is defined an invariant of a link in a three-manifold, which is denoted there $`\mathrm{\Omega }(M,L)`$. The invariant we use is related to the definition given there as follows:
$$Z^{LMO}(M,K)=\mathrm{\Omega }(M,K)\mathrm{\#}\nu ^1.$$
Thus, following Proposition 6.5 of \[LMO\], we have the restriction
(2.4.1)
$$Z^{LMO}(S^3,K)=\widehat{Z}(K).$$
## 3. A diagram-valued invariant of string links in the solid torus
###### Definition 3.0.1.
A skeleton is an oriented one-manifold whose boundary points are seperated into an ordered pair of ordered sets.
Our theory employs a certain enhancement of the familiar notion of uni-trivalent diagram, which will be called a winding diagram. These will possibly contain a certain new type of vertex which we will call a winding coupon. This may conveniently be thought of as some decoration of an underlying uni-trivalent diagram.
###### Definition 3.0.2.
A winding coupon is a bivalent vertex whose incoming edges are ordered.
This is depicted as follows. The edges are ordered so that the edge incoming at the base of $`t`$ is first, and the edge outgoing at the top of $`t`$ is second. We make the convention that a coupon labelled by $`t^1`$ is given the reverse orientation.
which is equivalent to
###### Definition 3.0.3.
A winding diagram on a skeleton $`SK`$ is a graph with univalent vertices, trivalent vertices and winding coupons, such that:
1. (Boundary points.) The set of univalent vertices is seperated into an ordered pair of ordered sets.
2. (Skeleton.) There are distinguished disjoint oriented cycles and oriented paths between univalent vertices which are labelled by components of the skeleton $`SK`$; forgetting not distinguished edges leaves one with $`SK`$.
3. (Internal verices.) Trivalent vertices which are not met by distinguished edges are vertex-oriented (have their incoming edges cyclically ordered).
Two winding diagrams are identified if there is a graph isomporhism between them respecting orientations at vertices (including winding coupons) and respecting skeleton information (orientations and labels of distinguished edges, ordering of boundary points).
###### Definition 3.0.4.
The grade of a winding diagram is half the number of trivalent vertices.
The space in question will be a quotient of the space of finite $``$-linear combinations of winding coupons of a fixed grade. The quotient will be by the span of the following classes of vectors. In the relations MULT and PUSH below, edges may be part of the skeleton.
$`AS:`$ $`\text{}+\text{}`$
$`IHX:`$ $`\text{}\text{}\text{}`$
$`STU:`$ $`\text{}\text{}+\text{}`$
$`MULT:`$ $`\text{}\text{}`$
$`PUSH:`$ $`\text{}\text{}`$
###### Remark 3.0.5.
Looking ahead to the definition of the map $`Thr^D`$, Definition 5.0.3, may give some feeling for these orientation conventions and relations.
###### Definition 3.0.6.
Let $`SK`$ be a skeleton. Let
$$A_m^{ST}(SK)=\frac{\left\{\begin{array}{c}\text{Finite }\text{-linear combinations of}\hfill \\ \text{degree }m\text{ winding diagrams on }SK\hfill \end{array}\right\}}{\text{span of above relations}}$$
Let $`A^{ST}(SK)`$ denote the completion of $`_{m=0}^{\mathrm{}}A^{ST}(SK)`$ with respect to degree.
###### Remark 3.0.7.
If $`SK`$ has $`\mu `$ components, then $`A_0^{ST}(SK)`$ is isomorphic, as a $``$-vector space, to $`[t_1^{\pm 1},\mathrm{},t_\mu ^{\pm 1}]`$.
###### Remark 3.0.8.
If the bottom boundary configuration (regarded as a word in the symbols $``$ and $``$, in the familiar sense) of the skeleton $`K`$ matches the top boundary configuration of a skeleton $`L`$, then an operation $`:A_n^{ST}(K)\times A_m^{ST}(L)A_{n+m}^{ST}(KL)`$ is obviously defined, and is extended to the completions.
###### Definition 3.0.9.
For some skeleton $`SK`$ let
$$\gamma :A(SK)A^{ST}(SK),$$
be the mapping defined by linearly extending the operation of mapping the element represented by some diagram in $`A(SK)`$ to the element that that diagram represents in $`A^{ST}(SK)`$.
### 3.1. Some notation
A coupon labelled by a polynomial represents an element of $`A^{ST}`$ via the following expansion. At this stage this is best regarded as a notation; later a space will be introduced ($`B^{QST}(X)`$, Definition 4.1) which will admit such labels in its definition.
If the polynomial is
$$p(t)=p_0+p_1t\mathrm{}+p_nt^n$$
then the following expansion is to be understood.
$$\text{}=p_0\text{}+p_1\text{}+\mathrm{}+p_n\text{}$$
It will also prove helpful to have the following diagrammatic. An oriented edge, labelled with a polynomial, indicates that a coupon with that label is to be introduced, in the sense just introduced. The orientation of the coupon is specified by the orientation of the edge.
$$\text{}\text{}$$
### 3.2. The invariant
###### Definition 3.2.1.
A 4-tuple $`(A,B,w_1,w_2)`$, where
* $`A`$ and $`B`$ are q-tangles,
* the bottom boundary word of $`A`$ is equal to the top boundary word of $`B`$ is equal to $`(w_1)(w_2)`$,
* the top boundary word of $`A`$ is equal to the bottom boundary word of $`B`$ is equal to $`(\mathrm{}())\mathrm{})`$,
is called a presentation for $`T`$, a $`\mu `$-string string link in the solid torus, if the result of composing $`A`$ with $`B`$, while drilling the hole at some point on the mutual bounding line between $`w_1`$ and $`w_2`$, gives $`T`$. It is clear that every string link in the solid torus has such a presentation.
For example, the string link in the solid torus associated to the previously considered surgery presentation of the figure of 8 knot, has the following presentation:
$$(\text{},\text{},,)$$
For a boundary word $`w`$, let $`G_w`$ be the winding diagram obtained from the identity diagram $`I_w`$ by attaching a winding coupon to each strand. This notation can be read as the gluing diagram. For example,
$$G_{()}=\text{}.$$
###### Definition 3.2.2.
Let $`^\mu `$ denote the skeleton underlying a $`\mu `$-string string link.
###### Definition 3.2.3.
If $`T`$ is a $`\mu `$-string string link in the solid torus, then let
(3.2.1)
$$Z^{ST}(T)=\gamma (\widehat{Z}(A_T))(I_{w_1}G_{w_2})\gamma (\widehat{Z}(B_T))A^{ST}(^\mu ),$$
where $`(A_T,B_T,w_1,w_2)`$ is a presentation for $`T`$.
The most pressing issue is, of course, to show that this is well-defined. For the time being, then, indicate the dependence on the presentation $`Z^{ST}(A_T,B_T,w_1,w_2)`$. The well-definedness will follow from the following, clear, observation.
###### Lemma 3.2.4.
Let $`AA^{ST}(SK)`$, take a word $`w`$ such that $`G_wA`$ is well-defined, and take a word $`w^{}`$ such that $`AG_w^{}`$ is well-defined. Then
$$G_wA=AG_w^{}.$$
###### Lemma 3.2.5.
$`Z^{ST}(A_T,B_T,w_1,w_2)`$ is independent of the choice of presentation, and hence is an invariant of $`T`$.
Proof.
It is clear that any two presentations can be related by a finite sequence of the following moves:
$$\begin{array}{cccc}(1)\hfill & (A(1_{w_1}C),B,w_1,w_2)\hfill & \hfill & (A,(1_{w_1}C)B,w_1,w_2^{}),\hfill \\ (2)\hfill & (A(C1_{w_2}),B,w_1,w_2)\hfill & \hfill & (A,(C1_{w_2})B,w_1^{},w_2).\hfill \end{array}$$
The lemma then follows from the functoriality of $`\widehat{Z}`$ and Lemma 3.2.4.
Thus we revert to the notation $`Z^{ST}(T)`$.
The normalisation of this invariant that is appropriate for surgery considerations is the following. This is the normalisation of \[LMMO\].
###### Definition 3.2.6.
Let $`T`$ be an $`\mu `$-string string link in a solid torus. Define:
$$\stackrel{ˇ}{Z}^{ST}(T)=\gamma ((\underset{\mu }{\underset{}{\nu \mathrm{}\nu }})\mathrm{\Delta }^{\mu 1}(\nu ))Z^{ST}(T),$$
in the space $`A^{ST}(^\mu )`$, recalling that $`\nu =\widehat{Z}(U)A()`$.
### 3.3. The co-product
We now equip $`A^{ST}(K)`$ with a co-product. The presence of winding coupons does not affect the following familiar definition.
###### Definition 3.3.1.
Take a diagram $`D`$ such that its dashed graph has connected components indexed by the set $`I`$. If $`JI`$ let $`D_J`$ indicate the diagram obtained by forgetting those components in the subset $`J`$. Then, define the mapping $`\mathrm{\Delta }`$ as the linear extension of
$$\mathrm{\Delta }(D)=\underset{JI}{}D_JD_{IJ}.$$
###### Remark 3.3.2.
1. If $`D`$ has an empty dashed graph then this operation is defined to be $`\mathrm{\Delta }(D)=DD`$.
2. This defines a co-product on the graded completions:
$$\mathrm{\Delta }:A^{ST}(K)A^{ST}(K)\widehat{}A^{ST}(K).$$
###### Lemma 3.3.3.
This is well-defined, co-commutative, co-associative, and commutes with compositions.
To see that it is well-defined we must show that relations are mapped to relations. The only novelty is a PUSH relation when two of the involved edges are part of the skeleton; this relation is easily checked. Observe that it commutes with compositions by construction. All other properties are standard.
###### Lemma 3.3.4.
For a string link in a solid torus $`T`$,
$$\mathrm{\Delta }(\stackrel{ˇ}{Z}^{ST}(T))=\stackrel{ˇ}{Z}^{ST}(T)\stackrel{ˇ}{Z}^{ST}(T).$$
Proof. This follows for the usual reasons: that is, from the corresponding property for $`\widehat{Z}`$, the corresponding property for the normalisation factors of Definition 3.2.6, from the obvious property that $`\mathrm{\Delta }(G_w)=G_wG_w`$ and from the commutation of composition with the co-product.
### 3.4. The Hopf algebra $`B^{ST}(X)`$
We turn to the case of special string links in the solid torus. These are, remember, string links in the solid torus such that a representative in general position with respect to the meridional disc has algebraic intersection zero with it. For such an $`\mu `$-string string link in the solid torus $`T`$, $`\stackrel{ˇ}{Z}^{ST}(T)`$ lies in a special subspace of $`A^{ST}(^\mu )`$.
###### Definition 3.4.1.
Let $`A^{ST,spec}(^\mu )`$ denote the subspace of $`A^{ST}(^\mu )`$ spanned by diagrams with the property that the product of all the labels on the winding coupons labelling some component of the skeleton is 1 (that is, using a factor of $`t^1`$ if some coupon is oriented against the orientation of that component); for each component.
###### Observation 3.4.2.
If $`T`$ is a special $`\mu `$-string string link in the solid torus then
$$\stackrel{ˇ}{Z}(T)A^{ST,spec}(^\mu ).$$
###### Remark 3.4.3.
If a diagram is in this subspace, then repeated applications of the $`PUSH`$ relation can be used to make all of the labels on the skeleton 1 (say, by pushing all the labels to one end). Then all the labels (all the winding) will be carried by the dashed graph.
We now introduce an isomorphic description of this subspace. This is an enhancement of the familiar algebra $`B`$. Let $`X=\{x_1,\mathrm{},x_\mu \}`$ denote a labelling set for the skeleton $`^\mu `$.
###### Definition 3.4.4.
Let a winding diagram on $`X`$ be a graph with oriented trivalent vertices, winding coupons, and univalent vertices labelled from $`X`$.
We may alternatively call this a symmetrised winding diagram on $`^\mu `$.
###### Definition 3.4.5.
Define
$$B_m^{ST}(X)=\frac{\left\{\begin{array}{c}\text{Finite }\text{-linear combinations of degree }m\\ \text{winding diagrams on }X\end{array}\right\}}{\text{span of AS, IHX, OR, MULT and PUSH relations}}$$
Let $`B^{ST}(X)`$ denote the graded completion of $`_{m=0}^{\mathrm{}}B_m^{ST}(X)`$. Equip this with the obvious analogs of the “disjoint-union” product, the “sum over partitions into two sets” co-product, the “empty set” unit and co-unit, and the “$`(1)`$ for every component” antipode.
###### Lemma 3.4.6.
$`B^{ST}(X)`$ is a commutative, co-commutative Hopf algebra.
###### Definition 3.4.7.
Let $`\chi :B^{ST}(X)A^{ST,spec}(^\mu )`$ be the operation defined on some symmetrised diagram of taking the average of all diagrams obtained by locating all univalent vertices labelled with $`x_1`$ on the first component, etc. ; linearly extended to each $`B_m^{ST}`$, and to $`B^{ST}`$.
###### Lemma 3.4.8.
The mapping $`\chi `$ describes a $``$-vector space isomorphism at each grade
$$B_m^{ST}(X)A_m^{ST,spec}(^\mu ),$$
commuting with coproducts
$$(\chi \chi )\mathrm{\Delta }=\mathrm{\Delta }\chi .$$
The inverse, $`\sigma :A_m^{ST,spec}(^\mu )B_m^{ST}(X)`$, is obtained by first pushing all the winding coupons onto the dashed graph (say by pushing all coupons on some component to the top of that component), and then applying Bar-Natan’s “formal PBW map” \[BN\]. It is straightforward to check (adapting \[BN\]) that these maps are well-defined and inverses of each other.
The conclusion of this development follows:
###### Lemma 3.4.9.
Let $`T`$ be a special $`\mu `$-string string link in the solid torus. Then $`\sigma (\stackrel{ˇ}{Z}^{ST}(T))`$ is a group-like element in the Hopf algebra $`B^{ST}(X)`$. Thus it is an exponential of a series of connected diagrams, a finite $``$-linear combination at each grade.
To see this, note that Lemma 3.3.4 indicates that $`\stackrel{ˇ}{Z}^{ST}(T)`$ is group-like in $`A^{ST,spec}`$. Thus its image in $`B^{ST}(X)`$ is also group-like because of the commuting of the co-product with the map $`\sigma `$. Thus it is an exponential of a primitive element (for example, \[Qu\], Appendix A): at each grade this will be a finite $``$-linear combination of connected diagrams.
### 3.5. The winding matrix
We now introduce $`W(T,t)M_\mu ([t,t^1]),`$ the winding matrix of $`T`$, where $`T`$ is a $`\mu `$-string string link in the solid torus. Number the components of $`T`$ and choose a diagram for $`T`$ that is in general position with respect to the meridional disc. We consider paths on this diagram. The “algebraic intersection of a path with the meridional disc” is the sum over all crossings of that path with the disc of: a plus one if the tangent vector of the path points in the direction of increasing $`y`$ at the intersection; and a minus one otherwise (according to the model of the solid torus described in Section 2.1).
###### Definition 3.5.1.
For a crossing $`c`$ of strands $`i`$ and $`j`$, let $`ϵ(i,j,c)`$ denote the algebraic intersection with the meridional disc of the path obtained by travelling from the base along $`i`$ to $`c`$, crossing to $`j`$, and then travelling to the top along $`j`$. (To be precise, if $`i=j`$ then change strands the first time the crossing is encountered).
###### Definition 3.5.2.
For $`T`$, an $`\mu `$-string string link in the solid torus, choose a (blackboard-framed) diagram for $`T`$ in general position with respect to the meridional disc. If $`c`$ denotes a crossing then let $`sgn(c)`$ denote the sign of that crossing. Let:
$$W_{ij}(T,t)=\{\begin{array}{cc}_{\text{c, a crossing of i and j}}\frac{1}{2}\text{sgn}(c)t^{ϵ(i,j,c)}\hfill & \text{if}ij,\hfill \\ & \\ _{\text{c, a self-crossing of }i}\frac{1}{2}\text{sgn}(c)(t^{ϵ(i,i,c)}+t^{ϵ(i,i,c)})\hfill & \text{otherwise.}\hfill \end{array}$$
###### Remark 3.5.3.
* The winding matrix is Hermitian: $`W_{ij}(T,t)=W_{ji}(T,t^1)`$.
* The winding matrix specialises to the linking matrix of the underlying tangle: $`W_{ij}(T,1)=\text{Lk}_{ij}(T)`$.
###### Lemma 3.5.4.
$`W_{ij}(T,t)`$ is an isotopy invariant of $`T`$, regarded as a framed string link in a solid torus.
We will return to the topological interpretation of $`W_{ij}`$ at the end of this section when we make an important connection with the Alexander polynomial. Besides, its isotopy invariance follows from its appearance in the following theorem (the unique Hermitian matrix with the following property).
Let $`X=\{x_1,\mathrm{},x_\mu \}`$ denote a labelling set for $`T`$.
###### Theorem 3.5.5.
If $`T`$ is a string link in a solid torus then
$$\sigma (\stackrel{ˇ}{Z}^{ST}(T))=\text{exp}_{}(\frac{1}{2}\underset{i,j}{}\text{})R,$$
where $`R`$ is a series of $`X`$-substantial diagrams.
Proof.
Lemma 3.4.9 indicated that $`\sigma (\stackrel{ˇ}{Z}^{ST}(T))`$ is of the form exp$`{}_{}{}^{}(S)`$ where $`S`$ is a series of connected diagrams. Thus to prove the theorem it is sufficient to calculate the degree one part of $`\stackrel{ˇ}{Z}^{ST}`$.
Take a diagram of $`T`$ in general position with respect to the meridional disc, and an associated presentation of $`T`$. Examining the definition of $`\stackrel{ˇ}{Z}^{ST}`$, we see that there will be a contribution of one chord from every crossing in this diagram. The introduced complexity is that there will be a distribution of winding coupons on the skeleton. (The reader is invited to consider the example that follows this proof).
Consider first a crossing between two different components, $`i`$ and $`j`$. Use the relation $`PUSH`$ to push the coupons onto the chord in the following manner: coupons that occur before the crossing as $`i`$ is traversed from the base should be pushed past the chord (where they will all cancel to 1); coupons that occur after the crossing as $`j`$ is traversed to the top should be pushed back past the chord (where they will all cancel to 1).
The reader can check that the label on the chord coming from that crossing is precisely as follows. Note that in this section we will use the notational convention described in Section 3.1.
Now consider self-crossings of components. Choose some self-crossing $`c`$ of some component $`x_i`$. Its contribution is as follows:
$$\frac{1}{2}\text{sgn}(c)\text{}$$
Now, use the following relation:
$$\text{}=\text{}\text{}$$
to write the contribution from this crossing as follows:
$$\chi (\frac{1}{2}\text{sgn}(c)(\text{}\frac{1}{2}\text{}))$$
In other words:
$$\chi (\frac{1}{2}\left(\frac{1}{2}\text{sgn}(c)(\text{}+\text{})\right)+r),$$
where $`r`$ is an $`X`$-substantial diagram.
###### Example 3.5.6.
Consider the example associated with the presentation of the figure of 8 knot.
$$\stackrel{ˇ}{Z}^{ST}(\text{})=1+\text{}\frac{3}{2}\text{}+r,$$
$$=1+\text{}+r^{},$$
where $`r`$ and $`r^{}`$ are series of diagrams that are either of grade greater than 1 or X-substantial.
### 3.6. A topological interpretation of $`W(T,t)`$
The attentive reader will have noticed the appearance of the Alexander polynomial for the figure of 8 knot in the previous calculation. Let us examine the meaning of the matrix $`W(T,t)`$ more closely.
Choose a base point close to the base of the strings, and choose paths from that basepoint to the bases of the strings (so that the bases, paths, and basepoint all lie in some ball). Then close the string link on the left, obtaining some link in the solid torus $`L`$, say with components $`\{K_1,\mathrm{},K_\mu \}`$, with a path from some basepoint to some point on each component.
Take the universal cyclic cover of the solid torus:
$$p:\stackrel{~}{ST}ST,$$
and lift the link $`\{K_1,\mathrm{},K_\mu \}`$ to $`\stackrel{~}{ST}`$. This can be done as we are restricting to special string links (that is, string links whose algebraic intersection with any meridional disc is zero). The group of translations is $``$: choose an action such that a path which starts at some point $`p`$; crosses the meridional disc in the direction on increasing $`y`$; and the returns to the $`p`$ (without again crossing the meridional disc) is lifted to a path starting at some $`p`$ and finishing at some $`tp`$. The lifted link $`\stackrel{~}{L}`$ can be identified as the set of translates
$$\{\mathrm{},t^1\stackrel{~}{K_1},\stackrel{~}{K_1},t\stackrel{~}{K_1},\mathrm{},t^1\stackrel{~}{K_\mu },\stackrel{~}{K_\mu },t\stackrel{~}{K_\mu },\mathrm{}\},$$
where the components $`\{\stackrel{~}{K_1},\mathrm{},\stackrel{~}{K_\mu }\}`$ can be fixed by choosing $`\stackrel{~}{K_1}`$ and then following the lifts of the arcs introduced when the closure was taken.
We define an invariant of $`T`$ using this arrangement, as follows.
###### Definition 3.6.1.
Let $`\stackrel{~}{\text{Lk}}(T)M_\mu ([t,t^1])`$ be defined by
$$\stackrel{~}{\text{Lk}}_{ij}(T)=\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}t^mlk(\stackrel{~}{K}_i,t^m\stackrel{~}{K}_j).$$
###### Lemma 3.6.2.
$$W(T,t)=\stackrel{~}{\text{Lk}}(T).$$
Proof.
Take a blackboard-framed diagram for $`T`$ that is in general position with respect to the projection of the meridional disc. A fundamental domain for a diagram for the pair $`(\stackrel{~}{ST},\stackrel{~}{L})`$ can be obtained by cutting such a diagram along the projection of the meridional disc, obtaining a rectangle, and then gluing a countable infinity of copies of that rectangle in the appropriate way.
Observe that a crossing between two different components, say $`x_i`$ and $`x_j`$, in the diagram for $`T`$ lifts to a crossing between $`K_i`$ and some translate of $`K_j`$. Which translate is decided by counting intersections with the meridional disc as in the definition of $`W`$.
The slightly different definition of $`W`$ along the diagonal is accounted for by the fact that self-crossings of a component $`x_i`$ in the diagram for $`T`$ will either lift to a self-crossing of $`K_i`$, or to a pair of crossings, between $`K_i`$ and some $`t^aK_i`$ and between $`K_i`$ and $`t^aK_i`$. Observe that the factors give the right weights in both situations.
###### Example 3.6.3.
Continuing the example of the figure of 8 knot:
Take $`T`$, a special string link in the solid torus, presenting some pair $`(M,K)`$. Remember that, according to Definition 1.0.3, $`A_{(M,K)}(t)`$ denotes the canonical Alexander polynomial of a knot in a $`HS^3`$.
###### Lemma 3.6.4.
$$A_{(M,K)}(t)=\pm \text{det}(W(T,t)).$$
Proof.
Surgery on the framed link $`\stackrel{~}{L}`$ in $`D^2\times `$ recovers the universal cyclic cover of the complement of $`K_T`$.
The Mayer-Vietoris sequence then indicates that the matrix $`\stackrel{~}{L}`$, and hence $`W(T,t)`$, is a presentation matrix for the $`[t]`$-module $`H_1(\stackrel{~}{MK};)`$.
The Alexander polynomial is defined as a generator of the order ideal of that module. In the situation at hand, given a square presentation matrix, this is calculated by the determinant of that matrix, as in the statement of the lemma. Note that this only specifies the polynomial up to multiplications of the form $`\pm t^n`$, and the statement of the lemma asserts that the recovered polynomial is symmetric under $`tt^1`$.
To see that this is true note that $`W(T,t)`$ is a Hermitian matrix, according to Remark 3.5.3.
$$\text{det}(W(T,t))=\text{det}(W(T,t)^{Tr.})=\text{det}(\overline{W(T,t)})$$
###### Remark 3.6.5.
$`W(T,t)`$ is, presumeably, the matrix referred to in Exercise C.13 of Rolfsen \[Rol\].
### 3.7. Alternative normalisations
We take this oppurtunity to draw attention to a certain subtle choice of normalisation that has been made in the construction of $`Z^{ST}`$ given here.
Let $`\alpha `$ denote a group-like element of $`A()`$, the space of uni-trivalent diagrams on a single strand.
###### Definition 3.7.1.
Let
$$Z^{ST}[\alpha ]:\left\{\begin{array}{c}Stringlinks\\ inthesolidtorus.\end{array}\right\}A^{ST}(^\mu ).$$
be defined in exactly the same way as given in Definition 3.2.3, except that in Equation 3.2.1 $`I_{\omega _1}G_{\omega _2}`$ should be replaced by
$$I_{w_1}(\mathrm{\Delta }_{\omega _2}(\alpha )G_{\omega _2}),$$
where $`\mathrm{\Delta }_{\omega _2}(\alpha )`$ is the paralleling operation across the strands described by $`\omega _2`$, applied to $`\alpha `$. Let $`\stackrel{ˇ}{Z}^{ST}[\alpha ]`$ denote the normalisation corresponding to Definition 3.2.6.
###### Remark 3.7.2.
Lemma 3.4.9 and Theorem 3.5.5 still hold if $`\stackrel{ˇ}{Z}^{ST}`$ is replaced by $`\stackrel{ˇ}{Z}^{ST}[\alpha ]`$.
Certain choices of $`\alpha `$ will prove appropriate for certain applications. For example, setting $`\alpha =\nu ^1`$ gives a normalisation that is better adapted to questions involving covering spaces.
Applying $`^{FGinST}𝑑X`$, followed by $`Thr^D`$, presumeably gives knots invariants which are different normalisations of the Kontsevich integral, in some sense. We will return to this question in the sequel.
## 4. Formal Gaussian integration
In this section we focus on the leftmost edge of Diagram 2.2.1.
###### Definition 4.0.1.
An element $`SB^{ST}(X)`$ is said to be integrable if it is of the following form.
(4.0.1)
$$S=\text{exp}_{}\left(\frac{1}{2}\underset{ij}{}\text{}\right)R,$$
1. $`W_{ij}(t)`$ is a Hermitian matrix ($`W_{ij}(t)=W_{ji}(t^1)`$) wih the property that det$`(W(1))=\pm 1`$,
2. $`R`$ is a series of $`X`$-substantial diagrams: that is, a diagram will have no chords (ignoring winding coupons).
###### Remark 4.0.2.
Observe that the above decomposition is unique. In particular, a Hermitian matrix satisfying Equation 4.0.1 (for some given $`S`$) will be unique. The matrix will be called the Gaussian matrix of $`S`$.
###### Definition 4.0.3.
Let $`B^{ST}(X)^{Int}`$ denote the subspace of integrable elements of $`B^{ST}(X)`$.
###### Remark 4.0.4.
As the associated Hermitian matrix is uniquely specified, we will freely apply any function of such matrices to the set $`B^{ST}(X)^{Int}`$ with the understanding that the function is to be applied to the Gaussian matrix of the element. An example of such use is the use of the functions $`\sigma _+`$ and det in the diagram above. These factors will be taken up again in Section 7.
### 4.1. Rational winding diagrams
The Aarhus calculation of the LMO invariant has an associated philosophy of formal Gaussian integration. The idea is that given an integrable element of $`B(X)^{Int}`$ one “integrates” by splitting off the quadratic part, inverting it, and contracting the result with the remainder.
We would like to introduce an analog of this in the situation at hand, that is, for the space $`B^{ST}(X)^{Int}`$. Whilst we have a clear definition for an integrable element, it remains for us to introduce a suitable space in which to “invert” the given matrix of polynomials.
The definition that follows generalises the definition of the space of winding diagrams, Definition 3.0.6. We will differentiate this space in our vocabulary by inserting the adjective “rational”.
###### Remark 4.1.1.
We should point out that for the proof of Theorem 2.2.2, and hence Rozansky’s conjecture, the following definition is unneccessary. The reader who feels the following is an unneccessarily cumbersome space may happily evade this space by observing that, in this work, every appearance of $`^{FGinST}𝑑X`$ is followed by a $`Thr^D`$. (See Remark 5.1). The point of including such a definition here is Conjecture 2.3.1.
###### Definition 4.1.2.
A rational winding coupon is a vertex of some even valency. The incoming edges are ordered up to reorderings of the form $`(\sigma ,\sigma )`$ for some $`\sigma \mathrm{\Sigma }_m`$ (if the valency is $`2m`$). A winding coupon is labelled from $`^1(t),`$ the ring of rational functions in a single variable which are non-singular at 1.
This is depicted as follows, the idea being that an edge leads through the coupon to the opposite edge. The ellipsis used as below will indicate the possibility of a number of other edges. The ordering of the edges is the bottom row, from left to right, followed by the top row, from left to right.
###### Definition 4.1.3.
A rational winding diagram labelled from the set $`X`$ is a graph with univalent vertices, trivalent vertices and winding coupons, such that:
1. (Univalent vertices.) Univalent vertices are labelled from $`X`$.
2. (Trivalent verices.) Trivalent vertices are vertex-oriented (have their incoming edges cyclically ordered).
Two rational winding diagrams are identified if there is a graph isomporhism between them respecting orientations at trivalent vertices, orientations and labels of winding coupons, and labels of univalent vertices.
###### Definition 4.1.4.
The grade of a rational winding diagram is half the number of trivalent vertices.
The space in question will be a quotient of the space of finite $``$-linear combinations of rational winding diagrams of a fixed grade. The quotient will be by the span of the following classes of vectors.
$`AS:`$ $`\text{}+\text{}`$
$`IHX:`$ $`\text{}\text{}\text{}`$
$`STU:`$ $`\text{}\text{}+\text{}`$
$`OR:`$ $`\text{}\text{}`$
$`ADD:`$ $`\text{}a\text{}b\text{}`$
$`MULT:`$ $`\text{}\text{}`$
$`SPLIT:`$ $`\text{}\text{}`$
$`COMM:`$ $`\text{}\text{}`$
$`PUSH1:`$ $`\text{}\text{}`$
$`PUSH2:`$ $`\text{}\text{}`$
###### Definition 4.1.5.
Let $`X`$ be a set of labels. Let
$$B_m^{QST}(X)=\frac{\left\{\begin{array}{c}\text{Finite }\text{-linear combinations of degree }m\hfill \\ \text{rational winding diagrams labelled from }X\hfill \end{array}\right\}}{\text{span of above relations}}$$
Let $`B^{QST}(X)`$ denote the completion of $`_{m=0}^{\mathrm{}}B^{QST}(X)`$ with respect to degree.
### 4.2. Formal Gaussian integration in the solid torus.
We can now introduce the map
$$^{FGinST}𝑑X:B^{ST}(X)^{Int}B^{QST}(\varphi ).$$
###### Definition 4.2.1.
If the unique decomposition of an element $`SB^{ST}(X)^{Int}`$ is
(4.2.1)
$$S=\text{exp}_{}\left(\frac{1}{2}\underset{ij}{}\text{}\right)R,$$
then
(4.2.2)
$$^{FGinST}𝑑XS=\text{exp}_{}\left(\frac{1}{2}\underset{ij}{}\text{}\right),R_XB^{QST}(\varphi ).$$
## 5. Threading
In this section we will introduce $`Thr^D`$, the operation of threading (rational) winding diagrams. This is used in the following square from the master diagram. In this section we will show it commutes.
###### Diagram 5.0.1.
* “Special string link in the solid torus” is defined in Definition 2.1.7, and “Special tangle” is defined in Definition 2.1.8.
* $`Thr`$ is the operation introduced in Definition 2.1.6 which threads the hole in the solid torus with a zero-framed unknot.
* The space $`B^{ST}(X)^{Int}`$ is defined in Definition 4.0.1, and $`B(X,\underset{¯}{k})`$ is recalled in Section 2.4.
* The operation $`Thr^D`$ will presently be introduced.
This operation will be defined on any of the spaces introduced so far which involve (possibly rational) winding coupons. If the labelling information is denoted $`L`$ (so maybe a skeleton and a set of labels, or possibly $`\varphi `$) then the map will be between spaces as follows. The $`Q`$ is in brackets because it may, or may not, be present.
$$Thr^D:B^{(Q)ST}(L)B(L,\underset{¯}{k}).$$
The definition of this map will be introduced via an intermediate construction, a generating diagram.
###### Definition 5.0.2.
A generating diagram is precisely the same as a winding diagram, except that its coupons are labelled with formal power series in a variable $`k`$. A generating diagram denotes a particular series of uni-trivalent diagrams.
1. (One edge) If there is only one edge going through the coupon, then the association is as follows.
If $`f(k)[[k]]`$ is written $`f_0+f_1k+f_2k^2+f_3k^3+\mathrm{}`$ then a coupon labelled as follows, is to be expanded as shown. See below for the meaning of the vertex.
$$\text{}=f_0\text{}+f_1\text{}+f_2\text{}+f_3\text{}\mathrm{}$$
2. (More than one edge) In this case, one takes the above expansion and then, for each diagram, takes the sum of diagrams obtained by lifting each introduced leg to each edge going through the coupon.
This vertex depends on whether the edge is internal or part of the skeleton, and then also on the orientation of the skeleton, as follows.
$$\text{}\text{};\text{}\text{};\text{}(1)\text{}$$
###### Definition 5.0.3.
The operation $`Thr^D`$ is defined by mapping a (rational) winding diagram to the series of uni-trivalent diagrams represented by making the substitution $`te^k`$ in the label of every winding coupon.
###### Lemma 5.0.4.
This is a well-defined operation.
For starters, the subsitution makes sense because we are restricting labels to the subspace of rational functions that are non-singular at $`t=1`$. Furthermore, the relations involving winding coupons (OR, ADD, MULT, SPLIT, COMM, PUSH1 and PUSH2) are easily checked.
###### Theorem 5.0.5.
Diagram 5.0.1 commutes.
This theorem depends crucially on a certain consequence of the Kontsevich integral proof of the “Wheels Conjectures” that has recently been given by Bar-Natan, Le and Thurston \[TW\]. See the forthcoming paper of Bar-Natan and Lawrence \[BNL\] for the following calculation. (See that paper also for references to other proofs of the “Wheeling Conjecture” that are in the literature.)
###### Theorem 5.0.6.
$$\widehat{Z}(\text{})=\text{exp}_{}(\text{})\nu (k)B(x,\underset{¯}{k}).$$
We use the following corollary. The tangle below is equipped with some choice of bracketting which is the same on both the top and the bottom boundary words.
###### Corollary 5.0.7.
$$\widehat{Z}(\text{})=\text{}\nu (k)$$
in the space $`A(\underset{r}{\underset{}{\mathrm{}}}\underset{s}{\underset{}{\mathrm{}}};\underset{¯}{k}).`$
Proof of corollary. This is proved with Le and Murakami’s paralleling formula \[LM2\]. The small point to observe is that an application of the paralleling formula gives
$$\text{exp}_{}(\text{})\mathrm{}\text{exp}_{}(\text{})\text{exp}_{}(\text{})\mathrm{}\text{exp}_{}(\text{})\nu (k).$$
This requires the averaged sum of orderings of legs on the components $`\{x_i\}`$, whereas the given statement requires the composition exponential. But the other ends of these chords are unordered. So the corollary follows as stated, for $`\widehat{Z}`$. (Observe that there are no problems with the choice of associator for this scenario).
Proof of Theorem 5.0.5. Take $`T`$, an $`\mu `$-string string link in the solid torus, given by some presentation $`(A_T,B_T,w_1,w_2)`$. Without loss of generality we can assume that $`w_2`$ is some bracketting of some word of the form $`\underset{r}{\underset{}{\mathrm{}}}\underset{s}{\underset{}{\mathrm{}}}`$.
Then, the functoriality of $`\widehat{Z}`$ indicates that $`\stackrel{ˇ}{Z}(Thr(T))`$ is equal to the following expression, in the space $`A(^\mu ,\underset{¯}{k}).`$
Alternatively, $`Thr^D(\stackrel{ˇ}{Z}^{ST}(T))\nu (k)`$ may be calculated, as follows.
$$Thr^D\left((\underset{\mu }{\underset{}{\nu \mathrm{}\nu }})\mathrm{\Delta }^\mu (\nu )\gamma (\widehat{Z}(A_T))(I_{w_1}G_{w_2})\gamma (\widehat{Z}(B_T))\right)\nu (k).$$
This is exactly the same thing, in the space $`A(^\mu ,\underset{¯}{k})`$. We have just proved the commutation of the top half of the following diagram. The precise statement of the theorem then follows from the commutation of the bottom half.
### 5.1. Evading $`B^{QST}`$
The reader who wishes to evade the space $`B^{QST}`$ may do so by redefining the map $`Thr^D^{FGinST}`$ as follows :
###### Alternative Definition 5.1.1.
If the unique decomposition of an element $`SB^{ST}(X)^{Int}`$ is
(5.1.1)
$$S=\text{exp}_{}\left(\frac{1}{2}\underset{ij}{}\text{}\right)R,$$
then
(5.1.2)
$$Thr^D\left(^{FGinST}𝑑X(S)\right)=\text{exp}_{}\left(\frac{1}{2}\underset{ij}{}\text{}\right),Thr^D(R)_XB(\underset{¯}{k}).$$
## 6. The LMO invariant
The LMO invariant was introduced by Thang Le, Jun Murakami and Tomotada Ohtsuki \[LMO\]. We will specialise our presentation to the setting in question, knots in integral homology three-spheres.
Let the pair of a knot $`K`$ in an integral homology three sphere $`M`$ be presented by some special framed tangle $`T`$. That is, $`T`$ has one closed component, and forgetting that closed component leaves the tangle a string link $`T^{}`$ whose components are to be surgered, after closure; the determinant of the linking matrix of those components is $`\pm 1`$. Let $`X=\{x_1,\mathrm{},x_\mu \}`$ be an index set for $`T^{}`$, and let $`lk(T^{})`$ denote the linking matrix of those components.
The LMO invariant is constructed as a sequence
$$Z_n^{LMO}(M,K)B_n(\underset{¯}{k})$$
with the property that
(6.0.1)
$$\text{Grad}_n(Z_m^{LMO}(M,K))=Z_n^{LMO}(M,K)B_n(\underset{¯}{k})$$
when $`mn`$. This sequence can thus be regarded as approximations to an invariant defined by
$`Z^{LMO}(M,K)`$ $`=`$ $`1+\text{Grad}_1(Z_1^{LMO}(M,K))+\text{Grad}_2(Z_2^{LMO}(M,K))+\mathrm{}B(\underset{¯}{k}).`$
It is convenient to introduce the operation that is the heart of this definition over a more general space.
###### Definition 6.0.1.
For some labelling set $`L`$, the space $`B^o(L)`$ is defined with exactly the same definition as $`B(L)`$, except that generating diagrams may also have a finite number of closed dashed loops.
###### Remark 6.0.2.
The space $`B(L)`$ clearly injects into $`B^o(L)`$.
###### Definition 6.0.3.
The mapping
$$^{(n)}𝑑X:B^o(X,X^{},\underset{¯}{k})B_n(X^{},\underset{¯}{k}),$$
is defined on a diagram $`D`$ as
$$\underset{i=1}{\overset{\mu }{}}\left(\frac{1}{n!}\left(\frac{1}{2}\text{}\right)^n\right),D,$$
followed by the exchange of each dashed loop component (some extra may arise) for a multiplicative factor of $`2n`$.
With this operation, $`Z_n^{LMO}(M,K)`$ is defined as follows. Let $`\sigma _\pm (lk(T^{}))`$ denote the number of positive (resp. negative) eigenvalues of $`lk(T^{})`$. Let $`U_\pm `$ be a $`\pm 1`$-framed unknot.
###### Definition 6.0.4.
The invariant $`Z_n^{LMO}(M,K)`$ is defined by
$$Z_n^{LMO}(M,K)=\frac{^{(n)}𝑑X\sigma (\stackrel{ˇ}{Z}(T))}{\left(^{(n)}𝑑U\sigma (\stackrel{ˇ}{Z}(U_+))\right)^{\sigma _+(lk(T^{}))}\left(^{(n)}𝑑U\sigma (\stackrel{ˇ}{Z}(U_{}))\right)^{\sigma _{}(lk(T^{}))}},$$
in the space $`B_n(\underset{¯}{k})`$.
###### Notation 6.0.5.
Let the numerator in the expression above be denoted $`Z_n^{LMO,o}(T)`$. Note that this $`o`$ is of a different nature to the $`o`$ used in Definition 6.0.1 above.
###### Theorem 6.0.6.
\[LMO, LeD\] This definition gives a well-defined invariant of knots in integral homology three-spheres, assmebled via Formula LABEL:LMOassem, with the specialisation
$$Z^{LMO}(S^3,K)=\widehat{Z}(K).$$
###### Remark 6.0.7.
Observe that our normalisation of the non-surgered components does not affect surgery formula.
## 7. The surgery formula
This section is the crux of the paper. It will prove that the following square from the master diagram commutes.
###### Diagram 7.0.1.
There are some components of this that are yet to be introduced.
###### Definition 7.0.2.
Let $`^1[t,t^1]`$ be the subring of $`[t,t^1]`$ of polynomials such that if $`p(t)^1[t,t^1]`$ then
1. $`p(t)=p(t^1),`$
2. $`p(1)=\pm 1.`$
Clearly $`det`$, evaluated on $`B^{ST}(X)^{Int}`$ (according to Remark 4.0.4) takes values in this ring.
###### Definition 7.0.3.
The mapping $`\sigma _+`$ on $`SB^{ST}(X)^{Int}`$ is the number of positive eigenvalues of $`W(1)`$, where $`W(t)`$ is the Gaussian matrix of $`S`$.
###### Definition 7.0.4.
The mapping
$$Wh^{}:^1[t,t^1]B(\underset{¯}{k}),$$
is defined by
$$Wh^{}(P(t))=\text{exp}_{}\left(\left[\frac{1}{2}\text{log}\left(\frac{P(e^h)}{P(1)}\right)\right]|_{h^{2n}\omega _{2n}}\right)\nu (k),$$
where the operation indicated is to expand the term inside the square brackets into a power series in $`h`$, and then to replace terms like $`ch^{2n}`$ by $`c\omega _{2n}`$, in exactly the same fashion as Definition 1.0.7.
Observe that the $`P(1)`$ factor just adjusts the sign of the polynomial.
###### Remark 7.0.5.
$$Wh(M,K)=Wh^{}(A_{(M,K)}(t)).$$
### 7.1. Translation by power series
Before we consider the details of this proof, we introduce a certain operation on diagrams. For $`F(x_1,\mathrm{},x_\mu ,k)B(X,\underset{¯}{k})`$, the notation $`F(x_1+z,x_2,\mathrm{},x_\mu ,k)`$, according to Aarhus, denotes the series in $`B(X,\underset{¯}{k},z)`$ obtained by replacing every diagram with $`l`$ legs labelled by $`x_1`$ by the $`2^l`$ diagrams obtained by relabelling each such leg with either $`x_1`$ or $`z`$. This may be extended linearly to simultaneous “translations” of other variables.
Now we introduce something novel. If $`f(k)[[k]]`$ then $`F(f(k)x_1,x_2,\mathrm{},x_\mu )`$ denotes the element obtained by replacing every diagram in the expansion of $`F`$ with a generating diagram obtained by labelling (simultaneously) every leg marked $`x_1`$ as follows. Observe that the added coupon is oriented according to the position of the univalent vertex.
$$\text{}\text{}$$
We will want to combine these operations. For example, take some $`f(k)[[k]]`$:
$$F(x_1+f(k)z,x_2,\mathrm{},x_\mu ,k)=F(x_1,x_2,\mathrm{},x_\mu ,k)+F(f(k)z,x_2,\mathrm{},x_\mu ,k).$$
###### Notation 7.1.1.
Take a matrix $`M(k)\text{M}_\mu ([[k]])`$. The notation $`F(\overline{x}+M(k)\overline{x}^{},k)`$ represents the element
$$F(x_1+\underset{i_1}{}M_{1i_1}(k)x_{i_1}^{},x_2+\underset{i_2}{}M_{2i_2}(k)x_{i_2}^{},\mathrm{},x_\mu +\underset{i_\mu }{}M_{\mu i_\mu }(k)x_{i_\mu }^{},k)$$
in $`B(X,X^{},\underset{¯}{k})`$.
### 7.2. Translation invariance
The core of our proof (our adaption of \[A3\]) is the following property.
###### Theorem 7.2.1.
Take some matrix of power series $`M(k)M_\mu ([[k]])`$.
$$^{(n)}𝑑XF(\overline{x},k)=^{(n)}𝑑XF(\overline{x}+M(k)\overline{x}^{},k)B_n(X^{},\underset{¯}{k}).$$
The proof of this lemma requires a certain subspace of $`B^o(X,X^{},\underset{¯}{k})`$. This is the subspace generated by $`C_n`$ vectors, which were introduced in \[A3\].
###### Definition 7.2.2 ($`C_n`$ vectors).
Consider some uni-trivalent diagram drawn to include some box with $`n`$ dashed edges attached (in an ordered fashion) to its top, and $`n`$ dashed edges attached (ditto.) to its base. If the box is labelled by some permutation $`\sigma \mathrm{\Sigma }_n`$, then this diagram is defined to represent the diagram obtained by joining up the edges according to $`\sigma `$. A $`C_n`$ relation is a linear combination of diagrams obtained from such a diagram by summing over all the diagrams obtained by labelling that box with all possible permutations.
###### Lemma 7.2.3 (\[LMO, A3\]).
A $`C_m`$ vector is in the kernel of $`^{(n)}𝑑X`$, if $`m2n+1`$.
###### Remark 7.2.4.
This is a slightly different viewpoint than is taken in the literature. LMO works with a different relation $`P_{n+1}`$. This was shown to be equivalent to $`C_{2n+1}`$ in \[A3\]. LMO introduces the relation $`P_{n+1}`$ and then shows that it is “redundant” in the image (that is, the quotient by diagrams of grade greater than $`n`$). Here we are in a more general situation than is explicitly found in the literature: namely, we are allowing non-surgered univalent vertices. The proof that $`P_{n+1}`$ is “redundant”, which is accessibly described in Section 2.5.4 of \[LeGr\], adapts immediately to this generality. Following the discussion there, one sees that every leg on a $`P_{n+1}`$ must still meet a seperate vertex.
Proof of Theorem 7.2.1.
Consider some diagram $`D(\overline{x},k)`$ appearing in the expansion of $`F(\overline{x},k)`$. There are three cases which cover the possibilities. Case 1: there is some component $`x_i`$ which has less than $`2n`$ legs labelled with it. In this case both $`D(\overline{x},k)`$ and $`D(\overline{x}+M\overline{x}^{},k)`$ are in the kernel of $`^{(n)}𝑑X`$. Case 2: each component of $`X`$ has exactly $`2n`$ legs labelled with it. In this case there is exactly one contributing diagram on the right hand side, because if any leg $`x_i`$ is relabelled $`M_{ij}x_j^{}`$ then there will then be less than $`2n`$ legs labelled with that component (such a diagram will be in the kernel.) The one contributing term is the same as that obtained from the left hand side. Case 3: some component has more than $`2n`$ legs labelled with it. This vanishes on the left hand side by definition. This also vanishes on the right hand side because this is expressible as a series of $`C_m`$ vectors, for some $`m2n+1`$.
### 7.3. Diagram 7.0.1 commutes
Take some element $`S(\overline{x})`$ of $`B^{ST}(X)^{Int}`$, with canonical decomposition:
$$S(\overline{x})=\text{exp}_{}\left(\frac{1}{2}\underset{ij}{}\text{}\right)R(\overline{x}).$$
We shall calculate $`^{(n)}𝑑X(Thr^D(S(\overline{x}))\nu (k))`$. All expressions in the following are to be evaluated in the space $`B_n(X)`$.
The first step is to split off the remainder, $`Thr^D(R(\overline{x}))`$. Denote this $`R^{}(\overline{x},k)`$. We can do the following because $`R^{}(\overline{x},k)`$ is a series of X-substantial diagrams.
Theorem 7.2.1 tells us that making the “translation”
$$x_ix_i\underset{j}{}W_{ij}^1(e^k)x_j^{}$$
inside the integrand will not affect the result. \[The matrix $`W(T,e^k)`$ is invertible in $`M_\mu ([[k]])`$ by assumption (Definition 4.0.1.) \]
Making that substitution transforms the integral into the following. \[Note that this requires the property that $`W^1(T,e^k`$) is Hermitian, which is again by assumption (Definition 4.0.1.)\]
$$^{(n)}𝑑X\text{exp}_{}\left(\frac{1}{2}\underset{i,j}{}\text{}\frac{1}{2}\underset{i,j}{}\text{}\right).$$
Substituting this expression back into the above gives the following.
(7.3.1)
evaluated in $`B_n(\underset{¯}{k})`$. The theorem follows from the calculation of the leading term that is performed in the following section.
### 7.4. The wheels line
Examining the structure of Equation 7.3.1, we see that the projection onto the wheels line of the element $`^{(n)}𝑑X(Thr^D(S(\overline{x}))\nu (k))`$ is precisely the following function of the entries of the matrix $`W(t)`$.
It is calculated as follows.
###### Theorem 7.4.1.
Let $`W(t)`$ be a Hermitian matrix satisfying det$`(W(1))=\pm 1`$. Then
$$\left(\nu (k)^{(n)}𝑑X\text{exp}_{}\left(\frac{1}{2}\underset{i,j}{}\text{}\right)\right)=(1)^{n\sigma _+(W(1))}Wh^{}(\text{det}(W(t))).$$
Proof.
We begin by observing that every such matrix $`W(t)`$ can be realised as the winding matrix of some special string link in the solid torus $`T`$, presenting some pair $`(M_T,K_T)`$. The calculation is performed by calculating the wheels line of $`Z^{LMO}(M_T,K_T)`$ in two different ways.
On the one hand, this “wheels line” has already been calculated by other means. Many authors contributed to this result. Let us highlight the original conjecture of Melvin–Morton \[MM\] and the Kontsevich integral proof given by Bar-Natan–Garoufalidis \[BNG\] (acknowledging other contemporaneous proofs \[RMM, KM\]). See Garoufalidis–Habegger \[GH\] for the following formula in the setting of $`HS^3`$s (recently extended to null-homologous knots in $`HS^3s`$, \[L\]).
###### Theorem 7.4.2.
Let $`T`$ be a $`\mu `$-string string link in the solid torus presenting some pair $`(M_T,K_T)`$. Then,
$$Z^{LMO}(M_T,K_T)=Wh^{}(\text{det}(W(T,t)))(1+R),$$
where $`R`$ is a series of diagrams whose dashed graphs have negative Euler charcteristic.
On the other hand, we have just seen (Equation 7.3.1) that $`Z^{LMO}(M_T,K_T)`$ may be calculated
$$\frac{[Wheelsbit][Therest]}{\left(^{(n)}𝑑U\sigma (\stackrel{ˇ}{Z}(U_+))\right)^{\sigma _+(W(1))}\left(^{(n)}𝑑U\sigma (\stackrel{ˇ}{Z}(U_{}))\right)^{\sigma _{}(W(1))}}$$
in the space $`B_n(\underset{¯}{k})`$, where
$`[Wheelsbit]`$ $`=`$ $`\left(\nu (k){\displaystyle ^{(n)}}𝑑X\text{exp}_{}\left({\displaystyle \frac{1}{2}}{\displaystyle \underset{i,j}{}}\text{}\right)\right),`$
$`[Therest]`$ $`=`$ $`\text{exp}_{}\left({\displaystyle \frac{1}{2}}{\displaystyle \underset{i,j}{}}\text{}\right),R^{}(\overline{x},k).`$
Given that \[LMO\]
$$^{(n)}𝑑U\sigma (\stackrel{ˇ}{Z}(U_{\pm 1}))=(1)^n+\text{terms of negative Euler characteristic,}$$
the theorem follows by equating wheels lines.
###### Exercise 7.4.3.
Prove Theorem 7.4.2 directly. That is, calculate
$$^{(n)}𝑑X\text{exp}_{}\left(\frac{1}{2}\underset{i,j}{}\text{}\right).$$
This is a recommended, non-trivial, exercise; it provides some perspective on what lies behind the commutativity of Diagram 7.0.1 and gives a particular sense in which wheels tangibly become cycles of something. Hint: Step A is the Aarhus formula; Step B is an appropriate identity from determinant theory, being careful with combinatorial factors.
## 8. Rozansky’s rationality conjecture
We start by recalling the surgery formula. Just say the pair $`(M,K)`$ is presented by some $`T`$, a $`\mu `$-string string link in the solid torus. If the associated decomposition is denoted
$$\sigma (\stackrel{ˇ}{Z}^{ST}(T))=\text{exp}_{}\left(\frac{1}{2}\underset{i,j}{}\text{}\right)R(\overline{x})B^{ST}(X),$$
where $`R(\overline{x})`$ is a series of $`X`$-substantial terms, then $`Z_n^{LMO}(M,K)`$ is equal to
$$\frac{\begin{array}{c}Wh(M,K)\text{exp}_{}(\frac{1}{2}_{i,j}\text{}),Thr^D(R(\overline{x}))\hfill \\ \\ \end{array}}{\left((1)^n^{(n)}𝑑U\sigma (\stackrel{ˇ}{Z}(U_+))\right)^{\sigma _+(W(1))}\left(^{(n)}𝑑U\sigma (\stackrel{ˇ}{Z}(U_{}))\right)^{\sigma _{}(W(1))}}B_n(\underset{¯}{k}).$$
We already know that this is a group-like element. In fact, we know that
$$Z^{LMO}(M,K)=Wh(M,K)\text{exp}(q),$$
where $`q`$ is a series of connected diagrams of degree greater than one whose dashed graphs have Euler characteristic less than zero. Thus we can read $`q`$ directly from the expression above. Write $`q=_{i=1}^{\mathrm{}}q^{(i)}`$ where $`q^{(i)}`$ is a series of diagrams of Euler characteristic minus $`i`$. Comparing the two expressions, it is clear that $`q^{(i)}`$ is calculated from precisely those diagrams appearing in $`R(\overline{x})`$ which have $`2i`$ trivalent vertices, plus some finite contribution from the signature corrections.
In other words, the term $`q^{(i)}`$ is a sum over those ways of gluing labelled chords into pairs of legs of some finite combination of polynomial generating diagrams which result in connected diagrams.
The correspondence is completed by noting that the determinant of the matrix $`W(T,e^k)`$ is $`\pm A_{(M,K)}(e^k)`$ so that the entries of the matrix $`W^1(T,e^k)`$ lie in $`L_{(M,K)}`$. Thus the edge-labels on the chords gluing the legs fall into that subspace, and the labels on the edges of the diagrams that are assembled according to the gluing formula also lie in this subspace.
It is an informative exercise to prove this theorem without appealing to the prior-known fact that the result is group-like.
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# Enhanced Spin Polarization of Conduction Electrons in Ni, explained by comparison with Cu
## I Spin–Polarized Currents in Magnetoelectronics
The rapidly growing field of magnetoelectronics is largely based on the manipulation of spin currents that are carried by electrons at the Fermi level $`ϵ_f`$. Examples are the application of giant magnetoresistance (GMR) in reading heads for hard disks and the use of spin–polarized tunneling and junction magnetoresistance (JMR) for a magnetic random access memory (MRAM). Spin–polarized tunneling is also being explored for high–resolution magnetic imaging by scanning tunneling microscopy (STM) . The magnitude of the magnetoresistance increases with the spin polarization of the currents, likewise the magnetic contrast in STM. A variety of efforts are directed towards designing new magnetic materials with higher spin–polarization, such as half–metallic compounds and nanostructures. For making systematic progress one first has to identify the electronic states that are responsible for the spin currents, then determine the fundamental parameters relevant for spin polarization, and eventually apply this knowledge to the design of new magnetic materials.
The character of the spin carriers in ferromagnets has been debated for some time. The initial puzzle has been whether s,p– or d–electrons dominate transport properties . The s,p–states have high group velocity, but low density of states and weak magnetism. The d–states carry the magnetic moment and have high density of states, but their group velocity is low. This dilemma can be resolved by looking at a realistic band structure where a free–electron–like s,p–band hybridizes with the magnetic d–levels close to the Fermi level . That allows the s,p–band to acquire a significant magnetic splitting .
The origin of the spin–polarization in is still under intense investigation . Various mechanisms have been proposed, such as a spin dependence of the density of states, spin–dependent electron scattering in the bulk and at interfaces, and a spin–dependent matrix element. A direct determination of the spin polarization from magneto–transport properties is difficult. An extensive set of values has been reported for spin–polarized tunneling into superconductors and Andreev reflection at point contacts to superconductors . The traditional explanation of such data has been the imbalance in the density of states at $`ϵ_f`$ for a magnetically–split free–electron band . It has been fairly successful for explaining the spin–polarization of Fe, but has failed for Ni where the observed spin–polarization of 23%–46% far exceeds the 6% spin–polarization expected from the density of states . A variety of more sophisticated approaches have been proposed for explaining spin–polarized tunneling . It is highly derirable to achieve high spin–polarization in tunneling from Ni alloys, such as permalloy (Ni<sub>0.8</sub>Fe<sub>0.2</sub>), since permalloy is the most common material in magnetoelectronics.
It is difficult to pinpoint the parameters relevant for achieving high spin–polarization from transport data alone. Such measurements integrate in k–space over the Fermi surface and involve additional parameters, such as scattering lengths. Angle–resolved photoemission is able to focus onto specific k–points and to separate out scattering lengths . Traditionally, this technique has lacked sufficient energy resolution for discerning the electronic states that are relevant for transport phenomena, that is those within a few thermal energies kT of $`ϵ_f`$. In our study, the energy resolution is 9 meV for electrons plus photons, compared to kT=25 meV at room temperature. Tunable synchrotron radiation allows us to map out the k component perpendicular to the surface independent of the parallel components. As non–magnetic reference material we use Cu. It has the same crystal structure as the adjacent Ni and a similar band topology. The main difference is an energy shift of the d–bands, which lie 2 eV lower in Cu than in Ni.
Our key observations are two intensity anomalies in the spin–split Ni conduction band: Below $`ϵ_f`$ the band loses intensity very rapidly in Ni but not in Cu. Furthermore, majority spins have larger photoemission intensity at $`ϵ_f`$ than minority spins. That creates an extra spin–polarization beyond the higher density of states for majority spins (which enters when integrating k over the Fermi surface). Several possible explanations are explored, such as increased electron scattering below $`ϵ_f`$, a photoemission matrix element that varies rapidly with E and k, and a transfer of spectral weight from single–electron excitations to many–electron excitations. Judging from our comparison with Cu and from simple matrix element calculations we assign the anomalies in Ni in large part to a rapid decrease of the matrix element at the point where the s,p–band becomes more d–like. Our finding suggests that a similar role can be expected from the matrix element in other phenomena, such as in spin–polarized tunneling.
## II Manybody States, Spectral Function, and Matrix Element
For encompassing the possibility of manybody interactions and electron scattering it is useful to start out with a very general characterization of electronic states in solids. That can be achieved by a spectral function $`A(𝐤,\omega )`$ which describes the spectral weight as a function of energy $`E=\mathrm{}\omega `$ and momentum $`𝐩=\mathrm{}𝐤`$. In a band structure model, where only single–electron excitations are possible, $`A(𝐤,\omega )`$ consists of sharp d–function peaks. The spectral function is far more general than this, however, and can describe all the many–electron effects measurable in a photoemission experiment. This generality is particularly useful for describing correlated electrons in the partially–filled 3d–shells of ferromagnets. Ni exhibits a broad satellite several eV below the single–hole states which may be viewed as a pair of correlated d–holes. The consequences of the two–hole satellite for the single–particle excitations in Ni are a reduction of the band width by 40% and a decrease of the magnetic splitting by a factor of 2–3 . These discrepancies have been associated with spectral weight shifting from the single–hole d–bands down to the two–hole satellite while preserving the center of gravity of the energy spectrum. A trading of spectral weight between single– and multi–electron excitations has been observed in adsorbates and oxides , too.
Such a manybody effect will be considered as one of two plausibe explanations for the drop of the photoemission intensity below $`ϵ_f`$ in Ni. The line of thought is the following: The conduction band in Ni is s,p–like above $`ϵ_f`$ but acquires more d–character below $`ϵ_f`$ as it starts hybridizing with the d–bands. The increasing d–character makes it prone to many–electron effects, such as a transfer of spectral weight to the two–hole satellite. Even the smaller intensity of the minority–spin peak at $`ϵ_f`$ would find a natural explanation in such a scenario, because the minority–spin d–bands lie higher in energy and hybridize more with the minority–spin conduction band at $`ϵ_f`$. For assessing this hypothesis we use Cu as reference material where many–electron effects weak. For example, the intensity of the two–hole satellite is 21% of the one–hole states in Ni, but only 2.5% in Cu .
The spectral function $`A(𝐤,\omega )`$ is related to the angle–resolved photoemission intensity $`I(\omega ,𝐤`$ by a matrix element $`M(𝐤,\omega )`$, that is specific to the photoemission process :
$`I(𝐤,\omega )=A(𝐤,\omega )\times M(𝐤,\omega )^2\times f(\omega )`$ (1)
The Fermi–Dirac function f(w) gives the occupancy. The spectral function itself has the form
$`A(𝐤,\omega )={\displaystyle \frac{1}{\pi }}\mathrm{}\left\{{\displaystyle \frac{1}{\omega ϵ_0(𝐤)\mathrm{\Sigma }(𝐤,\omega )}}\right\}`$ (2)
containing the complex self-energy $`\mathrm{\Sigma }(𝐤,\omega )`$ and the electron band dispersion $`ϵ_0(𝐤)`$. A fundamental sum rule for the spectral function implies a trade–off between single–electron and many–electron excitations:
$`{\displaystyle A(\omega ,𝐤)𝑑\omega }=1`$ (3)
The Fermi–Dirac function is absent, thus requiring an extrapolation of photoemission data above the Fermi level, or the inclusion of inverse photoemission data. One remaining piece in Eq. 1 to be determined is the matrix element $`M(𝐤)`$ for single–hole excitations from the $`\mathrm{\Sigma }_1`$ band:
$`M(𝐤)=Y_{final}(𝐤)\stackrel{}{A}_pY_{initial}(𝐤)`$ (4)
A is the vector potential of the photon and p the momentum operator. We have have performed an estimate of M(k) by using a combined interpolation scheme that takes the correct band width and splitting of the Ni d–bands into account .
## III The Photoemission Experiment
In photoemission, the parallel component $`𝐤_{}`$ is conserved and can be determined directly from the kinetic energy Ekin and the polar angle J of the photoelectrons. The perpendicular component $`𝐤_{}`$ varies with the photon energy hn and can be estimated using a free electron upper band with an inner potential In order to obtain a clear–cut spectral function we designed the experimental geometry such that it isolates a single band crossing the Fermi level with a high photoemission cross section. This is achieved by selecting the $`\mathrm{\Sigma }_1`$ conduction band along the direction in k–space. It crosses $`ϵ_f`$ about half–way between $`\mathrm{\Gamma }`$ and X and stays as far from the d–bands as possible . Dipole selection rules provide additional selectivity: The choice of p–polarized light with the electric field vector in the photoemission plane enhances the $`\mathrm{\Sigma }_1`$ band due to its even mirror symmetry and eliminates d–bands with odd symmetry. As a result, the photoemission data in Fig. 1 clearly show a single conduction band for Cu and a spin–split version of that band for Ni.
As a consistency check we map the same Fermi level crossing from two different surfaces, the (100) and the (110). For the (100) surface we reach the desired location with a photon energy $`h\nu `$=44 eV for Ni ($`h\nu `$=50 eV for Cu), combined with a polar angle of about $`20\mathrm{deg}`$ along the azimuth. The (110) surface probes the same k–point with a photon energy $`h\nu `$=27 eV and a polar angle of about 35$`\mathrm{deg}`$ along \[$`\overline{1}`$10 \]. For the (100) surface one starts at $`\mathrm{\Gamma }`$ for $`𝐤_{}`$ = 0 and reaches X at $`𝐤_{}=\sqrt{2.2}\pi /a=2.52\AA `$ in Ni ($`2.46\AA `$in Cu). For the (110) surface the bands are mapped in reverse, starting at X for $`𝐤_{}`$=0 and reaching $`\mathrm{\Gamma }`$ at $`𝐤_{}=\sqrt{2.2}\pi /a`$. This inverted k–scale shows up in Fig. 1 as approximate mirror symmetry of the (100) results (left) and the (110) results (right).
Comparing the intensities near $`ϵ_f`$ one finds opposite behavior for Ni and Cu (Fig. 1 top versus center). The Ni bands fade very quickly below $`ϵ_f`$, whereas the Cu band remains strong and even increases its intensity slightly. Losing oscillator strength so rapidly in Ni presents a puzzle: Where did the spectral weight go that ought to be there according to the sum rule in Eq. 3? Simple technical explanations fail. The k–acceptance of the analyzer would give equal trends for the intensity in Ni and Cu, contrary to the drop in Ni and increase in Cu. One could argue that the line width increases rapidly below $`ϵ_f`$ in Ni, thereby reducing the peak height. This lifetime broadening is due to the rapidly–increasing phase space for creating electron–hole pairs in the 3d bands of Ni. This hypothesis is discarded by fitting individual energy spectra at various k with Lorentzians in Fig. 2 and plotting the resulting peak areas in Fig. 3. The drop–off in Ni remains and contrasts with a slight increase in Cu.
The Lorentzian fit is equivalent to a simplified spectral function
$`A_0(𝐤,\omega )={\displaystyle \frac{1}{\pi }}{\displaystyle \frac{\mathrm{\Gamma }(𝐤)}{[wϵ(𝐤)]^2+\mathrm{\Gamma }(𝐤)^2}}`$ (5)
where the self–energy $`\mathrm{\Sigma }(𝐤,\omega )`$ is taken as functions of k only, not of $`\omega `$ . The real part of $`\mathrm{\Sigma }(𝐤)`$ is incorporated into the empirical band dispersion $`ϵ(𝐤)=ϵ_0(𝐤)+\mathrm{}\{\mathrm{\Sigma }(𝐤)\}`$. The imaginary part $`\mathrm{\Gamma }(𝐤)=\mathrm{}\{\mathrm{\Sigma }(𝐤)\}`$ describes a Lorentzian lifetime broadening. A small secondary electron background is added for fitting the data, which describes extrinsic energy losses of the photoelectrons on their way out. It consists of an integral over the Lorentzian line, which is equivalent to a step–like loss function.
In addition to the intensity drop below $`ϵ_f`$ there is a second anomaly in Ni. The area of the minority peak is smaller than that of the majority peak. This can be seen best from the k–distribution of the photoemission intensity at $`ϵ_f`$ in Fig. 4. The area ratio is $`I_{}/I_{}=1.8`$ for Ni(100) and $`I_{}/I_{}=1.2`$ for Ni(110). According to a single–electron band model one would expect very similar spectral weights for the two spin components, since they are so close together in k space. In fact, previous photoelectron spectra of the spin–split bands in Ni have usually been fitted with equal intensities for the two spins. We are able to unambiguously resolve the two components by measuring a k–distribution at $`ϵ_f`$, where the lifetime broadening is minimal. This spin asymmetry and the intensity drop in Ni are not sensitive to adsorbates such as a residual gas or a Cu overlayer, establishing them as pure bulk phenomena.
## IV Possible Explanations for the Anomalies in Ni
Within the framework established in Eqs. 1–4 there are two places where one can search for an explanation of the anomalous behavior of Ni relative to Cu. These are the matrix element $`M(w,𝐤)^2`$ and the spectral function $`A(𝐤,\omega )`$. As long as one wants to stay within the one–electron picture, the matrix element for excitation of single holes is the natural starting point. We have applied a combined interpolation scheme to the empirical band structures of Ni and Cu for obtaining estimates of $`M(w,𝐤)^2`$ . The result describes the (100) data qualitatively, including the opposite intensity trends for Ni and Cu. However, quantitative comparisons are fairly sensitive to the exact location of $`k_F`$, and the (110) data are not reproduced well. Clearly, more sophisticated calculations of the photoemission intensity are called for, such as the one–step model with evanescent surface wave functions. In the absence of quantitative calculations we use experimental results for explaining how the matrix element modifies the intensities in Ni and Cu. The key will be a rapid change in the hybridization between the s,p–band and the 3d–bands with energy.
While the $`\mathrm{\Sigma }_1`$ conduction band corresponds to the s,$`p_z`$ states in Cu, its symmetry allows for significant $`d_z^2`$ character in Ni. The Ni 3$`d_z^2`$ states lie close to $`ϵ_f`$ and strongly hybridize with the conduction band, whereas the Cu 3d states lie 2 eV lower. For finding a d–hybridization in Cu comparable to that of Ni one has to look 2 eV lower in energy, as shown in the bottom panels in Fig. 1. The group velocity, i. e. , the slope of the Cu conduction band is greatly reduced at this point and has become comparable to that of the Ni. This is the result of an avoided crossing with the $`d_z^2`$ level . Likewise, the intensity of the Cu band decreases strongly at these lower energies, similar to Ni below $`ϵ_f`$. The same situation is surveyed in k–space in Fig. 3. Ni and Cu behave similar if one shifts the Ni data to the point of comparable d–hybridization in Cu, i. e. , a shift to the left for (100) and to the right for (110). A calculation of the matrix element for (100) reproduces this effect qualitatively. From such similarities between the Ni bands at $`ϵ_f`$ and the Cu bands at 2 eV below $`ϵ_f`$ we conclude that the intensity changes in Ni are qualitatively consistent with a change in the matrix element due to increasing d–hybridization.
The imbalance between the two spin components can be explained in similar fashion. The minority spin conduction band is more d–like at $`ϵ_f`$ than its majority partner since it hybridizes with the higher–lying minority $`d_z^2`$–level. Therefore, its matrix element has decreased more than that of the majority band. The consequence is an enhanced spin–polarization at $`ϵ_f`$ which has implications for spin transport phenomena, such as spin–polarized tunneling and Andreev reflection at ferromagnetic point contacts . As mentioned above, the traditional density–of–states model fails to explain the high spin polarization observed in these experiments for Ni. The larger size of the majority spin Fermi surface in Ni would give only 6% spin polarization, compared to the observed 23%–46%. The extra spin–polarization that we find at $`ϵ_f`$ enhances the density–of–states effect and brings theory closer to experiment. For a quantitative comparison it will be necessary to map this polarization across the whole Fermi surface and to replace the photoemission matrix element by the tunneling matrix element.
Despite the qualitative success of the single–particle picture one ought to consider manybody effects in Ni. Excitations of two d–holes are well–documented in this material . Can they produce an effect similar to the decrease of the matrix element with increasing d–hybridization? There is a scenario where two–hole excitations steal spectral weight from the single–hole band, taking advantage of the sum rule in Eq. 3. It is not unreasonable to assume that the probability for exciting a pair of d–holes increases with the d–character of the band. Therefore, the same arguments as in the previous two paragraphs can be used, where increasing d–character of the band gives rise to a decreasing matrix element. It appears that only quantitative calculations of the matrix element can settle this issue. However, there are some interesting clues pointing towards a contribution of two–hole effects. The intensity drop in Ni is more abrupt than that in Cu at the point of comparable d–hybridization. This is particularly pronounced for the (110) surfaces (Fig. 3, right). One might expect a sharper drop–off for a two–hole process which scales like the square of the d–hybridization. An additional clue comes from the decreasing strength of two–hole excitations across the Periodic Table from Ni to Co and Fe . If manybody effects played a role in the spin polarization at $`ϵ_f`$, their influence would gradually fade from Ni to Co and Fe. Such a trend would nicely fit the results from spin–polarized tunneling, where the (one–electron) density–of–states model works best for Fe and worst for Ni .
## V Summary
In summary, we find a rapid loss of spectral weight in the conduction band of Ni below the Fermi level $`ϵ_f`$, which is opposite to the behavior of the analogous band in Cu. Possible mechanisms are considered, such as an increasing lifetime broadening, the single–hole matrix element, and many–hole excitations stealing spectral weight from single–hole excitations below $`ϵ_f`$. The comparison with Cu and a simple estimate of the matrix element indicate that the single–hole matrix element is able to give a qualitative explanation. An additional transfer of spectral weight to two–hole states is quite possible, however. The loss of spectral weight is larger for the minority spin band, thereby enhancing the spin–polarization at $`ϵ_f`$. The photoemission data suggest that similar enhancements of the spin–polarization might occur in magnetotransport and could be used in magnetoelectronic devices. For example, the spin–polarization observed in spin tunneling from Ni exceeds the traditional density–of–states model by a factor of five. The analogy with photoemission suggests that the tunneling matrix element might be responsible.
We would like to acknowledge stimulating discussions with J. C. Campuzano, R. Joynt, A. Chubukov, and J. Allen on many–electron effects and the spectral function. This work was supported by the NSF under Award Nos. DMR–9815416, DMR–9704196 and DMR–9809423. It was conducted at the SRC, which is supported by the NSF under Award No. DMR–9531009.
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# How Galaxies Disguise Their Ages
## 1 Introduction
For galaxies outside the Local Group, information on their stellar populations can only be obtained from integrated properties, such as colours and spectra. Unfortunately, integrated spectra and colours of composite stellar populations are known to be degenerate with respect to age and metallicity (e.g., Worthey 1994, for a recent discussion). For instance, the existence of a tight color-magnitude correlation for galaxies in nearby (Bower, Lucey & Ellis 1992) and distant (Stanford, Eisenhardt & Dickinson 1998) clusters is taken to imply that cluster galaxies form at high redshift over short timescales, whereas Schweizer & Seitzer (1992) are able to derive a tight $`UV`$ vs. $`V`$ correlation for field E/S0 galaxies, despite morphological and spectroscopic evidence for the existence of young stellar populations in these objects.
Worthey et al. (1994) have developed a system of narrowband spectrophotometric indices useful for separating age and metallicity effects in integrated spectra. Line strengths of Balmer line indices (H$`\beta `$, H$`\gamma `$ and H$`\delta `$) have been used to infer the presence of intermediate age populations in some early-type galaxies (e.g., Worthey 1997 and references therein). Jørgensen (1999) uses H$`\beta `$ indices to infer the existence of a correlation between mean age and metal abundance for galaxies in the Coma cluster: young galaxies may masquerade among older objects because of their higher metal abundance, allowing a wide range of ages to be consistent with the passive evolution of the observed colour-magnitude relation and the small intrinsic scatter (Ferreras, Charlot & Silk 1999).
On the other hand, the population synthesis models used to derive theoretical index grids do not paint a complete picture of the stellar populations of early-type galaxies. The spectral energy distributions of these objects exhibit an unexpected rise in flux shortward of 2500 Å, a phenomenon dubbed the far ultraviolet (FUV) upturn (see O’Connell 1999 for a review). The currently accepted explanation for the FUV upturn is that it is caused by evolution of metal-rich stars on to the extreme Horizontal Branch (HB) and their UV-bright progeny (Brown et al. 1997 and references therein), whereas standard models terminate their evolution on the red clump on the HB (e.g., Sweigart 1987). FUV sources are known to exist in the metal rich, old “open” cluster NGC6791 (Liebert, Saffer & Green 1994) and are identified with field subdwarfs O and B in our own Galaxy (Saffer et al. 1997). These objects show strong, broad Balmer lines and may contribute significantly to H$`\beta `$, H$`\gamma `$ and H$`\delta `$ indices. This is indeed the case for some metal poor globular clusters with long blue tails (de Freitas Pacheco & Barbuy 1995).
Nevertheless, these stars are not included in the population synthesis models of Worthey (1994), where HBs are treated as red clumps with a temperature offset. The isochrones of Bertelli et al. (1994) used by Vazdekis (1999) include an ‘AGB-manque’ phase for high metallicities, but only for low mass stars ($`M<0.60M_{}`$) at large ages ($`>20`$ Gyr), whereas the existence of hot blue stars in NGC6791 and the field sdO/B show that stars of $`1M_{}`$, are able to evolve on to the extreme HB at ages of $`10`$ Gyrs (Carraro, Girardi & Chiosi 1999). The models of Bressan, Chiosi & Tantalo (1996) predict FUV colours for galaxies of the appropriate ages and metallicities, but do not calculate Balmer line indices contributed from FUV sources, since fitting functions for stars of such high temperature and gravity are not yet available. Worthey, Dorman & Jones (1996) have computed the contribution to the integrated flux between 2000 and 2400 Å from a warm turnoff, and find that this may account for $`50\%`$ of the observed luminosity, but do not explore the effect of HB sources on spectrophotometric indices.
The purpose of this Letter is to consider, to a first approximation, how FUV sources affect integrated Balmer line indices and therefore whether the claims for younger mean ages in some early-type galaxies may not be better explained by variations in HB morphology. We find that these objects provide a significant amount of Balmer line absorption and may therefore affect the derivation of ages via the H$`\beta `$, H$`\gamma `$ and H$`\delta `$ indices.
## 2 Modelling
We adopt a semi-empirical approach in which we first estimate the fraction of extreme HB stars (and progeny) needed to produce the observed FUV colours, for two representative models at very different metallicities, and then compute the contribution from these objects to the total Balmer line absorption strength.
We use models by Dorman, O’Connell & Rood (1993a) in which a 10–30% fraction of HB stars evolves on to UV-bright phases. We use the two models for which detailed evolutionary calculations are presented by Dorman, Rood & O’Connell (1993b): one with \[Fe/H\]=+0.38, Y=0.292, $`M_c=0.464M_{}`$ (where $`M_c`$ is the core mass) and envelope masses of 0.003, 0.046 and 0.096 $`M_{}`$ and a model with \[Fe/H\]=–1.48, Y=0.247, $`M_c=0.485M_{}`$ and envelope masses of 0.003, 0.035 and 0.105 $`M_{}`$. For each of these models we follow the prescriptions of Dorman et al. (1993a) to calculate FUV colours and their contribution to the total $`V`$ band light.
We then use the evolutionary tracks presented in Dorman et al. (1993b) to estimate the fraction of total light at each $`T_{eff}`$ and $`\mathrm{log}g`$ step and the stellar atmospheres of Kurucz (1993) to estimate the equivalent widths of H$`\beta `$, H$`\gamma `$ and H$`\delta `$ at each stage. Since most of the models reach temperatures and gravities in excess of the grid calculated by Kurucz, we extrapolate H$`\beta `$,$`\gamma `$, $`\delta `$ equivalent widths to the appropriate T$`{}_{eff}{}^{}\mathrm{log}g`$ range by means of polynomial fits to the predictions for the existing grid, being unable to completely simulate the spectrum of these objects. These are then summed, scaling by the fraction of total light produced, to yield the equivalent widths of Balmer lines contributed from FUV sources during their lifetime and again scaled by the fraction of total $`V`$ band light to calculate the additional absorption line strength to integrated H$`\beta `$, H$`\gamma `$ and H$`\delta `$ indices, following the prescriptions of Freitas Pacheco & Barbuy (1995). Table 1 shows the models used. Here column 1 is the fraction of blue HB stars, column 2 the contribution to the $`V`$ band light, column 3 the 1550$`V`$ color, column 4 the 2500$`V`$ color and column 5 the extra H$`\beta `$ strength provided by these stars. A header at the top indicates the model parameters and mean temperature of the models.
## 3 Discussion
Figure 1 shows the excess equivalent width of H$`\beta `$ produced by HB stars and their progeny as a function of FUV colour of the host galaxy. We only plot the extra contribution to Balmer line strengths provided by these stars, without assuming any underlying model. We assume that the Balmer line strength produced by FUV sources can be added linearly to the chosen galaxy model from the Worthey (1994) compilation.
The range of FUV colours in the Burstein et al. (1988) sample is 2 to 4. From Figure 1, this corresponds to H$`\beta `$ equivalent widths of up to 0.6 Å, with a typical contribution of 0.3 Å, in agreement with earlier results on metal poor globular clusters (de Freitas Pacheco & Barbuy 1995) and ‘blue HB’ simulations of Buzzoni, Mantegazza & Gariboldi (1994), albeit at lower temperatures than those of FUV sources. In our simulations, metal-poor HB stars yield somewhat higher H$`\beta `$ strengths than those of metal-rich stars.
Figure 2 plots a single stellar population (SSP) model grid from Worthey (1994), and superposes the range of H$`\beta `$ strengths contributed by FUV sources to an underlying 12 Gyr old population. We add index strengths for the SSP and the HB stars linearly, as stated above. We also show some of the higher signal-to-noise measurements of Trager et al. (1998). It can be seen that spurious age differences of $`57`$ Gyrs can be introduced by FUV sources; conversely, assuming that more complex stellar populations can be represented by linear combinations of SSP models, a 10–20% burst of star formation observed $`13`$ Gyrs after star formation ceases may be explained by FUV sources. Larger bursts observed at later ages can also be accounted for in this manner.
Figure 1 shows that a range of H$`\beta `$ strengths is possible at any FUV colour. In turn, the spread in FUV colours at any age larger than 5 Gyr is significant (Tantalo et al. 1996). Taken at face value, this suggests some degeneracy in FUV colour, H$`\beta `$ strength and age for early-type galaxies.
Our simulations show that the contribution to H$`\gamma `$ and H$`\delta `$ are similar to H$`\beta `$. The model grid of Jones & Worthey (1995) spans about 0.2 Å, for a range of ages from 3 to 17 Gyrs. This suggests that FUV sources strongly affect these indices as well. On the other hand, the narrow H$`\gamma _{A,F}`$ and H$`\delta _{A,F}`$ indices of Worthey & Ottaviani (1997) appear to be much less sensitive to FUV source contributions, although in this case the narrowness of the index bandpasses probably requires more accurate modelling. The new H$`\gamma `$ index of Vazdekis & Arimoto (1999) varies by about 0.5 Å over ages of 1.6 to 17 Gyrs (depending on velocity dispersion). After correction for velocity dispersion, FUV source contribution can vary between 0 and 0.4 Å, spanning a sizeable portion of these newer grids.
Since Dorman et al. (1993b) do not present integrated fluxes for bands other then $`V`$, we are unable to estimate the effect of FUV sources on broadband colours. Nevertheless, the contribution to $`V`$ from high metallicity models is typically 5% and always less than 10%, which should not affect broadband colors. For low metallicity models HB stars may provide as much as 20% of the light and this may produce bluer than expected colours. For comparison, the two globular clusters with long blue tails, M13 and NGC6229, are seen to have too blue integrated $`UB`$ for their $`BV`$ colour (Reed, Hesser & Shawl 1988).
An useful consistency check is to compare $`2500V`$ colours for our models and observations. For high metallicity systems, Table 1 shows that we reproduce well the observed range of colours in the sample of Burstein et al. (1988), which is typically 3–4. Low metallicity models are far bluer, which is not surprising, since early-type galaxies are generally metal rich, but they are consistent with $`2500V`$ colours of globular clusters, which are typically 1.5–3.5 and span the appropriate metallicity range.
Figure 3 plots excess H$`\beta `$ emission (over an 18 Gyr old population) for galaxies with index measurements from Trager et al. (1998) and FUV colours from Burstein et al. (1988). This was calculated following Davies, Sadler & Peletier (1993) and includes some galaxies for which they provide excess H$`\beta `$ values. We find no significant correlation between $`\mathrm{\Delta }`$H$`\beta `$ and 1550$`V`$ colour, although some galaxies have both strong excess absorption and blue FUV colors. One of these galaxies, however, is NGC5102, the well known E+A galaxy.
This may imply that the effect of FUV sources on derived index strengths is weak. As shown in Table 1, the largest contribution to H$`\beta `$ comes from objects with high surface temperature. Brown et al. (1997) estimate that the largest contribution to the FUV flux comes from stars with surface temperatures lower than 25,000 K. These sources would not provide large H$`\beta `$ absorption, in agreement with the data presented in Figure 3. One caveat, is that FUV colors are measured through the large IUE aperture, which is generally wider than the slit sizes used to measure spectral indices: Ohl et al. (1998) show that FUV color gradients can be strong, making a comparison such as shown in Figure 3 possibly unrealistic.
We have chosen two of the Dorman et al. (1993b) tracks considered to best represent the FUV sources: it should be noted, however, that different objects, or different evolutionary tracks, may contribute to the FUV upturn: for instance post AGB stars are believed to be important in M31 (Bertola et al. 1995). The range of models used needs to be considerably expanded to include a wider range of objects to make these results more comprehensive.
There are a number of observational tests of our models: it is possible to measure indices for the ‘quiescent’ samples of Burstein et al. (1988) and Longo et al. (1989) in consistent apertures. Conversely, FUV strengths can be derived, from HST data, for Coma galaxies observed by Jørgensen (1999). It is also possible to explore the correlation of line strength indices with the distribution of FUV components observed by UIT (Brown et al. 1997). More accurate stellar atmospheres and fitting functions for hot, high gravity, metal rich stars are also necessary, to allow the FUV component effects to be accounted for in population synthesis models.
## 4 Conclusions
We have considered the FUV source contribution to Balmer line indices using a semi-empirical modelling technique. We find that this is significant and may lead to identification of spurious intermediate age populations or age gradients in galaxies. Nearly all Balmer line indices in use are potentially affected by FUV stars and consideration should be paid to their age sensitivity.
We would like to thank Adam Stanford for commenting on an earlier version of this paper, and Sabine Möhler for her help with Kurucz models. We would also like to thank the referee, Dr. Guy Worthey, for a number of useful suggestions. This work is supported by the Australian Research Council.
FIGURE CAPTIONS
Figure 1: FUV colours vs H$`\beta `$ strength produced by our models. The shaded region represents the range of FUV colours allowed by our simulations. A range of Balmer line strengths is possible at each colour.
Figure 2: A single stellar population grid of H$`\beta `$ vs. \[Fe/H\] for different ages and abundances. We superpose the range of possible contributions from FUV sources to a 12 Gyr old population (dark bars). We also plot indices for high goodness galaxies (open circles) from the sample of Trager et al. (1998). Age differences of a few Gyrs, for old populations, or bursts of 10–20% size observed 1–3 Gyrs after star formation ceases, can be accounted for by FUV sources.
Figure 3: Excess H$`\beta `$ emission vs. FUV colour, from the sample of galaxies in common to Trager et al. (1998), Davies et al. (1993) and Burstein et al. (1988). Excess H$`\beta `$ is defined as in Davies et al. (see text).
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# On Estimating the QSO Transmission Power Spectrum
## 1. Introduction
Recent theoretical research on the low column density ($`N_{\mathrm{HI}}{}_{}{}^{}{}_{}{}^{<}\mathrm{\hspace{0.17em}10}^{14.5}\mathrm{cm}^2`$) Lyman-alpha (Ly$`\alpha `$) forest at redshifts $`z24`$ points towards a picture in which the forest consists largely of mildly nonlinear fluctuations of a smooth intergalactic medium (e.g. Bi et al. (1992); Cen et al. (1994); Zhang et al. (1995); Reisenegger & Miralda-Escudé (1995); Hernquist et al. (1996); Miralda-Escudé et al. (1996); Muecket et al. (1996); Bi & Davidsen (1997); Bond & Wadsley (1997); Croft et al. (1997); Hui et al. (1997); Hui & Gnedin (1997); see Rauch (1998) for a review and further ref.). This provides the motivation to analyze the quasar (QSO) absorption spectrum as a continuous field with fluctuations, rather than as a collection of discrete absorption lines. The two-point correlation or its fourier transform, the power spectrum, comes to mind as a useful and common statistic used in other areas such as the microwave background or galaxy large-scale-structure. Indeed, its application to QSO spectrum has been discussed by a number of authors (Zuo & Bond (1994); Miralda-Escudé et al. (1996); Bi & Davidsen (1997); Cen et al. (1998)). Recently, Croft et al. Croft et al. (1998a, b) (see also Hui (1999); McDonald et al. (1999)) have shown that the mass power spectrum can be recovered from the QSO transmission power spectrum, from which one could further deduce cosmological parameters such as $`\mathrm{\Omega }_m`$ (Weinberg et al. (1998); Phillips et al. (2000)). There exist at present a large number of high quality QSO spectra (e.g. Hu et al. (1995); Lu et al. (1996); Kirkman & Tytler (1997); Cristiani et al. (1997); Kim et al. (1997); Rauch et al. (1997)) which makes this an exciting field of research. Upcoming quasar surveys such as the Sloan Digital Sky Survey (SDSS) and the Anglo-Australian Telescope Two Degree Field (AAT2DF) will enlarge the database significantly.
Here we take the view that the QSO transmission power spectrum / correlation is interesting in its own right, and focus on how to best measure it from the observed QSO spectra, independent of the underlying physical picture of the forest. The two major questions are: (1) what are the main sources of systematic errors and what are the best ways to bring them under control? (2) how to estimate the shot-noise, and to best combine data with different signal-to-noise (random errors)?
Let us start by defining the transmission power spectrum and correlation function. Two possibilities arise. One of them we call the un-normalized power spectrum $`P_{\mathrm{un}}`$ / two-point correlation $`\xi _{\mathrm{un}}`$ (Weinberg 1998, private communication; McDonald et al. (1999)):
$$\xi _{\mathrm{un}}(u)=f(u^{})f(u^{}+u),P_{\mathrm{un}}(k)=\xi _{\mathrm{un}}(u)e^{iku}𝑑u$$
(1)
where $`f`$ is the transmission defined by $`f=e^\tau `$ with $`\tau `$ being the optical depth (the absorption is then $`1f`$), $`u`$ or $`u^{}`$ is the observed velocity (or redshift or wavelength) along a line of sight, and $`k`$ is its fourier counterpart. The angular brackets denote ensemble averaging.
The other we call the normalized power spectrum $`P`$ / two-point correlation $`\xi `$ (e.g. Zuo & Bond (1994)):
$$\xi (u)=\delta _f(u^{})\delta _f(u^{}+u),P(k)=\xi (u)e^{iku}𝑑u$$
(2)
where $`\delta _f=(f\overline{f})/\overline{f}`$, with $`\overline{f}`$ being the mean transmission $`f`$.
We will almost exclusively focus on the latter, but will discuss at some point the pros and cons of the two, especially with regard to systematic errors. Unless otherwise stated, hereafter, power spectrum / correlation refers to the normalized version.
The layout of the paper is as follows. In §3, we provide a brief overview of how the raw data output (a two-dimensional CCD image) is reduced to a one-dimensional QSO spectrum. Note that the quantity $`f`$ above is never observed directly. It is important to have a description of how the whole data reduction procedure works, which is sometimes hard to find in the literature. We give some illustrations by showing simulated spectra with various realistic levels of noise.
In §4, we discuss the estimation of the power spectrum and two-point correlation, beginning with the introduction of the quadratic estimator in §4.1. An important point pertaining to the estimation of the two-point correlation is raised here – most estimators employed in the literature are not optimal; an alternative is given here which is an analogue of the Landy-Szalay estimator (Landy & Szalay (1993)) introduced originally for galaxy surveys. Aside from this point, we focus exclusively on the estimation of the power spectrum. In §4.2, we discuss three sources of systematic errors: continuum-fitting, gaps and metal absorption lines. Particular attention is paid to issues related to continuum-fitting. We advocate in §4.2.1 trend-removal to replace traditional continuum-fitting, which avoids the latter’s pitfalls. We further propose in §4.2.2 that the power spectrum of the continuum can be estimated using trend-removal as well, which offers us a way to measure accurately the transmission power spectrum on large scales where the continuum fluctuations might be important. In §4.2.3, we discuss the effects of gaps and un-removed metal lines. We then turn our attention in §4.3 to random errors. We emphasize here that the shot noise is not exactly Poisson distributed because of the particular way the data is reduced. We point out the importance of subtracting the shot-noise-bias correctly, and describe a systematic way of assigning error-bars to the power spectrum, and introduce minimum-variance-weighting techniques to combine data of different qualities. Some results here are stated without justification. The aim in this section is to summarize useful results for readers who might not be interested in details of the derivations, which are provided in the Appendices. The techniques used in the Appendices should be of broad interest e.g. the issue of generally non-Poissonian shot-noise might be relevant for galaxy power spectrum estimation.
Lastly, we conclude in §5. We summarize here our recipe for transmission power estimation – readers who would like a quick overview of our methods can skip directly to this section, and only refer back to the relevant sections to fill in the details. We give general advice on observing strategies, and discuss in particular analysis issues relevant for the Sloan Digital Sky Survey.
Before we begin, let us first make a few clarifying remarks about some of our notation and terminology.
## 2. Terminology and Notation: Averaging and Averages
In this paper, we refer to two different kinds of fluctuations which should be clearly distinguished. Take the observed photon (or electron) count from a quasar as an example, $`\widehat{N}_\mathrm{Q}`$. As one moves along a spectrum, the photon count fluctuates because of two very different reasons. First, it fluctuates because the universe is intrinsically inhomogeneous, giving rise to non-uniform absorption – we will refer to these as cosmic fluctuations. Second, it fluctuates because the observed photon count is a discrete realization of the underlying cosmic signal – Poisson fluctuation is the canonical example but not the only possible one, we will refer to these as discrete fluctuations.
We define two different kind of averages corresponding to these two different kind of fluctuations. The discrete average of the observed photon count is denoted by $`\widehat{N}_\mathrm{Q}_D\stackrel{~}{N}_\mathrm{Q}`$. In other words, $`\widehat{N}_\mathrm{Q}`$ constitutes a discrete realization of the underlying cosmic signal $`\stackrel{~}{N}_\mathrm{Q}`$. This signal $`\stackrel{~}{N}_\mathrm{Q}`$ itself suffers from cosmic fluctuations, and we will denote its ensemble average by $`\stackrel{~}{N}_\mathrm{Q}=\overline{N}_\mathrm{Q}`$. A fixed quasar continuum is assumed in this ensemble average i.e. it is the fluctuation in the spectrum caused by intervening absorption that constitutes the cosmic signal we are after. Finally, we will sometimes use $``$ to implicitly stand for $`_D`$ e.g. $`\widehat{N}_\mathrm{Q}`$ actually means $`\widehat{N}_\mathrm{Q}_D`$ which is the same as $`\stackrel{~}{N}_\mathrm{Q}=\overline{N}_\mathrm{Q}`$ <sup>1</sup><sup>1</sup>1Two exceptions in the use of $``$: in §4.2.2, we use $``$ to include averaging over an ensemble of different continua, and in §4.3 we use $`_{kk}`$ to denote averaging over a shell of Fourier modes.
To recap:
* $`\widehat{N}_\mathrm{Q}`$ is the directly observed quasar photon count.
* $`\stackrel{~}{N}_\mathrm{Q}\widehat{N}_\mathrm{Q}_D`$ refers to the idealized quasar photon count if one has infinite signal-to-noise (S/N).
* $`\overline{N}_\mathrm{Q}\stackrel{~}{N}_\mathrm{Q}=\widehat{N}_\mathrm{Q}_D`$ refers to the quasar photon count if one has infinite S/N and if one averages over all possible cosmic fluctuations keeping the continuum fixed e.g. by taking the same quasar and putting it at all possible orientations in the sky. For instance, if $`\stackrel{~}{N}_\mathrm{Q}=N_Ce^\tau `$ where $`N_C`$ is the true continuum and $`\tau `$ is the optical depth, then $`\overline{N}_\mathrm{Q}=N_Ce^\tau `$ where $`e^\tau `$ is the mean transmission.
Note that, when we use the term discrete, it is not implied that $`\widehat{N}_\mathrm{Q}`$ is necessarily an integer, although it is derived from some integral quantity such as the electron/photon count. We use the term discrete fluctuations instead of the usual Poisson fluctuations, because as we will see, $`\widehat{N}_\mathrm{Q}`$ is often not Poisson-distributed i.e. $`\widehat{N}_\mathrm{Q}^2_D\widehat{N}_\mathrm{Q}_D^2\stackrel{~}{N}_\mathrm{Q}=\widehat{N}_\mathrm{Q}_D`$ (see §3.1). The term shot-noise is often used to describe Poissionian discrete fluctuations, but in this paper we will use it more broadly to include non-Poissonian discrete fluctuations as well.
Finally, we note that we use the term ’quasar counts’ throughout this paper to refer to the photon counts from a quasar, rather than the number of quasars in a given patch of sky.
## 3. Data Reduction: From the CCD Image to the QSO Spectrum
### 3.1. A Brief Description
We discuss briefly here aspects of the data-processing necessary for understanding the noise properties of the reduced quasar spectrum. The reader is referred to Horne Horne (1986), Zuo & Bond Zuo & Bond (1994), Barlow & Sargent Barlow & Sargent (1997), Rauch et al. Rauch et al. (1997) and Cen et al. Cen et al. (1998) for more discussions.
The raw data consist of a two-dimensional (spatial and spectral) array of counts (data numbers or photon counts converted from them) from a CCD image. The one-dimensional array of estimated quasar counts in the spectral direction (as a function of velocity, redshift or wavelength) is obtained by collapsing the data in the spatial direction in the following fashion:
$$\widehat{N}_\mathrm{Q}^\alpha =\underset{i,\beta }{}W^{\alpha \beta }W^{i\beta }(\widehat{N}_{\mathrm{RAW}}^{i\beta }\stackrel{~}{N}_\mathrm{B}^{i\beta })$$
(3)
We have introduced and will stick with the following notations: the Latin index such as $`i`$ and the Greek index such as $`\alpha `$ denote the spatial and spectral coordinates respectively of a CCD pixel (there are in fact a few exceptions, which should be clear from the context); $`\widehat{N}_\mathrm{Q}^\alpha `$ is our estimated quasar count, $`\widehat{N}_{\mathrm{RAW}}^{i\beta }`$ is the raw count, and $`\stackrel{~}{N}_{\mathrm{B}}^{}{}_{}{}^{i\beta }`$ is the mean background count which includes the sky and the readout offset; $`W^{i\beta }`$ is a weighting of the spatial pixels for each spectral coordinate $`\beta `$, and $`W^{\alpha \beta }`$ represents a rebinning of the spectral pixels that is sometimes done to achieve, for instance, a linear wavelength scale. Note that the $`\alpha `$ and $`i`$ dimensions do not necessarily align with the two perpendicular axes of the CCD chip. The optical setup could be such that the spectrograph slit appears tilted at an angle to the CCD axes. We use $`\widehat{}`$ to emphasize the fact that the quantity of interest is a random variable with fluctuations. The $`\stackrel{~}{}`$ denotes a discrete average: e.g. $`\stackrel{~}{N}_{\mathrm{B}}^{}{}_{}{}^{i\beta }=\widehat{N}_{\mathrm{B}}^{}{}_{}{}^{i\beta }_D`$ where $`_D`$ denotes discrete averaging.
Implicitly assumed in the above formulation is that the discrete average $`\stackrel{~}{N}_{\mathrm{B}}^{}{}_{}{}^{i\beta }`$ is known, which is of course not strictly true, but since a typical slit covers a significant number of pixels that do not have any quasar photons, and since the background is often quite uniform, the discrete average can be estimated to high accuracy. Note also the readout offset can be measured using short exposures with closed shutters or from the CCD overscan region.
The weighting $`W^{i\beta }`$ typically has non-trivial spectral dependence ($`\beta `$) to remove at least two artifacts: variations in detector efficiency across the chip and a non-flat blaze. The former is usually estimated in a procedure called flat-fielding by shining a lamp into the detection system. The latter arises because of the non-trivial shape of a diffraction order. This can be partially estimated in the flat-fielding procedure, but is best done using a spectrophotometric standard star, usually a white dwarf. While the correction for the first artifact should be quite accurate, the blaze-removal is often approximate. Any residual that is not correctly removed will show up in the form of a non-trivial effective continuum. We will see in §4.2 perhaps some evidence of it. We assume in this paper that such artifacts show up as fluctuations on large scales (since the blaze itself is smoothly and slowly varying across a given order) but not on small scales (we will quantify the scales later on).
To make the above concrete, the raw count is given by:
$$\widehat{N}_{\mathrm{RAW}}^{i\beta }=\widehat{N}_\mathrm{B}^{i\beta }+\widehat{N}_\mathrm{Q}^{i\beta }$$
(4)
where the quasar contribution has the following discrete average:
$$\stackrel{~}{N}_\mathrm{Q}^{i\beta }\widehat{N}_\mathrm{Q}^{i\beta }_D=g_{\mathrm{ps}}^{i\beta }g_\mathrm{b}^\beta \stackrel{~}{N}_\mathrm{Q}^\beta $$
(5)
where $`g_{\mathrm{ps}}^{i\beta }`$ is the point-spread function which describes how the light from the quasar gets spread-out in the spatial direction $`i`$ at a given spectral coordinate $`\beta `$, and $`g_\mathrm{b}^\beta `$ accounts for the variation of the blaze and quantum efficiency as a function of wavelength. The symbol $`\stackrel{~}{N}_\mathrm{Q}^\beta `$ denotes the underlying quasar count (or cosmic signal i.e. discrete averaged).
Many different rebinning kernels $`W^{\alpha \beta }`$ (eq. ) are possible. The simplest choice is of course no rebinning with $`W_{\alpha \beta }=\delta _{\alpha \beta }`$.
There are several possible choices of the weighting $`W^{i\beta }`$ (eq. ), but any sensible choice has to satisfy the requirement that $`\widehat{N}_\mathrm{Q}^\alpha _D=\stackrel{~}{N}_\mathrm{Q}^\alpha `$, up to some constant normalization factor. This assumes that artificial fluctuations introduced by the blaze or detector efficiency are correctly taken out. If not, it shows up effectively as part of the continuum.
We give two examples of $`W^{i\beta }`$ here. The first is basically a uniform weighting over the spatial pixels that correspond to a given spectral coordinate:
$$W^{i\beta }=1/(g_\mathrm{b}^\beta \underset{j}{}g_{\mathrm{ps}}^{j\beta }),$$
(6)
where the range of $`i`$, or $`j`$, is chosen to lie within, say, some fraction of the quasar seeing disk. There is sometimes an additional complication due to cosmic ray hits, which will be discussed below.
The second is a minimum variance weighting (different from minimum variance weighting for measuring the power spectrum; §4.3) over the spatial pixels, introduced by Horne Horne (1986):
$$W^{i\beta }=(1/g_\mathrm{b}^\beta )(g_{\mathrm{ps}}^{i\beta }/V_{\mathrm{RAW}}^{i\beta })/(\underset{j}{}[g_{\mathrm{ps}}^{j\beta }]^2/V_{\mathrm{RAW}}^{j\beta })$$
(7)
where $`V_{\mathrm{RAW}}^{j\beta }`$ is the variance in the raw count:
$`V_{\mathrm{RAW}}^{j\beta }=(\widehat{N}_{\mathrm{RAW}}^{j\beta })^2_D\widehat{N}_{\mathrm{RAW}}^{j\beta }_D^2=\stackrel{~}{N}_\mathrm{Q}^{j\beta }+V_\mathrm{B}^{j\beta }`$ (8)
$`V_{\mathrm{B}}^{}{}_{}{}^{j\beta }=\stackrel{~}{N}_{\mathrm{S}}^{}{}_{}{}^{j\beta }+V_{\mathrm{R}.\mathrm{O}.}^{j\beta }`$
where $`V_\mathrm{B}^{j\beta }`$, the background variance, has two contributions, the sky variance $`\stackrel{~}{N}_{\mathrm{S}}^{}{}_{}{}^{j\beta }`$ and the readout variance $`V_{\mathrm{R}.\mathrm{O}.}^{j\beta }`$. A word of caution is necessary here regarding the second weighting. The raw variance, $`V_{\mathrm{RAW}}^{j\beta }`$, depends on the underlying cosmic signal or quasar count ($`\stackrel{~}{N}_\mathrm{Q}^{j\beta }`$ i.e. discrete averaged) which is not directly observable (the discrete averaged sky count and the true readout variance are also strictly speaking not directly observable, but they can be estimated quite accurately because they are relatively uniform and can be observed over a larger number of pixels). Modeling $`V_{\mathrm{RAW}}^{j\beta }`$ using the measured raw count (i.e. using $`\widehat{N}_\mathrm{Q}^{j\beta }`$ instead of $`\stackrel{~}{N}_\mathrm{Q}^{j\beta }`$ in eq. ) could lead to a biased estimation of $`\stackrel{~}{N}_\mathrm{Q}^\beta `$. Horne Horne (1986) suggested an iterative scheme to avoid this problem, but implementations of this weighting should be checked for a possible bias.
Pixels affected by cosmic-ray hits, which are usually easy to identify because of their wild fluctuations and spiky nature, are dealt with in two different ways, depending on the severity. For a given spectral coordinate, if only a small fraction of the corresponding spatial pixels are affected, the weighting in eq. (6) or eq. (7) is simply modified by allowing $`i`$ and $`j`$ to only range over the unaffected spatial pixels. However, if all or most corresponding spatial pixels are affected, then all recorded counts at that spectral coordinate are discarded, leaving a gap in the reduced quasar spectrum. Gaps could result also because of metal-line removal (an alternative would be to fit for the metal-line and subtract, instead of simply discarding the pixels) or defects in the CCD.
Finally, the (random) error-array output at the end of the data reduction corresponds to an estimate of
$$\sqrt{(\widehat{N}_\mathrm{Q}^\alpha \stackrel{~}{N}_{\mathrm{Q}}^{}{}_{}{}^{\alpha })^2_D}=\sqrt{\underset{i,\beta }{}(W^{\alpha \beta }W^{i\beta })^2V_{\mathrm{RAW}}^{i\beta }}$$
(9)
where we have assumed that the noise-fluctuations are independent among the pixels. We emphasize that in practice the error-array is only an estimate of the above quantity, because the true $`V_{\mathrm{RAW}}^{i\beta }`$ is unknown, but is estimated using the observed raw counts (using $`\widehat{N}_\mathrm{Q}^{j\beta }`$ instead of $`\stackrel{~}{N}_\mathrm{Q}^{j\beta }`$ in eq. ).
It is clear from the above discussion that in general fluctuations in $`\widehat{N}_\mathrm{Q}^\alpha `$ are non-Poissonian, in the sense that $`\widehat{N}_\mathrm{Q}^2_D\widehat{N}_\mathrm{Q}_D^2\widehat{N}_\mathrm{Q}_D=\stackrel{~}{N}_\mathrm{Q}`$. This is because of several reasons. First, $`\widehat{N}_\mathrm{Q}`$ suffers from additional discrete fluctuations from the background counts. Second, the weighting $`W^{\alpha \beta }`$ and $`W^{i\beta }`$ are in general non-trivial (i.e. not unity). A very simple example of the effect of non-unit weights is: suppose we multiply a Poisson variable by a factor of 2 and call the result $`\widehat{y}`$, it is easy to see that $`\widehat{y}^2_D\widehat{y}_D^2=2\widehat{y}_D\widehat{y}_D`$.
For the rest of this paper, we will pick for simplicity the weighting kernels $`W_{\alpha \beta }=\delta _{\alpha \beta }`$ and $`W_{i\beta }`$ as given by eq. (6). In the Appendices we will indicate where some of our expressions have to be modified to account for more general weightings.
### 3.2. Simulated QSO Spectra
For illustrations, and for later analyses, we have generated several different QSO spectra. The underlying noiseless (theoretical) transmission ($`f=e^\tau `$) is shown in the bottom panel of Fig. 1, and its associated power spectrum is shown in the top panel of the same figure. This is drawn from a SCDM (Standard Cold Dark Matter) simulation discussed in Gnedin Gnedin (1998) which made use of the Hydro-PM algorithm developed by Gnedin & Hui Gnedin & Hui (1998). The cell-size (comoving size of $`10\mathrm{h}^1\mathrm{kpc}`$) is small enough to resolve the effective Jeans scale, and so should retain all small scale structures. However, the box-size is unrealistically small (comoving size $`2.56\mathrm{h}^1\mathrm{Mpc}`$) which means a significant amount of large scale power is missing. For most of our investigations here, it is not necessary that the simulations are highly realistic, but our simulated transmission power spectrum is in fact broadly consistent with an observed one (Fig. 14). The long line of sight in Fig. 1 is generated by shooting a ray at some oblique angle through the simulation box and allowing it to wrap around the box several times, but never repeating itself. The mean redshift here is $`z=2.85`$. The ionizing background is chosen to give $`e^\tau =0.64`$ (Press et al. (1993)).
An example of a somewhat realistic reduced QSO spectrum ($`\widehat{N}_\mathrm{Q}^\alpha `$ in eq. ) and its error array (eq. ) can be found in the bottom two panels of Fig. 2 (ignore the other two panels for the moment). They are generated based on the prescriptions given in §3.1, assuming $`W^{\alpha \beta }=\delta ^{\alpha \beta }`$ and $`W^{i\beta }`$ is given by eq. (6). Briefly speaking, what we do is to first generate an array of $`g^\alpha `$ which represents $`g_\mathrm{b}^\alpha _jg_{\mathrm{ps}}^{j\alpha }`$ (i.e. we do not actually simulate the full two-dimensional CCD image; the spatial dimension is collapsed into $`g^\alpha `$); then, we create a Poisson realization of the (intermediate) quasar count $`g^\alpha N_C^\alpha e^{\tau _\alpha }`$ where $`N_C`$ is the continuum and $`e^{\tau _\alpha }`$ is predicted by our cosmological model; we similarly create a Poisson realization of the background count $`g^\alpha \times \mathrm{const}.`$ where the constant represents some fraction less than 1 and then subtract from it its Poisson mean, the end-result is then added to the above quasar count; lastly, we divide by $`g^\alpha `$ to obtain the reduced quasar count $`\widehat{N}_\mathrm{Q}^\alpha `$.
Note that the overall level of the reduced quasar count can be scaled up or down (because we are not interested in the absolute brightness of the quasar), provided the error array is scaled accordingly to conserve signal-to-noise (S/N; count divided by square root of the variance). This example resembles a high quality Keck HIRES spectrum, with S/N reaching up to 100 at certain pixels. It is composed of 12 echelle orders, $`50\AA `$ each (e.g. the instrument HIRES on the Keck telescope is an echelle spectrograph which consists of two diffraction gratings crossed at $`90^0`$ to each other; see Vogt et al. Vogt (1994)). The pixel size is $`0.05\AA `$ with a resolution Full-Width-Half-Maximum (FWHM) of $`0.125\AA `$. The example represents a case in which a relative calibration (but not necessarily absolute fluxing) between the orders has been attempted. The dashed line in the bottom panel shows the input continuum.
The error-array shows a lot of variations. A model of the random error as Gaussian distributed with uniform S/N that is sometimes used in the literature misses much structure. About $`3\%`$ of the spectrum consists of gaps which arise due to severe cosmic-ray hits. The spikes in the error-array correspond to wavelengths at which a fraction, but not all, of the corresponding spatial pixels are affected by cosmic-ray hits. They also take up $`3\%`$ of the spectrum. It is easy to see how these spikes arise from eq. (6) and (9). At wavelengths where some of the spatial pixels are thrown out because of cosmic-ray hits, $`W_{i\beta }`$ is enhanced because the sum over $`j`$ in its denominator is restricted to fewer pixels. Since it is the square of $`W_{i\beta }`$ that enters into the variance, a modest enhancement becomes a spike. Note how for each echelle order, the S/N drops towards the two ends. This is because of the blaze function which tends to suppress the flux at the ends. Note also that the S/N has a general decline towards the blue. This is due to a combination of a falling continuum, and decreasing detector efficiency.
Sometimes, a relative calibration between echelle orders is either difficult or simply not attempted. An example is shown in Fig. 3. Note how the continuum is broken into 12 discontinuous pieces.<sup>2</sup><sup>2</sup>2In some cases where the different echelle orders overlap, there could be two jumps at each order junction. These are taken from continuum-fits to actual data.
An example of data with much poorer quality is shown in Fig. 7. The pixel size is $`0.5\AA `$ and the FWHM is $`1.17\AA `$. The S/N is about 10 times worse than the two examples above. Such a spectrum could be the output of, say, a low dispersion single-grating spectrograph, which does not have the characteristic division into short pieces as in the case of the echelle spectrograph.
All other simulated data in this paper are slight variations of the above, which will be described in turn at the appropriate places.
## 4. Estimating the Power Spectrum / Two-point Correlation
### 4.1. The Quadratic Estimator
Given the one-dimensional array of estimated quasar counts $`\widehat{N}_\mathrm{Q}^\alpha `$ (eq. ), how should one go about estimating the two-point correlation or the power spectrum?
A common practice is to first continuum-fit, i.e. to estimate the continuum level $`\widehat{N}_\mathrm{C}^\alpha `$ and divide $`\widehat{N}_\mathrm{Q}^\alpha `$ by it to obtain an estimate of the transmission $`\widehat{f}=\widehat{N}_\mathrm{Q}^\alpha /\widehat{N}_\mathrm{C}^\alpha `$. Then, the estimators for the un-normalized two-point correlation and power spectrum (eq. ) are:
$$\widehat{\xi }_{\mathrm{un}}(u)=\underset{\alpha ,\beta }{}w^{\alpha \beta }(u)\widehat{f}^\alpha \widehat{f}^\beta ,\widehat{P}_{\mathrm{un}}(k)=\underset{\alpha ,\beta }{}w^{\alpha \beta }(k)\widehat{f}^\alpha \widehat{f}^\beta b(k)$$
(10)
where $`b(k)`$ subtracts out the shot-noise (i.e. a bias), and $`w^{\alpha \beta }(u)`$ and $`w^{\alpha \beta }(k)`$ are weighting kernels for which we will give some examples shortly (to be distinguished from $`W^{\alpha \beta }`$ in eq. ). These are commonly called quadratic estimators for the simple fact that they make use of the data $`\widehat{f}^\alpha `$ in quadratic combinations.
Defining the mean transmission to be $`\overline{f}`$, the obvious extensions of the above estimators, for the normalized two-point correlation and power spectrum (eq. ; unless otherwise stated, the two-point correlation or power spectrum with no qualifications refer to the normalized version), are
$$\widehat{\xi }_1(u)=\underset{\alpha ,\beta }{}w^{\alpha \beta }(u)(\widehat{f}^\alpha \overline{f})(\widehat{f}^\beta \overline{f})/\overline{f}^2,\widehat{P}_1(k)=\underset{\alpha ,\beta }{}w^{\alpha \beta }(k)(\widehat{f}^\alpha \overline{f})(\widehat{f}^\beta \overline{f})/\overline{f}^2b(k)$$
(11)
However, the form of the power spectrum or two-point correlation estimator given above suggests an interesting variation which allows us to avoid continuum-fitting altogether: $`(\widehat{f}^\alpha \overline{f})/\overline{f}`$ can be estimated instead by $`(\widehat{N}^\alpha \overline{N}^\alpha )/\overline{N}^\alpha `$ where $`\overline{N}^\alpha `$ is the mean count defined by $`\overline{N}^\alpha \widehat{N}^\alpha `$. Here $``$ denotes the cosmic average i.e. this corresponds to averaging out the cosmic fluctuations in $`f`$, for a fixed continuum ($`\widehat{N}^\alpha =N_C^\alpha \overline{f}`$ where $`N_C^\alpha `$ is the true continuum). <sup>3</sup><sup>3</sup>3It is implicitly assumed that discrete averaging has been carried out before cosmic averaging i.e. we use $``$ interchangeably with $`_D`$. See §2. Note that the mean count is dependent upon $`\alpha `$ because of the slowly varying continuum. We will discuss how to estimate $`\overline{N}^\alpha `$ shortly. The key here is that the absolute level of the continuum gets divided out by definition. Hence, let us define the following alternative estimators of the two-point correlation and power spectrum:
$`\widehat{\xi }_2(u)={\displaystyle \underset{\alpha ,\beta }{}}w^{\alpha \beta }(u)\widehat{\delta }_f^\alpha \widehat{\delta }_f^\beta ,\widehat{P}_2(k)={\displaystyle \underset{\alpha ,\beta }{}}w^{\alpha \beta }(k)\widehat{\delta }_f^\alpha \widehat{\delta }_f^\beta b(k)`$ (12a)
$`\widehat{\delta }_f^\alpha (\widehat{N}_\mathrm{Q}^\alpha \overline{N}_\mathrm{Q}^\alpha )/\overline{N}_\mathrm{Q}^\alpha `$ (12b)
This alternative power spectrum estimator is what we will focus on, but we will also briefly investigate the behavior of the estimators in eq. (10) and (11).
It remains to specify what $`w^{\alpha \beta }(u)`$, $`w^{\alpha \beta }(k)`$ and $`b(k)`$ are. The simplest choice is to use uniform weighting i.e. for the two-point correlation, it corresponds to:
$`w^{\alpha \beta }(u)=\mathrm{\Theta }^{\alpha \beta }(u)/{\displaystyle \underset{\mu \nu }{}}\mathrm{\Theta }^{\mu \nu }(u)`$ (13)
where $`\mathrm{\Theta }^{\alpha \beta }(u)`$ is equal to one if the two pixels $`\alpha `$ and $`\beta `$ are separated by a distance $`u`$ (or more generally, the distance falls into a bin that is centered around $`u`$ with some finite width), and zero otherwise. Using the above $`w^{\alpha \beta }(u)`$ corresponds to simply counting all pairs separated by a distance $`u`$, normalized by the total number of pairs.
With the above weighting, eq. (12) is analogous to an estimator of the two-point correlation introduced by Landy & Szalay Landy & Szalay (1993) for galaxy surveys: $`(DD2DR+RR)/RR`$, if one identifies $`DD`$ with $`_{\alpha \beta }w^{\alpha \beta }(u)\widehat{N}^\alpha \widehat{N}^\beta `$, $`DR`$ with $`_{\alpha \beta }w^{\alpha \beta }(u)\widehat{N}^\alpha \overline{N}^\beta `$ and $`RR`$ with $`_{\alpha \beta }w^{\alpha \beta }(u)\overline{N}^\alpha \overline{N}^\beta `$ and assumes $`\overline{N}^\alpha `$ varies very slowly with $`\alpha `$ on the scale of interest $`u`$ (the analogy becomes exact in the limit of a uniform $`\overline{N}^\alpha `$). As shown by Landy & Szalay Landy & Szalay (1993) (see also Szapudi & Szalay (1998); Dodelson et al. (1997)), a common alternative estimator $`DD/RR1`$ (equivalent to the estimator used by e.g. Zuo & Bond (1994); Cen et al. (1998)) is actually less desirable as it gives a larger variance compared to $`(DD2DR+RR)/RR`$.
With this being said, we are going to focus our attention on the power spectrum from now on, although most of our treatments below can be applied to the two-point correlation as well. For the power spectrum, the simplest choice of uniform weighting corresponds to:
$`w^{\alpha \beta }(k)=(/𝒩^2)R^{\alpha \beta }(k),R^{\alpha \beta }(k)(1/n_k){\displaystyle \underset{k^{}}{}}e^{ik^{}(u^\alpha u^\beta )}`$ (14a)
$`b(k)={\displaystyle \frac{}{𝒩^2}}{\displaystyle \underset{\alpha }{}}{\displaystyle \frac{q^\alpha \overline{N}_Q^\alpha +V_B^\alpha }{(\overline{N}_Q^\alpha )^2}},q^\alpha {\displaystyle \underset{i}{}}(W^{i\alpha })^2g_{\mathrm{ps}}^{i\alpha }g_\mathrm{b}^\alpha ,V_B^\alpha {\displaystyle \underset{i}{}}(W^{i\alpha })^2V_B^{i\alpha }`$ (14b)
where $``$ is the total length of the spectrum (in whatever units one prefers) and $`𝒩`$ is the total number of spectral pixels, $`R^{\alpha \beta }(k)`$ represents an average of the Fourier basis over some bin or band in k-space i.e. we estimate the power spectrum at a bin centered at $`k`$ by averaging over contributions from each $`k^{}`$ that belongs to the bin ($`n_k`$ is the total number of modes in it). <sup>4</sup><sup>4</sup>4See Seljak Seljak (1998) and Bond et al. Bond et al. (1997) for discussions on precautions one should take on binning. The symbol $`b(k)`$ represents the shot-noise contribution to the power that has to be subtracted off, $`g_{\mathrm{ps}}^{i\alpha }`$ and $`g_\mathrm{b}^\alpha `$ represent the point-spread function and the blaze as in eq. (5), and $`V_B^\alpha `$ is the background contribution to the shot-noise (including the sky and readout, see eq. ). Note how a part of the shot-noise depends on the reciprocal of the mean quasar count, as expected for Poisson fluctuations, quite analogous to the shot-noise of galaxy distributions. <sup>5</sup><sup>5</sup>5See e.g. Feldman et al. Feldman et al. (1994). However, the factor $`q^\alpha `$, which arises from non-trivial weighting of CCD pixels ($`W^{i\alpha }`$; eq. ), signifies that the shot-noise is not strictly Poisson distributed. Moreover, there are extra contributions to the shot-noise from the background counts, which are generally absent in the case of galaxy surveys. Derivations of the above statements are given in Appendix A.
The corresponding power spectrum estimator obeys:
$`\widehat{P}_2(k)={\displaystyle \frac{dk^{}}{2\pi }G(k,k^{})P(k^{})},G(k,k^{}){\displaystyle \underset{\alpha \beta }{}}w^{\alpha \beta }(k)e^{ik^{}(u^\alpha u^\beta )}`$ (15)
where $`G`$ is a window function that resembles, for $`k1/`$, a delta function centered at $`k=k^{}`$ with a width of the order of $`1/`$. The normalization of $`w^{\alpha \beta }(k)`$ in eq. (14a) ensures that $`𝑑k^{}G(k,k^{})/2\pi =1`$. See Appendix A for a derivation.
It should be emphasized, however, that the above statements are strictly true only if one ignores uncertainties in the mean count $`\overline{N}_Q^\alpha `$ – i.e. $`\overline{N}_Q^\alpha `$ is not known a priori but is instead estimated from the same data from which one is trying to measure correlations. We will discuss this further in §4.2.2. It suffices to say here that our results in this section remain valid as long as one stays away from scales comparable to the entire length of the quasar spectrum.
### 4.2. Systematic Errors
#### 4.2.1 Continuum-fitting versus Trend-removal
The power spectrum estimator $`\widehat{P}_2`$ in eq. (12), which we are going to focus most of our attention on, requires an estimate of mean count $`\overline{N}_Q^\alpha `$. The mean count is not strictly uniform because of a slowly fluctuating continuum i.e. $`\overline{N}_Q^\alpha =N_C^\alpha \overline{f}`$ where $`N_C^\alpha `$ is the continuum and $`\overline{f}`$ is the (flat) mean transmission. We assume $`\overline{N}^\alpha `$ has the following form:
$$\overline{N}_\mathrm{Q}^\alpha =\underset{a}{}C^ap^{a\alpha }$$
(16)
where $`p^0`$ is a constant, $`p^1`$ is the first order polynomial ($`p^{1\alpha }=u^\alpha `$), $`p^2`$ is the second order polynomial ($`p^{2\alpha }=(u^\alpha )^2`$), and so on. The coefficients $`C^a`$ need to be estimated from the quasar counts $`\widehat{N}_Q^\alpha `$. Note that most of our following arguments would go through for a different set of basis functions. One key assumption we will exploit is that $`\overline{N}^\alpha `$ is slowly fluctuating, so that we can truncate the above series at relatively low orders. Continuum-fitting in practice makes the same assumption.
To estimate $`C^a`$, we use a linear estimator:
$$C^a=\underset{\alpha }{}M^{a\alpha }\widehat{N}_\mathrm{Q}^\alpha $$
(17)
where $`𝐌`$ is a matrix to be specified. Comparing eq. (16) and (17), it is not hard to see that $`𝐌`$ has to satisfy $`_\alpha M^{a\alpha }p^{b\alpha }=\delta ^{ab}`$. The simplest choice is to adopt, in vector notation, $`𝐌=(\mathrm{𝐩𝐩}^𝐓)^\mathrm{𝟏}𝐩`$ where $`\mathrm{𝐩𝐩}^𝐓`$ in component-form is $`[pp^T]^{ab}=_\alpha p^{a\alpha }p^{b\alpha }`$. In summary, this means our estimator for the mean quasar count is
$$\overline{N}_\mathrm{Q}^\alpha =\underset{\beta }{}L^{\alpha \beta }\widehat{N}_\mathrm{Q}^\beta ,𝐋𝐩^𝐓(\mathrm{𝐩𝐩}^𝐓)^\mathrm{𝟏}𝐩$$
(18)
where $`𝐋`$ in component-form reads $`L^{\alpha \beta }_{ab}p^{a\alpha }\stackrel{~}{p}^{ab}p^{b\beta }`$ with $`\stackrel{~}{p}^{ab}`$ being the inverse of the matrix $`[pp^T]^{ab}`$. More sophisticated versions of the above can be found in Rybicki & Press (Rybicki & Press (1992)). Our experience is that the simple version given here suffices, because the shape of the true quasar continuum is quite uncertain anyway.
Note the crucial differences between traditional continuum-fitting and an estimation of the mean count as described above. The above method makes no reference to the absolute level of the continuum i.e. the count level where there is supposedly no absorption. Continuum-fitting in practice often involves human intervention (eye-balling) in the identification of such a level. In contrast, eq. (18) is straightforward to implement and automate. The mean count is then used to normalize the quasar count as in eq. (12b) before the power spectrum is estimated (eq. ). We call this procedure trend-removal to distinguish it from traditional continuum-fitting. Trend-removal is widely used in other areas (e.g. Press et al. (1992a); Rybicki & Press (1992); Tegmark et al. (1998)). It is akin to the estimation of, say, the long-term trend of the stock market in the midst of all its daily fluctuations.
Eq. (18), together with eq. (12) and (14), completely specifies the main power spectrum estimator we advocate. Several tests follow.
Test 1
In Fig. 2, we show the effect of different choices of the mean-transmission basis $`𝐩`$ (eq. ). The simulated spectrum is of Keck-HIRES quality, with a S/N as high as 100 at certain pixels, and it assumes one has a good relative calibration between the different echelle orders (12 in all) i.e. an almost ideal, state-of-the-art observed spectrum. The second panel from the top shows the recovery of the mean transmission. The solid line represents the true (input) mean. The rest shows the recovered mean for different bases: dotted line for a basis consisting of $`p^0`$ only (a constant i.e. the continuum or the mean is modeled as completely flat); short-dashed line for a basis consisting of polynomials up to the third order; long dashed line for also a basis of polynomials up to the third order but with coefficients fitted separately for each echelle order. The short-dashed line seems to give the best match to the true mean. However, none of them is perfect because the true mean does not, by choice, have a polynomial shape. This is what is likely to happen in practice – lacking a good understanding of the physics that determines the continuum shape of any given quasar, the best one can do is to pick a reasonable basis which contains enough freedom to describe the general features of the continuum, but not so much freedom that one overfits.
The important question is what impact the choice of basis has on power spectrum estimation. This is illustrated in the top panel of Fig. 2, where the fractional error in the measured power spectrum is shown. The one that gives the best match to the true power spectrum is indeed the one where a simple basis of $`p^0`$, …, $`p^3`$ is used for the whole length of the simulated spectrum. The biggest effect of underfitting (dotted line) or overfitting (long dashed) the mean transmission is on the power spectrum estimation on large scales. They cause respectively over- or under-estimation of the large scale power spectrum. An additional effect is that overfitting tends to introduce spurious power on small scales as well – witness the enhanced fluctuations in the error on small scales for the long-dashed line. We will see this more clearly in the next Test.
Without any prior knowledge of the intrinsic continuum shape of an observed quasar, how does one decide if one is overfitting or underfitting? One way is to look at the region of the observed spectrum redward of the Ly$`\alpha `$ emission line, which is free of the forest, and the continuum is therefore relatively easy to reconstruct. Assuming the general level of continuum-fluctuation is the same both redward and blueward of Ly$`\alpha `$, one can then gain an idea of what a good mean-transmission basis might be. Low redshift QSO spectra, where the continuum can be quite easily recovered even blueward of Ly$`\alpha `$, can also be used to gauge the scales at which the continuum fluctuates – naturally, one could also use such spectra to check the assumption that continuum-fluctuations have similar characteristics redward and blueward of Ly$`\alpha `$ (more on this below).
Test 2
In Fig. 3, we took the continuum-fits to an observed quasar spectrum and use them as the input continuum for our simulation. The simulated spectrum here represents a case in which no relative calibration between echelle orders have been attempted, which is often the case. This is why the continuum in the bottom panel is broken up into 12 pieces. The second panel from the top again illustrates the recovery of the mean transmission: dotted line for a flat model-continuum for each order, and dashed line for a basis of polynomials up to the third order, also separately for each order. The solid line is the true mean transmission. The top panel shows the accuracy of the corresponding power spectrum estimations. The assumption of a simple flat continuum for each order gives a power spectrum that is accurate to $`1\%`$. On the other hand, overfitting with up to third order polynomials not only causes an under-estimation of power on large scales, but also creates spurious power on small scales.
Combining Fig. 2 & 3 (note that they show the power spectrum estimation on different scales), the lessons are: 1. it is better to err on the side of underfitting the mean, which tends to over-estimates the power on large scales, but leave the power on small scales relatively unaffected (this relies crucially on the fact that the continuum has fluctuations only on large scales); 2. without sufficient prior knowledge of the true shape of the continuum, one can at least make conservative statements about the small-scale power, but the large-scale power is likely prone to systematics, unless some correction is made.
One additional comment: the input continuum in Fig. 3, which is taken from fits to actual data, certainly seems to suggest that the observed continuum has fluctuations on scales of an echelle order ($`50\AA `$). (We will quantify this better in §4.2.2.) It is unclear whether this is truly due to the intrinsic continuum, or whether it is an artifact of imperfect blaze removal or flat-fielding (see §3.1). If it is the former, then $`50\AA `$ represents a fundamental limit beyond which one cannot reliably measure the transmission power spectrum, at least not without some additional prior knowledge of the true continuum (which is what we will discuss in §4.2.2). If it is the latter, then in principle one should be able to do better and extend the range of reachable scales to larger ones. Which is the case remains to be seen.
Test 3
In Fig. 4, we show the effect of traditional continuum-fitting, which requires some degree of eye-balling. The same simulated spectrum as in Fig. 3 is given to an observer (one of the authors) with no knowledge of the input continuum. Note that the top second panel now shows the actual continuum level rather than the mean transmission level. The estimated continuum actually matches the true one surprisingly well. However, one can still see that the continuum is generally underestimated. In the top panel, we show the accuracy of two different power estimates. The long-dashed line corresponds to an estimate of the un-normalized power spectrum as defined in eq. (1) (the estimator is eq. ). There is clearly a $`5\%`$ positive bias here because of the underestimation of the continuum. One way to correct for it is of course to use simulations: applying exactly the same procedure to the observed data and the simulated data, and see how much bias results; but the size of the bias is likely to be model dependent. A simple alternative way to cure this problem is to measure the normalized power spectrum instead, using the continuum-fitted data, i.e. using the estimator in eq. (11). This is shown with a short-dashed line. It has an accuracy of $`1\%`$, comparable to the dotted line in Fig. 3. In view of this, it seems by-passing continuum-fitting altogether and proceeding simply with trend-removal is desirable.
Test 4
The failure of traditional continuum-fitting is more dramatic in cases where there is a lot of absorption e.g. at high redshifts. In Fig. 5 is shown a simulated spectrum with the ionizing background adjusted to give a mean transmission of 0.39, which is about the observed value at $`z=4`$ (Press et al. (1993)). The continuum is more seriously under-estimated leading to an overestimate of the un-normalized power spectrum by $`20\%`$ (the upper long-dashed line). The normalized power spectrum, estimated either using the continuum-fitted data (eq. ) or using directly the trend-removed data (eq. ), is much more accurately measured.
Test 5
Another example in which traditional continuum-fitting fails is shown in Fig. 6. This is based on the same spectrum as in Fig. 1, but convolved with a Gaussian of $`1.17\AA `$ FWHM and with much poorer S/N compared to the simulated spectra above. This is likely not the product of an echelle spectrograph, hence there is no division into 12 orders. We repeat the exercise of continuum-fitting and then power spectrum measurement as before. Interestingly, the significant discrete fluctuations due to the low S/N here cause an overestimation (unlike in Test 3 and 4) of the continuum level, and so an underestimation of the un-normalized power spectrum. Once again, the normalized power spectrum does not suffer from the same problem. Note the somewhat large fluctuations of the estimated power – this is largely due to the high level of shot-noise.
Fig. 7 shows the measurement of power spectrum through trend removal instead. A third order polynomial is used to estimate the mean transmission. The resulting (normalized) power spectrum estimate (eq. ) is of comparable accuracy to that using the continuum-fitted data. We also show in the top panel as a dotted line the power spectrum estimate without shot-noise subtraction (eq. ). Clearly, shot-noise introduces a bias of the order of $`10\%`$ here. We will have some more to say about this in §5.
Tests 4 and 5 above drive home the point that the bias of an estimate of the un-normalized power spectrum from continuum-fitted data is highly variable. It depends on the redshift, resolution as well as S/N of the data.
There have been in the literature discussions of an alternative method to normalize the quasar count: normalizing by the maximum value of the continuum-fitted count, instead of by the mean count (e.g. McDonald et al. (1999)). Note that this procedure is also sensitive to the S/N and resolution of the data. For instance, $`\mathrm{max}(\widehat{N}_\mathrm{Q}^\alpha /\widehat{N}_\mathrm{C}^\alpha )=1.4`$ in Fig. 6, while $`\mathrm{max}(\widehat{N}_\mathrm{Q}^\alpha /\widehat{N}_\mathrm{C}^\alpha )=1.12`$ in Fig. 4, where $`\widehat{N}_\mathrm{C}^\alpha `$ is the estimated continuum – they share exactly the same underlying cosmic signal but the former has a higher level of discrete fluctuations and poorer resolution – for reference, the true maximum transmission should be 0.99. This means one should take care to simulate the noise properties correctly (e.g. Rauch et al. (1997)).
Lastly, we should emphasize that while trend-removal seems to be more desirable than traditional continuum-fitting for the particular application here, continuum-fitting is still very useful for other purposes, which we will discuss in §4.2.2 and §5. But a fully automated procedure for continuum-fitting is clearly desirable.
#### 4.2.2 A Bonus of Trend-removal – Power Correction on Large Scales
As is clear from some of the previous tests in §4.2.1, the power spectrum measured on large scales (i.e. scales comparable to the typical scales where the continuum has fluctuations) could contain spurious contributions from the continuum, the size of which depends somewhat on the continuum/mean-shape-model one assumes. The strategy adopted in the §4.2.1 is a conservative one: assume a model for the continuum that is as simple (or smooth) as possible, perform trend-removal, and the resulting power spectrum would reflect the true transmission power spectrum at least on small scales, but not necessarily on large scales.
Can we do better? The answer is yes, under certain assumptions which we will make explicit shortly, and it illustrates an added benefit of trend-removal as introduced in the §4.2.1. Readers not interested in the details can skip directly to the end of this section where two examples of how our technique of power-correction works are given (Fig. 9 and 10).
Let us start by recalling the power spectrum estimator in eq. (12), but focusing now on the fact that the true $`\overline{N}^\alpha `$ is unknown, and has to be estimated using eq. (18), which assumes implicitly that the true mean count obeys eq. (16), which is of course only a reasonable guess. Let us use $`\widehat{\overline{N}_{\mathrm{Q}}^{}{}_{}{}^{\alpha }}`$ to denote the estimated mean count, which generally differs from the true mean count $`\overline{N}_\mathrm{Q}^\alpha `$. We have used $`\overline{N}^\alpha `$ somewhat sloppily before when we really meant $`\widehat{\overline{N}_{\mathrm{Q}}^{}{}_{}{}^{\alpha }}`$ e.g. eq. . In other words, eq. should be more accurately written as
$$\widehat{\overline{N}_\mathrm{Q}^\alpha }=\underset{\beta }{}L^{\alpha \beta }\widehat{N}_\mathrm{Q}^\beta ,𝐋𝐩^𝐓(\mathrm{𝐩𝐩}^𝐓)^\mathrm{𝟏}𝐩$$
(19)
where $`𝐩`$ represents the basis functions. Similarly, the estimator for the power spectrum in eq. (12) should be more accurately written as:
$`\widehat{\xi }_2(u)={\displaystyle \underset{\alpha ,\beta }{}}w^{\alpha \beta }(u)\widehat{d}_f^\alpha \widehat{d}_f^\beta ,\widehat{P}_2(k)={\displaystyle \underset{\alpha ,\beta }{}}w^{\alpha \beta }(k)\widehat{d}_f^\alpha \widehat{d}_f^\beta b(k)`$ (20a)
$`\widehat{d}_f^\alpha (\widehat{N}_\mathrm{Q}^\alpha \widehat{\overline{N}_\mathrm{Q}^\alpha })/\widehat{\overline{N}_{\mathrm{Q}}^{}{}_{}{}^{\alpha }}={\displaystyle \underset{\gamma }{}}D^{\alpha \gamma }\widehat{N}_\mathrm{Q}^\gamma /{\displaystyle \underset{\eta }{}}L^{\alpha \eta }\widehat{N}_\mathrm{Q}^\eta ,D^{\alpha \gamma }\delta ^{\alpha \gamma }L^{\alpha \gamma }`$ (20b)
We now assume the following quantities are small: $`\widehat{d}_f^\alpha `$ and $`(L^{\alpha \gamma }\widehat{N}_\mathrm{Q}^\gamma \overline{N}_\mathrm{Q}^\alpha )/\overline{N}_\mathrm{Q}^\alpha `$. The second quantity tells us how far off our estimate of the mean is from the true mean, while the first contains contributions both from the fluctuation in transmission and from the second quantity. Therefore, putting eq. (20b) and eq. (20a) together, the lowest order contributions to the expectation value of the estimator $`\widehat{P}_2(k)`$ are:
$`\widehat{P}_2(k)`$ $`=`$ $`{\displaystyle \underset{\alpha \beta \gamma \eta }{}}w^{\alpha \beta }(k)D^{\alpha \gamma }D^{\beta \eta }{\displaystyle \frac{\overline{N}_\mathrm{Q}^\gamma \overline{N}_\mathrm{Q}^\eta }{\overline{N}_\mathrm{Q}^\alpha \overline{N}_\mathrm{Q}^\beta }}(1+\widehat{\delta }_f^\gamma \widehat{\delta }_f^\eta )`$
$`=`$ $`P_\mathrm{C}(k)+{\displaystyle \frac{dk^{}}{2\pi }P(k^{})G_n(k,k^{})}`$
$`P_\mathrm{C}(k)`$ $``$ $`{\displaystyle \underset{\alpha \beta \gamma \eta }{}}w^{\alpha \beta }(k)D^{\alpha \gamma }D^{\beta \eta }{\displaystyle \frac{\overline{N}_\mathrm{Q}^\gamma \overline{N}_\mathrm{Q}^\eta }{\overline{N}_\mathrm{Q}^\alpha \overline{N}_\mathrm{Q}^\beta }}`$ (21b)
$`G_n(k,k^{})`$ $``$ $`{\displaystyle \underset{\alpha \beta \gamma \eta }{}}w^{\alpha \beta }(k)e^{ik^{}(u^\gamma u^\eta )}D^{\alpha \gamma }D^{\beta \eta }{\displaystyle \frac{\overline{N}_\mathrm{Q}^\gamma \overline{N}_\mathrm{Q}^\eta }{\overline{N}_\mathrm{Q}^\alpha \overline{N}_\mathrm{Q}^\beta }},`$ (21c)
where we have retained the old definition of $`\widehat{\delta }_f^\gamma `$ as $`(\widehat{N}_\mathrm{Q}^\gamma \overline{N}_\mathrm{Q}^\gamma )/\overline{N}_\mathrm{Q}^\gamma `$ (eq. ). The above gives an idea of how biased the estimator $`\widehat{P}_2(k)`$ is. Note that we have used $``$ here to include, in addition to the ensemble averaging as explained in §2, an averaging over the ensemble of possible continua (which changes $`\overline{N}_\mathrm{Q}^\alpha `$ because it is directly proportional to the continuum count $`N_\mathrm{C}^\alpha `$). We have assumed the fluctuations in the continuum are uncorrelated with fluctuations in the cosmic signal $`\widehat{\delta }_f^\gamma `$. We have also ignored the shot-noise contributions (e.g. $`b(k)`$) and will continue to do so for the rest of this section, because the scales where the continuum contamination could be a problem are typically large enough that shot-noise is subdominant.
The term $`P_\mathrm{C}(k)`$ can be viewed as the power spectrum of the continuum fluctuation. This is fluctuation in the sense of $`D^{\alpha \gamma }\overline{N}_\mathrm{Q}^\gamma =\overline{N}_\mathrm{Q}^\alpha L^{\alpha \gamma }\overline{N}_\mathrm{Q}^\gamma `$. This fluctuation would vanish if our trend-removal procedure were so accurate that the continuum shape is exactly reproduced. The term $`G_n(k,k^{})`$ is the effective window function, replacing the one in eq. (15), which does not take into account the error involved in trend-removal. The desirable normalization condition $`G_n(k,k^{})𝑑k^{}/(2\pi )=1`$ no longer holds with the choice of $`w^{\alpha \beta }`$ in eq. (14a). We have instead
$$G_n(k,k^{})\frac{dk^{}}{2\pi }=\frac{1}{𝒩}\underset{\alpha \beta \gamma }{}R^{\alpha \beta }(k)D^{\alpha \gamma }D^{\beta \gamma }\frac{(\overline{N}_\mathrm{Q}^\gamma )^2}{\overline{N}_\mathrm{Q}^\alpha \overline{N}_\mathrm{Q}^\beta }1+ϵ_G(k)$$
(22)
where $`R^{\alpha \beta }(k)`$ is defined in eq. (14a).
Assuming for now $`P_\mathrm{C}(k)`$ and $`ϵ_G(k)`$ can be measured from a suitable ensemble of continua, we propose the following alternative estimator to $`\widehat{P}_2(k)`$, which removes the bias due to continuum contamination:
$$\widehat{P}_3(k)=[\widehat{P}_2(k)P_\mathrm{C}(k)]/[1+ϵ_G(k)]$$
(23)
The above gives us an unbiased estimate of the windowed power spectrum. The window is effectively $`G_n(k,k^{})/[1+ϵ_G(k)]`$ which has the desirable normalization. We will not attempt further improvements such as deconvolution in this paper.
A useful alternative estimator, in cases where $`P_\mathrm{C}`$ dominates the bias in $`\widehat{P}_2(k)`$, is
$$\widehat{P}_4(k)=\widehat{P}_2(k)P_\mathrm{C}(k)$$
(24)
The above estimator gets rid of most of the bias in the estimator $`\widehat{P}_2(k)`$ if $`P_\mathrm{C}(k)/P(k)ϵ_G(k)`$. An interesting corollary is that, under such a condition, the bias in $`\widehat{P}_2(k)`$ is positive since $`P_\mathrm{C}(k)`$ is positive definite. Needless to say, this statement breaks down if $`P_\mathrm{C}(k)`$ is not the dominant source of bias, or if the fractional error in the mean-estimation is large (see e.g. Fig. 2).
It is interesting to compare our derivation above with the well-known one for the integral constraint bias in galaxy surveys (s.g. Peebles (1980); Landy & Szalay (1993); Bernstein (1994); Tegmark et al. (1998)). The integral constraint arises from the fact that the mean density of a galaxy survey has to be estimated from the same survey from which one is also measuring the power spectrum. The fact that the power spectrum estimator involves a non-trivial nonlinear combination of the data gives rise to a bias (see Hui & Gaztañaga (1999)), quite analogous to our derivation here. However, in the standard derivations, it is assumed the shape of the mean density is known (often taken to be uniform), and therefore $`P_\mathrm{C}(k)`$ effectively vanishes, whereas $`ϵ_G(k)`$ can be non-negligible on scales comparable to the size of the survey, but is otherwise small. The reader is referred to Bernstein Bernstein (1994) and Hui & Gaztañaga Hui & Gaztañaga (1999) for discussions on higher order contributions to the integral constraint.
How should one estimate $`P_\mathrm{C}(k)`$ and $`ϵ_G(k)`$? Given an ensemble of continua (with counts represented by $`N_\mathrm{C}^\alpha `$), our procedure is to replace $`\overline{N}_\mathrm{Q}^\alpha `$, which appears in the definitions of $`P_\mathrm{C}(k)`$ and $`ϵ_G(k)`$ (eq. \[21b\], ) with $`N_\mathrm{C}^\alpha `$, and compute the corresponding ensemble averages. Note that $`\overline{N}_\mathrm{Q}^\alpha =N_\mathrm{C}^\alpha e^\tau `$, but $`e^\tau `$, which is taken to be constant over the finite redshift range of interest, is divided out in the relevant definitions of $`P_\mathrm{C}(k)`$ and $`ϵ_G(k)`$.
The hard question is of course how to obtain a suitable ensemble of continua. The first thing one might try is to measure the power spectrum of the continuum-fits (i.e. $`P_\mathrm{C}(k)`$, or more generally, both $`P_\mathrm{C}(k)`$ and $`ϵ_G(k)`$) from exactly the same regions from which one attempts to measure the transmission power spectrum. While this can give us a crude idea of how significant the continuum power spectrum is, it is not entirely satisfactory because part of what has been ascribed to the continuum might actually be large scale fluctuations in the cosmic signal $`\delta _f^\alpha `$ that we are after, or vice versa.
The second option that comes to mind is to measure the continuum power spectrum from regions where the continuum determination is relatively secure. Two possibilities are A. low redshift quasar spectra where the forest is much less dense, and B. regions of spectra which lie redward of Ly$`\alpha `$ emission. The working hypothesis is that the continuum power spectrum in these two regions is the same as, or at least similar to, the one in the region where we attempt to estimate the transmission power spectrum (the forest of interest). There is no guarantee that the hypothesis is valid. For instance, regarding possibility A, the continuum power could systematically evolve with redshift. In fact it probably does: assuming that the statistical properties of the quasar continuum in rest frame do not evolve with redshift, the observed continuum power would evolve as $`P_\mathrm{C}[k,z_1]=P_\mathrm{C}[k(1+z_0)/(1+z_1),z_0]`$. One could in principle constrain such redshift evolution with a sufficiently large sample of low redshift quasar spectra. Regarding possibility B, it is not unreasonable to expect that the continuum power is higher on the red side compared to the blue side, because there are generally more broad emission lines on the red side (see e.g. Peterson (1997); Blandford, Netzer & Woltjer (1990); see below for caveats and a counter example, however). An upper bound on the blue continuum power is by itself interesting because one can then obtain a conservative estimate of how much spurious power is introduced by the continuum into one’s forest power measurements. Furthermore, systematic differences between the red and blue continuum power can be studied and quantified with a sufficiently large sample of low redshift quasars.
In Fig. 8, we show the continuum power spectrum measured from the continuum estimates on both sides of the Ly$`\alpha `$ emission of a quasar at z=3 (QSO 1157+3143). The continuum estimates are shown in the bottom two panels. After fitting a flat mean to each echelle order, we compute the continuum power spectrum just as if this were the forest, and the results from the red side and blue side are shown as solid and dotted lines respectively in the top panel. The two power spectra look similar. However, we emphasize that because of the lack of small scale power in the continuum, most of the power on small scales ($`k{}_{}{}^{}{}_{}{}^{>}\mathrm{\hspace{0.17em}\hspace{0.17em}1}\AA ^1`$) that we see in Fig. 8 is likely aliased from large scales. We will not attempt to perform a deconvolution to obtain the true small scale power; it suffices to note here that the true small scale power can only be smaller than what is shown in the figure. Also shown in the top panel is $`ϵ_G(k)`$, on both sides of Ly$`\alpha `$, which are basically indistinguishable from each other. Note that $`ϵ_G(k)1`$. The second panel from the top shows the fractional difference between $`P_\mathrm{C}`$ from the red and blue sides, which is about $`10\%`$, with the blue continuum power systematically higher than the red one. The results here, though drawn from admittedly a very small sample, are quite interesting for several reasons.
* The excess of the blue continuum power spectrum over the red one is consistent with the following hypothesis: that some of the fluctuations in the forest have been wrongly assigned to the continuum during the continuum-fitting process on the blue side. In other words, the true blue continuum power spectrum should be lower than the dotted line in the top panel of Fig. 8. An upper bound on the true blue continuum power spectrum is already very useful. One can use it to quantify how much, and on what scales, one should worry about spurious continuum power introduced into estimation of the transmission power. One can compare Fig. 8 with the theoretical expectation in Fig. 1, and see that the spurious power must be negligible for $`k{}_{}{}^{}{}_{}{}^{>}\mathrm{\hspace{0.17em}0.3}\AA ^1`$. This explains why the determination of the transmission power spectrum from both the continuum-fitted data or the trend-removed data is very accurate in examples like Fig. 5, as long as one considers the normalized power. Unfortunately, the pieces of continuum we examine are not long enough to yield useful information on larger scales or smaller $`k`$’s. If one takes a crude extrapolation, the continuum power spectrum (or more accurately its upper bound) might become non-negligible compared to the transmission power spectrum at $`k0.1\AA ^1`$. However, one must keep in mind that the theoretical transmission power spectrum in Fig. 1 is likely underestimated at small $`k`$’s because the simulation lacks large scale power. Nonetheless, there should be a genuine flattening of the transmission power spectrum at large scales. In any case, the first point to bear in mind is that an upper bound on the continuum power spectrum is useful as a conservative estimate of the possible spurious power. <sup>6</sup><sup>6</sup>6One should be aware of a possible pitfall of the above argument, however. It is not impossible that the opposite can happen, that one underfits the blue continuum, and ends up underestimating the blue continuum power. This is probably not the case here, where the data from which the blue continuum is estimated have high signal-to-noise and resolution (similar to the simulated spectrum in Fig. 4, where it can be seen that the continuum fit tends to have features that follow the forest). Underfitting the blue continuum is more likely for low resolution data, although even there, the situation is not clear: underfitting would result in underestimation of the continuum power on small scales, but not necessarily on large scales. Obviously, more tests are needed.
* Further, one can test the hypothesis that the excess in blue continuum power is due to contamination from the forest: if this is true, one expects the red and blue continuum power spectra to converge, as one goes to lower redshift quasars, because presumably, the blue continuum power spectrum should be less affected by the forest at lower redshifts. Even if their difference does not converge to zero (as suggested by the larger number of broad emission lines on the red side), but to some small but finite value, this is still a useful exercise because it gives us an idea of how different the red and blue continuum power spectra can be. If we can determine the blue continuum power spectrum to an accuracy of $`10\%`$ say, and use this to correct for the transmission power spectrum on large scales, this is already a significant improvement over not correcting for the large scale power, or simply throwing away the information on large scales altogether. For instance, if the blue continuum power does become comparable with the transmission power at $`k0.1\AA ^1`$, not subtracting off the spurious power would result in a fractional error of $`100\%`$, while subtracting off an approximate blue continuum power accurate to $`10\%`$ reduces the error by an order of magnitude.
Obviously, more testing using observed data is warranted, particularly on the estimation of red and blue continuum power as a function of redshift. This will be carried out in a separate paper. One natural question that might occur to the reader is whether a universal continuum power spectrum actually exists, given the large observed variations in the continuum from one quasar to another. It suffices to note that given an ensemble, the averaged power spectrum is always a well-defined quantity. The tricky part is to make sure the ensemble from which one estimates the continuum power spectrum has the same averaged continuum power spectrum as the ensemble of continua in the forest regions of interest. As a simple example: one might want to make sure the same proportion of radio-loud quasars are included in both ensembles. This is probably desirable if one uses the working hypothesis that low redshift blue continuum power is similar to high redshift blue continuum power, as suggested above. Alternatively, if the hypothesis that blue and red continuum power spectra resemble each other irrespective of redshift turns out to be a reasonable one, the simplest way to make sure one has the right ensemble is to use both sides of Ly$`\alpha `$ for any give quasar: use the blue side for its forest, and the red side for its continuum.
With all of the above caveats in mind, let us illustrate the technique of power correction with two idealized examples, where it is assumed the right ensemble of continua is in hand.
In Fig. 9, we show in the bottom panel a simulated spectrum with a somewhat unusual continuum (middle panel) with a fair amount of fluctuations. We generate a set of 10 different continua and impose each on our underlying cosmic absorption to obtain a set of 10 different simulated spectra (only one of which is shown in the figure). We compute the power spectrum using $`P_2(k)`$ as in eq. (20). The resulting fractional error from the true transmission power spectrum is shown as a solid line in the top panel. There is clearly a lot of spurious power on large scales due to the imperfectly estimated mean count, which reflects the wild fluctuations in the continuum. We then apply the power spectrum corrections: the dotted line shows $`\widehat{P}_3(k)`$ from eq. (23) while the dashed line shows $`\widehat{P}_4(k)`$ from eq. (24). One can see that subtracting the continuum power spectrum $`P_\mathrm{C}(k)`$ alone removes most of the spurious power.
To make the example realistic, we have multipled the continua in the forest region by a power law that goes like $`(u^\alpha )^{0.96}`$ (i.e. the ’blue’ continuum), and similarly multiplied the continua from which we actually estimate the continuum power by a power law of $`(u^\alpha )^{0.01}`$ (i.e. the hypothetical ’red’ continuum). This is meant to mimic a possible turn-over of the quasar continuum around Ly$`\alpha `$ (see e.g. Zheng et al. (1998) for evidence of a turn-over around Ly$`\beta `$). We have in mind a situation in which the continuum power spectrum $`P_\mathrm{C}`$ is estimated from the red side of Ly$`\alpha `$. Clearly, the fact that the mean trends on the blue and red are different does not present an obstacle.
In Fig. 10 we show a similar version of the above, but with much noisier data and poorer resolution, and a mean power-law of $`(u^\alpha )^{1.5}`$ and a mean of $`(u^\alpha )^{0.9}`$ have been imposed on the continua on the blue and red sides respectively. The same technique works here as well.
One last point we should make: when the quantities $`P_\mathrm{C}(k)`$ and $`ϵ_G(k)`$ are estimated from some ensemble of continua, they in general receive shot-noise contributions. We have ignored shot-noise here, assuming the scales where power correction is most interesting are sufficiently large that shot-noise is unimportant. This should be checked on a case by case basis.
#### 4.2.3 Gaps and Metal Absorption Lines
There are at least two other possible sources of systematic errors in addition to that due to continuum-fitting. Gaps are quite common in observed spectra due to defects in the CCD, incomplete spectral coverage, or cosmic ray hits. Fortunately, since they are at known locations, we can either consider only those parts of the spectrum that are between the gaps (for instance when the gaps are large), or interpolate to fill in the gaps (for instance when the gaps are small). The latter is what we have implicitly done in all of the tests mentioned in §4.2.1, where $`3\%`$ of the pixels are assumed discarded because of cosmic-ray hits. The hits are typically one to a few pixels wide, and we simply fill them in by linearly interpolating the counts from neighboring pixels. Clearly, we can recover the power spectrum quite well in spite of the need to interpolate.
A more challenging problem is possible systematics due to the presence of metal absorption lines. Shown in the bottom panel of Fig. 11 is a simulated spectrum with resolution and S/N very similar to that of Fig. 3 except that metal absorption lines as shown in the panel above have been added on top of the Ly$`\alpha `$ forest. This list of lines is taken from a quasar spectrum redward of Ly$`\alpha `$ (HS 1103+6416, z = 2.19). The mean transmission is estimated by assuming a flat trend for each echelle order as before. What is interesting is the dotted line in the top panel, demonstrating the creation of spurious power by the metal lines. The dashed / solid line shows fractional error in the power spectrum estimate if all metal lines with $`\tau >0.4`$ / $`\tau >1`$ are assumed “detected”, and therefore cut-out and treated as gaps as before (i.e. interpolated across). Such a procedure seems to eliminate much of the spurious power. In practice, sufficiently strong metal lines should be identifiable by their narrow widths.
Fig. 12 shows that metal absorption lines in data with lower resolution and poorer signal to noise have a relatively small effect on the accuracy of the power spectrum estimation.
### 4.3. Random Errors: Shot-noise-bias, Variance and Minimum Variance Weighting
Random errors arise firstly from (cosmic) sampling fluctuations, and secondly from electron/photon counting, which can be traced to fluctuations in the intrinsic quasar counts, the sky counts and the readout (see §3.1). We will summarily refer to the latter as shot-noise. Shot-noise affects two aspects of power spectrum estimation.
First, shot-noise introduces a bias which has to be subtracted off. This is the term $`b(k)`$ in eq. (14b). We will give here a more general expression for $`b(k)`$ suitable for different weightings ($`w^{\alpha \beta }`$). As we have demonstrated in Fig. 7, shot-noise-bias subtraction can be important for low S/N data. We will return to this point in §5.
Second, shot-noise, together with cosmic fluctuations, determines the variance of the power spectrum estimate. We will give the expression for the variance in this section, and then address the question of how to best combine data with different levels of S/N to minimize the variance.
The power spectrum estimator we will focus on is given in eq. (12). It is assumed trend removal as explained in §4.2.1 has been performed. We ignore uncertainties due to the unknown continuum in this section. Here we do not limit ourselves to the choice of uniform weighting (eq. ) as we have done so far. Let us rewrite $`w^{\alpha \beta }(k)`$ in eq. (12) as
$$w^{\alpha \beta }(k)=\overline{w}^{\alpha \beta }(k)R^{\alpha \beta }(k)$$
(25)
where $`R^{\alpha \beta }(k)`$ is given in eq. (14) and is an average of the Fourier basis over some bin centered at $`k`$, with width $`\mathrm{\Delta }k`$.
It can be shown that the variance of such a bin-averaged power estimate is given by (Appendix B)
$`V(k)[\widehat{P}_2(k)P(k)]^2={\displaystyle \frac{𝒩^3}{^2}}{\displaystyle \underset{\alpha }{}}[w^{\alpha \alpha }(k)]^2E^\alpha (k)`$ (26a)
$`E^\alpha (k){\displaystyle \frac{2}{n_{\overline{k}}}}\left[P(k)+{\displaystyle \frac{}{𝒩}}{\displaystyle \frac{q^\alpha \overline{N}_\mathrm{Q}^\alpha +V_\mathrm{B}^\alpha }{(\overline{N}_\mathrm{Q}^\alpha )^2}}\right]^2+{\displaystyle \frac{1}{}}T_{kk}`$ (26b)
$`+4B_{kk}{\displaystyle \frac{q^\alpha }{\overline{N}_\mathrm{Q}^\alpha }}{\displaystyle \frac{1}{𝒩}}+2P_{kk}\left[{\displaystyle \frac{q^\alpha }{\overline{N}_\mathrm{Q}^\alpha }}\right]^2{\displaystyle \frac{}{𝒩^2}}`$
$`+4P(k){\displaystyle \frac{q_{}^{}{}_{}{}^{\alpha }}{(\overline{N}_\mathrm{Q}^\alpha )^2}}{\displaystyle \frac{}{𝒩}}+{\displaystyle \frac{^2}{𝒩^3}}{\displaystyle \frac{q_{}^{\prime \prime }{}_{}{}^{\alpha }}{(\overline{N}_\mathrm{Q}^\alpha )^3}}`$
assuming the $`k`$ of interest satisfies $`k1/`$ and that the width of the bin $`\mathrm{\Delta }k`$ also satisfies $`\mathrm{\Delta }k1/`$, where $``$ is the length of the spectrum. This is sometimes referred to as the classical limit in the case of galaxy power spectrum measurement (Feldman et al. (1994); Hamilton (1997a)). We will not consider larger scales here, because measurements on such scales are likely dominated by systematic rather than random errors.
The symbol $`n_k`$ denotes the number of modes within the k bin of interest, $`𝒩`$ is the number of pixels in the length $``$; $`q^\alpha `$, $`\overline{N}_\mathrm{Q}^\alpha `$ and $`V_\mathrm{B}^\alpha `$ are as defined in eq. (14b). The quantities $`q^\alpha `$ and $`q^{\prime \prime \alpha }`$ are analogous to $`q^\alpha `$:
$$q^\alpha \underset{i}{}(W^{i\alpha })^3g_{\mathrm{ps}}^{i\alpha }g_\mathrm{b}^\alpha ,q^{\prime \prime \alpha }\underset{i}{}(W^{i\alpha })^4g_{\mathrm{ps}}^{i\alpha }g_\mathrm{b}^\alpha $$
(27)
The symbols $`T_{kk}`$, $`B_{kk}`$ and $`P_{kk}`$ represent the shell-averaged trispectrum, bispectrum and power spectrum respectively ($``$ here is to be distinguished from ensemble average discussed in §2):
$`T_{k_1k_2}`$ $`{\displaystyle \frac{1}{n_{k_1}n_{k_2}}}{\displaystyle \underset{k_1^{}}{}}{\displaystyle \underset{k_2^{}}{}}T(k_1^{},k_1^{},k_2^{},k_2^{})`$ (28a)
$`B_{k_1k_2}`$ $`{\displaystyle \frac{1}{n_{k_1}n_{k_2}}}{\displaystyle \underset{k_1^{}}{}}{\displaystyle \underset{k_2^{}}{}}B(k_1^{}k_2^{},k_1^{},k_2^{})`$ (28b)
$`P_{k_1k_2}`$ $`{\displaystyle \frac{1}{n_{k_1}n_{k_2}}}{\displaystyle \underset{k_1^{}}{}}{\displaystyle \underset{k_2^{}}{}}P(k_1^{}+k_2^{})`$ (28c)
where the sum over $`k_1^{}`$ extends over modes within the bin centered at $`k_1`$, and similarly for $`k_2^{}`$. The trispectrum $`T`$ and bispectrum $`B`$ are Fourier transforms of the four and three point correlation functions, defined in an analogous manner to eq. (2).
The variance as given in eq. (26) contains contributions from both cosmic fluctuations and discrete fluctuations (see §2). The terms such as $`P(k)^2`$ and $`T_{kk}`$ arise because of intrinsic fluctuations of the cosmic signal from one part of the universe to another – these terms are present even if one has data with arbitrarily high signal-to-noise. The terms containing $`\overline{N}_\mathrm{Q}^\alpha `$ arise because of discrete fluctuations – these we will loosely referred to as shot-noise.
As we have emphasized in §3.1 & 4.1, the shot noise contributions to the random error are not exactly Poisson-distributed. The shot noise contributions (ignoring cosmic sample fluctuations) in eq. (26b) would all be simply $`1/\overline{N}_\mathrm{Q}^\alpha `$ if $`\widehat{N}_\mathrm{Q}^\alpha `$ were strictly a Poisson variable. We have additional fluctuations in $`\widehat{N}_\mathrm{Q}`$ due to the background (sky and readout), and also due to non-unity weights used in reducing the data (eq. ; see also end of §3.1).
Given eq. (26), it is simple to derive a weighting $`\overline{w}^{\alpha \beta }(k)`$ that minimizes the variance $`V(k)`$, subject to the constraint that the effective window ($`G`$ as defined in eq. ) is properly normalized. This is most simply derived by minimizing the following Lagrangian:
$$L(k)=V(k)\lambda (G(k,k^{})\frac{dk^{}}{2\pi }1)$$
(29)
where $`\lambda `$ is a Lagrange multipler. Differentiating the above respect to $`\overline{w}^{\alpha \beta }(k)`$, and setting the result to zero, we obtain:
$$\overline{w}^{\alpha \beta }(k)=[E^\alpha (k)E^\beta (k)]^{\frac{1}{2}}/M(k),M(k)\underset{\mu }{}[E^\mu (k)]^1𝒩/$$
(30)
The corresponding shot-noise subtraction, instead of eq. (14b), would then be
$$b(k)=\underset{\alpha }{}\overline{w}^{\alpha \alpha }(k)\frac{q^\alpha \overline{N}_\mathrm{Q}^\alpha +V_\mathrm{B}^\alpha }{(\overline{N}_\mathrm{Q}^\alpha )^2}$$
(31)
where $`q^\alpha `$, $`N_Q^\alpha `$ and $`V_B^\alpha `$ are as defined in eq. (14b).
In summary, the minimum variance estimator of the power spectrum is
$$\widehat{P}_5(k)=\underset{\alpha ,\beta }{}R^{\alpha \beta }(k)[E^\alpha (k)^{\frac{1}{2}}\widehat{\delta }_f^\alpha ][E^\beta (k)^{\frac{1}{2}}\widehat{\delta }_f^\beta ]/M(k)b(k)$$
(32)
where $`b(k)`$ is given by eq. (31), $`E^\alpha (k)`$ and $`M(k)`$ are given in eq. (26) and (30), and $`R^{\alpha \beta }(k)`$ is as in eq. (14a).
The minimum variance estimator can be understood simply as follows: $`\widehat{\delta }_f^\alpha `$ at each pixel is weighed by $`1/\sqrt{E^\alpha (k)}`$ before the array is Fourier transformed, squared and grouped to form band power estimates. Note that the above estimator reduces to the one with uniform weighting (eq. ) if $`E^\alpha (k)`$ were independent of $`\alpha `$, e.g. when sample/cosmic variance is significantly larger than shot-noise ($`P(k)[/𝒩][q^\alpha \overline{N}_\mathrm{Q}^\alpha +V_\mathrm{B}^\alpha ]/[\overline{N}_\mathrm{Q}^\alpha ]^2`$). It is important to note that the weighting as a function of $`\alpha `$ is determined by $`\overline{N}_\mathrm{Q}^\alpha `$ rather than, say $`\widehat{N}_\mathrm{Q}^\alpha `$. Down-weighing pixels with a lot of absorption (hence relatively low $`\widehat{N}_\mathrm{Q}^\alpha `$) would be a wrong thing to do, since the fluctuations in absorption is the signal that we are after. The proper procedure is to down-weigh pixels with an overall lower mean count $`\overline{N}_\mathrm{Q}^\alpha `$.
Unfortunately, the minimum variance weighting given above is difficult to implement because one needs to specify simultaneously $`P`$, $`B`$ and $`T`$, in addition to the level of shot-noise. A common simplification is to use the Gaussian approximation in which $`E^\alpha (k)`$ is approximated as:
$$E^\alpha (k)\frac{2}{n_{\overline{k}}}\left[P(k)+\frac{}{𝒩}\frac{q^\alpha \overline{N}_\mathrm{Q}^\alpha +V_\mathrm{B}^\alpha }{(\overline{N}_\mathrm{Q}^\alpha )^2}\right]^2$$
(33)
(see e.g. Hamilton (1997a)). Note that in addition to ignoring the bispectrum and trispectrum terms, the above also ignores certain power spectrum terms which are mixed with shot-noise – the last three terms in eq. (26), which is equivalent to assuming that either the shot-noise or the correlation is sufficiently weak. With the above approximation, one can start with some initial $`P`$ and use the minimum variance weighting scheme to get a first measurement of $`P`$, and iterate subsequently (Bond et al. (1997)). Analogous (Gaussian) power spectrum estimators for galaxy-surveys and microwave background experiments have been widely discussed in the literature (e.g. Feldman et al. (1994); Vogeley & Szalay (1996); Tegmark et al. (1997); Hamilton (1997a); Tegmark et al. (1998); Bond et al. (1997); Seljak (1998)).
We will not attempt to address here the important question of how significant the non-Gaussian contributions are. A proper treatment will involve the analysis of a large number of simulations or a large data-set, which we hope to present in a future paper. It suffices to say that the very nonlinear mapping from the density field to $`e^\tau `$ will likely introduce some degree of non-Gaussianity, even if the initial density field is Gaussian.
The use of observed data to study this issue is particularly interesting, and deserves some more comments. In principle, since different QSO sightlines typically sample independent regions of the universe, one can estimate the variance of the transmission power spectrum, and hence infer the importance of the non-Gaussian contributions, using directly the fluctuations in power spectrum estimates from one sightline to another. However, one should keep in mind that shot-noise also contributes to the variance. Since different lines of sight generally have different S/N, the sightline-to-sightline fluctuations in power spectrum estimates should be interpreted with care. In a dataset of several quasars, it is possible that the quasar-to-quasar fluctuations are dominated by a few with low S/N, and their mean-square would give an overestimate of the true power spectrum variance.
We show in Fig. 13 an example in which the data consist of 6 high quality spectra (similar to Fig. 3) and 6 others with S/N about 20 times smaller. The bottom panel shows the power spectrum estimated with uniform weighting (eq. ) while the top panel represents the power spectrum estimated with minimum variance weighting using the Gaussian approximation. The (1 $`\sigma `$) error-bars are theoretical – they are estimated using eq. (26a) and (33). This illustrates how our weighting scheme can reduce the error bars at high $`k`$’s where shot-noise is important. Absent information on the non-Gaussian nature of the power spectrum variance, we advocate the Gaussian weighting scheme (eq. ) as a rational way to combine data with different levels of S/N to reduce the variance, even though it does not necessarily achieve minimum variance.
In combining the different spectra with different S/N, we have weighed the power spectrum estimate of each line of sight by its inverse variance, which is an obvious generalization of the minimum variance weighting introduced above. For instance, suppose we have two separate lines of sight $`A`$ and $`B`$, we could combine the two power spectrum estimate $`\widehat{P}_2^A`$ and $`\widehat{P}_2^B`$ in the following way, assuming the two lines of sight are independent:
$$\widehat{P}_2(k)=\left[\frac{\widehat{P}_2^A(k)}{V_A(k)}+\frac{\widehat{P}_2^B(k)}{V_B(k)}\right]/\left[\frac{1}{V_A(k)}+\frac{1}{V_B(k)}\right]$$
(34)
where $`V_A(k)`$ and $`V_B(k)`$ are estimated with the same $`P`$ but could have different levels of shot-noise. The noisier quasar spectrum is naturally down-weighed.
Lastly, we should emphasize the above discussion does not address the issue of cross-variance between power spectrum estimates at two different wave bands, which is introduced by the non-Gaussian terms (Meiksin & White (1999); Scoccimarro et al. (1999)). Hamilton Hamilton (1999) introduced a scheme which simultaneously diagonalizes the covariance and minimizes it. However, it makes specific assumptions about the form of the trispectrum and bispectrum, the validity of which for the forest remains to be checked.
## 5. Discussion
Our recipe for measuring the transmission power spectrum is summarized here.
* Given an array of reduced quasar counts $`\widehat{N}_\mathrm{Q}^\alpha `$, identified metal lines should be removed, especially the strong ones ($`\tau >1`$). Small gaps in the spectrum (e.g. due to cosmic-ray-hits removal) can be (linearly) interpolated over, while large gaps should be avoided (§4.2.3).
* The mean quasar counts ($`\overline{N}_\mathrm{Q}^\alpha `$) is estimated using eq. (18). The mean-basis (the functional form of the mean-trend) should be chosen to be as smooth as possible – underfitting is better than overfitting (see Test 1 and 2 of §4.2.1). In practice it appears that a flat mean suffices for short spectra ($`50\AA `$), while polynomials up to the third order can be used for longer spectra ($`500\AA `$). One can gain an idea of what a reasonable basis is using the red side of Ly$`\alpha `$ or low redshift QSO spectra where the continuum can be seen more clearly.
* Define the trend-removed and normalized fluctuation $`\widehat{\delta }_f^\alpha `$ according to eq. (12b), and the power spectrum is estimated using the quadratic estimator in eq. (12a). Different weightings are possible, the simplest that we recommend is given in eq. (14). A more sophisticated weighting scheme which can reduce the random error is given by eq. (30), (31), (32) and (33). If one is interested in the real-space correlation function instead, the recommended weighting is eq. (13) – this gives a smaller variance compared to other estimators commonly used in the literature. We emphasize that the shot-noise-bias ($`b(k)`$ in eq. \[14b\], or in eq. ) should be subtracted correctly, especially for noisy data.
* If better control over systematic errors on large scales ($`{}_{}{}^{}{}_{}{}^{>}\mathrm{\hspace{0.17em}\hspace{0.17em}30}\AA `$) introduced by the unknown continuum is desired, the techniques outlined in §4.2.2 can be used. The corresponding estimator is given in eq. (23), which requires an estimate of the continuum power spectrum $`P_\mathrm{C}`$ (eq. \[21b\]) and an additional correction factor $`ϵ_\mathrm{G}`$ (eq. ). This procedure requires the identification of an appropriate set of continua (see discussions in §4.2.2). Even if one is not interested in the power spectrum on large scales, we recommend this procedure as a consistency check that the spurious power introduced by the continuum is negligible on the scales of interest.
What implications does the above have for one’s observing strategy? To discuss this question, we need to take a closer look at the issue of shot-noise. The shot-noise enters in two different places in the above discussion. First, it contributes to the variance (random error) of the power spectrum estimate (eq. ). Second, it appears as a bias in the power spectrum estimate that we have to subtract off (e.g. $`b(k)`$ in eq. \[12a\] & ).
In the literature on power spectrum measurement, shot-noise subtraction has been largely ignored (e.g. Croft et al. (1998a); see McDonald et al. (1999) for an alternative approach where shot-noise is simulated rather than subtracted). Let us estimate how important it is. The expression in the simplest case of uniform weighting is given in eq. (14b) (see eq. for more complicated weightings), which can be rewritten as
$$b(k)=\frac{\mathrm{\Delta }u}{𝒩}\underset{\alpha }{}\frac{\underset{i}{}(W^{i\alpha })^2[g_{\mathrm{ps}}^{i\alpha }g_\mathrm{b}^\alpha \overline{N}_Q^\alpha +V_B^{i\alpha }]}{(\overline{N}_Q^\alpha )^2}$$
(35)
where $`\mathrm{\Delta }u`$ is the size of a pixel, $`𝒩`$ is the total number of pixels, and the rest of the symbols are as defined in §3.1: $`i`$ is the pixel-label in the spatial direction and $`\alpha `$ in the spectral direction, $`\overline{N}_\mathrm{Q}^\alpha `$ is the mean reduced quasar count, $`V_\mathrm{B}^{i\alpha }`$ is the background variance, $`W^{i\alpha }`$ is a weighting, and $`g_{\mathrm{ps}}^{i\alpha }`$ and $`g_\mathrm{b}^\alpha `$ are the point-spread function and blaze function respectively (eq. & ). An important observation is that the numerator within the summation is closely related to the variance array which is often given along with a spectrum (eq. )<sup>7</sup><sup>7</sup>7See paragraph after eq. (9) and §2 on the distinction between $`\stackrel{~}{N}_\mathrm{Q}^{i\alpha }`$ and $`\widehat{N}_\mathrm{Q}^{i\alpha }`$.:
$$\mathrm{var}(\alpha )=\underset{i}{}(W^{i\alpha })^2[\widehat{N}_\mathrm{Q}^{i\alpha }+V_B^{i\alpha }]$$
(36)
The quantity $`\widehat{N}_\mathrm{Q}^{i\alpha }`$ is of course different from $`g_{\mathrm{ps}}^{i\alpha }g_\mathrm{b}^\alpha \overline{N}_Q^\alpha `$ which we need to estimate the shot-noise, but since we are in practice interested in an average over all pixels, it turns out the following estimate of the shot-noise is accurate to within a percent for all cases we have tested:
$$b(k)\frac{\mathrm{\Delta }u}{𝒩}\underset{\alpha }{}\frac{\mathrm{var}(\alpha )}{(\overline{N}_Q^\alpha )^2}$$
(37)
Without the above approximation, an exact estimate of the shot-noise would require the knowledge of $`\widehat{N}_Q^\alpha `$, $`_i(W^{i\alpha })^2g_{\mathrm{ps}}^{i\alpha }g_\mathrm{b}^\alpha `$ and $`_i(W^{i\alpha })^2V_B^{i\alpha }`$.
Eq. (37) provides a useful means of estimating the shot-noise (see Appendix A on shot-noise estimation under more complicated circumstances i.e. with non-trivial rebinning or weighting). One can simplify further by making a crude approximation in relating $`b(k)`$ to the typical signal-to-noise ratio of the data (which is often quoted at the continuum) through
$$b(k)(\mathrm{\Delta }u/\overline{f})(N/S)^2$$
(38)
which can be justified if one ignores the part of the variance due to the sky and readout. We find that this simple rule of thumb generally provides an underestimate of the shot-noise (particularly at low $`S/N`$ where the background counts become important), but is accurate to within about a factor of $`2`$.
Fig. 14 summarizes some useful information for devising an observing strategy, based on our estimate of the shot-noise in eq. (38) above. The solid line shows the mean observed transmission power spectrum at $`z=3`$ taken from McDonald et al. McDonald et al. (1999).<sup>8</sup><sup>8</sup>8We divide the un-normalized power spectrum (eq. ) of McDonald et al. by the square of their measured mean transmission to obtain the normalized power spectrum given in Fig. 14. See tests 4 and 5 in §4.2.1 on the bias of the un-normalized power spectrum. The two horizontal dotted lines show the level of shot-noise expected for the 2 extremes of the kinds of observations we are likely to encounter – the bottom corresponds to very high signal to noise (S/N) observations with HIRES quality resolution (e.g. Hu et al. (1995); Kirkman & Tytler (1997); Rauch et al. (1997)), while the dotted line on the top corresponds to low S/N observations expected for a large number of quasars in the Sloan Digital Sky Survey (SDSS). We emphasize that the shot-noise level does not depend on the resolution per se, but on the pixel size for a given S/N. SDSS is expected to produce $`1000`$ QSO spectra at S/N $`=20`$ per pixel (QSO’s at $`z>2.7`$, where the redshift limit is determined by the blue limit of the spectrograph, $`3800\AA `$), $`10000`$ at S/N $`=15`$ and $`30000`$ at S/N $`=7`$, corresponding to $`i^{}`$ magnitude-cuts at $`18`$, $`19`$ and $`20`$ respectively (Fukugita et al. (1996); Fan (1999)). The pixel size of SDSS is quite uniform in velocity $`70\mathrm{km}/\mathrm{s}`$, which is equivalent to $`1.13\AA `$ at $`4864\AA `$ (Ly$`\alpha `$ at $`z=3`$).
Clearly, the importance of shot-noise depends on the scales at which one is interested in measuring the power spectrum. A few interesting scales are shown in Fig. 14. First, instrumental resolution imposes a high $`k`$ limit beyond which one cannot reliably measure the transmission power spectrum. The resolution window is often characterized by a Gaussian with a given FWHM. The effect of such a resolution window on the power spectrum can be represented by $`P(k)P(k)e^{k^2/k_\sigma ^2}`$ where $`k_\sigma =\sqrt{2/\mathrm{ln2}}/\mathrm{FWHM}1.7/\mathrm{FWHM}`$. Two representative $`k_\sigma `$’s are shown as long tickmarks at the top. Note that even at $`k=k_\sigma `$, the resolution window reduces the power by $`63\%`$ and so has a non-negligible effect. The Sloan FWHM is about $`2.1`$ pixels i.e. $`147\mathrm{km}/\mathrm{s}`$, or $`2.4\AA `$ at $`4864\AA `$.
On the other hand, the range of scales that is currently being used to infer the mass power spectrum is indicated by the interval near the bottom $`\mathrm{\Delta }k_{\mathrm{use}}`$. The high $`k`$ limit is set by the scales at which the shape of the power spectrum is preserved in the transformation from mass to transmission (i.e. linear biasing e.g. Croft et al. (1998a)). We can see that for high quality Keck spectra, information from a whole decade of measurable scales is unused for the recovery of the mass power spectrum – it would be very useful to push the current analysis techniques to these scales, since power on these scales is of particular interest in constraining e.g. neutrino properties (Hui et al. (1997); Croft et al. (1999)). Such an effort would require disentangling the effects of peculiar velocities and thermal broadening, however. At the other end, the low $`k`$ limit of currently usable scales is set by the scales at which the continuum fluctuates. This is indicated by the dashed line at the top, where the transmission power spectrum is unknown.
From the above discussion, we can distill a few tips for observing/analysis.
* To ensure that shot-noise is subdominant, one might want to achieve $`S/N{}_{}{}^{}{}_{}{}^{>}`$ $`\sqrt{10\mathrm{\Delta }u/\overline{f}/P(k_{\mathrm{int}.})}`$ where $`k_{\mathrm{int}.}`$ is the scale of interest, and $`\overline{f}`$ is the mean transmission. The factor of $`10`$ is somewhat arbitrary – this will ensure the shot-noise contribution to the power spectrum variance is no more than about $`20\%`$ (under the Gaussian approximation; see eq. ), or the $`1\sigma `$ error-bar on the power spectrum would only be increased by $`10\%`$ due to the contribution from shot-noise. An important question is what $`k_{\mathrm{int}.}`$ should be – that depends on at what scales one can usefully extract cosmologically interesting information. Current literature mainly focused on $`k_{\mathrm{int}.}{}_{}{}^{}{}_{}{}^{<}\mathrm{\hspace{0.17em}\hspace{0.17em}2}\AA ^1`$, where $`P0.06\AA `$, therefore $`S/N20\sqrt{\mathrm{\Delta }u/1\AA }`$ would be sufficient. Since $`P`$ rises with scale, shot-noise would be even less important at smaller $`k`$’s. Note that with very small $`\mathrm{\Delta }u`$ such as $`0.05\AA `$, $`S/N`$ as low as $`45`$ is acceptable. To give some examples, a $`S/N`$ of 8 per $`0.05\AA `$ can be achieved with an hour of exposure using Keck/HIRES for a $`V=19`$ quasar; on the other hand, a $`S/N`$ of about 15 per $`1.1\AA `$ is expected with just slightly under an hour of exposure using the SDSS spectrograph for a $`i^{}=19`$ quasar.
* A corollary of focusing on only $`k{}_{}{}^{}{}_{}{}^{<}\mathrm{\hspace{0.17em}\hspace{0.17em}2}\AA ^1`$ is that observations with $`k_\sigma {}_{}{}^{}{}_{}{}^{>}\mathrm{\hspace{0.17em}\hspace{0.17em}3}\times 2\AA ^1`$ or a resolution FWHM of $`0.3\AA `$ or $`R16000`$ at $`z=3`$ are acceptable. The factor of $`3`$ above (i.e. in $`3\times 2\AA ^1`$) is somewhat arbitrary – it ensures that at $`k=2\AA ^1`$, the resolution window does not reduce the power by more than $`10\%`$. If the resolution window is known accurately, or if one is willing to sacrifice information on the small scales close to $`k2\AA ^1`$, one could in principle consider lower resolutions. We would like to emphasize, however, that in principle, the modes at $`k>2\AA ^1`$ could still contain very interesting cosmological information, even though the current attempts at recovering the mass power spectrum ignored them.
* If shot-noise is subdominant compared to the power spectrum, the only other limiting factor to the size of the random error is the total size of one’s sample or the number of sightlines in it. Assuming all sightlines have similar coverage with length $``$, then the fractional error of a single k-mode (i.e. in a k-bin of $`2\pi /`$) is given by $`\delta P/P=C/\sqrt{N_{\mathrm{sight}}}`$ where $`N_{\mathrm{sight}}`$ is the number of sightlines assuming they are independent, and $`C=1`$ under the Gaussian approximation (eq. ), and a little larger than unity under more general circumstance (see e.g. Meiksin & White (1999); Scoccimarro et al. (1999)).
* How should one distribute one’s observing time among quasar targets to minimize the random error on the transmission power spectrum? There are many possible versions of this problem. We will discuss two, giving an explicit solution for the first, and only general expressions for the second. In the simplest case in which all the candidate quasar targets have similar magnitudes, given a finite amount of observing time, one can deduce the optimal total number of quasars one should target by
$$\mathrm{minimizing}N_{\mathrm{tot}.}^{}{}_{}{}^{1}[P(k_{\mathrm{int}.})+\frac{\mathrm{\Delta }u}{\overline{f}}\frac{A}{t}]^2\mathrm{subject}\mathrm{to}N_{\mathrm{tot}.}t=T_{\mathrm{tot}.}$$
(39)
where $`N_{\mathrm{tot}.}`$ is the total number of quasars targeted, $`k_{\mathrm{int}.}`$ is the scale of interest, $`T_{\mathrm{tot}.}`$ is the total amount of observing time one has, $`t`$ is the amount of time one spends on each quasar, and $`1/A`$ is equal to $`(S/N)^2`$ reached per unit exposure time. The above assumes eq. (33) and that the sightlines are independent. The solution is easy to deduce: $`N_{\mathrm{tot}.}=[P(k_{\mathrm{int}.})\overline{f}/\mathrm{\Delta }u][T_{\mathrm{tot}.}/A]`$, or $`t=A\mathrm{\Delta }u/\overline{f}/P(k_{\mathrm{int}.})`$. A typical value for $`1/A`$ is $`1/A1200\mathrm{hour}\mathrm{per}\mathrm{\AA }\times 10^{(19mag.)/2.5}\times [\mathrm{aperture}/100\mathrm{m}.^2]\times f_{\mathrm{throughput}}`$, where $`f_{\mathrm{throughput}}`$ is about unity for Keck/HIRES, and $`2.5`$ for the SDSS. Using again $`k_{\mathrm{int}.}2\AA ^1`$, for a $`19`$th magnitude quasar, with an aperture of $`6.25\mathrm{m}.^2`$ and assuming $`\mathrm{\Delta }u=1\AA `$ and $`f_{\mathrm{throughput}}=3`$, the exposure time is $`t=14`$ minutes only! The above prescription, however, only allows for just enough exposure time to reduce the shot-noise to a level comparable to the cosmic/sampling variance (i.e. $`P(k_{\mathrm{int}.})(\mathrm{\Delta }u/\overline{f})(N/S)^2`$) – the sole aim is to maximize the number of quasars observed within a given length of time to beat down the sampling variance. The prescription would certainly be different if one has, for instance, a finite number of quasar targets, or if one has other purposes in mind – such as measuring the mean decrement, etc (see earlier prescription for making shot-noise subdominant, equivalent to multiplying $`t`$ by about a factor of $`10`$). A more general version of the above problem deals with a case where the quasars span a range of magnitudes i.e. $`A`$ is no longer the same number for each quasar. A simple ansatz is to assume $`t=\alpha A\mathrm{\Delta }u/\overline{f}/P(k_{\mathrm{int}.})`$, in other words, spending more time for fainter quasars because it takes longer to beat down the shot-noise, except that we have a normalizing factor $`\alpha `$ which enforces the constraint of total observing time: $`\alpha =T_{\mathrm{tot}.}[\overline{f}P(k_{\mathrm{int}.})/\mathrm{\Delta }u]/_{A_{\mathrm{min}}}^AA^{}n(A^{})𝑑A^{}`$ where $`n(A)dA`$ is the number of quasars with $`A`$ falling in the given range, and $`A_{\mathrm{min}}`$ corresponds to the brightest quasar in one’s sample. Then, we can determine how many quasars one should include, starting from the brightest one, or how faint one should go by minimizing $`[_{A_{\mathrm{min}}}^An(A^{})𝑑A^{}]^1P(k_{\mathrm{int}.})^2(1+1/\alpha )^2`$ with respect to $`A`$.
The following is particularly relevant for SDSS or comparable observations.
* In addition to contributing to the power spectrum variance, shot-noise also contributes a bias which has to be subtracted off (see e.g. Fig. 7). This is quite important for SDSS because, with $`>10^4`$ sightlines, the survey has the capability of reducing the fractional error of the power spectrum to $`<1\%`$ per mode. Therefore, a bias of $`3100\%`$, depending on the scale of interest (as indicated by the top dotted line in Fig. 14), is not acceptable and should be subtracted off. We note that analyses so far in the literature (e.g. Croft et al. (1998b); McDonald et al. (1999)) focused on higher quality data where $`S/N30`$, with $`\mathrm{\Delta }u`$ ranging from about $`0.04\AA `$ to $`1\AA `$, and so according to eq. (38) and Fig. 14, the shot-noise bias was about $`1\%`$ of the power or smaller and therefore could be ignored, although a more careful check should be performed for some datasets with lower $`S/N`$.
* The low resolution of SDSS spectra implies that it would be difficult to obtain useful information on scales $`kk_\sigma 0.01(\mathrm{km}/\mathrm{s})^1`$ or $`0.7\AA ^1`$. On larger scales or smaller $`k`$’s, two problems have to be reckoned with. For $`k0.3k_\sigma k_\sigma `$, the resolution window suppresses the power by $`10\%`$ or more – therefore, one needs to have an accurate measure of the resolution window to recover the true transmission power spectrum. <sup>9</sup><sup>9</sup>9We thank Rupert Croft for useful discussions on this point. This can be achieved by using narrow metal lines or arc lines. There are relatively few sky lines in the relevant part of the spectrum.
* For scales $`k<0.004(\mathrm{km}/\mathrm{s})^1`$ or $`0.2\AA ^1`$, the effect of the continuum has to be properly taken into account, and the method of §4.2.2 can be useful here. From Fig. 14, it is clear that the range of scales accessible to SDSS would be quite limited unless a correction for continuum contamination is applied.
It is worth pointing out that while the above quoted numbers are all based on $`z=3`$, we do not expect them to change significantly for $`z=2`$ or $`z=4`$. This is in part because of the slow evolution of the transmission power spectrum – the growth of the mass spectrum with time is partially compensated by the lowering of the mean decrement (McDonald et al. (1999)).
At least three issues remain to be explored in future work. As we have emphasized, the concept of correcting for continuum contamination in the transmission power on large scales as laid out in §4.2.2 is an interesting one, but requires more testing. An important check is the measurement of continuum power spectrum as a function of redshift on both sides of Ly$`\alpha `$ emission. Second, counts-in-cells analysis (i.e. measuring moments of the one-point probability distribution function), just like power spectrum analysis, requires shot-noise subtraction, and since typically one considers cells at the limit of resolution, shot-noise is likely non-negligible except for high S/N spectra. Counts-in-cells analysis provides a very interesting way to test the gravitational instability paradigm (Gaztañaga & Croft (1998); Hui (1998); Nusser & Haehnelt (1999)), and should be done with care. Useful expressions will be presented elsewhere (Hui (2000)). Lastly, as is hopefully clear by now, the spirit of the methods presented in this paper is to avoid continuum-fitting and replace it with trend-removal. We have demonstrated that this is possible for measuring the transmission power spectrum. However, other quantities of interest such as the mean decrement requires an estimate of the continuum, by definition. Furthermore, to theoretically interpret the transmission power spectrum in terms of the mass fluctuation, current methods require the measurement of the mean decrement to fix a free parameter in one’s cosmological model, which is a combination of the ionizing background, the mean baryon density and the mean temperature. Therefore, in a sense, the technique of trend-removal only goes half-way in solving the problem of continuum-fitting. Although we still recommend our method over continuum-fitting because the transmission power spectrum is an unambiguous quantity which can and should be determined as accurately as possible (not to mention the fact that continuum-fitting is difficult with low S/N or low resolution, or at high $`z`$), there is clearly a need for an alternative method to bridge the gap between the measured transmission power and the theoretically interesting mass power. This will be explored in future publications (Hui & Burles (2000); Zaldarriaga et al. (2000)).
We thank for useful discussions Len Cowie, David Kirkman, Patrick Petitjean, Michael Rauch, Wal Sargent, Don Schneider, David Weinberg, Don York and the participants of the 1999 Haifa workshop on the intergalactic medium and large scale structure. We also thank Nick Gnedin for supplying a simulation. Special thanks are due to Matias Zaldarriaga for pointing out the importance of aliasing and for many interesting discussions. This work was supported in part by the DOE and the NASA grant NAG 5-7092 at Fermilab, and by the NSF grant PHY-9513835. LH thanks the IAS for the Taplin Fellowship.
## Appendix A
Our main aim here is to derive eq. (15) for the estimator $`\widehat{P}_2(k)`$ which is given by eq. (12) and eq. (14), with an eye towards generalization to $`W^{i\beta }`$ different from eq. (6) and $`W^{\alpha \beta }\delta ^{\alpha \beta }`$. Derivations of results in §4.3 on estimator variance and the minimum variance power spectrum estimator can be found in Appendix B. We will ignore the integral constraint and assume $`\overline{N}^\alpha `$ is known to high accuracy.
We need first of all the correlation matrix $`\widehat{\delta }_f^\alpha \widehat{\delta }_f^\beta `$. We will do it in 2 steps, first let us derive $`\widehat{\delta }_f^\alpha \widehat{\delta }_f^\beta _D`$. Rewriting $`\widehat{\delta }_f`$ (eq. \[12b\]) as $`(\widehat{N}^\alpha \stackrel{~}{N}^\alpha )/\overline{N}^\alpha +(\stackrel{~}{N}^\alpha \overline{N}^\alpha )/\overline{N}^\alpha `$ where $`\stackrel{~}{N}^\alpha =\widehat{N}^\alpha _D`$, it can be shown that
$`\widehat{\delta }_f^\alpha \widehat{\delta }_f^\beta _D`$ $`=`$ $`\delta _f^\alpha \delta _f^\beta +{\displaystyle \frac{1}{\overline{N}_\mathrm{Q}^\alpha \overline{N}_\mathrm{Q}^\beta }}(\widehat{N}^\alpha \stackrel{~}{N}^\alpha )(\widehat{N}^\beta \stackrel{~}{N}^\beta )_D`$
$`=`$ $`\delta _f^\alpha \delta _f^\beta +{\displaystyle \frac{1}{\overline{N}_\mathrm{Q}^\alpha \overline{N}_\mathrm{Q}^\beta }}{\displaystyle \underset{\gamma i}{}}W^{\alpha \gamma }W^{\beta \gamma }(W^{i\gamma })^2(\widehat{N}_\mathrm{Q}^{i\gamma }_D+V_\mathrm{B}^{i\gamma })`$
where $`\delta _f^\alpha `$ is to be distinguished from $`\widehat{\delta }_f^\alpha `$ in that it has only cosmic or sample fluctuations (eq. ), $`\widehat{N}_\mathrm{Q}^{i\gamma }`$ is a strictly Poisson variable with an average given by eq. (5), and $`V_\mathrm{B}^{i\gamma }`$ is the variance contributions from the sky and readout (eq. 8).
Taking the cosmic mean of the above, we obtain the correlation matrix
$$\widehat{\delta }_f^\alpha \widehat{\delta }_f^\beta =\xi (u^\alpha u^\beta )+\frac{1}{\overline{N}_\mathrm{Q}^\alpha \overline{N}_\mathrm{Q}^\beta }\underset{\gamma }{}W^{\alpha \gamma }W^{\beta \gamma }\underset{i}{}[(W^{i\gamma })^2g_{\mathrm{ps}}^{i\gamma }g_\mathrm{b}^\gamma \stackrel{~}{N}_Q^\gamma +(W^{i\gamma })^2V_\mathrm{B}^{i\gamma }]$$
(41)
The second term on the right is the shot-noise contribution which has to be subtracted off. Using the estimator in eq. (12), with $`w^{\alpha \beta }(k)`$ given in eq. (14), the correct shot-noise subtraction is:
$$b(k)=(/𝒩^2)\underset{\alpha \beta \gamma }{}W^{\alpha \gamma }W^{\beta \gamma }\underset{i}{}[(W^{i\gamma })^2g_{\mathrm{ps}}^{i\gamma }g_\mathrm{b}^\gamma \stackrel{~}{N}_\mathrm{Q}^\gamma +(W^{i\gamma })^2V_\mathrm{B}^{i\gamma }]/(\overline{N}_\mathrm{Q}^\alpha \overline{N}_\mathrm{Q}^\beta )$$
(42)
where we have made use of the assumption that $`W^{\alpha \gamma }`$ is non-zero only for $`\alpha `$ and $`\gamma `$ on very small separations, and that the $`k`$ of interest satisfies $`k(u^\alpha u^\gamma )1`$ on such separations.
Note that for weightings such as the one given in eq. (7), $`W^{i\gamma }`$ depends on $`\stackrel{~}{N}_\mathrm{Q}^\gamma `$ which makes an estimation of the shot-noise non-trivial. However, simplification results in two extreme limits: in the signal dominated regime where $`\stackrel{~}{N}_\mathrm{Q}^{j\beta }V_\mathrm{B}^{j\beta }`$, $`W^{i\gamma }`$ reduces to uniform weighting as in eq. (6); in the background dominated regime where $`\stackrel{~}{N}_\mathrm{Q}^{j\gamma }V_\mathrm{B}^{j\gamma }`$, $`W^{i\gamma }`$ reduces to $`(1/g_\mathrm{b}^\gamma )(g_{\mathrm{ps}}^{i\gamma }/V_\mathrm{B}^{i\gamma })/(_j(g_{\mathrm{ps}}^{i\gamma })^2/V_\mathrm{B}^{i\gamma })`$ which is independent of $`\stackrel{~}{N}_\mathrm{Q}^\gamma `$. In such cases, and for $`W^{\alpha \gamma }=\delta ^{\alpha \gamma }`$, the above $`b(k)`$ reduces to the one given in eq. (14b).
For $`W^{\alpha \gamma }\delta ^{\alpha \gamma }`$ (i.e. rebinning has bee done; we continue to assume $`W^{i\gamma }`$ is roughly independent of the signal), we can write $`b(k)`$ as
$$b(k)(/𝒩^2)\underset{\alpha \beta \gamma }{}\frac{W^{\alpha \gamma }W^{\beta \gamma }}{\overline{N}_\mathrm{Q}^\alpha \overline{N}_\mathrm{Q}^\beta }\underset{i}{}(W^{i\gamma })^2[g_{\mathrm{ps}}^{i\gamma }g_\mathrm{b}^\gamma \stackrel{~}{N}_\mathrm{Q}^\gamma +V_\mathrm{B}^{i\gamma }]$$
(43)
where the term under summation of $`i`$ is simply the variance array of the pre-rebinned data (eq. ), and one can replace $`\overline{N}_\mathrm{Q}^\alpha \overline{N}_\mathrm{Q}^\beta `$ by $`(\overline{N}_\mathrm{Q}^\alpha )^2`$, since $`W^{\alpha \gamma }W^{\beta \gamma }`$ is non-zero only if $`\alpha `$ and $`\beta `$ are close together. Eq. (37) is therefore replaced by the following if the data have been rebinned:
$$b(k)\frac{\mathrm{\Delta }u}{𝒩}\underset{\alpha \beta }{}\frac{\mathrm{var}(\alpha ,\beta )}{\overline{N}_Q^\alpha \overline{N}_Q^\beta }$$
(44)
where $`\mathrm{var}(\alpha ,\beta )`$ is the variance matrix of the rebinned data.
To complete our derivation, we need to show that the choice of $`w^{\alpha \beta }(k)`$ given in eq. (14) has the correct normalization such that the window function satisfies $`𝑑k^{}G(k,k^{})/2\pi =1`$ (eq. ). Putting eq. (14) into eq. (12), and using the correlation matrix given in eq. (41) together with the relation between the two point function and the power spectrum in eq. (2), it is not hard to see that $`\widehat{P}_2(k)`$ satisfies eq. (15) with $`G(k,k^{})`$ given by $`_{\alpha \beta }w^{\alpha \beta }(k)e^{ik^{}(u_\alpha u_\beta )}`$. Using $`𝑑k^{}e^{ik^{}(u_\alpha u_\beta )}=2\pi \delta ^{\alpha \beta }𝒩/`$ then completes the derivation. One might want to explore more complicated data windowing (e.g. Press et al. (1992b); Hamilton (1997b)), but since in practice uncertainties in the large scale power estimate, where the survey window matters most, are likely dominated by the continuum, the simple choice we have adopted is probably adequate.
## Appendix B
We derive here the band-power variance given in eq. (26). The power spectrum estimator is given in eq. (12) with the matrix $`w^{\alpha \beta }(k)`$ given by eq. (25). We ignore here the uncertainty in the estimation of the mean count $`\overline{N}_\mathrm{Q}^\alpha `$.
The band-power covariance can be written compactly as
$$C(k_1,k_2)\mathrm{\Delta }\widehat{P}^\alpha (k_1)\mathrm{\Delta }\widehat{P}^\beta (k_2)=\underset{\alpha \beta \gamma \eta }{}w^{\alpha \beta }(k_1)w^{\gamma \eta }(k_2)(\widehat{\delta }_f^\alpha \widehat{\delta }_f^\beta \widehat{\delta }_f^\gamma \widehat{\delta }_f^\eta \widehat{\delta }_f^\alpha \widehat{\delta }_f^\beta \widehat{\delta }_f^\gamma \widehat{\delta }_f^\eta )$$
(45)
The band-power variance is simply the diagonal piece of this matrix: $`C(k,k)`$.
We can work out $`\widehat{\delta }_f^\alpha \widehat{\delta }_f^\beta \widehat{\delta }_f^\gamma \widehat{\delta }_f^\eta \widehat{\delta }_f^\alpha \widehat{\delta }_f^\beta \widehat{\delta }_f^\gamma \widehat{\delta }_f^\eta `$ using the same methodology as used in Appendix A for $`\widehat{\delta }_f^\alpha \widehat{\delta }_f^\beta `$: rewrite $`\widehat{\delta }_f^\alpha `$ as $`(\widehat{N}^\alpha \stackrel{~}{N}^\alpha )/\overline{N}^\alpha +(\stackrel{~}{N}^\alpha \overline{N}^\alpha )/\overline{N}^\alpha `$ where $`\stackrel{~}{N}^\alpha =\widehat{N}^\alpha _D`$, and as before, take the discrete-ensemble average $`_D`$ before taking the cosmic average $``$. The result is
$`C(k_1,k_2)={\displaystyle \underset{\alpha \beta \gamma \eta }{}}w^{\eta \gamma }(k_1)w^{\alpha \beta }(k_2)(`$ (46)
$`\delta _f^\alpha \delta _f^\beta \delta _f^\gamma \delta _f^\eta +\delta _f^\alpha \delta _f^\eta \delta _f^\beta \delta _f^\gamma +\delta _f^\alpha \delta _f^\gamma \delta _f^\beta \delta _f^\eta +\delta _f^\alpha \delta _f^\beta \delta _f^\gamma \delta _f^\eta _c`$
$`+{\displaystyle \frac{1}{\overline{N}_\mathrm{Q}^\eta \overline{N}_\mathrm{Q}^\gamma }}{\displaystyle \underset{\sigma i}{}}W^{\eta \sigma }W^{\gamma \sigma }(W^{i\sigma })^2[g_{\mathrm{ps}}^{i\sigma }g_\mathrm{b}^\sigma \overline{N}_\mathrm{Q}^\sigma (1+\delta _f^\sigma )+V^{i\sigma }]\delta _f^\alpha \delta _f^\beta `$
$`+1\mathrm{other}\mathrm{perm}.:(\alpha \gamma ,\beta \eta )`$
$`+{\displaystyle \frac{1}{\overline{N}_\mathrm{Q}^\beta \overline{N}_\mathrm{Q}^\gamma }}{\displaystyle \underset{\sigma i}{}}W^{\beta \sigma }W^{\gamma \sigma }(W^{i\sigma })^2(g_{\mathrm{ps}}^{i\sigma }g_\mathrm{b}^\sigma \overline{N}_\mathrm{Q}^\sigma (1+\delta _f^\sigma )+V_\mathrm{B}^{i\sigma })\delta _f^\alpha \delta _f^\eta `$
$`+3\mathrm{other}\mathrm{perm}.:(\beta \alpha ,\gamma \eta ),(\gamma \eta ),(\beta \alpha )`$
$`+{\displaystyle \frac{1}{\overline{N}_\mathrm{Q}^\beta \overline{N}_\mathrm{Q}^\gamma \overline{N}_\mathrm{Q}^\eta }}{\displaystyle \underset{\sigma i}{}}W^{\beta \sigma }W^{\eta \sigma }W^{\gamma \sigma }(W^{i\sigma })^3[g_{\mathrm{ps}}^{i\sigma }g_\mathrm{b}^\sigma \overline{N}_\mathrm{Q}^\sigma (1+\delta _f^\sigma )+\stackrel{~}{N}_\mathrm{S}^{i\sigma }]\delta _f^\alpha `$
$`+3\mathrm{other}\mathrm{perm}.:(\alpha \beta ),(\alpha \gamma ),(\alpha \eta )`$
$`+{\displaystyle \frac{1}{\overline{N}_\mathrm{Q}^\alpha \overline{N}_\mathrm{Q}^\beta \overline{N}_\mathrm{Q}^\gamma \overline{N}_\mathrm{Q}^\eta }}{\displaystyle \underset{\sigma i\chi j}{}}W^{\alpha \sigma }W^{\beta \sigma }W^{\eta \chi }W^{\gamma \chi }(W^{i\sigma })^2(W^{j\chi })^2(g_{\mathrm{ps}}^{i\sigma }g_\mathrm{b}^\sigma \overline{N}_\mathrm{Q}^\sigma (1+\delta _f^\sigma )+V_\mathrm{B}^{i\sigma })`$
$`(g_{\mathrm{ps}}^{j\chi }g_\mathrm{b}^\chi \overline{N}_\mathrm{Q}^\chi (1+\delta _f^\chi )+V_\mathrm{B}^{j\chi })`$
$`+{\displaystyle \frac{1}{\overline{N}_\mathrm{Q}^\alpha \overline{N}_\mathrm{Q}^\beta \overline{N}_\mathrm{Q}^\gamma \overline{N}_\mathrm{Q}^\eta }}{\displaystyle \underset{\sigma i\chi j}{}}W^{\alpha \sigma }W^{\eta \sigma }W^{\beta \chi }W^{\gamma \chi }(W^{i\sigma })^2(W^{j\chi })^2(g_{\mathrm{ps}}^{i\sigma }g_\mathrm{b}^\sigma \overline{N}_\mathrm{Q}^\sigma (1+\delta _f^\sigma )+V_\mathrm{B}^{i\sigma })`$
$`(g_{\mathrm{ps}}^{j\chi }g_\mathrm{b}^\chi \overline{N}_\mathrm{Q}^\chi (1+\delta _f^\chi )+V_\mathrm{B}^{j\chi })+1\mathrm{other}\mathrm{perm}.:(\alpha \beta )`$
$`+{\displaystyle \frac{1}{\overline{N}_\mathrm{Q}^\alpha \overline{N}_\mathrm{Q}^\beta \overline{N}_\mathrm{Q}^\gamma \overline{N}_\mathrm{Q}^\eta }}{\displaystyle \underset{\sigma i}{}}W^{\alpha \sigma }W^{\beta \sigma }W^{\eta \sigma }W^{\gamma \sigma }(W^{i\sigma })^4g_{\mathrm{ps}}^{i\sigma }g_\mathrm{b}^\sigma \overline{N}_\mathrm{Q}^\sigma (1+\delta _f^\sigma )`$
$`\widehat{\delta }_f^\alpha \widehat{\delta }_f^\beta \widehat{\delta }_f^\gamma \widehat{\delta }_f^\eta )`$
The first set of terms (second line) arise from the shot-noise- free part of $`\widehat{\delta }_f^\alpha `$, namely $`(\stackrel{~}{N}^\alpha \overline{N}^\alpha )/\overline{N}^\alpha `$. The next two sets of terms (third + fourth and fifth + sixth lines) come from combinations of $`\widehat{\delta }_f^\alpha `$ involving products of two shot-noise terms with two shot-noise-free terms. The next set of terms (seventh + eighth lines) arises from products of three shot-noise terms and a noise-free one. The next set of terms (nineth to thirteenth lines) comes from products of four shot-noise terms. The last term corresponds to what has to be subtracted off to compute the covariance.
To make further progress, we assume $`W^{i\alpha }`$ is independent of $`\delta _f^\alpha `$, which is strictly correct for $`W^{i\alpha }`$ given by eq. (6), but only roughly so for eq. (7) (see discussion in Appendix A). Then, taking the small wavelength limit in the sense that $`k,\mathrm{\Delta }k>2\pi /`$ ($`\mathrm{\Delta }k`$ is the size of a $`k`$-bin), and making use of the fact that $`W^{\alpha \sigma }W^{\beta \sigma }`$ is only non-zero at separations $`u^\alpha u^\beta `$ much less than $`1/k`$ where $`k`$ is the wavenumber of interest, we obtain
$`C(k_1,k_2)={\displaystyle \underset{\sigma }{}}(\overline{w}^{\sigma \sigma })^2{\displaystyle \frac{𝒩^3}{^2}}(`$ (47)
$`{\displaystyle \frac{2}{n_{k_1}}}[P(k_1)+{\displaystyle \frac{}{𝒩}}{\displaystyle \underset{i}{}}{\displaystyle \frac{1}{(\overline{N}_\mathrm{Q}^\sigma )^2}}(W^{i\sigma })^2(g_{\mathrm{ps}}^{i\sigma }g_\mathrm{b}^\sigma \overline{N}_\mathrm{Q}^\sigma +V_\mathrm{B}^{i\sigma }){\displaystyle \underset{\beta \gamma }{}}W^{\beta \sigma }W^{\gamma \sigma }]^2\delta ^{k_1k_2}+{\displaystyle \frac{T_{k_1k_2}}{}}`$
$`+4{\displaystyle \frac{1}{𝒩}}B_{k_1k_2}{\displaystyle \underset{i}{}}{\displaystyle \frac{1}{\overline{N}_\mathrm{Q}^\sigma }}(W^{i\sigma })^2g_{\mathrm{ps}}^{i\sigma }g_\mathrm{b}^\sigma {\displaystyle \underset{\beta \gamma }{}}W^{\beta \sigma }W^{\gamma \sigma }`$
$`+2(P(k_1)+P(k_2)){\displaystyle \frac{}{𝒩^2}}{\displaystyle \underset{i}{}}{\displaystyle \frac{1}{(\overline{N}_\mathrm{Q}^\sigma )^2}}(W^{i\sigma })^3g_{\mathrm{ps}}^{i\sigma }g_\mathrm{b}^\sigma {\displaystyle \underset{\beta \eta \gamma }{}}W^{\beta \sigma }W^{\eta \sigma }W^{\gamma \sigma }`$
$`+2P_{k_1k_2}{\displaystyle \frac{}{𝒩^2}}\left[{\displaystyle \frac{1}{\overline{N}_\mathrm{Q}^\sigma }}(W^{i\sigma })^2g_{\mathrm{ps}}^{i\sigma }g_\mathrm{b}^\sigma {\displaystyle \underset{\alpha \eta }{}}W^{\alpha \sigma }W^{\eta \sigma }\right]^2`$
$`+{\displaystyle \frac{^2}{𝒩^3}}{\displaystyle \frac{1}{(\overline{N}_\mathrm{Q}^\sigma )^4}}[(W^{i\sigma })^4(g_{\mathrm{ps}}^{i\sigma }g_\mathrm{b}^\sigma \overline{N}_\mathrm{Q}^\sigma +\stackrel{~}{N}_S^{i\sigma })]{\displaystyle \underset{\alpha \beta \gamma \eta }{}}W^{\alpha \sigma }W^{\beta \sigma }W^{\eta \sigma }W^{\gamma \sigma })`$
A few comments are in order. The terms in the third + fourth + nineth + tenth lines of eq. (46) are canceled by the last term of eq. (46). The terms in the third + fourth lines of eq. (46) contain contributions proportional to $`B(k,k,0)`$ which vanish. There is also a term from the nineth line of eq. (46) that is proportional to $`P(0)`$ which vanishes also. The shot-noise terms in the second line of eq. (47) come respectively from the fifth + sixth + nineth + tenth lines of eq. (46). The rest of the terms in eq. (47) basically follow the order they are presented in eq. (46).
Lastly, setting $`W^{\alpha \beta }=\delta ^{\alpha \beta }`$ and $`k_1=k_2`$ then recovers eq. (26).
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# Time-Dependence of the Mass Accretion Rate in Cluster Cooling Flows
## 1 Introduction
X-ray and optical data strongly suggest that cooling accretion is taking place in more than half of galaxy clusters (for a review, see Fabian 1994). In non-cooling flow clusters, X-ray surface brightness profiles obtained from imaging data are usually well modeled by a single-temperature gas in hydrostatic equilibrium with a background isothermal potential (Jones & Forman 1984). In contrast, cooling flow clusters exhibit significant excess central emission compared to that derived from a fit to an isothermal. This excess emission is thought to be due to gas losing its thermal energy and condensing out of the intracluster medium (ICM) at rates that often exceed $`100\mathrm{M}_{}\mathrm{yr}^1`$. Over a cluster lifetime, the central dominant galaxies in these clusters may accrete up to $`10^{12}\mathrm{M}_{}`$ of cooled gas. (We assume a Hubble constant of $`50\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ and that a cluster lifetime is comparable to a Hubble time.)
X-ray surface brightness profiles of cooling flow clusters can be used to infer mass accretion rates $`\dot{M}_{\mathrm{surf}}`$ which increase roughly linearly with radius from the centers of cooling flows (e.g., Stewart et al. 1984; Thomas, Fabian, & Nulsen 1987). If the flows are steady, this in turn implies that matter must be cooling and condensing out of the ICM over the full range of radii ($`100\mathrm{kpc}`$) for which cooling flow emission is observed. If there were no distributed mass deposition, then gas would be deposited only at the center. The resulting X-ray surface brightness profiles would be significantly more centrally peaked than those observed.
X-ray spectra of cooling flows imply cooling rates $`\dot{M}_{\mathrm{spec}}`$ within a factor of two of those inferred from the imaging data (e.g., Canizares, Markert, & Donahue 1988; Allen et al. 2000). This result is an independent confirmation of the accretion rates derived from imaging data. The dynamical accretion rate in a spherical system, $`\dot{M}=4\pi r^2\rho u`$ (where $`\rho `$ is the gas density and $`u`$ is the flow velocity), is not observed directly. Although the dynamical and X-ray-derived accretion rates should be the same if the flow is in a steady state, the radial velocities in the accretion flow are typically on the order of tens of kilometers per second, well below both the velocity dispersion of a single galaxy and the spectral resolution of current X-ray instruments.
Although the final repository for the cooling gas remains ambiguous, the evidence is quite strong that it is cooling and decoupling from the flow (Fabian 1994; White, Jones, & Forman 1997; Peres et al. 1998; Markevitch et al. 1998; Allen et al. 2000; White 2000). Heating processes are unlikely to balance cooling since their functional dependences on density and temperature do not match those for radiative losses. Indeed, previous numerical investigations have shown that alternatives to large cooling accretion rates, such as heat conduction (Bregman & David 1988; Meiksin 1988), supernova heating, and drag heating by orbiting galaxies (Bregman & David 1989), are able to reproduce the data only for narrow ranges in the respective free parameters, if at all. Moreover, it is difficult to imagine a process other than cooling that can power an excess emission rate as high as $`10^{44}\mathrm{erg}\mathrm{s}^1`$ without producing other noticeable signatures. Finally, most alternatives to cooling accretion fail to account for observed soft X-ray lines.
Models that assume a steady accretion flow with mass deposition can reproduce the observed X-ray features of cluster cooling flows (e.g., White & Sarazin 1987, hereafter WS; Fabian 1994). WS calculated steady-state models with a simple mass deposition formation law, but such models provide no information on the evolution of cooling flows.
In this paper we use a suite of time-dependent spherical models to assess the evolution of cooling flows in static and evolving gravitational potentials. In addition to providing predictions of the spatial structure and spectral properties of relatively relaxed cluster cooling flows, the spherical model also has implications for their long term evolution. The time evolution of the accretion rate has consequences for the total amount of accreted material in cooling flows, and it can provide a test for the models when compared to X-ray observations of high-redshift clusters. Two free parameters likely to affect the time-variation of $`\dot{M}`$ in the spherical models are the mass deposition efficiency and the ratio $`\beta `$ of gravitational binding energy to thermal energy in the gas. The parameter $`\beta `$ determines the shape of the gas distribution in clusters when the gas is in hydrostatic equilibrium.
Self-similar models can offer insight into the evolutionary effects, and they provide a useful testing ground for fully time-dependent calculations (Chevalier 1987, 1988; Bertschinger 1989; Lufkin & Hawley 1993). However, they are limited in that they allow either a narrow range of initial conditions or a limited number of physical processes. As a result, such models have not been successful in attaining detailed agreement with the observations. Simple scaling arguments are potentially useful (Bertschinger 1988; White 1988), and we review these in § 2 below. Our numerical simulations of cooling flow evolution in static gravitational potentials are presented in § 3.
As a final point of inquiry, we examine the consequences of continued cluster evolution by investigating the effects of a secularly deepening cluster potential and continued accretion of gas from the Hubble flow. We present a simple physical argument in § 2.3 which shows that adiabatic compression of gas in a deepening gravitational potential does not inhibit cooling flows. Our numerical simulations of cooling flow evolution in an evolving gravitational potential are described in § 4 and compared to similar work by Meiksin (1990, hereafter M90). Our conclusions differ from those in M90, which found that cooling flows are slow to reach steady state and mass dropout rates were strongly reduced by time-dependent effects in an evolving gravitational potential. We do not find these effects in our simulations of cooling flows evolving in the same deepening gravitational potential as used in M90, which is consistent with the simple arguments we present in § 2.3. Our results are summarized in § 5.
## 2 Physical Arguments
Before proceeding to the numerical models, we first discuss some physical arguments regarding the temporal behavior of cooling flow clusters. These include definitions of the various physical regions in the flow, a review of analytic predictions of the time dependence of the cooling accretion rate $`\dot{M}`$, and a discussion of whether adiabatic compression can extend the cooling time. These physical arguments form a basis for the discussion of the numerical models.
The most important process in models for cooling flows is obviously the radiative cooling of the gas. Over essentially all of the region of the flow, this occurs on a timescale which is much longer than the dynamical time of the gas, and as a result the flows and compression which occur are relatively slow. We do not include thermal conduction or other diffusive processes in our calculations. The only other processes which affect the temperature of the gas are shocks and adiabatic compression or expansion of the gas. In the absence of cluster mergers, the flow velocities are low and shocks are not important in our models. However, adiabatic compression of the gas will occur as a result of the cooling and inflow of the gas in the gravitational potential, or through the slow cosmological growth in the cluster potential. We discuss the effects of adiabatic heating on the cooling time scale of the gas in § 2.3 below.
### 2.1 Timescales
The general structure of cluster cooling flows in a static potential is determined by three characteristic timescales: the age of the cluster $`t_{\mathrm{age}}`$, the cooling time $`t_{\mathrm{cool}}`$ and the dynamical time $`t_{\mathrm{dyn}}`$. The instantaneous isobaric cooling time is
$$t_{\mathrm{cool}}=\frac{5}{2}\frac{kT}{\mu m_p\rho \mathrm{\Lambda }(T)}5\times 10^9\left(\frac{T}{10^8\mathrm{K}}\right)\left(\frac{\mathrm{\Lambda }}{5\times 10^{24}\mathrm{ergs}\mathrm{cm}^3\mathrm{s}^1\mathrm{g}^2}\right)^1\left(\frac{n}{10^2\mathrm{cm}^3}\right)^1\mathrm{yr},$$
(1)
where $`\mu `$ is the mean molecular weight in units of the proton mass $`m_p`$ (we assume $`\mu =0.6`$ in this work), and $`n=\rho /\mu m_p`$ is the total particle number density. The function $`\mathrm{\Lambda }(T)`$ is the cooling rate such that $`\rho ^2\mathrm{\Lambda }`$ has the units $`\mathrm{erg}\mathrm{s}^1\mathrm{cm}^3`$. We assume the cooling function of WS for half-solar abundances. The dynamical time is
$$t_{\mathrm{dyn}}10^8\left(\frac{\sigma _c}{1000\mathrm{km}\mathrm{s}^1}\right)^1\left(\frac{r_c}{100\mathrm{kpc}}\right)\mathrm{yr},$$
(2)
where $`r_c`$ is the cluster core radius and $`\sigma _c`$ is the velocity dispersion of the cluster galaxies (also nearly equal to the initial sound speed in the gas prior to cooling).
The cooling radius $`r_{\mathrm{cool}}`$ is defined as the point where the instantaneous isobaric cooling time equals the system age. Inside the cooling radius, the cooling time is less than the age of the system, and a cooling flow occurs. As long as the cooling time exceeds the dynamical time, the gas flows subsonically into the center, regulated by cooling. If the gas cools completely before reaching the center and the rate of mass drop out is not too large ($`q3`$), the flow generally passes through a sonic radius $`r_s`$, where the flow speed equals the local sound speed (Sarazin & Graney 1991). For $`r_{\mathrm{cool}}rr_s`$ the system is expected to be reasonably well-described as being in steady-state, an assumption which we test explicitly in § 3.3. The validity of the steady-state approximation at $`r_{\mathrm{cool}}`$ cannot be assessed except with time dependent models.
The above definition for the cooling radius is straightforward, as well as familiar, but it may not be the most useful. Since the gas is flowing inward as it cools, we consider whether a more relevant timescale may be the integrated isobaric cooling time
$$t_{\mathrm{int}}=\frac{5}{2P}_0^\theta \frac{\theta ^{}d\theta ^{}}{\mathrm{\Lambda }(\theta ^{})},$$
(3)
where $`\theta kT/\mu m_p`$. Typically, the integrated cooling time $`t_{\mathrm{int}}`$ is about a factor of two shorter than the instantaneous cooling time $`t_{\mathrm{cool}}`$. Consequently, gas may actually be taking part in the flow out to a considerably larger radius $`r_{\mathrm{int}}`$, where $`t_{\mathrm{int}}=t_{\mathrm{age}}`$. We comment more on this in § 3.2.
### 2.2 $`\dot{M}`$ as a Function of Time
The time-dependent calculations of §§ 3 and 4 solve explicitly for the accretion rate $`\dot{M}`$ as a function of time. However, it may also be possible to estimate $`\dot{M}(t)`$ based on the assumption that gas outside the cooling radius remains unaffected by cooling in the interior. Writing the dynamical accretion rate as $`\dot{M}=4\pi r^2\rho u`$, one might expect that the velocity can be approximated by the propagation rate of the cooling radius, so that $`\dot{M}`$ is interpreted as the rate at which gas is swept over by the cooling radius. This amounts to making the substitution $`u=dr_{\mathrm{cool}}/dt|_{t_{\mathrm{cool}}}`$. With this approximation, White (1988) showed that the cooling radius evolves with time as $`r_{\mathrm{cool}}t^\eta `$, where $`\eta =[(1\mathrm{\Delta }_T\mathrm{\Lambda })\mathrm{\Delta }_rT\mathrm{\Delta }_r\rho ]^1`$, and where we have used the notation $`\mathrm{\Delta }_x(d\mathrm{ln}/d\mathrm{ln}x)`$. The exponent of the time dependence in the accretion rate $`\dot{M}`$ is then given by
$$\mathrm{\Delta }_t\dot{M}=(3+\mathrm{\Delta }_r\rho )\eta 1\xi .$$
(4)
If the gas is initially isothermal, then $`\xi <(>)0`$ when $`\mathrm{\Delta }_r\rho <(>)3/2`$. In other words, cooling flow accretion rates will increase (decrease) with time when the cooling radius is in a region where the gas density profile is shallower (steeper) than $`r^{1.5}`$.
The gas density profiles of clusters of galaxies are often fit with the isothermal “beta” model. In this model, the gas density distribution is given by
$$\rho (r)=\rho _0\left[1+\left(\frac{r}{r_c}\right)^2\right]^{3\beta /2},$$
(5)
where $`\rho _0`$ is the central gas density (e.g., Jones & Forman 1984). We will sometimes give $`n_0\rho _0/\mu m_p`$, which is the corresponding total particle number density. If the galaxies in a cluster are isothermal with a distribution given by the analytic King approximation to an isothermal sphere, then $`\beta `$ is determined by the ratio of the velocity dispersion of the galaxies to that of the gas,
$$\beta =\frac{\mu m_p\sigma _c^2}{kT}.$$
(6)
Typically, clusters have $`\beta 2/3`$, and the density profile tends to steepen from $`\mathrm{\Delta }_r\rho 0`$ near the center to $`\mathrm{\Delta }_r\rho 2`$ at large radii. Thus, there will be a transition radius $`r_{\mathrm{tran}}`$ where $`\xi \mathrm{\Delta }_t\dot{M}`$ changes sign from positive to negative. For cases where $`r_{\mathrm{cool}}>r_{\mathrm{tran}}`$, one would infer from equation (4) that $`\dot{M}`$ is decreasing. However, there appears to be some variation in the value of $`\beta `$ from cluster to cluster, from 0.5 to about 1.2. This variation translates into a variation in the gas scale height relative to cluster core radii, and hence in the value of $`\xi `$ at the cooling radius. There is a correlation between cluster temperatures (both the gas temperatures at large radii and the dynamical temperatures \[velocity dispersions\] of galaxies in clusters) and the overall slope of the X-ray surface brightness profiles (White 1991), implying that hot clusters ($`T_{\mathrm{gas}}7\times 10^7\mathrm{K}`$) tend to have density profiles scaling as $`\rho r^2`$ at large radii, while cool clusters ($`T_{\mathrm{gas}}5\times 10^7\mathrm{K}`$) have shallower profiles. There is therefore a range of cases, with both increasing and decreasing mass accretion rates inferred via equation (4). In cases where imaging data imply a cooling flow, the slope of the X-ray surface brightness profile typically does not make a sudden inflection at either a core or cooling radius. Consequently, determining the transition radius from imaging alone is particularly difficult. Thus, in the absence of spatially resolved spectroscopy, there is some a priori ambiguity in the sign and magnitude of $`d\dot{M}/dt`$. In § 3, we investigate this question for a plausible range in $`\beta `$.
### 2.3 Adiabatic Compression
One obvious limitation to the analysis of § 2.2 is its failure to account for ongoing dynamical evolution of the cluster itself. As subclusters merge, and as matter continues to accrete from the Hubble flow, the cooling flow may be altered, or disrupted altogether. A 1-D calculation can take these processes into account only in the limit of quasi-static deepening of the cluster potential well. Observations (Bird 1994) and numerical simulations (Evrard 1990) of hierarchical structure formation show that the growth of galaxy clusters is only roughly approximated by a spherically symmetric deepening of a cluster’s gravitational potential. If the changes in the cluster potential result from violent subcluster mergers, shocks may be driven into the gas. Such mergers cannot be easily modeled in a 1-D simulation. Nonetheless, following work by M90, we consider a simple model in which the growth of the cluster is described by a slow deepening of the cluster potential and the resulting adiabatic compression of the intracluster gas. Although this is a simplistic approximation, we can still address the question of whether such adiabatic compression is sufficient to change the qualitative nature of the steady flow compared to that of an isolated, static cluster, as has been suggested by M90.
A simple physical argument shows that adiabatic compressive heating is not sufficient to balance radiative cooling in intracluster gas. The point is that, whatever the cause of such compression, it is a reversible process and consequently does not change the entropy of the gas. Thus, the only changes in the entropy are those due to radiative cooling. Consider a parcel of gas initially at temperature $`T_i`$ and density $`\rho _i`$, with an adiabatic index $`\gamma `$. For the purpose of this argument, we approximate the cooling function as a power law, $`\mathrm{\Lambda }T^\alpha `$. Adiabatic compression to density $`\rho _f`$ will heat the gas to a temperature of $`T_f=T_i(\rho _f/\rho _i)^{\gamma 1}`$. The cooling time scales as $`T^{1\alpha }\rho ^1`$, so, after compression, the cooling time is
$$t_{\mathrm{cool},f}=t_{\mathrm{cool},i}\left(\frac{\rho _f}{\rho _i}\right)^{\gamma +\alpha \gamma \alpha 2}.$$
(7)
For $`\gamma =5/3`$ the effect of compressional heating is exactly balanced by the increased cooling rate for a critical value $`\alpha =\alpha ^{}=1/2`$. We note that $`\alpha 1/2`$ for almost all temperatures in the range $`10^6T10^9`$, except for a narrow range between $`T10^7\mathrm{K}`$ and $`T3\times 10^7\mathrm{K}`$ (e.g., Raymond, Cox, & Smith 1976; WS). Thus, adiabatic compression generally should decrease the cooling time for a given parcel of gas. This result is independent of the physical processes responsible for the compression.
## 3 Time-Dependent Flow in a Static External Potential
### 3.1 Assumptions and Equations
The models in this section are designed to resemble those of WS for the purpose of assessing time-dependent effects in existing models with mass deposition. We therefore include radiative cooling and a mass deposition term, but ignore the effects of conduction and external heating, referring discussion of the latter to the published literature. Neglecting angular momentum, the corresponding spherically symmetric fluid equations are
$`{\displaystyle \frac{d\rho }{dt}}+{\displaystyle \frac{\rho }{r^2}}{\displaystyle \frac{}{r}}(r^2u)=\dot{\rho }`$ (8)
$`{\displaystyle \frac{du}{dt}}+{\displaystyle \frac{1}{\rho }}{\displaystyle \frac{P}{r}}+{\displaystyle \frac{\mathrm{\Phi }}{r}}=0`$ (9)
$`P{\displaystyle \frac{d}{dt}}\mathrm{ln}(P\rho ^\gamma )=(\gamma 1)\rho ^2\mathrm{\Lambda }(T),`$ (10)
where $`P`$ and $`T`$ are the pressure and temperature respectively, $`\rho ^2\mathrm{\Lambda }`$ is the cooling rate in $`\mathrm{erg}\mathrm{s}^1\mathrm{cm}^3`$, $`\mathrm{\Phi }`$ is the total gravitational potential and
$$\frac{d}{dt}\frac{}{t}+u\frac{}{r}.$$
For the mass dropout rate $`\dot{\rho }`$ we adopt the cooling-time law of WS,
$$\dot{\rho }=q\frac{\rho }{t_{\mathrm{cool}}},$$
(11)
where $`q`$ is an efficiency factor of order unity and $`t_{\mathrm{cool}}`$ is defined in equation (1). This mass deposition law is based on the ansatz that thermal instabilities due to cooling lead to growing perturbations which leave the flow over a range of radii. (It is often expedient in numerical simulations to cut off the cooling function at some floor temperature $`T_{\mathrm{floor}}`$, the value of which is arbitrary provided that the gas temperature does not reach $`T_{\mathrm{floor}}`$ unless the flow becomes supersonic.)
We solve equations (8)–(10) using a time-explicit hydrodynamics code in spherical symmetry (Lufkin & Hawley 1993). The code is similar to existing second-order Eulerian codes except that it solves the fluid equations in Lagrangian mode with a globally conservative remap to the fixed grid. The code has been tested on a variety of problems, including self-similar cooling flows and cosmological accretion (for details, see Lufkin & Hawley 1993).
#### 3.1.1 The Assumed Mass Deposition Efficiency $`q`$
All of the cooling flow models considered here have a finite mass deposition efficiency coefficient $`q`$. If one assumes no mass deposition, the mass accretion rate is constant with radius in a steady-state flow, and the gas cools catastrophically at the center (WS; Meiksin 1988, 1990), i.e., the cooling is sufficiently strong for the gas to pass through a sonic point as it approaches the center. The catastrophe is prevented (i.e. the flow remains subsonic all the way to the center) for $`q3.4`$. The dynamical mass accretion rate $`\dot{M}`$ is also found to be roughly proportional to radius for $`q34`$. Although a smaller $`q`$ cannot be ruled out observationally in all cases, a value as small as $`q=0`$ implies an X-ray surface brightness profile that is much too sharply peaked in the center. We therefore restrict our attention to the two cases $`q=1`$ (below which the central density is too sharply peaked) and $`q=4`$, and examine the effect of this variation on model flows.
For these applications, the numerical grid covers the full dynamic range from scales of $`1\mathrm{kpc}`$ near the center out to a radius of several Mpc, with a reflecting inner boundary ($`u0`$ as $`r0`$). In the $`q=1`$ case there is generally a sonic radius near $`r1\mathrm{kpc}`$ in the steady-state models. The sonic radius is therefore unresolved by the grid spacing of $`1\mathrm{kpc}`$ near the center in the time-dependent models. If an outflow boundary inner condition is used, the grid must resolve the region between the sonic radius and the origin, in order to have supersonic outflow at the inner grid boundary (Lufkin & Hawley 1993). Unfortunately, timestep constraints prevent the use of an outflow boundary in this case, owing to gas velocities in excess of $`100\mathrm{km}\mathrm{s}^1`$ inward from the sonic radius.
#### 3.1.2 The Background Cluster
The initial conditions for the time-dependent simulations are chosen such that the gas is isothermal and in hydrostatic equilibrium with the external potential $`\mathrm{\Phi }`$. Integrating the hydrostatic equation we obtain for the initial density profile
$$\rho (r)_{t=0}=\rho _0\mathrm{exp}\left\{[\mathrm{\Phi }(r)\mathrm{\Phi }_0]/(kT_{\mathrm{}}/\mu m_p)\right\},$$
(12)
where $`T_{\mathrm{}}`$ is the asymptotic gas temperature at large radii, and $`\rho _0`$ and $`\mathrm{\Phi }_0`$ are the density and potential at $`r=0`$. The background potential for the cluster gas is assumed be that of a massive central galaxy $`\mathrm{\Phi }_g`$ plus a smooth cluster component $`\mathrm{\Phi }_c`$. The cluster mass distribution is taken as a King approximation to an isothermal sphere:
$$\rho _c(r)=\rho _{c0}\left[1+(r/r_c)^2\right]^{3/2},$$
(13)
where $`r_c`$ is the cluster core radius. The potential for this distribution is given by
$$\mathrm{\Phi }_c(r)=\mathrm{\Phi }_{c0}\frac{\mathrm{ln}\left[r/r_c+\sqrt{1+(r/r_c)^2}\right]}{r/r_c}.$$
(14)
The central potential $`\mathrm{\Phi }_{c0}=4\pi G\rho _{c0}r_c^2`$ and is related to the line-of-sight velocity dispersion $`\sigma _c`$ by $`\mathrm{\Phi }_{c0}=9\sigma _c^2`$. We assume $`\sigma _c=1054\mathrm{km}\mathrm{s}^1`$ ($`\mathrm{\Phi }_{c0}=10^{17}\mathrm{erg}\mathrm{g}^1`$) for the hot cluster runs, and $`\sigma _c=527\mathrm{km}\mathrm{s}^1`$ for the cool cluster runs. The corresponding cluster masses inside of 250 kpc are $`10^{14}`$ and $`5\times 10^{13}\mathrm{M}_{}`$, respectively.
We note that our adopted potential has a finite depth, whereas a true isothermal is infinitely deep. Consequently, the density profile for the gas becomes sufficiently shallow at large radii that the total X-ray luminosity diverges. This actually follows from assuming that the gas at large radii is hydrostatic and isothermal, which cannot be true on spatial scales where the sound crossing time is greater than a Hubble time ($`r10\mathrm{Mpc}`$). In the evolving cluster models of § 4, the outer parts of the cluster are in free fall from the cold Hubble flow, with an accretion shock at a radius of 3–10 Mpc. Because the cooling flow is confined to the inner 100 kpc or so, the solution is not sensitive to the structure of the cluster much beyond this radius. We therefore quote total luminosities within one half and two Mpc. Furthermore, because gas densities far from the center are not well constrained observationally, we omit discussion of the structure of the ICM at very large radii.
#### 3.1.3 The Central Galaxy
All of the models here and below assume the presence of a massive, stationary central galaxy. However, in choosing a form for the galactic potential we are guided by the desire to extend these models to higher dimensions, with the possibility of allowing a galaxy to orbit in the cluster (e.g., Lufkin, Balbus, & Hawley 1995). The primary considerations are computational efficiency and accuracy when the galaxy’s position changes with time. The potential should be smooth and fully resolved, with $`\mathrm{\Phi }_g0`$ as $`r0`$ to avoid grid noise. It should also be expressible in terms of simple analytic functions. A potential of the form of equation (14) would be one possibility. However, the mass in this model diverges logarithmically, and it is generally preferable to have a galaxy potential corresponding to a finite mass. While it is tempting to truncate an analytic King model at some finite radius, this can generate anomalous sound waves for an orbiting galaxy on a finite difference grid.
A simple form which satisfies our numerical criteria is Plummer’s model (Binney & Tremaine 1987):
$$\mathrm{\Phi }_g=\mathrm{\Phi }_{g0}\left(1+r^2/r_g^2\right)^{1/2},$$
(15)
where $`r_g`$ is a characteristic radius. As it stands, the density distribution which gives rise to equation (15) does not approximate real galaxies well. However, we have found that a superposition of suitably weighted components can give a good approximation to a King model. For example, the addition of two potentials, one for the center and one for an extended distribution, can be written as
$$\mathrm{\Phi }_g=\mathrm{\Phi }_{g0}\left[\frac{(1y)}{(1+r^2/r_g^2)^{1/2}}+\frac{y}{(1+r^2/r_t^2)^{1/2}}\right],$$
(16)
where $`0<y<1`$ and $`r_t`$ is the characteristic radius of the extended component. We have chosen to specify the potential directly, rather than starting from the underlying density distribution. This is because the gas interacts with the galaxy only through gradients in the potential. The analytic form we have chosen contains functions which can be calculated at high speed, whereas transcendental functions (such as power laws, logarithms and arctangents) can increase the run time significantly. The square root function is an exception. It uses a divide-and-average algorithm, which converges very fast. A tabulated function, although not particularly expensive, can be rather unwieldy, especially when a galaxy’s position changes with time.
At large radii, ($`rr_t`$), the potential tends to that of a point mass, so that the total mass has the finite value
$$M_g=\frac{[(1y)r_g+yr_t]\mathrm{\Phi }_{g0}}{G}.$$
(17)
For $`y=0.15`$ and $`r_t^2=40r_g^2`$, the gradient of this potential closely approximates that of an analytic King model truncated at $`r=20r_g`$.
We note that in equations (13) and (14), $`r_c`$ is a true core radius (i.e., it corresponds to the radius at which the surface density falls to one half of its central value). In the model galaxy chosen here, this occurs at a radius of about $`0.63r_g`$. Figure 1 shows $`g=\mathrm{\Phi }_g/r`$ for a King model in which the central potential is $`\mathrm{\Phi }_{g0}=7\sigma _g^2`$ (solid line; see Binney & Tremaine 1987), compared with truncated (short dash) and nontruncated (dotted) King approximations to an isothermal sphere. These are plotted along with a model galaxy with a potential of the form of equation (16), for $`y=0.15,r_g=1.58r_c,r_t=10r_c`$ (long dash). The vertical scale is in arbitrary units, and can be set either by the central potential or the total mass of the galaxy $`M_g`$.
In the simulations below, the core radii $`r_c`$ are 250 kpc for the cluster and 4.41 kpc for the central galaxy ($`r_g=7\mathrm{kpc}`$). Assuming a total mass for the galaxy of $`2\times 10^{12}\mathrm{M}_{}`$ implies a line-of-sight velocity dispersion of $`288\mathrm{km}\mathrm{s}^1`$ at the center and $`251\mathrm{km}\mathrm{s}^1`$ averaged over a circular aperture subtending a radius of $`10\mathrm{kpc}`$, in the manner of Bailey & MacDonald (1981). The composite density distribution can also be reasonably approximated by an NFW profile (Navarro, Frenk, & White 1997) with a scaling radius of $`r_s200`$ kpc.
### 3.2 Results
Table 1 lists the input parameters and some of the calculated global characteristics for ten simulations. Each of five model clusters are evolved for $`15\mathrm{Gyr}`$ for $`q=1`$ and $`q=4`$. The first model, (runs 1 and 2) is taken to be typical of a strong cooling flow in a hot cluster, and is discussed in detail in §§ 3.2.1. The second model (runs 3 and 4) corresponds to a relatively cool cluster and is discussed in § 3.2.2. The remaining three cluster models, also discussed in § 3.2.2, are used to explore a range of values in $`\beta `$. The results of runs 1–4 are plotted in Figures 2 – 9, and are compared to steady-state calculations in § 3.2.3.
The physical assumptions, as stated above, are identical to those in WS, with the exception of the form of the external potential. Self-gravity is neglected for the static models, but is included in the evolving cluster models in § 4 below. Because the cooling flow is primarily pressure-driven, the inclusion of self-gravity has little effect.
#### 3.2.1 Hot Cluster Runs
Shown in Figure 2 are the density, temperature, mass accretion rate and cooling times as functions of radius at various times in the evolution, for the hot cluster run with $`q=1`$. The initial conditions, given in Table 1, are chosen so that the models have reasonable global properties at $`t=10^{10}\mathrm{yr}`$ in the evolution. In discussing the properties of the models, we use an age of $`10^{10}`$ yr as a fiducial time for comparison to observed clusters. The solid curves in the figures refer to this fiducial age. More generally, the dotted, short-dashed, long-dashed, solid and dot-dashed lines correspond to ages of $`0,3,6,10\mathrm{and}15\mathrm{Gyr}`$, respectively. The curve at 15 Gyr (beyond the fiducial age) is included to show that a steady state has been reached.
Figure 3 shows the time evolution of both the instantaneous and integrated cooling radii, $`r_{\mathrm{cool}}`$ and $`r_{\mathrm{int}}`$. We also give the dynamical mass accretion rate $`\dot{M}=4\pi r^2\rho u`$ as determined at each of these cooling radii. The X-ray-derived mass accretion rate (or cooling rate) $`\dot{M}_X`$ is determined by assuming that all of the X-ray luminosity from within the cooling radius is due to gas cooling within that radius. This is the standard assumption made in analyses of the X-ray surface brightness profiles to derive the total cooling or accretion rates $`\dot{M}_{\mathrm{surf}}`$ (e.g., Thomas et al. 1987). (Recall that $`\dot{M}_X=\dot{M}`$ only for steady state flows.) As expected, there is close agreement between $`\dot{M}`$ and $`\dot{M}_X`$ at each radius, although the value of $`\dot{M}`$ at $`r_{\mathrm{int}}`$ is nearly twice that at $`r_{\mathrm{cool}}`$. The panel at lower right shows the value of $`\xi \mathrm{\Delta }_t\dot{M}`$ as derived from local gradients in the simulation (eq. 4). We find $`\xi 0.29`$ at $`r_{\mathrm{int}}`$ and $`\xi 0.59`$ at $`r_{\mathrm{cool}}`$. The value at $`r_{\mathrm{cool}}`$ compares reasonably to a value of $`\mathrm{\Delta }_t\dot{M}0.34`$ measured directly from the simulation by least-squares fitting of a straight line to $`\mathrm{log}\dot{M}(t)`$ as measured at $`r_{\mathrm{cool}}`$ and dumped every $`10^8`$ years.
We determine the transition radius (where $`\xi =0`$ in eq. ) by examining the numerical results directly: $`r_{\mathrm{tran}}=234\mathrm{kpc}`$ at $`t=10^{10}\mathrm{yr}`$. The measured cooling radius at this time is $`r_{\mathrm{cool}}=100\mathrm{kpc}<r_{\mathrm{tran}}`$, so we would infer from equation (4) that the accretion rate is increasing, as Figure 3 confirms. There are several sources of error in the estimate of $`r_{\mathrm{cool}}`$. The grid spacing is about $`5\mathrm{kpc}`$ in this region, although this can be reduced if necessary by running the problem at higher resolution. More serious is the assumption that the relevant characteristic radius is $`r_{\mathrm{cool}}`$, not $`r_{\mathrm{int}}`$ (see eqs. 1 and 3). At $`t=10^{10}\mathrm{yr}`$ in the $`q=1`$ case, the location where $`t_{\mathrm{int}}=t_{\mathrm{age}}`$ is $`r_{\mathrm{int}}=151\mathrm{kpc}<r_{\mathrm{tran}}`$. Thus we see that the scaling relation of § 2.2 is confirmed in this case for both measures of the cooling radius.
Several other features are noticeable from the figures. Gas in the core has cooled completely between 6 and 10 Gyr, with little change in the flow thereafter. The initial cooling time at the center is 10.8 Gyr, but the integrated isobaric cooling time is 6.43 Gyr. This time period corresponds to the phase in which the cooling radius and mass accretion rates slow to a steady rate of increase (Figure 3). The particular shape of the temperature profile is determined by the cooling function and the shape of the potential, while the density and accretion rate profiles are nearly featureless. Because the sonic radius is unresolved, the innermost grid zones reflect low-amplitude sound waves into the flow. These are visible in the accretion rate at 15 Gyr, but do not significantly affect the flow.
One might expect an inflection point in the density profile as the center begins to cool, but this is not observed. The entire inner region cools almost simultaneously, but isobarically, resulting in a smooth transition to the steeper density profile seen at later times.
The hot cluster $`q=4`$ solution (run 2) is shown in Figures 4 and 5. This model has a higher rate of mass drop out at larger radii. In this case, the initial density required to recover an accretion rate of about $`300\mathrm{M}_{}\mathrm{yr}^1`$ at $`t=10\mathrm{Gyr}`$ is sufficiently high ($`n_0=0.033\mathrm{cm}^3`$) that the center cools in about $`5\mathrm{Gyr}`$. After this time we find that $`\dot{M}`$ is slowly decreasing, even though both cooling radii ($`r_{\mathrm{cool}}=74\mathrm{kpc};r_{\mathrm{int}}=128\mathrm{kpc}`$) lie well inside $`r_{\mathrm{tran}}=354\mathrm{kpc}`$. Moreover, the logarithmic rates of change of $`\dot{M}_X`$ measured at both $`r_{\mathrm{int}}`$ and $`r_{\mathrm{cool}}`$ are approximately equal $`\xi 0.45`$, whereas the value predicted by equation (4) goes no lower than $`\xi =0.2`$ for this time frame. The value of $`\eta d\mathrm{ln}r_{\mathrm{cool}}/d\mathrm{ln}t`$ is also discordant; we measure $`\eta 0.5`$ from the simulations, whereas the density and temperature profiles would suggest a value of $`\eta 1.0`$. Thus, it appears that the scaling argument detailed in § 2.2, which assumes $`q=0`$, breaks down for large values of $`q`$. We will comment more on this below. Because there is no sonic point in this model, there is no need for high resolution in the innermost regions of the flow.
#### 3.2.2 Cool Cluster Runs
We perform a second pair of simulations in which all the input parameters are the same as those above except that the cluster has a central velocity dispersion of $`\sigma _c=527\mathrm{km}\mathrm{s}^1`$ and the asymptotic gas temperature is $`T_{\mathrm{}}=3\times 10^7\mathrm{K}`$. These parameters imply $`\beta =0.6`$ (run 3 and 4 in Table 1), compared with $`\beta =0.75`$ in the hot cluster simulations. The results for runs with $`q=1`$ and $`q=4`$ are shown in Figures 6–9. The same general comments regarding the hot cluster models also apply in this case. The $`q=1`$ case shows an accretion rate $`\dot{M}`$ that is increasing with time, in agreement with the arguments in § 2.2, while the $`q=4`$ case shows a nearly constant $`\dot{M}`$ even though the predicted $`\xi `$ is clearly positive near the cooling radius and nowhere negative.
There is some variation in the particular values of $`\xi `$ and $`\eta `$ compared with the hot cluster runs. To see if these quantities correlate with the value of $`\beta `$ in these models, we ran a series of cool cluster models in which only the initial temperature was allowed to vary. All other external parameters are the same as those in runs 3 and 4. These are listed as runs 5–10 in Table 1, and the derived values of $`\eta `$ and $`\xi `$ (taken at $`r_{\mathrm{cool}}`$ at $`t=10^{10}\mathrm{yr}`$) are plotted versus $`\beta `$ in Figure 10. Other than the extreme case of $`\beta =0.36`$, (which is included only to show that we recover the expected limiting case $`\xi ,\eta \mathrm{}`$ as $`\beta 0`$), the variation with $`\beta `$ is slight compared to the dependence of $`\xi `$ and $`\eta `$ on $`q`$. In every case the values of $`\xi `$ and $`\eta `$ derived by measuring the time derivatives directly from the models are much closer to the expected values when $`q=1`$ than when $`q=4`$. We must therefore conclude that the scaling argument of § 2.2 applies only in the limit of small $`q`$.
The reason for the inability of equation (4) to predict the flow behavior for large $`q`$ is not immediately clear, although there are several ways in which the flow may be affected by a large value of $`q`$. One is that the rapid mass deposition can significantly reduce the gas density in the interior, thereby reducing the pressure, so that gas slumps to the center from radii well beyond the cooling radius. Unfortunately, there does not appear to be a well-defined position at which to evaluate $`\xi `$ from observed density and temperature profiles. This is because $`q`$ is not known in real clusters, and it may be that the relevant cooling radius should be redefined again to account for the high rate of mass deposition in this model. Another possible cause for concern is that if the mass deposition formula (eq. 11) is applied throughout the flow, it implies that matter can condense out even if the local cooling time exceeds the age. To guard against this, we ran models 5–10 with a cutoff to the mass deposition ($`q0`$) for $`t_{\mathrm{cool}}>t_{\mathrm{age}}`$. We also repeated the first four runs with this modification, but there were no differences in the results other than a leveling-off in the dynamical accretion rate at the cooling radius.
### 3.3 Correspondence to Steady-State Models
We have seen that the dynamical mass accretion rate is a good approximation to the accretion rate one would infer, under the assumption of steady flow, from X-ray observations of the model clusters described above. A further test of the assumption of steady flow is to compute steady-state models corresponding to a fixed time in the evolutionary models. The equations for steady flow are given by WS, and are obtained by setting all partial time derivatives to zero in equations (8)–(10). The resulting ordinary differential equations are then solved numerically subject to appropriate boundary conditions. Because the time-dependent calculation used here and the ODE solver are independent, this is also a further test of the hydrodynamics code.
There are a number of possible ways to specify the three boundary conditions on the steady-state models. Following WS, we choose to fix the age ($`t_{\mathrm{age}}=10^{10}\mathrm{yr}`$), and values for the temperature $`T|_{r_{\mathrm{cool}}}`$ and the dynamical mass accretion rate $`\dot{M}|_{r_{\mathrm{cool}}}`$ at the cooling radius. For direct comparison with the time-dependent models, we can take the values from Table 1. The solution to the steady flow equations is plotted in Figures 11 (hot cluster) and 12 (cool cluster). These curves can be compared directly to Figures 2 and 4 at $`t=10\mathrm{Gyr}`$ (solid lines). The agreement is very close, with only a small departure near $`r=0`$ in the models with $`q=1`$. This difference is not unexpected given the artificial boundary conditions $`u|{}_{r=0}{}^{}=0`$ imposed in the time-dependent calculations on flows with sonic radii.
## 4 Time-Dependent Flow in a Static Evolving Potential
One possible limitation to the preceding models is the neglect of cluster evolution. This can consist of dynamical evolution in the form of continued mergers or two-body encounters which act to deepen the cluster potential, or of continued infall of gas from the Hubble flow. The gas itself may contribute to the change in the total potential. M90 included all of these effects in a series of simulations in spherical symmetry. Two classes of models were considered: one with a centrally dominant galaxy (because all cooling flows have one) and one without a central galaxy. We reexamine the effects of cluster evolution on cooling flows and use the same cluster models as M90 so we may compare results. We are interested in whether the flows reach a steady state with or without mass deposition.
### 4.1 Model Characteristics
In the models of M90, the cluster potential is as described above in equation (14), with an initial velocity dispersion of about $`500\mathrm{km}\mathrm{s}^1`$. The depth of the potential is assumed to increase gradually with time until, at the end of each run, the velocity dispersion is about $`1000\mathrm{km}\mathrm{s}^1`$. This is done to mock the dynamical evolution of the cluster from a redshift of $`z=1.5`$ to $`z=0`$, corresponding to a time span running from 3.3 Gyr to 13.1 Gyr (for $`\mathrm{\Omega }_0=1`$ and $`H_0=50\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$). The gas was assumed to trace the dark matter initially and the central gas temperature put as much energy per unit mass in the gas as in the dark matter. The initial temperature profile was assumed to be adiabatic.
Models without a central galaxy assumed that the initial gas velocity follows the same profile as that for collisionless infall onto a region of overdensity equivalent to that given by the total cluster mass interior to each point. The gas is therefore at rest only at $`r=0`$ and at the turnaround radius far outside the cluster, with infall interior to the turnaround radius. Initially the gas in the cluster is thus not in hydrostatic equilibrium, but instead falls inward, then bounces adiabatically. The gas in the center relaxes on a dynamical timescale, which is much less than the cooling time, so the precise initial conditions do not seriously affect the subsequent evolution of the flow. M90 found that gas will cool catastrophically at the center of such flows if $`q=0`$, but that for $`q=1/2`$ the flow is stabilized against runaway cooling. (Note that M90 employs a different definition of $`q`$, and the value of $`q=1/2`$ in M90 corresponds to $`q=5/6`$ in our notation.)
We are primarily interested in the models of M90 which contain a central galaxy in the initial conditions. In this case, the initial velocity profile is the same as in the case without a central galaxy except that $`u=0`$ for $`r<250\mathrm{kpc}`$, resulting in a discontinuous initial velocity profile.
### 4.2 Numerical Simulations
Figure 13 shows our simulations using the PLPC code (Lufkin & Hawley 1993) for the fiducial model of M90 with a central galaxy and without mass dropout ($`q=0`$). All of the parameters given in M90, including the potential for the central galaxy and the self gravity of the gas, are duplicated here. (We use a different cooling function, which we comment on below.) Up until the second time frame, our results and M90 agree in detail. We note that the discontinuous velocity profile results in an impulsive squeeze on the initially hydrostatic atmosphere interior to $`r_c`$. This strong pressure wave travels into the center, bounces off the reflecting inner boundary, and travels back out into the infalling gas beyond $`r_c`$. This is evident in both the velocity and temperature profiles at $`t=4\mathrm{Gyr}`$ (dashed line in Figure 13). We find, using fine time resolution in the data dumps, that two strong bounces occur before the gas relaxes at the center. (The relevant plots in M90 also show evidence for these bounces.) The run begins at $`t=3.3\mathrm{Gyr}`$, when the central instantaneous isobaric cooling time is $`9.0\mathrm{Gyr}`$.
For time slices later than 4.0 Gyr, our results differ from those in M90. We find that the gas cools in less than an initial cooling time, with an ensuing cooling catastrophe at the center ($`t=9.9\mathrm{Gyr}`$, dot–dash line of Figure 13). For the same model, M90 found no cooling catastrophe.
We have made a number of tests to search of the source of this discrepancy with the results of M90. First, the calculations were repeated using exactly the same grid of 150 zones as used in M90:
$`r_i=`$ $`1.88\mathrm{kpc}(i1)`$ $`i21`$
$`r_i=`$ $`1.046r_{i1}`$ $`i>21.`$
For comparison, we also made plots which were 1-2-1 smoothed (as in M90), but detected no variation from those shown in Figure 13. Second, we considered the effects of differences in the cooling function in the two sets of models. The cooling function adopted in M90 is that of Gaetz & Salpeter (1983), whereas we have used that of WS, which is somewhat stronger. We therefore ran the case with a central galaxy with the strength of the cooling function reduced by 20% to compensate. The result was virtually unchanged. This is not surprising, since the cooling rate is proportional to the square of the density. It therefore appears unlikely that the disparate results are caused by small differences in the cooling functions. Finally, we have run the fiducial model with two other hydro codes, the direct-Eulerian MONO scheme of Hawley, Smarr, & Wilson (1984) and a modified version of the implicit scheme of Ruppell & Cloutman (1975). Although the exact time of cooling collapse varies by about 20% (this measure is very sensitive to the value of the central density, which in turn is sensitive to the numerical diffusion as the impulsive compression waves bounce off the center early in the simulation), the density, velocity and temperature profiles agree in detail with the calculation of Figure 13.
One possible reason for this discrepancy is the grid noise (zone-to-zone variations in the fluid variables) present in the numerical calculations of M90. This could cause sound waves to be amplified or converted into heat energy, perhaps via artificial viscosity. We also note that the crests in the temperature and velocity profiles are sharper in our calculation than in the figures of M90. This is true even when our model output is smoothed identically to that present in M90. This suggests that the features in the plots in M90 are spread over several grid zones, which might indicate that a significant level of numerical diffusion was present. Such diffusion may have transported heat energy into the core.
We find that a cooling catastrophe occurs either with or without the central galaxy. We do reproduce the result in M90 that the catastrophe occurs sooner in the case without a central galaxy even though the initial cooling time is longer. As suggested in M90, this is likely due to the homologous collapse of the constant-density core.
When mass deposition is included, the flow reaches a steady state on an initial cooling timescale. The details in the flow variables appear to be affected by transient features left over from the discontinuous initial conditions ($`q=1`$ and 4; Figures 14–17). In particular, there are discontinuities in the position of the cooling radius due to the shock waves that bounced off the center at the beginning of the runs. As expected, the temperature beyond 100 kpc increases gradually due to continued shock heating from secondary infall.
The agreement between the cooling and dynamical accretion rates is also quite close, as in the static models of § 3. Thus, we conclude that flows with mass deposition do reach a steady state when cluster evolution and secondary infall are included in the spherical model. For both runs, $`\dot{M}`$ stays constant to within a factor of two after the flows reach a steady state. Evolutionary effects in spherical symmetry therefore do not appear to affect the time evolution of the mass accretion rate relative to that found for static clusters.
## 5 Discussion and Conclusions
We have solved for the time-dependent behavior of cooling gas in a variety of spherically symmetric cluster models, with particular emphasis on the evolution of the mass accretion rate.
We find that the steady-state approximation is valid after the initial onset of cooling at the center for cluster flows in both static and evolving external potentials. This result is insensitive to the inclusion of self-gravity. For the models considered here the accretion rate either increases or stays about the same with time. While the rates of increase or decrease that we see in the simulations would be difficult to infer from imaging observations, spatially resolved spectra can be a useful diagnostic (e.g., Wise & Sarazin 1993). The difference arises from the fact that while the age is an assumed parameter in the models, it can be measured by calculating the cooling time at the radius of extent of soft X-ray lines. Moreover, the cooling rates derived from imaging are more model-dependent, than those derived from spectra, which can in principle be obtained from a measure of the mass cooling through a single line.
A decreasing accretion rate with time is seen only for cases where the temperature of the gas corresponds closely to the cluster velocity dispersion ($`\beta =1`$). In a number of clusters, especially poorer ones, the gas may be hotter by a factor of two than the virial temperature measured from galaxy motions. Thus, although we cannot rule out a decreasing mass accretion rate on the basis of these experiments, nearly constant or increasing accretion rates are favored for most clusters, unless they have been strongly affected by mergers. In the absence of mergers, the spherical model with mass deposition would predict that cooling flows were not much more vigorous in the recent past than they are today. Attempts to include additional evolutionary effects via secular deepening of the cluster potential well and continued accretion from the Hubble flow in spherical symmetry do not alter these results. However, it is likely that cluster mergers can have a strong effect on cooling flows (e.g., McGlynn & Fabian 1984; Gómez et al. 2000). Further study of the effect of cluster evolution on cooling flows will require numerical simulations in three dimensions with high spatial resolution and including cooling.
We would like to thank Steve Balbus, Joel Bregman, John Hawley, Brian McNamara, Avery Meiksin, Bill Sparks and Mike Wise for numerous informative discussions. E. A. L. acknowledges the support of NSF grants PHY90-18251 and AST-8919180. C. L. S. was supported in part by NASA Astrophysical Theory Program grant NAGW–2376. R. E. W. III was supported in part by the NSF and the State of Alabama through EPSCoR II and by a National Research Council Senior Research Associateship at NASA GSFC.
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# UVES observations of QSO 0000–2620: oxygen and zinc abundances in the Damped Ly𝛼 galaxy at z{_{𝑎𝑏𝑠}}=3.3901 1footnote 11footnote 1Based on observations made with the ESO 8.2m KUEYEN telescope operated on Paranal Observatory, Chile.
## 1 Introduction
Damped Ly$`\alpha `$ (DLA) systems ($`\mathrm{log}`$ N(HI) $``$ 20.3 cm<sup>-2</sup>) are seen in absorption in the QSO spectra up to the highest redshift. They likely trace the field population of galaxies with no a priori selection based on morphology, since biases apply to the background sources at first approximation and not to the intervening systems. The identification of the galaxy population associated with the DLAs has important bearings on cosmological elemental evolution and has been actively debated.
Wolfe et al. (1995) and Prochaska & Wolfe (1997,1999) suggested a close connection between the DLA galaxies and the progenitors of present-day spiral galaxies. The connection is based on the fact that the comoving gas density of the DLA galaxies at high redshift roughly equals the local density provided by the luminous matter in nearby spiral galaxies. In addition, the kinematics of the metal lines has been interpreted as having formed in massive rotating disks. However, Rao & Turnshek (2000) did not find a decrease in the comoving gas density of DLA galaxies at low redshifts as expected if the gas had been converted into stars, and Haehnelt, Steinmetz, & Rauch (1998) and Ledoux et al. (1998) did not find the kinematical argument very compelling.
DLA systems offer the possibility to measure abundances accurately for a wide range of elements in a variety of normal, gas-rich galaxies spanning a considerable look-back time. The metallicity of DLA galaxies is generally inferred from the absolute abundance of the volatile element Zinc which is found at $``$ 10% of solar, with no clear evolution over the redshift interval from $`z_{abs}`$ $``$ 0.5 up to 3 (Pettini et al. 1999). The apparent lack of evolution has been considered a problem for the interpretation of DLAs as proto-spirals, in particular at low redshift where abundances much closer to the solar ones are expected. The models of Fritze et al. (1999) for Sa and Sd galaxies bracket the Zn data from the high and low metallicity sides over the entire redshift range, but require that DLAs are biased against early type spirals at low redshift. A possible bias is represented by the low gas and the high dust contents of early type galaxies at low redshift. On the same line, Prantzos and Boisser (2000) found the observed Zn abundances consistent with a spiral origin once the empirical relation 18.8$``$ \[Zn/H\]+$`\mathrm{log}N(HI)`$ 21 is imposed to the models. Jimenez et al. (1999) found that the population of Low Surface Brighness galaxies (LSB) better matches the observed zinc metallicities. Imaging of DLA galaxies at relatively low redshift revealed that the population of DLA galaxies is not dominated by spirals but rather composed by a variety of morphological types which include LSB, dwarf irregulars and late type spirals. (Le Brun et al. 1997; Rao & Turnshek 1998).
Relative elemental abundances are an important complement to the information provided by absolute abundances such as zinc metallicity. In particular, ratios of elements which are produced by Type I and Type II supernovae in different proportions provide independent insights into defining the kind of chemical evolution. Despite the efforts of different authors, it has been difficult to establish the \[$`\alpha `$/Fe\] abundance pattern in DLA systems because of the different level of depletion of refractory elements onto dust grains. In the compilation of Savaglio et al. (2000) the average of 37 measures is $`<`$\[Si/Fe\]$`>`$ +0.43 $`\pm `$ 0.18 <sup>2</sup><sup>2</sup>2Using the customary definition \[X/H\]= log (X/HI) - log (X/H)☉. At face value this is consistent with what is observed in the Galactic halo stars, which has been interpreted as evidence for an origin of DLAs in proto-spirals (Lu et al. 1996, Prochaska & Wolfe 1999).
Evidence that at least in some DLAs relative abundances do not conform to those of the Galactic halo stars was first provided by Molaro et al. (1996) in a first study of the abundances in the DLA at $`z_{abs}`$=3.3901 towards QSO 0000–2620, which is considered again in this paper. Molaro, Centurión, & Vladilo (1998) and Centurión et al. (2000) also found the ratio of the undepleted elements Sulphur and Zinc to be approximately solar in a sample of half a dozen DLAs. Vladilo (1998) corrected the abundances of Si and Fe for the differential elemental depletion and obtained approximate solar ratios (\[Si/Fe\]$``$0) in all cases investigated, thus suggesting that solar ratios might be rather common among DLAs. Pettini et al. (2000) considering DLAs with low dust depletion found \[Si/Fe\] ratios broadly in line with Galactic stellar values although there are also examples of near-solar \[Si/Fe\] ratios at \[Fe/H\]$`<`$–1. Savaglio et al. (2000), after correcting for dust, found evidence for overabundance of a factor 2 in Si in about 1/4 of the DLA systems.
The \[O/Zn\] ratio is probably the best diagnostic tool for the \[$`\alpha `$/Fe\] ratio we have for DLAs since Oxygen is produced by TypeII SNe and is essentially undepleted in the interstellar medium. Unfortunately, it is difficult to obtain the oxygen abundance with the required accuracy. In the present study of the z=3.3901 DLA system towards QSO 0000–2620 we present an accurate oxygen abundance determination based on unsaturated lines, together with a Zn detection at the highest redshift ever obtained, and discuss the \[O/Zn\] ratio in connection with the chemical evolution of the associated galaxy.
## 2 Observations and data reduction
QSO 0000–2620 is a bright QSO (V=17.5) with z<sub>em</sub>=4.108 discovered by C. Hazard, which shows a damped Ly$`\alpha `$ system at z<sub>abs</sub> = 3.3901. This is one of the few damped systems at redshift greater than 3 known so far and is well suited for a detailed study of chemical abundances in primordial galaxies. Measurements of abundances of this system have been obtained by Savaglio et al. (1994) and Molaro et al. (1996) from EMMI observations at the ESO NTT, and by Lu et al. (1996) and Prochaska $`\&`$ Wolfe (1999) from HIRES observations at KeckI.
A few spectra of QSO 0000–2620 were obtained as test observations during the first commissioning of the Ultraviolet-Visual Echelle Spectrograph (UVES) on the Nasmyth focus of the ESO 8.2m KUEYEN telescope at Paranal, Chile, in October 1999, which were released for public use. In this paper we have used the data of higher quality: two exposures of 4000 sec and one of 4500 sec secured on October 13, when the seeing, as given by the telescope guide probe, was between 0.35 and 0.5 arcsec FWHM. The observations were all obtained in the standard dichroic $`\mathrm{\#}`$2 mode which includes the spectral region from 3700 to 5000 Å in the blue arm and from 6700 Å to 10500 Å in the red arm. The CCD in the blue arm is a 2Kx4K, 15 $`\mathrm{\mu m}`$ pixel size thinned, anti-reflection coated EEV CCD-44, while in the red the detector is a mosaic of an EEV CCD-44 of the same type of the blue arm and a MIT/LL CCID-20. The average pixel scale in the direction of the dispersion is 0.22 and 0.16 arcsec/pixel for the blue and red arm, respectively. The slit width was set at 0.9 arcsec. The CCDs were read-out in 2x2 pixel binned mode, which gives a r.o.n. of 2.1 e<sup>-</sup> rms for the blue arm CCD, and 2.1 and 3.4 e<sup>-</sup> rms for the two red arm CCDs. More details about the instrument can be found in Dekker, D’Odorico and Kaufer (2000). The full width at half maximum of the instrumental profile, $`\mathrm{\Delta }\lambda _{\mathrm{instr}}`$, was measured from the emission lines of the Thorium-Argon lamp. The resulting resolving power is R= $`\lambda `$/$`\mathrm{\Delta }\lambda _{instr}`$ $``$ 48000, which corresponds to a velocity resolution of $``$ 6 km s<sup>-1</sup>.
The data reduction was performed using the ECHELLE context routines implemented in the ESO MIDAS package. Flat-fielding, cosmic ray removal, sky subtraction, and wavelength calibration were performed on each spectrum separately. Typical r.m.s of the wavelength calibrations are $``$0.6 km s<sup>-1</sup>.
The observed wavelength scale of each spectrum was then transformed into vacuum, heliocentric wavelength scale. The spectra were then added together by using their S/N as weights. Finally, the local continuum was determined in the average spectrum by using a spline to smoothly connect the regions free from absorption features. The continuum for the Ly$`\alpha `$ forest region was fitted by using the small regions deemed to be free of absorptions and by interpolating between these regions with a spline.
With respect to the previous ESO NTT observations, UVES spectra have a higher extension into the blue as well as a more than double resolution and higher S/N. With respect to the Keck data, which cover the range 5100–7660 Å, UVES spectra extend further towards both the blue and the red sides. In the region of overlap, which extends from $``$ 6700 to 7660 Å the data have similar resolution and comparable S/N as can be seen from a comparison of our figures with Fig.1 of Prochaska $`\&`$ Wolfe (1999).
## 3 Measurements of the Column Densities
In this work we restricted our analysis to the metal lines falling in the Ly $`\alpha `$ forest which have never been observed at both comparable resolution and S/N, and to lines in the reddest part of the spectrum. Measurements are presented in Table 1. Column densities have been obtained by fitting theoretical Voigt profiles to the observed absorption lines via $`\chi ^2`$ minimization. The fit was performed using the FITLYMAN package within MIDAS (Fontana & Ballester 1995). During the fitting procedure the theoretical profiles were convolved with the instrumental point spread function modeled from the analysis of the emission lines of the arcs. Portions of the profiles recognized as contaminated by intervening Ly $`\alpha `$ clouds were excluded from the analysis.
The FITLYMAN routines determine the redshift, the column density, and the broadening parameter ($`b`$-value)<sup>3</sup><sup>3</sup>3The broadening parameter is defined as $`b=2^{1/2}\sigma _v`$, where $`\sigma _v`$ is the one-dimensional gaussian velocity dispersion of ions along the line of sight. of the absorption components, as well as the fit errors for each quantity. The atomic data were taken from the compilation of Morton (1991) with the revisions by Cardelli & Savage (1995), Bergeson & Lawler (1993a,b) and Fedchak and Lawler (1999), as specified in Table 1.
Since Ly$`\alpha `$ is not included in our range, when calculating the abundances the hydrogen column density is taken as $`\mathrm{log}N(HI)`$=21.41$`\pm `$0.08 from Lu et al. (1996), after having verified the value consistent with the Ly$`\beta `$ absorption. With such a large column density, the ion species are from the dominant ionization stages in HI gas and do not require ionization corrections (Viegas 1995). The damped structure is known to include minor additional components on both sides of the main absorption (Molaro et al. 1996, Lu et al. 1996), but our analysis is restricted to the unsaturated lines were these weak components do not appear and do not affect the abundances.
The derived column densities and $`b`$-values, together with the redshifts, are reported in Table 1. The difficulty in determining accurate OI column densities stems from the fact that the only transition available redwards the Ly$`\alpha `$ emission is OI $`\lambda `$ 1302Å which has always been found heavily saturated in DLAs. On the other hand OI $`\lambda `$ 1356 Å has never been detected and has provided only loose upper limits. Upon searching the Ly$`\alpha `$ forest we detected several oxygen lines, namely OI $`\lambda `$925.446, $`\lambda `$929.517, $`\lambda `$ 936.629, $`\lambda `$ 948.685, $`\lambda `$950.884, $`\lambda `$988.773 , $`\lambda `$ 1039.230 Å which are relatively free from contamination caused by hydrogen clouds. The OI $`\lambda `$ 925 Å and OI $`\lambda `$950 Å lines have a very low f value (3.5$`\times ^4`$ and 15.7$`\times ^4`$), which is 140 and 25 times lower than the OI $`\lambda `$ 1302 Å,̃ respectively. The lines are not saturated thus allowing accurate determination of the oxygen abundance. The two oxygen lines are shown in Fig. 1. Owing to the presence of minor contaminants we analyzed the two lines independently. The two lines yield similar column densities, where that obtained from OI $`\lambda `$ 925 Å is slightly higher. We notice that the OI $`\lambda `$ 925 Å is somewhat broader ($`b`$=14 km s<sup>-1</sup>) than the other metal lines, thus suggesting the presence of some blend which may explain the small difference in the column densities. In fitting OI $`\lambda `$ 950 the presence of the wing of the Ly$`\delta `$ line on the blue side is accounted for. The adopted abundance for OI is taken as the straight average of the two measurements, i.e. $`\mathrm{log}`$N(OI)=16.42 (or \[O/H\]=–1.85$`\pm `$0.1) with an error of $`\pm `$0.1 which accounts for the uncertainty in the continuum placement. The other $`\alpha `$-element, Silicon, is obtained from the unsaturated SiII $`\lambda `$ 1808 Å line also shown in Fig. 1. We obtained a column density of $`\mathrm{log}`$ N(SiII)=15.06, which is in good agreement with the value of 15.086 derived by Prochaska & Wolfe (1999).
Several N transitions belonging to the NI $`\lambda `$952, $`\lambda `$953, $`\lambda `$956, $`\lambda `$1134 Å multiplets have been detected for the first time in a DLA. In each multiplet, however, only few lines resulted relatively free from contaminations, namely NI $`\lambda `$ 952.415, $`\lambda `$ 953.415, $`\lambda `$ 953.655, NI $`\lambda `$ 963.990, and $`\lambda `$ 965.041 Å . In deriving the nitrogen abundance we further restricted the analysis to those lines which appear totally free from contaminations, which comprise the three lines of the NI $`\lambda `$ 1134 Å multiplet and the NI $`\lambda `$ 953.6 Å transition and are shown in Fig 2. The N column density is found $`\mathrm{log}`$ (NI) =14.70, which is in excellent agreement with the $`\mathrm{log}`$ N(NI)= 14.68 $`\pm `$ 0.14 derived by Molaro et al. (1996) from the analysis of the NI $`\lambda `$ 1134 Å multiplet.
The iron abundance is obtained from the unsaturated FeI $`\lambda `$ 1611 Å line. The abundance is determined at \[Fe/H\]=–2.04, which is 0.3 dex higher than the value obtained by Molaro et al. (1996), and 0.2 dex lower than that obtained by Prochaska and Wolfe (1999) by means of the same transition. The FeII line together with the other iron-peak elements are shown in Fig. 3.
Chromium is clearly detected through the CrII $`\lambda `$ 2056 Å and $`\lambda `$ 2062 Å transitions, while the CrII $`\lambda `$ 2066 Å line is obscured by a strong sky emission line. The two lines yield almost identical Cr column densities with a mean value of $`\mathrm{log}`$N(CrII)=13.10.
Zinc abundance is obtained from the stronger component of the doublet Zn II $`\lambda `$ 2026.1360 Å while the weaker ZnII $`\lambda `$ 2062.664 Å remains below the detection threshold. The value found here is \[Zn/H\]=–2.07 which is consistent with the upper limits set at \[Zn/H\]$``$–1.76 (3$`\sigma `$) by Pettini et al. (1995).
We have clearly detected the NiII $`\lambda `$ 1709 Å while the NiII $`\lambda `$ 1751 Å showed to be partially blended with an O<sub>2</sub> atmospheric line. The abundances of NiII $`\lambda `$ 1709 Å yields \[Ni/H\]=–2.27, which is also consistent with the value of NiII $`\lambda `$ 1751 Å when the telluric contamination is accounted for. We note that the abundance of Ni, despite the recent upwards revision in the $`f`$ values by Fedchak and Lawler (1999), remains slightly lower than the other iron-peak elements.
Other transitions present in our spectrum such as AlII and CIV are either saturated or are not significantly different from the values published in the quoted references. In this paper they are not examined again.
## 4 Discussion
### 4.1 Zinc Abundance and Dust Content
Zinc abundances have been used to trace back the cosmological chemical evolution of Damped Ly$`\alpha `$ galaxies. Instead of the expected growth of metallicity with cosmic time, which is a natural outcome of stellar evolution, Pettini et al. (1999) could not find any trend of the \[Zn/H\] ratio with z<sub>abs</sub>. The DLA towards QSO 0000–2620 is the only one at z$``$ 3 for which a Zn abundance measurement is available at present. The zinc abundance is found at \[Zn/H\]=–2.07, which is the lowest level found so far in any DLA at any redshift. The low metallicity level shows that this Ly$`\alpha `$ galaxy is indeed in the early stages of its chemical evolution. When compared with extant data, this Zn abundance suggests that a mild chemical evolution is present in the DLA population. Observational biases may affect this conclusion, as discussed in a separate paper (Vladilo et al. 2000).
A second important point is that zinc abundance at \[Zn/H\]=–2.07 is fully consistent with both iron and chromium abundances, which are \[Fe/H\]=–2.04 and \[Cr/H\]=–2.0 respectively. It is well established that in the DLAs Zn is consistently more abundant than Cr and Fe (Pettini et al. 1994, 1997), while in Galactic halo stars both Cr and Zn track the Fe content down to very low metallicities. The different behavior of Fe and Zn in the DLAs is currently interpreted as the result of differential depletion of these elements from the gas phase onto dust grains in analogy with the nearby interstellar medium. The fact that in this DLA the abundances of the refractory elements Fe and Cr are found close to the non-refractory element Zn suggests that dust grains have not yet formed in this protogalaxy. Such a circumstance is extremely rare among DLAs, and even in the cases with low depletion considered by Pettini et al. (2000), a difference between Zn and the other refractory elements Fe, Cr and Ni has been always observed. The absence of dust in this damped system frees the measurements of relative abundances from the uncertainties deriving from the unknown fraction of elements which condensate onto grains, thus making this DLA of particular value in the study of chemical composition patterns.
### 4.2 Abundances of $`\alpha `$-capture versus iron-peak elements
Information on $`\alpha `$-capture elements is important because in the early stages of the chemical evolution of galaxies the abundances are likely dominated by Type II SNe products which are richer in $`\alpha `$-elements than Type I SNe which entered into the game only later on. Thus in the early stages the $`\alpha `$-capture elements are expected to be relatively more abundant than the iron-peak elements, and the ratio should reverse later on during evolution. The exact timing of the turn-over depends on both the SFR and IMF, which differ for the different types of galaxies. The \[$`\alpha `$/Fe\] ratio is therefore an indicator of the chemical evolution history which can be used to understand the nature of galaxies when data on their morphology and colors are lacking. Silicon is the most accessible $`\alpha `$ element in DLAs and the compilation of Savaglio et al. (2000), who consider 37 measures, obtain a mean value of $`<`$\[Si/Fe\]$`>`$=0.43$`\pm `$0.18. The problem here is that we observe a similar ratio (\[Si/Fe\]=+0.36) in the Warm Phase of the Galactic interstellar medium, due to the fact that the gaseous silicon and iron abundances are lowered by different amounts by dust depletion (Savage & Sembach 1996). After accounting for dust, Savaglio et al. (2000) conclude that Si is intrinsically enhanced in comparison to iron only in about 25% of the DLAs with SiII determined. Outram et al. (1999) in two DLA systems inferred the absence of dust from the similarity of Si and S abundances and argued for an intrinsic overabundance of \[$`\alpha `$/Fe\] of about 0.5 dex. However, as noted in Matteucci et al. (1997), Si might be produced in Type I SNe, which could explain the complicate observational pattern. Therefore, Si should be considered an indicator of $`\alpha `$ element abundances less reliable than O or S.
Evidence that relative abundances in some DLAs do not conform to those of the galactic halo stars has been provided by the \[S/Zn\] ratio (Molaro et al. 1996, 1998, Centurión et al. 2000) in a sample of 6 DLAs but with one case, the DLA at z<sub>abs</sub>=2.476 towards Q0841+129, consistent with \[S/Zn\]$`>`$0.2.
The oxygen abundance we derive here is the first one obtained by means of unsaturated lines. Since the transitions lie in the forest, the uncertainty is dominated by possible Ly$`\alpha `$ interlopers. However, if these hydrogen lines do actually contaminate the spectrum at the wavelength of the oxygen lines, the real oxygen abundance should be even lower than the value found here.
Oxygen is a typical product of type II SNe and in the atmospheres of Galactic metal-poor stars it shows the notorious enhancement of \[O/Fe\] $``$ +0.5. Claims for an even more extreme oxygen overabundance with \[O/Fe\] reaching +1.0 dex at \[Fe/H\]=–3.0 have been made recently by Israelian et al. (1998) and Boesgaard et al. (1999).
Oxygen has a rather small condensation temperature (T<sub>c</sub> $``$ 180 K) and its depletion factor is of 0.00$`{}_{0.31}{}^{}{}_{}{}^{+0.18}`$ dex in the warm gas phase towards $`\zeta `$ Ophiuchi (Savage & Sembach 1996). Therefore, the \[O/Zn\] ratio is an excellent dust-free diagnostic tool for the \[$`\alpha `$/iron-peak\] abundance ratio.
From the average of the two transitions we derive \[O/H\]=–1.86, which implies \[O/Zn\]=0.21, a relatively modest, if any, $`\alpha `$ enhancement. Si abundance is almost the same as oxygen and we obtain \[Si/H\]=–1.91, which corresponds to \[Si/Zn\]=0.16. We note also here that Lu et al. (1996) tentatively identified the SII 1253.811 Å line associated with the damped, and derived a column density $`\mathrm{log}`$ N(SII) = 14.70 $`\pm `$ 0.03, i.e. an abundance of \[S/H\]=–1.91, which implies \[S/Zn\]=0.16. Thus, the S abundance of Lu et al. (1996) is also consistent with our O and Si abundances in the system.
For the first time in the DLA studied here we have at our disposal a set of 3 indicators for the $`\alpha `$ element abundances and three indicators for the iron peak elemental abundances, which do not need any dust correction. The \[O,Si,S/Zn,Fe,Cr\] ratios, in whatever combination the elements are chosen, are in the range 0.09–0.22 dex. This range of values is significantly lower than analogous ratios in Galactic stars with comparable metallicities, and the difference is more stringent in this case for the low value of the metallicity, which is one order of magnitude lower than the average value in DLAs. The lack of significant \[$`\alpha `$/Fe\] enhancement corroborates previous results on this system as well as the suggestion that, on the basis of their relative elements, at least some DLA galaxies seem to undergo a chemical evolution which differs from that of the Milky–Way. Centurión et al. (2000) showed that the \[S/Zn\] ratios in DLAs are consistent with a decreasing trend with increasing metallicity. These ratios show also a correlation, although less clear-cut, with cosmic time, where the lower \[S/Zn\] values occur mostly at low redshifts. The \[S/Zn\]=0.16 and \[O/Zn\]=0.19 for our DLA add a new point which supports the presence of such trends. Considering the extant data of \[S/Zn\] and the value \[O/Zn\]=0.20 we derived here, the regression analysis yields the correlation \[S,O/Zn\] = –0.36($`\pm `$0.11)\[Zn/H\] – 0.52. The $`\alpha `$/iron-peak ratio becomes solar at \[Fe/H\] $``$-1.4 and further decreases at higher metallicities.
Chemical evolutionary models predict solar \[$`\alpha `$/Fe\] at low metallicities, when star formation proceeds in bursts separated by relatively long quiescent periods or whenever star formation is slower than that of our Galaxy. These conditions are met in dwarf galaxies, LSB galaxies and in the outer regions of disks (Matteucci et al. 1997, Jiménez et al. 1999). In these galaxies the metal enrichment is so slow that Type Ia supernovae have enough time to evolve and enrich the gas with iron-peak elements, in such a way that reduced $`\alpha `$ over iron ratios are attained at low metallicity.
In the Milky Way few metal poor stars with solar \[$`\alpha `$/Fe-peak\] ratios have been found (Carney et al. 1997, Nissen & Schuster 1997). These exceptional stars show large apogalactic distances and it was suggested that they belonged to a satellite galaxy which experienced a different chemical evolution history than the Milky Way, and were subsequently accreted by the latter.
Incidentally, we note that the interpretation of the low \[O/Fe\] ratio at low metallicities as an intrinsic effect of TypeI SNe does not support the suggestion by Umeda et al. (1999) that Type I SNe may be suppressed at metallicities lower that \[Fe/H\] $`<`$–1.0.
### 4.3 The Nitrogen abundance
In our system the nitrogen abundances relative to Zn and O are \[N/Zn\]=–0.54 and \[N/O\]=–0.74. Both values are somewhat lower than previously reported by Molaro et al. (1996) because of the upwards revision in oxygen and iron abundances. The \[N/O\] is not so extreme as previously thought, but the N abundance still requires a primary production of the element according to the evolutionary models discussed in Matteucci et al (1997). N abundance in this DLA, but also in other DLAs, remain lower than the corresponding values in the Milky Way, where N is found to trace Fe in lockstep (Israelian et al. 2000), which again suggests a different pattern than the Milky Way. Nitrogen abundances have been discussed by Lu, Sargent, & Barlow (1998) and by Centurión et al. (1998), who pointed out the presence of a real scatter among DLA systems with some values very close to a pure secondary behavior and other values which require a primary nucleosynthesis. The different values of N/O observed at a given O/H may be understood in terms of the delayed delivery of primary N with respect to O when star formation proceeds in bursts. During the quiescence period N is deposited in the ISM and the N/O ratio increases while O/H remains constant. Therefore a relatively high \[N/O\], such as that observed here, nicely fits with the picture of a delayed N release after a quiescence period, which is also required by the low $`\alpha `$ over iron ratios.
### 4.4 Conclusions
UVES observations of the QSO 0000-2620 resulted in the detection of the ZnII 2026.136 Å transition originated in the damped system at z<sub>abs</sub>=3.3901. The abundance derived is \[Zn/H\]=-2.06 which is presently the lowest among DLAs. This low metallicity level shows that the galaxy is in the early stages of its chemical evolution. When compared with the larger abundances observed in DLA galaxies at lower redshift, this measurement provides a first hint of a cosmological chemical evolution in which abundances increase with decreasing redshift (see Vladilo et al 2000)
The abundances of Cr and Fe were also obtained and found similar to the abundance of Zn. This coincidence in the abundances between refractory and non refractory elements is interpreted as dust has not formed yet.
The OI 925 Å and 950 Å lines were detected although they fall within the Lyman$`\alpha `$ forest. These lines have oscillator strengths much smaller than that of the OI 1302 Å line which is strongly saturated. These detections allow, for the first time, a rather accurate measurement of the oxygen abundance in a DLA system. The oxygen abundance, \[O/H\]=-1.85, is remarkably similar to the sulphur abundance derived by Lu et al (1996), \[S/H\]=-1.91, and both abundances are similar to that of silicon, \[Si/H\]=-1.91, which is consistent with the absence of dust in the system as implied by Zn and Fe abundances.
In this DLA the relative abundances of $`\alpha `$ and iron-peak elements are mildly overabundant compared to solar and do not show the enhancement expected for a progenitor of an early type spiral such as our own Galaxy. The paradigm of DLAs as progenitors of the Milky-Way type spirals is not supported by the study of the dust-independent chemical evolution indicators such as \[O/Zn\] and \[S/Zn\]. The almost solar abundance ratios measured at the low metallicity of \[Fe/H\]$``$=–2 suggest that this DLA should be associated with objects with low, or episodic, star formation rates such as LSB or dwarf galaxies. However, these results are not characteristic of this DLA alone and similar low \[S/Zn\] ratios are also found in other DLAs (Centurión et al 2000).
## 5 Acknowledgements
The new high resolution spectra analyzed in this paper are of unique quality and were obtained during the first nights of commissioning of a new instrument at a new telescope. For these results we are indebted to all ESO staff involved in the VLT construction and in UVES commissioning.
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# 1 SM Estimates of 𝐷⁰-𝐷̄⁰ Oscillations
## 1 SM Estimates of $`D^0\overline{D}^0`$ Oscillations
### 1.1 General Features
Oscillations in general are described by two dimensionless quantities:
$$x=\frac{\mathrm{\Delta }M}{\overline{\mathrm{\Gamma }}},y=\frac{\mathrm{\Delta }\mathrm{\Gamma }}{2\overline{\mathrm{\Gamma }}}$$
(1)
where
$$\mathrm{\Delta }MM_2M_1,\mathrm{\Delta }\mathrm{\Gamma }\mathrm{\Gamma }_1\mathrm{\Gamma }_2,\overline{\mathrm{\Gamma }}=\frac{1}{2}\left(\mathrm{\Gamma }_1+\mathrm{\Gamma }_2\right).$$
(2)
The leading contributions to $`\mathrm{\Delta }M_K`$ as well as $`\mathrm{\Delta }M_B`$ are obtained from the well-known quark box diagram. The fields in the internal loop – $`c`$ in the former and $`t`$ quarks in the latter case besides $`W`$ bosons – are much heavier than the external quark fields. Those heavy degrees of freedom can be integrated out leading to $`\mathrm{\Delta }M_{K,B}`$ being described by the expectation value of a local operator.<sup>1</sup><sup>1</sup>1It turns out that $`\mathrm{\Delta }M_K`$ is not completely dominated by local contributions reflecting dynamics operating around the scale $`m_c`$: a significant fraction is due to long distance dynamics characterized by low scales $`\mathrm{\Lambda }_{\mathrm{QCD}}`$. This would change dramatically if $`m_c`$ were larger. This has been a highly successful ansatz: within the present theoretical uncertainties the SM can reproduce the observed values of $`\mathrm{\Delta }M_K`$ and $`\mathrm{\Delta }M_{B_d}`$ without forcing any parameter.
While the size of $`\mathrm{\Delta }\mathrm{\Gamma }_K`$ is naturally understood as due to $`K3\pi `$ being severely phase space restricted, it cannot be inferred from such a local operator, as described above. For neutral $`B`$ mesons, on the other hand $`\mathrm{\Delta }\mathrm{\Gamma }_B`$ can be estimated through short distance dynamics .
On very general grounds one expects $`D^0\overline{D}^0`$ oscillations to be quite slow within the SM, since two structural reasons combine to make $`x_D`$ and $`y_D`$ small:
* While the bulk of charm decays is Cabibbo allowed, the amplitude for $`D^0\overline{D}^0`$ transitions is necessarily twice Cabibbo suppressed – as is therefore the ratio between oscillation and decay rate: $`\mathrm{\Delta }M_D/\mathrm{\Gamma }_D,\mathrm{\Delta }\mathrm{\Gamma }_D/\mathrm{\Gamma }_D𝒪(\mathrm{sin}^2\theta _C)`$. The amplitudes for $`K^0\overline{K}^0`$ and $`B_d\overline{B}_d`$ are Cabibbo and KM suppressed, respectively – yet so are their decay widths allowing the oscillation and decay rate to be quite comparable.
* Due to the GIM mechanism one has $`\mathrm{\Delta }M=\mathrm{\Delta }\mathrm{\Gamma }=0`$ in the limit of flavor symmetry. Yet flavor symmetry breaking driving $`K^0\overline{K}^0`$ is characterized by $`m_c^2m_u^2`$ and therefore no real suppression arises. On the other hand $`SU(3)`$ breaking controlling $`D^0\overline{D}^0`$ is typified by $`m_s^2m_d^2`$ (or, in terms of hadrons, $`M_K^2M_\pi ^2`$) as compared to the scale $`M_D^2`$; it provides a significant reduction.
Having two Cabibbo suppressed classes of decays one concludes on these very general grounds:<sup>2</sup><sup>2</sup>2One can argue that because of the $`\mathrm{\Delta }C=\mathrm{\Delta }Q`$ rule in semileptonic charm decays one should write the nonleptonic rather than the total width in Eq. (3); yet this difference is small for $`D^0`$ and certainly is in the theoretical ‘noise’.
$$\mathrm{\Delta }M_D,\mathrm{\Delta }\mathrm{\Gamma }_DSU(3)\mathrm{breaking}\times 2\mathrm{sin}^2\theta _C\times \mathrm{\Gamma }_D$$
(3)
The proper description of $`SU(3)`$ breaking thus becomes the central issue.
With typical nonleptonic $`D`$ decay channels exhibiting sizeable $`SU(3)`$ breaking – see our discussion in Sect. 2 – a priori one cannot count on this suppression to amount to more than a factor of about two or three in the $`D`$ width difference. There are reasons to believe that a larger reduction may occur for the mass difference $`\mathrm{\Delta }M_D`$ driven by virtual intermediate states. Yet an order of magnitude reduction, in particular in $`\mathrm{\Delta }\mathrm{\Gamma }_D`$ would seem unjustifiably pessimistic. Thus
$$\frac{\mathrm{\Delta }M_D}{\mathrm{\Gamma }_D}\stackrel{<}{}\frac{\mathrm{\Delta }\mathrm{\Gamma }_D}{\mathrm{\Gamma }_D}\stackrel{<}{}\frac{1}{3}\times 2\mathrm{sin}^2\theta _C\mathrm{few}\times 0.01$$
(4)
represents a conservative bound for overall mixing based on very general features of the SM; for the mass difference this estimate can actually be seen on the cautious side.
The following line of arguments is usually employed: (i) Quark-level contributions are estimated by the usual quark box diagrams and yield only insignificant contributions to $`\mathrm{\Delta }M_D`$ and $`\mathrm{\Delta }\mathrm{\Gamma }_D`$ (see below). (ii) Various schemes employing contributions of selected hadronic states are invoked to estimate the impact of long distance dynamics; the numbers typically resulting are $`x_D,y_D\mathrm{\hspace{0.17em}10}^410^3`$ . (iii) These findings lead to the following widely embraced conclusions: An observation of $`x_D>10^3`$ would reveal the intervention of New Physics beyond the SM, while $`y_Dy_D|_{SM}10^3`$ has to hold since New Physics has hardly a chance to contribute to it.
Beyond the general property that both $`\mathrm{\Delta }M_D`$ and $`\mathrm{\Delta }\mathrm{\Gamma }_D`$ have to vanish in the $`SU(3)`$ limit, the dynamics underlying them have different features: $`\mathrm{\Delta }M_D`$ receives contributions from virtual intermediate states whereas $`\mathrm{\Delta }\mathrm{\Gamma }`$ is generated by on-shell transitions. Therefore the former is usually considered to represent a more robust quantity than the latter; actually it has often been argued that quark diagrams cannot be relied upon to even estimate $`\mathrm{\Delta }\mathrm{\Gamma }`$. A folklore has arisen that theoretical evaluations of the two quantities rest on radically different grounds.
Yet we note that despite these differences there is no fundamental distinction in the theoretical treatment of $`\mathrm{\Delta }M_D`$ and $`\mathrm{\Delta }\mathrm{\Gamma }_D`$: both can be described through an operator product expansion, and its application relies on local quark-hadron duality for both $`\mathrm{\Delta }M_D`$ and $`\mathrm{\Delta }\mathrm{\Gamma }_D`$. Only the numerical aspects differ, as does the sensitivity to New Physics.
### 1.2 Operator Product Expansion
Following the general treatment of inclusive weak transitions, see Refs. , we can describe $`D^0\overline{D}^0`$ oscillations by considering a correlator
$$\widehat{T}_{D\overline{D}}(\omega )=\frac{1}{2}\mathrm{d}^4x\mathrm{e}^{i\omega t}iT\{_W(x)_W(0)\},T_{D\overline{D}}(\omega )=\frac{1}{2M_D}\overline{D}|\widehat{T}_{D\overline{D}}(\omega )|D$$
(5)
as a function of a complex variable $`\omega `$. Here $`_W`$ is the $`\mathrm{\Delta }C=1`$ Hamiltonian density. With the mixing amplitude of interest
$$A(\omega )=2\underset{n}{}\frac{1}{2M_D}\frac{\overline{D}|_W|n(\stackrel{}{k}=0)n(\stackrel{}{k}=0)|_W|D}{E_nM_D+\omega +iϵ}\mathrm{\Delta }\stackrel{~}{M}_D(\omega )+\frac{i}{2}\mathrm{\Delta }\stackrel{~}{\mathrm{\Gamma }}_D(\omega )$$
(6)
one has $`4T_{D\overline{D}}(\omega )=A(\omega )+A(\omega )`$, whereas $`\mathrm{\Delta }M_D=\mathrm{\Delta }\stackrel{~}{M}_D(0)`$, $`\mathrm{\Delta }\mathrm{\Gamma }_D=\mathrm{\Delta }\stackrel{~}{\mathrm{\Gamma }}_D(0)`$. The summation above runs over all intermediate states $`|n`$ with energies $`E_n`$ and vanishing spacelike momentum. $`\mathrm{\Delta }M_D`$ can be expressed by a dispersive integral over $`\mathrm{\Delta }\stackrel{~}{\mathrm{\Gamma }}_D(\omega )`$ evaluated through the principal value prescription:
$$\mathrm{\Delta }M_D=\frac{1}{2\pi }\mathrm{V}.\mathrm{P}.d\omega \frac{\mathrm{\Delta }\stackrel{~}{\mathrm{\Gamma }}_D(\omega )}{\omega }.$$
(7)
Applying the operator product expansion (OPE) to Eq. (5) provides us with a consistent evaluation of the transition rates through an expansion in powers of $`1/m_c`$. With the charm quark mass exceeding the typical scale of strong interactions $`\mu _{\mathrm{hadr}}`$ by a modest amount only, one cannot count on obtaining a reliable quantitative description in this way; yet it still yields a useful classification of various effects. This is briefly reviewed below.
The leading term for $`\mathrm{\Delta }C=2`$ transitions comes from dimension-6 four-fermion operators of the generic form $`(\overline{u}c)(\overline{u}c)`$ with the corresponding Wilson coefficient receiving contributions from different sources.
(a) Effects due to intermediate $`b`$ quarks are most simply calculated since they are highly virtual and Euclidean:
$$\mathrm{\Delta }M_D^{(b\overline{b})}\frac{G_F^2m_b^2}{8\pi ^2}\left|V_{cb}^{}V_{ub}\right|^2\frac{1}{2M_D}\overline{D}^0|(\overline{u}\gamma _\mu (1\gamma _5)c)(\overline{u}\gamma _\mu (1\gamma _5)c)|D^0;$$
(8)
however they are highly suppressed by the tiny KM mixing with the third generation. Using factorization to estimate the matrix element one finds:
$$x_D^{(b\overline{b})}\mathrm{few}\times 10^7.$$
(9)
Loops with one $`b`$ and one light quark likewise are suppressed.
(b) For the light intermediate quarks the momentum scale is set by the external mass $`m_c`$, and the corresponding factor is given by $`G_F^2m_c^2/8\pi ^2\mathrm{sin}^2\theta _C\mathrm{cos}^2\theta _C`$ (from now on we will often omit the KM factors when they are obvious). However, it is highly suppressed by the GIM factor $`\left(\frac{m_s^2m_d^2}{m_c^2}\right)^2`$ leading to <sup>3</sup><sup>3</sup>3 This contribution is obviously saturated at the momentum scale $`m_c`$, and thus refers to the Wilson coefficient of the $`D=6`$ operator. We disagree with statements that these are long-distance contribution simply because they are proportional to $`m_s^2`$.
$$\widehat{T}_{D\overline{D}}(\omega )=\frac{G_F^2}{16\pi ^2}|V_{cs}^{}V_{us}|^2(\frac{\left(m_s^2m_d^2\right)^2}{(m_c\omega )^2}+\frac{\left(m_s^2m_d^2\right)^2}{(m_c+\omega )^2})\times $$
$$\left[(\overline{u}\gamma _\mu (1\gamma _5)c)(\overline{u}\gamma _\mu (1\gamma _5)c)+2(\overline{u}(1+\gamma _5)c)(\overline{u}(1+\gamma _5)c)\right].$$
(10)
Hence we read off for its contribution to the mass difference
$$\mathrm{\Delta }M_D^{(s,d)}\frac{G_F^2m_c^2}{4\pi ^2}\left|V_{cs}^{}V_{us}\right|^2\frac{\left(m_s^2m_d^2\right)^2}{m_c^4}$$
$$\times \frac{1}{2M_D}\overline{D}^0|(\overline{u}\gamma _\mu (1\gamma _5)c)(\overline{u}\gamma _\mu (1\gamma _5)c)+2(\overline{u}(1+\gamma _5)c)(\overline{u}(1+\gamma _5)c)|D^0.$$
(11)
(This expression differs from what is usually quoted in the literature, see, e.g. Ref. .) As follows from Eq. (10) bare quark loops do not contribute to $`\mathrm{\Delta }\mathrm{\Gamma }_D`$ at this order. The latter is suppressed by additional powers of $`m_s/m_c`$, or by $`\alpha _s/\pi `$ when gluon corrections are accounted for (e.g., through the anomalous dimension of the light quark mass).
The GIM suppression by two powers of $`m_s/m_c`$ for each quark line is inevitable for left-handed weak vertices. This feature persists for Penguin operators, albeit in a slightly different way. Numerically one finds:
$$\mathrm{\Delta }\mathrm{\Gamma }_D^{\mathrm{box}}<\mathrm{\Delta }M_D^{\mathrm{box}}\mathrm{few}\times 10^{17}\mathrm{GeV}\widehat{=}x_D^{\mathrm{box}}\mathrm{few}\times 10^5$$
(12)
However, since the leading Wilson coefficient is highly suppressed, one has to consider also the contributions from higher dimensional operators. It turns out that the $`SU(3)`$ GIM suppression is in general not as severe as $`(m_s^2m_d^2)/m_c^2`$ per fermion line: it can be merely $`m_s/\mu _{\mathrm{hadr}}`$ if the fermion line is soft. In the so-called practical version of the OPE this is described by condensates contributing to the next terms in the $`1/m_c`$ expansion.
There is a simple rule of thumb: cutting a quark line, we pay the price of a power suppression $`\mu _{\mathrm{hadr}}^3/m_c^3`$; yet the GIM suppression now becomes only $`m_s/\mu _{\mathrm{hadr}}`$. Altogether this yields a factor $`4\pi ^2\mu _{\mathrm{hadr}}^2/(m_sm_c)`$ which can result in an enhancement. In particular, we keep in mind that $`SU(3)`$ breaking effects in condensates are not significantly suppressed, and the ratio between $`\mu _{\mathrm{hadr}}`$ and $`m_c`$ is not much smaller than unity.
An example is given by the diagram in Fig. 1a. It yields a six-fermion operator of the generic form $`(\overline{u}c)(\overline{u}c)(\overline{d}d\overline{s}s)`$ with the Wilson coefficient $`\frac{m_s^2}{m_c^2}G_F^2m_c^1.`$ This contribution thus scales like $`4\pi ^2m_s^3\mu _{\mathrm{hadr}}^3/m_c^5`$ compared to the “standard” factor $`G_F^2f_D^2m_c^2M_D`$. Explicitly, this contribution (neglecting gluon corrections) is given by
$$\mathrm{\Delta }M_D^{(D=9)}=2\mathrm{sin}^2\theta _C\mathrm{cos}^2\theta _C\frac{G_F^2m_s^2}{m_c^3}\times $$
(13)
$$\frac{1}{2M_D}\overline{D}|\overline{u}^i\gamma _\mu (1\gamma _5)c^j\overline{u}^k\gamma _\nu (1\gamma _5)\gamma _0c^i\left(\overline{s}^j\gamma ^\mu \gamma _0\gamma ^\nu (1\gamma _5)s^k\overline{d}^j\gamma ^\mu \gamma _0\gamma ^\nu (1\gamma _5)d^k\right)|D.$$
Here $`i,j,k`$ are color indices. Note, that the $`c`$ quark operators here are normalized at the low momentum scale. Therefore, $`\frac{1+\gamma _0}{2}c(x)`$ describes only annihilation of $`c`$ quark, and $`\frac{1\gamma _0}{2}c(x)`$ only creation of charmed antiquark. This implies, for example, that the relation $`\gamma _0c(x)\times c(y)=c(x)\times \gamma _0c(y)`$ always holds for the considered matrix elements, since only two combinations $`\frac{1\pm \gamma _0}{2}c(x)\times \frac{1\gamma _0}{2}c(y)`$ survive. It is interesting that for the “neutral current” type color flow in the both weak vertices, the contributions of Figs. 1a (proportional to $`a_2^2`$) vanish to the leading order in $`1/m_c`$.
$`SU(3)`$ suppression can be further softened by cutting both fermion lines. To transfer a large momentum one has to add a gluon, like in Fig. 1b (for this reason another loop factor of $`4\pi ^2`$ is replaced by $`4\pi \alpha _s`$). These yield eight-fermion operators with the flavor structure <sup>4</sup><sup>4</sup>4There are actually additional contributions not obtained merely by cutting quark lines; they are reminiscent of the soft part of Penguin contributions, see below. $`(\overline{u}c)(\overline{u}c)\left[(\overline{d}d)(\overline{d}d)+(\overline{s}s)(\overline{s}s)(\overline{d}d)(\overline{s}s)(\overline{s}s)(\overline{d}d)\right].`$ With the $`SU(3)`$ suppression in the matrix element due to double antisymmetrization between $`s`$ and $`d`$ only $`m_s^2/\mu _{\mathrm{hadr}}^2`$, this contribution scales like $`\frac{4\pi ^2m_s^2\mu _{\mathrm{hadr}}^4}{m_c^6}G_F^2f_D^2m_c^2M_D`$.
It is interesting to note that, in principle, the $`SU(3)`$ suppression can be as mild as only the first power of $`m_s`$. Namely, if we schematically define
$$\overline{D}|(\overline{u}c)(\overline{u}c)[(\overline{d}d)(\overline{d}d)+(\overline{s}s)(\overline{s}s)(\overline{d}d)(\overline{s}s)(\overline{s}s)(\overline{d}d)]|D\zeta \overline{D}|(\overline{u}c)(\overline{u}c)(\overline{d}d)(\overline{d}d)|D$$
then, at small $`m_s`$, the $`SU(3)`$ suppression factor $`\zeta `$ can scale as $`m_s/\mu _{\mathrm{hadr}}`$. Indeed, if the matrix element, as a function of the two quark masses is given by
$$\overline{D}|(\overline{u}c)(\overline{u}c)(\overline{q}_1q_1)(\overline{q}_2q_2)|D=A+B(m_1+m_2)\mathrm{ln}\frac{\mu _{\mathrm{hadr}}}{m_1+m_2}$$
then $`\zeta 2Bm_s\mathrm{ln}2`$. More accurately, for the actual operator the leading, linear in $`m_s`$ term in the “soft GIM” factor $`\zeta `$ is determined by the matrix element
$$\frac{M_K^2}{16\pi ^2f_\pi ^2}<\overline{D}^0|(\overline{u}_L\mathrm{\Gamma }c_L)^2\overline{d}_L\mathrm{\Gamma }s_L\overline{s}_L\mathrm{\Gamma }d_L|D^0>$$
(14)
in the limit of massless $`s`$ and $`d`$ and is due to the chiral $`\overline{K}^0\eta (\pi ^0)K^0`$ loop. The explicit structure of the Lorentz and color matrices $`\mathrm{\Gamma }`$ above follows from the operator given below in Eq. (15).
Let us sketch this step. At small $`m_s`$ (for simplicity we put $`m_d=0`$, but $`m_u`$ can be arbitrary) it is convenient to use the current quark fields $`s^{}`$, $`d^{}`$ instead of mass eigenstates, and perform an expansion around the symmetry point $`m_s=m_d`$. Then the transition operator has the simple flavor structure $`(\overline{u}c)^2(\overline{s}^{}d^{})^2`$. Its matrix element between $`\overline{D}^0`$ and $`D^0`$ vanishes to zeroth order since the operator has $`\mathrm{\Delta }S^{}=\mathrm{\Delta }D^{}=2`$. The mass perturbation, however, has a $`\mathrm{\Delta }S^{}=\mathrm{\Delta }D^{}=1`$ piece $`\mathrm{sin}\theta _C\mathrm{cos}\theta _Cm_s\overline{d}^{}s^{}`$. The nonvanishing second-order correction to the matrix element $`\overline{D}^0|(\overline{u}c)^2(\overline{s}^{}d^{})^2|D^0`$ is then generally proportional to $`m_s^2`$. The only exception comes from the pseudo-goldstone loop of Fig. 2 shaped by
$$\frac{\mathrm{d}^4k}{(2\pi )^4i}\frac{1}{(M^2k^2)^3}=\frac{1}{32\pi ^2M^2}$$
with $`M^2m_s`$. This contribution is proportional to the zero-momentum amplitude $`\overline{D}^0\overline{K}^0\overline{K}^0|(\overline{u}c)^2(\overline{s}d)^2|D^0`$ which, by PCAC is related to the chiral limit matrix element of the double commutator of the operator with the $`\overline{d}s`$ axial charge. This yields the stated equation.
The above estimate serves only as an existence proof. Most probably, such infrared effects involving pion loops are not the dominant source. Even without such nonanalytic terms the matrix elements typically depend strongly on the quark masses, and in the actual world the double subtraction present in the above operators, can result in only a mild suppression factor in spite of being formally of the order of $`m_s^2`$.
Estimates of the actual size of these contributions at present suffer from considerable uncertainties, primarily in the matrix elements. Direct computation of the bare diagram yields the following cumbersome result:
$$2\widehat{T}(0)=\frac{8\pi \alpha _sG_F^2}{m_c^4}\mathrm{sin}^2\theta _C\mathrm{cos}^2\theta _C𝒢O^{(12)},$$
(15)
$$O^{(12)}=\{\overline{u}\mathrm{\Gamma }^\alpha d\overline{s}\mathrm{\Gamma }_0t^ac\overline{u}\mathrm{\Gamma }_0d\overline{s}\mathrm{\Gamma }^\alpha t^ac+iϵ^{0\mu \alpha \beta }\overline{u}\mathrm{\Gamma }_\beta t^ad\overline{s}\mathrm{\Gamma }_\mu c\}\times $$
$$\left\{\overline{u}\mathrm{\Gamma }_\alpha s\overline{d}\mathrm{\Gamma }_0t^ac\overline{u}\mathrm{\Gamma }_0s\overline{d}\mathrm{\Gamma }_\alpha t^ac+iϵ_{0\nu \alpha \gamma }\overline{u}\mathrm{\Gamma }^\gamma t^as\overline{d}\mathrm{\Gamma }^\nu c\right\}$$
where $`\mathrm{\Gamma }_\mu =\gamma _\mu (1\gamma _5)`$ and $`t^a=\frac{\lambda ^a}{2}`$ are color matrices. The symbol $`𝒢`$ denotes GIM-type subtraction; its action on a generic operator with four $`s`$ and $`d`$ quarks like in $`O^{(12)}`$ is defined as
$$𝒢(\overline{s}Ud)(\overline{d}Vs)=$$
(16)
$$(\overline{s}Ud)(\overline{d}Vs)(\overline{s}Us)(\overline{s}Vs)(\overline{d}Ud)(\overline{d}Vd)+(\overline{d}Us)(\overline{s}Vd)+(\overline{s}Us)(\overline{d}Vd)+(\overline{d}Ud)(\overline{s}Vs)$$
(the last two terms do not have counterparts in the box diagram). The exact coefficients and the color structure in $`O^{(12)}`$ are modified by straightforward renormalization of the weak decay operators, yet the major uncertainty lies in the matrix elements.
In the spirit of the “Educated dimensional analysis” of Ref. we estimate the magnitude of the matrix element of $`𝒢O^{(12)}`$ as
$$\zeta f_D^2M_D(0.3\text{GeV})^6\zeta 710^5\text{GeV}^9$$
(17)
where $`\zeta `$ accounts for the $`SU(3)`$ GIM suppression,
$$\overline{D^0}|𝒢O^{(12)}|D^0=\zeta \overline{D^0}|O^{(12)}|D^0.$$
Note that $`\alpha _s`$ enters at the charm mass scale, and for consistency must be evaluated in the $`V`$\- rather than the $`\overline{\mathrm{MS}}`$-scheme, which is obvious in the BLM approximation. Numerically we end up with
$$\delta ^{(12)}x_D𝒪(10^3).$$
(18)
It is actually conceivable that medium-size instantons yield enhanced contributions to the corresponding matrix element of the eight-fermion operators in question, due to the induced t’ Hooft vertex of the form $`(\overline{u}_Lu_R)(\overline{d}_Ld_R)(\overline{s}_Ls_R)`$. Note that it can yield the effect $`m_s^2`$ which is not formally related to the spontaneous symmetry breaking since would violate only the anomalous singlet $`U_A(1)`$.
Diagrams in Figs. (1) and (2) literally do not produce an absorptive part and, therefore contribute to $`\mathrm{\Delta }M_D`$, but not $`\mathrm{\Delta }\mathrm{\Gamma }_D`$. Yet the latter can be generated, for example, through a cut across the gluon propagator in Fig. 2 if it is dressed; in $`\mathrm{\Delta }M_D`$ this would contribute to the anomalous dimension. This is the leading contribution in the BLM approximation. It amounts to replacing $`\alpha _s`$ by $`\frac{9}{4}\text{ }i\alpha _s^2/(1+81/16\alpha _s^2)`$ in Eq. (15). Such approximation is justified if other contributions to the anomalous dimension can be neglected. At the charm scale these modifications do not seem to lead to a particular numerical suppression of $`\mathrm{\Delta }\mathrm{\Gamma }_D`$ compared to $`\mathrm{\Delta }M_D`$. Therefore, we arrive at
$$x_D,y_D𝒪(10^3)$$
(19)
In summary: We have shown that the high degree of $`SU(3)`$ invariance and related GIM suppression $`(m_s^2m_d^2)^2/m_c^2`$ exhibited by quark box diagrams for $`D^0\overline{D}^0`$ is not typical for the process. Terms in the $`D^0\overline{D}^0`$ amplitude can be proportional to $`m_s^2`$ or even $`m_s^1`$. Such contributions arise naturally in the OPE through condensate contributions containing higher dimensional operators. While those are formally suppressed by powers in the heavy quark mass, this does not constitute a very significant factor for the case of charm. It had been noted long ago (see, for example, Ref. ) that estimates of the absorptive part of the $`D^0\overline{D}^0`$ amplitude are very sensitive to low energy parameters: evaluating it on the quark level one encountered much more effective $`SU(3)`$ cancellations than in its potential hadronic counterparts $`M_K^4/m_c^4`$. Application of the OPE treatment allows to clarify and justify those, rather tentative suggestions identifying the possible source of such enhanced contributions in the framework of the $`1/m_c`$ expansion. Moreover, it becomes clear that the same applies to the mass difference $`\mathrm{\Delta }M_D`$ as well.
Our numerical estimates are rather uncertain. However we note that the OPE naturally accommodates the size of mixing previously discussed in the literature as a possible effect of general long-distance dynamics. At the same time, we think that $`x_D`$, $`y_D`$ exceeding $`510^3`$ cannot be attributed to the OPE contributions in the framework of standard assumptions.
Regardless of the size of the matrix elements involved, we still can state that $`D^0\overline{D}^0`$ mixing must be suppressed in the SM whenever it makes sense to speak of it in the $`1/m_c`$ expansion. The natural yardstick for the unsuppressed level is the overall nonleptonic decay width (up to inherent CKM mixing factors). As illustrated above, smallness of the higher-order terms compared to the formally leading in $`1/m_c`$ effect for mixing cannot serve as a valid universal criterion. Yet a certain suppression of the higher-order terms compared to the (parton) decay width per se is a necessary condition for applying sensibly the $`1/m_c`$ expansion. Since the identified effects with mild GIM cancellation emerge in relatively high order in $`1/m_c`$, the minimal suppression must amount to a noticeable factor, as asserted above.
### 1.3 Quark-Hadron Duality
Beyond the question about the size of the higher order corrections in the $`1/m_c`$ expansion there is the more fundamental one about local quark-hadron duality; i.e., to which degree of accuracy does a quark level result derived from the OPE describe an inclusive quantity involving hadrons? In other words, how well knowledge of $`\mathrm{\Delta }\stackrel{~}{M}(\omega )`$, $`\mathrm{\Delta }\stackrel{~}{\mathrm{\Gamma }}(\omega )`$ in Eq. (6) at large (compared to $`\mathrm{\Lambda }_{\mathrm{QCD}}`$) complex $`\omega `$ based on the short-distance expansion, determines their values at $`\omega =0`$ measured in experiment.
It was pointed out in Ref. that with heavy quarks one often has to deal with a novel aspect, referred to as global duality: a Euclidean dispersive integral will reproduce the contributions coming from all cuts in the Minkowski domain, while some of them are unphysical for the considered decay process. One then has to filter out the contributions to the integral that correspond to the individual process of interest. Since this can be done only to a certain accuracy, it introduces a further source of theoretical uncertainty.
Even though the cuts in $`T_{D\overline{D}}(\omega )`$ describe the same physical channel, this complication still persists in a certain form for flavor oscillation processes: with $`T_{D\overline{D}}(\omega )=A(\omega )+A(\omega )`$, it has two cuts which overlap. However, since one of them starts at $`\omega m_c`$ and another at $`\omega m_c`$, they can be disentangled at $`|\omega |m_c`$ in the $`1/m_c`$ expansion of ‘practical’ OPE in the same way as in the usual heavy quark decay widths . The central issue is again the validity of local duality, both for $`\mathrm{\Delta }\mathrm{\Gamma }`$ and $`\mathrm{\Delta }M`$. Its more dedicated explanation and its implementations can be found in various reviews (see, for example, Ref. ).
That also $`\mathrm{\Delta }M_D`$ is sensitive to duality violations is readily seen by imagining the presence of a narrow resonance of appropriate quantum numbers close to the $`D`$ meson mass: it would significantly affect the size of $`\mathrm{\Delta }M_D`$. Yet, practically, it is natural to expect local duality violations to be smaller in $`\mathrm{\Delta }M_D`$ than in $`\mathrm{\Delta }\mathrm{\Gamma }_D`$; i.e., the onset of duality should occur at a lower scale for the former than for the latter. For $`\mathrm{\Delta }\mathrm{\Gamma }_D`$ is directly given by the discontinuity in the corresponding $`D\overline{D}`$ transition amplitude, whereas $`\mathrm{\Delta }M_D`$ can be represented by the principal value of the dispersion integral, Eq. (7). The latter provides a measure of averaging that reduces the sensitivity to the resonances or thresholds and thus local duality violations. This has been illustrated by a simple model with a single resonance in Ref. .
While no general proof has been given for the validity of local duality, considerable evidence has been accumulated over the years that it does apply for sufficiently heavy flavors. Detailed studies of OPE and duality have been performed recently within model field theories, in particular the t’ Hooft model which is QCD in 1+1 dimensions with $`N_c\mathrm{}`$. Analytical analyses showed that the OPE in the inverse powers of heavy quark mass holds for the heavy quark decay widths in spite of certain doubts which had been voiced. The above papers did not consider the width difference between the two neutral heavy meson eigenstates. Yet using the technique developed there, it is not difficult to establish the similar correspondence between the hadronic saturation and the quark box diagrams at least through the next-to-leading order in $`1/m_Q`$.
Concluding duality to be valid asymptotically – for $`m_Q\mathrm{}`$ – one turns to the question at how low a scale duality emerges to apply with some accuracy. Most authors would expect it to be valid for $`m_Qm_b`$; yet assuming duality to hold already at the charm scale even in a semiquantitative fashion would appear to be a rather iffy proposition for semileptonic transitions, let alone for nonleptonic ones. As argued above, experimental observation of a stronger suppression of $`D^0\overline{D}^0`$ oscillations compared to the phenomenological estimate Eq.(4), would be an indication of a relatively low onset of duality for the inclusive decay widths. The width difference $`\mathrm{\Delta }\mathrm{\Gamma }_D`$ is an even more sensitive, undiluted probe for duality violations than $`\mathrm{\Delta }M_D`$. A scenario with a sizeable $`\mathrm{\Delta }\mathrm{\Gamma }_D`$ and a somewhat smaller value of $`\mathrm{\Delta }M_D`$ could still imply the nearby onset of local duality.
One reservation has to be made though, due to a notorious complication peculiar to local duality violation. Since the duality violating component ‘oscillates’ (as a function of $`m_c`$), it can actually vanish for certain mass values. Determining the size of such effects at a single scale cannot yield a definite conclusion since an accidental vanishing at that scale cannot be ruled out. Yet we have two measures for mixing, namely $`\mathrm{\Delta }M_D`$ and $`\mathrm{\Delta }\mathrm{\Gamma }_D`$, and their oscillatory dependence on $`m_c`$ in general will be out of phase.
## 2 Contributions to $`D^0\overline{D}^0`$ from Exclusive Channels
We have stated above that the OPE expectation of in particular $`y_D𝒪(10^3)`$ is highly remarkable since a priori one would estimate it to be an order of magnitude larger, see Eq. (4). We will illustrate this point by considering transitions to two pseudoscalar mesons, which are common to $`D^0`$ and $`\overline{D}^0`$ decays and can thus communicate between them:
$$D^0\stackrel{CS}{}K^+K^{},\pi ^+\pi ^{}\stackrel{CS}{}\overline{D}^0,$$
(20)
$`D^0`$ $`\stackrel{CA}{}`$ $`K^{}\pi ^+\stackrel{CS^2}{}\overline{D}^0`$
$`D^0`$ $`\stackrel{CS^2}{}`$ $`K^+\pi ^{}\stackrel{CA}{}\overline{D}^0.`$ (21)
where $`CA`$, $`CS`$ and $`CS^2`$ denotes the channel as Cabibbo allowed, Cabibbo suppressed and doubly Cabibbo suppressed, respectively.
In the $`SU(3)`$ limit one obviously has $`\mathrm{\Delta }\mathrm{\Gamma }(D^0K\overline{K},\pi \pi ,K\pi ,\pi \overline{K})=0`$ since the amplitudes for Eqs. (21) would then be equal in size and opposite in sign to those of Eq. (20). Yet the measured branching ratios
$`\mathrm{BR}(D^0K^+K^{})`$ $`=`$ $`(4.27\pm 0.16)10^3`$ (22)
$`\mathrm{BR}(D^0\pi ^+\pi ^{})`$ $`=`$ $`(1.53\pm 0.09)10^3`$ (23)
$`\mathrm{BR}(D^0K^{}\pi ^+)`$ $`=`$ $`(3.85\pm 0.09)10^2`$ (24)
$`\mathrm{BR}(D^0K^+\pi ^{})`$ $`=`$ $`(2.8\pm 0.9)10^4`$ (25)
show very considerable $`SU(3)`$ breakings:
$`{\displaystyle \frac{\mathrm{BR}(D^0K^+K^{})}{\mathrm{BR}(D^0\pi ^+\pi ^{})}}`$ $``$ $`2.8\pm 0.2`$ (26)
$`{\displaystyle \frac{\mathrm{BR}(D^0K^+\pi ^{})}{\mathrm{BR}(D^0K^{}\pi ^+)}}`$ $``$ $`(3\pm 1)\mathrm{tan}^4\theta _C`$ (27)
compared to ratios of unity and $`\mathrm{tan}^4\theta _C`$, respectively, in the symmetry limit.
One would then conclude that the $`K\overline{K},\pi \pi ,K\pi ,\pi \overline{K}`$ contributions to $`\mathrm{\Delta }\mathrm{\Gamma }`$ should be merely Cabibbo suppressed with flavor $`SU(3)`$ providing only moderate further reduction – similar to the general expectation of Eq. (4):
$$\frac{\mathrm{\Delta }\mathrm{\Gamma }}{\mathrm{\Gamma }}|_{DK\overline{K},\pi \pi ,K\pi ,\pi \overline{K}}𝒪(0.01).$$
(28)
Yet despite these large $`SU(3)`$ breakings an almost complete cancellation takes place between their contributions to $`D^0\overline{D}^0`$ oscillations:
$$\mathrm{BR}(D^0K^+K^{})+\mathrm{BR}(D^0\pi ^+\pi ^{})2\sqrt{\mathrm{BR}(D^0K^{}\pi ^+)\mathrm{BR}(D^0K^+\pi ^{})}$$
$$\left(8_{10}^{+12}\right)10^4$$
(29)
to be compared to
$$\mathrm{BR}(D^0K^{}\pi ^+)+\mathrm{BR}(D^0K^+K^{})+\mathrm{BR}(D^0\pi ^+\pi ^{})+\mathrm{BR}(D^0K^+\pi ^{})$$
$$(4.46\pm 0.01)10^2$$
(30)
In principle, a note of caution should be sounded here: In writing down Eq. (29) we have ignored the possibility that $`SU(3)`$ breaking final state interactions can generate a strong phase shift $`\delta _{K\pi }`$ between $`D^0K^{}\pi ^+`$ and $`D^0K^+\pi ^{}`$ amplitudes. If this happens, the last interference term then gets multiplied by a factor $`\mathrm{cos}\delta _{K\pi }`$. Here and in what follows we neglect this phase shift as motivated by the naive quark level diagrams.<sup>5</sup><sup>5</sup>5This fully conforms to the spirit of the OPE description whose prediction we are trying to examine at the level of hadrons. Whether it is small as suggested by some is not clear ; in any case it could have a significant impact on the cancellations among the different terms. Yet we meant this discussion only as a qualitative illustration of our argument on the relation between $`SU(3)`$ symmetry and duality.
There is evidence that Eq. (26) overstates the amount of $`SU(3)`$ breaking in inclusive transitions: the data on Cabibbo suppressed four body modes read
$`\mathrm{BR}(D^0K^+K^{}\pi ^+\pi ^{})`$ $`=`$ $`(2.52\pm 0.24)10^3`$
$`\mathrm{BR}(D^0\pi ^+\pi ^{}\pi ^+\pi ^{})`$ $`=`$ $`(7.4\pm 0.6)10^3;`$ (31)
i.e., again these exclusive channels exhibit very sizeable $`SU(3)`$ breaking
$$\frac{\mathrm{BR}(D^0K^+K^{}\pi ^+\pi ^{})}{\mathrm{BR}(D^0\pi ^+\pi ^{}\pi ^+\pi ^{})}0.34\pm 0.04;$$
(32)
Yet adding these two- and four-body modes then leads to a result which is quite compatible with equality of the combined rates:
$$\frac{\mathrm{BR}(D^0K^+K^{},K^+K^{}\pi ^+\pi ^{})}{\mathrm{BR}(D^0\pi ^+\pi ^{},\pi ^+\pi ^{}\pi ^+\pi ^{})}0.8\pm 0.1.$$
(33)
While this sum cannot be unambiguously related to the violation of $`SU(3)`$ in the fully inclusive rates $`\mathrm{\Gamma }(cs\overline{s}u)`$ vs. $`\mathrm{\Gamma }(cd\overline{d}u)`$, the observed trend is at least suggestive.
To summarize the discussion in this Section:
* The $`SU(3)`$ breaking in exclusive nonleptonic channels is naturally expected to be sizeable, and this is indeed what is observed, see Eq. (26). The deviations from the symmetric case are actually substantially larger than what had been anticipated by most authors.
* Quark based calculations lead to the prediction that inclusive $`D`$ decays exhibit $`SU(3)`$ invariance to a high degree, since the symmetry breaking in described by $`m_s^2/m_c^2𝒪(0.01)`$.
* Emerging data provide the first indication that $`SU(3)`$ breaking is quite reduced when one sums up over various nonleptonic channels, see Eq. (33).
* Likewise the overall contributions to $`\mathrm{\Delta }\mathrm{\Gamma }`$ from channels with two pseudoscalar mesons in the final state appear to be considerably reduced, see Eq. (29).
## 3 Experimental Bounds and Lessons on Duality
The present experimental landscape can be portrayed by the following numbers inferred from various analyses of $`D^0K^+K^{}`$ vs. $`D^0K^{}\pi ^+`$ and $`D^0K^+\pi ^{}`$ vs. $`D^0K^{}\pi ^+`$. From general bounds on mixing one can infer:
$$|x_D|,|y_D|0.028,\mathrm{\hspace{0.33em}\hspace{0.33em}95}\%\mathrm{C}.\mathrm{L}.\mathrm{CLEO}\text{[16]}$$
(34)
Targeting more specifically width differences one finds
$$0.04y_D0.06,\mathrm{\hspace{0.33em}\hspace{0.33em}90}\%\mathrm{C}.\mathrm{L}.\mathrm{E791}\text{[17]}$$
(35)
$$0.058y_D^{}0.01,\mathrm{\hspace{0.33em}\hspace{0.33em}95}\%\mathrm{C}.\mathrm{L}.\mathrm{CLEO}\text{[16]}$$
(36)
The CLEO study analyzes the time evolution of $`D^0(t)K^+\pi ^{}`$ and is thus sensitive to
$$y_D^{}=y_D\mathrm{cos}\delta _{K\pi }x_D\mathrm{sin}\delta _{K\pi }$$
(37)
where $`\delta _{K\pi }`$ denotes the strong phase between $`D^0K^+\pi ^{}`$ and $`\overline{D}^0K^+\pi ^{}`$. A very recent and still preliminary FOCUS study compares the lifetimes for $`DK^+K^{}`$ and $`DK\pi `$:
$$y_D=0.0342\pm 0.0139\pm 0.0074\mathrm{FOCUS}\text{[18]}$$
(38)
At this point we want to summarize and draw the following conclusions:
* Based on general grounds one expects
$$\mathrm{\Delta }M_D,\mathrm{\Delta }\mathrm{\Gamma }_DSU(3)\mathrm{breaking}\times 2\mathrm{sin}^2\theta _C\times \mathrm{\Gamma }_D$$
(39)
The observation of large deviations from $`SU(3)`$ invariance in nonleptonic $`D`$ decays suggests a conservative estimate
$$\frac{\mathrm{\Delta }\mathrm{\Gamma }_D}{\mathrm{\Gamma }_D}\mathrm{few}\times 0.01$$
(40)
with $`\mathrm{\Delta }M_D/\mathrm{\Gamma }_D`$ being somewhat smaller.
* Specific dynamical features have to intervene to suppress $`D^0\overline{D}^0`$ below these levels. Such features arise naturally in a quark level treatment of $`SU(3)`$ symmetry breaking as it arises in an OPE. Assuming local duality one obtains from the OPE the prediction
$$x_D,y_D𝒪(10^3)$$
(41)
without invoking additional long distance contributions. The main uncertainty in this prediction rests in the size of the relevant hadronic matrix elements. The OPE allows for rather mild $`SU(3)`$ GIM cancellations. Consequently, for sufficiently small values of $`m_c`$ there could be unsuppressed contributions to the oscillation rate. Yet for the actual charm mass such contributions should be reasonably suppressed compared to Eq. (40) since they emerge in higher orders in $`1/m_c`$. Therefore, we would consider the degree of suppression of $`x_D`$, $`y_D`$ below the $`1\%`$ level as a measure of applicability of local duality. From this perspective, a stronger suppression of $`x_D`$ compared to $`y_D`$ seems the natural situation.
* The data have reached the general bound of Eq. (4). Any further reduction in the experimental bound on $`y_D`$ means that $`D^0\overline{D}^0`$ oscillations proceed more slowly than can be understood on the basis of general selection rules (a “symmetry level”).
* There is some tentative evidence that inclusive decays might exhibit the effective $`SU(3)`$ invariance expected to arise on the quark level.
* If the suggestion coming from the FOCUS data is confirmed that actually $`y_D𝒪(0.01)`$ holds then one of two conclusions can be drawn: Either $`\mathrm{\Delta }M_D`$ is just “around the corner”, i.e. a moderate improvement in experimental sensitivity should reveal a nonvanishing value for it without establishing the intervention of New Physics. This would mean we had seriously underestimated the size of the relevant matrix elements. Or it would represent a clear-cut violation of local quark-hadron duality at the charm scale.
Note added: After submitting this paper we were informed about the publication (N.U. is grateful to A. Petrov for bringing our attention to it) where it was first proposed that contributions due to chiral symmetry breaking that are subleading in $`1/m_c`$ can generate a moderately larger value for $`\mathrm{\Delta }M_D`$ than the SM box estimate. That paper focussed on an analogue of the OPE for the $`\mathrm{\Delta }C=1`$ transitions rather than directly for $`\mathrm{\Delta }C=2`$. Our analysis differs in a number of conceptual as well as technical aspects. We analyze the OPE for $`\mathrm{\Delta }C=2`$ transitions and obtain the formally leading term linear in $`m_s`$ that had been missed in . We also address the difference between $`\mathrm{\Delta }M_D`$ and $`\mathrm{\Delta }\mathrm{\Gamma }_D`$ explicitly.
Acknowledgments: This work has been supported by the National Science Foundation under grant number PHY96-05080 and by RFFI grant \# 99-02-18355. We thank D. Asner and A. Petrov for helpful comments, and M. Lublinsky for discussions. N.U. gratefully acknowledges the hospitality of Physics Department of the Technion and the support of the Lady Davis grant during completion of this paper.
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# References
Integrable generalised spin ladder models
Angela Foerster<sup>1</sup>, Jon Links<sup>1,2</sup> and Arlei Prestes Tonel<sup>1</sup>
<sup>1</sup>Instituto de Física da UFRGS
Av. Bento Gonçalves 9500, Porto Alegre, RS - Brazil
<sup>2</sup>Department of Mathematics
University of Queensland, Queensland, 4072, Australia
## Abstract
We present two new integrable spin ladder models which posses three general free parameters besides the rung coupling $`J`$. Wang’s systems based on the $`SU(4)`$ and $`SU(3|1)`$ symmetries can be obtained as special cases. The models are exactly solvable by means of the Bethe ansatz method.
Recently there has been a great interest on spin ladder systems from both theoretical and experimental point of view for their relevance to some quasi-one dimensional materials, which under hole doping may exhibit superconductivity . These systems are reasonably well approximated by Heisenberg ladders, which take into account only couplings along the legs and the rungs. Although these systems are not exactly solvable, a variety of solvable ladder models have been found , , . Of particular interest is a general 2-leg spin ladder system with biquadratic interactions proposed by Wang , which by suitable choices of the interchain and interrung coupling originates two integrable spin ladder models based on the $`SU(4)`$ and $`SU(3/1)`$ symmetries. In these cases the rung interactions appear as chemical potentials which break the underlying symmetries of the models. Subsequently other generalised integrable spin ladders have been proposed in the literature . In these cases, no (or few) free parameters are present due to the strict conditions of integrability.
The purpose of this paper is to present two new integrable generalized spin ladders with three extra parameters without violating integrability. These models are exactly solvable by the Bethe ansatz and they reduce to Wang’s models for a special limit of these extra parameters.
Let us begin by introducing the first generalised spin ladder model, whose Hamiltonian reads
$$H^{(1)}=\underset{j=1}{\overset{L}{}}[h_{j,j+1}+{\scriptscriptstyle \frac{1}{2}}J(\stackrel{}{\sigma _j}.\stackrel{}{\tau _j}1)]$$
(1)
where
$`h_{j,j+1}`$ $`=\sigma _j^+\sigma _{j+1}^{}\left[{\displaystyle \frac{t_1^2}{4}}(1+\tau _j^z)(1+\tau _{j+1}^z)+{\displaystyle \frac{t_2^2}{4}}(1\tau _j^z)(1\tau _{j+1}^z)+t_3^2\tau _j^+\tau _{j+1}^{}+\tau _j^{}\tau _{j+1}^+\right]`$
$`+\sigma _j^{}\sigma _{j+1}^+\left[{\displaystyle \frac{t_1^2}{4}}(1+\tau _j^z)(1+\tau _{j+1}^z)+{\displaystyle \frac{t_2^2}{4}}(1\tau _j^z)(1\tau _{j+1}^z)+\tau _j^+\tau _{j+1}^{}+t_3^2\tau _j^{}\tau _{j+1}^+\right]`$
$`+{\scriptscriptstyle \frac{1}{4}}(1+\sigma _j^z)(1+\sigma _{j+1}^z)\left[{\scriptscriptstyle \frac{1}{2}}(1+\tau _j^z\tau _{j+1}^z)+t_1^2\tau _j^+\tau _{j+1}^{}+t_1^2\tau _j^{}\tau _{j+1}^+\right]`$
$`+{\scriptscriptstyle \frac{1}{4}}(1\sigma _j^z)(1\sigma _{j+1}^z)\left[{\scriptscriptstyle \frac{1}{2}}(1+\tau _j^z\tau _{j+1}^z)+t_2^2\tau _j^+\tau _{j+1}^{}+t_2^2\tau _j^{}\tau _{j+1}^+\right].`$
Above $`\stackrel{}{\sigma _j}`$ and $`\stackrel{}{\tau _j}`$ are Pauli matrices acting on site $`j`$ of the upper and lower legs, respectively, $`J`$ is the strength of the rung coupling and $`t_1,t_2,t_3`$ are general independent parameters. $`L`$ is the number of rungs and periodic boundary conditions are imposed. By setting $`t_1,t_2,t_31`$ in equation (1), Wang’s model based on the $`SU(4)`$ symmetry <sup>1</sup><sup>1</sup>1strictly speaking, it is SU(4) in the absence of the rung coupling can be recovered.
The integrability of this model can be shown by the fact that it can be mapped to the Hamiltonian below, which can be derived from an $`R`$matrix obeying the Yang-Baxter algebra for $`J=0`$, while for $`J0`$ the rung interactions take the form of a chemical potential term.
$$\widehat{H}^{(1)}=\underset{j=1}{\overset{L}{}}\left[\widehat{h}_{j,j+1}2JX_j^{00}\right]$$
(2)
where
$`\widehat{h}_{j,j+1}`$ $`={\displaystyle \underset{\alpha =0}{\overset{3}{}}}X_j^{\alpha \alpha }X_{j+1}^{\alpha \alpha }+X_j^{20}X_{j+1}^{02}+X_j^{02}X_{j+1}^{20}`$
$`+t_1^2\left(X_j^{10}X_{j+1}^{01}+X_j^{12}X_{j+1}^{21}\right)+t_2^2\left(X_j^{30}X_{j+1}^{03}+X_j^{32}X_{j+1}^{23}\right)+t_3^2X_j^{31}X_{j+1}^{13}`$
$`+t_1^2\left(X_j^{01}X_{j+1}^{10}+X_j^{21}X_{j+1}^{12}\right)+t_2^2\left(X_j^{03}X_{j+1}^{30}+X_j^{23}X_{j+1}^{32}\right)+t_3^2X_j^{13}X_{j+1}^{31}.`$
Above $`X_j^{\alpha \beta }=|\alpha _j><\beta _j|`$ are the Hubbard operators with $`|\alpha _j>`$ the orthogonalised eigenstates of the local operator $`\stackrel{}{\sigma _j}.\stackrel{}{\tau _j}`$, as in Wang’s case .
The following $`R`$-matrix
$$R=\left(\begin{array}{cccccccccccccccccccc}a& 0& 0& 0& |& 0& 0& 0& 0& |& 0& 0& 0& 0& |& 0& 0& 0& 0& \\ 0& t_1^2b& 0& 0& |& c& 0& 0& 0& |& 0& 0& 0& 0& |& 0& 0& 0& 0& \\ 0& 0& b& 0& |& 0& 0& 0& 0& |& c& 0& 0& 0& |& 0& 0& 0& 0& \\ 0& 0& 0& t_2^2b& |& 0& 0& 0& 0& |& 0& 0& 0& 0& |& c& 0& 0& 0& \\ & & & & & & & & & & & & & & & & & & & \\ 0& c& 0& 0& |& t_1^2b& 0& 0& 0& |& 0& 0& 0& 0& |& 0& 0& 0& 0& \\ 0& 0& 0& 0& |& 0& a& 0& 0& |& 0& 0& 0& 0& |& 0& 0& 0& 0& \\ 0& 0& 0& 0& |& 0& 0& t_1^2b& 0& |& 0& c& 0& 0& |& 0& 0& 0& 0& \\ 0& 0& 0& 0& |& 0& 0& 0& t_3^2b& |& 0& 0& 0& 0& |& 0& c& 0& 0& \\ & & & & & & & & & & & & & & & & & & & \\ 0& 0& c& 0& |& 0& 0& 0& 0& |& b& 0& 0& 0& |& 0& 0& 0& 0& \\ 0& 0& 0& 0& |& 0& 0& c& 0& |& 0& t_1^2b& 0& 0& |& 0& 0& 0& 0& \\ 0& 0& 0& 0& |& 0& 0& 0& 0& |& 0& 0& a& 0& |& 0& 0& 0& 0& \\ 0& 0& 0& 0& |& 0& 0& 0& 0& |& 0& 0& 0& t_2^2b& |& 0& 0& c& 0& \\ & & & & & & & & & & & & & & & & & & & \\ 0& 0& 0& c& |& 0& 0& 0& 0& |& 0& 0& 0& 0& |& t_2^2b& 0& 0& 0& \\ 0& 0& 0& 0& |& 0& 0& 0& c& |& 0& 0& 0& 0& |& 0& t_3^2b& 0& 0& \\ 0& 0& 0& 0& |& 0& 0& 0& 0& |& 0& 0& 0& c& |& 0& 0& t_2^2b& 0& \\ 0& 0& 0& 0& |& 0& 0& 0& 0& |& 0& 0& 0& 0& |& 0& 0& 0& a& \end{array}\right),$$
(3)
with
$$a=x+1,b=x,c=1,$$
obeys the Yang-Baxter algebra
$$R_{12}(xy)R_{13}(x)R_{23}(y)=R_{23}(y)R_{13}(x)R_{12}(xy)$$
(4)
and originates the Hamiltonian (2) for $`J=0`$ by the standard procedure
$$\widehat{h}_{j,j+1}=P\frac{d}{dx}R(x)|_{x=0},$$
where $`P`$ is the permutation operator.
The model can be solved exactly by the Bethe ansatz method and the Bethe ansatz equations read
$`t_1^{2(LM_3)}t_2^{2M_3}t_3^{2M_3}\left({\displaystyle \frac{\lambda _ji/2}{\lambda _j+i/2}}\right)^L`$ $`=`$ $`{\displaystyle \underset{lj}{\overset{M_1}{}}}{\displaystyle \frac{\lambda _j\lambda _li}{\lambda _j\lambda _l+i}}{\displaystyle \underset{\alpha =1}{\overset{M_2}{}}}{\displaystyle \frac{\lambda _j\mu _\alpha +i/2}{\lambda _j\mu _\alpha i/2}}`$
$`t_1^{2(LM_3)}t_2^{2M_3}t_3^{2M_3}{\displaystyle \underset{\beta \alpha }{\overset{M_2}{}}}{\displaystyle \frac{\mu _\alpha \mu _\beta i}{\mu _\alpha \mu _\beta +i}}`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{M_1}{}}}{\displaystyle \frac{\mu _\alpha \lambda _ji/2}{\mu _\alpha \lambda _j+i/2}}{\displaystyle \underset{\delta =1}{\overset{M_3}{}}}{\displaystyle \frac{\mu _\alpha \nu _\delta i/2}{\mu _\alpha \nu _\delta +i/2}}`$ (5)
$`t_1^{2(M_2M_1)}t_2^{2(LM_1+M_2)}t_3^{2(M_1M_2}{\displaystyle \underset{\gamma \delta }{\overset{M_3}{}}}{\displaystyle \frac{\nu _\delta \nu _\gamma i}{\nu _\delta \nu _\gamma +i}}`$ $`=`$ $`{\displaystyle \underset{\alpha =1}{\overset{M_2}{}}}{\displaystyle \frac{\nu _\delta \mu _\alpha i/2}{\nu _\delta \mu _\alpha +i/2}}`$
The energy eigenvalues of the Hamiltonian (2) are given by
$$E=\underset{j=1}{\overset{M_1}{}}\left(\frac{1}{\lambda _j^2+1/4}2J\right)+\frac{3}{4}\left(12J\right)L$$
(6)
where $`\lambda _j`$ are solutions to the Bethe ansatz equations (5).
Now let us introduce a second integrable spin ladder model with three extra parameters, whose Hamiltonian reads
$$=\underset{j=1}{\overset{L}{}}[k_{j,j+1}\frac{1}{4}(12J)(\stackrel{}{\sigma _j}.\stackrel{}{\tau _j}1)]$$
(7)
where
$$k_{j,j+1}=h_{j,j+1}\frac{1}{8}(\stackrel{}{\sigma _j}.\stackrel{}{\tau _j})(\stackrel{}{\sigma _{j+1}}.\stackrel{}{\tau _{j+1}})+\frac{1}{4}(\stackrel{}{\sigma _j}.\stackrel{}{\tau _j}1)$$
(8)
and $`h_{j,j+1}`$ is given by eq. (1).
The solvability of Hamiltonian above lies in the fact that it can be mapped, as before, to a Hamiltonian which can be derived for an $`R`$-matrix satisfying the Yang-Baxter algebra for $`J=0`$, while for $`J0`$ the rung interactions take the form of a chemical potential term
$`\widehat{}={\displaystyle \underset{j=1}{\overset{L}{}}}\left[\widehat{k}_{j,j+1}+(12J)X_j^{00}\right]`$ (9)
where
$`\widehat{k}_{j,j+1}`$ $`={\displaystyle \underset{\alpha =0}{\overset{3}{}}}X_j^{\alpha \alpha }X_{j+1}^{\alpha \alpha }2X_j^{00}X_{j+1}^{00}+X_j^{20}X_{j+1}^{02}+X_j^{02}X_{j+1}^{20}`$ (10)
$`+t_1^2\left(X_j^{10}X_{j+1}^{01}+X_j^{12}X_{j+1}^{21}\right)+t_2^2\left(X_j^{30}X_{j+1}^{03}+X_j^{32}X_{j+1}^{23}\right)+t_3^2X_j^{31}X_{j+1}^{13}`$
$`+t_1^2\left(X_j^{01}X_{j+1}^{10}+X_j^{21}X_{j+1}^{12}\right)+t_2^2\left(X_j^{03}X_{j+1}^{30}+X_j^{23}X_{j+1}^{32}\right)+t_3^2X_j^{13}X_{j+1}^{31}.`$
For $`J=0`$ the model is derived by standard methods from the following $``$-matrix
$$=\left(\begin{array}{cccccccccccccccccccc}w& 0& 0& 0& |& 0& 0& 0& 0& |& 0& 0& 0& 0& |& 0& 0& 0& 0& \\ 0& t_1^2b& 0& 0& |& c& 0& 0& 0& |& 0& 0& 0& 0& |& 0& 0& 0& 0& \\ 0& 0& b& 0& |& 0& 0& 0& 0& |& c& 0& 0& 0& |& 0& 0& 0& 0& \\ 0& 0& 0& t_2^2b& |& 0& 0& 0& 0& |& 0& 0& 0& 0& |& c& 0& 0& 0& \\ & & & & & & & & & & & & & & & & & & & \\ 0& c& 0& 0& |& t_1^2b& 0& 0& 0& |& 0& 0& 0& 0& |& 0& 0& 0& 0& \\ 0& 0& 0& 0& |& 0& a& 0& 0& |& 0& 0& 0& 0& |& 0& 0& 0& 0& \\ 0& 0& 0& 0& |& 0& 0& t_1^2b& 0& |& 0& c& 0& 0& |& 0& 0& 0& 0& \\ 0& 0& 0& 0& |& 0& 0& 0& t_3^2b& |& 0& 0& 0& 0& |& 0& c& 0& 0& \\ & & & & & & & & & & & & & & & & & & & \\ 0& 0& c& 0& |& 0& 0& 0& 0& |& b& 0& 0& 0& |& 0& 0& 0& 0& \\ 0& 0& 0& 0& |& 0& 0& c& 0& |& 0& t_1^2b& 0& 0& |& 0& 0& 0& 0& \\ 0& 0& 0& 0& |& 0& 0& 0& 0& |& 0& 0& a& 0& |& 0& 0& 0& 0& \\ 0& 0& 0& 0& |& 0& 0& 0& 0& |& 0& 0& 0& t_2^2b& |& 0& 0& c& 0& \\ & & & & & & & & & & & & & & & & & & & \\ 0& 0& 0& c& |& 0& 0& 0& 0& |& 0& 0& 0& 0& |& t_2^2b& 0& 0& 0& \\ 0& 0& 0& 0& |& 0& 0& 0& c& |& 0& 0& 0& 0& |& 0& t_3^2b& 0& 0& \\ 0& 0& 0& 0& |& 0& 0& 0& 0& |& 0& 0& 0& c& |& 0& 0& t_2^2b& 0& \\ 0& 0& 0& 0& |& 0& 0& 0& 0& |& 0& 0& 0& 0& |& 0& 0& 0& a& \end{array}\right),$$
(11)
with
$$a=x+1,b=x,c=1,w=x+1,$$
which obeys the Yang-Baxter algebra. The above Hamiltonian has a similar algebraic structure as that of an SU(3/1) supersymmetric t-J model. Using the algebraic nested Bethe ansatz method this model can be solved and the Bethe ansatz equation are given by
$`t_1^{2(LM_3)}t_2^{2M_3}t_3^{2M_3}\left({\displaystyle \frac{\lambda _ji/2}{\lambda _j+i/2}}\right)^L`$ $`=`$ $`{\displaystyle \underset{\alpha =1}{\overset{M_2}{}}}{\displaystyle \frac{\lambda _j\mu _\alpha i/2}{\lambda _j\mu _\alpha +i/2}}`$
$`t_1^{2(LM_3)}t_2^{2M_3}t_3^{2M_3}{\displaystyle \underset{\beta \alpha }{\overset{M_2}{}}}{\displaystyle \frac{\mu _\alpha \mu _\beta i}{\mu _\alpha \mu _\beta +i}}`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{M_1}{}}}{\displaystyle \frac{\mu _\alpha \lambda _ji/2}{\mu _\alpha \lambda _j+i/2}}{\displaystyle \underset{\delta =1}{\overset{M_3}{}}}{\displaystyle \frac{\mu _\alpha \nu _\delta i/2}{\mu _\alpha \nu _\delta +i/2}}`$ (12)
$`t_1^{2(M_2M_1)}t_2^{2(LM_1+M_2)}t_3^{2(M_1M_2)}{\displaystyle \underset{\gamma \delta }{\overset{M_3}{}}}{\displaystyle \frac{\nu _\delta \nu _\gamma i}{\nu _\delta \nu _\gamma +i}}`$ $`=`$ $`{\displaystyle \underset{\alpha =1}{\overset{M_2}{}}}{\displaystyle \frac{\nu _\delta \mu _\alpha i/2}{\nu _\delta \mu _\alpha +i/2}}`$
The eigenenergy of the Hamiltonian (9) is given by
$$E=\underset{j=1}{\overset{M_1}{}}\left(\frac{1}{\lambda _j^2+1/4}+2J1\right)2JL$$
(14)
above $`\lambda _j`$ are solutions of Bethe-ansatz equations (12).
To summarize, we have introduced a new generalization of Wang’s spin ladder models based on the $`SU(4)`$ and $`SU(3|1)`$ symmetries. This was achieved by introducing three extra parameters into the system without violating integrability. The Bethe ansatz equations as well as the energy expressions of the models were presented. The physics of the integrable models presented here is expected to be of interest, since the presence of these extra parameters may turn the phase diagram of the models much richer.
Acknowledgements
JL thanks the Fundação de Amparo a Pesquisa do Estado do Rio Grande do Sul and Australian Research Council for financial support. He also thanks the Instituto de Física da UFRGS for their kind hospitality. AF and APT thank CNPq-Conselho Nacional de Desenvolvimento Científico e Tecnológico for financial support.
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# Strong field limit of the Born-Infeld 𝑝-form electrodynamics
## 1 Introduction
The duality between strong and weak coupling regimes of the underlying theory has played in recent years very prominent role (see e.g. review in ). In the present paper we study the weak-strong field limit correspondence for the $`p`$-form Born-Infeld theory. Recently, Born-Infeld nonlinear electrodynamics (BI) has found a beautiful applications in string theory and $`p`$-brane physics . The motivation to study the corresponding $`p`$-form version of the BI theory comes also from the string theory where one considers extended objects ($`p`$-branes) coupled to a $`p`$-form gauge potential .
In the weak field limit of $`p`$-form BI theory one obtains a linear $`p`$-form Maxwell theory. The corresponding strong field limit is not so well known. It was studied in for $`p=1`$ under the name Ultra Born-Infeld theory (UBI). In the present paper we find the corresponding $`p`$-form UBI and study its properties. It turns out that this theory being nonlinear possesses very instructive features: it is invariant under the full conformal group in (2$`p`$+2)-dimensional Minkowski space-time. Us usual the parity of $`p`$ plays a crucial role. However, in both cases (i.e. for $`p`$ odd and even) the corresponding $`p`$-form UBI displays the full canonical symmetry group of the underlying canonical structure, i.e. $`SO(2,1)`$ and $`SO(1,1)\times Z_2`$ symmetry for odd and even $`p`$ respectively. Moreover, for odd $`p`$ it has an infinite hierarchy of conservation laws. Therefore, the strong field limit of BI theory is even more symmetric than the Maxwell theory which is also conformally invariant and being linear has an infinite hierarchy of constants of motion.
## 2 Born-Infeld $`p`$-form theory
### 2.1 A general theory
Consider a general nonlinear $`p`$-form electrodynamics defined in $`D=2p+2`$ dimensional Minkowski space-time $`^{2p+2}`$ with the signature of the metric tensor $`(,+,\mathrm{},+)`$. The corresponding field tensor $`F=dA`$ ($`A`$ denotes a $`p`$-form gauge potential) gives rise to the following relativistic and gauge invariants:
$`S_p`$ $`=`$ $`{\displaystyle \frac{1}{2(p+1)!}}F_{\mu _1\mathrm{}\mu _{p+1}}F^{\mu _1\mathrm{}\mu _{p+1}},`$ (2.1)
$`P_p`$ $`=`$ $`{\displaystyle \frac{1}{2(p+1)!}}F_{\mu _1\mathrm{}\mu _{p+1}}F^{\mu _1\mathrm{}\mu _{p+1}},`$ (2.2)
where the Hodge star operation in $`^{2p+2}`$ is defined by:
$`F^{\mu _1\mathrm{}\mu _{p+1}}={\displaystyle \frac{1}{(p+1)!}}\eta ^{\mu _1\mathrm{}\mu _{p+1}\nu _1\mathrm{}\nu _{p+1}}F_{\nu _1\mathrm{}\nu _{p+1}}`$ (2.3)
and $`\eta ^{\mu _1\mu _2\mathrm{}\mu _{2p+2}}`$ is the covariantly constant volume form in the Minkowski space-time. The Hodge star satisfies $`=(1)^p`$ which implies the crucial difference between $`p`$-form theories with different parities of $`p`$. Having a Lagrangian $`L_p=L_p(S_p,P_p)`$ one introduces a $`G`$-tensor
$`G^{\mu _1\mathrm{}\mu _{p+1}}=(p+1)!{\displaystyle \frac{L_p}{F_{\mu _1\mathrm{}\mu _{p+1}}}}.`$ (2.4)
Eq. (2.4) defines the constitutive relation for the underlying $`p`$-forms electrodynamics. In the Maxwell theory one has $`G(F)=F`$ but in the general case this relation may be highly nonlinear.
Now one may define the electric and magnetic intensities and inductions in the obvious way:
$`E_I`$ $`=`$ $`F_{I0},`$ (2.5)
$`B_I`$ $`=`$ $`{\displaystyle \frac{1}{(p+1)!}}ϵ_{IJk}F^{Jk},`$ (2.6)
$`D_I`$ $`=`$ $`G_{I0},`$ (2.7)
$`H_I`$ $`=`$ $`{\displaystyle \frac{1}{(p+1)!}}ϵ_{IJk}G^{Jk},`$ (2.8)
where we introduced a $`p`$-index $`I=(i_1i_2\mathrm{}i_p)`$ with $`i_k=1,2,\mathrm{},2p+1`$. $`ϵ_{i_1\mathrm{}i_pj_1\mathrm{}i_pk}=ϵ_{IJk}`$ denotes the Lévi-Civita tensor in $`2p+1`$ dimensional Euclidean space, i.e. a space-like hyperplane $`\mathrm{\Sigma }`$ in the Minkowski space-time. Note, that $`ϵ^{i_1\mathrm{}i_{2p+1}}:=\eta ^{0i_1\mathrm{}i_{2p+1}}`$.
In terms of $`(E,B,D,H)`$ the field equations $`dF=0`$ and $`dG=0`$ read:
$`_0B^I`$ $`=`$ $`(1)^p{\displaystyle \frac{1}{p!}}ϵ^{IkJ}_kE_J,`$ (2.9)
$`_0D^I`$ $`=`$ $`{\displaystyle \frac{1}{p!}}ϵ^{IkJ}_kH_J.`$ (2.10)
They are supplemented by the following constraints (Gauss laws):
$`_iB^i\mathrm{}=_iD^i\mathrm{}=0.`$ (2.11)
The dynamical properties of the $`p`$-form electrodynamics is fully described by the energy-momentum tensor $`T^{\mu \nu }`$:
$`T^{\mu \nu }={\displaystyle \frac{1}{p!}}F^{\mu I}G_I^\nu +g^{\mu \nu }L_p,`$ (2.12)
or in components:
$`T^{00}`$ $`=`$ $`{\displaystyle \frac{1}{p!}}EDL_p,`$ (2.13)
$`T^{0k}`$ $`=`$ $`(1)^{p+1}(E\times H)^k,`$ (2.14)
$`T^{k0}`$ $`=`$ $`(1)^{p+1}(D\times B)^k,`$ (2.15)
$`T^{kl}`$ $`=`$ $`{\displaystyle \frac{1}{(p1)!}}\left(E^k\mathrm{}D_{\mathrm{}}^l+H^k\mathrm{}B_{\mathrm{}}^l\right)+\delta ^{kl}\left({\displaystyle \frac{1}{p!}}HBL_p\right),`$ (2.16)
where we introduced a convenient notation: $`ED:=E^ID_I`$. Moreover,
$`(E\times H)^k={\displaystyle \frac{1}{(p!)^2}}ϵ^{kIJ}E_IH_J.`$ (2.17)
Note, that
$`(E\times H)^k=(1)^p(H\times E)^k.`$
Now, observe that dynamical equations (2.9)-(2.10) have already canonical form. The Hamiltonian $`_p=T^{00}`$ and the corresponding Poisson bracket for the canonical variables $`(D^I,B^J)`$ reads:
$$\{D^I(𝐱),B^J(𝐲)\}_p=ϵ^{IkJ}_k\delta ^{(2p+1)}(𝐱𝐲),$$
(2.18)
(all other brackets vanish).
### 2.2 Canonical symmetries
#### 2.2.1 $`p`$ odd
The $`p`$-form theory based on $`L_p=L_p(S_p,P_p)`$ is obviously relativistically invariant. As is well known in the Hamiltonian framework this invariance is equivalent to the symmetry of the energy-momentum tensor. Let us introduce the following scalar quantities built out of the canonical variables $`(D^I,B^J)`$:
$`\alpha `$ $`=`$ $`{\displaystyle \frac{1}{2p!}}(DD+BB),`$ (2.19)
$`\beta `$ $`=`$ $`{\displaystyle \frac{1}{2p!}}(DDBB),`$ (2.20)
$`\gamma `$ $`=`$ $`{\displaystyle \frac{1}{p!}}DB.`$ (2.21)
Now, the condition $`T^{0k}=T^{k0}`$ which is equivalent to
$$(E\times H)^k=(D\times B)^k,$$
(2.22)
implies the following equation for $`_p`$:
$$(_\alpha _p)^2(_\beta _p)^2(_\gamma _p)^2=1.$$
(2.23)
Eq. (2.23) has a hyperbolic $`SO(2,1)`$ symmetry. It turns out that the group $`SO(2,1)`$ has a natural representation on the level of a canonical structure for a general $`p`$-form theory.
For $`p`$ odd the canonical structure defined in (2.18) is invariant under:
1) duality $`SO(2)`$ rotations:
$`D^I`$ $``$ $`D^I\mathrm{cos}\phi B^I\mathrm{sin}\phi ,`$
$`B^I`$ $``$ $`D^I\mathrm{sin}\phi +B^I\mathrm{cos}\phi ,`$ (2.24)
2) hyperbolic $`SO(1,1)`$ rotations:
$`D^I`$ $``$ $`D^I\mathrm{cosh}\phi +B^I\mathrm{sinh}\phi ,`$
$`B^I`$ $``$ $`D^I\mathrm{sinh}\phi +B^I\mathrm{cosh}\phi ,`$ (2.25)
3) $`R^{}`$-scaling
$`D^I`$ $``$ $`e^\lambda D^I,`$
$`B^I`$ $``$ $`e^\lambda B^I.`$ (2.26)
The easiest way to find the corresponding generators for (2.2.1)-(2.2.1) is to use a two potential formulation , . Let us introduce a $`p`$-form potential $`Z^I`$ for $`D^I`$:
$$D^I=ϵ^{IkJ}_kZ_J,$$
(2.27)
in analogy to
$$B^I=ϵ^{IkJ}_kA_J,$$
(2.28)
where both $`A^I`$ and $`Z^I`$ are in the transverse gauge, i.e. $`_iA^i\mathrm{}=_iZ^i\mathrm{}=0`$. Defining $`A_{(a)}^I=(A^I,Z^I)`$ and $`B_{(a)}^I=(B^I,D^I)`$ the corresponding generators may be written as follows:
$`G_1`$ $`=`$ $`{\displaystyle \frac{1}{2p!}}{\displaystyle d^{2p+1}x\left(A_{(1)}B_{(1)}+A_{(2)}B_{(2)}\right)},`$ (2.29)
$`G_2`$ $`=`$ $`{\displaystyle \frac{1}{2p!}}{\displaystyle d^{2p+1}x\left(A_{(1)}B_{(1)}A_{(2)}B_{(2)}\right)},`$ (2.30)
$`G_3`$ $`=`$ $`{\displaystyle \frac{1}{2p!}}{\displaystyle d^{2p+1}x\left(A_{(1)}B_{(2)}+A_{(2)}B_{(1)}\right)}.`$ (2.31)
For the Maxwell $`p`$-form theory one has $`_p^{Maxwell}=\alpha `$ and (2.23) is trivially satisfied. As is well know Maxwell theory is invariant under (2.2.1), i.e. $`\{_p^{Maxwell},G_1\}=0`$ and, therefore, $`G_1`$ defines a constant of motion. Its physical meaning will be clarified in subsection 2.3.
#### 2.2.2 $`p`$ even
For even $`p`$ the situation is very different. The invariant $`P_p`$ defined in (2.2) vanishes and therefore the general Hamiltonian $`_p`$ does not depend upon $`\gamma `$ defined in (2.21). Eq. (2.23) reduces now to
$$(_\alpha _p)^2(_\beta _p)^2=1.$$
(2.32)
Eq. (2.32) displays the $`SO(1,1)`$ symmetry which is now realized by (2.2.1). Neither (2.2.1) nor (2.2.1) are implementable as canonical transformations and (2.2.1) is generated by:
$`G_4`$ $`=`$ $`{\displaystyle \frac{1}{2p!}}{\displaystyle d^{2p+1}x\left(A_{(1)}B_{(2)}A_{(2)}B_{(1)}\right)}.`$ (2.33)
However, in this case we have two additional discrete $`Z_2`$ symmetries:
$$D^IB^I,B^ID^I,$$
(2.34)
and
$$D^IB^I,B^ID^I.$$
(2.35)
Now, Maxwell theory is only $`Z_2\times Z_2`$-invariant, i.e. w.r.t. (2.34) and (2.35).
### 2.3 $`p`$–photons
For odd $`p`$ the canonical generator $`G_1`$ defined in (2.29) is constant in time (in the Maxwell theory). To find the physical interpretation of this quantity let us introduce a complex notation:
$`F^I`$ $`=`$ $`D^I+iB^I,`$
$`G^I`$ $`=`$ $`E^I+iH^I.`$
The dynamical equations (2.9)-(2.10) rewritten in terms of $`(F^I,G^J)`$ have the following form:
$$i_0F^I=\frac{1}{p!}ϵ^{IkJ}_kG_J.$$
(2.36)
Moreover, let
$$V^I=Z^I+iA^I,$$
i.e. $`F^I=ϵ^{IkJ}_kV_J`$. Now, introduce a Fourier representation for $`V(𝐱)`$:
$$\stackrel{~}{V}_I(𝐤)=d^{2p+1}xe^{i\mathrm{𝐤𝐱}}V_I(𝐱).$$
(2.37)
The transverse gauge $`_iV^i\mathrm{}(𝐱)=0`$ implies $`k_l\stackrel{~}{V}^l\mathrm{}(𝐤)=0`$. Let $`\omega ^{(1)},\omega ^{(2)},\mathrm{},\omega ^{(N_p)}`$ form an orthonormal basis of $`p`$-forms on $`2p`$-dimensional plane perpendicular to $`𝐤`$, i.e. $`\omega ^{(i)}\omega ^{(j)}=\delta ^{ij}`$. The number $`N_p`$ reads
$`N_p=\left(\begin{array}{c}2p\\ p\end{array}\right),`$ (2.40)
and it is equal to the number of degrees of freedom for a $`p`$-form theory . Now, for any $`\alpha \{1,2,\mathrm{},N_p\}`$ there exists exactly one index $`\alpha ^{}`$ such that
$$(\omega ^{(\alpha )}\times \omega ^{(\alpha ^{})})^l=\frac{k^l}{k}.$$
(2.41)
Therefore, let us define a complex basis:
$$e^{(\alpha )}=\frac{1}{\sqrt{2}}(\omega ^{(\alpha )}+i\omega ^{(\alpha ^{})}),$$
(2.42)
where $`\alpha `$ runs from 1 to $`N_p/2`$. Note, that the complex basis $`e^{(\alpha )}`$ satisfies:
$`ϵ^{IlJ}k_le_J^{(\alpha )}=ike^{(\alpha )I}`$ (2.43)
and obviously: $`\overline{e^{(\alpha )}}e^{(\beta )}=\delta ^{\alpha \beta }`$ and $`e^{(\alpha )}e^{(\beta )}=0`$, where $`\overline{a}`$ denotes a complex conjugation of $`a`$.
Now, decomposing the Fourier transform of $`V_I(𝐱)`$:
$`\stackrel{~}{V}_I(𝐤)={\displaystyle \underset{\alpha =1}{\overset{N_p/2}{}}}\left(e_I^{(\alpha )}f_+^{(\alpha )}(𝐤)+\overline{e^{(\alpha )}}_If_{}^{(\alpha )}(𝐤)\right),`$ (2.44)
and inserting into (2.29) one obtains:
$`G_1`$ $`=`$ $`{\displaystyle \frac{1}{2p!}}{\displaystyle d^{2p+1}xV_I(𝐱)ϵ^{IkJ}_k\overline{V_J(𝐱)}}`$ (2.45)
$`=`$ $`{\displaystyle \frac{1}{2p!(2\pi )^{2(2p+1)}}}{\displaystyle d^{2p+1}kk\underset{\alpha =1}{\overset{N_p/2}{}}\left(|f_+^{(\alpha )}(𝐤)|^2|f_{}^{(\alpha )}(𝐤)|^2\right)}.`$
Note, that the energy of the Maxwell field reads:
$`_p`$ $`=`$ $`{\displaystyle \frac{1}{2p!}}{\displaystyle d^{2p+1}x\left(DD+BB\right)}={\displaystyle \frac{1}{2p!}}{\displaystyle d^{2p+1}xF(𝐱)\overline{F(𝐱)}}`$ (2.46)
$`=`$ $`{\displaystyle \frac{1}{2p!}}{\displaystyle d^{2p+1}xϵ^{IkJ}_kV_J(𝐱)ϵ_{IjL}^j\overline{V^L(𝐱)}}`$
$`=`$ $`{\displaystyle \frac{1}{2p!(2\pi )^{2(2p+1)}}}{\displaystyle d^{2p+1}kk^2\underset{\alpha =1}{\overset{N_p/2}{}}\left(|f_+^{(\alpha )}(𝐤)|^2+|f_{}^{(\alpha )}(𝐤)|^2\right)}.`$
Now, $`_p`$ measures the sum of intensities of all possible polarizations. On the other hand $`G_1`$ measures the intensity of right-handed polarizations minus the intensity of left-handed polarizations. In the quantum theory of a $`p`$-form Maxwell field they correspond to different polarization states of “$`p`$-photons” ($`N_p/2`$ left-handed defined by $`e^{(\alpha )}`$ and $`N_p/2`$ right-handed defined by $`\overline{e^{(\alpha )}}`$).
Remarkably, (2.46) is valid in the linear (Maxwell) case only but (2.45) holds for any $`p`$-form theory and defines a constant of motion for any duality invariant theory. Neither $`G_2`$ nor $`G_3`$ have any clear physical interpretation.
### 2.4 Born-Infeld model
The Born-Infeld $`p`$-form theory is defined by the following Lagrangian:
$$L_p^{(BI)}=\frac{b^2}{p!}\left(1\sqrt{12b^2p!S_p(b^2p!P_p)^2}\right),$$
(2.47)
where $`b`$ denotes a generalized fundamental parameter of Born and Infeld . In terms of $`E^I`$ and $`B^I`$ our two basic invariants read:
$`S_p`$ $`=`$ $`{\displaystyle \frac{1}{2p!}}(EEBB),`$
$`P_p`$ $`=`$ $`\{\begin{array}{cc}\frac{1}{p!}EB,\hfill & \text{for odd}p,\hfill \\ 0\hfill & \text{for even}p.\hfill \end{array}`$ (2.50)
The corresponding $`D^I`$ and $`H^I`$ fields read:
$`D^I`$ $`=`$ $`{\displaystyle \frac{1}{l_p}}\left(E^I+b^2P_pB^I\right),`$
$`H^I`$ $`=`$ $`{\displaystyle \frac{1}{l_p}}\left(B^Ib^2P_pE^I\right),`$
with $`l_p=\sqrt{12b^2p!S_p(b^2p!P_p)^2}`$. From (2.13) one easily gets the corresponding Hamiltonian:
$`_p^{(BI)}={\displaystyle \frac{b^2}{p!}}\left(\sqrt{1+b^2(DD+BB)+b^4\left[(DD)(BB)ϵ_p(DB)^2\right]}1\right),`$ (2.51)
where
$`ϵ_p=\{\begin{array}{cc}1\hfill & \text{odd}p\hfill \\ 0\hfill & \text{even}p\hfill \end{array}.`$ (2.54)
Note, that for any $`p`$-forms $`D`$ and $`B`$
$`(p!)^2|D\times B|^2=(DD)(BB)+(1)^p(DB)^2.`$
Therefore, for odd $`p`$
$`_p^{(BI)}={\displaystyle \frac{b^2}{p!}}\left(\sqrt{1+b^2(DD+BB)+b^4(p!)^2|D\times B|^2}1\right).`$ (2.55)
In terms of $`(\alpha ,\beta ,\gamma )`$ the BI Hamiltonian (2.51) reads:
$`_p^{(BI)}=b^2\left(\sqrt{1+b^2\alpha +b^4(\alpha ^2\beta ^2ϵ_p\gamma ^2)}1\right),`$ (2.56)
and satisfies (2.23) or (2.32) for odd or even $`p`$ respectively.
It is easy to see that $`p`$-form BI is invariant under:
* dual $`SO(2)`$ rotations (2.2.1) for $`p`$ odd,
* $`Z_2\times Z_2`$ transformations (2.34) and (2.35) for $`p`$ even,
exactly as Maxwell theory.
## 3 Strong field limit
The crucial property of the Born-Infeld model is that the magnitude of $`E^I`$ and $`B^I`$ fields is bounded by the value of the critical parameter $`b`$, i.e. $`|F^{\mu _1\mathrm{}\mu _{p+1}}|<b`$. Therefore, we are not able to study the strong field limit of this model in the Lagrangian framework. Actually, performing the $`b0`$ limit in the Born-Infeld Lagrangian (2.47) one obtains $`|P_p|`$ which defines a trivial theory. However, in the Hamiltonian framework the $`D^I`$ field is not bounded and the question about the strong field limit is well posed. The same property displays the relativistic particle’s dynamics: the particle’s velocity is always bounded $`|𝐯|<c`$ contrary to its momentum $`𝐩`$ (or energy $`E`$) and the ultrarelativistic limit is defined by $`|𝐩|mc`$ (or $`Emc^2`$). The particle’s Hamiltonian
$$H(𝐪,𝐩)=c\sqrt{𝐩^2+(mc)^2}+V(𝐪),$$
(3.1)
tends in the ultrarelativistic limit to
$$H^U(𝐩)=c|𝐩|$$
(3.2)
which implies
$$\dot{𝐩}=0,𝐯=c𝐧,$$
(3.3)
with $`𝐧=𝐩/|𝐩|`$, i.e. one has a free evolution of photons. By analogy we call the strong field limit of BI theory the Ultra Born-Infeld theory (UBI) .
### 3.1 $`p`$ odd
Performing $`b0`$ limit in (2.51) one gets
$$_p^{(UBI)}=|D\times B|.$$
(3.4)
Therefore, the constitutive relations are given by:
$`E_I`$ $`=`$ $`p!{\displaystyle \frac{\delta _p^{(UBI)}}{\delta D^I}}=ϵ_{kIJ}n^kB^J,`$ (3.5)
$`H_I`$ $`=`$ $`p!{\displaystyle \frac{\delta _p^{(UBI)}}{\delta B^I}}=ϵ_{kIJ}n^kD^J,`$ (3.6)
where $`n^k`$ stands for the unit (2p+1)-vector in the direction of the generalized Poynting vector:
$`n^k={\displaystyle \frac{(D\times B)^k}{|D\times B|}}.`$ (3.7)
The dynamical equations (2.9)-(2.10) have the following form:
$`_0B^I`$ $`=`$ $`\delta _{[jJ]}^{kI}_k(n^jB^J),`$ (3.8)
$`_0D^I`$ $`=`$ $`\delta _{[jJ]}^{kI}_k(n^jD^J),`$ (3.9)
with $`\delta _{kl\mathrm{}}^{ij\mathrm{}}=\delta _k^i\delta _l^j\mathrm{}`$. The remarkable property of $`p`$-form UBI is a structure of the energy-momentum tensor. One easily finds:
$`T^{0k}`$ $`=`$ $`T^{k0}=_p^{(UBI)}n^k,`$ (3.10)
$`T^{kl}`$ $`=`$ $`_p^{(UBI)}n^kn^l.`$ (3.11)
These relations were already derived by Białynicki-Birula for $`p=1`$. But they hold for any odd $`p`$. In terms of $`(\alpha ,\beta ,\gamma )`$ the UBI Hamiltonian (3.4) reads:
$`_p^{(UBI)}=\sqrt{\alpha ^2\beta ^2\gamma ^2},`$ (3.12)
and displays the full $`SO(2,1)`$ symmetry of (2.23), i.e. $`p`$-form UBI is invariant under: (2.2.1), (2.2.1) and (2.2.1).
Moreover, for any $`p`$ the trace of $`T^{\mu \nu }`$ for UBI vanishes and, therefore, UBI is invariant under the conformal group in $`^{2p+2}`$. This last property follows from the fact that UBI does not contain a dimensional parameter.
Note, that after performing the Legendre transformation the UBI Lagrangian vanishes: $`L_p^{(UBI)}=DE_p^{(UBI)}0`$. It does not mean, however, that the theory is trivial (we already know that it is not). Vanishing of $`L_p^{(UBI)}`$ denotes the presence of Lagrangian constraints. It is easy to see that the constitutive relations imply $`S_p=0`$ and $`P_p=0`$. Therefore, the UBI action has the following form:
$$W_p^{(UBI)}=d^{2p+2}x\left(\mathrm{\Lambda }_1S_p+\mathrm{\Lambda }_2P_p\right),$$
(3.13)
where $`\mathrm{\Lambda }_1`$ and $`\mathrm{\Lambda }_2`$ are Lagrange multipliers. Now, (3.13) implies
$`D^I=\mathrm{\Lambda }_1E^I+\mathrm{\Lambda }_2B^I,`$
and hence $`|D\times B|=\mathrm{\Lambda }_1BB`$ (since $`EE=BB`$), which gives $`\mathrm{\Lambda }_1=|D\times B|/BB`$. Moreover, $`DB=\mathrm{\Lambda }_2BB`$ and hence $`\mathrm{\Lambda }_2=DB/BB`$. Inserting
$`E^I`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Lambda }_1}}D^I{\displaystyle \frac{\mathrm{\Lambda }_2}{\mathrm{\Lambda }_1}}B^I,`$
$`H^I`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Lambda }_1^2+\mathrm{\Lambda }_2^2}{\mathrm{\Lambda }_1^2}}B^I{\displaystyle \frac{\mathrm{\Lambda }_2}{\mathrm{\Lambda }_1}}D^I,`$
into field eqs. (2.9)–(2.9) one gets (3.8)–(3.9).
### 3.2 $`p`$ even
Now, due to (2.51) the UBI Hamiltonian simplifies to
$`_p^{(UBI)}={\displaystyle \frac{1}{p!}}|D||B|,`$ (3.14)
giving rise to the following constitutive relations:
$`E^I`$ $`=`$ $`{\displaystyle \frac{|B|}{|D|}}D^I,`$ (3.15)
$`H^I`$ $`=`$ $`{\displaystyle \frac{|D|}{|B|}}B^I.`$ (3.16)
The corresponding stress tensor reads:
$$T^{kl}=_p^{(UBI)}\left(\delta ^{kl}p\left(|D|^2D^k\mathrm{}D_{\mathrm{}}^l+|B|^2B^k\mathrm{}B_{\mathrm{}}^l\right)\right).$$
(3.17)
Now, dynamical equations (2.9)-(2.10) read:
$`_0B^I`$ $`=`$ $`{\displaystyle \frac{1}{p!}}ϵ^{IkJ}_k\left({\displaystyle \frac{|B|}{|D|}}D_J\right),`$ (3.18)
$`_0D^I`$ $`=`$ $`{\displaystyle \frac{1}{p!}}ϵ^{IkJ}_k\left({\displaystyle \frac{|D|}{|B|}}B_J\right).`$ (3.19)
Note, that (3.14) rewritten in terms of $`(\alpha ,\beta )`$ has $`SO(1,1)`$–invariant form:
$`_p^{(UBI)}=\sqrt{\alpha ^2\beta ^2}.`$ (3.20)
The theory displays the full symmetry of the canonical structure, i.e. it is invariant under (2.2.1) generated by $`G_4`$ and under $`Z_2\times Z_2`$ defined in (2.34)–(2.35). Obviously, $`T_\mu ^\mu =0`$ and the theory is conformally invariant. The Lagrangian structure may be derived from the following action
$$W_p^{(UBI)}=d^{2p+2}x\mathrm{\Lambda }S_p,$$
(3.21)
($`\mathrm{\Lambda }`$ – Lagrange multiplier) in analogy to (3.13).
Note, that for Cauchy data satisfying
$`DB=0.`$ (3.22)
one has
$`|D\times B|={\displaystyle \frac{1}{p!}}|D||B|,`$ (3.23)
and the Hamiltonian (3.14) has the same form as in (3.4). The constitutive relations (3.15)–(3.16) are equivalent to:
$`E_I`$ $`=`$ $`ϵ_{kIJ}n^kB^J,`$ (3.24)
$`H_I`$ $`=`$ $`ϵ_{kIJ}n^kD^J,`$ (3.25)
with $`n_k`$ defined in (3.7) and the stress tensor (3.17) may be rewritten as in (3.11).
## 4 Fluid dynamics and new constants of motion
In this section we generalize the observation made in for any odd $`p`$. Observe that due to (3.10)-(3.11), $`T^{\mu \nu }`$ may be written in the following form:
$$T^{\mu \nu }=_p^{(UBI)}U^\mu U^\nu ,$$
(4.1)
where the (2$`p`$+2)-velocity $`U^\mu =(1,n^k)`$ satisfies $`U^\mu U_\mu =0`$ (for even $`p`$ it holds for the Cauchy data satisfying (3.22)). Such a theory describes a dust of particles moving with the speed of light in $`^{2p+2}`$ – “$`p`$-photons”. It is easy to show that both the continuity equation
$$_\mu (_p^{(UBI)}U^\mu )=0,$$
(4.2)
and the Euler equation
$$U^\nu _\nu U^\mu =0$$
(4.3)
are satisfied. Moreover, one easily proves that due to (4.2) and (4.3) the following infinite set of continuity equations hold:
$$_\mu \left(_p^{(UBI)}U^\mu U^{i_1}U^{i_2}\mathrm{}U^{i_k}\right)=0.$$
(4.4)
They give rise to the following hierarchy of conserved quantities:
$$𝐊^{i_1\mathrm{}i_k}=d^{2p+1}x_p^{(UBI)}U^{i_1}U^{i_2}\mathrm{}U^{i_k}.$$
(4.5)
All these quantities are in involution, i.e.
$$\{𝐊^{i_1\mathrm{}i_k},𝐊^{j_1\mathrm{}j_l}\}_p=0.$$
(4.6)
The only exception is $`p=0`$. In this case $`|u^1|=1`$ and all $`𝐊^{1\mathrm{}1}`$ are aqual (up to a sign). But now the conformal group is infinite dimensional and one has still an infinite number of constants of motion.
For other relation between Born-Infeld theory and fluid dynamics see e.g. .
## 5 Self-dual field
Note, that the Maxwell eqs. for even $`p`$ have the following form:
$`_0B^I`$ $`=`$ $`{\displaystyle \frac{1}{p!}}ϵ^{IkJ}_kD_J,`$ (5.1)
$`_0D^I`$ $`=`$ $`{\displaystyle \frac{1}{p!}}ϵ^{IkJ}_kB_J,`$ (5.2)
and differ from the field eqs. of UBI (3.18)–(3.19) by the presence of the scaling parameter $`|D|/|B|`$.
Now, let us introduce chiral and anti-chiral combinations:
$`V_\pm ^I={\displaystyle \frac{1}{\sqrt{2}}}(D^I\pm B^I).`$ (5.3)
The Maxwell Hamiltonian rewritten in terms of $`V_\pm `$ has the following form
$$_p^{Maxwell}=\frac{1}{2p!}\left(V_+V_++V_{}V_{}\right),$$
(5.4)
and the Maxwell eqs. read:
$$V_\pm ^I=\pm \frac{1}{p!}ϵ^{IkJ}_kV_{J\pm },$$
(5.5)
i.e. both components decouple and evolve independently. This property holds for the Maxwell theory only. However, if the initial condition is such that $`V_{}=0`$, i.e. $`D=B`$ or $`V_{}=0`$, i.e. $`D=B`$, then for any $`t`$, $`V_{}(t)=0`$ or $`V_+(t)=0`$ also for the UBI theory defined by (3.18)-(3.19). This property holds for any $`p`$-form theory defined by a Hamiltonian satisfying (2.32) and invariant under $`Z_2\times Z_2`$ given by (2.34)-(2.35). Therefore, for chiral (anti-chiral) data the dynamics always reduces to
$$B^I=\pm \frac{1}{p!}ϵ^{IkJ}_kB_J,$$
(5.6)
i.e. it corresponds to (anti)self-dual field . For (anti)self-dual data both Maxwell and UBI Hamiltonians read:
$$_p^{self}=\frac{1}{p!}BB,$$
(5.7)
and (5.6) defines the Hamiltonian system with respect to the following canonical structures:
$$\{B^I(𝐱),B^J(𝐲)\}=\pm ϵ^{IkJ}_k\delta ^{(2p+1)}(𝐱𝐲).$$
(5.8)
## 6 Conclusions
Note, that $`p`$-form UBI is uniquely defined by the symmetry group. Namely, the UBI Hamiltonian is a maximally symmetric solution of (2.23) or (2.32) for odd and even $`p`$ respectively. Consider e.g. (2.23). Its $`SO(2,1)`$-symmetric solution is a function of one variable $`f=f(\tau )`$ with
$`\tau =\alpha ^2\beta ^2\gamma ^2.`$
One immediately shows that $`f=\sqrt{\tau }`$ in agreement with (3.12).
For odd $`p`$ the $`p`$-form theory is duality invariant (i.e. invariant under $`SO(2)`$ rotations (2.2.1)) iff
$$(_\eta L_p)^2(_\xi L_p)^2=1,$$
(6.1)
where
$`\eta `$ $`=`$ $`\sqrt{S_p^2+P_p^2},`$
$`\xi `$ $`=`$ $`S_p.`$
Eq. (6.1) has the same form as (2.32). Now, the maximally symmetric solution to (6.1) reads $`L_p=|P_p|`$, i.e. the corresponding theory is trivial. Note, that performing $`b0`$ limit in the BI Lagrangian (2.47) we have also obtained $`|P_p|`$.
The presence of an infinite hierarchy of constants of motion often implies complete integrability of the theory. This question for $`p=1`$ UBI was posed in . The answer is not known. It would be also interesting to find the corresponding quantum version of this theory.
Note added: While this paper was being completed I received unpublished notes , in which similar results were obtained. Moreover, in the particles fluid of section 4 was generalized to $`p`$-brane fluids.
## Acknowledgements
I thank prof. I. Białynicki-Birula for pointing out the Ref. anf prof. G. W. Gibbons for sending me his unpublished notes . This work was partially supported by the KBN Grant No. 2 P03A 047 15.
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# A Time–Dependent Born–Oppenheimer Approximation with Exponentially Small Error Estimates
## 1 Introduction
In this paper we construct exponentially accurate approximate solutions to the time–dependent Schrödinger equation for a molecular system. The small parameter that governs the approximation is the usual Born–Oppenheimer expansion parameter $`ϵ`$, where $`ϵ^4`$ is the ratio of the electron mass divided by the mean nuclear mass. The approximate solutions we construct agree with exact solutions up to errors whose norms are bounded by $`C\mathrm{exp}\left(\gamma /ϵ^2\right)`$, for some $`C`$ and $`\gamma >0`$, under analyticity assumptions on the electron Hamiltonian.
The Hamiltonian for a molecular system with $`K`$ nuclei and $`NK`$ electrons moving in $`l`$ dimensions has the form
$$H(ϵ)=\underset{j=1}{\overset{K}{}}\frac{ϵ^4}{2M_j}\mathrm{\Delta }_{X_j}\underset{j=K+1}{\overset{N}{}}\frac{1}{2m_j}\mathrm{\Delta }_{X_j}+\underset{i<j}{}V_{ij}(X_iX_j).$$
Here $`X_j\mathrm{I}\mathrm{R}^l`$ denotes the position of the $`j^{\text{th}}`$ particle, the mass of the $`j^{\text{th}}`$ nucleus is $`ϵ^4M_j`$ for $`1jK`$, the mass of the $`j^{\text{th}}`$ electron is $`m_j`$ for $`K+1jN`$, and the potential between particles $`i`$ and $`j`$ is $`V_{ij}`$. For convenience, we assume each $`M_j=1`$. We set $`d=Kl`$ and let $`X=(X_1,X_2,\mathrm{},X_K)\mathrm{I}\mathrm{R}^d`$ denote the nuclear configuration vector. We can then decompose $`H(ϵ)`$ as
$$H(ϵ)=\frac{ϵ^4}{2}\mathrm{\Delta }_X+h(X).$$
The first term on the right hand side represents the nuclear kinetic energy, and the second is the “electron Hamiltonian” that depends parametrically on $`X`$. For each fixed $`X`$, $`h(X)`$ is a self-adjoint operator on the Hilbert space $`_{\text{el}}=L^2(\mathrm{I}\mathrm{R}^{(NK)l})`$.
The time–dependent Schrödinger equation we approximately solve in $`L^2(\mathrm{I}\mathrm{R}^d,_{\text{el}})`$ as $`ϵ0`$ is
$$iϵ^2\frac{\psi }{t}=\frac{ϵ^4}{2}\mathrm{\Delta }_X\psi +h(X)\psi $$
Asymptotic expansions in powers of $`ϵ`$ of certain solutions to this equation are derived in . We obtain our construction by truncating these expansions after an $`ϵ`$–dependent number of terms, in an effort to minimize the norm of the error. Similar strategies have been used to obtain exponentially accurate results for adiabatic approximations and semiclassical approximations , both of which play roles in the Born–Oppenheimer approximation we are studying here.
Roughly speaking, the time–dependent Born–Oppenheimer approximation says the following for small $`ϵ`$: The electrons move very rapidly and adjust their state adiabatically as the more slowly moving nuclei change their positions. If the electrons start in a discrete energy level of $`h(X)`$, they will remain in that level. In the process, the electron states create an effective potential in which the motion of the heavy nuclei is well described by a semiclassical approximation. The asymptotic expansions show that this intuition is valid up to errors of order $`ϵ^k`$ for any $`k`$.
In Born–Oppenheimer approximations, adiabatic and semiclassical limits are being taken simultaneously, and they are coupled. Analysis of errors for the adiabatic and semiclassical approximations shows that they are each accurate up to errors whose bounds have the form $`C\mathrm{exp}\left(\gamma /ϵ^2\right)`$ . Non-adiabatic transitions are known in some systems to be of this order, and tunnelling in semiclassical approximations makes contributions of this order. Thus, one cannot expect to do better than approximations of this type because of two well-known physical phenomena that Born–Oppenheimer approximations do not take into account.
In some systems, tunnelling might dominate the error. In some, non-adiabatic electronic transitions may dominate. In others, the two effects can be of comparable magnitude.
One of the motivations for our work is to generate a “good” basis upon which to build a “surface hopping model” that would accurately describe non-adiabatic electronic transitions. Prior authors (see, e.g., ) have proposed such models based on the zeroth order time–dependent Born–Oppenheimer approximation. Using the zeroth order states as a basis of the surface hopping model, the non-adiabatic transitions appear at order $`ϵ^2`$. This is huge compared to the exponentially small physical phenomenon one would like to study, and we believe interference between transitions that occur at different times is responsible for the exponential smallness of the physically interesting quantity. Our view is that by choosing a much better set of states on which to base the model, one will obtain a much more useful approximation. Sir Michael Berry has advocated such ideas for the somewhat simpler adiabatic approximation (which does not have the complications of the nuclear motion). These ideas have been used in to prove the accuracy of certain results for non–adiabatic transitions that are exponentially small.
Remarks:
1. There are some other exponentially accurate results in the general topic of Born–Oppenheimer approximations. The prior results come from study of the time–independent Schrödinger equation and depend on global properties of the system. Our results are time–dependent and make use of local information.
Klein and Martinez show that resonances associated with predissociation processes have exponentially long lifetimes. Benchaou and Martinez also show that certain S–matrix elements associated with non–adiabatic transitions are exponentially small.
2. The papers cited in the previous remark obtain estimates that depend on the global structure of the electron energy levels. The results we obtain depend on a particular classical path. When the path stays away from the nuclear configurations where the gap between relevant electonic levels is minimized, one would expect the non–adiabatic errors from our approximation to be smaller, i.e., both results would obtain errors of order $`\mathrm{exp}(\mathrm{\Gamma }/ϵ^2)`$, but we would obtain a larger value of $`\mathrm{\Gamma }`$.
We expect this because in our case, the Landau–Zener formula predicts that our $`\mathrm{\Gamma }`$ should come from the minimum gap between eigenvalues on the classical path, rather than the global minimum gap.
3. From a mathematical point of view, the optimal truncation procedure in this context was first stated for the adiabatic approximation for two component systems of ODE’s by Berry . It was first proved to yield exponentially accurate results for Hilbert space valued ODE’s by Nenciu . In we used this idea for the semiclassical approximation, which is a complex valued PDE setting. The present paper can be viewed as extending these ideas to a Hilbert space valued PDE setting.
### 1.1 Hypotheses
We assume that the electron Hamiltonian $`h(X)`$ satisfies the following analyticity hypotheses:
* + For any $`X\mathrm{I}\mathrm{R}^d`$, $`h(X)`$ is a self-adjoint operator on some dense domain $`𝒟_{\text{el}}`$, where $`_{\text{el}}`$ is the electronic Hilbert space. We assume the domain $`𝒟`$ is independent of $`X`$ and $`h(X)`$ is bounded from below uniformly in $`\mathrm{I}\mathrm{R}^d`$.
+ There exists a $`\delta >0`$, such that for every $`\psi 𝒟`$, the vector $`h(X)\psi `$ is analytic in
$`S_\delta =\{z\mathrm{I}\mathrm{C}^d:|\text{Im}(z_j)|<\delta ,j=1,\mathrm{},d\}`$.
* There exists an open set $`\mathrm{\Xi }\mathrm{I}\mathrm{R}^d`$, such that for all $`X\mathrm{\Xi }`$, there exists an isolated, multiplicity one eigenvalue $`E(X)`$ of $`h(X)`$ associated with a normalized eigenvector $`\mathrm{\Phi }(X)_{\text{el}}`$. We assume without loss that the origin belongs to $`\mathrm{\Xi }`$.
Remarks: 1. Hypothesis H<sub>0</sub> implies that the family of operators $`\{h(X)\}_{XS_\delta }`$ is a holomorphic family of type A.
2. It follows from H<sub>0</sub> and H<sub>1</sub> that there exists $`\delta ^{}(0,\delta )`$ and $`\mathrm{\Xi }^{}\mathrm{\Xi }`$ such that the complex and vector valued functions $`E()`$ and $`\mathrm{\Phi }()`$ admit analytic continuations on the set
$`\mathrm{\Sigma }_\delta ^{}=\{z\mathrm{I}\mathrm{C}^d:\text{Re}(z)\mathrm{\Xi }^{}\text{and}|\text{Im}(z_j)|<\delta ^{},j=1,\mathrm{},d\}`$.
### 1.2 Summary of the Main Results
Our main results are stated precisely as Theorem 4.1 in Section 4. Two generalizations of this result are presented in Section 8.
Roughly speaking, Theorem 4.1 states the following:
Under hypotheses H<sub>0</sub> and H<sub>1</sub>, we construct $`\mathrm{\Psi }_{}(X,t,ϵ)`$ (that depends on a parameter $`g`$) for $`t[0,T]`$. For small values of $`g`$, there exist $`C(g)`$ and $`\mathrm{\Gamma }(g)>0`$, such that in the limit $`ϵ0`$,
$$\text{e}^{itH(ϵ)/ϵ^2}\mathrm{\Psi }_{}(X,0,ϵ)\mathrm{\Psi }_{}(X,t,ϵ)_{L^2(\text{IR}^d,_{\text{el}})}C(g)\text{e}^{\mathrm{\Gamma }(g)/ϵ^2}$$
In the state $`\mathrm{\Psi }_{}(X,t,ϵ)`$, the electrons have a high probability of being in the electron state $`\mathrm{\Phi }(X)`$. For any $`b>0`$ and sufficiently small values of $`g`$, the nuclei are localized near a classical path $`a(t)`$ in the sense that there exist $`c(g)`$ and $`\gamma (g)>0`$, such that in the limit $`ϵ0`$,
$$\left(_{|Xa(t)|>b}\mathrm{\Psi }_{}(X,t,ϵ)_{_{\text{el}}}^2𝑑x\right)^{1/2}c(g)\text{e}^{\gamma (g)/ϵ^2}.$$
The mechanics of the nuclear configuration $`a(t)`$ is determined by classical dynamics in the effective potential $`E(X)`$.
Two theorems in Section 8 generalize this result. The first allows the time interval to grow as $`ϵ`$ tends to zero. The second allows more general initial conditions.
## 2 Coherent States and Classical Dynamics
In the construction of our approximation to the solution of the molecular Schrödinger equation, we need wave packets that describe the semiclassical dynamics of the heavy nuclei. In the present context, the semiclassical parameter is $`\mathrm{}=ϵ^2`$. We make use of a convenient set of coherent states (also called generalized squeezed states), that we express here in terms of the semiclassical parameter $`\mathrm{}`$.
We recall the definition of the coherent states $`\phi _j(A,B,\mathrm{},a,\eta ,X)`$ that are described in detail in . A more explicit, but more complicated definition is given in .
We adopt the standard multi-index notation. A multi-index $`j=(j_1,j_2,\mathrm{},j_d)`$ is a $`d`$-tuple of non-negative integers. We define $`|j|=_{k=1}^dj_k`$, $`X^j=X_1^{j_1}X_2^{j_2}\mathrm{}X_d^{j_d}`$,
$`j!=(j_1!)(j_2!)\mathrm{}(j_d!)`$, and $`D^j=\frac{^{|j|}}{(X_1)^{j_1}(X_2)^{j_2}\mathrm{}(X_d)^{j_d}}`$.
Throughout the paper we assume $`a\mathrm{I}\mathrm{R}^d`$, $`\eta \mathrm{I}\mathrm{R}^d`$ and $`\mathrm{}>0`$. We also assume that $`A`$ and $`B`$ are $`d\times d`$ complex invertible matrices that satisfy
$`A^tBB^tA`$ $`=`$ $`0,`$
$`A^{}B+B^{}A`$ $`=`$ $`2I.`$ (2.1)
These conditions guarantee that both the real and imaginary parts of $`BA^1`$ are symmetric. Furthermore, $`\text{Re}BA^1`$ is strictly positive definite, and $`\left(\text{Re}BA^1\right)^1=AA^{}`$.
Our definition of $`\phi _j(A,B,\mathrm{},a,\eta ,X)`$ is based on the following raising operators that are defined for $`m=1,\mathrm{\hspace{0.17em}2},\mathrm{},d`$.
$$𝒜_m(A,B,\mathrm{},a,\eta )^{}=\frac{1}{\sqrt{2\mathrm{}}}\left[\underset{n=1}{\overset{d}{}}\overline{B}_{nm}(X_na_n)i\underset{n=1}{\overset{d}{}}\overline{A}_{nm}(i\mathrm{}\frac{}{X_n}\eta _n)\right].$$
The corresponding lowering operators $`𝒜_m(A,B,\mathrm{},a,\eta )`$ are their formal adjoints.
These operators satisfy commutation relations that lead to the properties of the
$`\phi _j(A,B,\mathrm{},a,\eta ,X)`$ that we list below. The raising operators $`𝒜_m(A,B,\mathrm{},a,\eta )^{}`$ for
$`m=1,\mathrm{\hspace{0.17em}2},\mathrm{},d`$ commute with one another, and the lowering operators $`𝒜_m(A,B,\mathrm{},a,\eta )`$ commute with one another. However,
$$𝒜_m(A,B,\mathrm{},a,\eta )𝒜_n(A,B,\mathrm{},a,\eta )^{}𝒜_n(A,B,\mathrm{},a,\eta )^{}𝒜_m(A,B,\mathrm{},a,\eta )=\delta _{m,n}.$$
Definition For the multi-index $`j=0`$, we define the normalized complex Gaussian wave packet (modulo the sign of a square root) by
$`\phi _0(A,B,\mathrm{},a,\eta ,X)=\pi ^{d/4}\mathrm{}^{d/4}(det(A))^{1/2}`$
$`\times \mathrm{exp}\left\{(Xa),BA^1(Xa)/(2\mathrm{})+i\eta ,(Xa)/\mathrm{}\right\}.`$
Then, for any non-zero multi-index $`j`$, we define
$`\phi _j(A,B,\mathrm{},a,\eta ,)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{j!}}}\left(𝒜_1(A,B,\mathrm{},a,\eta )^{}\right)^{j_1}\left(𝒜_2(A,B,\mathrm{},a,\eta )^{}\right)^{j_2}\mathrm{}`$
$`\times \left(𝒜_d(A,B,\mathrm{},a,\eta )^{}\right)^{j_d}\phi _0(A,B,\mathrm{},a,\eta ,).`$
Properties 1. For $`A=B=I`$, $`\mathrm{}=1`$, and $`a=\eta =0`$, the $`\phi _j(A,B,\mathrm{},a,\eta ,)`$ are just the standard Harmonic oscillator eigenstates with energies $`|j|+d/2`$.
2. For each admissible $`A`$, $`B`$, $`\mathrm{}`$, $`a`$, and $`\eta `$, the set $`\{\phi _j(A,B,\mathrm{},a,\eta ,)\}`$ is an orthonormal basis for $`L^2(\mathrm{I}\mathrm{R}^d)`$.
3. In , the state $`\phi _j(A,B,\mathrm{},a,\eta ,X)`$ is defined as a normalization factor times
$$_j(A;\mathrm{}^{1/2}|A|^1(Xa))\phi _0(A,B,\mathrm{},a,\eta ,X).$$
Here $`_j(A;y)`$ is a recursively defined $`|j|^{\text{th}}`$ order polynomial in $`y`$ that depends on $`A`$ only through $`U_A`$, where $`A=|A|U_A`$ is the polar decomposition of $`A`$.
4. By scaling out the $`|A|`$ and $`\mathrm{}`$ dependence and using Remark 3 above, one can show that $`_j(A;y)\text{e}^{y^2/2}`$ is an (unnormalized) eigenstate of the usual Harmonic oscillator with energy $`|j|+d/2`$.
5. When the dimension $`d`$ is $`1`$, the position and momentum uncertainties of the
$`\phi _j(A,B,\mathrm{},a,\eta ,)`$ are $`\sqrt{(j+1/2)\mathrm{}}|A|`$ and $`\sqrt{(j+1/2)\mathrm{}}|B|`$, respectively. In higher dimensions, they are bounded by $`\sqrt{(|j|+d/2)\mathrm{}}A`$ and $`\sqrt{(|j|+d/2)\mathrm{}}B`$, respectively.
6. When we approximately solve the Schrödinger equation, the choice of the sign of the square root in the definition of $`\phi _0(A,B,\mathrm{},a,\eta ,)`$ is determined by continuity in $`t`$ after an arbitrary initial choice.
The following simple but very useful lemma is proven in .
###### Lemma 2.1
Let $`P_{|j|n}`$ denote the projection onto the span of the $`\phi _j(A,B,\mathrm{},a,\eta ,)`$ with $`|j|n`$.
$$(Xa)^mP_{|j|n}=P_{|j|n+|m|}(Xa)^mP_{|j|n},$$
(2.2)
and
$$(Xa)^mP_{|j|n}\left(\sqrt{2\mathrm{}}dA\right)^{|m|}\left(\frac{(n+|m|)!}{n!}\right)^{1/2}.$$
(2.3)
In the Born-Oppenheimer approximation, the semiclassical dynamics of the nuclei is generated by an effective potential given by a chosen isolated electronic eigenvalue $`E(X)`$ of the electronic hamiltonian $`h(X)`$, $`X\mathrm{I}\mathrm{R}^d`$. For a given effective potential $`E(X)`$ we describe the semiclassical dynamics of the nuclei by means of the time dependent basis constructed as follows:
By assumption H<sub>1</sub>, the potential $`E:\mathrm{\Xi }\mathrm{I}\mathrm{R}^d\mathrm{I}\mathrm{R}`$ is smooth and bounded below. Associated to $`E(X)`$, we have the following classical equations of motion:
$`\dot{a}(t)`$ $`=`$ $`\eta (t),`$
$`\dot{\eta }(t)`$ $`=`$ $`E(a(t)),`$
$`\dot{A}(t)`$ $`=`$ $`iB(t),`$ (2.4)
$`\dot{B}(t)`$ $`=`$ $`iE^{(2)}(a(t))A(t),`$
$`\dot{S}(t)`$ $`=`$ $`{\displaystyle \frac{\eta (t)^2}{2}}E(a(t)),`$
where $`E^{(2)}`$ denotes the Hessian matrix for $`E`$. We always assume the initial conditions $`A(0)`$, $`B(0)`$, $`a(0)`$, $`\eta (0)`$, and $`S(0)=0`$ satisfy (2).
The matrices $`A(t)`$ and $`B(t)`$ are related to the linearization of the classical flow through the following identities:
$`A(t)`$ $`=`$ $`{\displaystyle \frac{a(t)}{a(0)}}A(0)+i{\displaystyle \frac{a(t)}{\eta (0)}}B(0),`$
$`B(t)`$ $`=`$ $`{\displaystyle \frac{\eta (t)}{\eta (0)}}B(0)i{\displaystyle \frac{\eta (t)}{a(0)}}A(0).`$
Because $`E`$ is smooth and bounded below, there exist global solutions to the first two equations of the system (2.4) for any initial condition if $`\mathrm{\Xi }=\mathrm{I}\mathrm{R}^d`$. From this, it follows immediately that the remaining three equations of the system (2.4) have global solutions. If $`\mathrm{\Xi }\mathrm{I}\mathrm{R}^d`$, for any initial conditions, there exists a $`0<T\mathrm{}`$ so that solutions to the system (2.4) exist for any time $`t[0,T]`$. $`T`$ is finite if and only if the solution $`a(t)`$ corresponding to the chosen initial condition leaves the set $`\mathrm{\Xi }`$ in finite time.
Furthermore, it is not difficult to prove that conditions (2) are preserved by the flow.
The usefulness of our wave packets stems from the following important property . If we decompose the potential as
$$E(X)=W_a(X)+Z_a(X)W_a(X)+(E(X)W_a(X)),$$
where $`W_a(X)`$ denotes the second order Taylor approximation (with the obvious abuse of notation)
$$W_a(X)E(a)+E^{(1)}(a)(Xa)+E^{(2)}(a)(Xa)^2/2$$
then for all multi-indices $`j`$,
$`i\mathrm{}{\displaystyle \frac{}{t}}\left[\text{e}^{iS(t)/\mathrm{}}\phi _j(A(t),B(t),\mathrm{},a(t),\eta (t),X)\right]`$
$`=`$ $`\left({\displaystyle \frac{\mathrm{}^2}{2}}\mathrm{\Delta }+W_{a(t)}(X)\right)\left[\text{e}^{iS(t)/\mathrm{}}\phi _j(A(t),B(t),\mathrm{},a(t),\eta (t),X)\right],`$
if $`A(t)`$, $`B(t)`$, $`a(t)`$, $`\eta (t)`$, and $`S(t)`$ satisfy (2.4). In other words, our semiclassical wave packets $`\phi _j`$ exactly take into account the kinetic energy and quadratic part $`W_{a(t)}(X)`$ of the potential when propagated by means of the classical flow and its linearization around the classical trajectory selected by the initial conditions.
In the rest of the paper, whenever we write $`\phi _j(A(t),B(t),\mathrm{},a(t),\eta (t),X)`$, we tacitly assume that $`A(t),B(t),a(t),\eta (t)`$, and $`S(t)`$ are solutions to (2.4) with initial conditions satisfying (2).
## 3 The Born–Oppenheimer Expansion in Powers of $`ϵ`$
In this section we derive an explicit formal expansion in $`ϵ`$ for the solution to the molecular Schrödinger equation by means of a multiple scales analysis. This asymptotic analysis is similar to that performed, e.g., in . We discuss this in detail because we need more detailed information on the structure of successive terms in the expansion.
We start with the molecular Schrödinger equation for $`d`$ nuclear configuration dimensions,
$$iϵ^2\frac{\mathrm{\Psi }}{t}=\frac{ϵ^4}{2}\mathrm{\Delta }_X\mathrm{\Psi }+h(X)\mathrm{\Psi }.$$
(3.1)
We consider the isolated, multiplicity one, smooth eigenvalue $`E(X)`$ of $`h(X)`$ of hypothesis H<sub>1</sub>. For the moment we assume $`E(X)`$ is well defined on all of $`\mathrm{I}\mathrm{R}^d`$ rather than just on a subset $`\mathrm{\Xi }\mathrm{I}\mathrm{R}^d`$. Later we introduce a cut-off function to take care of the general case. We consider the solution $`a(t)`$, $`\eta (t)`$, $`A(t)`$, $`B(t)`$, and $`S(t)`$ to the system (2.4) of ODE’s. Then, we choose the phase of the eigenfuction $`\stackrel{~}{\mathrm{\Phi }}(X,t)`$ so that
$$\stackrel{~}{\mathrm{\Phi }}(X,t),(i\frac{}{t}+i\eta (t)_X)\stackrel{~}{\mathrm{\Phi }}(X,t)_{_{el}}=0.$$
(3.2)
This can always be done. See, e.g., .
The multiple scales analysis consists of separating the two length scales that are important in the nuclear variable $`X`$. The electron wave function is sensitive on an $`O(1)`$ scale in this variable, so $`X`$, or equivalently, $`Xa(t)`$ is relevant. The quantum mechanical fluctuations of nuclear wave function occur on an $`O(ϵ)`$ length scale, so $`(Xa(t))/ϵ`$ is also relevant. We replace the variable $`X`$ by both $`w=Xa(t)`$ and $`y=w/ϵ`$, and consider them as independent variables. This leads to the new problem of studying
$`iϵ^2{\displaystyle \frac{\widehat{\mathrm{\Psi }}}{t}}`$ $`=`$ $`[{\displaystyle \frac{ϵ^4}{2}}\mathrm{\Delta }_wϵ^3_w_y{\displaystyle \frac{ϵ^2}{2}}\mathrm{\Delta }_y+iϵ^2\eta (t)_w+iϵ\eta (t)_y`$ (3.3)
$`+[h(a(t)+w)E(a(t)+w)]+E(a(t)+ϵy)]\widehat{\mathrm{\Psi }}.`$
We easily check that if $`\widehat{\mathrm{\Psi }}(w,y,t)`$ solves (3.3) then $`\widehat{\mathrm{\Psi }}(Xa(t),(Xa(t))/ϵ,t)`$ solves (3.1).
We define $`\mathrm{\Phi }(w,t)=\stackrel{~}{\mathrm{\Phi }}(X,t)`$. Then (3.2) becomes
$$\mathrm{\Phi }(w,t),i\frac{}{t}\mathrm{\Phi }(w,t)_{_{el}}=0.$$
(3.4)
We seek solutions to (3.3) of the form
$$\widehat{\mathrm{\Psi }}(w,y,t)=e^{iS(t)/ϵ^2}e^{i\eta (t)y/ϵ}\varphi (w,y,t).$$
This requires $`\varphi (w,y,t)`$ to satisfy
$`iϵ^2{\displaystyle \frac{\varphi }{t}}`$ $`=`$ $`[{\displaystyle \frac{ϵ^4}{2}}\mathrm{\Delta }_wϵ^3_w_y+({\displaystyle \frac{ϵ^2}{2}}\mathrm{\Delta }_y+{\displaystyle \frac{ϵ^2}{2}}E^{(2)}(a(t))y^2)`$
$`+\left[h(a(t)+w)E(a(t)+w)\right]`$
$`+(E(a(t)+ϵy)E(a(t))ϵE^{(1)}(a(t))yϵ^2E^{(2)}(a(t)){\displaystyle \frac{y^2}{2!}})]\varphi ,`$
where here and below we make use of the shorthand notation
$$E^{(m)}(x)\frac{y^m}{m!}=\underset{\{k:|k|=m\}}{}\frac{(D^kE)(x)y^k}{k!},$$
in the usual multi-index notation. We next assume that $`\varphi (w,y,t)`$ has an expansion of the form
$$\varphi (w,y,t)=\varphi _0(w,y,t)+ϵ\varphi _1(w,y,t)+ϵ^2\varphi _2(w,y,t)+\mathrm{}$$
We further decompose each $`\varphi _n`$ as
$$\varphi _n(w,y,t)=g_n(w,y,t)\mathrm{\Phi }(w,t)+\varphi _n^{}(w,y,t),$$
by projecting into the $`\mathrm{\Phi }(w,t)`$ direction and into the orthogonal directions in $`_{el}`$.
We substitute this expansion into (3) and equate terms of the corresponding powers of $`ϵ`$.
Order 0. The zeroth order terms require
$$\left[h(a(t)+w)E(a(t)+w)\right]\varphi _0(w,y,t)=0.$$
This forces
$$\varphi _0^{}(w,y,t)=0.$$
Order 1. The first order terms require
$$\left[h(a(t)+w)E(a(t)+w)\right]\varphi _1(w,y,t)=0.$$
This forces
$$\varphi _1^{}(w,y,t)=0.$$
Order 2. The second order terms require
$$i\frac{\varphi _0}{t}=\left(\frac{1}{2}\mathrm{\Delta }_y+E^{(2)}(a(t))\frac{y^2}{2!}\right)\varphi _0+\left[h(a(t)+w)E(a(t)+w)\right]\varphi _2.$$
We separately examine the components of this equation in the $`\mathrm{\Phi }`$ direction and in the orthogonal directions. By (3.4), this yields the two conditions
$$i\frac{g_0}{t}=\left(\frac{1}{2}\mathrm{\Delta }_y+E^{(2)}(a(t))\frac{y^2}{2!}\right)g_0,$$
(3.6)
and
$$\left[h(a(t)+w)E(a(t)+w)\right]\varphi _2=ig_0\frac{\mathrm{\Phi }}{t}.$$
(3.7)
We arbitrarily choose $`g_0`$ to be the following $`w`$–independent particular solution of (3.6):
$$g_0(w,y,t)=ϵ^{d/2}\underset{|j|J}{}c_{0,j}\phi _j(A(t),B(t),1,0,0,y),$$
(3.8)
where $`c_{0,j}=c_j`$ is determined by the initial conditions.
We let the Hilbert space $`_{\text{el}}^{}`$ be the subspace of $`_{\text{el}}`$ orthogonal to $`\mathrm{\Phi }(w,t)`$. The restriction of $`\left[h(a(t)+w)E(a(t)+w)\right]`$ to $`_{\text{el}}^{}`$ is invertible, and we denote the inverse by $`r(w,t)=\left[h(a(t)+w)E(a(t)+w)\right]_r^1`$. With this notation, equation (3.7) forces
$`\varphi _2^{}(w,y,t)`$ $`=`$ $`ig_0(w,y,t)r(w,t){\displaystyle \frac{\mathrm{\Phi }}{t}}(w,t)`$ (3.9)
$`=`$ $`ϵ^{d/2}{\displaystyle \underset{|j|J}{}}d_{2,j}(w,t)\phi _j(A(t),B(t),1,0,0,y),`$
where
$$d_{2,j}(w,t)=c_{0,j}r(w,t)\frac{\mathrm{\Phi }}{t}(w,t)$$
(3.10)
is $`_{\text{el}}`$–valued.
Order 3. The third order terms require
$`i{\displaystyle \frac{\varphi _1}{t}}`$ $`=`$ $`\left({\displaystyle \frac{1}{2}}\mathrm{\Delta }_y+E^{(2)}(a(t)){\displaystyle \frac{y^2}{2!}}\right)\varphi _1`$
$`_w_y\varphi _0+E^{(3)}(a(t)){\displaystyle \frac{y^3}{3!}}\varphi _0+\left[h(a(t)+w)E(a(t)+w)\right]\varphi _3.`$
We separately examine the components of this equation in the $`\mathrm{\Phi }`$ direction and in the orthogonal directions. By (3.4), this yields the two conditions
$`i{\displaystyle \frac{g_1}{t}}\left({\displaystyle \frac{1}{2}}\mathrm{\Delta }_y+E^{(2)}(a(t)){\displaystyle \frac{y^2}{2!}}\right)g_1`$
$`=(_yg_0)\mathrm{\Phi },_w\mathrm{\Phi }+E^{(3)}(a(t)){\displaystyle \frac{y^3}{3!}}g_0,`$ (3.11)
and
$$\left[h(a(t)+w)E(a(t)+w)\right]\varphi _3=ig_1\frac{\mathrm{\Phi }}{t}+(_yg_0)(P_{}_w\mathrm{\Phi }),$$
(3.12)
where $`P_{}(w,t)`$ is the projection in $`_{el}`$ onto $`_{\text{el}}^{}`$.
The solution to (3.11) with $`g_1(w,y,0)=0`$ can be written as
$$g_1(w,y,t)=ϵ^{d/2}\underset{|j|J+3}{}c_{1,j}(w,t)\phi _j(A(t),B(t),1,0,0,y),$$
for some coefficients $`c_{1,j}(w,t)`$.
Equation (3.12) determines
$`\varphi _3^{}(w,y,t)`$ $`=`$ $`r(w,t)\left(ig_1(w,y,t){\displaystyle \frac{\mathrm{\Phi }}{t}}(w,t)+(_yg_0)(w,y,t)(P_{}(w,t)_w\mathrm{\Phi }(w,t))\right)`$
$`=`$ $`ϵ^{d/2}{\displaystyle \underset{|j|J+3}{}}d_{3,j}(w,t)\phi _j(A(t),B(t),1,0,0,y),`$
where
$`d_{3,j}(w,t)`$ $`=`$ $`i\left(r(w,t)\dot{\mathrm{\Phi }}(w,t)\right)c_{1,j}(w,t)`$
$`+{\displaystyle \underset{|q|J}{}}r(w,t)(P_{}_w\mathrm{\Phi })(w,t)\phi _j,_y\phi _qc_{0,q}(w,t).`$
Here and below $`\dot{}{\displaystyle \frac{}{t}}`$.
Order $`𝐧`$. The $`n^{\text{th}}`$ order terms require
$`i{\displaystyle \frac{\varphi _{n2}}{t}}`$ $`=`$ $`\left({\displaystyle \frac{1}{2}}\mathrm{\Delta }_y+E^{(2)}(a(t)){\displaystyle \frac{y^2}{2!}}\right)\varphi _{n2}{\displaystyle \frac{1}{2}}\mathrm{\Delta }_w\varphi _{n4}_w_y\varphi _{n3}`$
$`+{\displaystyle \underset{m=3}{\overset{n}{}}}E^{(m)}(a(t)){\displaystyle \frac{y^m}{m!}}\varphi _{nm}+\left[h(a(t)+w)E(a(t)+w)\right]\varphi _n.`$
The components of this equation in the $`\mathrm{\Phi }(w,t)`$ direction require
$`i{\displaystyle \frac{g_{n2}}{t}}\left({\displaystyle \frac{1}{2}}\mathrm{\Delta }_y+{\displaystyle \frac{1}{2!}}E^{(2)}(a(t))y^2\right)g_{n2}`$
$`={\displaystyle \frac{1}{2}}\mathrm{\Delta }_wg_{n4}\mathrm{\Phi },_w\mathrm{\Phi }(_wg_{n4}){\displaystyle \frac{1}{2}}\mathrm{\Phi },\mathrm{\Delta }_w\mathrm{\Phi }g_{n4}`$
$`_w_yg_{n3}\mathrm{\Phi },_w\mathrm{\Phi }(_yg_{n3})+{\displaystyle \underset{m=3}{\overset{n}{}}}E^{(m)}(a(t)){\displaystyle \frac{y^m}{m!}}g_{nm}`$
$`{\displaystyle \frac{1}{2}}\mathrm{\Phi },\mathrm{\Delta }_w\varphi _{n4}^{}\mathrm{\Phi },_w_y\varphi _{n3}^{}+i{\displaystyle \frac{\mathrm{\Phi }}{t}},\varphi _{n2}^{}.`$ (3.13)
Note that the last term has been transformed from $`i\mathrm{\Phi },{\displaystyle \frac{\varphi _{n2}^{}}{t}}`$ to $`i{\displaystyle \frac{\mathrm{\Phi }}{t}},\varphi _{n2}^{}`$. The equivalence of these expressions follows from differentiation of $`\mathrm{\Phi },\varphi _{n2}^{}=\mathrm{\hspace{0.17em}0}`$ with respect to $`t`$.
The components orthogonal to $`\mathrm{\Phi }(w,t)`$ require
$`\left[h(a(t)+w)E(a(t)+w)\right]\varphi _n`$
$`=P_{}\left(i{\displaystyle \frac{\varphi _{n2}^{}}{t}}\right)+\left({\displaystyle \frac{1}{2}}\mathrm{\Delta }_y{\displaystyle \frac{1}{2!}}E^{(2)}(a(t))y^2\right)\varphi _{n2}^{}{\displaystyle \underset{m=3}{\overset{n}{}}}E^{(m)}(a(t)){\displaystyle \frac{y^m}{m!}}\varphi _{nm}^{}`$
$`+{\displaystyle \frac{1}{2}}P_{}\mathrm{\Delta }_w\varphi _{n4}^{}+(P_{}_w\mathrm{\Phi })(_wg_{n4})+{\displaystyle \frac{1}{2}}(P_{}\mathrm{\Delta }_w\mathrm{\Phi })g_{n4}`$
$`+P_{}_w_y\varphi _{n3}^{}+(P_{}_w\mathrm{\Phi })(_yg_{n3})+i{\displaystyle \frac{\mathrm{\Phi }}{t}}g_{n2}.`$ (3.14)
Equation (3.14) determines $`\varphi _n^{}(w,y,t)`$ by an application of $`\left[h(a(t)+w)E(a(t)+w)\right]_r^1`$.
It is easily checked that the solution to (3.13) with $`g_{n2}(w,y,0)=\mathrm{\hspace{0.17em}0}`$ has the form
$$g_{n2}(w,y,t)=ϵ^{d/2}\underset{|j|J+3n6}{}c_{n2,j}(w,t)\phi _j(A(t),B(t),1,0,0,y),$$
(3.15)
for some coefficients $`c_{n2,j}(w,t)`$, and that the $`y`$ dependence of the vector $`\varphi _n^{}`$ has the same form, with other coefficients depending on $`(w,t)`$, i.e.,
$$\varphi _n^{}(w,y,t)=ϵ^{d/2}\underset{|j|J+3n6}{}d_{n,j}(w,t)\phi _j(A(t),B(t),1,0,0,y),$$
(3.16)
where the $`d_{n,j}(w,t)`$ take their values in the electronic Hilbert space.
Equations (3.13) and (3.14) determine $`c_{n2,j}`$ and $`d_{n,j}`$. When recursively solving these equations, we must determine $`d_{n,j}`$ before $`c_{n,j}`$ because the right hand side of (3.13) (with $`n2`$ replaced by $`n`$) contains $`\varphi _n^{}`$.
The solution to (3.14) in terms of the $`d_{n,j}`$, is
$$d_{n,j}(w,t)=\underset{i=1}{\overset{8}{}}\mathrm{\Delta }_i(w,t),$$
(3.17)
where
$`\mathrm{\Delta }_1(w,t)`$ $`=`$ $`ir(w,t)P_{}(w,t)\dot{d}_{n2,j}(w,t)`$
$`\mathrm{\Delta }_2(w,t)`$ $`=`$ $`{\displaystyle \underset{3|m|n}{}}{\displaystyle \frac{(D^mE)(a(t))}{m!}}{\displaystyle \underset{|q|J+3(n|m|2)}{}}\phi _j,y^m\phi _qr(w,t)d_{n|m|,q}(w,t)`$
$`\mathrm{\Delta }_3(w,t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}r(w,t)P_{}(w,t)(\mathrm{\Delta }_wd_{n4,j})(w,t)`$
$`\mathrm{\Delta }_4(w,t)`$ $`=`$ $`r(w,t)P_{}(w,t)(_w\mathrm{\Phi })(_wc_{n4,j})(w,t)`$
$`\mathrm{\Delta }_5(w,t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}r(w,t)P_{}(w,t)(\mathrm{\Delta }_w\mathrm{\Phi })c_{n4,j}(w,t)`$
$`\mathrm{\Delta }_6(w,t)`$ $`=`$ $`{\displaystyle \underset{|q|J+3(n5)}{}}r(w,t)P_{}(w,t)\phi _j,_y\phi _q(_wd_{n3,q})(w,t)`$
$`\mathrm{\Delta }_7(w,t)`$ $`=`$ $`{\displaystyle \underset{|q|J+3(n3)}{}}r(w,t)P_{}(w,t)(_w\mathrm{\Phi })\phi _j,_y\phi _qc_{n3,q}(w,t)`$
$`\mathrm{\Delta }_8(w,t)`$ $`=`$ $`ir(w,t)P_{}(w,t)\dot{\mathrm{\Phi }}(w,t)c_{n2,j}(w,t).`$
Similarly, the solution to (3.13) in terms of the $`c_{n,j}`$ is obtained by integration with respect to $`t`$ of $`i\dot{c}_{n,j}(w,t)`$, where
$$i\dot{c}_{n,j}(w,t)=\underset{i=1}{\overset{9}{}}\mathrm{\Gamma }_i(w,t),$$
(3.18)
where
$`\mathrm{\Gamma }_1(w,t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\mathrm{\Delta }_wc_{n2,j})(w,t)`$
$`\mathrm{\Gamma }_2(w,t)`$ $`=`$ $`\mathrm{\Phi },_w\mathrm{\Phi }(_wc_{n2,j})(w,t)`$
$`\mathrm{\Gamma }_3(w,t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Phi },\mathrm{\Delta }_w\mathrm{\Phi }c_{n2,j}(w,t)`$
$`\mathrm{\Gamma }_4(w,t)`$ $`=`$ $`{\displaystyle \underset{|q|J+3(n1)}{}}\phi _j,_y\phi _q(_wc_{n1,q})(w,t)`$
$`\mathrm{\Gamma }_5(w,t)`$ $`=`$ $`{\displaystyle \underset{|q|J+3(n1)}{}}\mathrm{\Phi },_w\mathrm{\Phi }\phi _j,_y\phi _qc_{n1,q}(w,t)`$
$`\mathrm{\Gamma }_6(w,t)`$ $`=`$ $`{\displaystyle \underset{3|m|n+2}{}}{\displaystyle \underset{|q|J+3(n+2m)}{}}{\displaystyle \frac{(D^mE)(a(t))}{m!}}\phi _j,y^m\phi _qc_{n+2m,q}(w,t)`$
$`\mathrm{\Gamma }_7(w,t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Phi },(\mathrm{\Delta }_wd_{n2,j})(w,t)`$
$`\mathrm{\Gamma }_8(w,t)`$ $`=`$ $`{\displaystyle \underset{|q|J+3(n3)}{}}\phi _j,_y\phi _q\mathrm{\Phi },(_wd_{n2,q})(w,t)`$
$`\mathrm{\Gamma }_9(w,t)`$ $`=`$ $`i\dot{\mathrm{\Phi }},d_{n,j}(w,t).`$
## 4 The Main Result
We introduce a $`C^{\mathrm{}}`$ real valued cut-off function $`F:\mathrm{I}\mathrm{R}^d\mathrm{I}\mathrm{R}`$ that equals $`1`$ in a neighborhood of the origin and equals zero away from the origin. More precisely, we choose $`0<b_0<b_1<\mathrm{}`$, such that
$$\text{supp}(_{w_i}F)(w)\{w\mathrm{I}\mathrm{R}^d:b_0<|w|<b_1\},$$
for any $`i\{1,\mathrm{},d\}`$, and such that for any $`t\mathrm{\Omega }`$, all quantities appearing in the above expansion are well defined for $`w\mathrm{I}\mathrm{R}^d`$ with $`|w|<b_1`$. Here $`\mathrm{\Omega }`$ is a particular simply connected open complex neighborhood of the real interval $`[0,T]`$ that we construct in Section 5 under hypotheses H<sub>0</sub> and H<sub>1</sub>.
We define our approximate solution to (3.1) at order $`N`$ by the following expression:
$`\widehat{\mathrm{\Psi }}_N(w,y,t)`$
$`=`$ $`F(w)\text{e}^{iS(t)/ϵ^2}\text{e}^{i\eta (t)y/ϵ}\left({\displaystyle \underset{n=0}{\overset{N}{}}}ϵ^ng_n(w,y,t)\mathrm{\Phi }(w,t)+{\displaystyle \underset{n=2}{\overset{N+2}{}}}ϵ^n\varphi _n^{}(w,y,t)\right).`$
We prove in Section 7.2 that this quantity agrees with an exact solution up to an error whose norm is bounded by $`ϵ^N`$ for $`t[0,T]`$.
We emphasize that once the molecular hamiltonian $`h(X)`$ and its spectral data $`E(X)`$, $`\mathrm{\Phi }(X)`$ are given, the only arbitrary input of the above derived expansion consists of the set of coefficients $`c_{0,j}`$, $`|j|J`$. We note that at time $`t=0`$, we have $`c_{n,j}(0,w)0`$ for all $`n1`$. Thus, at $`t=0`$, the approximation reduces to
$`\widehat{\mathrm{\Psi }}_N(w,y,0)`$
$`=`$ $`F(w)\text{e}^{iS(0)/ϵ^2}\text{e}^{i\eta (0)y/ϵ}\left(g_0(0,y,0)\mathrm{\Phi }(w,0)+{\displaystyle \underset{n=2}{\overset{N+2}{}}}ϵ^n\varphi _n^{}(w,y,0)\right).`$
This expression is completely determined by $`g_0(0,y,0)`$, the nuclear part of the wave function parallel to the chosen electronic level at time $`0`$.
As is usual in the study of adiabatic problems, in order to get accurate information on the evolution of an initial wave function associated with a specific electronic level, one needs to include a higher order component perpendicular to that electronic level. This higher order part is completely determined by the parallel part. Here it is given (up to phase and cut-off functions) at time $`0`$ by $`{\displaystyle \underset{n=2}{\overset{N+2}{}}}ϵ^n\varphi _n^{}(w,y,0)`$.
We now state our main Theorem:
###### Theorem 4.1
Assume hypotheses H<sub>0</sub> and H<sub>1</sub> and consider the above construction. For all sufficiently small choices of $`g>0`$, there exist $`C(g)>0`$ and $`\mathrm{\Gamma }(g)>0`$ such that, for $`N(ϵ)=[[g^2/ϵ^2]]`$, the vector $`\mathrm{\Psi }_{}(X,t,ϵ)=\widehat{\mathrm{\Psi }}_{N(ϵ)}(Xa(t),(Xa(t))/ϵ,t)`$ satisfies
$$\text{e}^{itH(ϵ)/ϵ^2}\mathrm{\Psi }_{}(X,0,ϵ)\mathrm{\Psi }_{}(X,t,ϵ)_{L^2(\text{IR}^d,_{\text{el}})}C(g)\text{e}^{\mathrm{\Gamma }(g)/ϵ^2},$$
for all $`t[0,T]`$, as $`ϵ0`$.
Moreover, we have the following exponential localization result. For any $`b>0`$ and a sufficiently small choice of $`g>0`$ (that depends on $`b`$), there exist $`c(g)`$ and $`\gamma (g)>0`$, such that
$$\left(_{|xa(t)|>b}\mathrm{\Psi }_{}(X,t,ϵ)_{_{\text{el}}}^2𝑑x\right)^{1/2}c(g)\text{e}^{\gamma (g)/ϵ^2},$$
for all $`t[0,T]`$, as $`ϵ0`$.
The strategy of the proof is as follows: We consider the approximation $`\widehat{\mathrm{\Psi }}_N(Xa(t),(Xa(t))/ϵ,t)`$ and the exact solution to the Schrödinger equation with the same initial conditions. We estimate the norm of the error (that is the difference between these two quantities) as a function of both $`N`$ and $`ϵ`$. Apart from some subtleties, the norm of the error is bounded by $`Cϵ^N(\tau N^{1/2})^N`$, for some constants $`C`$ and $`\tau >0`$. We minimize the error estimate over all choices of $`N`$. This yields $`Ng^2/ϵ^2`$, for sufficiently small $`g>0`$, and an estimate of order $`\text{e}^{\mathrm{\Gamma }(g)/ϵ^2}`$ for the norm of the error.
We prove two extensions of this result in Section 8. In the first extension, we consider the validity of our approximation on the Ehrenfest time scale, i.e., when $`T=T(ϵ)\mathrm{ln}(1/ϵ)`$. In the second extension, we study the dependence of our construction on $`J`$, in order to extend our main result to a wider class of initial conditions. We refer the reader to Section 8 for the precise statements.
## 5 Analyticity Properties
Our estimates depend on analyticity in $`t\mathrm{\Omega }`$ of the vectors $`c_n(w,t)l^2(\text{N}^d,\mathrm{I}\mathrm{C})`$ and $`d_n(w,t)l^2(\text{N}^d,_{\text{el}})`$, where $`\mathrm{\Omega }`$ is the particular simply connected open complex neighborhood of the real interval $`[0,T]`$ mentioned at the beginning of Section 4.
To construct $`\mathrm{\Omega }`$, we begin with several observations. Our hypotheses imply that the eigenvalue $`E(X)`$ is analytic in $`\mathrm{\Sigma }_\delta ^{}`$, so the solutions $`a(t)`$, $`\eta (t)`$, $`A(t)`$, $`B(t)`$, and $`S(t)`$ are well defined for all $`t[0,T]`$. Moreover, by standard arguments , these functions all have analytic continuations from $`[0,T]`$ to a simply connected open set $`\mathrm{\Omega }_1`$ that contains $`[0,T]`$. We assume without loss of generality that $`\mathrm{\Omega }_1=\overline{\mathrm{\Omega }_1}`$, where $`\overline{\mathrm{\Omega }_1}`$ denotes the conjugate of $`\mathrm{\Omega }_1`$.
We note that $`A^{}(t)`$ and $`B^{}(t)`$ also have analytic continuations from $`[0,T]`$ to $`\mathrm{\Omega }_1`$. To see this for $`A^{}(t)`$, note that for $`t[0,T]`$, $`A^{}(t)=A^{}(\overline{t})`$, and $`A^{}(\overline{t})`$ has an analytic continutation to $`\overline{\mathrm{\Omega }_1}`$. The argument for $`B^{}(t)`$ is similar.
It now follows easily from the definitions that for each $`X`$, $`\phi _j(A(t),B(t),ϵ^2,a(t),\eta (t),X)`$ and $`\overline{\phi _j(A(t),B(t),ϵ^2,a(t),\eta (t),X)}`$ have analytic continuations from $`[0,T]`$ to some simply connected open set $`\mathrm{\Omega }_2`$. For $`t[0,T]`$, the real part of $`B(t)A(t)^1`$ is strictly positive. This positivity will remain true for the real part of the analytic continuation of $`B(t)A(t)^1`$ on some simply connected subset $`\mathrm{\Omega }\mathrm{\Omega }_1\mathrm{\Omega }_2`$ that contains $`[0,T]`$. We assume without loss of generality that $`\mathrm{\Omega }=\overline{\mathrm{\Omega }}`$ and we can assume that $`\mathrm{\Omega }`$ has the form $`\{t:a<\text{Re}t<b\text{ and }|\text{Im}t|<c\}`$ where $`a>c>0`$ and $`b>T+c`$. It follows that for $`t\mathrm{\Omega }`$, both $`\phi _j(A(t),B(t),ϵ^2,a(t),\eta (t),x)`$ and $`\overline{\phi _j(A(t),B(t),ϵ^2,a(t),\eta (t),x)}`$ have analytic continuations from $`[0,T]`$ to $`\mathrm{\Omega }`$ as elements of $`L^2(\mathrm{I}\mathrm{R}^d)`$.
Using these results and carefully examining the constructions of the vectors $`c_n(w,t)`$ and $`d_n(w,t)`$, we see that they are analytic in $`t`$ for $`t\mathrm{\Omega }`$, also.
Our hypotheses on $`h()`$ and the above results also show that each of the following quantities is analytic in $`t`$ for $`t\mathrm{\Omega }`$ and each fixed $`w\mathrm{\Sigma }_\delta \mathrm{I}\mathrm{C}^d`$, for sufficiently small $`\delta `$:
$`r(w,t)=\left[h(a(t)+w)E(a(t)+w)\right]_r^1,`$
$`\mathrm{\Phi }(w,t),`$
$`(D_w^\alpha \mathrm{\Phi })(w,t),\text{for }|\alpha |2,`$
$`D_w^\alpha E(a(t)),\text{for all}\alpha `$
$`P_{}(w,t).`$
By explicit computation of the phase corresponding to (3.2) it is easy to check that $`\mathrm{\Phi }(w,t)`$ and its derivatives are also analytic for $`t\mathrm{\Omega }`$.
Moreover, if $`f_i(w,t)`$, ($`i`$ in some finite set) represents any of these quantities, $`f_i`$ is analytic in $`w\mathrm{\Sigma }_\delta `$, for any fixed $`t\mathrm{\Omega }`$. Thus, by the Cauchy integral formula, we can assume that the following bounds hold (with the appropriate norm in each case):
$$(D_w^\alpha f_i)(w,t)c_iG_i^{|\alpha |}\frac{\alpha !}{(1+|\alpha |)^{d+1}},$$
(5.1)
for some $`c_i`$, $`G_i`$, $`w\mathrm{\Sigma }_\delta `$, and $`\alpha `$ ranges over all multi-indices.
We can assume here that all $`G_iD_2`$ for some constant
$$D_21,$$
and we associate the prefactors $`c_i`$ in (5.1) with the different functions according to the rules
$$\begin{array}{ccccccc}c_1& & rP_{}& & c_2& & \mathrm{\Phi }\\ c_3& & \dot{\mathrm{\Phi }}& & c_4& & _w\mathrm{\Phi }\\ c_5& & \mathrm{\Delta }_w\mathrm{\Phi }& & c_6& & E\\ c_4^{}& & \mathrm{\Phi },_w\mathrm{\Phi }& & c_5^{}& & \mathrm{\Phi },\mathrm{\Delta }_w\mathrm{\Phi }.\end{array}$$
## 6 Structure and Estimates of the $`c_n(w,t)`$ and $`d_n(w,t)`$
In this section, we decompose the functions $`g_n`$ and $`\varphi _n^{}`$ of Section 3 into pieces, each of which satisfies various estimates.
Throughout this section, all $`w`$–dependent quantities are defined for $`w`$ in the support of the cut-off function $`F`$. Furthermore, all the results of this section are claimed to hold only on the support of $`F`$.
Our decompositions of $`g_n(w,y,t)`$ and $`\varphi _n^{}(w,y,t)`$ have the following forms:
$`g_n(w,y,t)`$
$`=`$ $`ϵ^{d/2}{\displaystyle \underset{\beta _{n,1}}{}}{\displaystyle \underset{pn}{}}{\displaystyle \underset{|l|+kp+\frac{n}{2}}{}}{\displaystyle \underset{|j|J+n+2(p|l|k)}{}}c_{n,p,l,k,\beta ,j}(w,t)\phi _j(A(t),B(t),1,0,0,y).`$
and
$`\varphi _n^{}(w,y,t)=`$ (6.2)
$`ϵ^{d/2}{\displaystyle \underset{\beta _{n,2}}{}}{\displaystyle \underset{pn1}{}}{\displaystyle \underset{|l|+kp+\frac{n1}{2}}{}}{\displaystyle \underset{|j|J+(n1)+2(p|l|k)}{}}d_{n,p,l,k,\beta ,j}(w,t)\phi _j(A(t),B(t),1,0,0,y).`$
In (6), $`n`$, $`k`$ and $`p`$ are non-negative integers; $`j`$ and $`l`$ are multi-indices; and the index $`\beta `$ runs over a finite set $`_{n,1}`$. The number $`J`$ is fixed by the initial conditions. Each $`c_{n,p,l,k,\beta ,j}`$ is a complex valued function.
In (6.2), $`n2`$, $`k`$ and $`p`$ are non-negative integers; $`j`$ and $`l`$ are multi-indices; and the index $`\beta `$ runs over a finite set $`_{n,2}`$. Each $`d_{n,p,l,k,\beta ,j}(w,t)`$ takes values in $`_{el}`$.
We let $`c_{n,p,l,k,\beta }(w,t)`$ and $`d_{n,p,l,k,\beta }(w,t)`$ respectively denote vectors in $`l^2(\text{N}^d,\mathrm{I}\mathrm{C})`$ and
$`l^2(\text{N}^d,_{el})`$ whose components are $`c_{n,p,l,k,\beta ,j}(w,t)`$ and $`d_{n,p,l,k,\beta ,j}(w,t)`$.
The crucial step in the proof of Theorem 4.1 is the following:
###### Proposition 6.1
There is a recursive construction of the coefficients $`c_{n,p,l,k,\beta ,j}(w,t)`$ and $`d_{n,p,l,k,\beta ,j}(w,t)`$ for $`w`$ on the support of $`F`$.
The indices for $`c_{n,p,l,k,\beta ,j}(w,t)`$ are non-negative and satisfy
$`\beta `$ $``$ $`_{n,1}`$
$`p`$ $``$ $`n,`$
$`|l|+k`$ $``$ $`p+{\displaystyle \frac{n}{2}},`$
$`|j|`$ $``$ $`J+n+2(p|l|k).`$
The indices for $`d_{n,p,l,k,\beta ,j}(w,t)`$ are non-negative and satisfy
$`n`$ $``$ $`2`$
$`\beta `$ $``$ $`_{n,2},`$
$`p`$ $``$ $`n1`$
$`|l|+k`$ $``$ $`p+{\displaystyle \frac{n1}{2}}`$
$`|j|`$ $``$ $`J+(n1)+2(p|l|k).`$
Moreover, the following conditions are satisfied:
i) For any $`n>0`$, $`c_{n,0,l,k,\beta ,j}(w,t)=\mathrm{\hspace{0.17em}0}`$.
ii) There exists $`K_0>0`$, such that the number of terms in both of the sums (6) and (6.2) is bounded by $`\text{e}^{K_0n}`$.
iii) For $`t\mathrm{\Omega }`$, let $`\text{dist}(t)`$ be the distance from $`t`$ to the complement of $`\mathrm{\Omega }`$. The coefficients $`c_{n,p,l,k,\beta }(w,t)`$ and $`d_{n,p,l,k,\beta }(w,t)`$ are analytic for $`t\mathrm{\Omega }`$, and there exist constants $`D_1`$ and $`D_2`$, such that
$`(D_w^\alpha c_{n,p,l,k,\beta })(w,t)`$
$``$ $`D_1D_2^{|\alpha |+|l|+4n}{\displaystyle \frac{(\alpha +l)!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p}{p!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}\left[{\displaystyle \frac{(J+n+2(p|l|k))!}{J!}}\right]^{1/2},`$
and
$`(D_w^\alpha d_{n,p,l,k,\beta })(w,t)`$ (6.4)
$`D_1D_2^{|\alpha |+|l|+4(n1)}{\displaystyle \frac{(\alpha +l)!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p}{p!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}\left[{\displaystyle \frac{(J+(n1)+2(p|l|k))!}{J!}}\right]^{1/2}.`$
Remark The complicated estimates (6.1) and (6.4) are motivated by estimates used in semiclassical approximations and adiabatic approximations. The factors on the right hand sides that explicitly involve $`J`$, $`n`$, and $`p`$ occur in the semiclassical paper . The factors that involve $`\alpha `$ and $`l`$ appear in the adiabatic paper . The factors that involve $`k`$ occur in a proof of the adiabatic results of that are based on Cauchy estimates instead of Nenciu’s lemma (that we generalize below as Lemma 6.4). We were unable to prove Proposition 6.1 without using a combination of all of these techniques. We estimate adiabatic error terms by using Nenciu’s approach in the $`w`$ variable and Cauchy estimates in the $`t`$ variable.
### 6.1 The Toolbox
To prove Proposition 6.1, we repeatedly use the following very handy lemmas, whose proofs are given in Section 9. The first two lemmas deal with basic properties of analytic functions of one variable and are consequences of the Cauchy integral formula.
###### Lemma 6.1
For $`k=0`$, define $`k^k=1`$. Suppose $`g`$ is an analytic vector–valued function on the strip $`S_\delta =\{t:|\text{Im}t|<\delta \}`$. If $`g`$ satisfies
$$g(t)Ck^k(\delta |\text{Im}t|)^k,$$
for some $`k0`$, then $`g^{}`$ satisfies
$$g^{}(t)C(k+1)^{k+1}(\delta |\text{Im}t|)^{k1},$$
for all $`tS_\delta `$.
Lemma 6.1 has a generalization to regions other than infinite strips. The generalization is needed if one wishes to study problems where analyticity holds only in a neighborhood of a finite time interval. The proof of the generalized lemma is similar to that of Lemma 6.1, but involves slightly more complicated geometry. The precise statement is the following:
###### Lemma 6.2
For $`k=0`$, define $`k^k=1`$. Suppose $`g`$ is an analytic vector–valued function in an open region $`\mathrm{\Omega }IC`$. For $`t\mathrm{\Omega }`$, let $`\text{dist}(t)`$ be the distance from $`t`$ to $`\mathrm{\Omega }^C`$, the complement of $`\mathrm{\Omega }`$. If $`g`$ satisfies
$$g(t)Ck^k(\text{dist}(t))^k,$$
for all $`t\mathrm{\Omega }`$ and some $`k0`$, then $`g^{}`$ satisfies
$$g^{}(t)C(k+1)^{k+1}(\text{dist}(t))^{k1},$$
for all $`t\mathrm{\Omega }`$.
The next lemma gives estimates on indefinite integrals of certain analytic functions under stronger assumptions on the domain $`\mathrm{\Omega }`$.
###### Lemma 6.3
Suppose $`f`$ is an analytic vector–valued function in an open region $`\mathrm{\Omega }IC`$. For $`t\mathrm{\Omega }`$, let $`\text{dist}(t)`$ be the distance from $`t`$ to $`\mathrm{\Omega }^C`$. We assume the domain is star-shaped with respect to the origin and that the origin is the most distant point to $`\mathrm{\Omega }^C`$, i.e., $`\text{dist}(0)\text{dist}(t)`$, for all $`t\mathrm{\Omega }`$. Moreover, we assume that $`\text{dist}(t)`$ is monotone decreasing along any line emanating from the origin. If $`f`$ satisfies
$$f(t)C|t|^p(\text{dist}(t))^k,$$
for all $`t\mathrm{\Omega }`$ and some $`k0`$, then $`_0^tf(s)𝑑s`$ satisfies
$$_0^tf(s)𝑑sC\frac{|t|^{p+1}}{p+1}(\text{dist}(t))^k,$$
for all $`t\mathrm{\Omega }`$.
Remark: In our situation, examples of sets $`\mathrm{\Omega }`$ we can use that satisfy the conditions of Lemma 6.3 are infinite symmetrical horizontal strips or the rectangular regions chosen in Section 5.
A fourth tool we repeatedly use below is a multidimensional generalization of a lemma used in . We warn the reader that the symbol for a norm means different things in different contexts, e.g., for scalar–valued, operator–valued, and vector–valued functions, it respectively means absolute value, operator norm, and vector space norm.
###### Lemma 6.4
The quantity
$$\nu =\underset{\alpha }{sup}(1+|\alpha |)^{d+1}\underset{\{l:\mathrm{\hspace{0.17em}0}l_i\alpha _i\}}{}\frac{1}{(1+|l|)^{d+1}}\frac{1}{(1+|\alpha l|)^{d+1}}.$$
(6.5)
is finite.
Let $`\mathrm{\Sigma }`$ be an open subset of $`IC^d`$. Suppose $`M()C^{\mathrm{}}(\mathrm{\Sigma })`$ is scalar–valued or operator–valued, and $`N()C^{\mathrm{}}(\mathrm{\Sigma })`$ is either operator–valued or vector–valued. Assume these functions satisfy
$`\left(D^\alpha M\right)(x)`$ $``$ $`m(x)a(x)^{|\alpha +p|}{\displaystyle \frac{(\alpha +p)!}{(1+|\alpha |)^{d+1}}}`$ (6.6)
$`\left(D^\alpha N\right)(x)`$ $``$ $`n(x)a(x)^{|\alpha +q|}{\displaystyle \frac{(\alpha +q)!}{(1+|\alpha |)^{d+1}}}`$ (6.7)
for $`x\mathrm{\Sigma }`$, all multi-indices $`\alpha `$, and some fixed multi-indices $`p`$ and $`q`$. Then
$$\left(D^\alpha (MN)\right)(x)m(x)n(x)\nu a(x)^{|\alpha +p+q|}\frac{(\alpha +p+q)!}{(1+|\alpha |)^{d+1}}$$
(6.8)
for each multi-index $`\alpha `$, where $`\nu `$ is defined by (6.5).
### 6.2 Proof of Proposition 6.1
We prove Proposition 6.1 by induction and begin with the case $`n=0`$. We construct $`c_{0,0,0,0,\beta ,j}c_{0,j}`$ with $`\beta =1_{0,1}\{1\}`$. We note that there is no $`d_{n,p,l,k,\beta }(w,t)`$ for $`n1`$; the inequalities for its indices in the conclusion to the proposition cannot be satisfied by non-negative integers. Whenever $`d_{n,p,l,k,\beta }(w,t)`$ with $`n1`$ appears in any of the formal calculations below, it is understood to be zero.
We now assume that the estimates (6.1) and (6.4) on $`c_{m,p,l,k,\beta }(w,t)`$ and $`d_{m,p,l,k,\beta }(w,t)`$ are true for all $`mn1`$ and prove they still hold for $`m=n`$.
Our strategy is to show that each contribution $`\mathrm{\Delta }_i`$ and $`\mathrm{\Gamma }_i`$ consists of a finite sum of terms that satisfy the required estimate. We estimate the number of terms by a separate argument. Our main tools are Lemmas 2.1, 6.2, 6.3, and 6.4.
The index $`\beta `$ must be considered when counting the number of terms, but it plays no role in the estimates of the individual terms. To simplify the notation, we drop it while estimating the terms.
The Term $`\mathrm{\Delta }_1`$
We begin by considering the contribution to (6.4) from the term $`\mathrm{\Delta }_1`$ in (3.17).
By induction, each $`d_{n2,p,l,k,\beta }(w,t)`$ is analytic for $`t\mathrm{\Omega }`$ and has a $`p^{\text{th}}`$ order zero at $`t=0`$. It follows that $`d_{n2,p,l,k,\beta }(w,t)=t^pf(t)`$, where $`f`$ is analytic in $`\mathrm{\Omega }`$. When we take the time derivative, we obtain two terms, $`pt^{p1}f(t)`$ and $`t^p\dot{f}(t)`$. These, respectively, give rise to two terms $`d_{n,p1,l,k,\beta ^{}}(w,t)`$ and $`d_{n,p,l,k+1,\beta ^{\prime \prime }}(w,t)`$.
We consider all $`w`$-derivatives of $`\mathrm{\Delta }_1(w,t)`$. We apply the induction hypothesis, Lemma 6.2, and Lemma 6.4 to obtain
$`D_w^\alpha \left(rP_{}(w,t)\dot{d}_{n2,p,l,k}(w,t)\right)`$
$`c_1\nu D_1D_2^{|\alpha |+|l|+4(n3)}{\displaystyle \frac{(\alpha +l)!}{(1+|\alpha |)^{d+1}}}\sqrt{{\displaystyle \frac{(J+(n3)+2(p|l|k))!}{J!}}}`$
$`\times \left({\displaystyle \frac{|t|^{p1}}{(p1)!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}+{\displaystyle \frac{t^p}{p!}}{\displaystyle \frac{(k+1)^{k+1}}{\text{dist}(t)^{k+1}}}\right)`$
$`=c_1\nu D_1D_2^8D_2^{|\alpha |+|l|+4(n1)}{\displaystyle \frac{(\alpha +l)!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p^{}}{p^{}!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}\sqrt{{\displaystyle \frac{(J+(n1)+2(p^{}|l|k))!}{J!}}}`$
$`+c_1\nu D_1D_2^8D_2^{|\alpha |+|l|+4(n1)}{\displaystyle \frac{(\alpha +l)!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p}{p!}}{\displaystyle \frac{k_{}^{}{}_{}{}^{k^{}}}{\text{dist}(t)^k^{}}}\sqrt{{\displaystyle \frac{(J+(n1)+2(p|l|k^{}))!}{J!}}},`$
with $`p^{}=p1`$ and $`k^{}=k+1`$. We check that
$`p`$ $``$ $`n3<n1,`$
$`p^{}`$ $``$ $`n4<n1,`$
$`|l|+k`$ $``$ $`p+(n3)/2=p^{}+(n1)/2,`$
$`|l|+k^{}`$ $``$ $`p+(n3)/2+1=p+(n1)/2,`$
and the ranges of the components of each vector satisfy
$`|j|`$ $``$ $`J+(n1)+2(p^{}|l|k),`$
$`|j|`$ $``$ $`J+(n1)+2(p|l|k^{}),`$
as required. Hence, we get the desired bound for each of the two contributions, provided
$$D_2^8c_1\nu .$$
The Term $`\mathrm{\Delta }_2`$
In the analysis of this term, we encounter an infinite matrix that represents multiplication by $`y^m`$ in the basis of semiclassical wave packets. We denote this matrix by $`\phi ,y^m\phi `$. Its entries are $`\phi _j,y^m\phi _q(t)`$, for multi-indices $`m,j,q\text{N}^d`$. We recall that Lemma 2.1 gives bounds for these matrix elements and also states that $`\phi _j,y^m\phi _q(t)=0`$ if $`\left||j||q|\right|>|m|`$.
We adopt the analogous notation for the infinite matrix $`\phi ,D_y^m\phi `$ that represents the operator $`D_y^m`$ in the basis of semiclassical wave packets.
We define $`d_0=\sqrt{2}d`$. Then, using (5.1), Lemmas 2.1, 6.3, 6.2, and 6.4, and some algebra, we obtain
$`D_w^\alpha {\displaystyle \underset{\stackrel{~}{m}=3}{\overset{n}{}}}{\displaystyle \underset{|m|=\stackrel{~}{m}}{}}{\displaystyle \frac{D^mE(a(t))}{m!}}\phi ,y^m\phi (rP_{})(w,t)d_{n\stackrel{~}{m},p,l,k}(w,t)`$
$``$ $`{\displaystyle \underset{\stackrel{~}{m}=3}{\overset{n}{}}}{\displaystyle \underset{|m|=\stackrel{~}{m}}{}}(d_0A)^{\stackrel{~}{m}}{\displaystyle \frac{c_6c_1\nu D_2^{\stackrel{~}{m}}m!}{(1+\stackrel{~}{m})^{d+1}m!}}\sqrt{{\displaystyle \frac{(J+(n1)+2(p|l|k))!}{J!}}}`$
$`\times D_1D_2^{|\alpha |+|l|+4(n1\stackrel{~}{m})}{\displaystyle \frac{(\alpha +l)!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p}{p!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}`$
$``$ $`{\displaystyle \underset{\stackrel{~}{m}=3}{\overset{n}{}}}{\displaystyle \underset{|m|=\stackrel{~}{m}}{}}D_1c_6c_1\nu {\displaystyle \frac{(d_0A)^{\stackrel{~}{m}}}{D_2^{3\stackrel{~}{m}}}}D_2^{|\alpha |+|l|+4(n1)}{\displaystyle \frac{(\alpha +l)!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p}{p!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}`$
$`\times \sqrt{{\displaystyle \frac{(J+n1+2(p|l|k))!}{J!}}}.`$
We also verify the constraints on the parameters and components of the vectors:
$`p`$ $``$ $`n1\stackrel{~}{m}n1,`$
$`|l|+k`$ $``$ $`p+(n1\stackrel{~}{m})/2p+(n1)/2,`$
$`|j|`$ $``$ $`J+(n\stackrel{~}{m}1)+2(p|l|k)+\stackrel{~}{m}J+(n1)+2(p|l|k).`$
Hence, we see that each contribution from $`\mathrm{\Delta }_2(w,t)`$ satisfies the required bound, provided the following two conditions are fulfilled
$`D_2^3`$ $``$ $`d_0A,`$
$`D_2^9`$ $``$ $`(d_0A)^3c_6c_1\nu .`$
There are
$$\underset{\stackrel{~}{m}=3}{\overset{n}{}}\underset{|m|=\stackrel{~}{m}}{}1\underset{|m|n}{}1=\left(\begin{array}{c}n+d\\ d\end{array}\right)\sigma _0\text{e}^{\sigma n}$$
such contributions, where $`\sigma >0`$ can be chosen arbitrarily small (see ).
The Term $`\mathrm{\Delta }_3`$
For this term we make the laplacian explicit and write
$$\mathrm{\Delta }_wd_{n4,p,l,k}(w,t)=\underset{i=1}{\overset{d}{}}(D_{w_i}^2d_{n4,p,l,k})(w,t).$$
We introduce $`l_{i,2}=l+(0,0,\mathrm{}0,2,0,\mathrm{},0)`$, where the $`2`$ sits in the $`i`$th column.
We then estimate
$`D_w^\alpha {\displaystyle \frac{1}{2}}\left(rP_{}(w,t)\mathrm{\Delta }_wd_{n4,p,l,k}(w,t)\right)`$
$``$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{d}{}}}c_1\nu D_1D_2^{|\alpha |+|l|+2+4(n5)}{\displaystyle \frac{(\alpha +l_{i,2})!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p}{p!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}`$
$`\times \sqrt{{\displaystyle \frac{(J+(n5)+2(p|l|k))!}{J!}}}`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \frac{c_1\nu D_1}{2D_2^{16}}}D_2^{|\alpha |+|l_{i,2}|+4(n1)}{\displaystyle \frac{(\alpha +l_{i,2})!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p}{p!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}`$
$`\times \sqrt{{\displaystyle \frac{(J+(n1)+2(p|l_{i,2}|k))!}{J!}}}.`$
Again, the constraints are satisfied since
$`p`$ $``$ $`n5<n1`$
$`|l_{i,2}|+k`$ $``$ $`p+(n5)/2+2=p+(n1)/2`$
$`|j|`$ $``$ $`J+(n5)+2(p|l|k)=J+(n1)+2(p|l_{i,2}|k)`$
and each of the $`d`$ contributions stemming from $`\mathrm{\Delta }_3(w,t)`$ satisfies the required estimate, provided
$$D_2^{16}c_1\nu /2.$$
We estimate each of the remaining terms $`\mathrm{\Delta }_i(w,t)`$$`i=4,\mathrm{},8`$, in the same fashion, using the same tools. Since this is straightforward, we only outline the arguments.
The Term $`\mathrm{\Delta }_4`$
We expand the dot product
$$(_w\mathrm{\Phi })(_wc_{n4,p,l,k})=\underset{i=1}{\overset{d}{}}(D_{w_i}\mathrm{\Phi })(D_{w_i}c_{n4,p,l,k})$$
and use the definition
$$l_{i,1}=l+(0,0,\mathrm{}0,1,0,\mathrm{},0),$$
where the $`1`$ sits at the $`i`$th column. Recall that the estimates on the $`c_{m,p,l,k}`$’s differ from those on the $`d_{m,p,l,k}`$’s by a shift of $`1`$ in the $`m`$ dependence. We have
$`D_w^\alpha \left(rP_{}(w,t)_w\mathrm{\Phi }(w,t)_wc_{n4,p,l,k}(w,t)\right)`$
$``$ $`{\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \frac{c_1c_4\nu ^2D_1}{D_2^{12}}}D_2^{|\alpha |+|l_{i,1}|+4(n1)}{\displaystyle \frac{(\alpha +l_{i,1})!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p}{p!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}`$
$`\times \sqrt{{\displaystyle \frac{(J+(n1)+2(p|l_{i,1}|k))!}{J!}}},`$
with all constraints on $`|j|,p,|l_{i,1}|,k`$ satisfied. Thus each of the $`d`$ contributions stemming from $`\mathrm{\Delta }_4(w,t)`$ satisfies the required estimate, provided
$$D_2^{12}c_1c_4\nu ^2.$$
The Term $`\mathrm{\Delta }_5`$
This term is similar to the previous one. We obtain
$`D_w^\alpha \left({\displaystyle \frac{1}{2}}rP_{}(w,t)(\mathrm{\Delta }_w\mathrm{\Phi })(w,t)c_{n4,p,l,k}(w,t)\right)`$
$``$ $`{\displaystyle \frac{c_1c_5\nu ^2D_1}{2D_2^{12}}}D_2^{|\alpha |+|l|+4(n1)}{\displaystyle \frac{(\alpha +l)!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p}{p!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}`$
$`\times \sqrt{{\displaystyle \frac{(J+(n1)+2(p|l|k))!}{J!}}},`$
with all constraints on $`|j|,p,|l|,k`$ satisfied. Thus the contribution stemming from $`\mathrm{\Delta }_5(w,t)`$ satisfies the required estimate, provided
$$D_2^{12}c_1c_5\nu ^2/2.$$
The Term $`\mathrm{\Delta }_6`$
At this point the matrices $`\phi ,D_{y_i}\phi `$ play a role that we control by the momentum space analog of Lemma 2.1. Expanding the dot product and introducing the matrices $`\phi ,D_{y_i}\phi `$$`i=1,\mathrm{},d`$ we have the following estimate for this term:
$`D_w^\alpha {\displaystyle \underset{i=1}{\overset{d}{}}}(rP_{})(w,t)\phi ,D_{y_i}\phi D_{w_i}d_{n3,p,l,k}(w,t)`$
$``$ $`{\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \frac{d_0c_1\nu BD_1}{D_2^{12}}}D_2^{|\alpha |+|l_{i,1}|+4(n1)}{\displaystyle \frac{(\alpha +l_{i,1})!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p}{p!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}`$
$`\times \sqrt{{\displaystyle \frac{(J+(n1)+2(p|l_{i,1}|k))!}{J!}}}`$
with all constraints on $`|j|,p,|l_{i,1}|,k`$ satisfied. Thus, each of the $`d`$ contributions stemming from $`\mathrm{\Delta }_6(w,t)`$ satisfies the required estimate, provided
$$D_2^{12}d_0c_1\nu B.$$
The Term $`\mathrm{\Delta }_7`$
Similarly,
$`D_w^\alpha {\displaystyle \underset{i=1}{\overset{d}{}}}(rP_{})(w,t)\phi ,D_{y_i}\phi (D_{w_i}\mathrm{\Phi })(w,t)c_{n3,p,l,k}(w,t)`$
$``$ $`{\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \frac{d_0c_1c_4\nu ^2BD_1}{D_2^8}}D_2^{|\alpha |+|l|+4(n1)}{\displaystyle \frac{(\alpha +l)!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p}{p!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}`$
$`\times \sqrt{{\displaystyle \frac{(J+(n1)+2(p|l|k))!}{J!}}},`$
with all constraints on $`|j|,p,|l|,k`$ satisfied. Thus, each of the $`d`$ contributions stemming from $`\mathrm{\Delta }_7(w,t)`$ satisfies the required estimate, provided
$$D_2^{12}d_0c_1c_4\nu ^2B.$$
The Term $`\mathrm{\Delta }_8`$
Finally,
$`D_w^\alpha \left(rP_{}(w,t)\dot{\mathrm{\Phi }}(w,t)c_{n2,p,l,k}(w,t)\right)`$
$`{\displaystyle \frac{c_1c_3\nu ^2D_1}{D_2^4}}D_2^{|\alpha |+|l|+4(n1)}{\displaystyle \frac{(\alpha +l)!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p}{p!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}\sqrt{{\displaystyle \frac{(J+(n1)+2(p|l|k))!}{J!}}},`$
with all constraints on $`|j|,p,|l|,k`$ satisfied. Thus the contribution stemming from $`\mathrm{\Delta }_8(w,t)`$ satisfies the required estimate, provided
$$D_2^4c_1c_3\nu ^2.$$
We now perform a similar analysis for the quantities $`\mathrm{\Gamma }_i(w,t)`$ that appear in the expression for $`\dot{c}_{n,p,l,k}(w,t)`$. We integrate these terms with respect to $`t`$ and apply Lemma 6.3. According to the lemma, integration of a term with a given value of $`p`$ gives rise to a term with $`p^{}=p+1`$ in the estimates. We also note that the estimates we want to prove for the $`c`$’s differ from those for the $`d`$’s by the replacement of $`n1`$ by $`n`$.
The Term $`_0^t\mathrm{\Gamma }_1`$
We use the same techniques above to obtain
$`D_w^\alpha {\displaystyle _0^t}{\displaystyle \frac{1}{2}}(\mathrm{\Delta }_wc_{n2,p,l,k})(w,s)𝑑s`$
$``$ $`{\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \frac{D_1}{2D_2^8}}D_2^{|\alpha |+|l_{i,2}|+4n}{\displaystyle \frac{(\alpha +l_{i,2})!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{t^p^{}}{p^{}!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}\sqrt{{\displaystyle \frac{(J+n+2(p^{}|l_{i,2}|k))!}{J!}}}.`$
We check that the constraints are satisfied:
$`p^{}`$ $``$ $`n1<n,`$
$`|l_{i,2}|+k`$ $``$ $`p+(n2)/2+2=p^{}+n/2`$
$`|j|`$ $``$ $`J+(n2)+2(p|l|k)=J+n+2(p^{}|l_{i,2}|k).`$
Thus, each of the $`d`$ contributions stemming from $`\mathrm{\Gamma }_1(w,t)`$ satisfies the required estimate provided
$$D_2^81/2.$$
The Term $`_0^t\mathrm{\Gamma }_2`$
Similarly, with $`p^{}=p+1`$,
$`D_w^\alpha {\displaystyle _0^t}{\displaystyle \underset{i=1}{\overset{d}{}}}\mathrm{\Phi },D_{w_i}\mathrm{\Phi }(w,s)D_{w_i}c_{n2,p,l,k}(w,s)ds`$
$``$ $`{\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \frac{c_4^{}\nu D_1}{D_2^8}}D_2^{|\alpha |+|l_{i,1}|+4n}{\displaystyle \frac{(\alpha +l_{i,1})!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p^{}}{p^{}!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}\sqrt{{\displaystyle \frac{(J+n+2(p^{}|l_{i,1}|k))!}{J!}}},`$
with all constraints on $`|j|,p^{},|l_{i,1}|,k`$ satisfied. Thus each of the $`d`$ contributions stemming from $`\mathrm{\Gamma }_2(w,t)`$ satisfies the required estimate, provided
$$D_2^8c_4^{}\nu .$$
The Term $`_0^t\mathrm{\Gamma }_3`$
Again, with $`p^{}=p+1`$,
$`D_w^\alpha {\displaystyle _0^t}{\displaystyle \frac{1}{2}}\mathrm{\Phi },(\mathrm{\Delta }_w\mathrm{\Phi })(w,s)c_{n2,p,l,k}(w,s)𝑑s`$
$``$ $`{\displaystyle \frac{c_5^{}\nu D_1}{2D_2^8}}D_2^{|\alpha |+|l|+4n}{\displaystyle \frac{(\alpha +l)!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p^{}}{p^{}!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}\sqrt{{\displaystyle \frac{(J+n+2(p^{}|l|k))!}{J!}}},`$
with all constraints on $`|j|,p^{},|l|,k`$ satisfied. Thus, the contribution stemming from $`\mathrm{\Gamma }_3(w,t)`$ satisfies the required estimate, provided
$$D_2^8c_5^{}\nu /2.$$
The Term $`_0^t\mathrm{\Gamma }_4`$
Recall that the matrices $`\phi ,D_{y_i}\phi `$ are controlled by an analog of Lemma 2.1.
$`D_w^\alpha {\displaystyle _0^t}{\displaystyle \underset{i=1}{\overset{d}{}}}\phi ,D_{y_i}\phi D_{w_i}c_{n1,p,l,k}(w,s)ds`$
$`{\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \frac{d_0BD_1}{D_2^4}}D_2^{|\alpha |+|l_{i,1}|+4n}{\displaystyle \frac{(\alpha +l_{i,1})!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p^{}}{p^{}!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}\sqrt{{\displaystyle \frac{(J+n+2(p^{}|l_{i,1}|k))!}{J!}}},`$
with all constraints on $`|j|,p^{},|l_{i,1}|,k`$ satisfied. Thus each of the $`d`$ contributions stemming from $`\mathrm{\Gamma }_4(w,t)`$ satisfies the required estimate, provided
$$D_2^4d_0B.$$
The Term $`_0^t\mathrm{\Gamma }_5`$
For this term we obtain
$`D_w^\alpha {\displaystyle _0^t}{\displaystyle \underset{i=1}{\overset{d}{}}}\mathrm{\Phi },(D_{w_i}\mathrm{\Phi })(w,s)\phi ,D_{y_i}\phi c_{n1,p,l,k}(w,s)ds`$
$`{\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \frac{c_4^{}\nu d_0BD_1}{D_2^4}}D_2^{|\alpha |+|l|+4n}{\displaystyle \frac{(\alpha +l)!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p^{}}{p^{}!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}\sqrt{{\displaystyle \frac{(J+n+2(p^{}|l|k))!}{J!}}},`$
with all constraints on $`|j|,p^{},|l|,k`$ satisfied. Thus each of the $`d`$ contributions stemming from $`\mathrm{\Gamma }_5(w,t)`$ satisfies the required estimate, provided
$$D_2^4c_4^{}\nu d_0B.$$
The Term $`_0^t\mathrm{\Gamma }_6`$
In this term, we encounter the sum over all previous $`c`$’s. As in the similar contribution from $`\mathrm{\Delta }_2`$, we obtain
$`D_w^\alpha {\displaystyle \underset{\stackrel{~}{m}=3}{\overset{n+2}{}}}{\displaystyle \underset{|m|=\stackrel{~}{m}}{}}{\displaystyle _0^t}{\displaystyle \frac{D^mE(a(t))}{m!}}\phi ,y^m\phi c_{n+2\stackrel{~}{m},p,l,k}(w,s)`$
$``$ $`{\displaystyle \underset{\stackrel{~}{m}=3}{\overset{n}{}}}{\displaystyle \underset{|m|=\stackrel{~}{m}}{}}D_1c_6D_2^8{\displaystyle \frac{(d_0A)^{\stackrel{~}{m}}}{D_2^{3\stackrel{~}{m}}}}D_2^{|\alpha |+|l|+4n}{\displaystyle \frac{(\alpha +l)!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p^{}}{p^{}!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}`$
$`\times \sqrt{{\displaystyle \frac{(J+n+2(p^{}|l|k))!}{J!}}}.`$
We check that the constraints on the parameters and components of the vectors are satisfied
$`p^{}`$ $``$ $`n\stackrel{~}{m}+3n`$
$`|l|+k`$ $``$ $`p+(n\stackrel{~}{m}+2)/2p^{}+n/2`$
$`|j|`$ $``$ $`J+n+2+2(p|l|k)=J+n+2(p^{}|l|k).`$
Hence we see that each contribution from $`\mathrm{\Delta }_2(w,t)`$ satisfies the required bound, provided the following two conditions are fulfilled
$`D_2^3`$ $``$ $`d_0A`$
$`D_2`$ $``$ $`(d_0A)^3c_6.`$
There are $`{\displaystyle \underset{\stackrel{~}{m}=3}{\overset{n}{}}}{\displaystyle \underset{|m|=\stackrel{~}{m}}{}}1\sigma _0\text{e}^{\sigma n}`$ such contributions, where $`\sigma >0`$ can be chosen arbitrarily small.
The Term $`_0^t\mathrm{\Gamma }_7`$
This terms depends on the $`d`$’s. Recall the estimates are a little different for them.
$`D_w^\alpha {\displaystyle _0^t}{\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \frac{1}{2}}\mathrm{\Phi }(w,s),(D_{w_i}^2d_{n2,p,l,k})(w,s)ds`$
$``$ $`{\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \frac{c_2\nu D_1}{2D_2^{12}}}D_2^{|\alpha |+|l_{i,2}|+4n}{\displaystyle \frac{(\alpha +l_{i,2})!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p^{}}{p^{}!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}\sqrt{{\displaystyle \frac{(J+n+2(p^{}|l_{i,2}|k))!}{J!}}},`$
with all constraints on $`|j|,p^{},|l_{i,2}|,k`$ satisfied. Thus each of the $`d`$ contributions stemming from $`\mathrm{\Gamma }_7(w,t)`$ satisfies the required estimate, provided
$$D_2^{12}c_2\nu /2.$$
The Term $`_0^t\mathrm{\Gamma }_8`$
Similarly,
$`D_w^\alpha {\displaystyle _0^t}{\displaystyle \underset{i=1}{\overset{d}{}}}\phi ,D_{y_i}\phi \mathrm{\Phi }(w,s),(D_{w_i}d_{n1,p,l,k})(w,s)ds`$
$`{\displaystyle \underset{i=1}{\overset{d}{}}}{\displaystyle \frac{c_2\nu d_0BD_1}{D_2^8}}D_2^{|\alpha |+|l_{i,1}|+4n}{\displaystyle \frac{(\alpha +l_{i,1})!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p^{}}{p^{}!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}\sqrt{{\displaystyle \frac{(J+n+2(p^{}|l_{i,1}|k))!}{J!}}},`$
with all constraints on $`|j|,p^{},|l_{i,1}|,k`$ satisfied. Thus, each of the $`d`$ contributions stemming from $`\mathrm{\Gamma }_8(w,t)`$ satisfies the required estimate provided
$$D_2^8c_2\nu d_0B.$$
The Term $`_0^t\mathrm{\Gamma }_9`$
Finally,
$`D_w^\alpha {\displaystyle _0^t}\dot{\mathrm{\Phi }}(w,s),d_{n,p,l,k}(w,s)𝑑s`$
$``$ $`{\displaystyle \frac{c_3\nu D_1}{D_2^4}}D_2^{|\alpha |+|l|+4n}{\displaystyle \frac{(\alpha +l)!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p^{}}{p^{}!}}{\displaystyle \frac{k^k}{\text{dist}(t)^k}}\sqrt{{\displaystyle \frac{(J+n+2(p^{}|l|k))!}{J!}}},`$
with all constraints on $`|j|,p^{},|l|,k`$ satisfied. Thus, the contribution stemming from $`\mathrm{\Gamma }_9(w,t)`$ satisfies the required estimate, provided
$$D_2^4c_3\nu .$$
By choosing $`D_2`$ large enough, all conditions are satisfied. This completes the induction for part iii) of Proposition 6.1.
The integration required to construct the $`c`$’s shows that we obtain non-zero results for $`c_{n,p,l,k,\beta }`$ for $`n>0`$ only when $`p1`$. This proves part i) of Proposition 6.1.
We now turn to the proof of part ii) of Proposition 6.1.
### 6.3 Counting the Number of Terms that Occur in Our Expansion
In our Born–Oppenheimer expansion, the $`n^{\text{th}}`$ order term has the form
$`\varphi _n(w,y,t)=g_n(w,y,t)\mathrm{\Phi }(w,t)+\varphi _n^{}(w,y,t)`$. The way we compute $`g_n(w,y,t)`$ and $`\varphi _n^{}(w,y,t)`$, they decompose naturally as sums over the parameter $`\beta `$. We define $`u_n`$ to be the number of such terms in $`g_n(w,y,t)`$ and $`v_n`$ to be the number of terms in $`\varphi _n^{}(w,y,t)`$.
An examination of our construction shows that $`u_n`$ and $`v_n`$ satisfy the recursive estimates
$`u_{n+1}`$ $``$ $`{\displaystyle \underset{j=0}{\overset{3}{}}}a_ju_{nj}+{\displaystyle \underset{j=0}{\overset{3}{}}}b_jv_{nj}+{\displaystyle \underset{j=0}{\overset{n}{}}}c_1\gamma _1^ju_{nj}+{\displaystyle \underset{j=0}{\overset{n}{}}}c_2\gamma _2^jv_{nj}+v_{n+1}`$ (6.9)
$`v_{n+1}`$ $``$ $`{\displaystyle \underset{j=0}{\overset{3}{}}}d_ju_{nj}+{\displaystyle \underset{j=0}{\overset{3}{}}}e_jv_{nj}+{\displaystyle \underset{j=0}{\overset{n}{}}}c_3\gamma _3^ju_{nj}+{\displaystyle \underset{j=0}{\overset{n}{}}}c_4\gamma _4^jv_{nj},`$ (6.10)
where $`a_i`$, $`b_i`$, $`c_i`$, $`d_i`$, $`e_i`$ and $`\gamma _i`$ are fixed numbers. The exponentials $`\gamma _i^j`$ arise from an estimate (proven in the proof of Lemma 5.2 of ) for the number of Taylor series terms of any given order in the expansion of $`E(a(t)+ϵy)`$.
We substitute (6.10) for the last term in (6.9) and add the result to (6.10). By some simple estimates this leads to a recursive estimate for the single quantity $`z_n=u_n+v_n`$ of the form
$$z_{n+1}\underset{j=0}{\overset{3}{}}\stackrel{~}{a}_jz_{nj}+\underset{j=0}{\overset{n}{}}\stackrel{~}{c}\stackrel{~}{\gamma }^jz_{nj}.$$
An easy induction on $`n`$ shows that this implies that $`z_n`$ grows at most like $`e^{kn}`$ for a sufficiently large value of $`k`$. The quantity $`z_n`$ is the number of terms in $`\varphi _n(w,y,t)`$, so this proves the assertion.
Proposition 6.1 now follows easily.
## 7 Exponential Error Bounds
In this section, we prove Theorem 4.1.
### 7.1 The Explicit Error Term
We use the following abstract lemma, whose proof is an easy application of Duhamel’s formula (see e.g. ).
###### Lemma 7.1
Suppose $`H(\mathrm{})`$ is a family of self-adjoint operators for $`\mathrm{}>0`$. Suppose $`\psi (t,\mathrm{})`$ belongs to the domain of $`H(\mathrm{})`$, is continuously differentiable in $`t`$, and approximately solves the Schrödinger equation in the sense that
$$i\mathrm{}\frac{\psi }{t}(t,\mathrm{})=H(\mathrm{})\psi (t,\mathrm{})+\xi (t,\mathrm{}),$$
where $`\xi (t,\mathrm{})`$ satisfies
$$\xi (t,\mathrm{})\mu (t,\mathrm{}).$$
Then, for $`t>0`$,
$$\text{e}^{itH(\mathrm{})/\mathrm{}}\psi (0,\mathrm{})\psi (t,\mathrm{})\mathrm{}^1_0^t\mu (s,\mathrm{})𝑑s.$$
The analogous statement holds for $`t<0`$.
We substitute our approximate solution (4)
$`F\text{e}^{iS/ϵ^2}\text{e}^{i\eta y/ϵ}\left({\displaystyle \underset{n=0}{\overset{N}{}}}ϵ^n\varphi _n+ϵ^{N+1}\varphi _{N+1}^{}+ϵ^{N+2}\varphi _{N+2}^{}\right)`$ into the Schrödinger equation and compute the residual term $`\xi _N`$.
It is more convenient to write this term in the multiple scales notation. We also use the notation $`ϵ^m{\displaystyle \frac{E^{(m)}(a(t))}{m!}}y^m`$ to denote the Taylor series term $`{\displaystyle \underset{|j|=m}{}}ϵ^{|j|}{\displaystyle \frac{(D^jE)(a(t))}{j!}}y^j`$.
In this notation, the residual $`\xi _N(w,y,t)`$ is given, up to a phase factor, by two sums of terms. The first one contains all terms that do not involve derivatives of the cut-off. The second contains all terms that do involve derivatives of the cut-off.
The first sum is $`F(w)`$ times the following:
$`{\displaystyle \frac{ϵ^{N+3}}{2}}\left(\mathrm{\Delta }_wg_{N1}\right)\mathrm{\Phi }`$ (7.1)
$`+`$ $`{\displaystyle \frac{ϵ^{N+4}}{2}}\left(\mathrm{\Delta }_wg_N\right)\mathrm{\Phi }`$ (7.2)
$`+`$ $`ϵ^{N+3}\left(_wg_{N1}\right)\left(_w\mathrm{\Phi }\right)`$ (7.3)
$`+`$ $`ϵ^{N+4}\left(_wg_N\right)\left(_w\mathrm{\Phi }\right)`$ (7.4)
$`+`$ $`{\displaystyle \frac{ϵ^{N+3}}{2}}g_{N1}\left(\mathrm{\Delta }_w\mathrm{\Phi }\right)`$ (7.5)
$`+`$ $`{\displaystyle \frac{ϵ^{N+4}}{2}}g_N\left(\mathrm{\Delta }_w\mathrm{\Phi }\right)`$ (7.6)
$`+`$ $`{\displaystyle \frac{ϵ^{N+3}}{2}}\left(\mathrm{\Delta }_w\varphi _{N1}^{}\right)`$ (7.7)
$`+`$ $`{\displaystyle \frac{ϵ^{N+4}}{2}}\left(\mathrm{\Delta }_w\varphi _N^{}\right)`$ (7.8)
$`+`$ $`ϵ^{N+3}\left(_w_yg_N\right)\mathrm{\Phi }`$ (7.9)
$`+`$ $`ϵ^{N+3}\left(_yg_N\right)\left(_w\mathrm{\Phi }\right)`$ (7.10)
$`+`$ $`ϵ^{N+3}\left(_w_y\varphi _N^{}\right)`$ (7.11)
$`+`$ $`iϵ^{N+3}\dot{\varphi }_{N+1}^{}`$ (7.12)
$`+`$ $`iϵ^{N+4}\dot{\varphi }_{N+2}^{}`$ (7.13)
$`+`$ $`{\displaystyle \frac{ϵ^{N+5}}{2}}\left(\mathrm{\Delta }_w\varphi _{N+1}^{}\right)`$ (7.14)
$`+`$ $`{\displaystyle \frac{ϵ^{N+6}}{2}}\left(\mathrm{\Delta }_w\varphi _{N+2}^{}\right)`$ (7.15)
$`+`$ $`ϵ^{N+4}\left(_w_y\varphi _{N+1}^{}\right)`$ (7.16)
$`+`$ $`ϵ^{N+5}\left(_w_y\varphi _{N+2}^{}\right)`$ (7.17)
$`+`$ $`{\displaystyle \frac{ϵ^{N+3}}{2}}\left(\mathrm{\Delta }_y\varphi _{N+1}^{}\right)`$ (7.18)
$`+`$ $`{\displaystyle \frac{ϵ^{N+4}}{2}}\left(\mathrm{\Delta }_y\varphi _{N+2}^{}\right)`$ (7.19)
$``$ $`{\displaystyle \frac{ϵ^{N+3}}{2}}E^{(2)}(a(t))y^2\varphi _{N+1}^{}`$ (7.20)
$``$ $`{\displaystyle \frac{ϵ^{N+4}}{2}}E^{(2)}(a(t))y^2\varphi _{N+2}^{}`$ (7.21)
$``$ $`{\displaystyle \underset{n=0}{\overset{N}{}}}ϵ^{Nn}\left(E(a(t)+ϵy){\displaystyle \underset{m2+n}{}}ϵ^m{\displaystyle \frac{E^{(m)}(a(t))}{m!}}y^m\right)g_{Nn}\mathrm{\Phi }`$ (7.22)
$``$ $`{\displaystyle \underset{n=0}{\overset{N}{}}}ϵ^{Nn}\left(E(a(t)+ϵy){\displaystyle \underset{m2+n}{}}ϵ^m{\displaystyle \frac{E^{(m)}(a(t))}{m!}}y^m\right)\varphi _{Nn}^{}`$ (7.23)
$``$ $`ϵ^{N+1}\left(E(a(t)+ϵy){\displaystyle \underset{m2}{}}ϵ^m{\displaystyle \frac{E^{(m)}(a(t))}{m!}}y^m\right)\varphi _{N+1}^{}`$ (7.24)
$``$ $`ϵ^{N+2}\left(E(a(t)+ϵy){\displaystyle \underset{m2}{}}ϵ^m{\displaystyle \frac{E^{(m)}(a(t))}{m!}}y^m\right)\varphi _{N+2}^{}.`$ (7.25)
The second sum arises from terms in which the cut-off $`F(w)`$ is differentiated. It is
$`{\displaystyle \underset{n=0}{\overset{N}{}}}{\displaystyle \frac{ϵ^{n+4}}{2}}(\mathrm{\Delta }_wF)g_n\mathrm{\Phi }`$ (7.26)
$`+`$ $`{\displaystyle \underset{n=0}{\overset{N+2}{}}}{\displaystyle \frac{ϵ^{n+4}}{2}}(\mathrm{\Delta }_wF)\varphi _n^{}`$ (7.27)
$`+`$ $`{\displaystyle \underset{n=0}{\overset{N}{}}}ϵ^{n+4}(_wF)(_wg_n)\mathrm{\Phi }`$ (7.28)
$`+`$ $`{\displaystyle \underset{n=0}{\overset{N}{}}}ϵ^{n+4}g_n(_wF)(_w\mathrm{\Phi })`$ (7.29)
$`+`$ $`{\displaystyle \underset{n=0}{\overset{N+2}{}}}ϵ^{n+4}(_wF)(_w\varphi _n^{})`$ (7.30)
$`+`$ $`{\displaystyle \underset{n=0}{\overset{N}{}}}ϵ^{n+3}(_wF)(_yg_n)\mathrm{\Phi }`$ (7.31)
$`+`$ $`{\displaystyle \underset{n=0}{\overset{N+2}{}}}ϵ^{n+3}(_wF)(_y\varphi _n^{})`$ (7.32)
### 7.2 Optimal Truncation
Each error term in the first sum (7.1)–(7.25) can be written as a uniformly bounded function times one of the following two forms:
$`𝒜`$ $`=`$ $`\mathrm{\Psi }(w,t){\displaystyle \underset{r}{}}{\displaystyle \underset{|j|\rho (r)}{}}c_{r,j}(w,t)\phi _j(y,t)`$
$``$ $`=`$ $`{\displaystyle \underset{r^{}}{}}{\displaystyle \underset{|j|\rho ^{}(r^{})}{}}d_{r^{},j}(w,t)\phi _j(y,t),`$
where $`\mathrm{\Psi }(w,t)_{\text{el}}`$, $`\phi _j(y,t)=ϵ^{d/2}\phi _j(A(t),B(t),1,0,0,y)`$, $`r,r^{}`$ denote a collective set of indices that belong to some finite set, and $`\rho (r)`$ and $`\rho ^{}(r^{})`$ limit the number of multi-indices $`j`$ allowed in the second sum.
The error term $`\xi (w,y,t)_{\text{el}}`$ needs to be estimated for $`t\mathrm{I}\mathrm{R}`$, in the following norm
$`\xi (t)`$ $`=`$ $`\{{\displaystyle _{\text{IR}^d}}\xi (xa(t),(xa(t))/ϵ,t)_{_{\text{el}}}^2dx\}^{1/2}`$
$`=`$ $`\{{\displaystyle _{\text{IR}^d}}\xi (w,w/ϵ,t)_{_{\text{el}}}^2dw\}^{1/2}.`$
With that norm, using the Cauchy–Schwarz inequality and the $`L^2(\mathrm{I}\mathrm{R}^d)`$ orthonormality of the $`\phi _j(y,t)`$, we obtain the following estimate for the norm of $`𝒜`$ in terms of the norm of vector $`c_r(w,t)l^2(\text{N}^d,\mathrm{I}\mathrm{C})`$,
$$𝒜\underset{r}{}\underset{w\text{supp}F}{sup}\mathrm{\Psi }(w,t)_{_{\text{el}}}\underset{w\text{supp}F}{sup}c_r(w,t)\left(\underset{|j|\rho (r)}{}\mathrm{\hspace{0.17em}1}\right)^{1/2}.$$
(7.33)
By similar arguments we get the following estimate for the norm of $``$ in terms of the norm of the vector $`d_r^{}(w,t)l^2(\text{N}^d,_{\text{el}})`$,
$$\underset{r^{}}{}\underset{w\text{supp}F}{sup}d_r^{}(w,t)\left(\underset{|j|\rho ^{}(r^{})}{}\mathrm{\hspace{0.17em}1}\right)^{1/2}.$$
(7.34)
Note also that
$$\underset{|j|\rho ^{}(r^{})}{}\mathrm{\hspace{0.17em}1}\left(\begin{array}{c}\rho ^{}(r^{})+d\\ d\end{array}\right),$$
(7.35)
which grows at most polynomially with $`\rho ^{}(r^{})`$.
###### Lemma 7.2
For $`t[0,T]`$, and for any $`\alpha \text{N}^d`$ and $`\gamma \text{N}^d`$, there exist $`C_0>0`$ and $`\tau _0>0`$, such that
$`{\displaystyle \underset{\beta _{n,1}}{}}{\displaystyle \underset{pn}{}}{\displaystyle \underset{k+|l|p+\frac{n}{2}}{}}{\displaystyle \underset{|j|J+n+2(p|l|k)}{}}\mathrm{\Psi }(w,t)D_w^\alpha D_y^\gamma c_{n,p,l,k,\beta ,j}(w,t)\phi _j(y,t)`$ (7.36)
$``$ $`C_0\left\{n^{1/2}\tau _0\right\}^n`$
and
$`{\displaystyle \underset{\beta _{n,2}}{}}{\displaystyle \underset{pn1}{}}{\displaystyle \underset{k+|l|p+\frac{n1}{2}}{}}{\displaystyle \underset{|j|J+n1+2(p|l|k)}{}}D_w^\alpha D_y^\gamma d_{n,p,l,k,\beta ,j}(w,t)\phi _j(y,t)`$ (7.37)
$``$ $`C_0\left\{n^{1/2}\tau _0\right\}^n`$
If the operator $`D_y^\gamma `$ is replaced by the operator $`y^\gamma `$, the same bounds are valid.
Proof: We begin with (7.36). We have
$`{\displaystyle \underset{|j|J+n+2(p|l|k)}{}}D_w^\alpha c_{n,p,l,k,\beta ,j}(w,t)D_y^\gamma \phi _j(y,t)`$
$`=`$ $`{\displaystyle \underset{|\stackrel{~}{k}|J+|\gamma |+n+2(p|l|k)}{}}\left(\phi ,D_y^\gamma \phi D_w^\alpha c_{n,p,l,k,\beta }(w,t)\right)_{\stackrel{~}{k}}\phi _{\stackrel{~}{k}}(y,t).`$
We know that the vector $`\phi ,D_y^\gamma \phi D_w^\alpha c_{n,p,l,k,\beta }(w,t)`$ satisfies the estimate
$`\phi ,D_y^\gamma \phi D_w^\alpha c_{n,p,l,k,\beta }(w,t)D_1D_2^{|\alpha |+|l|+4n}{\displaystyle \frac{(\alpha +l)!}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{|t|^p}{p!}}{\displaystyle \frac{k^k}{\delta ^k}}(Bd_0)^{|\gamma |}`$
$`\times \sqrt{{\displaystyle \frac{(J+|\gamma |+n+2(p|l|k))!}{J!}}}.`$ (7.38)
Here $`\delta >0`$ is the distance in the complex plane from $`[0,T]`$ to the complement of $`\mathrm{\Omega }`$.
Since the number of indices in $`_{n,1}`$ is bounded by $`\text{e}^{K_0n}`$, $`D_21`$, and $`(\alpha +l)!(|\alpha |+|l|)!`$, we can estimate the sum (7.36) by
$`{\displaystyle \frac{D_1D_2^{|\alpha |}(Bd_0)^{|\gamma |}}{\sqrt{J!}(1+|\alpha |)^{d+1}}}\text{e}^{K_0n}D_2^{11n/2}{\displaystyle \underset{pn}{}}{\displaystyle \frac{|t|^p}{p!}}{\displaystyle \underset{|l|+kp+n/2}{}}\left({\displaystyle \frac{k}{\delta }}\right)^k`$
$`\times (|\alpha |+|l|)!\sqrt{(J+|\gamma |+n+2(p|l|k))!}.`$ (7.39)
Then, using $`a!b!(a+b)!`$, the fact that $`(a+2p)!/(p!)^2`$ is increasing in $`p`$, and $`pn`$, we have
$$\frac{(J+|\gamma |+n+2(p|l|k))!((|\alpha |+|l|)!)^2}{(p!)^2}\frac{(J+|\gamma |+2|\alpha |+3n2k)!}{(n!)^2},$$
so that (7.39) is bounded by
$`{\displaystyle \frac{D_1D_2^{|\alpha |}(Bd_0)^{|\gamma |}}{n!\sqrt{J!}(1+|\alpha |)^{d+1}}}\text{e}^{K_0n}D_2^{11n/2}{\displaystyle \underset{pn}{}}|t|^p{\displaystyle \underset{k=0}{\overset{p+n/2}{}}}\left({\displaystyle \frac{k}{\delta }}\right)^k`$
$`\times \sqrt{(J+|\gamma |+2|\alpha |+3n2k)!}{\displaystyle \underset{|l|p+n/2k}{}}1.`$
The last term is bounded by $`\left(\begin{array}{c}[[3n/2]]+d\\ d\end{array}\right)\sigma _0\text{e}^{3\sigma n/2}`$, where $`[[x]]`$ denotes the integer part of $`x`$.
Using $`k^{2k}(2k)^{2k}`$, $`a^ab^b(a+b)^{a+b}`$ and $`a!a^a`$ we have
$$(J+|\gamma |+2|\alpha |+3n2k)!k^{2k}(J+|\gamma |+2|\alpha |+3n)^{J+|\gamma |+2|\alpha |+3n}.$$
Since we can assume without loss that $`\delta <1`$, this implies
$`{\displaystyle \underset{k=0}{\overset{p+n/2}{}}}\left({\displaystyle \frac{k}{\delta }}\right)^k\sqrt{(J+|\gamma |+2|\alpha |+3n2k)!}`$ (7.40)
$``$ $`(J+|\gamma |+2|\alpha |+3n)^{\frac{J+|\gamma |+2|\alpha |+3n}{2}}{\displaystyle \underset{k=0}{\overset{p+n/2}{}}}\delta ^k`$
$``$ $`(J+|\gamma |+2|\alpha |+3n)^{\frac{J+|\gamma |+2|\alpha |+3n}{2}}{\displaystyle \frac{K_1}{\delta ^{n/2}}}\delta ^p,`$
for some constant $`K_1`$ that satisfies $`\delta ^11K_1^1\delta ^1`$. Together with
$$\underset{pn}{}(t/\delta )^p\{\begin{array}{ccc}K_2(t/\delta )^n& \text{if}& t/\delta >1\\ K_2^n& \text{if}& t/\delta =1\\ K_2& \text{if}& t/\delta <1,\end{array}$$
(7.41)
where $`K_2`$ is constant, we get (in the first case above)
$`{\displaystyle \underset{\beta _{n,1}}{}}{\displaystyle \underset{pn}{}}{\displaystyle \underset{k+|l|p+\frac{n}{2}}{}}\phi ,D_y^\gamma \phi D_w^\alpha c_{n,p,l,k,\beta }(w,t)`$
$``$ $`{\displaystyle \frac{\sigma _0K_1K_2D_1D_2^{|\alpha |}(Bd_0)^{|\gamma |}}{(1+|\alpha |)^{d+1}}}{\displaystyle \frac{\text{e}^{(K_0+3\sigma /2)n}D_2^{11n/2}t^n}{\sqrt{J!}n!\delta ^{3n/2}}}(J+|\gamma |+2|\alpha |+3n)^{\frac{J+|\gamma |+2|\alpha |+3n}{2}}.`$
We postpone the study of the dependence of our estimates on $`t`$ and $`J`$ to Section 8. So, using the above,
$$(J+|\gamma |+2|\alpha |+3n)^{\frac{J+|\gamma |+2|\alpha |+3n}{2}}(J+|\gamma |+2|\alpha |+3n)^{\frac{J+|\gamma |+2|\alpha |}{2}}((J+|\gamma |+2|\alpha |+3)n)^{\frac{3n}{2}},$$
and the existence of $`0<a<b`$, such that $`a^nn^nn!b^nn^n`$, we learn the existence of positive constants (i.e., independent of $`n`$) $`K_3`$, $`K_4`$ and $`K_5`$, such that
$$\underset{\beta _{n,1}}{}\underset{pn}{}\underset{k+|l|p+\frac{n}{2}}{}\phi ,D_y^\gamma \phi D_w^\alpha c_{n,p,l,k,\beta }(w,t)K_3K_4^n\frac{n^{3n/2}}{a^nn^n}K_3(K_5n^{1/2})^n.$$
This yields the result with $`C_0=K_3`$ and $`\tau _0=K_5`$.
The second sum is dealt with in the same manner, since the vectors $`d_{n,p,l,k,\beta }(w,t)`$ satisfy the same bounds as $`c_{n,p,l,k,\beta }(w,t)`$ does with $`n`$ replaced by $`n1`$.
Finally, the replacement of $`D_y^\gamma `$ by the operator $`y^\gamma `$ means that the matrix $`\phi |D_y^\gamma \phi `$ must be replaced by the matrix $`\phi |y^\gamma \phi `$. But the latter has the same properties as the former; the bounds above remain true with $`B`$ replaced by $`A`$ from (7.2) onward. This affects the definition of $`C_0`$ only.
The following lemma is the key to the proof of exponential accuracy of our approximation by means of optimal truncation.
###### Lemma 7.3
For sufficiently small $`g>0`$, there exist $`\mathrm{\Gamma }(g)`$ and $`C(g)>0`$ such that the choice $`N(ϵ)=[[g^2/ϵ^2]]`$ implies that the norm of the error term $`\xi _{N(ϵ)}(t)`$ given by (7.1) satisfies
$$\xi _{N(ϵ)}(t)C(g)\text{e}^{\mathrm{\Gamma }(g)/ϵ^2}.$$
Proof: The previous lemma, formulas (7.33), (7.34) and (7.35) show that all terms in the first sum defining $`\xi _N`$ except (7.12), (7.13), (7.22), and (7.23) are exponentially small, once we prove
$$C_0ϵ^{N(ϵ)}\left\{N(ϵ)^{1/2}\tau _0\right\}^{N(ϵ)}C\text{e}^{\mathrm{\Gamma }/ϵ^2}.$$
(7.42)
Because $`g^2/ϵ^21Ng^2/ϵ^2`$, if we choose $`0<g<1/\tau _0`$, the left hand side of this inequality is bounded by
$`C_0\left\{ϵN^{1/2}\tau _0\right\}^NC_0\left\{g\tau _0\right\}^NC_0\text{e}^{|\mathrm{ln}(g\tau _0)|N}C_0\text{e}^{|\mathrm{ln}(g\tau _0)|}\text{e}^{|\mathrm{ln}(g\tau _0)|g^2/ϵ^2},`$ (7.43)
which gives
$$C(g)=C_0\text{e}^{|\mathrm{ln}(g\tau _0)|}\text{and}\mathrm{\Gamma }(g)=|\mathrm{ln}(g\tau _0)|g^2.$$
The terms (7.12) and (7.13) can be dealt with in a similar fashion once we have computed
$`\dot{\varphi }_{N+1}^{}`$ $`=`$ $`{\displaystyle \underset{\beta _{N+1,2}}{}}{\displaystyle \underset{pN}{}}{\displaystyle \underset{k+|l|p+\frac{N}{2}}{}}{\displaystyle \underset{|j|J+N+2(p|l|k)}{}}\dot{d}_{N+1,p,l,k,\beta ,j}(w,t)\phi _j(y,t)`$
$`+d_{N+1,p,l,k,\beta ,j}(w,t)\dot{\phi }_j(y,t),`$
where the second term equals
$`{\displaystyle \underset{\beta _{N+1,2}}{}}{\displaystyle \underset{pN}{}}{\displaystyle \underset{k+|l|p+\frac{N}{2}}{}}{\displaystyle \underset{|\stackrel{~}{k}|J+N+2+2(p|l|k)}{}}(({\displaystyle \frac{i}{2}}\phi ,\mathrm{\Delta }_y\phi `$
$`{\displaystyle \frac{iE^{(2)}(a(t))}{2}}\phi ,y^2\phi )d_{N+1,p,l,k,\beta }(w,t))_{\stackrel{~}{k}}\phi _{\stackrel{~}{k}}(y,t).`$
Lemma 6.2 shows that $`\dot{d}_{N+1,p,l,k,\beta }`$ satisfies bounds similar to those satisfied by $`d_{N+1,p,l,k,\beta ,j}`$ and the term above is taken care of by lemma 7.2. Similar statements are true for $`\dot{\varphi }_{N+2}^{}`$, and the analysis above also applies to these error terms.
Next consider (7.22). By the mean value theorem, there exists $`\zeta _q(y,t,ϵ)=a(t)+\theta _q(y,t,ϵ)ϵy`$, where $`q\text{N}^d`$ and $`\theta _q(y,t,ϵ)(0,1)`$, such that
$$E(a(t)+ϵy)\underset{m2+n}{}ϵ^m\frac{E^{(m)}(a(t))}{m!}y^m=\underset{|q|=2+n+1}{}ϵ^{|q|}\frac{D^qE(\zeta _q(y,t,ϵ))}{q!}y^q.$$
Hence, we need to estimate
$$\underset{n=0}{\overset{N}{}}ϵ^{N+3}\underset{|q|=2+n+1}{}\frac{D^qE(\zeta _q(y,t,ϵ))}{q!}y^q\underset{\beta ,p,k,l,j}{}c_{Nn,p,l,k,\beta ,j}(w,t)\phi _j(y,t)\mathrm{\Phi },$$
(7.44)
with the following restrictions
$`|j|`$ $``$ $`J+(Nn)+2(pk|l|)`$
$`k+|l|`$ $``$ $`p+(Nn)/2`$
$`p`$ $``$ $`Nm`$
$`\beta `$ $``$ $`_{1,Nn}.`$ (7.45)
We take a fixed value of $`n[0,N]`$, and consider the vectors
$$\frac{D^qE(\zeta _q(y,t,ϵ))}{q!}\phi ,y^q\phi (t)c_{Nn,p,l,k,\beta }(w,t)$$
we have to estimate. Due to the presence of the cut-off function $`F`$ (which we have omitted in the notation), we have
$$\frac{|D^qE(\zeta _q(y,t,ϵ))|}{q!}\frac{c_6D_2^{|q|}}{(1+|q|)^{d+1}},$$
and with our bounds on the matrix $`\phi ,y^q\phi (t)`$ and on the vector $`c_{Nn,p,l,k,\beta }(w,t)`$, we can write
$`{\displaystyle \frac{D^qE(\zeta _q(y,t,ϵ))}{q!}}\phi ,y^q\phi (t)c_{Nn,p,l,k,\beta }(w,t)`$
$``$ $`{\displaystyle \frac{c_6D_2^{n+3}(d_0A)^{n+3}}{(1+(n+3))^{d+1}}}{\displaystyle \frac{\sqrt{(J+3+N+2(pk|l|))!}}{\sqrt{J!}}}D_1D_2^{|l|+4(Nn)}l!{\displaystyle \frac{|t|^p}{p!}}{\displaystyle \frac{k^k}{\delta ^k}}.`$
Then we use similar estimates to the above and the restrictions (7.2) to get
$`k^{2k}l!l!{\displaystyle \frac{(J+3+N+2(pk|l|))!}{p!p!}}(2k)^{2k}{\displaystyle \frac{(J+3+N+2(pk))!}{p!p!}}`$
$``$ $`(2k)^{2k}{\displaystyle \frac{(J+3+3N2n2k)!}{(Nn)!(Nn)!}}{\displaystyle \frac{(J+3+3N2n)^{J+3+3N2n}}{(Nn)!(Nn)!}}.`$
Using this and $`|l|3(Nn)/2`$, we see that (7.2) is bounded above by
$$\frac{c_6D_1(D_2d_0A)^3}{4^{d+1}\sqrt{J!}}D_2^{11N/2}\frac{(d_0A)^n}{D_2^{9n/2}}\frac{|t|^p}{(Nn)!\delta ^k}(J+3+3N2n)^{(J+3+3N)/2}.$$
Finally, with
$$\underset{p=0}{\overset{Nn}{}}\underset{k=0}{\overset{p+\frac{Nm}{2}}{}}|t|^p\delta ^kK_1K_2\delta ^{(Nm)/2}\left(\frac{t}{\delta }\right)^{Nn},$$
(see (7.40), (7.41)), the bounds $`_{|l|p+(Nn)/2}\mathrm{\hspace{0.17em}1}\sigma _0\text{e}^{3\sigma (Nn)/2}`$, $`_{|q|n+3}\mathrm{\hspace{0.17em}1}\sigma _0\text{e}^{\sigma (n+3)}`$, and $`|_{1,Nn}|\text{e}^{K_0(Nn)}`$, we get (with the conditions (7.2) on the summations)
$`{\displaystyle \underset{|q|=2+n+1}{}}{\displaystyle \underset{\beta ,p,k,l,j}{}}ϵ^N{\displaystyle \frac{D^qE(\zeta _q(y,t,ϵ))}{q!}}\phi ,y^q\phi (t)c_{Nn,p,l,k,\beta }(w,t)`$ (7.47)
$``$ $`{\displaystyle \frac{\sigma _0^2\text{e}^{3\sigma }K_1K_2c_6D_1(D_2d_0A)^3}{4^{d+1}\sqrt{J!}}}\text{e}^{5\sigma N/2}D_2^N(d_0A)^Nϵ^N(J+3+3N)^{(J+3+N)/2}`$
$`\times \left({\displaystyle \frac{D_2^{9/2}}{\delta ^{1/2}d_0A}}\right)^{Nn}\left({\displaystyle \frac{t}{\delta }}\right)^{Nn}{\displaystyle \frac{(J+3+3N)^{Nn}}{(Nn)!}}.`$
Postponing the study of the $`t`$ and $`J`$ dependence of our estimates, we use the bound
$`(J+3+3N)^{(J+3+N)/2}N^{N/2}(J+3+3N)^{(J+3)/2}(J+6)^{N/2}`$ to establish the existence of constants $`L_0,L_1,L_2`$, independent of $`N`$ and $`n`$, such that (7.47) is bounded above by
$$L_0L_1^Nϵ^NN^{N/2}\frac{(L_2N)^{Nn}}{(Nn)!}.$$
It remains for us to sum over $`n`$ and use (7.33) to bound (7.44) by
$`ϵ^3\sqrt{\sigma _0}\text{e}^{3\sigma N/2}L_0L_1^Nϵ^NN^{N/2}{\displaystyle \underset{n=0}{\overset{N}{}}}{\displaystyle \frac{(L_2N)^{Nn}}{(Nn)!}}`$
$``$ $`ϵ^3\sqrt{\sigma _0}L_0(\text{e}^{3\sigma /2}L_1)^Nϵ^NN^{N/2}{\displaystyle \underset{s=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(L_2N)^s}{s!}}`$
$``$ $`ϵ^3\sqrt{\sigma _0}L_0(\text{e}^{3\sigma /2}L_1\text{e}^{L_2})^Nϵ^NN^{N/2}.`$
If we choose $`g<1/(L_1\text{e}^{L_2+3\sigma /2})`$, we can apply the analysis (7.43) to obtain an exponentially small bound on (7.44) by the optimal truncation $`N(ϵ)=[[g^2/ϵ^2]]`$.
Since the estimates we have on the $`d`$’s are similar to those we have on the $`c`$’s, with the replacement of $`n`$ by $`n1`$, the same exponential bound is valid for (7.23), (see (7.34)) and the analysis of the the first collection of error terms is completed.
We now need to take into account the error terms (7.26) to (7.32) arising from the derivatives of the cut-off function $`F`$. Choose $`F_0>0`$ that satisfies
$$\mathrm{max}\{|\mathrm{\Delta }_wF(w)|,_wF(w)\}F_0,$$
uniformly in $`w`$, and recall that for any $`i=1,\mathrm{},d`$,
$$\text{supp}_{w_i}F(w)\{w\mathrm{I}\mathrm{R}^d:b_0<|w|<b_1\}$$
(7.48)
for some $`0<b_0<b_1<\mathrm{}`$. Now consider (7.26). We express $`g_n`$ in terms of the $`c`$’s to see that the norm of $`{\displaystyle \frac{2}{ϵ^4}}`$ times (7.26) (in $`L^2(\mathrm{I}\mathrm{R}^d,_{\text{el}})`$) can be bounded as follows:
$`{\displaystyle \underset{n=0}{\overset{N}{}}}\mathrm{\Delta }_wF(w)g_n(w,y,t)ϵ^n\mathrm{\Phi }(w,t)`$ (7.49)
$``$ $`{\displaystyle \underset{n=0}{\overset{N}{}}}ϵ^n\sqrt{{\displaystyle _{\text{IR}^d}}\left|\mathrm{\Delta }_wF(w){\displaystyle \underset{|j|J+3n}{}}c_{n,j}(w,t)ϵ^{d/2}\phi _j(w/ϵ,t)\right|^2\mathrm{\Phi }(w,t)_{_{\text{el}}}^2𝑑w}`$
$``$ $`F_0{\displaystyle \underset{n=0}{\overset{N}{}}}ϵ^n\underset{w\text{supp}F\text{IR}^d}{sup}c_n(w,t)\sqrt{{\displaystyle \underset{|j|J+3n}{}}{\displaystyle _{|w|b_0}}\left|ϵ^{d/2}\phi _j(w/ϵ,t)\right|^2𝑑w}.`$
We know from Section 7 of that there exists a constant $`0<\beta _d`$ depending on the dimension $`d`$ only, such that
$$\sqrt{2|j|+d}<b_0/(Aϵ),\text{for all}|j|J+3N$$
and $`|j|J+3N`$ imply
$$\sqrt{_{|w|b_0}\left|ϵ^{d/2}\phi _j(w/ϵ,t)\right|^2𝑑w}\text{e}^{\beta _d|j|}\text{e}^{(b_0^2)/(12A^2ϵ^2)}.$$
All the conditions here will be satisfied if $`N(ϵ)=[[g^2/ϵ^2]]`$, provided we choose $`g`$ and $`ϵ`$ to satisfy
$$ϵ^2(d+2J)+6g^2<b_0^2/A^2.$$
For such a choice, using $`_{|j|J+3n}\text{e}^{2\beta _d|j|}\sigma _0\text{e}^{(\sigma +2\beta _d)(J+3n)}`$, we get
$$\sqrt{\underset{|j|J+3n}{}_{|w|b_0}\left|ϵ^{d/2}\phi _j(w/ϵ,t)\right|^2𝑑w}\sqrt{\sigma _0}\text{e}^{(\sigma +2\beta _d)(J+3n)/2}\text{e}^{(b_0^2)(12A^2ϵ^2)}.$$
Moreover, by means of manipulations that by now are familiar,
$`c_n(w,t){\displaystyle \underset{\beta _{n,1}}{}}{\displaystyle \underset{pn}{}}{\displaystyle \underset{k+|l|p+\frac{n}{2}}{}}c_{n,p,l,k,\beta }(w,t)`$
$``$ $`{\displaystyle \underset{\beta _{n,1}}{}}{\displaystyle \underset{pn}{}}{\displaystyle \underset{k+|l|p+\frac{n}{2}}{}}D_1D_2^{|l|+4n}l!{\displaystyle \frac{|t|^p}{p!}}{\displaystyle \frac{k^k}{\delta ^k}}{\displaystyle \frac{(J+n+2(p|l|k))}{\sqrt{J!}}}`$
$``$ $`{\displaystyle \underset{\beta _{n,1}}{}}{\displaystyle \underset{pn}{}}{\displaystyle \underset{k+|l|p+\frac{n}{2}}{}}{\displaystyle \frac{D_1D_2^{11n/2}}{\sqrt{J!}}}{\displaystyle \frac{|t|^p}{\delta ^k}}{\displaystyle \frac{(J+3n)^{(J+3n)/2}}{n!}}`$
$``$ $`\text{e}^{K_0n}\sigma _0\text{e}^{3\sigma n/2}K_1K_2{\displaystyle \frac{1}{\delta ^n}}\left({\displaystyle \frac{|t|}{\delta }}\right)^n{\displaystyle \frac{D_1D_2^{11n/2}}{\sqrt{J!}}}(J+3n)^{J/2}{\displaystyle \frac{(J+3N)^{3n/2}}{n!}}.`$
Combining these estimates, we get the existence of positive constants $`M_0`$ and $`M_1`$, such that for $`N=[[g^2/ϵ^2]]`$,
$`{\displaystyle \underset{n=0}{\overset{N}{}}}\mathrm{\Delta }_wF(w)g_n(w,y,t)ϵ^n\mathrm{\Phi }(w,t)\text{e}^{(b_0^2)/(12A^2ϵ^2)}{\displaystyle \underset{n=0}{\overset{N}{}}}M_0{\displaystyle \frac{(ϵM_1N^{3/2})^n}{n!}}`$
$``$ $`\text{e}^{(b_0^2)/(12A^2ϵ^2)}M_0\text{e}^{ϵM_1N^{3/2}}\text{e}^{(b_0^2)/(12A^2ϵ^2)}M_0\text{e}^{M_1g^3/ϵ^2}M_0\text{e}^{(b_0^2)/(24A^2ϵ^2)},`$
provided
$$M_1g^3<b_0^2/(24A^2).$$
All other terms in the list (7.26) to (7.32) can be estimated in a similar fashion under a similar condition on $`g`$.
This concludes the proof of our lemma.
Remark: It is not difficult to check that if we keep $`N`$ fixed, then our approximation (4) $`\widehat{\psi }(w,y,t)`$ is accurate up to an error of order $`ϵ^N`$, as expected.
A by–product of our estimates on the terms stemming from the introduction of the cutoff is that our approximation is exponentially localized in a ball centered at $`a(t)`$ of any radius $`b_0`$, as stated in the second part of Theorem 4.1.
Hence, we have completed the proof of Theorem 4.1.
## 8 Generalizations
As in , under some mild supplementary assumptions, we can extend our results to allow $`0tT(ϵ)`$ with $`T(ϵ)\mathrm{ln}(1/ϵ^2)`$. This proves the validity of our construction up to the Ehrenfest time scale.
###### Theorem 8.1
In addition to the assumptions of Theorem 4.1, assume that a classical solution to the equation (2.4) exists for all $`tIR`$. Moreover, assume that for all $`z`$ in a complex neighborhood of $`\mathrm{\Xi }`$, the following bound is satisfied
$$|E(z)|N\text{e}^{M|z|},$$
and that $`E(x)`$ is bounded below. Suppose also that there exist $`L`$ and $`\lambda >0`$, such that for all $`tIR`$
$$A(t)+B(t)L\text{e}^{\lambda t}.$$
Then, there exist $`\tau ^{}`$, $`C^{}`$, $`T^{}>0`$, and $`0<\sigma ,\sigma ^{}<2`$ such that the approximation defined by choosing $`N(ϵ)1/ϵ^\sigma `$ is accurate up to an error whose norm is bounded by $`C^{}\text{e}^{\tau ^{}/ϵ^\sigma ^{}}`$, uniformly for all times $`0tT^{}\mathrm{ln}(1/ϵ^2)`$.
Proof: It is enough to mimick the proof of the corresponding result for the semiclassical propagation of the Schrödinger equation in , since our hypotheses imply that nothing can happen on the adiabatic side of the problem. By the conservation of energy, the exponential bound on $`E(z)`$ and the assumed existence of a Liapunov exponent, we easily see from the proof of Lemmas 7.2 and 7.3, that the behavior in $`t`$ of all constants (independent of $`N`$) is at worst exponential in $`t`$. From the conditions $`D_2\text{e}^{KT}`$, with $`K`$ some constant, we need to take $`g(T)g_0\text{e}^{g_1t}`$ so that the optimal truncation procedure yields an error of the order $`\text{e}^{K_0T}\text{e}^{g_0^2\text{e}^{2g_1}/ϵ^2}`$. The choice $`T(ϵ)T^{}\mathrm{ln}(1/ϵ^2)`$, with $`T^{}>0`$ sufficiently small, gives the desired result.
Similarly, we can extend our results to allow initial conditions in a wider class of vectors. Indeed, we have been careful to make explicit the $`J`$ dependence in all estimates so that we can control the error term as a function of $`J`$. Recall that $`J`$ is fixed arbitrarily in (3.8) which gives the expansion in the basis $`\phi _j(A(0),B(0),ϵ^2,a(0),\eta (0),x)`$ of the nuclear part of the wave function that we take as an initial condition.
As in , for $`(a,\eta )\mathrm{I}\mathrm{R}^{2d}`$, we introduce the operator $`\mathrm{\Lambda }_ϵ(a,\eta )`$ such that
$$(\mathrm{\Lambda }_ϵ(a,\eta )f)(x)=ϵ^d\text{e}^{i\eta (xa)/ϵ^2}f((xa)/ϵ).$$
We define a dense set $`𝒞`$ in $`L^2(\mathrm{I}\mathrm{R}^d)`$, that is contained in the set $`𝒮`$ of Schwartz functions, by
$`𝒞`$ $`=`$ $`\{f(x)={\displaystyle \underset{j}{}}c_j\phi _j(\mathrm{I}\mathrm{I},\mathrm{I}\mathrm{I},\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}0},x)𝒮,\text{ such that}`$ (8.1)
$`\text{there exists}K>0\text{ with }{\displaystyle \underset{|j|>J}{}}|c_j|^2\text{e}^{KJ},\text{ for large }J\}.`$
Remark It is easy to check that the inequality in (8.1) is equivalent to the requirement that the coefficients of $`f`$ satisfy
$$|c_j|\text{e}^{K|j|},$$
for large $`|j|`$. Another equivalent definition of $`𝒞`$ is
$$𝒞=_{t>0}\text{e}^{tH_{ho}}𝒮,$$
where $`H_{ho}=\mathrm{\Delta }/2+x^2/2`$ is the harmonic oscillator Hamiltonian. The set $`𝒞`$ is also called the set of analytic vectors for the harmonic oscillator Hamiltonian.
Let $`f𝒞`$. We set
$`f_J(y,t)`$ $`=`$ $`{\displaystyle \underset{|j|J}{}}c_j\phi _j(A(t),B(t),ϵ^2,0,0,y),\text{and}`$
$`f(y,t)`$ $`=`$ $`{\displaystyle \underset{j}{}}c_j\phi _j(A(t),B(t),ϵ^2,0,0,y)`$
where the classical quantities $`a(t)`$, $`\eta (t)`$, $`A(t)`$, $`B(t)`$, and $`S(t)`$ correspond to the initial conditions $`a(0)`$, $`\eta (0)`$, $`A(0)=B(0)=\mathrm{I}\mathrm{I}`$, and $`S(0)`$. We consider the construction described in Section 4 corresponding to the initial condition $`g_0(0,y,t)=f_J(y,t)`$, making explicit the dependence on $`J`$ in the notation:
$`\widehat{\mathrm{\Psi }}_{J,N}(w,y,t)`$
$`=`$ $`F(w)\text{e}^{iS(t)/ϵ^2}\text{e}^{i\eta (t)y/ϵ}\left({\displaystyle \underset{n=0}{\overset{N}{}}}ϵ^ng_{n,J}(w,y,t)\mathrm{\Phi }(w,t)+{\displaystyle \underset{n=2}{\overset{N+2}{}}}ϵ^n\varphi _{n,J}^{}(w,y,t)\right).`$
Recall that
$`\widehat{\mathrm{\Psi }}_{J,N}(w,y,0)`$
$`=`$ $`F(w)\text{e}^{iS(0)/ϵ^2}\text{e}^{i\eta (0)y/ϵ}\left(f_J(y,0)\mathrm{\Phi }(w,0)+{\displaystyle \underset{n=2}{\overset{N+2}{}}}ϵ^n\varphi _{n,J}^{}(w,y,0)\right).`$
Let $`\nu >0`$, and consider $`N(ϵ)=[[g^2/ϵ^2]]`$ and $`J(ϵ)=\nu N(ϵ)`$. We define our more general initial conditions as
$`\widehat{\mathrm{\Psi }}_f(w,y,0)`$
$`=`$ $`F(w)\text{e}^{iS(0)/ϵ^2}\text{e}^{i\eta (0)y/ϵ}\left(f(y,0)\mathrm{\Phi }(w,0)+{\displaystyle \underset{n=2}{\overset{N(ϵ)+2}{}}}ϵ^n\varphi _{n,J(ϵ)}^{}(w,y,0)\right),`$
which corresponds, when we get back to the variables $`(X,t)`$, to an initial state $`\widehat{\mathrm{\Psi }}_f(Xa(0),(Xa(0))/ϵ,0)`$ whose projection along the electronic eigenvector $`\stackrel{~}{\mathrm{\Phi }}(X,0)`$ yields a nuclear wave packet of the form $`(\mathrm{\Lambda }_ϵ(a(0),\eta (0))f)(X)`$. Note that the component of the initial state perpendicular to $`\stackrel{~}{\mathrm{\Phi }}(X,0)`$ necessary to achieve exponential accuracy depends on $`ϵ`$. This component is determined by the coefficients of the function $`f`$.
We can now state our result for such general initial conditions
###### Theorem 8.2
Assume the hypotheses of Theorem 4.1 and consider the above constructions. There exist sufficiently small $`g>0`$ and positive constants $`C(g)`$, $`\mathrm{\Gamma }(g)`$, such that with the definition
$$\mathrm{\Psi }_{}(X,t,ϵ)=\widehat{\mathrm{\Psi }}_{J(ϵ),N(ϵ)}(Xa(t),(Xa(t))/ϵ,t),$$
we have
$$\text{e}^{itH(ϵ)/ϵ^2}\mathrm{\Psi }_f(X,0,ϵ)\mathrm{\Psi }_{}(X,t,ϵ)_{L^2(\text{IR}^d,_{\text{el}})}C(g)\text{e}^{\mathrm{\Gamma }(g)/ϵ^2},$$
for all $`t[0,T]`$, as $`ϵ0`$.
Moreover, the result for times $`T\mathrm{ln}(1/ϵ^2)`$ corresponding to Theorem 8.1 is also true for these initial conditions.
Proof: We have
$`\text{e}^{itH(ϵ)/ϵ^2}\mathrm{\Psi }_f(X,0,ϵ)`$
$`=`$ $`\text{e}^{itH(ϵ)/ϵ^2}(\mathrm{\Psi }_f(X,0,ϵ)\mathrm{\Psi }_{}(X,0,ϵ))+\text{e}^{itH(ϵ)/ϵ^2}\mathrm{\Psi }_{}(X,0,ϵ)`$
$`=`$ $`\mathrm{\Psi }_{}(X,t,ϵ)+O(\text{e}^{itH(ϵ)/ϵ^2}\mathrm{\Psi }_{}(X,0,ϵ)\mathrm{\Psi }_{}(X,t,ϵ)_{L^2(\text{IR}^d,_{\text{el}})})`$
$`+O\left(\mathrm{\Psi }_f(X,0,ϵ)\mathrm{\Psi }_{}(X,0,ϵ)_{L^2(\text{IR}^d,_{\text{el}})}\right).`$
By our choice of function $`f`$, the last term is exponentially small in $`1/ϵ^2`$. The remaining norm to estimate corresponds to the situation of Theorem 4.1 in which we let the parameter $`J`$ grow as $`1/ϵ^2`$, according to our choice of $`J(ϵ)`$. But, as in the proof of Theorem 3.6 in for the corresponding result in semiclassical dynamics, we have made the dependence in $`J`$ of all the key estimates explicit. It is enough to go through the proof of theorem 4.1 to check that with $`J=\nu N`$, all arguments can be repeated to get the same $`N`$ and $`ϵ`$ behavior for the estimates on the error terms, (see for details). Hence, we see that for sufficiently small $`g`$, we can approximate the solution corresponding to these generalized initial conditions up to an error of order $`\text{e}^{\mathrm{\Gamma }(g)/ϵ^2}`$. The Ehrenfest time regime is dealt with similarly.
## 9 Technicalities
In this section we give the proofs of the auxiliary lemmas we used in the course of the main argument.
Proof of Lemma 6.1: We first consider the case $`k1`$. By Cauchy’s formula, we can write
$$g^{}(t)=\frac{1}{2\pi i}_\mathrm{\Gamma }\frac{g(s)}{(ts)^2}𝑑s,$$
(9.2)
where $`\mathrm{\Gamma }`$ is the circular contour with center $`t`$ and radius $`{\displaystyle \frac{1}{k+1}}(\delta |\text{Im}t|)`$.
For $`s`$ on $`\mathrm{\Gamma }`$, we have $`(\delta |\text{Im}s|){\displaystyle \frac{k}{k+1}}(\delta |\text{Im}t|)`$. Thus,
$$g(s)Ck^k(\delta |\text{Im}s|)^kCk^k\left[\frac{k}{k+1}(\delta |\text{Im}t|)\right]^k$$
So, by putting the norm inside the integral in (9.2), we have
$`g^{}(t)`$ $``$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle \frac{2\pi }{k+1}}(\delta |\text{Im}t|)Ck^k\left[{\displaystyle \frac{k}{k+1}}(\delta |\text{Im}t|)\right]^k\left[{\displaystyle \frac{1}{k+1}}(\delta |\text{Im}t|)\right]^2`$
$`=`$ $`C(k+1)^{k+1}(\delta |\text{Im}t|)^{k1}.`$
For $`k=0`$ we use the same argument with the radius of $`\mathrm{\Gamma }`$ replaced by $`\alpha (\delta |\text{Im}t|)`$ for any $`\alpha <1`$. This yields the bound
$$g^{}(t)C\alpha ^1(\delta |\text{Im}t|)^1.$$
The lemma follows because $`\alpha <1`$ is arbitrary.
Proof of Lemma 6.4: To prove the quantity $`\nu `$ is finite, we estimate
$`{\displaystyle \underset{\{l:\mathrm{\hspace{0.17em}0}l_i\alpha _i\}}{}}{\displaystyle \frac{1}{(1+|l|)^{d+1}}}{\displaystyle \frac{1}{(1+|\alpha l|)^{d+1}}}`$ (9.8)
$`=`$ $`{\displaystyle \underset{\begin{array}{c}\{l:\mathrm{\hspace{0.17em}0}l_i\alpha _i\}\\ |l|[[\frac{|\alpha |}{2}]]\end{array}}{}}{\displaystyle \frac{1}{(1+|l|)^{d+1}}}{\displaystyle \frac{1}{(1+|\alpha l|)^{d+1}}}`$
$`+{\displaystyle \underset{\begin{array}{c}\{l:\mathrm{\hspace{0.17em}0}l_i\alpha _i\}\\ |l|>[[\frac{|\alpha |}{2}]]\end{array}}{}}{\displaystyle \frac{1}{(1+|l|)^{d+1}}}{\displaystyle \frac{1}{(1+|\alpha l|)^{d+1}}}`$
$``$ $`{\displaystyle \frac{2}{\left(\mathrm{\hspace{0.17em}1}+[[\frac{|\alpha |}{2}]]\right)^{d+1}}}{\displaystyle \underset{\begin{array}{c}\{l:\mathrm{\hspace{0.17em}0}l_i\alpha _i\}\\ |l|[[\frac{|\alpha |}{2}]]\end{array}}{}}{\displaystyle \frac{1}{(1+|l|)^{d+1}}}`$ (9.11)
$``$ $`{\displaystyle \frac{2^{d+2}}{\left(\mathrm{\hspace{0.17em}1}+|\alpha |\right)^{d+1}}}{\displaystyle \underset{\begin{array}{c}\{l:\mathrm{\hspace{0.17em}0}l_i\alpha _i\}\\ |l|[[\frac{|\alpha |}{2}]]\end{array}}{}}{\displaystyle \frac{1}{(1+|l|)^{d+1}}}`$
$``$ $`{\displaystyle \frac{2^{d+2}}{\left(\mathrm{\hspace{0.17em}1}+|\alpha |\right)^{d+1}}}{\displaystyle \underset{l}{}}{\displaystyle \frac{1}{(1+|l|)^{d+1}}}.`$
Thus, $`\nu 2^{d+2}{\displaystyle \underset{l}{}}(1+|l|)^{d1}`$.
To see that the right hand side of this inequality is finite, we note that the number of multi-indices $`l`$ with $`|l|=L`$ is the binomial coefficient $`\left(\begin{array}{c}L+d1\\ d1\end{array}\right)^{}`$, with the convention that $`\left(\begin{array}{c}0\\ 0\end{array}\right)=1`$. Thus,
$`\nu `$ $``$ $`2^{d+2}{\displaystyle \underset{L=0}{\overset{\mathrm{}}{}}}\left(\begin{array}{c}L+d1\\ d1\end{array}\right){\displaystyle \frac{1}{(1+L)^{d+1}}}`$
$`=`$ $`{\displaystyle \frac{2^{d+2}}{(d1)!}}{\displaystyle \underset{L=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(L+d1)(L+d2)\mathrm{}(L+1)}{(L+1)^{d+1}}}.`$
For large $`L`$, $`{\displaystyle \frac{(L+d1)(L+d2)\mathrm{}(L+1)}{(L+1)^{d+1}}}`$ is asymptotic to $`L^2`$, so $`\nu `$ is finite.
Since $`D^\alpha (MN)={\displaystyle \underset{\{l:\mathrm{\hspace{0.17em}0}l_i\alpha _i\}}{}}\left[{\displaystyle \underset{j=1}{\overset{d}{}}}\left(\begin{array}{c}\alpha _j\\ l_j\end{array}\right)\right]\left(D^lM\right)\left(D^{(\alpha l)}N\right)`$, we have
$`\left(D^\alpha (MN)\right)(x)`$ (9.25)
$``$ $`{\displaystyle \underset{\{l:\mathrm{\hspace{0.17em}0}l_i\alpha _i\}}{}}\left[{\displaystyle \underset{j=1}{\overset{d}{}}}\left(\begin{array}{c}\alpha _j\\ l_j\end{array}\right)\right]m(x)n(x)a(x)^{|\alpha +p+q|}{\displaystyle \frac{(l+p)!}{(1+|l|)^{d+1}}}{\displaystyle \frac{(\alpha l+q)!}{(1+|\alpha l|)^{d+1}}}`$
$`=`$ $`m(x)n(x)a(x)^{|\alpha +p+q|}(\alpha +p+q)!`$
$`\times {\displaystyle \underset{\{l:\mathrm{\hspace{0.17em}0}l_i\alpha _i\}}{}}\left[{\displaystyle \underset{j=1}{\overset{d}{}}}\left(\begin{array}{c}\alpha _j\\ l_j\end{array}\right)\left(\begin{array}{c}\alpha _j+p_j+q_j\\ l_j+p_j\end{array}\right)^1\right]{\displaystyle \frac{1}{(1+|l|)^{d+1}(1+|\alpha l|)^{d+1}}}.`$
Since $`\left(\begin{array}{c}\alpha _j+p_j+q_j\\ l_j+p_j\end{array}\right)\left(\begin{array}{c}\alpha _j+q_j\\ l_j\end{array}\right)\left(\begin{array}{c}\alpha _j\\ l_j\end{array}\right)`$, we therefore have
$`\left(D^\alpha (MN)\right)(x)`$ $``$ $`m(x)n(x)a(x)^{|\alpha +p+q|}(\alpha +p+q)!`$
$`\times {\displaystyle \underset{\{l:\mathrm{\hspace{0.17em}0}l_i\alpha _i\}}{}}{\displaystyle \frac{1}{(1+|l|)^{d+1}(1+|\alpha l|)^{d+1}}}.`$
$``$ $`m(x)n(x)\nu a(x)^{|\alpha +p+q|}{\displaystyle \frac{(\alpha +p+q)!}{(1+|\alpha |)^{d+1}}}.\text{ }`$
Proof of Lemma 6.3: If $`f(t)`$ satisfies $`f(t)C|t|^p\text{dist}(t)^k`$, for all $`t\mathrm{\Omega }`$, there exists $`g(t)`$ analytic in $`\mathrm{\Omega }`$, such that $`f(t)=t^pg(t)`$ and $`g(t)C\text{dist}(t)^k`$. We use the integration path from $`0`$ to $`t\mathrm{\Omega }`$ parametrized by $`\gamma (u)=tu`$, with $`u[0,1]`$, to compute
$`{\displaystyle _0^t}f(s)𝑑s={\displaystyle _0^1}f(tu)𝑑u={\displaystyle _0^1}t(tu)^pg(tu)𝑑u`$ (9.26)
$``$ $`C|t|^{p+1}{\displaystyle _0^1}{\displaystyle \frac{u^p}{\text{dist}(tu)^k}}𝑑uC{\displaystyle \frac{|t|^{p+1}}{p+1}}\text{dist}(t)^k,`$
since, by assumption, $`\text{dist}(ut)`$ is a decreasing function of $`u`$.
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# May 3, 2000OAT Int. Rep. 71/00OAT Pub. Num. 2140Submitted to PASP Data Streams from the Low Frequency Instrument On-Board the Planck Satellite: Statistical Analysis and Compression Efficiency
## 1 Introduction and Scanning Strategy
The Planck satellite (formerly COBRAS/SAMBA, Bersanelli et al. (1996)), which is planned to be launched in 2007, will produce full sky CMB maps with high accuracy and resolution over a wide range of frequencies (Mandolesi et al. (1998a); Puget et al. (1998)). Table 1 summarizes the basic properties of LFI aboard Planck. The reported sensitivities per resolution element – i.e. a squared pixel with side equal to the Full Width at Half Maximum (FWHM) extent of the beam –, in terms of antenna temperature, represents the goals of LFI for 14 months of routine scientific operations) as recently revised by the LFI Consortium (Mandolesi et al. (1999)).
The limited bandwidth reserved to the downlink of scientific data calls for huge lossless compression, theoretical upper limit being about four (Maris et al. (1999)). Careful simulations are demanded to quantify the capability of true compressors for “realistic” synthetic data and improve the theoretical analysis, including CMB signal (monopole, dipole and anisotropies), foregrounds and instrumental noise.
During the data acquisition phase the Planck satellite will rotate at a rate of one circle per minute around a given spin axis that changes its direction every hour (of 2.5 on the ecliptic plane in the case of simple scanning strategy), thus observing the same circle on the sky for 60 consecutive times (Mandolesi et al. (1998a, b)). LFI will produce continuous data streams of temperature differences between the microwave sky and a set of on-board reference sources; both differential measurements and reference source temperatures must be recorded.
The LFI Proposal assumes a sampling time $`\tau _\mathrm{s}7`$ msec for each detector (Mandolesi et al. (1998a)), thus calling for a typical data rate of $`260`$ Kb/sec, while the allocated bandwidth to download Planck data to ground is in total $`60`$ Kb/sec. Assuming the total bandwidth to be equally split between instruments, $`30`$ Kb/sec on the average would be assigned to LFI asking for a compression of about a factor $`8.4`$. Data have to be downloaded without information losses and by minimizing scientific processing on board.
A possible solution would be to adapt the sampling rate to the angular resolution specific for each frequency. This should allow to save about up to a factor $`9`$ for the 30 GHz channel, but since only $`7\%`$ of the samples come from such channel (see table 1) the overall reduction in the final data rate would be $`17\%`$.
On the other hand, it is unlikely that the bandwidth for the downlink channel may be enhanced to solve the bandwidth problem, since the ground facilities are shared between different missions and there is the need to minimize possible cross-talks between the instrument and the communication system.
With the aim of optimizing of the transmission bandwidth dedicated to the downlink of LFI data from the Planck spacecraft to the FIRST/Planck Ground Segment, we analyze in detail the role that can be played by lossless compression of LFI data before they are sent to Earth.
We apply different compression algorithms to suitable sets of Planck LFI simulated data streams generated by considering different combinations of astrophysical and instrumental signals and for different instrumental characteristics and detection electronics.
The first considered contribution is that introduced by receiver noise: we consider here the case of pure white noise and of white noise coupled to $`1/f`$ noise with different knee frequencies. The reference load temperature is assumed to be 20 K for present tests; because of the strong dependence of the $`1/f`$ noise on the load temperature, this can be considered a worst case, since the actual baseline reference load is of 4 K.
Different sky signal sources are subsequently added to the receiver noise: CMB fluctuations, CMB dipole, Galaxy emission and extragalactic point sources. The signal from the different sky components are convolved with the corresponding antenna pattern shapes, assumed to be symmetric and gaussian with the FWHM reported in Table 1.
We generate simulated data streams at the two extreme frequency channels, 30 GHz and 100 GHz and consider data streams with different time lengths.
Regarding the detection electronics, we explore different signal offset and scaling.
The large number of above combinations was systematically explored using an automated program generator as described by Maris & Staniszkis (1998).
In Section 2 we characterize quantitatively the LFI signal component by component. Section 3 we discuss how the acquisition chain is modeled to perform compression simulations. A theoretical analysis of the compression efficiency is presented in section 4. While section 5 is devoted to the analysis of the signal statistics. The subject of quantization error is illustrated in section 6. The experimental protocol and results about compression are reported in section 7. Further constraints on the on-board data compression are reported in section 8. A proposal for an alternative coding method is made in section 9. The overall compression rate is estimated in section 10. Conclusions are in section 11. Appendix A is included to further illustrate the estimation of the overall compression rate.
## 2 Characterization of Planck/LFI signal components
The simulated cosmological and astrophysical components are generated according to the methods described in Burigana et al. (1998b) and the data stream and noise generation as in Burigana et al. (1997b), Seiffert et al. (1997) and Maino et al. (1999). We summarize here below the basic points.
$``$ Modeling the CMB pattern – The CMB monopole and dipole have been generated by using the Lorentz invariance of photon distribution functions, $`\eta `$, in the phase space (Compton–Getting effect): $`\eta _{obs}(\nu _{obs},\stackrel{}{n})=\eta _{CMB}(\nu _{CMB})`$, where $`\nu _{obs}`$ is the observation frequency, $`\nu _{CMB}=\nu _{obs}(1+\stackrel{}{\beta }\times \stackrel{}{n})/\sqrt{1\beta ^2}`$ is the corresponding frequency in the CMB rest frame, $`\stackrel{}{n}`$ is the unit vector of the photon propagation direction and $`\stackrel{}{\beta }=\stackrel{}{v}/c`$ the observer velocity. A blackbody spectrum at $`T_0=2.725`$ K (Mather et al. (1999)) is assumed for $`\eta `$. For gaussian models, the CMB anisotropies at $`l2`$ can be simulated by following the standard spherical harmonic expansion (see, e.g., Burigana et al. (1998a) or by using FFT (Fast Fourier Transform) techniques which take advantage of equatorial pixelisations (Muciaccia et al. (1997))).
$``$ Modeling the Galaxy emission – The Haslam map at 408 MHz (Haslam et al. (1982)) is the only full-sky map currently available albeit large sky areas are sampled at 1420 MHz (Reich & Reich (1986)) and at 2300 MHz (Jonas et al. (1998)). To clean these maps from free-free emission we use a 2.7 GHz compilation of $``$ 7000 HII sources (Witebsky (1978)), private communication) at resolution of $``$ 1. They are subtracted for modelling the diffuse components and then re-added to the final maps. We use a spectral index $`\beta _{ff}=2.1`$ from 2.7 to 1 GHz and $`\beta _{ff}=0`$ below 1 GHz. We then combine the synchrotron maps producing a spectral index map between 408-2300 MHz with a resolution of $`\stackrel{_<}{_{}}2^{}÷3^{}`$ ($`<\beta _{sync}>2.8`$). This spectral index map is used to scale the synchrotron component down to $``$ 10 GHz. In fact, for typical (local) values of the galactic magnetic field ($`2.5\mu `$G), the knee in the electron energy spectrum in cosmic rays ($``$ 15 Gev) corresponds to $``$ 10 GHz (Platania et al. (1998)). From the synchrotron map obtained at 10 GHz and the DMR 31.5 GHz map we derive a high frequency spectral index map for scaling the synchrotron component up to Planck frequencies. These maps have a poor resolution and the synchrotron structure needs to be extrapolated to Planck angular scales. An estimate of the synchrotron angular power spectrum and of its spectral index, $`\gamma `$ ($`C_ll^\gamma `$), has been provided by Lasenby et al. (1998); we used $`\gamma =3`$ for the angular structure extrapolation (Burigana et al. (1998a)). Schlegel (Schlegel et al. (1998)) provided a map of dust emission at 100$`\mu `$m merging the DIRBE and IRAS results to produce a map with IRAS resolution ($`7^{}`$) but with DIRBE calibration quality. They also provided a map of dust temperature, $`T_d`$, by adopting a modified blackbody emissivity law, $`I_\nu B_\nu (T_d)\nu ^\alpha `$, with $`\alpha =2`$. This can be used to scale the dust emission map to Planck frequencies using the dust temperature map as input for the $`B_\nu (T_d)`$ function. Unfortunately the dust temperature map has a resolution of $`1^{}`$; again, we use an angular power spectrum $`C_ll^3`$ to scale the dust skies to the Planck proper resolution. Merging maps at different frequencies with different instrumental features and potential systematics may introduce some internal inconsistencies. More data on diffuse galactic emission, particularly at low frequency, would be extremely important.
$``$ Modeling the extragalactic source fluctuations – The simulated maps of point sources have been created by an all–sky Poisson distribution of the known populations of extragalactic sources in the $`10^5<S(\nu )<10`$ Jy flux range exploiting the number counts of Toffolatti et al. (1998) and neglecting the effect of clustering of sources. The number counts have been calculated by adopting the Danese et al. (1987) evolution model of radio selected sources and an average spectral index $`\alpha =0`$ for compact sources up to $`200`$ GHz and a break to $`\alpha =0.7`$ at higher frequencies (see Impey & Neugebauer (1988); De Zotti & Toffolatti (1998)), and by the model C of Franceschini et al. (1994) updated as in Burigana et al. (1997a), to account for the isotropic sub-mm component estimated by Puget et al. (1996) and Fixsen et al. (1996). At bright fluxes, far–IR selected sources should dominate the number counts at High Frequency Instrument (HFI) channels for $`\nu \stackrel{_>}{_{}}300`$ GHz, whereas radio selected sources should dominate at lower frequencies (Toffolatti et al. (1998)).
$``$ Instrumental noise – The white noise depends on instrumental performances (bandwidth $`\mathrm{\Delta }\nu `$, system temperature $`T_{sys}`$), on the observed sky signal, $`T_{sky}`$, dominated by CBM monopole, and on the considered integration time, $`\tau `$, according to:
$$\mathrm{\Delta }T_{\mathrm{wn}}=\frac{\sqrt{2}(T_{sys}+T_{sky})}{\sqrt{\mathrm{\Delta }\nu \tau }}.$$
(1)
Under certain idealistic assumptions, Burigana et al. (1997b) and Seiffert et al. (1997) provide analytical estimates for the knee frequency, $`f_k`$, of LFI radiometers; it is predicted to critically depend also on the load temperature, $`T_{load}`$, according to:
$$f_k=\frac{A^2\mathrm{\Delta }\nu }{8}(1r)^2\left(\frac{T_{sys}}{T_{sys}+T_{sky}}\right)^2,$$
(2)
where $`r=(T_{sky}+T_{sys})/(T_{load}+T_{sys})`$ and $`A`$ is a constant, depending on the state of art of radiometer technology, which has to be minimized for reducing via hardware the knee frequency (current estimates are $`A1.8\times 10^5`$ for 30 and 44 GHz radiometers and $`A2.5\times 10^5`$ for 70 and 100 GHz).
Recent experimental results from Seiffert (private communication Seiffert (1999)) show knee frequency values of this order of magnitude, confirming that the present state of art of the radiometer technology is close to reach the ideal case.
A pure white noise stream can be easily generated by employed well tested random generator codes and normalizing their output to the white noise level $`\mathrm{\Delta }T_{\mathrm{wn}}`$. A noise stream which takes into account both white noise and $`1/f`$ noise can be generated by using FFT methods. After generating a realisation of the real and imaginary part of the Fourier coefficients with spectrum defined $`S_{noise}(f)(1+f_k/f)`$, we transform them and obtain a real noise stream which has to be normalized to the white noise level $`\mathrm{\Delta }T_{\mathrm{wn}}`$ (Maino et al. (1999)).
$``$ Modeling the observed signal – We produce full sky maps, $`T_{sky}`$, by adding the antenna temperatures from CMB, Galaxy emission and extragalactic source fluctuations. Planck will perform differential measurements and not absolute temperature observations; we then represent the final observation in a given $`i`$-$`th`$ data sample in the form $`T_i=R_i(T_{sky,i}+N_iT_{x,i}^r)`$, where $`N_i`$ is the instrumental noise generated as described above.
$`T_{x,i}^r`$ is a reference temperature subtracted in the differential data and $`R_i`$ is a constant which accounts for the calibration. Of course, the uncertainty on $`R_i`$ and the non reduced time variation of $`T_{x,i}^r`$ have to be much smaller than the Planck nominal sensitivity. Thus, we generate the “observed” map assuming a constant value, $`T_x^r`$, of $`T_{x,i}^r`$ for all the data samples. We note that possible constant small off-sets in $`T_x^r`$ could be in principle accepted, not compromising an accurate knowledge of the anisotropy pattern. We arbitrarily generate the “observed” map with $`R_i=R=1`$ for all the data samples.
## 3 A model of Acquisition Chain
To test rigorously the efficiency of different compressors the best solution is to generate a realistically simulated signal for different mission hypotheses and apply to them the given compressors. To be realistical the simulation of the signal generation should contain both astrophysical and instrumental effects. It would be helpful that the final simulation would be able to given a hint about the influence of the various signal components and their variance. Of course it is useless to reproduce in full detail the LFI to obtain a signal simulation accurate enough to test compressors. A simplified model of the LFI, its front-end electronics and its operations will be enough.
At the base of the simplified model is the concept of acquisition pipeline. This pipeline is composed by all the modules which process the astrophysical signal: from its collection to the production of the final data streams which are compressed and then sent to Earth. In the real LFI, the equivalent of the acquisition pipeline may be obtained following the flow of the astrophysical information, from the telescope through the front-end electronics and the main Signal Processing Unit (SPU) to the memory of the Data Processing Unit (DPU) which is in charge to downlink it to the computer of the spacecraft and then to Earth. The acquisition pipeline is represented in figure 1. Since its purpose is to describe the signal processing and its parameters it must not be regarded as a representation of the true on-board electronics since some functionalities may be shared between different real modules. In this scheme Front End operations of the true LFI are assigned to the first simulation level, while on-board processing and compression to the second one.
The simulated microwave signal from the sky is collected and compared with the temperature of a reference load which, in our simulations, is supposed to have exactly the CMB temperature $`T_0=2.725`$ K (Mather et al. (1999)) <sup>1</sup><sup>1</sup>1Alternatively, sky the reference-load signals may be sampled separately and then $`\mathrm{\Delta }T`$ may be compute numericaly by the DPU.. The difference $`\mathrm{\Delta }T`$ expressed in $`\mu `$K is sampled along a scan circle producing a data stream of 60 scan circles with 8640 samples (pointings).
Signal detection is simulated by Bersanelli et al. (1996); Maris et al. (1998, 1999)
$$V_{\mathrm{out}}=\mathrm{AFO}+\mathrm{VOT}\mathrm{\Delta }T,$$
(3)
where $`V_{\mathrm{out}}`$ is the detection chain output in Volts, $`\mathrm{VOT}`$ is the antenna temperature to the detector voltage conversion factor ($`0.5`$V/K $`VOT+1.5`$ V/K) while $`\mathrm{AFO}`$ is a detection chain offset ($`5`$V $`\mathrm{AFO}+5`$V). Of course in our simulation this offset takes into account all offset sources, including variations of the reference temperature, and not only of the electrical offset. Similarly the $`\mathrm{VOT}`$ factor takes into account also differences among the different detectors which affect the calibration of the temperature/voltage relation. The range for $`\mathrm{VOT}`$ and $`\mathrm{AFO}`$ is large enough to include the whole set of nominal instrumental configurations, allowing also for somewhat larger and smaller values.
The analog to digital conversion (ADC) is described by the formula:
$$V_{\mathrm{out}}^{\mathrm{adu}}(\mathrm{adu})=\mathrm{trunc}\left(2^{N_{\mathrm{bits}}}\frac{V_{\mathrm{out}}V_{\mathrm{min}}}{V_{\mathrm{max}}V_{\mathrm{min}}}\right),$$
(4)
where $`\mathrm{trunc}(.)`$ is the decimal truncation operator, $`N_{\mathrm{bits}}`$ is the number of quantization bits produced by the ADC, while $`V_{\mathrm{min}}`$ and $`V_{\mathrm{max}}`$ are the lower and upper limits of the voltage scale accepted in input by the ADC. In our case: $`N_{\mathrm{bits}}=16`$ bits, $`V_{\mathrm{min}}=10`$ V, $`V_{\mathrm{max}}=+10`$ V. So the quantization unit “adu” (analog/digital unit) is
$$1\mathrm{adu}=\frac{V_{\mathrm{max}}V_{\mathrm{min}}}{2^{N_{\mathrm{bits}}}}$$
(5)
or in terms of antenna temperature the quantization step is
$$\mathrm{\Delta }=\frac{V_{\mathrm{max}}V_{\mathrm{min}}}{2^{N_{\mathrm{bits}}}\mathrm{VOT}}$$
(6)
for a typical $`\mathrm{VOT}=1`$ V/K, $`N_{\mathrm{bits}}=16`$ bits, $`1\mathrm{\Delta }3\times 10^4`$ K/adu. After digitization the simulated signal is written into a binary file of 16 bits integers and sent to the compression pipeline.
The simplified LFI is composed of four acquisition pipelines, one for each frequency, each one being representative of the set of devices which form the full detection channel for the given frequency. The overall data-rate after loss-less compression for LFI should be obtained summing the contribution expected from each detector. Since in the real device each radiometer for a given frequency channel, will be characterized by different values of $`\mathrm{VOT}`$ and $`\mathrm{AFO}`$, the distribution of these parameters has to be taken in account computing the overall compression efficiency. In particular a greater attention should be devoted to the distribution of the $`\mathrm{VOT}`$ parameter since the compression efficiency is particularly sensitive to it. However, since the distribution of operating conditions and instrumental parameters are not yet fully defined, we assumed that all the detectors belongin to a given frequency channel are identical <sup>2</sup><sup>2</sup>2But see section 10 and the appendix for a more detailed discussion. and located at the telescope focus.
## 4 An Informal Theoretical Analysis About the Compression Efficiency
An informal theoretical analysis may be helpful to evaluate the maximum lossless compression efficiency expected from LFI and to discuss the behaviour of the different compressors. For further details we remind the reader to Nelson & Gailly (1996).
Data compression is based on the partition of a stream of bits into short chunks, represented by strings of bits of fixed length $`N_{\mathrm{bits}}`$, and to code each string of bits $`S_{\mathrm{In}}`$ into another string $`S_{\mathrm{Out}}`$ whose length $`N_{\mathrm{bits}}^{\mathrm{out}}`$ is variable and, in principle, shorter than $`S_{\mathrm{In}}`$. In this scheme, when the string of bits represents a message, the possible combinations of bits in $`S_{\mathrm{In}}`$ represents the symbols by which the message is encoded. From this description the compression operation is equivalent to map the input string set $`\{S_{\mathrm{In}}\}`$ into an output string set $`\{S_{\mathrm{Out}}\}`$ through a compressing function $`_{\mathrm{Comp}}`$. A compression algorithm is called lossless when it is possible to reverse the compression process reconstructing the $`S_{\mathrm{In}}`$ string from $`S_{\mathrm{Out}}`$ through a decompression algorithm. So the condition for a compression programs to be lossless is that the related $`_{\mathrm{Comp}}`$ is a one-to-one application of $`\{S_{\mathrm{In}}\}`$ into $`\{S_{\mathrm{Out}}\}`$. In this case the decompressing algorithm is the inverse function of $`_{\mathrm{Comp}}`$. Of course in the general case it is not possible to have at the same time lossless compression and $`N_{\mathrm{bits}}>N_{\mathrm{bits}}^{\mathrm{out}}`$ for any string in the input set. The problem is solved assuming that the discrete distribution $`P(S_{\mathrm{In}})`$ of strings belonging to the input stream of bits is not flat but that a most probable string exists. So a good $`_{\mathrm{Comp}}`$ will assign the shortest $`S_{\mathrm{Out}}`$ to the most probable $`S_{\mathrm{In}}`$ and, the least probable the input string, the longest the output string. In the worst case output strings longer than the input string will be assigned to those strings of $`\{S_{\mathrm{In}}\}`$ which are least probable. With this statistical tuning of the compression function the final length of the compressed stream will be shorter than the original length, the averaged length of $`S_{\mathrm{Out}}`$ being:
$$\overline{N_{\mathrm{bits}}^{\mathrm{out}}}=\underset{S_{\mathrm{In}}\{S_{\mathrm{In}}\}}{}P(S_{\mathrm{In}})N_{\mathrm{bits}}^{\mathrm{out}}(_{\mathrm{Comp}}(S_{\mathrm{In}})).$$
(7)
Several factors affect the efficiency of a given compressor, in particular best performances are obtained when the compression algorithm is tuned on the specific distribution of symbols. Since the symbol distribution depends on $`N_{\mathrm{bits}}`$ and on the specific input stream, an ideal general-purpose self-adapting compressor should be able to perform the following operations: i) acquire the full bit stream (in the hypothesis it has a finite length) and divide it in chunks of length $`N_{\mathrm{bits}}`$, ii) perform a frequency analysis of the various symbols, iii) create an optimized coding table which associates to each $`S_{\mathrm{In}}`$ a specific $`S_{\mathrm{Out}}`$, iv) perform the compression according to the optimized coding table, v) send the coding table to the uncompressing program together with the compressed bit stream. The uncompressing program will restore the original bit stream using the associated optimized coding table.
In practice in most cases the chunks size $`N_{\mathrm{bits}}`$ is hardwired into the compressing code (typically $`N_{\mathrm{bits}}=8`$ or 16 bits), also the fine tuning of the coding table for each specific bit stream is too expensive in terms of computer resources to be performed in this way, and the same holds for coding table transmission. So there are compressors which work as if the coding table or, equivalently, the compression function is fixed. In this way the bit stream may be compressed chunk by chunk by the compressing algorithm which will act as a filter. Other compressors perform the statistical tuning on a small set of chunks taken at the beginning of the stream, and then apply the same coding table to the full input stream. In this case the compression efficiency will be sensitive to the presence of correlations between difference parts of the input stream. In this respect self-adaptive codes may be more effective than non-adaptive ones, if their adapting strategy is sensitive to the kind of correlations in the input stream.
On the other hand other solutions may be adopted to obtain a good compromise between computer resources and compression optimization. For example all of the previous compressors are called static since the coding table is fixed in one way or the other at the beginning of the compression process and then used all over the input stream. Another big class of self-adaptive codes is represented by dynamical self-adaptive compressors, which gain the statistical knowledge about the signal as the compression proceeds changing time by time the coding table. Of course these codes compress worse at the beginning and better at the end of the data stream, provided its statistical properties are stationary. They are also able to self-adapt to remarkable changes in the characteristics of the input stream, but only if these changes may be sensed by the adapting code. Otherwise the compressor will behave worse than a well-tuned static compressor. Moreover, if the signal changes frequently, it may occur that the advantage of the dynamical self adaptability is compensated by the number of messages added to the output stream to inform the decompressing algorithm of the changes occurred to the coding table. Last but not least, if some error occurs during the transmission of the compressed stream and the messages about changes in the coding table are lost, it will be impossible to correctly restore it at the receiving station. This problem may be less severe for a static compressor since, as an example, it is possible to split the output stream in packets putting stop codes and storing the coding table on-board until a confirmation message from the receiving station is sent back to confirm the correct transmission.
It is then clear that each specific compression algorithm is statistically optimized for a given kind of input stream with its own statistical properties. So to obtain an optimized compressor for LFI it is important to properly characterize the statistics of the signal to be compressed and to test different existing compressors in order to map the behaviour of different compression schemes using realistically simulated signals and, as soon as possible, the true signals produced by the LFI electrical model.
In order to evaluate the performances of different compression scheme we considered the Compression Rate $`C_\mathrm{r}`$ defined as:
$$C_\mathrm{r}=\frac{L_u}{L_c}$$
(8)
where $`L_u`$ is the length of the input string in bytes and $`L_c`$ is the length of the output string in bytes <sup>3</sup><sup>3</sup>3Often compressors are evaluated looking at the compression efficiency $`\eta _c=1/C\mathrm{r}`$ but we considered $`C_\mathrm{r}`$ more effective for our purposes.. Other important estimators of to evaluate the performances of a given compression code are the memory allocation and the compression time. Both of them must be evaluated working on the final model of the on board computer. Since this component is not fully defined for the Planck/LFI mission, in this work we neglect these aspects of the problem.
The measure represented by one of the 8640 samples which form one scan circle is white noise dominated, the r.m.s. $`\sigma _T`$ being about a factor of ten higher then the CMB fluctuations signal. If so, at the first approximation it is possible to assume the digitized data stream from the front-end electronics as a stationary time serie of independent samples produced by a normal distributed white noise generator. In such situation symbols are represented by the quantized signal levels, and it is easy to infer the best coding table and by the information theory the expected compression rate for an optimized compressor is promptly estimated (Gaztñaga et al. (1998)). In our notation, for a zero average signal:
$$C_\mathrm{r}^{\mathrm{Th}}=\frac{N_{\mathrm{bits}}\mathrm{ln}2}{\mathrm{ln}(\sqrt{2\pi e}\sigma _l/\mathrm{adu})+\mathrm{ln}\mathrm{VOT}}$$
(9)
where $`\sigma _l`$ is the r.m.s. of the sampled signal <sup>4</sup><sup>4</sup>4It has to be noted that eq. (9) is an approximated formula which is rigorously valid when $`\sigma _l/\mathrm{adu}1`$..
From Eq. (9) it is possible to infer that the higher is the $`\mathrm{VOT}`$, (i.e. higher is the $`\mathrm{\Delta }T`$ resolution) the worse is the compression rate, as already observed in Maris et al. (1998), Maris et al. (1999). The reason being the fact that as $`\mathrm{VOT}`$ is increased the number of quantization levels (i.e. of symbols) to be coded is increased and their distribution becomes more flat increasing $`\overline{N_{\mathrm{bits}}^{\mathrm{out}}}`$. Assuming that all the white noise is thermal in origin $`\sigma _l\sigma _T2\times 10^3`$ K. With the $`\mathrm{adu}`$ defined in equation (5) together with the typical values of $`V_{\mathrm{min}}`$ and $`V_{\mathrm{max}}`$ assumed therein and $`N_{\mathrm{bits}}=16`$ bits we have $`C_\mathrm{r}^{\mathrm{Th}}11.09/(3.30+\mathrm{ln}\mathrm{VOT})`$. In conclusion, for $`\mathrm{VOT}=0.5`$, $`1.0`$, $`1.5`$ V/K the $`C_\mathrm{r}^{\mathrm{Th}}`$ is respectively $`4.26`$, $`3.36`$, $`3.00`$. In addition figure 2 represents the effect of a reduction of $`N_{\mathrm{bits}}`$ on $`C_\mathrm{r}^{\mathrm{Th}}`$ compared to $`C_\mathrm{r}^{\mathrm{Th}}`$ for $`N_{\mathrm{bits}}=16`$.
## 5 Statistical Signal Analysis
A realistic estimation of the compression efficiency must be based on a quantitative analysis of the signal statistics, which includes: statistics of the binary representation (section 5.1), entropy section 5.2) and normality tests (section 5.3).
### 5.1 Binary Statistics
Most of the off-the-shelf compressors considered here do not handle 16 bits words, but 8 bits words. The 16 bits samples produced by the adc unit are splitted into two consecutive 8 bits (1 byte) words labeled: most significant bits (MSB) word and least significant bits (LSB) word. To properly understand the compression efficiency limits it is important to understand the statistical distribution of 8 bits words composing the quantized signal from LFI.
Figure 3 represents the frequency distribution of symbols when the full data stream of 60 scan circles is divided into 8 bits words. Since for most of the samples the range spans over $`64`$ levels (5 bits) only the bytes corresponding to the MSB words assume a limited range of values producing the narrow spike in the figure. The belt shaped distribution at the edges is due to the set of LSB words. The distributions are quite sensitive to the quantization step, but do not change too much with the signal composition, the largest differences coming from the cosmological dipole contribution.
From the distribution in figure 3 one may wonder if it would not be possible to obtain a more effective compression splitting the data stream into two substreams: the MSB substream (with compression efficiency $`C_\mathrm{r}^{\mathrm{MSB}}`$) and the LSB substream (with compression efficiency $`C_\mathrm{r}^{\mathrm{LSB}}`$). Since the two components are so different in their statistics, with the MSB substream having an higher level of redundancy than the original data stream, it would be reasonable to expect that the final compression rate $`2/(1/C_\mathrm{r}^{\mathrm{MSB}}+1/C_\mathrm{r}^{\mathrm{LSB}})`$ be greater than the compression rate obtained compressing directly the original data stream. We tested this procedure taking some of the compressors considered for the final test. From these tests It is clear that $`C_\mathrm{r}^{\mathrm{MSB}}>>C_\mathrm{r}`$ but since most of the redundancy of the original data stream is contained in the MSB substream the LSB substream can not be compressed in an effective way, as a result $`C_\mathrm{r}^{\mathrm{LSB}}<C_\mathrm{r}`$ and $`2/(1/C_\mathrm{r}^{\mathrm{MSB}}+1/C_\mathrm{r}^{\mathrm{LSB}})\stackrel{_<}{_{}}C_\mathrm{r}`$. So the best way to perform an efficient compression is to apply the compressor to the full stream without performing the MSB / LSB separation. Apart from these theoretical considerations, we performed some tests with our simulated data stream confirming these result.
### 5.2 Entropy Analysis
Equation (9) is valid in the limit of a continuous distribution of quantization levels. Since in our case the quantization step is about one tenth of the signal rms this is no longer true. To properly estimate the maximum compression rate attainable from these data we evaluate the entropy of the discretized signal using different values of the $`\mathrm{VOT}`$.
Our entropy evaluation code takes the input data stream and determines the frequency $`f_s`$ of each symbol $`s`$ in the quantized data stream and computing the entropy as: $`_sf_s\mathrm{log}_2f_s`$ where $`s`$ is the symbol index. In our simulation we take both 8 and 16 bits symbols ($`s`$ spanning over $`0`$, $`\mathrm{}`$, $`255`$ and $`0`$, $`\mathrm{}`$, $`65535`$). Since in our scheme the ADC output is 16 bits, we considered 8 bits symbols entropy both for the LSB and MSB 8 bits word and 8 bits entropy after merging the LSB and MSB significant bits set.
As expected, since $`\mathrm{AFO}`$ merely shifts the quantized signal distribution, entropy does not depend on $`\mathrm{AFO}`$. For this reason we take $`\mathrm{AFO}=0`$ V, i.e., no shift.
Table 2 reports the 16 bits entropy as a function of $`\mathrm{VOT}`$, composition and frequency. As obvious entropy, i.e. information content, increases increasing $`\mathrm{VOT}`$ i.e. quantization resolution. The entropy $`H`$ distribution allows to evaluate the $`C_\mathrm{r}`$ r.m.s. espected from different data streams realizations:
$$\mathrm{RMS}(C_\mathrm{r})C_\mathrm{r}\frac{\mathrm{RMS}(H)}{H}.$$
(10)
Since data will be packed in chuncks of finite length it is important not only to study the entropy distribution for the entire data-stream, which will give an indication of the overall compressibility of the data stream as a wall, but also the entropy distribution for short packets of fixed length. So each data stream was splitted into an integer number of chunks of fixed length $`l_{\mathrm{chunck}}`$. For each chunck the entropy was measured, and the corresponding distribution of entropies for the given $`L_{\mathrm{chunck}}`$ as its mean and rms was obtained. We take $`l_{\mathrm{chunk}}=16`$, $`32`$, $`64`$, $`135`$, $`8640`$, $`17280`$ 16-bits samples, so each simulated $`8640\times 60`$ data stream will be splitted into 32400, 16200, 8100, 3840, 60, 30 chuncks. Small chunck sizes are introduced to study the entropy distribution as seen by most of the true compressors which do not compress one circle (8640 samples) at a time. Long chuncks distributions are usefull to understand the entropy distribution for the overall data-stream.
The entropy distribution per chunck is approximately described by a normal distribution (see figure 4), so the mean entropy and its r.m.s. are enough to characterize the results. Not however that the corresponding distribution of compression rates is not exactly normally distributed, however for the sake of this analysis we will assume that even the $`C_\mathrm{r}`$ distribution is normally distributed.
The mean entropy measured over one scan circle ($`l_{\mathrm{chunk}}=8640`$ samples) coincides with the entropy measured for the full set of 60 scan circles, the entropy r.m.s. being of the order of $`10^2`$ bits. Consequently the expected r.m.s. for $`C_\mathrm{r}`$ compressing one or more circles at a time will be less than $`1\%`$.
The mean entropy and its rms are not independent quantities. Averaged entropy decreases as $`L_{\mathrm{chunck}}`$ decreases, but correspondingly the entropy r.m.s. increases. As a consequence the averaged $`C_\mathrm{r}`$ decreases decreasing $`L_{\mathrm{chunck}}`$, but the fraction of chunks in which the compressor performs significantly worst than in average increases. The overall compression rate, i.e. the $`C_\mathrm{r}`$ referred to the full mission, beeing affected by them.
### 5.3 Normality Tests
Since normal distribution of signals is assumed in 4 it would be interesting to fix how much the digitized signal distribution deviates from the normality. Also it would be important to characterize the influence of the 1/f noise and of the other signal components, especially the cosmic dipole, in the genesis of such deviations. To obtain an efficient compression it would be important that the samples are as more as possible statistically uncorrelated and normally distributed. In addition one should make sure that the detection chain does not cause any systematic effect which will introduce spurious non normal distributed components. This is relevant not only for the compression problem itself, which is among the data processing operations the least sensitive to small deviations from the normal distribution, but also in view of the future data reduction, calibration and analysis. For them the hypothesis of normality in the signal distribution is very important in order to allow a good separation of the foreground components. Last but not least, the hypothesis of conservation of normality along the detection chain, is important for the scientific interpretation of the results, since the accuracy expected from the Planck/LFI experiment should allow to verify if really the distribution of the CMB fluctuations at $`l\stackrel{_>}{_{}}14`$ is normal, as predicted by the standard inflationary models, or as seems suggested by recent 4 years COBE/DMR results (Bromley & Tegmark (1999); Ferreira, Górsky, Magueijo (1999)).
For this reason a set of normality tests was applied to the different components of the simulated signal before and after digitization in order to characterize the signal statistics and its variation along the detection process. Of course this work may be regarded as a first step in this direction, a true calibration of the signal statistics will be possible only when the front end electronics simulator will be available. Those tests have furthermore the value of a preparation to the study of the true signal.
Normality tests were applied on the same data streams used for data compression. Given on board memory limits, it is unlikely that more than a few circles at a time can be stored before compression, so statistical tests where performed regarding each data stream for a given pointing, as a collection of 60 independent realizations of the same process. Of course this is only approximately true. The 1/f noise correlates subsequent scan circles, but since its r.m.s. amplitude per sample is typically about one-tenth of the white noise r.m.s. or less, these correlations can be neglected in this analysis.
Starting from the folded data streams a given normality test was applied to each set of 60 realizations for each one of the 8640 samples, transforming the stream of samples in a stream of test results for the given test. The cumulative distribution of frequency was then computed over the 8640 test results. Since 60 samples does not represent a large statistics, significant deviations from theorethically evaluated confidence levels are expected resulting in an excessive rejection or acceptation rates. For this reason each test was calibrated applying it to the undigitized white noise data stream. Moreover, in order to analyze how the normality evolves increasing the signal complexity, tests was repeated increasing the information content of the generated data stream.
To simplify the discussion we considered as a reference test the usual Kolmogorow - Smirnov D test from Press et al. (1986) and we fix a $`95\%`$ acceptance level. The test was “calibrated” using the MonteCarlo white noise generator of our mission simulator in order to fix the threshold level $`D_{\mathrm{th}}`$ as the $`D`$ value for which more than $`95\%`$ of our samples show $`DD_{\mathrm{th}}`$. From Table 3 the quantization effect is evident, at twice the nominal quantization step ($`\mathrm{VOT}=2`$ V/K) in $`30\%`$ of the samples (i.e. 2592 samples) the distribution of realizations deviates from a normal distribution ($`D>D_{\mathrm{th}}`$). Since the theoretical compression rate from eq. (9) is for a continuous distribution of levels ($`\sigma \mathrm{\Delta }`$) a smaller $`C_\mathrm{r}`$ should is expected. Since the deviation from the normal distribution is a systematic effect, for the sake of cosmological data analysis one may tune the D test to take account of the quantization. As an example, the third line in Tab. 3 reports the threshold for the quantized signal $`D_{\mathrm{th}}^\mathrm{Q}`$ for which $`95\%`$ of the quantized white noise samples are accepted as normal distributed. The line below represents the success rate for the full quantized signal. After the recalibration the test is able to recognize that in $`95\%`$ of the cases the signal is drawn from a normal distribution, but at the cost of a growth in the threshold $`D`$ which now is a function of the quantization step $`\mathrm{\Delta }`$.
As for the entropy distribution and the binary statistics, even in this case most of the differences between the results obtained for a pure white noise signal and the full signals are explained by the presence of the cosmological dipole. However these simulations are not accurate enough to draw any quantitative conclusions about the distortion in the sampling statistics induced by digitization, but they suggest that to approximate the instrumental signal as a quantized white noise plus a cosinusoidal term associated to the cosmic dipole is more than adequate in order to understand the optimal loss-less compression rate achievable in the case of the Planck/LFI mission.
## 6 Quantization and Quantization Error
A possible solution to solve the bandwidth problem is to reduce the amount of information of the sampled signal i.e. its entropy. Independently from the way in which this is performed, the final compression strategy will be lossy, and the final reconstructed (uncompressed) signal will be corrupted with respect to the original one, degrading in some regard the experimental performances. In this regard, any sort of lossy compression may be seen as a kind of signal rebinning with a coarser resolution (quantization step) in $`\mathrm{\Delta }T/T`$.
There are at least six aspects in Planck/LFI operations which may be affected by a coarser quantization:
1. $`C_l`$ and periodical signals reconstruction;
2. destriping;
3. foreground separation;
4. point like sources detection;
5. variable sources characterization;
6. tests for normality of CMB fluctuations.
Since the non linear nature of the quantization process, all of them are hard to be analytically evaluated and for this reason a specific simulation task is in progress for the Planck/LFI collaboration (White & Seiffert (1999), Maris et al. (2000)). However an heuristic evaluation for the point (1) by analytical means is feasible.
Quantization operates a convolution of the normal distribution of the input signal with the quantization operator $`(x:\mathrm{\Delta })=`$sign$`(x)\mathrm{\Delta }floor(|x/\mathrm{\Delta }|)`$. If the quantization error: $`(x(x:\mathrm{\Delta }))`$ is uniformly distributed its expectation is $`\mathrm{\Delta }/2`$ and its variance is $`\mathrm{\Delta }/\sqrt{12}`$ (Kollár (1994)). Quantization over a large amount of samples may be regarded as an extra source of noise which will enhance the variance per sample. If the quantization error is statistically independent from the input quantized signal and if it may be added in quadrature to the white noise variance $`\sigma _{WN}`$, the total variance per sample will be $`\sigma _{WN}^2\left(1+\frac{\mathrm{\Delta }^2/\sigma _{WN}^2}{12}\right)`$. So for $`\mathrm{\Delta }\stackrel{_<}{_{}}\sigma _{WN}`$ the expected quantization r.m.s. is $`\stackrel{_<}{_{}}4\%`$. From error propagation the relative error on the $`C_l`$ is (Maino (1999)):
$$\frac{\delta C_l}{C_l}=\sqrt{\frac{4\pi }{A}}\sqrt{\frac{2}{(2l+1)f_{sky}}}\left[1+\frac{\sigma ^2\theta ^2}{B_l^2C_l}\right]$$
(11)
so that the quantization contribution to the overall error will be small and dominated by the cosmic variance for a large set of $`l`$. However the application of such encouraging result must be considered carefully in a true experimental framework. Apart from the assumptions, it has to be demonstrated indeed that a large quantization error like this will not harm significantly the aforementioned aspects, moreover the impact of signal quantization will depend on how and in which point of the detection chain it will be performed.
## 7 Experimental Evaluation of Off-The-Shelf Compressors
This section describes the evaluation protocol and the experimental results of the compression of simulated data streams for Planck/LFI.
### 7.1 Evaluation Protocol
First tests were performed on a HP-UX workstation on four compressors (Maris et al. (1998)) but given the limited number of off-the-shelf compression codes for such platform, we migrated the compression pipeline on a Pentium III based Windows/NT workstation.
As described in section 2 the signal composition is defined by many components, both astrophysical and instrumental in origin. In particular, it is important to understand how each component or instrumental parameter, introducing deviations to the pure white noise statistics, affects the final compression rate.
To scan systematically all the relevant combinations of signal compositions and off-the-shelf compressors, a Compression Pipeline was created. The pipeline is based on five main components: the signal quantization pipeline, the signal database, the compression pipeline, the compression data base, the post-processing pipeline. The signal quantization pipeline performs the operations described in the upper part of figure 1. The simulated astrophysical signals are hold in a dedicated section of the signal archive, they are processed by the quantization pipeline and stored back in a reserved section of the signal archive. So quantized data streams are generated for each relevant combination of the quantization parameters, signal composition and sky pointing.
Each compressor is then applied by the compression pipeline to the full set of quantized signals in the signal archive. Results, in terms of compression efficiency as a function of quantization parameters are stored in the compression database. The statistical analysis of section 5 are performed with a similar pipeline.
Finally the post-processing pipeline scans the compression data base in order to produce plots, statistics, tables and synthetic fits. Its results are usually stored into one of the two databases.
The pipeline is managed by PERL 5.004 script files which drive FORTRAN, C, IDL programs or on-the-shelf utilities gluing and coordinating their activities. Up to $`\mathrm{75\hspace{0.17em}000}`$ lines of repetitive code are required per simulation run. They are generated by a specifically designed Automated Program Generator (APG) written in IDL (Maris & Staniszkis (1998)). The APG takes as an input a table which specifies: the set of compressors to be tested, the set of quantization parameters to be used, the order in which to perform the scan of each parameter/compressor, the list of supporting programs to be used, other servicing parameters. The program linearizes the resulting parameter space and generates the PERL simulation code or, alternatively, performs other operations such as: to scan the results data base to produce statistics, plots, tables, and so on. The advantage of this method is that a large amount of repetitive code, may be quickly produced, maintained or replaced with a minor effort each time a new object (compressor, parameter or analysis method) is added to the system.
### 7.2 Experimental Results
Purpose of these compression tests is to give an upper limit to the lossless compression efficiency for LFI data and to look for an optimal compressor to be proposed to the LFI consortium.
A decision about the final compression scheme for Planck/LFI has not been taken yet and only future studies will be able to decide if the best performing one will be compatible with on-board operations (constrained by: packet independence and DPU capabilities) and will be accepted by the Planck/LFI collaboration.
For this reason up to now only off-the-shelf compressors and hardware where considered. To test any reasonable compression scheme a wide selection of lossless compression algorithms, covering all the known methods, was applied to our simulated data. Lacking a comprehensive criteria to fix a final compressor, as memory and CPU constrains, we report in a compact form the results related to all the tested compressors. We are confident that in the near future long duration flight balloon experiments as on-board electronics prototypes will provide us with a more solid base to test and improve the final compression algorithms looking at real data.
Tables 4, 5 list the selected compression programs. Since the behaviour (and efficiency) of each compressor is determined by a set of parameters one or more macro file operating a given combination of compressor code plus parameters is defined. It has to be noted that uses is a space qualified algorithm, based on Rice compression method, for which space qualified dedicated hardware already exists.
To evaluate the performances of each compressor, figures of merit are drawn like the one in figure 5 which shows the results for the best performing compressor: arith-n1. Looking at such figures it is possible to note as the compression efficiency does not depend much on the signal composition. This is true even when large, impulsive signals, as planets, affecting few samples over thousands are introduced. Again, this is a consequence of the fact that white noise dominates the signal, being the most important component to affect the compression efficiency. In this regard it has been speculated that the 1/f component should improve the correlation between neighborhood samples affecting the compression efficiency (Maris et al. (1998)) no relevant effect may be detected into our simulations. As an example from figure 5 for the 30 GHz signal the addition of the 1/f noise to the white noise data stream affects the final $`C_\mathrm{r}`$ for less than $`0.5\%`$.
The only noticeable (i.e. some $`6\%`$) effect due to an increase in the signal complexity, occurs when the cosmic dipole is added. In the present signal the dipole amplitude is comparable with the white noise amplitude ($`3`$ mK) so its effect is to distort the sample distribution, making it leptocurtic. As a consequence compressors, which usually work best for a normal distributed signal, becomes less effective. Since the dipole introduces correlations over one full scan circle, i.e. some $`10^3`$ samples, while compressors establish the proper coding table observing the data stream for a small set of consecutive samples (from some tens to some hundred samples), even a self adaptive compressor will likely loose the correlation introduced by the dipole. A proper solution to this problem is suggested in section 9. The other signal components do not introduce any noticeable systematic effect. The small differences shown by the figures of merit may be due to the compression variance and depend strongly on the compressor of choice. As an example a given compressor may be more effective to compress the simulated data stream with the full signal than the associated simpler data stream containing only white noise, 1/f noise, CMB and dipole. At the same time another compressor may show an opposite behaviour.
As shown by Figure 6, and as expected from eq. (9) increasing $`\mathrm{VOT}`$, i.e. increasing the quantization step, increases the compression rate. In addition $`C_\mathrm{r}`$ increases increasing $`N_\mathrm{c}`$ up to an $`20\%`$. The increase is noticeable for $`N_\mathrm{c}<15`$ and saturates after $`N_\mathrm{c}=30`$. On the contrary its dependence on the offset ($`\mathrm{AFO}`$) is negligible (less than $`1\%`$). For these reasons in the subsequent analysis the AFO dependency is neglected and the corresponding simulations are averaged.
### 7.3 Synthetical Description
The full data base of simulated compression results takes about 14 MBytes, for practical purposes it is possible to synthesize all this information using a phenomenological relation which connects $`C_\mathrm{r}`$ with $`N_\mathrm{c}`$ and $`\mathrm{VOT}`$ whose free parameters may be fitted using the data obtained from the simulations. In short:
$$C_\mathrm{r}^{\mathrm{Fit}}(\mathrm{VOT},N_\mathrm{c})=\frac{C_{\mathrm{r},1}}{(N_\mathrm{c})+𝒮(N_\mathrm{c})\mathrm{ln}\left[\frac{\mathrm{VOT}}{1\mathrm{V}/\mathrm{K}}\right]}$$
(12)
where $`C_{\mathrm{r},1}`$ is the $`C_\mathrm{r}`$ for $`N_\mathrm{c}=1`$, $`\mathrm{VOT}=1.0`$ V/K, while $`(N_\mathrm{c})`$ and $`𝒮(N_\mathrm{c})`$ describe the $`C_\mathrm{r}`$ dependence on $`N_\mathrm{c}`$. In particular the relation is calibrated for any compressor imposing that $`C_\mathrm{r}(\mathrm{VOT}=1\mathrm{V}/\mathrm{K},N_\mathrm{c}=1)=C_{\mathrm{r},1}`$.
The linear dependency of $`1/C_\mathrm{r}^{\mathrm{Fit}}`$ over $`\mathrm{ln}\mathrm{VOT}`$ is a direct consequence of equation (9), and is confirmed by a set of tests performed over the full set of our numerical results for the compression efficiency, the r.m.s. residual between the best fit (12) and simulated data being less than $`1.5\%`$, in almost the $`92\%`$ of the cases and less than $`1\%`$ in $`72\%`$ of the cases. The dependencies of its parameters $``$ and $`\mathrm{S}`$ over $`N_\mathrm{c}`$ are obtained by a test-and-error method performed on our data set and we did not investigate further on their nature. For all practical purposes our analysis shows that these functions are well approximated by a series expansion:
$$(N_\mathrm{c})\mathrm{exp}\left(\underset{k=1}{\overset{2}{}}A_k(\mathrm{ln}N_\mathrm{c})^k\right),$$
(13)
$$𝒮(N_\mathrm{c})𝒮_1\mathrm{exp}\left(\underset{k=1}{\overset{5}{}}B_k(\mathrm{ln}N_\mathrm{c})^k\right).$$
(14)
here $`𝒮_1`$, $`A_k`$ and $`B_k`$ are free parameters obtained by fitting the simulated data, in particular $`𝒮_1`$ is the slope for $`N_\mathrm{c}=1`$.
Since an accuracy of some percent in determining the free parameters of $`C_\mathrm{r}^{\mathrm{Fit}}(\mathrm{VOT},N_\mathrm{c})`$ is enough, the fitting procedure was simplified as follow. For a given compressor, signal component, swap status, and $`N_\mathrm{c}`$ value $``$ and $`𝒮`$ where determined by a $`\chi ^2`$ fitting procedure. The list of $``$ and $`𝒮`$ as a function of $`N_\mathrm{c}`$ have been fitted by using relations (13) and (14) respectively. The fitting algorithm tests different degrees of the polynomial in the aforementioned relations (up to 2 for $`(N_\mathrm{c})`$, up to 5 for $`𝒮(N_\mathrm{c})`$) stopping when the maximum deviation of the fitted relation respect to the data is smaller than $`0.5\%`$ for $``$ or 0.0001 for $`𝒮`$, or when the maximum degree is reached.
Tables 6, 7, 8, 9 report the results of the compression exercise ordered for decreasing $`C_{\mathrm{r},1}`$. The first column is the name of compression macro (i.e. a given compression program with a defined selection of modificators and switches) as listed in tables: 4, 5. The third and fourth column are the fitted $`C_{\mathrm{r},1}`$and $`𝒮_1`$ as defined in: (12), (13), (14). From the $`5^{\mathrm{th}}`$ to the $`7^{\mathrm{th}}`$ columns and from the $`8^{\mathrm{th}}`$ to the $`13^{\mathrm{th}}`$ columns the polynomial degree and the expansion parameters for (13) and (14) are reported.
Many compressors are sensitive to the ordering of the Least and Most Significant Bytes of a 16 bits word in the computer memory and files. Two ordering conventions are assumed: UnSwapped i.e. Least Significant Byte is stored First or Swapped i.e. Most Significant Byte is stored First. As in Digital VAX/VMS Operating System, Microsoft Windows/NT operating system convention is Most Significant Byte first. For this reason each test was repeated twice, one time with the original data stream file with swapped bytes and the other after unswapping bytes. If the gain in $`C_{\mathrm{r},1}`$ after unswapping is bigger than some percent, unswapped compression is reported, otherwise the swapped one is reported. These two cases are distinct by the second column of tables 6, 7, 8, 9 which is marked with a y if unswapping is applied before compressing. It is interesting to note that not only 16 bits compressors, such as uses, are sensitive to swapping. Also many 8 bits compressors are sensitive to it, maybe that this is due to the fact that if the most probable 8 bits symbol is presented first at the compressor a slightly better balanced coding table is built.
It should be noted that the coefficients reported here are obtained compressing one or more full scan circles at a time, so their use to extrapolate $`C_\mathrm{r}`$ when each scan circle is divided in small chunks which are separately compressed has to be performed carefully, especially for $`\mathrm{VOT}0.5`$ V/K where some extrapolated $`C_\mathrm{r}`$ grows instead of to decrease for a decreasing $`N_\mathrm{c}`$ as in most of the cases. However we did not investigate further the problem because the time required to perform all the tests over all the compressors increases decreasing $`N_\mathrm{c}`$, and because up to now a final decision about the packet length has not been made yet. Moreover, short data chunks introduce other constrains which are not accounted for by eq. (9) but which are discussed in section 8.
Apart from the choice of the best compressor, Tables 6, 7, 8, 9 allows interesting comparisons.
The performances of the arithmetic compression arith are very sensitive to changes in the coding order $`n=0`$, $`\mathrm{}`$, 7. The computational weight grows with $`n`$, while $`C_\mathrm{r}`$ is minimal at $`n=0`$, maximal for $`n=1`$ and decreases increasing $`n`$ further.
Both non-Adaptive Huffman (huff-c) and Adaptive Huffman (ahuff-c) are in the list of the worst compressors, considering both the pure white noise signal and the full signal.
We implemented the space-qualified uses compressor with a wide selection of combinations of its control parameters: the number of coding bits, the number of samples per block, the possibility to search for correlations between neighborhood samples. We report the tests for 16 bits coding only, changing the other parameters. Uses is very sensitive to byte unswapping, when not performed uses does not compress at all. On the other hand, opposite to arith the sensitivity of the final $`C_\mathrm{r}`$ to the various control parameters is small or negligible. In most cases $`C_{\mathrm{r},1}`$ differs of less than $`0.01`$ for changing the combination of control parameters, such changes are not displayed by the two digits approximation in the tables, but they are accounted for by the sorting procedure which fixes the table ordering. At 30 GHz most of the tested compressors cluster around $`C_{\mathrm{r},1}=2.67`$ and at this level arith-n3 is as good as uses. At 100 GHz the best uses macros clusters around $`C_{\mathrm{r},1}=2.43`$ \- $`2.44`$, equivalent to arith-n2 performances. In our tests uses performs worst at 8 samples per block without correlation search, but apart from it, in our case the correlation search does not improve significantly the compression performances. Some commercial programs such as boa, bzip compress better than uses.
## 8 Further Constrains: Packet Independence and Packet Length
As an example of global constrains to the on-board compression we discuss the problems related to Packets Independence and Packets Length.
Data from the LFI must be packetized before being sent to Earth. Packets independence is considered to be a requirement, then each packet must be self-consistent, its loss or its erroneous transmission must not interfere with the data retrieval from subsequent packets. More over each packet must carry in “clear” format (i.e. uncompressed) all the information needed to decode its content. That is: each packet must contain its own decoding table or decoding information. A typical packet length is about some hundred of bytes, but smaller length may be planned if required; at the same time a typical decoding table holds something less than a hundred bytes leaving limited room for data.
In addition, for a fixed length $`L_\mathrm{u}`$ of a random input stream (expressed in bits) the output $`L_\mathrm{c}`$ will not be a constant but will change in time with respect to the averaged length $`L_\mathrm{u}/C_\mathrm{r}`$. Of course, it is not possible to predict in advance what will be the final length of a given bit stream. So either $`L_\mathrm{u}`$ is held fixed, loosing in compression efficiency, or $`L_\mathrm{u}`$ is adapted with some interactive method, maximizing the compression efficiency but at the cost of a significant slowing of the compression process.
In conclusion, the packets independence plus limited packet length prevents from sending the decoding table, leaving only two possibilities open: i) send the relevant bytes only (Maris (1999a)), ii) to use a predefined coding table (Maris (1999b)), both methods are described in the next section.
## 9 Proposed Coding and Compression Scheme
The basic principle of the first method named Least Significant Bits Packing (LSBP) is to send only those bits of the 16 bits output from the ADC which are affected by the signal and the noise. This is effective for the nominal mission since with the planned quantization step of 0.3 mK/adu, at one sigma the noise will fill about 21 levels, this will require at least 5 bits over 16 and it is reasonable to expect a final data flow equivalent to $`C_{\mathrm{r},1}<3`$. It is not possible to improve much the compression rate by compressing the resulting 5 bits data stream, since its entropy would be $`H<5.4`$ bits and $`C_\mathrm{r}\stackrel{_<}{_{}}1.08`$.
In order to ensure the compression to be lossless all the samples exceeding the \[$`\sigma `$, $`+\sigma `$\] (5 bits) range have to be sent separately coding at the same time: their position (address) in the stream vector and their value. So, for $`N_{\mathrm{bits}}<16`$ bits corresponding to a threshold $`x_{\mathrm{th}}=2^{N_{\mathrm{bits}}}`$, each group of samples stored into a packet is partitioned into two classes accordingly with their value $`x`$:
| Regular Samples (RS) | $``$def$``$ | all those samples for which: | $`|x|x_{\mathrm{th}}`$, |
| --- | --- | --- | --- |
| Spike Samples (SS) | $``$def$``$ | all those samples for which: | $`|x|x_{\mathrm{th}}.`$ |
The coding process then consists of two main steps: i) to split the data stream in Regular and Spike Samples preserving the original ordering in the stream of Regular Samples, ii) to store (send) the first $`N_{\mathrm{bits}}`$ bits of the regular samples and, in a separated area, the 16 bits values and the location in the original data stream of each Spike Sample, i.e. Spike Samples will require more space to be stored than regular ones. The decoding process will be the reverse of this packing process.
In this scheme each packet will be divided into two main areas: the Regular Samples Area (RSA) which hold the stream of Regular Samples, the Spike Sample Area (SSA) which hold the stream of Spike Samples, plus a number of fields which will contain packing parameters such as: the number of samples, the number of regular samples, the offset, etc. Since the number of samples in each area will change randomly it will be not possible to completely fill a packet. The filling process will leave an empty area in the packet in average smaller than $`N_{\mathrm{bits}}`$.
In Maris (1999a) a first evaluation for the 30 GHz channel is given assuming that the signal is composed only of white noise plus the CMB dipole. As noticed in section 7.2 the cosmological dipole affects the compression efficiency reducing it of a small amount. To deal with it a possible solution would be to subdivide each data stream in packets, subtract to each measure of a given packet the integer average of samples (computed as a 16 bits integer number) and then compress the residuals. Each integer average will be sent to Earth together with the related packet where the operation will be reversed. Since all the numbers are coded as 16 bits integers all the operations are fully reversible and no round off error occurs. However it cannot be excluded that the computational cost of such operation will compensate the gain in $`C_\mathrm{r}`$.
Two schemes are proposed to perform the cosmological dipole self-adaptement. In Scheme A the average of samples in the packet are subtracted before coding and then sent separately. In Scheme B $`x_{\mathrm{th}}`$ is varied proportionally to the dipole contribution. Both of them assumes that the dipole contribution is about a constant over a packet length. From this assumption: $`L_p\stackrel{_<}{_{}}200`$ samples i.e. $`L_p<512`$ bytes, since for $`L_p>512`$ bytes the cosmic dipole contribution can not be considered as a time constant. For larger packets a better modeling (i.e. more parameters) will be required in order not to degrade the compression efficiency.
A critical point is to fix the best $`x_{\mathrm{th}}`$, i.e. $`N_{\mathrm{bits}}`$, for a given signal statistics, coding scheme and packet length $`L_p`$. Even here $`C_\mathrm{r}`$ grows with the packet length but it does not change monotously with $`x_{\mathrm{th}}`$. An increase in $`x_{\mathrm{th}}`$ ($`N_{\mathrm{bits}}`$) decreases the number of spike samples, but increases the size of each regular sample. While the opposite occurs when $`x_{\mathrm{th}}`$ is decreased, and when $`N_{\mathrm{bits}}<4`$ bits $`C_\mathrm{r}<1`$. For both the schemes the optimality is reached for $`N_{\mathrm{bits}}=6`$ bits, but Scheme A is better than B, with: $`C_\mathrm{r}(\text{Scheme A}`$, $`L_p=512\text{ bytes})=2.61`$, $`C_\mathrm{r}(\text{Scheme B}`$, $`L_p=512\text{ bytes})=2.29`$.
Compared with arith-n1, this compression rate is smaller of about a 14 - 30%. This is due to two reasons: i) coding by a threshold cut is less effective than to apply an optimized compressor; ii) the results reported in tables 6, 7, 8, 9 refer to the compression of a full circle of data instead of a small packet, resulting in a higher efficiency. However, the efficiency of this coding method is similar to the efficiency of the bulk of the other true loss-less compressors tested up to now, and when the need to send a decoding table is considered, is even higher.
The second possible solution to the packeting problem is to use one or more standardized coding tables for the compression scheme of choice (Maris (1999b)). In this case the coding table would be loaded into the on-board computer before launch or time by time in flight and the table should be known in advance at Earth. Major advantages would be: 1. the coding table has not to be sent to Earth; 2. the compression operator will be reduced to a mapping operator which may be implement as a tabular search, driven by the input 8 or 16 bits word to be compressed; 3. any compression scheme (Huffman, arithmetic, etc.) may be implemented replacing the coding table without changes to the compression program; 4. the compression procedure may be easily written in C or the native assembler language for the on-board computer or, alternatively, a simple, dedicated hardware may be implemented and interfaced to the on-board computer. The disadvantages of this scheme are: 1. each table must reside permanently in the central computer memory unless a dedicated hardware is interfaced to it; 2. it is difficult to use adaptive schemes in order to tune the compressor to the input signal, as a consequence the $`C_\mathrm{r}`$ may be somewhat smaller than in the case of a true self-adapting compressor code.
The first problem may be circumvented limiting the length of the words to be compressed. In our case the data streams may be divided in chunks of 8 bits and the typical table size would be $`\stackrel{_<}{_{}}1`$ Kbyte. Precomputed coding tables may be accurately optimized by Monte-Carlo simulations on ground or using signals from ground tests of true hardware.
The second problem may be overcome by using a preconditioning stage, reducing the statistics of the input signal to the statistics for which the pre-calculated table is optimized. In addition more tables may reside in the computer memory and selected looking to the signal statistics. With a simple reversible statistical preconditioner, about ten tables per frequency channel would be stored in the computer memory, so that the total memory occupation would be less than about 40 Kbytes. It cannot be excluded that the two methods just outlined cannot be merged.
## 10 Estimation of the Overall Compression Rate
The overall compression rate (efficiency) is the average of $`C_\mathrm{r}`$ ($`\eta _\mathrm{c}`$) over the full set of detectors. Appendix A illustrates the mathematical aspects of such average. From (A4):
$$\overline{C_\mathrm{r}}(N_\mathrm{c})=\left[\underset{\nu }{}\frac{f_\nu }{C_{\mathrm{r}}^{}{}_{,\nu }{}^{}(N_\mathrm{c})}\right]^1.$$
(15)
We will limit ourselves to the most probable case $`N_\mathrm{c}=1`$ and to the most effective compressor arith-n1. The compression parameters $`C_{\mathrm{r},1}`$ and $`𝒮_1`$ at 30 GHz and 100 GHz are derived from our simulations, while $`C_{\mathrm{r},1}`$ and $`𝒮_1`$ at 44 GHz and 70 GHz are obtained by linear interpolation of the simulated values as a function of $`\mathrm{ln}\sigma _\nu `$. After that we obtain:
$$\overline{C_\mathrm{r}}\frac{2.66}{1+0.271\times \mathrm{ln}\mathrm{VOT}}.$$
(16)
As expected the overall compression rate is dominated by the 100 GHz channel. Taking in account the conservative $`\mathrm{VOT}`$ distribution considered in equation (A8) the overall compression rate becomes: $`\overline{C_\mathrm{r}}2.63`$ which represents a $`2\%`$ correction only. It is likely that this correction will be even smaller, since the amplifiers gain will be adjusted in order to cover a smaller $`\mathrm{VOT}`$ interval. So this $`2\%`$ correction represents our greatest uncertainty in our estimation of the expected compression rate, and we may conservatively conclude that:
$$\overline{C_\mathrm{r}}_{,\mathrm{arith}\mathrm{n1}}2.65\pm 0.02$$
(17)
Recently a new evaluation of the expected instrumental sensitivity leads to some change in the expected white noise r.m.s.. These changes affect in particular the 30 GHz channel, but does not change significantly the 100 GHz channel so that the overall compression rate will be practically unaffected.
## 11 Conclusions
The expected data rate from the Planck Low Frequency Instrument is $`260`$ kbits/sec. The bandwidth for the scientific data download currently allocated is just $`60`$ kbit/sec. Assuming an equal subdivision of the bandwidth between the two instruments on-board Planck, an overall compression rate of a factor 8.7 is required to download all the data.
In this work we perform a full analysis on realistically simulated data streams for the 30 GHz and 100 GHz channels in order to fix the maximum compression rate achievable by loss-less compression methods, without considering explicitly other constrains such as: the power of the on-board Data Processing Unit, or the requirements about packet length limits and independence, but taking in account all the instrumental features relevant to data acquisition, i.e.: the quantization process, the temperature / voltage conversion, number of quantization bits and signal composition.
As a complement to the experimental analysis we perform in parallel a theoretical analysis of the maximum compression rate. Such analysis is based on the statistical properties of the simulated signal and is able to explain quantitatively most of the experimental results.
Our conclusions about the statistical analysis of the quantized signal are: I) the nominally quantized signal has an entropy $`h5.5`$ bits at 30GHz and $`h5.9`$ bits at 100GHz, which allows a theoretical upper limit for the compression rate $`2.9`$ at 30 Ghz and $`2.7`$ at 100 GHz. II) Quantization may introduce some distortion in the signal statistics but the subject requires a deepest analysis.
Our conclusions about the compression rate are summarized as follows: I) the compression rate $`C_\mathrm{r}`$ is affected by the quantization step, since greater is the quantization step higher is $`C_\mathrm{r}`$ (but worse is the measure accuracy). II) $`C_\mathrm{r}`$ is affected also by the stream length $`L_u`$, i.e. more circles are compressed better then few circles. III) the dependencies on the quantization step and $`L_u`$ for each compressor may be summarized by the empirical formula (12). A reduced compression rate $`C_{\mathrm{r},1}`$ is correspondingly defined. IV) the $`C_\mathrm{r}`$ is affected by the signal composition, in particular, by the white noise r.m.s. and by the dipole contribution, the former being the dominant parameter and the latter influencing $`C_\mathrm{r}`$ for less than $`6\%`$. The inclusion of the dipole contribution reduces the overall compression rate. The other components (1/f noise, CMB fluctuations, the galaxy, extragalactic sources) have little or no effect on $`C_\mathrm{r}`$. In conclusion, for the sake of compression rate estimation, the signal may be safely represented by a sinusoidal signal plus white noise. V) since the noise r.m.s. increases with the frequency, the compression rate $`C_\mathrm{r}`$ decreases with the frequency, for the LFI $`\mathrm{\Delta }C_\mathrm{r}/C_\mathrm{r}\stackrel{_<}{_{}}10\%`$. VI) the expected random r.m.s. in the overall compression rate is less than $`1\%`$. VII) we tested a large number of off-the-shelf compressors, with many combinations of control parameters so to cover every conceivable compression method. The best performing compressor is the arithmetic compression scheme of order 1: arith-n1, the final $`C_{\mathrm{r},1}`$ being 2.83 at 30 GHz and 2.61 at 100 GHz. This is significantly less than the bare theoretical compression rate (9) but when the quantization process is taken properly into account in the theoretical analysis, this discrepancy is largely reduced. VIII) taking into account the data flow distribution among different compressors the overall compression rate for arith-n1 is:
$$\overline{C_\mathrm{r}}_{,\mathrm{arith}\mathrm{n1}}2.65\pm 0.02$$
This result is due to the nature of the signal which is noise dominated and clearly excludes the possibility to reach the required data flow reduction through loss-less compression only.
Possible solutions deal with the application of lossy compression methods such as: on-board averaging, data rebinning, or averaging of signals from duplicated detectors, in order to reach an overall lossy compression of about a factor 3.4, which coupled with the overall loss-less compression rate of about 2.65 should allow to reach the required final compression rate $`8.7`$. However each of these solutions will introduce heavy constraints and important reduction of performances in the final mission design, so that careful and deep studies will be required in order to choose the best one.
Another solution to the bandwidth problem would be to apply a coarser quantization step. This has however the drawback of reducing the signal resolution in terms of $`\mathrm{\Delta }T/T`$.
Lastly the choice of a given compressor cannot be based only on its efficiency obtained from simulated data, but also on the on-board available CPU and on the official ESA space qualification: tests with this hardware platform and other compressors will be made during the project development. Moreover, in the near future long duration flight balloon experiments and ground experiments (see Lasenby et al. (1998); De Bernardis & Masi (1998)) will provide a solid base to test and improve compression algorithms. In addition the final compression scheme will have to cope with requirements about packet length and packet independence. We discuss briefly this problems recalling two proposals (Maris (1999b), Maris (1999a)) which suggest solutions to cope with these constrains.
## Appendix A Appendix: Formulation of the Final Data Flow
In this appendix we will discuss how to account for the distribution of the acquisition parameters between the different detectors in the computation of the overall compression rate. Since the formalism is simpler we will develop expressions for $`\eta _\mathrm{c}=1/C_\mathrm{r}`$ instead of $`C_\mathrm{r}`$.
We have pointed out in 5.2 that the compression efficiency is a random variable, whose distribution is a function of all those parameters which are relevant to fix the statistical distribution of the input signal. In our case: $`\nu `$, $`\mathrm{VOT}`$, $`\mathrm{AFO}`$, $`N_{\mathrm{circ}}`$ are the relevant parameters, so that the conditioned probability to have a compression efficiency in the range $`\eta _\mathrm{c}`$, $`\eta _\mathrm{c}+d\eta _\mathrm{c}`$ is:
$$𝒫_{\nu ,N_{\mathrm{Circ}}}\left(\eta _\mathrm{c}|\mathrm{AFO},\mathrm{VOT}\right)d\eta _\mathrm{c}.$$
(A1)
This probability may be obtained by our MonteCarlo simulations for different combinations of $`\mathrm{AFO}`$, $`\mathrm{VOT}`$, $`N_{\mathrm{circ}}`$ and $`\nu `$. Then the averaged compression efficiency is:
$$\overline{\eta }_{\mathrm{c}\nu ,N_{\mathrm{Circ}}}(\mathrm{AFO},\mathrm{VOT})=_0^+\mathrm{}𝑑\eta _\mathrm{c}\eta _\mathrm{c}𝒫_{\nu ,N_{\mathrm{Circ}}}\left(\eta _\mathrm{c}|\mathrm{AFO},\mathrm{VOT}\right).$$
(A2)
Of course we assumed that for any $`\nu `$, $`\mathrm{VOT}`$, $`\mathrm{AFO}`$, $`N_{\mathrm{circ}}`$ the probability distribution is integrable and normalized to 1, while the integration limits $`0`$, $`+\mathrm{}`$ are to be intended as formal. There are several detectors for any frequency channel, each one having its own $`\mathrm{AFO}`$ and $`\mathrm{VOT}`$, so distributions of $`\mathrm{AFO}`$ and $`\mathrm{VOT}`$ values may be guessed among the different detectors. Assuming they are integrable and normalized to 1 as well it is possible to compute the most probable $`\overline{\eta }_{\mathrm{c}}^{}{}_{\nu ,N_{\mathrm{Circ}}}{}^{}`$ as <sup>5</sup><sup>5</sup>5Here
$$_{\mathrm{AFO}_{\mathrm{min}}}^{\mathrm{AFO}_{\mathrm{max}}}𝑑\mathrm{AFO}𝒫_\nu (\mathrm{AFO})=1,_{\mathrm{VOT}_{\mathrm{min}}}^{\mathrm{VOT}_{\mathrm{max}}}𝑑\mathrm{VOT}𝒫_\nu (\mathrm{VOT})=1$$
. ;
$$\overline{\eta }_{\mathrm{c}}^{}{}_{\nu ,N_{\mathrm{Circ}}}{}^{}=_{\mathrm{AFO}_{\mathrm{min}}}^{\mathrm{AFO}_{\mathrm{max}}}𝑑\mathrm{AFO}𝒫_\nu (\mathrm{AFO})_{\mathrm{VOT}_{\mathrm{min}}}^{\mathrm{VOT}_{\mathrm{max}}}𝑑\mathrm{VOT}𝒫_\nu (\mathrm{VOT})\overline{\eta }_{\mathrm{c}}^{}{}_{\nu ,N_{\mathrm{Circ}}}{}^{}(\mathrm{AFO},\mathrm{VOT}).$$
(A3)
With this definition the final overall compression efficiency is:
$$\overline{\eta }_{\mathrm{c}}^{}{}_{N_{\mathrm{Circ}}}{}^{}=\underset{\nu =30,44,70,100\mathrm{G}\mathrm{H}\mathrm{z}}{}f_\nu \overline{\eta }_{\mathrm{c}}^{}{}_{\nu ,N_{\mathrm{Circ}}}{}^{}$$
(A4)
where $`f_\nu `$ is the partition function for the data flow through the different detectors, if $`n_{\mathrm{dtc},\nu }`$ is the number of detectors for the frequency channel $`\nu `$ (see Tab. I), $`n_{\mathrm{dtc}}=_{\nu =30,44,70,100\mathrm{G}\mathrm{H}\mathrm{z}}n_{\mathrm{dtc},\nu }=112`$, is the total number of detectors and if the number of samples for frequency is a constant, then:
$$f_\nu =\frac{n_{\mathrm{dtc},\nu }}{112},$$
(A5)
so that for $`\nu =30`$, $`44`$, $`70`$ and 100 GHz respectively: $`f_\nu =0.0714`$, $`0.1071`$, $`0.2143`$ and $`0.6071`$, finally the expect data rate for each set of 60 circles is:
$$\overline{R}_{N_{\mathrm{Circ}}}=16\mathrm{bits}\times \mathrm{\hspace{0.33em}60}\mathrm{circles}\times \mathrm{\hspace{0.33em}8640}\mathrm{samples}\times \mathrm{\hspace{0.33em}112}\mathrm{detectors}\times \overline{\eta }_\mathrm{c}^{N_{\mathrm{Circ}}}.$$
(A6)
Presently there are no data to know in advance the distribution of $`\mathrm{VOT}`$ and $`\mathrm{AFO}`$ values between the different detectors. For this reason in this work we assumed simply flat distributions, identical for each frequency for such parameters. More over, the $`\mathrm{AFO}`$ contribution is negligible, so that the variance introduced by this parameter is neglected. From (9) we assumed that the compression efficiency is approximately a linear function of $`\mathrm{ln}\mathrm{VOT}`$ or:
$$\overline{\eta }_{\mathrm{c}}^{}{}_{\nu ,N_{\mathrm{Circ}}}{}^{}(\mathrm{VOT})\overline{\eta }_{\mathrm{c}}^{}{}_{\nu ,N_{\mathrm{Circ}},1}{}^{}+\dot{\overline{\eta }}_{\mathrm{c}}^{}{}_{\nu ,N_{\mathrm{Circ}}}{}^{}\mathrm{ln}\mathrm{VOT}$$
(A7)
where $`\dot{\overline{\eta }}_{\mathrm{c}}^{}{}_{\nu ,N_{\mathrm{Circ}}}{}^{}`$ is the first derivative of $`\overline{\eta }_{\mathrm{c}}^{}{}_{\nu ,N_{\mathrm{Circ}}}{}^{}(\mathrm{VOT})`$ with respect to $`\mathrm{ln}\mathrm{VOT}`$ computed for $`\mathrm{VOT}=1`$ V/K, $`\overline{\eta }_{\mathrm{c}}^{}{}_{\nu ,N_{\mathrm{Circ}},1}{}^{}\overline{\eta }_{\mathrm{c}}^{}{}_{\nu ,N_{\mathrm{Circ}}}{}^{}(\mathrm{VOT}=1\mathrm{V}/\mathrm{K})`$. As an example, at 30 GHz for arith-n1 the full signal compression rate is $`\overline{\eta }_{\mathrm{c}}^{}{}_{\nu ,N_{\mathrm{Circ}}}{}^{}(\mathrm{VOT})0.3534+0.287\times \mathrm{ln}\mathrm{VOT}(\mathrm{K}/\mathrm{V})`$ with one interpolation error less than $`0.2\%`$. With these approximations eq. (A3) becomes
$$\overline{\eta }_{\mathrm{c}}^{}{}_{\nu ,N_{\mathrm{Circ}}}{}^{}\overline{\eta }_{\mathrm{c}}^{}{}_{\nu ,N_{\mathrm{Circ}},1}{}^{}+\dot{\overline{\eta }}_{\mathrm{c}}^{}{}_{\nu ,N_{\mathrm{Circ}}}{}^{}_{0.5\mathrm{V}/\mathrm{K}}^{1.5\mathrm{V}/\mathrm{K}}𝑑\mathrm{VOT}\frac{\mathrm{ln}\mathrm{VOT}}{1.0\mathrm{V}/\mathrm{K}}$$
(A8)
and after integration we obtain the final formula
$$\overline{\eta }_{\mathrm{c}}^{}{}_{\nu ,N_{\mathrm{Circ}}}{}^{}\overline{\eta }_{\mathrm{c}}^{}{}_{\nu ,N_{\mathrm{Circ}},1}{}^{}0.045229\dot{\overline{\eta }}_{\mathrm{c}}^{}{}_{\nu ,N_{\mathrm{Circ}}}{}^{}$$
(A9)
for the case in the previous example: $`\overline{\eta }_{\mathrm{c}}^{}{}_{\nu ,N_{\mathrm{Circ}}=2}{}^{}0.3404`$ which is equivalent to a compression efficiency $`2.94`$.
To understand the influence of the error in the $`\mathrm{VOT}`$ determination over the distribution on the final predictions the computation is made for a truncated (i.e. zero outside the $`\mathrm{VOT}`$ range of interest) normal distribution of $`\mathrm{VOT}`$. The r.m.s. for the $`\mathrm{VOT}`$ distribution is chosen in the $`\mathrm{VOT}`$ range \[0.5, 1.5\] V/K we obtain respectively $`\overline{\eta }_{\mathrm{c}}^{}{}_{\nu ,N_{\mathrm{Circ}}}{}^{}0.3494`$, $`0.3439`$, $`0.3420`$; which corresponds to compression efficiencies: 2.86, 2.91, 2.92 respectively. Similar results are obtained with a quadratic $`\mathrm{VOT}`$ distribution. In conclusion these predictions are robust against the shape of the $`\mathrm{VOT}`$ distribution, at least for distributions which are symmetric around the nominal $`\mathrm{VOT}=1`$ V/K value.
We warmly acknowledge a number of people which actively support this work with fruitful discussions, in particular F. Argüeso, M. Bersanelli, L. Danese, G. De Zotti, E. Gaztñaga, J. Herrera, N. Mandolesi, P. Platania, A. Romeo, M. Seiffert and L. Toffolatti and K. Gorski and all people involved in the construction of the Healpix pixelisation tools, largely employed in this work, and Dr. G. Lombardi from Siemens Bochun - Germany and Dr. G. Maris from ETNOTEAM - Milano for fruitful discussions about compression principles and their practical application.
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