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warning/0005/cond-mat0005299.html	ar5iv	text	# Wigner crystallization in the two electron quantum dot ## I Introduction During the last ten years quantum dots have attracted a lot of interest both experimentally and theoretically . Recently, much attention was paid to the investigation of quantum dots in strong perpendicular magnetic fields where a rich structure of cusps and steps in the chemical potential $`\mu (N,B)`$ (as a function of confined electron number $`N`$ and applied magnetic field strength $`B`$) was observed . The above structure is related to changes in the ground state electron density in the presence of a magnetic field, and it initiated numerous theoretical investigations after Chamon and Wen proposed the quantum dot edge reconstruction. The electron density in a quantum dot is the result of the interplay of the repulsive character of the electron interaction and the attractive forces caused by the confining potential, the magnetic field and the electron exchange interaction. If the magnetic field is strong enough, the overlap between the electron wave functions becomes less, the electron interaction will dominate, and a ring of electrons at the dot edge is formed. It was also found that if the magnetic field is increased further the above electron ring becomes unstable, and a ground state with a broken rotational symmetry appears. The possibility of the appearance of spin waves and charge density waves have also been reported. A quantum dot with a density profile consisting of rings with lumps reminds us to the Wigner crystal which is the ground state of the classical electron system in a dot in the absence of a magnetic field. In the latter case the Wigner crystal occurs when the potential energy (the inter-electron interaction and the confinement potential) dominates over the electron kinetic energy. This is just the classical limit in a quantum problem. Therefore, classical or quasi-classical methods should be adequate for the description of the Wigner crystal. But from a first sight it appears that such a classical limit is not reached when a strong magnetic field is applied. The above mentioned electron density reconstruction was revealed when the magnetic field is so strong that the electrons occupy the lowest Landau level. Thus, the electron motion quantization is essential, and consequently, the kinetic energy exceeds the potential energy due to the Coulomb interaction. Nevertheless, due to the degeneracy of the Landau level this large electron kinetic energy is actually frozen out, and Wigner crystallization results from the same potential energy as in the classical case without magnetic field. Actually the above crystallization is a result of different energy scales in the electron system under consideration. The purpose of the present paper is, by using an exactly solvable model of two electrons in a dot, to illustrate the conditions under which a Wigner crystal can be formed in the case of a strong magnetic field ($`B`$), and to compare this quantum Wigner crystal with the classical zero magnetic field one. The paper is organized as follows. In Sect. II we present our model. Sect. III gives an introduction to the quasi-classical approach for the $`B=0`$ case. Sect. IV discusses the general $`B0`$ case. Our conclusions are presented in Sect. V. ## II Model We consider two electrons with effective masses $`m^{}`$ moving in the $`z=0`$ plane and which are confined by a two-dimensional harmonic potential of characteristic frequency $`\omega _0`$. A magnetic field is applied in the $`z`$ direction and described by the vector potential chosen in the symmetric gauge $`𝐀=[𝐁\times 𝐫]/2`$. The corresponding Hamiltonian can be separated into two parts $$H=H_R+H_r,$$ (1) (see, for instance ) which represents the center-of-mass and relative motion (with corresponding coordinates $`𝐑=(𝐫_1+𝐫_2)/2`$ and $`𝐫=𝐫_1𝐫_2`$). In dimensionless form those parts can be written as follows: $`H_R`$ $`=`$ $`{\displaystyle \frac{1}{4}}_R^2+\left\{1+{\displaystyle \frac{B^2}{4}}\right\}R^2{\displaystyle \frac{iB}{2}}[𝐑\times _R]_z,`$ (3) $`H_r`$ $`=`$ $`_r+{\displaystyle \frac{1}{4}}\left\{1+{\displaystyle \frac{B^2}{4}}\right\}r^2{\displaystyle \frac{iB}{2}}[𝐫\times _r]_z+{\displaystyle \frac{\lambda }{r}}.`$ (4) The system energy is measured in $`\mathrm{}\omega _0`$ units, and the coordinates in $`a_0=\sqrt{\mathrm{}/m^{}\omega _0}`$ units. The symbol $`\lambda =a_0/a_B`$ stands for the electron interaction coupling constant which is the ratio of the characteristic dot dimension $`a_0`$ and the Bohr radius $`a_B=ϵ\mathrm{}^2/m^{}e^2`$. The magnetic field strength $`B`$ is measured in $`\mathrm{\Phi }_0/\pi a_0^2`$ units where $`\mathrm{\Phi }_0=\pi \mathrm{}c/e`$ is the magnetic flux quantum. We will concentrate ourselves to the study of the coordinate wave function part of the ground singlet state $$\mathrm{\Psi }(𝐫_1,𝐫_2)=\mathrm{\Phi }(𝐑)\psi (𝐫),$$ (5) the corresponding electron density $`\rho (𝐫)`$ $`=`$ $`{\displaystyle d^2r_1d^2r_2\|\mathrm{\Psi }(𝐫_1,𝐫_2)\|^2\widehat{\rho }(𝐫)}`$ (6) $`=`$ $`{\displaystyle d^2r_1d^2r_2\|\mathrm{\Psi }(𝐫_1,𝐫_2)\|^2\underset{n=1}{\overset{2}{}}\delta (𝐫𝐫_n)}`$ (7) $`=`$ $`2{\displaystyle d^2r_1\|\mathrm{\Phi }(𝐫+𝐫_1/2)\|^2\|\psi (𝐫_1)\|^2}`$ (8) and the correlation function $`K(𝐫,𝐫^{})`$ $`=`$ $`{\displaystyle d^2r_1d^2r_2\|\mathrm{\Psi }(𝐫_1,𝐫_2)\|^2}`$ (10) $`\times \left\{\widehat{\rho }(𝐫)\widehat{\rho }(𝐫^{})\delta (𝐫𝐫^{})\widehat{\rho }(𝐫)\right\}`$ $`=`$ $`2\|\mathrm{\Phi }(\{𝐫+𝐫^{}\}/2)\|^2\|\psi (𝐫𝐫^{})\|^2.`$ (11) It is the latter function which enables one to determine whether, or not the system is in the Wigner crystal state. ## III Zero magnetic field case We consider first the simplest case in which no magnetic field is applied. Then the center-of-mass motion part is trivial. It just represents the two dimensional harmonic oscillator motion which has no relation to the Wigner crystallization problem. Its ground state energy is $`E_R=1`$ with the corresponding wave function $$\mathrm{\Phi }(𝐑)=\sqrt{\frac{2}{\pi }}\mathrm{exp}(R^2).$$ (12) Due to the cylinder symmetry the relative motion wave function part can be written as $$\psi (𝐫)=\frac{1}{\sqrt{2\pi }}\mathrm{exp}(im\phi )R(r),$$ (13) where the radial function $`R(r)`$ is obtained by solving the one dimensional eigenvalue problem as determined by the Hamiltonian $`H_r={\displaystyle \frac{1}{r}}{\displaystyle \frac{d}{dr}}r{\displaystyle \frac{d}{dr}}+V(r,\lambda ),`$ (15) $`V(r,\lambda )={\displaystyle \frac{m^2}{r^2}}+{\displaystyle \frac{1}{4}}r^2+{\displaystyle \frac{\lambda }{r}}.`$ (16) As was already mentioned strong electron correlation (and the Wigner crystal as well) occurs when the potential energy dominates over the electron kinetic energy, i. e. when $`\lambda \mathrm{}`$. In this interesting limit the eigenvalue problem is strongly simplified and can be solved by analytical means. Indeed, in the case of $`\lambda \mathrm{}`$ the potential (16) has a minimum close to the point $`r_0=(2\lambda )^{1/3}`$ (see solid curve in Fig. 1). The potential can be expanded into a $`(rr_0)`$ series (dashed curve in Fig. 1) $$V(r,\lambda )\frac{3}{4}(2\lambda )^{2/3}+\frac{3}{4}(rr_0)^2+\frac{m^2}{(2\lambda )^{2/3}},$$ (17) which actually coincides with the quasi-classical expansion into a $`\lambda ^{2/3}`$ series . The form of the above expansion clearly indicates the energy scales of the different excitations in the quantum dot. The first term ($`\lambda ^{2/3}`$) is just the classical dot energy which can also be obtained using the hydrodynamic approximation . In the $`\lambda \mathrm{}`$ limit this energy dominates. Thus, we are in the region of the classical Wigner crystal. The solution of the Schrödinger equation with the harmonic term in the expansion of the potential (17) leads to $`\lambda `$-independent equidistant electron ring vibration excitations (shown in Fig. 1 by the horizontal dashed lines). Note that the separation of those vibration levels are of order $`1`$, and consequently, much less than the classical potential energy. The last term in the expansion (17) describes the rotation energy. It leads to a small (of order $`\lambda ^{2/3}`$) splitting of the rotation levels as is shown by the encircled part in Fig. 1 which is enlarged. The spectrum is thus similar to the one of molecules with bands of rotation levels, attached to each vibration level. For the expanded potential (17) we obtain the following quasi-classical radial ground state wave function, $$P(r)=N\mathrm{exp}\left\{a(rr_0)^2\right\},$$ (18) where $`a=\sqrt{3}/4`$ and the normalization $`N=(2a/\pi )^{1/4}r_0^{1/2}`$. Inserting the above relative motion wave function (18) together with the center-of-mass motion function (12) into expressions (6,10) we obtain the electron density $`\rho (r,\phi )={\displaystyle \frac{2N^2}{\pi ^2}}{\displaystyle _0^{\mathrm{}}}𝑑r^{}r^{}\mathrm{exp}(2r^2r_{}^{}{}_{}{}^{2}/2)`$ (20) $`\times \mathrm{exp}\left\{2a(r^{}r_0)^2\right\}{\displaystyle _0^{2\pi }}𝑑\phi ^{}\mathrm{exp}\left\{2rr^{}\mathrm{cos}(\phi \phi ^{})\right\}`$ $`=`$ $`{\displaystyle \frac{4N^2}{\pi }}{\displaystyle _0^{\mathrm{}}}𝑑r^{}r^{}\mathrm{exp}\left\{2r^2r_{}^{}{}_{}{}^{2}/22a(r^{}r_0)^2\right\}`$ (22) $`\times I_0(2rr^{}){\displaystyle \frac{N^2}{\pi }}\sqrt{{\displaystyle \frac{2\gamma }{a}}}\mathrm{exp}\{\gamma (rr_0/2)^2\},`$ and the correlation function $`K(𝐫,𝐫^{})`$ $`=`$ $`{\displaystyle \frac{4N^2}{\pi }}\mathrm{exp}\left\{(𝐫+𝐫^{})^2/2\right\}`$ (24) $`\times \mathrm{exp}\left\{2a(\|𝐫𝐫^{}\|r_0)^2\right\}.`$ In expression (20), $`I_0`$ is the Bessel function and $`\gamma =8a/(1+4a)1.268`$. A schematic picture of the above functions is shown in Fig. 2 by the shadowed regions. In Fig. 2(a) the electron density is depicted which is mainly concentrated on a thin ring. Notice that it has no lumps, and exhibits the cylindrical symmetry of the system Hamiltonian. Wigner crystal state can be seen in the correlation function which is plotted in Fig. 2(b). The latter is actually the same as the conditional probability distribution. One electron is fixed in position $`𝐫^{}`$ (indicated by a cross in Fig. 2(b)) and the density of the other electron is then mainly concentrated in the opposite position which clearly indicates the presence of Wigner crystallization in this quasi-classical limit. ## IV Strong magnetic field case Due to the symmetric gauge which preserves the cylindrical symmetry, both (i. e. center-of-mass and relative) wave functions can be written as a product of the orbital exponent and the radial part as in expression (13). The corresponding radial Hamiltonians (for center-of-mass and relative motions) can be presented as follows: $`H_R`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \frac{d}{dR}}R{\displaystyle \frac{d}{dR}}+{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{M}{R}}BR\right)^2+R^2,`$ (26) $`H_r`$ $`=`$ $`{\displaystyle \frac{1}{r}}{\displaystyle \frac{d}{dr}}r{\displaystyle \frac{d}{dr}}+\left({\displaystyle \frac{m}{r}}{\displaystyle \frac{Br}{4}}\right)^2+{\displaystyle \frac{1}{4}}r^2+{\displaystyle \frac{\lambda }{r}}.`$ (27) Because we are interested in the asymptotic limit $`B\mathrm{}`$, it is convenient to scale the variables and the Hamiltonian as follows $$𝐫𝐫B^{1/2},𝐑𝐑B^{1/2},HHB,$$ (28) in order to have the expansion parameter $`B^1`$ explicitly in our problem. Inserting (28) into expression (IV) we arrive at the following Hamiltonian: $`H_R`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \frac{d}{dR}}R{\displaystyle \frac{d}{dR}}+{\displaystyle \frac{1}{4}}\left({\displaystyle \frac{M}{R}}R\right)^2+V_R(R,B),`$ (30) $`H_r`$ $`=`$ $`{\displaystyle \frac{1}{r}}{\displaystyle \frac{d}{dr}}r{\displaystyle \frac{d}{dr}}+\left({\displaystyle \frac{m}{r}}{\displaystyle \frac{r}{4}}\right)^2+V_r(r,B),`$ (31) where in the asymptotic region $`B\mathrm{}`$ the terms $`V_R(R,B)`$ $`=`$ $`{\displaystyle \frac{R^2}{B^2}},`$ (33) $`V_r(r,B)`$ $`=`$ $`{\displaystyle \frac{r^2}{4B^2}}+{\displaystyle \frac{\lambda }{r\sqrt{B}}}`$ (34) can be treated as small perturbations. The solution of the Schrödinger equations for the zero order Hamiltonians (without the terms $`V_R`$ and $`V_r`$) leads to the degenerate Landau levels where the ground state has energy $`E_R=E_r=1/2`$ labeled by integer positive momentum ($`M`$ and $`m`$) values. Mathematically this degeneracy is a consequence of the equivalence of the zero center-of-mass Hamiltonian potential $`(M/RR)^2/4`$ which is shown in Fig. 3(a) where the orbital momentum is indicated by the corresponding integers. For any momentum $`M`$ the potential curve has a minimum equal to zero at the position $`R_{\mathrm{min}}=\sqrt{M}`$. A similar potential is obtained for the zero order relative motion equation, and will therefore not be discussed. Next we take the perturbation terms into account. For the case of center-of-mass motion the potential is shown in Fig. 3(b). The potential is composed of the same curves as in the case without perturbation, but they are now moved up by the amount $`R^2/B^2`$ shown in Fig. 3(b) by the dotted curve. The lowest minimum is obtained for $`M=0`$ which implies that the center-of-mass motion wave function is the same as in the case without magnetic field (12). This is not the case for the relative motion where all potential curves are shifted up by the amount $`V_r(r,B)`$ as shown in Fig. 3(c) by the dotted curve. According to Eq. (34) this curve increases for both $`r0`$ and $`r\mathrm{}`$ and reaches a minimum value at $$r_{\mathrm{min}}=(2\lambda )^{1/3}\sqrt{B}.$$ (35) Minimizing the potentials with respect to the relative angular momentum for $`r=r_{\mathrm{min}}`$ we obtain the corresponding orbital momentum $$m=\left(\frac{\lambda }{4}\right)^{2/3}B.$$ (36) Consequently, the relative motion wave function is peaked on the ring of radius $`r_0`$. When the magnetic field strength increases the radius $`r_0`$ tends to infinity, while the orbital number is growing as well. Expanding the potential with the lowest minimum in the vicinity of the equilibrium point (35) we find $$\left(\frac{m}{r}\frac{r}{4}\right)^2\frac{(rr_{\mathrm{min}})^2}{4},$$ (37) which means that the thickness of the ring remains constant. The layout of the energy spectrum in the large magnetic field limit is shown in Fig. 4. Notice that the spectrum has two different energy scales as in the $`B=0`$ case. But the physical meaning of those scales is quite different. The largest energy scale is the electron kinetic energy. And the rotation levels are split due to the interplay of the Coulomb interaction between the electrons and the confinement potential. Going back to our original units before the scaling (28) we find that the ring radius is $`r_{\mathrm{min}}=(2\lambda )^{1/3}`$, which is identical to the zero magnetic field case. But now the thickness of the ring is $`B^{1/2}`$, and it tends to zero as the magnetic field strength approaches infinity. Surprisingly, we arrived at the same situation as was found for the case without magnetic field. The relative motion wave function is located on a ring whose diameter greatly exceeds its thickness. Thus the conclusions of previous section concerning Wigner crystallization in the quantum dot are also valid in the strong magnetic field case. While for $`B=0`$ the Wigner crystal state is reached for $`\lambda 1`$ we find that a sufficiently strong magnetic field can crystallize the system for any $`\lambda 0`$. ## V Conclusions The occurrence of Wigner crystallization depends on the ratio of the distance between the electrons ($`l`$) and the characteristic dimension of the single electron wave packet ($`a`$). Actually this ratio is a measure of the electron density in the quantum dot. In our case $`l`$ is given by the radius $`r_{\mathrm{min}}`$ of the ring in the correlation function, and $`a`$ coincides with its thickness. This ratio is $$\chi =\frac{l}{a}=\frac{r_{\mathrm{min}}}{a}=\frac{\lambda ^{1/3}}{B^{1/2}}=\lambda ^{1/3}\sqrt{B},$$ (38) which in real units reads $$\chi =\left(\frac{a_0}{a_B}\right)^{1/3}\left(\frac{a_0}{l_B}\right).$$ (39) Here $`l_B=\sqrt{\mathrm{}c/eB}`$ is the magnetic length. The larger this ratio the more pronounced the Wigner crystal state is. Thus, both a strong magnetic field and strong interaction favors Wigner crystallization, but in a different way. The electron interaction makes the system less dense by enlarging the inter-particle distance. The magnetic field also makes the system effectively less dense but by compressing the single electron wave packages. Also we would like to note that mathematically in both cases the static Wigner crystal properties can be calculated by the same method, namely by the minimization of the classical potential energy which is composed of the electron interaction and the confinement potential. Nevertheless the physical meaning of that calculation is quite different. In the strong electron interaction case the potential energy dominates making the whole problem quasi-classical, while in the case of strong magnetic fields, the problem is essentially quantum mechanical (the Landau level energy is dominating). But due to the degeneracy of the problem the system is guided by the same potential energy as for the $`B=0`$ case but with different energy scales. Therefore, although the Wigner crystal configuration is the same, one can expect different dynamics in the two limiting cases. ## Acknowledgments This work is supported by the Flemish Science Foundation, IUAP-VI and the ”Bijzonder Onderzoeksfonds van de Universiteit Antwerpen”. One of us (F. M. P.) is a research director with FWO-Vl. We acknowledge discussions with B. Partoens, S. Reimann and D. Pfannkuche.
warning/0005/gr-qc0005030.html	ar5iv	text	# Axiomatic approach to radiation reaction of scalar point particles in curved spacetime ## I Introduction There has been much recent interest in calculating the motion of astrophysical systems which emit gravitational waves in anticipation of data from a new generation of detectors. Full three-dimensional numerical simulations are required in order to produce useful results for many of the most promising observational candidates, such as colliding black holes. However, there also exists a large class of systems which can be accurately modelled by a small isolated body moving in the fixed background created by a much larger body (e.g., a solar mass star falling into a supermassive black hole). For such a system, we might hope to produce useful results by treating the smaller object as a point particle and introducing the effects of its fields and internal structure as perturbations to the background geodesic orbit. The perturbations due to the particle’s own fields, commonly called “radiation reaction” or “self force” effects, are particularly important because they include the forces responsible for the decay of the body’s orbit. If both the background spacetime and the unperturbed orbit of the body possess enough symmetry, it is possible to infer the effects of these forces on the orbit from global conservation principles: one calculates the energy and/or angular momentum radiated to infinity by a particle in geodesic motion, and then modifies the orbit to reflect this energy and angular momentum loss in a time-averaged fashion. (Obviously, this procedure can be iterated if greater accuracy is required.) Some justification for this method is provided by Quinn and Wald . However, in the absence of such symmetries, it is necessary to directly calculate the effects of the local fields in the neighborhood of the particle. Unfortunately, this problem is ill-posed, since the fields diverge in the neighborhood of the particle’s world line, so that any such local calculation must include a rule for extracting the appropriate finite part of these divergent fields. There is an extensive literature devoted to this regularization problem. In 1938, Dirac reproduced the force expression (originally given by Abraham ) for a point particle coupled to an electromagnetic field in Minkowski spacetime by imposing local energy conservation on a tube surrounding the particle’s world line and subtracting the infinite contributions to the force through a “mass renormalization” scheme. In 1960, Dewitt and Brehme generalized this approach to an arbitrary curved background spacetime. (A trivial calculational error in their paper was later corrected by Hobbs .) More recently, Mino et al. further adapted this approach to produce a force expression for a point particle coupled to a linearized gravitational field on a vacuum background spacetime, and Quinn and Wald rederived both the electromagnetic and gravitational forces using an axiomatic approach which, in effect, regularizes the forces by comparing forces in different spacetimes. There has emerged from this work a consensus regarding the correct equation of motion for a particle coupled to electromagnetic fields on an arbitrary curved background and for a particle coupled to linearized gravitational fields on a vacuum background. In principle, the latter equation allows one to calculate the dynamics of the astrophysical systems of interest described above. In practice, however, very little progress has been made in applying either equation of motion to concrete physical examples for two reasons. First, given a world line in an arbitrary spacetime, the calculation of the associated retarded fields is a complex and difficult problem. Second, once these fields are calculated, the equations of motion require one to identify that portion of the retarded field at each point of the world line which arises from source contributions interior to the light cone. This part of the field is often called the “tail term,” and most approximation schemes for calculating the retarded field entangle the tail and non-tail contributions to the field. Nevertheless, some progress has been made, notably in the electromagnetic case. In 1964, DeWitt and DeWitt calculated the tail term for an electromagnetic particle in a circular orbit on a Schwarzschild background to leading order in the background curvature and the velocity of the particle. In 1980, Smith and Will calculated the force on an electromagnetic particle held static on a Schwarzschild background, essentially by repeating DeWitt and Brehme’s local stress-energy conservation argument. Neither result has been generalized to the case of a massive particle coupled to gravitational fields, nor has there been any direct progress on the more complex systems which are of interest to the gravitational wave astronomy community. However, several new ideas have emerged in recent years which may lead to further progress. Ori has suggested an alternative regularization scheme involving averaging of multipole moments which is better adapted to concrete calculations, while others have suggested a hybrid scheme in which the tail term is calculated through a combination of Hadamard expansion techniques for small distances and multipole techniques for larger distances . It is clearly important to test these ideas. In particular, we must know whether these schemes are equivalent to the equations of motion discussed above. Because of its mathematical simplicity, one natural system in which to explore all of these questions is that of a point particle coupled to a scalar field. Motivated by this, several researchers have begun to apply the ideas discussed above to the scalar system. In particular, Ori’s method has been applied to the motion of scalar particles in the Kerr spacetime and in the Schwarzschild spacetime , and Wiseman has adapted the calculation of Smith and Will in order to calculate the force on a scalar particle held static in the Schwarzschild spacetime. In the present paper, we generalize the axiomatic approach of Quinn and Wald in order to produce the general equation of motion for a point particle coupled to a scalar field on an arbitrary background spacetime. It is hoped that this general expression will be useful in evaluating the validity of the calculational schemes described above for the scalar case, and that this comparison will ultimately help to clarify the relationship between the various methods which have been proposed for the electromagnetic and gravitational cases. In Sec. II, we derive an expression for the force on a particle following an arbitrary trajectory in curved spacetime. Then, in Sec. III, we explore the equations of motion which follow from this expression in the absence of non-scalar forces. ## II The scalar force Given a spacetime containing a particle world line and a Klein-Gordon field sourced by the particle, we wish to define the total scalar force $`f_\mathrm{S}^a`$ on the particle at each point of the world line, including so-called self-force or radiation reaction effects. For an electromagnetic point particle in flat spacetime, an expression of this sort was first given by Abraham in 1905, was later rederived in a relativistic context by Dirac , and is often found in textbooks (e.g., Jackson ). However, since there are no classical point particles in nature, and the theoretical status of such objects is problematic at best, it is important to ask how any such prescription is constrained by physics. Our view is that the force law should reflect the force on an extended body coupled to a scalar field in the limit of small spatial extent. In particular, fix the background spacetime and consider a family of extended bodies and corresponding scalar fields parameterized by $`ϵ`$, the spatial size of the bodies. For each body in the family, we define a center of mass world line $`z(\tau )`$ (e.g., by the methods of Beiglböck ) and calculate the charge, $`q`$, and mass, $`m`$, of the body with respect to this world line,<sup></sup><sup></sup>*Because the scalar charge density is a scalar quantity, the total charge that one calculates for an extended body depends upon the spacelike surface used to slice the body. This is in contrast to electromagnetism, where the charge density is the time component of a conserved vector field and the total charge is independent of slice. as well as the force $`f_\mathrm{S}^a[ϵ]`$ exerted by the scalar field on the body. (For the definition of the force exerted on a small body by a field to which it is coupled, see Quinn and Wald .) We further require that $`m`$ and $`q`$ vanish as $`ϵ`$ goes to zero. For such a one-parameter family, Quinn and Wald argue that it is possible to specify some set of conditions on the internal structure and composition of the extended bodies such that the limit of $`f_\mathrm{S}^a[ϵ]`$ for small $`ϵ`$ is independent of their internal details. We would like our expression for $`f_\mathrm{S}^a`$ to correctly the describe the order $`q`$ and $`q^2`$ contributions to $`f_\mathrm{S}^a[ϵ]`$ which are independent of the internal structure of the body under these conditions. (Other corrections which arise from the internal structure, such as multipole effects and spin effects, have been derived elsewhere and should simply contribute additively at this order.) Unfortunately, the limit described above is quite delicate, and the task of specifying conditions to ensure its convergence appears to be formidable. (The analysis of Dixon demonstrates the degree of complexity which arises even without considering self-field effects.) Nevertheless, certain properties of this limit are strongly suggested by the nature of the divergences in the scalar field. Following Quinn and Wald , we will introduce these properties as axioms, and then give the unique prescription for $`f_\mathrm{S}^a`$ which satisfies these axioms. In the next subsection, we will motivate our crucial Comparison Axiom by considering the point particle limit described above and develop the expansions required to state the axiom. Then, in the following subsection, we state both axioms and give the unique prescription for $`f_\mathrm{S}^a`$ that satisfies them, which is the main result of this paper. ### A Motivation for the Comparison Axiom Consider a spacetime $`(M,g_{ab})`$ containing a spatially compact body characterized by stress-energy $`T_{\mathrm{body}}^{ab}`$ and scalar charge density $`\rho `$, a smooth Klein-Gordon field $`\varphi `$, and possibly some other set of fields which are coupled to the body, characterized by $`T_{\mathrm{ext}}^{ab}`$. The Klein-Gordon field $`\varphi `$ satisfies the equation $$^a_a\varphi =4\pi \rho $$ (1) with stress-energy $$T_\mathrm{S}^{ab}=\frac{1}{4\pi }(^a\varphi ^b\varphi \frac{1}{2}g^{ab}g_{cd}^c\varphi ^d\varphi ).$$ (2) Assuming that the total stress-energy is conserved, so that $$_b(T_{\mathrm{body}}^{ab}+T_\mathrm{S}^{ab}+T_{\mathrm{ext}}^{ab})=0,$$ (3) then the force density exerted on the body by the scalar field is given by $$_bT_{\mathrm{body}}^{ab}+_bT_{\mathrm{ext}}^{ab}=_bT_\mathrm{S}^{ab}=\rho ^a\varphi .$$ (4) Therefore, naively taking the point particle limit, we would expect the force on a scalar particle of charge $`q`$ to be given by $$f_\mathrm{S}^a=q^a\varphi .$$ (5) Unfortunately, this expression is meaningless as it stands, since $`_a\varphi `$ diverges on the world line of the particle. (The situation is exactly the same with the Lorentz force law $`f_{\mathrm{EM}}^a=qF^{ab}u_b`$.) However, if we consider two points $`P`$ and $`\stackrel{~}{P}`$ along the world lines of two different particles in two different spacetimes (each with charge $`q`$), and we identify the neighborhoods of $`P`$ and $`\stackrel{~}{P}`$, then we might hope that, under some conditions, the difference $`^a\varphi \stackrel{~}{}^a\stackrel{~}{\varphi }`$ will be finite even as the two individual fields diverge. Under such conditions, it seems reasonable to expect that the difference between the forces on the particles will be given by the (finite) difference between the field gradients. That is, $$f_\mathrm{S}^a\stackrel{~}{f}_\mathrm{S}^a=\underset{r0}{lim}q^a\varphi \stackrel{~}{}^a\stackrel{~}{\varphi }_r.$$ (6) (Here, the average over a sphere of radius $`r`$, denoted by $`_r`$, is introduced to allow for the possibility that the $`r0`$ limit of the difference is finite, but direction-dependent.) Quinn and Wald give plausibility arguments which suggest that the counterpart of Eq. (6) is indeed a property of the point particle limit in the electromagnetic and gravitational cases. These arguments generalize straightforwardly to the scalar case, so we will not give the details here. Instead, we will simply impose Eq. (6) as an axiom and investigate the consequences for $`f_\mathrm{S}^a`$. This idea will be the basis of our crucial Comparison Axiom in the next subsection. However, first we must find out what conditions to impose on the spacetimes, the world lines near $`P`$ and $`\stackrel{~}{P}`$, and the identification of their neighborhoods in order to ensure that the difference in the field gradients be finite as $`r0`$. In order to answer this question, we will now examine in detail the singularity structure of the scalar field in the neighborhood of the world line. Consider a scalar field satisfying Eq. (1) in a spacetime $`(M,g_{ab})`$ with a point particle source $$\rho (x)=q\delta ^4(x,z(\tau ))𝑑\tau .$$ (7) In contrast to the electromagnetic case, the Klein-Gordon equation does not require conservation of charge. For simplicity, we shall assume throughout our analysis that the charge $`q`$ is constant along the world line. We wish to expand $`\varphi `$ in $`r`$, the spatial distance from the world line $`z(\tau )`$. We are primarily interested in the divergent contributions to $`\varphi `$, characterized by the negative powers of $`r`$ in the expansion, since these divergent contributions will determine the conditions required for convergence of the limit in Eq. (6). It follows from the general theory of propagation of singularities (see theorem 26.1.1 of Hormander ) that every solution of Eq. (1) which is smooth away from the world line will have the same singularity structure near the world line, so we are free to choose any convenient solution for our expansion. Later, when we wish to produce an explicit expression for $`f_\mathrm{S}^a`$, we will want to write $`\varphi `$ in terms of the advanced and retarded solutions. Therefore, these are the solutions which we will analyze in the following expansion. Given any point $`x`$ in a spacetime $`(M,g_{ab})`$, there exists a convex normal neighborhood $`C(x)`$ containing $`x`$ \[i.e. a neighborhood $`C(x)`$ such that there exists a unique geodesic connecting any two points within $`C(x)`$\]. For $`x^{}C(x)`$, the Hadamard elementary solution of Eq. (1) can be written in the form $$G^{(1)}(x,x^{})=\frac{1}{\pi }\left[\frac{U(x,x^{})}{\sigma (x,x^{})}+V(x,x^{})\mathrm{ln}\|\sigma (x,x^{})\|+W(x,x^{})\right],$$ (8) with corresponding advanced (+) and retarded (-) Green’s functions $$G_\pm (x,x^{})=\theta _\pm (x,x^{})\left[U(x,x^{})\delta \left(\sigma (x,x^{})\right)V(x,x^{})\theta \left(\sigma (x,x^{})\right)\right].$$ (9) Here, $`\sigma (x,x^{})`$ is the biscalar of squared geodesic distanceThe biscalar of squared geodesic distance $`\sigma (x,x^{})`$ is equal to half of the squared length of the geodesic connecting $`x`$ and $`x^{}`$: negative for timelike separated events, positive for spacelike separated events, and zero for null separated events. It is defined only when there is a unique geodesic connecting $`x`$ and $`x^{}`$. and $`U`$, $`V`$, and $`W`$ are all smooth biscalar fields. (For an explanation of the bitensor formalism, see Dewitt and Brehme .) The scalar function $`\theta _\pm (x,x^{})`$ is unity when $`x^{}`$ is in the causal future/past of $`x`$ and vanishes otherwise. For $`x`$ near the world line $`z(\tau )`$, let $`\tau _\mathrm{\Sigma }`$ be the proper time along the world line which is simultaneous with $`x`$ in the sense that the spatial surface $`\mathrm{\Sigma }`$ generated by geodesics perpendicular to $`u^a`$ at $`z(\tau _\mathrm{\Sigma })`$ intersects $`x`$. In particular, let $`x`$ lie a proper distance $`r`$ along the geodesic generated by unit spatial vector $`\widehat{r}^a`$ at $`z(\tau _\mathrm{\Sigma })`$, and let $`z(\tau _+)`$ and $`z(\tau _{})`$ be the intersection of the world line with the future and past light cones of $`x`$, respectively. We require that $`x`$ be close enough to the world line that $`z(\tau _\mathrm{\Sigma })`$, $`z(\tau _+)`$, and $`z(\tau _{})`$ all lie within the neighborhood $`C(x)`$, and we denote the future and past intersections of the world line with the boundary of $`C(x)`$ by $`z(T_+)`$ and $`z(T_{})`$, respectively. This is illustrated in Fig. 1. For the retarded field $`\varphi _{}`$, we then have $`\varphi _{}(x)`$ $`=`$ $`{\displaystyle G_{}(x,x^{})\rho (x^{})\sqrt{g}d^4x^{}}`$ (10) $`=`$ $`{\displaystyle G_{}(x,x^{})\left(q\delta ^4(x^{},z(\tau ))𝑑\tau \right)\sqrt{g}d^4x^{}}`$ (11) $`=`$ $`q{\displaystyle G_{}(x,z(\tau ))𝑑\tau }`$ (12) $`=`$ $`q{\displaystyle _T_{}^{T_+}}\theta _{}[x,z(\tau )]\left[U(x,z(\tau ))\delta \left(\sigma (x,z(\tau ))\right)V(x,z(\tau ))\theta \left(\sigma (x,z(\tau ))\right)\right]𝑑\tau `$ (14) $`+q{\displaystyle _{\mathrm{}}^T_{}}G_{}(x,z(\tau ))𝑑\tau `$ $`=`$ $`q{\displaystyle _T_{}^{\tau _\mathrm{\Sigma }}}\left[U\delta (\sigma )V\theta (\sigma )\right]𝑑\tau +q{\displaystyle _{\mathrm{}}^T_{}}G_{}𝑑\tau `$ (15) In the last line and hereafter, we suppress the spacetime dependence for all biscalars, since each depends upon $`x`$ in its first argument and $`z(\tau )`$ in its second argument. For a bitensor $`A`$, we introduce the notation $$\dot{A}\frac{d}{d\tau }A(x,z(\tau ))=u^a^{}_a^{}A(x,z(\tau )).$$ (16) We have $$d\tau =\frac{d\tau }{d\sigma }d\sigma =\left(\frac{d\sigma }{d\tau }\right)^1d\sigma =\dot{\sigma }^1d\sigma ,$$ (17) which gives us $$\varphi _{}=q\left[\left\{\dot{\sigma }^1U\right\}_{\tau =\tau _{}}_T_{}^\tau _{}V𝑑\tau \right]+q_{\mathrm{}}^T_{}G_{}𝑑\tau .$$ (18) We now wish to produce the corresponding expression for $`_a\varphi _{}`$. Note that the right side of Eq. (18) depends upon $`x`$ in two ways: explicitly through the first argument of each biscalar and implicitly through $`\tau _{}`$. We have $`_a\varphi _{}`$ $`=`$ $`q_a\left[\left\{\dot{\sigma }^1U\right\}_{\tau =\tau _{}}{\displaystyle _T_{}^\tau _{}}V𝑑\tau \right]+q{\displaystyle _{\mathrm{}}^T_{}}_aG_{}d\tau `$ (19) $`=`$ $`q[\{\dot{\sigma }^2_a\dot{\sigma }U+\dot{\sigma }^1_aU\}_{\tau =\tau _{}}+\{\dot{\sigma }^2\ddot{\sigma }U+\dot{\sigma }^1\dot{U}\}_{\tau =\tau _{}}_a\tau _{}`$ (21) $`{\displaystyle _T_{}^\tau _{}}_aVd\tau \{V\}_{\tau =\tau _{}}_a\tau _{}]q{\displaystyle }_{\mathrm{}}^T_{}_aG_{}d\tau `$ Since $`\sigma (x,z(\tau _{}))=0`$, we have $$_a\{\sigma \}_{\tau =\tau _{}}=\{_a\sigma \}_{\tau =\tau _{}}+\{\dot{\sigma }\}_{\tau =\tau _{}}_a\tau _{}=0,$$ (22) so that $$_a\tau _{}=\{\dot{\sigma }^1_a\sigma \}_{\tau =\tau _{}}.$$ (23) Therefore, we have $`_a\varphi _{}`$ $`=`$ $`q[\{\dot{\sigma }^2_a\dot{\sigma }U+\dot{\sigma }^1_aU+\dot{\sigma }^3\ddot{\sigma }U_a\sigma \dot{\sigma }^2\dot{U}_a\sigma +\dot{\sigma }^1V_a\sigma \}_{\tau =\tau _{}}`$ (25) $`{\displaystyle _T_{}^\tau _{}}_aVd\tau ]+q{\displaystyle }_{\mathrm{}}^T_{}_aG_{}d\tau `$ In Eqs. (18) and (LABEL:gradphimuv), we would like to combine the integrals which appear on the right side. For $`T_{}\tau <\tau _{}`$, we have $`G(x,z(\tau ))=V(x,z(\tau ))`$. Furthermore, since $`V`$ is a smooth biscalar, $$_T_{}^\tau _{}V𝑑\tau =\underset{ϵ0}{lim}_T_{}^{\tau _{}ϵ}V𝑑\tau $$ (27) and $$_T_{}^\tau _{}_aVd\tau =\underset{ϵ0}{lim}_T_{}^{\tau _{}ϵ}_aVd\tau .$$ (28) Therefore, combining the integrals, we have $$\varphi _{}=q\left\{\dot{\sigma }^1U\right\}_{\tau =\tau _{}}+\underset{ϵ0}{lim}q_{\mathrm{}}^{\tau _{}ϵ}G_{}𝑑\tau .$$ (29) and $`_a\varphi _{}`$ $`=`$ $`q\left\{\dot{\sigma }^2_a\dot{\sigma }U+\dot{\sigma }^1_aU+\dot{\sigma }^3\ddot{\sigma }U_a\sigma \dot{\sigma }^2\dot{U}_a\sigma +\dot{\sigma }^1V_a\sigma \right\}_{\tau =\tau _{}}`$ (31) $`+\underset{ϵ0}{lim}q{\displaystyle _{\mathrm{}}^{\tau _{}ϵ}}_aG_{}d\tau `$ In order to investigate the singularity structure of $`\varphi _{}`$ and $`_a\varphi _{}`$ near the world line, we need expansions for the expressions in brackets on the right sides of Eqs. (29) and (LABEL:gradphimuvtail) which are valid to $`O[r^0]`$. (The integrals in these equations make smooth contributions to the fields.) The required small distance expansions for $`U`$, $`V`$, $`\sigma `$, and their derivatives can all be found in DeWitt and Brehme or derived straightforwardly from expressions given therein. Switching the roles of the primed and unprimed indices for notational simplicity and including the corresponding results for the advanced field, $`\varphi _+`$, we have $$\varphi _\pm (x^{})=q\left(r^1\frac{1}{2}a^a\widehat{r}_a\right)\pm \underset{ϵ0}{lim}q_{\tau _\pm \pm ϵ}^\pm \mathrm{}G_\pm (x^{},z(\tau ))𝑑\tau +O[r]$$ (33) and $`_a^{}\varphi _\pm (x^{})`$ $`=`$ $`q\overline{g}_{a^{}a}(r^2\widehat{r}^a{\displaystyle \frac{1}{2}}r^1a^a+{\displaystyle \frac{1}{2}}r^1(a^b\widehat{r}_b)\widehat{r}^a{\displaystyle \frac{3}{8}}(a^b\widehat{r}_b)^2\widehat{r}^a+{\displaystyle \frac{3}{4}}(a^b\widehat{r}_b)a^a`$ (38) $`{\displaystyle \frac{1}{6}}R_{bdce}u^bu^c\widehat{r}^d\widehat{r}^e\widehat{r}^a{\displaystyle \frac{1}{8}}a^2\widehat{r}^a{\displaystyle \frac{1}{12}}R_{bc}\widehat{r}^b\widehat{r}^c\widehat{r}^a+{\displaystyle \frac{1}{2}}(\dot{a}^b\widehat{r}_b)u^a`$ $`+{\displaystyle \frac{1}{12}}R_{bc}u^bu^c\widehat{r}^a+{\displaystyle \frac{1}{6}}R_{bc}u^b\widehat{r}^cu^a+{\displaystyle \frac{1}{3}}R^a{}_{cbd}{}^{}u_{}^{b}u^c\widehat{r}^d\pm {\displaystyle \frac{1}{3}}a^2u^a{\displaystyle \frac{1}{3}}\dot{a}^a`$ $`{\displaystyle \frac{1}{6}}R_{bc}u^bu^cu^a+{\displaystyle \frac{1}{6}}R^{ab}\widehat{r}_b{\displaystyle \frac{1}{12}}R\widehat{r}^a{\displaystyle \frac{1}{6}}R^{ab}u_b\pm {\displaystyle \frac{1}{12}}Ru^a)`$ $`\pm \underset{ϵ0}{lim}q{\displaystyle _{\tau _\pm \pm ϵ}^\pm \mathrm{}}_aG_\pm (x^{},z(\tau ))𝑑\tau +O[r],`$ where $`\overline{g}_{a^{}a}`$ is the bivector of geodetic parallel displacement, defined by DeWitt and Brehme . We began this calculation in order to investigate what conditions we need to impose on the spacetime neighborhoods and trajectories of scalar particles in different spacetimes and on our identification of these neighborhoods in order to ensure that the subtraction of field gradients in Eq. (6) is finite, and Eq. (38) provides the answer to this question. Since the divergent terms in Eq. (38) depend only upon the four-velocity and four-acceleration of the particle (and not, for example, on higher derivatives of the motion or the local curvature), the subtraction in Eq. (6) will be finite as long as the magnitudes of the four-accelerations of the two particles are equal and we identify the local spacetime neighborhoods in such a way that the four-velocities and four-accelerations, the geodesic distances from the world lines, and the parallel transport defined by $`g_{aa^{}}`$ all coincide up to $`O[r^0]`$. Given points $`P`$ and $`\stackrel{~}{P}`$ on two world lines such that $`a^aa_a=\stackrel{~}{a}^a\stackrel{~}{a}_a`$, we can achieve this by identifying the spacetime neighborhoods of $`P`$ and $`\stackrel{~}{P}`$ with their respective tangent spaces $`T_P`$ and $`T_{\stackrel{~}{P}}`$ via the exponential mapThe exponential map identifies $`v^aT_P`$ with the spacetime point which lies unit affine parameter along the geodesic generated by $`v^a`$., and then identifying $`T_P`$ and $`T_{\stackrel{~}{P}}`$ via any linear map which takes $`u^a`$ to $`\stackrel{~}{u}^a`$ and $`a^a`$ to $`\stackrel{~}{a}^a`$. Under this identification, it is clear that four-velocities, four-accelerations, and geodesic distances will coincide exactly, so we need only check that parallel transport will also agree up to the appropriate order. One way to see this is to write out Eq. (38) explicitly in coordinates adapted to our identification map, so that each point in the neighborhood of $`P`$ is mapped to the point with the same coordinates in the neighborhood of $`\stackrel{~}{P}`$. (Using such coordinates, our map identifies a vector field in the neighborhood of $`P`$ with the vector field in the neighborhood of $`\stackrel{~}{P}`$ having the same coordinate components.) One such coordinate system is Riemann normal coordinates<sup>§</sup><sup>§</sup>§In order to construct Riemann normal coordinates for a neighborhood of point $`P`$, identify points in the neighborhood with points in $`T_P`$ via the exponential map, and then pick any orthonormal basis for $`T_P`$.. In these coordinates, the coordinate components of $`g_{a^{}a}`$ are given by $$\overline{g}_{\alpha \beta }=g_{\alpha \beta }+\frac{1}{6}r^2R_{\alpha \gamma \beta \delta }\widehat{r}^\gamma \widehat{r}^\delta +O[r^3].$$ (39) (We have dropped the primed indices completely since expression relates components rather than tensors.) Comparing this to Eq. (38), we see that $`\overline{g}_{\alpha \beta }`$ simply acts as the identity at this order in $`r`$. \[The term $`(1/6)qR_{\alpha \gamma \beta \delta }\widehat{r}^\gamma \widehat{r}^\delta \widehat{r}^\alpha `$, which arises from the multiplication of the $`r^2`$ term in Eq. (38) and the $`r^2`$ term in Eq. (39), vanishes by the symmetries of the Riemann tensor.\] Therefore, the divergent terms will indeed cancel under the identification we have described. This provides the basis of our crucial Comparison Axiom in the next subsection. ### B The axiomatic approach We are now prepared to give our prescription for $`f_\mathrm{S}^a`$, the total scalar force acting on the particle. We have seen that the subtraction of field gradients in Eq. (6) will be finite as long as the two particles’ four-accelerations have the same magnitude and we identify the spacetime neighborhoods via the exponential map as described above. We now elevate this property to the status of an axiom that any prescription for $`f_\mathrm{S}^a`$ must satisfy. Axiom 1 (Comparison Axiom) Consider two points, $`P`$ and $`\stackrel{~}{P}`$, each lying on timelike world lines in possibly different spacetimes which contain Klein-Gordon fields $`\varphi `$ and $`\stackrel{~}{\varphi }`$ sourced by particles of charge $`q`$ on the world lines. If the four-accelerations of the world lines at $`P`$ and $`\stackrel{~}{P}`$ have the same magnitude, and if we identify the neighborhoods of $`P`$ and $`\stackrel{~}{P}`$ via the exponential map such that the four-velocities and four-accelerations are identified, then the difference between the scalar forces $`f_\mathrm{S}^a`$ and $`\stackrel{~}{f}_\mathrm{S}^a`$ is given by the limit as $`r0`$ of the field gradients, averaged over a sphere at geodesic distance $`r`$ from the world line at $`P`$. $$f_\mathrm{S}^a\stackrel{~}{f}_\mathrm{S}^a=\underset{r0}{lim}q^a\varphi \stackrel{~}{}^a\stackrel{~}{\varphi }_r$$ (40) Since the Comparison Axiom requires only that the four-accelerations of the particles agree, we now need only fix the dependence of $`f_\mathrm{S}^a`$ on acceleration in some arbitrary spacetime in order to uniquely determine $`f_\mathrm{S}^a`$. Motivated by the time-reflection symmetry of the half-advanced, half-retarded solution for a uniformly accelerating trajectory in flat spacetime, we impose the following axiom, which should be familiar from electromagnetism. Axiom 2 (Flat spacetime axiom) If $`(M,g_{ab})`$ is Minkowski spacetime, the world line is uniformly accelerating, and $`\varphi `$ is the half-advanced, half-retarded solution, $`\varphi =\frac{1}{2}(\varphi _++\varphi _{})`$, then $`f_\mathrm{S}^a=0`$ at every point on the world line. We will now show that, if there exists a prescription for $`f_\mathrm{S}^a`$ satisfying these two axioms, it must be unique. Consider a point $`P`$ on the world line of a scalar particle of charge $`q`$ in some spacetime, and let the particle have acceleration $`a^a`$ at point $`P`$. Let $`f_\mathrm{S}^a`$ and $`g_\mathrm{S}^a`$ be two prescriptions for the scalar force, both satisfying the axioms given above. Now consider a uniformly accelerating particle with the same charge $`q`$ and the same acceleration $`a^a`$ in a flat spacetime $`(^4,\eta _{ab})`$, and construct the half-advanced, half-retarded solution $`\stackrel{~}{\varphi }=\frac{1}{2}(\stackrel{~}{\varphi }_++\stackrel{~}{\varphi }_{})`$ for this particle. By our second axiom, we know that $`\stackrel{~}{f}_\mathrm{S}^a=\stackrel{~}{g}_\mathrm{S}^a=0`$ at every point $`\stackrel{~}{P}`$ along the world line of this uniformly accelerating particle. Therefore, identifying the neighborhoods of $`P`$ and $`\stackrel{~}{P}`$ as in the Comparison Axiom above, we have $$f_\mathrm{S}^ag_\mathrm{S}^a=(f_\mathrm{S}^a\stackrel{~}{f}_\mathrm{S}^a)(g_\mathrm{S}^a\stackrel{~}{g}_\mathrm{S}^a)=\underset{r0}{lim}q^a\varphi \stackrel{~}{}^a\stackrel{~}{\varphi }_r\underset{r0}{lim}q^a\varphi \stackrel{~}{}^a\stackrel{~}{\varphi }_r=0.$$ (41) This argument establishes uniqueness, but it also demonstrates existence by providing a prescription which is guaranteed to satisfy the axioms. Namely, given a point $`P`$ along the world line of a scalar particle with charge $`q`$ in any spacetime, we simply construct the half-advanced, half-retarded solution $`\stackrel{~}{\varphi }`$ for a uniformly accelerating particle in flat spacetime with the same charge and acceleration. The scalar force $`f_\mathrm{S}^a`$ is then given by $$f_\mathrm{S}^a=\underset{r0}{lim}q^a\varphi \stackrel{~}{}^a\stackrel{~}{\varphi }_r.$$ (42) This is the prescription for the total scalar force which we set out to find at the beginning of this section. Writing $`\varphi `$ as $`\varphi =\varphi _{\mathrm{in}}+\varphi _{}`$, we can use Eq. (38) to turn this prescription into an explicit formula for $`f_\mathrm{S}^a`$. The result is $`f_\mathrm{S}^a`$ $`=`$ $`q^a\varphi _{\mathrm{in}}+q^2\left({\displaystyle \frac{1}{3}}(\dot{a}^aa^2u^a)+{\displaystyle \frac{1}{6}}(R^{ab}u_b+R_{bc}u^bu^cu^a){\displaystyle \frac{1}{12}}Ru^a\right)`$ (44) $`+\underset{ϵ0}{lim}q^2{\displaystyle _{\mathrm{}}^{\tau ϵ}}^aG_{}(z(\tau ),z(\tau ^{}))𝑑\tau ^{}.`$ This expression, which is the main result of the paper, allows us to calculate $`f_\mathrm{S}^a`$ for any trajectory $`z(\tau )`$ in any spacetime. As stated at the beginning of the section, the physical significance of this expression is that it should correctly describe the order $`q`$ and $`q^2`$ contributions to the force on a nearly spherical extended body in the point particle limit. ## III The equations of motion We now wish to consider the special case in which no non-scalar forces are present, so that the evolution of the world line $`z(\tau )`$ is determined by the scalar field. In the next subsection, we derive equations of motion for $`z(\tau )`$ in this case. Then, in the following subsection, we explore one of the consequences of these equations of motion: that the mass of particle varies with time. ### A Derivation of the equations of motion Consider once again the extended body described in Sec. II. In the absence of non-scalar fields, conservation of stress-energy dictates that $$_bT_{\mathrm{body}}^{ab}=_bT_\mathrm{S}^{ab}.$$ (45) According to the arguments of Quinn and Wald , in the point particle limit, the center of mass world line $`z(\tau )`$ will therefore satisfy $$u^b_b(mu^a)=\frac{dm}{dt}u^a+ma^a=f_\mathrm{S}^a,$$ (46) where $`f_\mathrm{S}^a`$ is the limiting force we derived in Sec. II. Inserting our expression for $`f_\mathrm{S}^a`$ from Eq. (44) and separating the components parallel to $`u^a`$ and perpendicular to $`u^a`$, we have $`a^a`$ $`=`$ $`{\displaystyle \frac{1}{m}}(f_\mathrm{S}^a+u^ag_{bc}u^bf_\mathrm{S}^c)`$ (47) $`=`$ $`{\displaystyle \frac{q}{m}}(^a\varphi _{\mathrm{in}}+u^au^b_b\varphi _{\mathrm{in}})+{\displaystyle \frac{q^2}{m}}\left({\displaystyle \frac{1}{3}}(\dot{a}^aa^2u^a)+{\displaystyle \frac{1}{6}}(R^{ab}u_b+R_{bc}u^bu^cu^a)\right)`$ (49) $`+\underset{ϵ0}{lim}{\displaystyle \frac{q^2}{m}}{\displaystyle _{\mathrm{}}^{\tau ϵ}}(^aG_{}+u^ag_{bc}u^b^cG_{})𝑑\tau ^{}`$ and $$\frac{dm}{d\tau }=f_\mathrm{S}^au_a=qu^a_a\varphi _{\mathrm{in}}\frac{1}{12}q^2R\underset{ϵ0}{lim}q^2_{\mathrm{}}^{\tau ϵ}u_a^aG_{}d\tau ^{}.$$ (50) We now note three important features of these equations. First, for each point along the world line, the integrals in these expressions represent that portion of $`^a\varphi _{}`$ which arises from source contributions interior to the past light cone of the point. This contribution to the force, often called the “tail term,” is a direct consequence of the failure of Huygen’s principle in curved spacetime, and can be understood as the result of scalar radiation backscattering from the background curvature and re-intersecting the particle world line. The presence of this tail term is the primary obstacle to applying these equations in physically realistic situations, since most methods for calculating the retarded field of an arbitrary world line irretrievably mix the tail and non-tail portions of the field. Secondly, we can provide further insight into the nature of Eq. (50) by tracing the origin of the Ricci scalar term in the Hadamard expansion of the field given in Sec. II. This term arises directly from the $`\{V\}_{\tau =\tau _{}}_a\tau _{}`$ term in Eq. (21). In particular, we have $$\underset{\tau ^{}\tau }{lim}G_{}(z(\tau ),z(\tau ^{}))=\frac{1}{12}R,$$ (51) so that we can rewrite Eq. (50) as $$\frac{dm}{d\tau }=qu^a_a(\varphi _{\mathrm{in}}+\varphi _{\mathrm{tail}}),$$ (52) where $`\varphi _{\mathrm{tail}}`$ is defined by $$\varphi _{\mathrm{tail}}=\underset{ϵ0}{lim}q_{\mathrm{}}^{\tau ϵ}G_{}𝑑\tau ^{}.$$ (53) The implications of Eq. (52) for global energy conservation are explored by Quinn and Wald . Finally, owing to the presence of the Abraham-Lorentz $`\dot{a}^a`$ term, these equations share the unphysical “runaway” solutions which have been discussed thoroughly in the electromagnetic case. (See Jackson for one such discussion.) In order to interpret these solutions, it is important to remember that we view the force law given by Eq. (44) as an approximate expression for the force on an extended body, valid to $`O[q^2]`$, rather than a fundamental description of a point particle. Therefore, we can eliminate these unphysical solutions through the reduction of order technique. This technique is discussed in detail by Flanagan and Wald , but the basic idea is simple. Recall that we wish Eq. (49) to describe the limiting motion of a one-parameter family of extended bodies in which both the charge and the mass of the bodies vanish as the parameter goes to zero. For concreteness, let us assume that the charge and mass are given by $`q=aϵ`$ and $`m=bϵ`$. In order to apply the reduction of order technique to Eq. (49), we simply insert the entire right side of the equation in place of $`a^a`$ in the $`\dot{a}^a`$ and $`a^2u^a`$ terms and discard any resulting terms which are $`O[ϵ^2]`$ or higher. The result is $`a^a`$ $`=`$ $`{\displaystyle \frac{q}{m}}(^a\varphi _{\mathrm{in}}+u^au_b^b\varphi _{\mathrm{in}})`$ (56) $`+{\displaystyle \frac{1}{3}}{\displaystyle \frac{q^2}{m}}\left({\displaystyle \frac{q}{m}}\left(u^b_b^a\varphi _{\mathrm{in}}+u^au^bu^c_b_c\varphi _{\mathrm{in}}\right){\displaystyle \frac{q^2}{m^2}}\left(^b\varphi _{\mathrm{in}}_b\varphi _{\mathrm{in}}+(u^b_b\varphi _{\mathrm{in}})^2\right)u^a\right)`$ $`+{\displaystyle \frac{1}{6}}{\displaystyle \frac{q^2}{m}}(R^{ab}u_b+R_{bc}u^bu^cu^a)+\underset{ϵ0}{lim}{\displaystyle \frac{q^2}{m}}{\displaystyle _{\mathrm{}}^{\tau ϵ}}(^aG_{}+u^ag_{bc}^bG_{}u^c)𝑑\tau ^{},`$ which is free of the unphysical runaway solutions. ### B Time variation of the mass In stark contrast to the electromagnetic case, $`f_\mathrm{S}^a`$ includes contributions which point along the four-velocity of the particle, resulting in a time-varying mass. This is not a special feature of the self force, nor of curved spacetime. Rather, it reflects a fundamental difference between the two continuum theories. Consider a small body in Minkowski spacetime with a center of mass world line $`z(\tau )`$. The rest mass of such a body is given by $$m=_\mathrm{\Sigma }u_aT_{\mathrm{body}}^{ab}ϵ_{bcde},$$ (57) where $`u^a`$ is the four-velocity of $`z(\tau )`$ (defined away from the world line by global parallelism), $`\mathrm{\Sigma }`$ is the surface perpendicular to $`u^a`$, and $`ϵ_{abcd}`$ is the volume element compatible with the (flat) metric. Therefore, we have $`{\displaystyle \frac{dm}{d\tau }}`$ $`=`$ $`{\displaystyle \frac{d}{d\tau }}{\displaystyle _\mathrm{\Sigma }}u_aT_{\mathrm{body}}^{ab}ϵ_{bcde}`$ (58) $`=`$ $`{\displaystyle _\mathrm{\Sigma }}\mathrm{\pounds }_w[u_aT_{\mathrm{body}}^{ab}ϵ_{bcde}]`$ (59) $`=`$ $`{\displaystyle _\mathrm{\Sigma }}u_a_bT_{\mathrm{body}}^{ab}w^cϵ_{cdef},`$ (60) where $`w^a`$ is the vector field which connects successive time slices $`\mathrm{\Sigma }(\tau )`$. For a body coupled to a scalar field, we have $`_bT_{\mathrm{body}}^{ab}=_bT_\mathrm{S}^{ab}=\rho ^a\varphi `$, so that $$\frac{dm}{d\tau }=_\mathrm{\Sigma }\rho u_a^a\varphi w^cϵ_{cdef},$$ (61) which is clearly, in general, nonvanishing. By contrast, in the electromagnetic case, we have $`_bT_{\mathrm{body}}^{ab}=_bT_{\mathrm{EM}}^{ab}=F^{ab}j_b`$, so that $$\frac{dm}{d\tau }=_\mathrm{\Sigma }u_aF^{ab}j_b\varphi w^cϵ_{cdef}.$$ (62) For typical models of charged matter, $`j^a`$ and $`u^a`$ will become collinear as we take the point particle limit, and $`dm/dt`$ will vanish. Perhaps because it is tempting to generalize from the more familiar electromagnetic case, this time variation of the mass in the scalar case has largely been ignored in the literature. Some authors use the equation of motion $`ma^a=q^a\varphi `$ (e.g., Shapiro and Teukolsky ). This equation is clearly inconsistent, and therefore in general has no solutions, since $`a^a`$ is perpendicular to the four-velocity while $`^a\varphi `$, in general, is not. Others explicitly project $`^a\varphi `$ perpendicular to the four-velocity as in Eq. (49) above in order to obtain the acceleration of the particle, but then simply ignore the component of $`^a\varphi `$ which points along $`u^a`$ and assume that the mass is constant (e.g., Ori ). While such an equation of motion is mathematically consistent, it violates global conservation of stress-energy. (See Quinn and Wald .) In the discussion above, we have motivated our point particle equations of motion by imposing local stress-energy conservation on continuum matter and taking the point particle limit, using our axioms to extract the appropriate finite part of the divergent fields. The time variation of the mass arises as a direct consequence of this local stress-energy conservation. In the literature on point particles, one sometimes sees an alternative derivation which makes no reference to the continuum theory. Instead, the author defines an action for the point particle system and then formally minimizes this action with respect to variations of the fields and the world line in order to produce equations of motion. For completeness, we give such a derivation here, paying particular attention to the time dependence of the particle’s mass. Fix a globally hyperbolic spacetime $`(M,g_{ab})`$ and two Cauchy surfaces for the spacetime, $`C_1`$ and $`C_2`$. Let $`\varphi `$ be a smooth scalar field and $`z(\tau )`$ be a smooth world line in the region $`V`$ between $`C_1`$ and $`C_2`$. We fix the value of $`\varphi `$ and the position of $`z(\tau )`$ on $`C_1`$ and $`C_2`$ and define the action, $`S`$, as $`S`$ $`=`$ $`{\displaystyle _V}\left[{\displaystyle \frac{1}{8\pi }}\left(g^{ab}_a\varphi _b\varphi \right)+{\displaystyle \frac{1}{2}}{\displaystyle mg_{ab}u^au^b\delta ^4(xz(\tau ))𝑑\tau }+{\displaystyle q\varphi \delta ^4(xz(\tau ))𝑑\tau }\right]ϵ_{abcd}.`$ (63) Formally minimizing this action with respect to variations of $`\varphi `$, we arrive at $$^a_a\varphi =4\pi q\delta ^4(xz(\tau ))𝑑\tau ,$$ (64) while minimization with respect to variations of $`z(\tau )`$ yields $$\frac{dm}{d\tau }u^a+ma^a=q^a\varphi .$$ (65) These are the same equations we arrived at by considering the point particle limit of the continuum theory. Of course, here we have assumed $`\varphi `$ and $`z(\tau )`$ to be smooth in order to define the action, while the solutions of Eq (64) are clearly distributional. Therefore, no solutions of these equations exist. However, we may view this as a formal derivation of our equations from an action principle. Note that, if we had assumed from the outset that $`m`$ was constant, the only change to the equations would have been to set $`dm/d\tau =0`$ in Eq. (65). Clearly, the resulting equation is inconsistent, since $`^a\varphi `$ does not, in general, point along the four-velocity. Still, one might wonder, despite the stress-energy conservation arguments given above, whether the above action can be modified to produce the equation of motion $$\frac{dm}{d\tau }u^a+ma^a=q(^a\varphi +u^ag_{bc}u^b^c\varphi ),$$ (66) since this equation would have the immediate consequence that $`dm/d\tau =0`$, as in the electromagnetic case. Wiseman has considered a large class of possible coupling terms and has found that, within this class, one cannot produce Eq. (66) without introducing a nonlinear coupling on the right side of Eq. (64). Based on this work and the stress-energy considerations discussed above, we conjecture that there exists no action which produces Eq. (66) while preserving Eq. (64). ## Acknowledgements I would like to thank Robert Wald for his expert guidance over the course of the last several years. I would also like to thank Alan Wiseman for many useful discussions. This research was supported in part by NSF grant PHY95-14726 to the University of Chicago.
warning/0005/math0005292.html	ar5iv	text	# Flat Lorentz 3-manifolds and cocompact Fuchsian groups ## 1. Introduction Consider Minkowski $`2+1`$-space $`𝔼`$ and let $`G\mathrm{SO}(2,1)^0`$ be a discrete subgroup. Suppose that a group of affine isometries of $`𝔼`$ with linear part $`G`$ acts properly and freely on $`𝔼`$. In a remarkable preprint , Geoffrey Mess proved the following theorem: ###### Theorem. $`G`$ is not cocompact in $`\mathrm{SO}(2,1)^0`$. Mess deduces this result as part of a general theory of domains of dependence in constant curvature Lorentzian 3-manifolds. We give an alternate proof, using an invariant introduced by Margulis and Teichmüller theory. We thank Scott Wolpert for helpful conversations concerning Teichmüller theory. We also wish to thank Paul Igodt and the Algebra Research Group at the Katholieke Universiteit Leuven at Kortrijk, Belgium for their hospitality at the “Workshop on Crystallographic Groups and their Generalizations II”, where these results were obtained. ## 2. Background Let $`^{2,1}`$ be a 3-dimensional real vector space with inner product $$𝔹(x,y)=x_1y_1+x_2y_2x_3y_3.$$ The group of linear isometries of $`^{2,1}`$ will be denoted by $`\mathrm{SO}(2,1)`$. Let $`\mathrm{Isom}(^{2,1})`$ denote the group of affine isometries, that is, the group of all transformations of the form $`h:^{2,1}`$ $`^{2,1}`$ $`x`$ $`g(x)+u`$ where $`g\mathrm{O}(2,1)`$ and $`u^{2,1}`$. We write $`g=𝕃(h)`$ and $`h=(g,u)`$. Evidently $`\mathrm{Isom}(^{2,1})`$ is isomorphic to the semidirect product $`\mathrm{O}(2,1)^{2,1}`$ where $`^{2,1}`$ denotes the vector group of translations of $`𝔼`$. Let $`G\mathrm{O}(2,1)`$ be a subgroup. An affine deformation of $`G`$ is a homomorphism $`\varphi :G\mathrm{Isom}(^{2,1})`$ such that $`𝕃(\varphi (g))=g`$. An affine deformation $`\varphi `$ is proper if the resulting action of $`G`$ by affine transformations on $`^{2,1}`$ is a proper action. Write $$\varphi (g)=(g,u(g)).$$ The condition that $`\varphi `$ be a homomorphism is that the map $`u=u_\varphi :G^{2,1}`$ satisfy the cocycle condition (1) $$u_\varphi (g_1g_2)=u_\varphi (g_1)+g_1u_\varphi (g_2).$$ A map $`u:G^{2,1}`$ satisfying (1) is called a cocycle and the vector space of cocycles is denoted by $`Z^1(G,^{2,1})`$. If $`\varphi _1,\varphi _2`$ are affine deformations of $`G`$ which are conjugate by translation by $`v^{2,1}`$, then the difference $`u_{\varphi _1}u_{\varphi _2}`$ is the cocycle $$\delta v:gvg(v).$$ Such a cocycle is called a coboundary. The subspace of coboundaries is denoted by $`B^1(G,^{2,1})`$. We say that $`\varphi _1,\varphi _2`$ are translationally conjugate. Translational conjugacy classes of affine deformations of $`G`$ correspond to elements in the cohomology group $$H^1(G,^{2,1})=Z^1(G,^{2,1})/B^1(G,^{2,1}).$$ Suppose that $`\varphi :G\mathrm{Isom}(^{2,1})`$ is a proper affine deformation. By Fried-Goldman , the group $`G`$ is solvable or the linear part $$𝕃\varphi :G\mathrm{O}(2,1)$$ is an isomorphism onto a discrete subgroup of $`\mathrm{O}(2,1)`$. (Indeed, this conclusion is obtained for any proper affine action on $`^3`$.) The solvable groups are easily classified by embedding them as lattices in Lie subgroups which themselves act properly. When $`G`$ is not solvable, then interesting examples do exist (Margulis ). Furthermore every torsionfree non-cocompact discrete subgroup $`G\mathrm{O}(2,1)`$ for which $`H^1(G;^{2,1})0`$ admits proper affine deformations (Drumm ). Recall that an element of $`\mathrm{O}(2,1)`$ is hyperbolic if it has three distinct real eigenvalues. A subgroup $`G\mathrm{O}(2,1)`$ is purely hyperbolic if every element is hyperbolic. A cocompact discrete subgroup contains a purely hyperbolic subgroup of finite index . ## 3. An invariant of affine isometries In , Margulis defines an invariant $`\alpha _\varphi :G`$ of an affine deformation $`\varphi `$ of a purely hyperbolic subgroup $`G\mathrm{O}(2,1)`$ as follows. We assume that $`G\mathrm{SO}(2,1)^0`$. Choose a component $`𝒩_+`$ of the complement of $`0`$ in the lightcone. Since any element $`g`$ of $`G`$ is hyperbolic its three eigenvalues are distinct positive real numbers $$\lambda (g)<1<\lambda (g)^1.$$ Choose an eigenvector $`𝗑^{}(g)𝒩_+`$ for $`\lambda (g)`$ and an eigenvector $`𝗑^+(g)𝒩_+`$ for $`\lambda (g)^1`$, respectively. Then there exists a unique eigenvector $`𝗑^0(g)`$ for $`g`$ with eigenvalue $`1`$ such that: * $`𝔹(𝗑^0(g),𝗑^0(g))=1`$; * $`(𝗑^{}(g),𝗑^+(g),𝗑^0(g))`$ is a positively oriented basis. Notice that $`𝗑^0(g^1)=𝗑^0(g)`$. If $`\varphi `$ is an affine deformation corresponding to a cocycle $`u`$, then $`\alpha _\varphi `$ is defined as: (2) $`\alpha _\varphi :G`$ $``$ $`g`$ $`𝔹(𝗑^0(g),u(g)).`$ More generally, $`\alpha _\varphi (g)=𝔹(𝗑^0(g),\varphi (g)(x)x))`$ for any $`x𝔼`$. Furthermore $`\alpha _\varphi `$ is a class function on $`G`$ and recently Drumm-Goldman have proved that the mapping $`H^1(G,^{2,1})`$ $`^G`$ $`[u]`$ $`\alpha _\varphi `$ is injective, that is, $`\alpha `$ is a complete invariant of the conjugacy class of the affine deformation. In , Margulis proved the following theorem (see also Drumm ): ###### Theorem 1 (Margulis). Suppose that $`G\mathrm{SO}(2,1)^0`$ is purely hyperbolic and let $`\varphi :G\mathrm{Isom}(^{2,1})`$ be an affine deformation. If there exist $`g_1,g_2G`$ such that $`\alpha _\varphi (g_1)>0>\alpha _\varphi (g_2)`$, then $`\varphi `$ is not proper. Affine deformations defining free actions correspond to cocycles for which $`\alpha (g)0`$ for $`g𝕀`$. We shall say that a cocycle $`u`$ is positive (respectively negative) if $`\alpha (g)>0`$ (respectively $`\alpha (g)<0`$) whenever $`𝕀gG`$. Clearly $`u`$ is positive if and only if $`u`$ is negative. We conjecture a converse to Theorem 1: an affine deformation is proper if and only if its cocycle is positive or negative. ## 4. Deformation-theoretic interpretation of $`\alpha `$ We reduce the proof of Mess’s theorem to facts about deformations of hyperbolic Riemann surfaces. Let $`M`$ be a surface with a complete hyperbolic structure and $`\pi =\pi _1(M)`$ its fundamental group. A representation $`\varphi :\pi \mathrm{SO}(2,1)^0`$ is Fuchsian if it is an embedding onto a discrete subgroup of $`\mathrm{SO}(2,1)^0`$. When $`M`$ is a closed surface, the space of conjugacy classes of Fuchsian representations $`\varphi :\pi \mathrm{SO}(2,1)^0`$ is an open subset of the space of conjugacy classes of all representations, which identifies with the Teichmüller space $`𝔗(M)`$ of $`M`$. (See Weil , §VI of Raghunathan for the general theory and Goldman for the case of surface groups.) Its tangent space identifies with the cohomology group $`H^1(G,^{2,1})`$ where $`G=\varphi (\pi )`$. Since the classical theory of Fuchsian groups is usually phrased in terms of $`\mathrm{SL}(2,)`$ (rather than $`\mathrm{SO}(2,1)`$), and since $`2\times 2`$ matrices are more tractable than $`3\times 3`$ matrices, we work with $`\mathrm{SL}(2,)`$. The Lie groups $`\mathrm{SL}(2,)`$ and $`\mathrm{SO}(2,1)`$ are locally isomorphic, but not globally isomorphic. One model for the local isomorphism is the adjoint representation, as follows. The trace form of any nontrivial representation (for example the Killing form) provides the Lie algebra $`𝔰𝔩(2,)`$ with a Lorentzian inner product invariant under the adjoint representation. Thus $`𝔰𝔩(2,)`$ is isometric to $`^{2,1}`$; we give an explicit orthogonal basis. In this way the adjoint representation $`\mathrm{Ad}:\mathrm{SL}(2,)\mathrm{Isom}(𝔰𝔩(2,))`$ defines a local isomorphism $`\rho :\mathrm{SL}(2,)\mathrm{SO}(2,1)`$ of Lie groups. The local isomorphism $`\rho :\mathrm{SL}(2,)\mathrm{O}(2,1)`$ is not injective — its kernel consists of the center $`\{\pm 𝕀\}`$ of $`\mathrm{SL}(2,)`$. Nor is $`\rho `$ surjective — its image is the identity component $`\mathrm{SO}^0(2,1)`$ of $`\mathrm{O}(2,1)`$. Neither issue is problematic here, since purely hyperbolic discrete subgroups of $`\mathrm{SO}(2,1)`$ lift to subgroups of $`\mathrm{SL}(2,)`$ (Abikoff , Culler , Kra ). Let $`G`$ be a purely hyperbolic subgroup of $`\mathrm{SO}(2,1)`$, with inclusion $`\iota :G\mathrm{SO}(2,1)`$. Then there exists a representation $`\stackrel{~}{\iota }:G\mathrm{SL}(2,)`$ such that $`\iota =\rho \stackrel{~}{\iota }`$. Furthermore composition with the local isomorphism $`\rho `$ induces a covering space $$\mathrm{Hom}(G,\mathrm{SL}(2,))\mathrm{Hom}(G,\mathrm{Isom}^0(^{2,1})).$$ Thus smooth paths in $`\mathrm{Hom}(G,\mathrm{Isom}^0(^{2,1}))`$ lift to $`\mathrm{Hom}(G,\mathrm{SL}(2,))`$. Henceforth we suppress $`\stackrel{~}{\iota }`$ (identifying $`G`$ with its image $`\stackrel{~}{\iota }(G)`$ in $`\mathrm{SL}(2,)`$) and consider paths in $`\mathrm{Hom}(G,\mathrm{SL}(2,))`$. ## 5. $`𝔰𝔩(2,)`$ and $`^{2,1}`$ For the calculations later, we now give a detailed description of the local isomorphism $`\rho `$ derived from the adjoint representation. For convenience, consider the Lie algebra $`𝔰𝔩(2,)`$ with inner prouct (3) $$𝔹(X,Y):=\frac{1}{2}\mathrm{tr}(XY).$$ The basis $$e_1=\left[\begin{array}{cc}1& 0\\ 0& 1\end{array}\right],e_2=\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right],e_3=\left[\begin{array}{cc}0& 1\\ 1& 0\end{array}\right].$$ is orthogonal with respect to $`𝔹`$ and satisfies $$𝔹(e_1,e_1)=𝔹(e_2,e_2)=1,𝔹(e_3,e_3)=1.$$ This provides an isometry of Lorentzian vector spaces $`\psi :𝔰𝔩(2,)`$ $`^{2,1}`$ $`\left[\begin{array}{cc}v_1& v_2\\ v_3& v_1\end{array}\right]`$ $`\left[\begin{array}{c}v_1\\ (v_2+v_3)/2\\ (v_2+v_3)/2\end{array}\right].`$ With respect to this isometry the adjoint representation defines a local isomorphism $`\rho :\mathrm{SL}(2,)\mathrm{O}(2,1)`$ satisfying: $$\psi (\mathrm{Ad}(g)v)=\rho (g)\psi (v)$$ whenever $`g\mathrm{SL}(2,)`$ and $`v𝔰𝔩(2,)`$. (In other words, $`\psi :𝔰𝔩(2,)_{\mathrm{Ad}}^{2,1}`$ is $`\rho `$-equivariant.) Explicitly, $`\mathrm{SL}(2,)`$ $`\stackrel{\rho }{}\mathrm{O}(2,1)`$ $`\left[\begin{array}{cc}a& b\\ c& d\end{array}\right]`$ $`\left[\begin{array}{ccc}1+2bc& ac+bd& ac+bd\\ ab+cd& (a^2b^2c^2+d^2)/2& (a^2b^2+c^2+d^2)/2\\ ab+cd& (a^2+b^2c^2+d^2)/2& (a^2+b^2+c^2+d^2)/2\end{array}\right]`$ (where $`adbc=1`$). Differentiation at $`𝕀\mathrm{SL}(2,)`$ (that is, at $`a=d=1`$, $`b=c=0`$) gives the Lie algebra isomorphism $`𝔰𝔩(2,)`$ $`𝔬(2,1)`$ $`\left[\begin{array}{cc}v_1& v_2\\ v_3& v_1\end{array}\right]`$ $`\left[\begin{array}{ccc}0& v_3v_2& v_2+v_3\\ v_2v_3& 0& 2v_1\\ v_2+v_3& 2v_1& 0\end{array}\right].`$ An element $`g\mathrm{SL}(2,)`$ is hyperbolic if it has two real distinct eigenvalues, which are necessarily reciprocal. If $`g`$ has eigenvalues $`\mu ,\mu ^1`$ with $`\|\mu \|<1`$, then $`\rho (g)`$ has eigenvalues $`\lambda =\mu ^2,1,\mu ^2`$. In particular $`g\mathrm{SL}(2,)`$ is hyperbolic if and only if $`\rho (g)`$ is hyperbolic. There exists $`f\mathrm{SL}(2,)`$ such that $$fgf^1=g_0$$ where $$g_0=\pm \left[\begin{array}{cc}\mu & 0\\ 0& \mu ^1\end{array}\right]$$ and $$0<\mu <1<\mu ^1.$$ The eigenvectors of $`g_0=\rho (g_0)`$ are: $`𝗑^{}(g_0)`$ $`=\psi \left(\left[\begin{array}{cc}0& 2\\ 0& 0\end{array}\right]\right)=\left[\begin{array}{c}0\\ 1\\ 1\end{array}\right]`$ $`𝗑^+(g_0)`$ $`=\psi \left(\left[\begin{array}{cc}0& 0\\ 2& 0\end{array}\right]\right)=\left[\begin{array}{c}0\\ 1\\ 1\end{array}\right]`$ $`𝗑^0(g_0)`$ $`=\psi \left(\left[\begin{array}{cc}1& 0\\ 0& 1\end{array}\right]\right)=\left[\begin{array}{c}1\\ 0\\ 0\end{array}\right].`$ The eigenvectors for $`g`$ are the images of the eigenvectors of $`g_0`$ under $`f`$. Now we derive a formula for $`\alpha (g)`$ for an affine deformation $`\varphi `$ which is of the form $`h=(\rho (g),\psi (v)(g))`$ where $`gG\mathrm{SL}(2,)`$ and $`v𝔰𝔩(2,)`$. Suppose that $`g\mathrm{SL}(2,)`$ is hyperbolic. We use the embedding $`\mathrm{SL}(2,)𝔤𝔩(2,)`$. Orthogonal projection $`𝔤𝔩(2,)`$ $`\stackrel{\mathrm{\Pi }}{}𝔰𝔩(2,)`$ $`g`$ $`g{\displaystyle \frac{\mathrm{tr}(g)}{2}}𝕀.`$ maps $`g_0`$ to a diagonal matrix of trace zero. Dividing $`\mathrm{\Pi }(g_0)`$ by $$\mathrm{sgn}(\mathrm{tr}(g))\sqrt{det(\mathrm{\Pi }(g_0))}$$ gives the diagonal matrix corresponding to $`𝗑^0(g_0)^{2,1}`$ (where $`\mathrm{sgn}(x)`$ denotes the sign of a nonzero real number $`x`$). Since $`\mathrm{tr}(g_0)=\pm (\mu +\mu ^1)`$, $$det(\mathrm{\Pi }(g_0))=(\mu \mu ^1)^2=(\mathrm{tr}(g_0)^24)/4$$ so $`\mathrm{sgn}(\mathrm{tr}(g_0))`$ $`\mathrm{\Pi }(g_0)/\sqrt{det(\mathrm{\Pi }(g_0))}`$ $`=\mathrm{sgn}(\mathrm{tr}(g_0))\left(g_0{\displaystyle \frac{\mathrm{tr}(g_0)}{2}}𝕀\right)/\left({\displaystyle \frac{\sqrt{\mathrm{tr}(g_0)^24}}{2}}\right)`$ $`=\left[\begin{array}{cc}1& 0\\ 0& 1\end{array}\right]`$ corresponds to $`𝗑^0(g)`$. Conjugation by $`f`$ gives the general formula (4) $$\psi :\mathrm{sgn}(\mathrm{tr}(g))\left(g\frac{\mathrm{tr}(g)}{2}𝕀\right)/\left(\frac{\sqrt{\mathrm{tr}(g)^24}}{2}\right)𝗑^0(g)$$ From (4) follows a formula for $`\alpha (g)`$ in terms of traces. Suppose that $`G\mathrm{SL}(2,)`$ is purely hyperbolic and $`uZ^1(G,𝔰𝔩(2,))Z^1(G,^{2,1})`$. Taking the trace of the product of (4) with $`u(g)`$, and applying (2) and (3) yields: (5) $$\alpha (g)=\mathrm{sgn}(\mathrm{tr}(g))\frac{\mathrm{tr}\left(u(g)g\right)}{\sqrt{\mathrm{tr}(g)^24}}$$ ## 6. Trace and displacement length Let $`\mathrm{Hyp}`$ denote the subset of $`\mathrm{SL}(2,)`$ consisting of hyperbolic elements. The image of the trace function $`\mathrm{tr}:\mathrm{Hyp}`$ consists of the disjoint two intervals $`(\mathrm{},2)`$ and $`(2,\mathrm{})`$. Furthermore hyperbolic elements $`g\mathrm{Hyp}`$ are determined up to conjugacy by their trace. In terms of hyperbolic geometry, $`\mathrm{tr}(g)`$ relates to the displacement length $`\mathrm{}(g)`$, that is, the minimum distance $`g`$ moves a point $`x\mathrm{H}_{}^2`$. This minimum is realized when $`x`$ lies in the $`g`$-invariant geodesic, which is necessarily unique. Equivalently $`\mathrm{}(g)`$ is the length of the shortest homotopically nontrivial closed curve in the quotient $`\mathrm{H}_{}^2/g`$. Such a shortest curve is necessarily a simple closed geodesic. Let $`\stackrel{~}{g}\mathrm{SL}(2,)`$ be a lift of $`g\mathrm{Isom}(^{2,1})`$ to $`\mathrm{SL}(2,)`$, that is, $`g=\rho (\stackrel{~}{g})`$. Displacement length of $`g`$ relates to $`\mathrm{tr}(\stackrel{~}{g})`$ and the eigenvalue $`0<\mu <1`$ by: $`\mathrm{}(g)`$ $`=2\mathrm{log}\mu `$ $`\|\mathrm{tr}(\stackrel{~}{g})\|`$ $`=2\mathrm{cosh}(\mathrm{}(g)/2)`$ (the sign of $`\mathrm{tr}(\stackrel{~}{g})`$ is ambiguous since $`\mathrm{ker}(\rho )=\{\pm 𝕀\}`$). Since (6) $$\frac{d\|\mathrm{tr}\|}{d\mathrm{}}=\mathrm{sinh}(\mathrm{}/2)>0$$ trace depends monotonically on displacement length. Associated to a cocycle $`uZ^1(G,^{2,1})`$ are real analytic paths $`\stackrel{~}{\iota }_t`$ in $`\mathrm{Hom}(G,\mathrm{SL}(2,))`$ of the form $$\stackrel{~}{\iota }_t(g)=g\mathrm{exp}\left(tu(g)+O(t^2)\right)$$ where $`t`$ is defined in an open interval $`I_g`$ containing zero. (In general $`I_g`$ may depend on $`g`$.) We say that the cocycle $`u`$ is tangent to the path $`\stackrel{~}{\iota }_t`$. Given a path $`\stackrel{~}{\iota }_t\mathrm{Hom}(G,\mathrm{SL}(2,))`$ where $`\stackrel{~}{\iota }_t(G)\mathrm{Hyp}`$, consider the two functions $`\tau _g:I_g`$ $``$ $`t`$ $`\left\|\mathrm{tr}\left(\stackrel{~}{\iota }_t(g)\right)\right\|`$ and $`L_g:I_g`$ $``$ $`t`$ $`\mathrm{}\left(\stackrel{~}{\iota }_t(g)\right).`$ When $`\stackrel{~}{\iota }_t`$ corresponds to a path $`\mu (t)`$ in $`𝔗(M)`$, then $`L_g=\mathrm{}_g\mu `$ where $`\mathrm{}_g:𝔗(M)`$ is the geodesic length function associated to $`g`$. ###### Lemma 2. Let $`\varphi `$ be an affine deformation of $`G`$ corresponding to the cocycle $`uZ^1(G,^{2,1})`$ and let $`gG`$. Suppose that $`\mu (t)`$ is a path in $`𝔗(M)`$ tangent to $`u`$. Then (7) $$\alpha _\varphi (g)=L_g^{}(0).$$ Furthermore $`\alpha _\varphi (g)`$ and $`\tau _g^{}(0)`$ have the same sign. ###### Proof. Let $`\stackrel{~}{\iota }_t:G\mathrm{SL}(2,)`$ be a smooth path of representations starting at the inclusion $`\iota `$ corresponding to $`\mu (t)`$. $`\tau _g^{}(0)`$ $`={\displaystyle \frac{d}{dt}}\|_{t=0}\|\mathrm{tr}\stackrel{~}{\iota }_t(g)\|`$ $`=\pm {\displaystyle \frac{d}{dt}}\|_{t=0}\mathrm{tr}\left(g(\mathrm{exp}(tu(g)+O(t^2)))\right)`$ $`=\pm {\displaystyle \frac{d}{dt}}\|_{t=0}\mathrm{tr}\left(g(𝕀+tu(g)+O(t^2))\right)`$ $`=\pm \mathrm{tr}\left(gu(g)\right)`$ where the sign equals $`\mathrm{sgn}(\mathrm{tr}(\stackrel{~}{\iota }_t(g)))=\mathrm{sgn}(\mathrm{tr}(\stackrel{~}{\iota }_0(g)))`$. Applying (5) to the last expression gives (8) $$\tau _g^{}(0)=\frac{\sqrt{\mathrm{tr}(g)^24}}{2}\alpha (g).$$ Thus $`\tau _g^{}(0)`$ has the same sign as $`\alpha (g)`$ as claimed. To prove (7), apply (6) and the chain rule to obtain: (9) $$\tau _g^{}(0)=\mathrm{sinh}\left(\frac{L_g(0)}{2}\right)L_g^{}(0).$$ Since $$\mathrm{sinh}\left(\frac{L_g(0)}{2}\right)=\frac{\sqrt{\mathrm{tr}(g)^24}}{2},$$ (7) follows from (8) and (9). ∎ Thus a cocycle is positive (respectively negative) in the sense of Theorem 1 if and only if the corresponding deformation in $`𝔗(M)`$ increases (respectively decreases) lengths of closed curves, to first order. ## 7. Reduction to Teichmüller theory Suppose that $`G\mathrm{SL}(2,)`$ and $`\varphi :G\mathrm{Isom}(^{2,1})`$ is a proper affine deformation. By Theorem 1, the corresponding cocycle $`uZ^1(G,^{2,1})`$ is either positive or negative; by replacing $`u`$ by $`u`$ if necessary, we assume that $`u`$ is positive. By Fried-Goldman , $`G`$ is necessarily discrete and is isomorphic to its image in the group of affine isometries. Suppose that $`G`$ is cocompact. By passing to a subgroup of finite index, we may assume that $`G`$ is torsionfree. Then $`G`$ acts freely on the real hyperbolic plane $`\mathrm{H}_{}^2`$ and since $`G`$ is discrete and cocompact, $`\mathrm{H}_{}^2/G`$ is a closed hyperbolic surface $`M`$. Furthermore $`G`$ is isomorphic to the fundamental group $`\pi _1(M)`$. The representation $`\stackrel{~}{\iota }`$ corresponds to a point $`O`$ in the Teichmüller space $`𝔗(M)`$ and the cohomology class $`[u]H^1(G,^{2,1})`$ corresponds to a tangent vector $`\upsilon `$ to $`𝔗(M)`$ at $`O`$. ###### Lemma 3. There exists a path $`\mu (t)`$ in $`𝔗(M)`$, defined for all $`0t<\mathrm{}`$ starting at $`O𝔗(M)`$ with tangent vector $`\upsilon T_O𝔗(M)`$: (10) $`\mu (0)`$ $`=O`$ $`\mu ^{}(0)`$ $`=\upsilon `$ such that, for each $`gG`$, the geodesic length function $`\mathrm{}_g`$ is convex along $`\mu (t)`$. Assuming Lemma 3 and that $`u`$ is positive, we obtain a contradiction. Since $`\alpha (g)>0`$, the directional derivative $$\mu ^{}(0)\mathrm{}_g=\upsilon \mathrm{}_g=L_g^{}(0)>0$$ by Lemma 2. Convexity implies that $`\mu ^{}(t)\mathrm{}_g`$ cannot decrease as $`t+\mathrm{}`$. Thus $$(\mathrm{}_g\mu )^{}(t)=\mu ^{}(t)\mathrm{}_g\mu ^{}(0)\mathrm{}_g=\alpha (g)>0$$ for all $`t0`$. In particular $`\mathrm{}_g\mu `$ is monotone. Furthermore (11) $$\mathrm{}_g(\mu (t))+\mathrm{}\text{ as }t+\mathrm{},$$ that is, each closed geodesic on the hyperbolic surface $`\mu _t`$ lengthens as $`t+\mathrm{}`$. Such a path $`\mu `$ cannot exist for closed hyperbolic surfaces. Let $`N>0`$. Then for only finitely many conjugacy classes $`F=\{[g_1],\mathrm{},[g_m]\}`$ in $`G\pi _1(M)`$, the corresponding closed geodesics in $`M`$ have length $`<N`$. (Here $`[g]`$ denotes the conjugacy class of $`gG`$.) For any $`gG`$ with $`[g]F`$, the length function $`L_g(t)>L_g(0)N`$. Now consider $`[g_i]F`$. Let $$\alpha _0=\underset{1im}{\mathrm{min}}\alpha (g_i)>0.$$ Convexity, together with (7) implies that $$L_{g_i}(t)L_{g_i}(0)+t\alpha (g_i)t\alpha _0.$$ Hence, for $`t>N/\alpha _0`$, $$L_g(t)=\mathrm{}_g(\mu _t)>N$$ for all $`gG\{𝕀\}`$. However, for any closed hyperbolic surface $`M`$ there exists a simple closed geodesic of length at most $`2\mathrm{log}(22\chi (M))`$ (Lemma 5.2.1 of Buser ). Taking $`N>2\mathrm{log}(22\chi (M))`$, we obtain the desired contradiction. ∎ ###### Proof of Lemma 3. Here are two constructions for $`\mu `$, the first based on the Riemannian geometry of $`𝔗(M)`$ with the Weil-Petersson metric and the second based on Thurston’s earthquake flows. Let $`\mu (t)`$ be the Weil-Petersson geodesic satisfying (10). By Corollary 4.7 of Wolpert , the geodesic length function $`\mathrm{}_g`$ is strictly convex along $`\mu (t)`$ and the directional derivative $`\upsilon \mathrm{}_g>0`$, for any $`gG\{1\}`$. Therefore $`\mathrm{}_g\mu (t)`$ is monotonically increasing for $`t>0`$. However, in general the Weil-Petersson metric is geodesically incomplete (Chu , Wolpert ), so that $`\mu (t)`$ is only defined for $`t_1<t<t_2`$ where $`t_1<0<t_2`$. We show this is impossible under our assumptions on $`\mu ^{}(0)=\upsilon `$. By Mumford’s compactness theorem (Mumford , Harvey , 2.5.1 or Buser , 6.6.5), the subspace of moduli space consisting of hyperbolic surfaces whose injectivity radius is larger than any positive constant is compact. An incomplete geodesic on a Riemannian manifold must leave every compact set. Therefore, if the Weil-Petersson geodesic $`\mu (t)`$ cannot be extended to $`t_2<\mathrm{}`$, then $$\underset{tt_2}{lim}\underset{gG\{𝕀\}}{inf}\mathrm{}_g(\mu (t))=0,$$ contradicting monotonicity of $`\mathrm{}_g`$. Hence $`\mu (t)`$ is defined for all $`t<\mathrm{}`$. As above, convexity implies (11). Alternatively, take $`\mu `$ to be the earthquake path introduced by Thurston (see Kerckhoff and Thurston ). For the given tangent vector $`\upsilon `$, there exists a unique measured geodesic lamination $`\lambda `$ such that the corresponding earthquake path $`\mu (t)=_\lambda (t)`$ satisfies (10) (Kerckhoff , Proposition 2.6). By Kerckhoff (see also Wolpert ), each length function $`\mathrm{}_g`$ is convex along the earthquake path $`_\lambda `$, implying (11). Indeed, $`\mathrm{}_g`$ is strictly convex along $`\mu `$ since the lamination $`\lambda `$ fills up $`M`$ — that is, every nonperipheral simple closed curve $`\sigma `$ intersects $`\lambda `$. For otherwise $`\mathrm{}_\sigma `$ would be constant along $`\mu `$, contradicting $$\frac{d}{dt}\|_{t=0}\mathrm{}_\sigma \mu (t)>0.$$ ###### Remark. Another proof, closer in spirit to the proof in , involves the density of simple closed curves in the projective measured lamination space. Let $`𝒮`$ denote the set of isotopy classes of simple closed curves on $`M`$ and let $`𝒫(M)`$ denote Thurston’s space of projective equivalence classes of measured geodesic laminations on $`M`$. Since $`(M)`$ $`T_O𝔗(M)`$ $`\lambda `$ $`_\lambda ^{}(0)`$ is a homeomorphism (Proposition 2.6 of ), there exist $`\lambda (M)`$ satisfying $`_\lambda ^{}(0)=\upsilon 0`$. Theorem 5.1 of implies $`𝒫(M)`$ $`T_O^{}𝔗(M)`$ $`[\lambda ]`$ $`d\mathrm{log}\mathrm{}_\lambda `$ is an embedding onto a convex sphere in $`T_O^{}𝔗(M)`$ (where $`\mathrm{}_\lambda (N)`$ denotes the length of the lamination $`\lambda `$ as measured in $`N`$). Since $`𝒮`$ is dense in $`𝒫(M)`$, there exist $`\gamma _1,\gamma _2𝒮`$ such that $`(d\mathrm{log}\mathrm{}_{\gamma _1})(\lambda )`$ $`>0`$ $`(d\mathrm{log}\mathrm{}_{\gamma _2})(\lambda )`$ $`<0.`$ Let $`g_1,g_2\pi _1(M)`$ correspond to $`\gamma _1,\gamma _2`$ respectively. Then $$\upsilon L_{g_1}>0,\upsilon L_{g_2}<0,$$ contradicting Theorem 1 and Lemma 2. ###### Remark. Mess’s original proof uses Lorentzian geometry, and in particular the theory of domains of dependence in constant curvature Lorentzian space forms developed in and Scannell . As part of his general theory, Mess shows that any affine deformation sufficiently near the holonomy of a complete flat Lorentz 3-manifold is the holonomy of a complete flat Lorentz 3-manifold, that is, the nearby action is also proper and free. The cocycle $`u`$ corresponds to the velocity vector to an earthquake path $`_\lambda `$ along a measured geodesic lamination $`\lambda `$, and $`\lambda `$ is approximated by a finite measured geodesic lamination, that is, a disjoint union of simple closed geodesics. However for a finite lamination, the corresponding group action is not free (elements of $`G`$ corresponding to curves disjoint from $`\lambda `$ have fixed points), a contradiction.
warning/0005/cond-mat0005392.html	ar5iv	text	# Systematic study of Ga1-xInxAs self-assembled quantum wires with different interfacial strain relaxation ## I Introduction Optical properties of III-V semiconductor nanostructures have attracted a great deal of interest recently for their applications in optical communications that involve switching, amplification, and signal processings. $`Ga_xIn_{1x}As`$ and $`Ga_xIn_{1x}P`$ are among the most important ternary III-V compound semiconductors. The band gap of Ga<sub>x</sub>In<sub>1-x</sub>As covers both the 1.3 and 1.55 $`\mu m`$ wavelengths, which are preferred in long distance fiber communications. However, long-wavelength photonic devices based on lattice-matched Ga<sub>0.47</sub>In<sub>0,53</sub>As/InP heterostructures suffer from strong Auger recombination and inter-valence band absorption processes\[1-3\]. Recently, to improve the performance of long-wavelength semiconductor lasers, long wavelength($`1.55\mu `$m) Ga<sub>x</sub>In<sub>1-x</sub>As quantum-wire (QWR) lasers have been grown by a single step molecular beam epitaxy technique\[4-6\]. An important optical property is the change in the emission of light (in energy, polarization, and intensity) that results from phase-space filling of carriers in one- and two-dimensionally (2D) confined systems. As the dimensionality of the quantum confinement increases from 1D to 2D, the sharpening of the density of states gives rise to a lower excitation threshold, thereby yielding potentially enhanced optical effects. It is found that the QWR laser structure is a promising candidate for optical communication device because of the many predicted benefits, such as higher gain, reduced temperature sensitivity, higher modulation bandwidths, and narrower spectral linewidths. The fabrication of quantum wires via the strain-induced lateral-layer ordering (SILO) process starts with the growth of short-period superlattices (SPS) \[e.g. (GaAs)<sub>2</sub>/(InAs)<sub>2.25</sub>\] along the direction. The excess fractional InAs layer leads to stripe-like islands during the MBE growth. The presence of stripes combined with the strain effect lead to a natural phase separation as additional layers of GaAs or InAs are deposited. A self-assembled quantum wire heterostructure can then be created by sandwiching the composition modulated layer between barrier materials such as Al<sub>0.24</sub>Ga<sub>0.24</sub>In<sub>0.52</sub>As (quarternary), Al<sub>0.48</sub>In<sub>0.52</sub>As(ternary), or InP (binary)\[4-6\]. It was found that different barrier materials can lead to different degree of lateral composition modulation, and the consequent optical properties are different. Besides the self-assembled lateral ordering, it is believed that the strain also plays a key role in the temperature stability and optical anisotropy for the QWR laser structure. Much work has been undertaken which predicts that by using strained-layer superlattice to form the active region of a quantum-wire laser, the threshold current can be decreased by one order of magnitude, and the optical loss due to intervalence-band absorption and Auger recombination will also be greatly reduced\[4-6\]. Also, the temperature sensitivity is reduced by an order of magnitude compared with strain-free structures. Temperature sensitivity of the lasing wavelength for typical GaInAs/InP lasers is around $`5\AA /{}_{}{}^{0}C`$. By using a distributed-feedback structure, the temperature dependence of the lasing wavelength can be reduced to $`1\AA /{}_{}{}^{0}C`$. With the self-assembled GaInAs QWR, the temperature sensitivity is smaller than $`0.1\AA /{}_{}{}^{0}C`$. Thus, it represents an important improvement on the current technology in fabricating long wavelength lasers for fiber communication. A few theoretical studies of Ga<sub>x</sub>In<sub>1-x</sub>As self-assembled QWRs based on a simplified uniform strain model have been reported.\[8-10\] In these calculations the SPS region is modeled by a Ga<sub>x</sub>In<sub>1-x</sub>As alloy with a lateral modulation of the composition $`x`$. Although these calculations can explain the QWR band gap and optical anisotropy qualitatively, it does not take into account the detailed SPS structure and the microscopic strain distribution. The understanding of these effects is important if one wish to have a full design capability of the self-assembled QWR optoelectronic devices. On the other hand, microscopic strain distributions in several SPS structures have been calculated, and the electronic properties of these SPS structures have been studied via an empirical pseudopotential method. However, the effects of sandwiching the SPS structures between barrier materials so as to form QWRs have not been explored. It is also worth pointing out that the effective masses obtained in the pseudopotential method for the constituent bulk materials are 0.032$`m_0`$ and 0.092$`m_0`$ for InAs and GaAs, respectively. These values are 30% larger than the actual values (0.024$`m_0`$ and 0.067$`m_0`$); thus, the energy levels of quantum confined states obtained by this method are subject to similar uncertainty. In this paper we present a systematic study of the effects of microscopic strain distribution on the electronic and optical properties of the Ga<sub>1-x</sub>In<sub>x</sub>As self-assembled QWRs via the combination of the effective bond-orbital model (EBOM) for electronic states and the valence-force-field (VFF) model for microscopic strain distribution. Both clamped and unclamped SPS structures with different degrees of lateral alloy mixing are considered. The clamped structure has a SPS layer sandwiched between the substrate and a thick capping layer so that the interface between the top monolayer of the SPS and the barrier region is atomically flat. The unclamped structures correspond to a SPS layer sandwiched between the substrate and a thin capping layer, which allows more flexible relaxation of interface atoms so the strain energy can be relieved. This leads to a wavy interface structure. Most self-assembled QWR samples reported to date are closer to unclamped structures with a wavy interface, although it should be possible to produce self-assembled QWRs closer to the clamped structures by using a thick capping layer. Our theoretical studies show that there are profound differences in the electronic and optical properties between the clamped and unclamped self-assembled QWR structures. In particular, we found that in clamped structures the electron and hole are confined in the Ga-rich region and the polarization of Photoluminescence (PL) can go from predominantly along for SPS with abrupt change in In/Ga composition to predominantly along \[$`1\overline{1}0`$\] for SPS with smooth change in In/Ga composition. On the other hand, for unclamped structures the electron and hole are confined in the In-rich region, and the optical polarization is always predominately along \[$`1\overline{1}0`$\] with a weak dependence on the lateral alloy mixing. This also implies that by changing the degree of strain relaxation at the interface between the SPS and the capping layer, one can tailor the optical properties of self-assembled QWRs. Thus, through this theoretical study, we can gain a better understanding of the strain engineering of self-assembled QWR structures which may find applications in fiber-optical communication. ## II Model structures For both clamped and unclamped SPS structures, we consider different degrees of lateral alloy mixing and examine the effects of the microscopic strain distribution on the electronic and optical properties. Two example QWR model structures considered in the present paper (prior to alloy mixing) are depicted in Figure 1. The supercell of the first model structure consists of 8 pairs of (001) (GaAs)<sub>2</sub>/(InAs)<sub>2</sub> SPS with a total thickness of $`94\AA `$ (quantum well region) followed by a Al<sub>0.24</sub>Ga<sub>0.24</sub>In<sub>0.52</sub>As layer (barrier region) with thickness $`60\AA `$ (20 diatomic layers). The supercell of the second model structure consists of 8 pairs of (001) (GaAs)<sub>2</sub>/(InAs)<sub>2.25</sub> SPS with a total thickness of $`100\AA `$ (quantum well region) followed by a Al<sub>0.24</sub>Ga<sub>0.24</sub>In<sub>0.52</sub>As layer ( barrier region). We assume an arrangement of alternating stripe like islands due to strain induced lateral ordering. In the diagram, no lateral alloy mixing is shown. In our calculations of strain distribution and electronic structures, different degrees of lateral alloy mixing will be considered. Experimentally, self-assembled InGaAs QWRs were usually grown with the (2/2.25) SPS structure. However, the migration of excess In atoms during MBE growth could lead a structure somewhere between (2/2) SPS and (2/2.25) SPS structures. In both model structures, we can divide the self-assembled QWR into two regions with the left half being Ga rich and the right half In rich. During growth, varying degrees of lateral alloy mixing of these islands with the surround atoms is likely to occur. In the atomic layers with In-rich alloy filled the right half of the unit cell (layer 2 in structure 1 and layers 7 and 9 in structure 2), the In composition $`x_{In}`$ is assumed to vary in the direction (or $`y^{}`$ direction) according to the relation $$x_{In}=\{\begin{array}{cc}x_m[1\mathrm{sin}(\pi y^{}/2b)]/2\hfill & \text{ for }y^{}<b\hfill \\ 0\hfill & \text{ for }b<y^{}<L/2b\hfill \\ x_m\{1+\mathrm{sin}[\pi (y^{}L/2)/2b]\}/2\hfill & \text{ for }L/2b<y^{}<L/2+b\hfill \\ x_m\hfill & \text{ for }L/2+b<y^{}<Lb\hfill \\ x_m\{1\mathrm{sin}[\pi (y^{}L)/2b]\}/2\hfill & \text{ for }y^{}>Lb,\hfill \end{array}$$ (1) where $`x_m`$ is the maximum In composition in the layer, $`2b`$ denotes the width of lateral composition grading, and $`L`$ is the period of the lateral modulation in the direction. In layer 4 of structure 1, which contains Ga-rich alloy in the left half of the unit cell, the Ga composition ($`x_{Ga}`$) as a function of $`y^{}`$ is given by a similar equation with the sign of the sine function reversed as compared to Eq. (1). In structure 2, there are a few atomic layers that contain 0.25 monolayer of In (layers 3 and 13) or Ga( layers 5 and 11). The lateral alloy modulation in layers 3 and 13 is described by $$x_{In}=\{\begin{array}{cc}0\hfill & \text{ for }0<y^{}<5L/8b\hfill \\ x_m\{1+\mathrm{sin}[\pi (y^{}5L/8)/2b]\}/2\hfill & \text{ for }5L/8b<y^{}<5L/8+b,\hfill \\ x_m\hfill & \text{ for }5L/8+b<y^{}<7L/8b\hfill \\ x_m\{1\mathrm{sin}[\pi (y^{}7L/8)/2b]\}/2\hfill & \text{ for }7L/8b<y^{}<7L/8+b\hfill \\ 0\hfill & \text{ for }7L/8+b<y^{}<L.\hfill \end{array}$$ (2) Similar equation for $`x_{Ga}`$ in layers 5 and 11 can be deduced from the above. By varying the parameters $`x_m`$ and $`b`$, we can get different degrees of lateral alloy mixing. Typically $`x_m`$ is between 0.6 and 1, and $`b`$ is between zero and $`15a_{[110]}62\AA `$. A valence force field (VFF) model\[13-15\] is used to find the equilibrium atomic positions in the self-assembled QWR structure by minimizing the lattice energy. The strain tensor at each atomic (In or Ga) site is then obtained by calculating the local distortion of chemical bonds. We found that different local arrangement of atoms can lead to very different strain distribution. In particular, the shear strain in the clamped structure can change substantially when the In/Ga composition modulation is changed. Consequently, the optical anisotropy can be reversed due to the change in the strength of the shear strain caused by the intermixing of In and Ga atoms. ## III Theoretical Approach The method used in this paper for calculating the strained QWR band structure is based on the effective bond-orbital model (EBOM). A detailed description of this method has been published elsewhere. EBOM is a tight binding-like model with minimum set of localized basis (the bonding or anti-bonding orbitals). The interaction and optical parameters are obtained by a correspondence with the $`𝐤𝐩`$ theory, which can be cast into analytic forms. Thus, the model can be viewed as a spatially discretized version of the $`𝐤𝐩`$ method, while retaining the virtues of LCAO (linear combination of atomic orbitals) method. The $`𝐤𝐩`$ model is the most popular one for treating electronic structures of semiconductor quantum wells or superlattices. However, when applied to complex structures such as self-assembled quantum wires\[4-9\] or quantum dots, the method becomes very cumbersome if one wish to implement the correct boundary conditions that take into account the differences in $`𝐤𝐩`$ band parameters for different materials involved. EBOM is free of this problem, since different material parameters are used at different atomic sites in a natural way. For simple structures, when both EBOM and $`𝐤𝐩`$ model are equally applicable, the results obtained are essentially identical. The optical matrix elements for the QWR states are computed in terms of elementary optical matrix elements between the valence-band bond orbitals and the conduction-band orbitals. The present calculation includes the coupling of the four spin-3/2 valence bands and the two spin-1/2 conduction bands closest to the band edges. Thus, it is equivalent to a 6-band $`𝐤𝐩`$ model. For our systems studied here, the band-edge properties are relatively unaffected by the split-off band due to the large spin-orbit splitting as discussed in our previous paper. Hence, the split-off bands are ignored here. The bond-orbitals for the GaAs and InAs needed in the expansion of the superlattice states contain the following: four valence-band bond orbitals per bulk unit cell, which are p-like orbitals coupled with the spin to form orbitals with total angular momentum J=3/2 plus two conduction-band bond orbitals with J=1/2. They are written as $$\|𝐑,u_{JM}>=\underset{\alpha ,\sigma }{}C(\alpha ,\sigma ,J,M)\|\stackrel{}{R},\alpha >\psi _\sigma ,$$ (3) where $`J=1/2`$ and 3/2 for the conduction and valence bands, respectively, and $`M=J,\mathrm{}J`$, $`\psi _\sigma `$, designates the electron spinor($`\sigma `$=1/2,-1/2), and $`\|𝐑,\alpha >`$ denote an $`\alpha `$-like($`\alpha =s,x,y,z`$) bond orbital, located at unit cell $`𝐑`$. $`C(\alpha ,\sigma ,J,M)`$ are the coupling coefficients obtainable by group theory. All these bond orbitals are assumed to be sufficiently localized so that the interaction between orbitals separated farther than the nearest-neighbor distance can be ignored. The effect of strain is included by adding a strain Hamiltonian $`H^{st}`$ to the EBOM Hamiltonian. The matrix elements of $`H^{st}`$ in the bond-orbital basis can be obtained by the deformation-potential theory of Bir and Pikus. We use the valence-force field (VFF) model of Keating and Martin to calculate the microscopic strain distribution. This model has been shown to be successful in fitting and predicting the elastic constants of elastic continuum theory, calculating strain distribution in a quantum well, and determining the atomic structure of III-V alloys. It was also used in the calculation of the strain distribution in self-assembled quantum dots. The VFF model is a microscopic theory which includes bond stretching and bond bending, and avoids the potential failure of elastic continuum theory in the atomically thin limit. The total energy of the lattice is taken as $$V=\frac{1}{4}\underset{ij}{}\frac{3}{4}\alpha _{ij}(d_{ij}^2d_{0.ij}^2)^2/d_{0,ij}^2+\frac{1}{4}\underset{i}{}\underset{jk}{}\frac{3}{4}\beta _{ijk}(\stackrel{}{d_{ij}}.\stackrel{}{d_{ik}}+d_{0,ij}d_{0,ik}/3)^2/d_{0,ij}d_{0,ik}$$ (4) where $`i`$ runs over all the atomic sites, $`j,k`$ run over the nearest-neighbor sites of $`i`$, $`\stackrel{}{d_{ij}}`$ is the vector joining the sites i and j, $`d_{ij}`$ is the length of the bond, $`d_{0,ij}`$ is the corresponding equilibrium length in the binary constituents, and $`\alpha _{ij}`$ and $`\beta _{ijk}`$ are the bond stretching and bond bending constants, respectively. $`\alpha `$ and $`\beta `$ are from Martin’s calculations. For the bond-bending parameter $`\beta `$ of In-As-Ga, we take $`\beta _{ijk}=\sqrt{\beta [ij]\beta [ik]}`$ following Ref. . To find the strain tensor in the $`InAs/GaAs`$ self-assembled QWR, we start from ideal atomic positions and minimize the system energy using the Hamiltonian given above. Minimization of the total energy requires one to solve a set of coupled equations with 3N variables, where N is the total number of atoms. Direct solution of these equations is impractical in our case, since the system contains more than 6,000 atoms. We use an approach taken by several authors which has been shown to be quite efficient. In the beginning of the simulation all the atoms are placed on the InP lattice, we allow atoms to deviate from this starting positions and use periodic boundary conditions in the plane perpendicular to the growth direction, while keeping atoms in the planes outside the SPS region at their ideal atomic positions for a InP lattice (since the self-assembled QWR is grown epitaxially on the InP substrate). In each iteration, only one atomic position is displaced and other atom positions are held fixed. The direction of the displacement of atom $`i`$ is determined according to the force $`f_i=V/x_i`$ acting on it. All atoms are displaced in sequence. the whole sequence is repeated until the forces acting on all atoms become zero at which point the system energy is a minimum. Once the positions of all the atoms are known, the strain distribution is obtained through the strain tensor calculated according to the method described in Ref. . Let $`R^0`$ be the position matrix without strain: $`R^0=\left(\begin{array}{ccc}R_{12x}^0& R_{23x}^0& R_{34x}^0\\ R_{12y}^0& R_{23y}^0& R_{34y}^0\\ R_{12z}^0& R_{23z}^0& R_{34z}^0\end{array}\right),`$ where $`𝐑_{ij}^0=𝐑_j^0𝐑_i^0;i,j=1,4`$ and $`𝐑_i^0`$ $`(i=1,4)`$ denote positions of four As atoms surrounding a Ga or In atom. Here we choose $`𝐑_{12}^0=(1,1,0)a/2`$, $`𝐑_{23}^0=(1,0,1)a/2`$, and $`𝐑_{34}^0=(1,1,0)a/2`$. $`a`$ is the lattice constant of GaAs or InAs, depending on the site. Let $`R`$ be the corresponding position matrix with strain. It was shown in Ref. that the strain tensor is given by $$ϵ=RR_0^11.$$ (5) After getting the strain tensor, the strain Hamiltonian is given by Bir and Pikus $$H^{st}=\left(\begin{array}{ccc}\mathrm{\Delta }V_H+D_1& \sqrt{3}de_{xy}& \sqrt{3}de_{xz}\\ \sqrt{3}de_{xy}& \mathrm{\Delta }V_H+D_2& \sqrt{3}de_{yz}\\ \sqrt{3}de_{xz}& \sqrt{3}de_{yz}& \mathrm{\Delta }V_H+D_3\end{array}\right),$$ (6) where $`e_{ij}=(ϵ_{ij}+ϵ_{ji})/2`$, and $$\mathrm{\Delta }V_H=(a_1+a_2)(ϵ_{xx}+ϵ_{yy}+ϵ_{zz}),D_1=b(2ϵ_{xx}ϵ_{yy}ϵ_{zz}),D_2=b(2ϵ_{yy}ϵ_{xx}ϵ_{zz}),D_3=b(2ϵ_{zz}ϵ_{yy}ϵ_{xx}),$$ The strain potential on the s states is given by $$\mathrm{\Delta }V_c=c_1(ϵ_{xx}+ϵ_{yy}+ϵ_{zz}),$$ The strain Hamiltonian in the bond-orbital basis $`\|JM>`$ can be easily found by using the coupling constants, i.e, $$<JM\|H^{st}\|J^{}M^{}>=\underset{\alpha ,\alpha ^{},\sigma }{}C(\alpha ,\sigma ;J,M)^{}C(\alpha ^{},\sigma ;J^{},M^{})H_{\alpha \alpha ^{}}^{st}$$ (7) The elastic constants $`C_{12}`$ and $`C_{11}`$ for GaAs, InAs and AlAs can be found in Ref. . The deformation potentials $`a_1,a_2,b,c_1,d`$ can be found in Ref.\[23-26\]. The linear interpolation and virtual crystal approximations used to obtain the corresponding parameters for the barrier material (Al<sub>0.24</sub>Ga<sub>0.24</sub>In<sub>0.52</sub>As). The above strain Hamiltonian is derived locally for the each cation atom in the self-assembled QWR considered. To calculate the electronic states of the self-assembled QWR for model structure 1\[see Fig. 1(a)\], we first construct a zeroth-order Hamiltonian for a superlattice structure which contains in each period 8 pairs of (001) (GaAs)<sub>2</sub>/(InAs)<sub>2</sub> SPS layers (with a total thickness around 100 Å) and 20 diatomic layers of Al<sub>0.24</sub>Ga<sub>0.24</sub>In<sub>0.52</sub>As (with thickness around 60 Å). So, the superlattice unit cell for the zeroth-order model contains 52 diatomic layers. The appropriate strain Hamiltonian for the the (GaAs)<sub>2</sub>/(InAs)<sub>2</sub> SPS on InP is also included. For model structure 2 \[see Fig. 1(b)\], the zero-order superlattice consistes of two additional monolayers of InAs inserted into the 8 pairs of (2/2) SPS: one between the 2nd and 3rd pair of (2/2) SPS, the other between the 5th and 6th pair of (2/2) SPS. The eigen-states for the zero-th order Hamiltonian for different values of $`k_2`$ (separated by the SL reciprocal lattice vectors in the direction) are then used as the basis for calculating the self-assembled QWR electronic states. The difference in the Hamiltonian (including strain effects) caused by the intermixing of Ga and In atoms at the interfaces is then added to the zeroth-order Hamiltonian, and the electronic states of the full Hamiltonian is solved by diagonalizing the Hamiltonian matrix defined within a truncated set of eigen-states of the zeroth-order Hamiltonian. A total of $`500`$ eigen-states of the zero-th order Hamiltonian (with 21 different $`k_{110}`$ points) were used in the expansion. The subbands closest to the band edge are converged to within 0.1 meV. ## IV Results and discussions A. Strain distributions In this section we discuss strain distributions in our model structures as described in section I with varying degree of lateral alloy mixing. Both clamped (with thick capping layer) and unclamped (with thin capping layer) situations are considered. To model the alloy structure with composition modulation, we use a super-cell which contains 72 atoms in the ($`y^{}`$) direction, 36 atoms in the \[$`1\overline{1}0`$\] ($`x^{}`$) direction and 64 or 68 atomic planes along the ($`z`$) direction \[8 pairs of (2/2) or (2/2.25) SPS\] plus a GaAs capping layer. In the atomic planes which consist of alloy structure, we first determine the In composition at a given $`y^{}`$ according to Eq. (1) and then use a random number generator to determine the atomic species along the $`x^{}`$ direction. The bottom layer of atoms are bonded to the InP substrate with the substrate atoms fixed at their ideal atomic positions. The calculated strain distributions are then averaged over the $`x^{}`$ coordinate. For the clamped case, the GaAs capping layer is assumed to be lattice matched to InP with a flat surface. For the unclamped case, the capping layer is allowed to relax freely, thus giving rise to a wavy surface structure. The diagonal strains in the four atomic layers that constitute the (2/2) SPS in structure 1 for the clamped and unclamped cases are shown in Figs. 2 and 3, respectively. The lateral alloy modulation considered is described by Eq. (1) with $`x_m=0.8`$ and $`b=7a_{[110]}`$. For best illustration, we show diagonal strains in a rotated frame, in which $`x^{}`$ is \[1$`\overline{1}`$0\], $`y^{}`$ is , and $`z^{}`$ is . The layer number in the figure labels the atomic layers in Fig. 1(a), starting from the bottom layer as layer 1. There are two main features of the microscopic strain distribution. First, in the ideal situation (without atomic relaxation), one would predict $`ϵ_{y^{}y^{}}`$ to be the same as $`ϵ_{x^{}x^{}}`$ due to symmetry. However, with atomic relaxation, all Ga(In) atoms tend to shift in a direction so as to reduce the strain in the Ga(In)-rich region. Thus, the magnitude of $`ϵ_{y^{}y^{}}`$ on Ga(In) sites (dashed lines) is lower(higher) than $`ϵ_{x^{}x^{}}`$ in Ga(In)-rich region. Second, the $`z`$ component strain (dash-dotted lines) tend to compensate the other two components such that the volume of each bulk unit cell is closer to that for the unstrained bulk. Thus, we see that $`ϵ_{z^{}z^{}}`$ has an opposite sign compared to $`ϵ_{x^{}x^{}}`$ or $`ϵ_{y^{}y^{}}`$ at all atomic sites. Comparing Fig. 2 with Fig. 3, we see that the main difference between the clamped and unclamped case is that in the atomic layers with lateral composition modulation (layers 2 and 4), the hydrostatic strain (sum of all three diagonal strain components) is much smaller in the unclamped structure than in the clamped structure. This would lead to a major difference in the band-edge profile for the two structures to be discussed below. The difference in $`ϵ_{x^{}x^{}}`$ and $`ϵ_{y^{}y^{}}`$ leads to a nonzero shear strain $`e_{xy}=(ϵ_{y^{}y^{}}ϵ_{x^{}x^{}})/2`$ in the original coordinates. For the clamped case, the shear strain is found to be particularly strong near the boundary where the alloy composition begins to change, and it is sensitive to the degree of lateral alloy mixing. For $`x_m=0.8`$ and $`b=7a_{[110]}`$, the maximum value of $`ϵ_{xy}`$ is around 0.4% For abrupt composition modulation ($`x_m=1`$ and $`b=0`$) the maximum $`ϵ_{xy}`$ value increases five-fold to around 2%. The other shear strain components ($`ϵ_{xz}`$ and $`ϵ_{yz}`$) are found to have similar magnitude. For the unclamped case, the shear strain is strong even in regions away from the boundary, and the maximum shear strain is larger than its counterpart in the clamped structure by about 30%. The diagonal strains of structure 2 for the unclamped case with $`x_m=1`$ and $`b=7a_{[110]}`$ are shown in Fig. 4. The layer number in the figure labels the atomic layers in Fig. 1(b), starting from the bottom layer as layer 1. Only four representative atomic layers (3,4,5, and 7) are shown. We note that the average magnitude of the hydrostatic strain in the In-rich region in model structure 2 is comparable to that for unclamped structure 1. B. Electronic structures In order to understand the aspect of lateral quantum confinement due to composition modulation and strain, we examine the band-edge energies of a strained quantum well structure whose well material is the same structure as appeared in the SILO QWR with a fixed value of $`y^{}`$ and the barrier is 60Åthick Al<sub>0.24</sub>Ga<sub>0.24</sub>In<sub>0.52</sub>As layer. For structure 1 depicted in Fig. 1(a), the well material consists of 8 pairs of (GaAs)<sub>2</sub>/(InAs)<sub>2</sub> SPS. The conduction band minimum and valence band maximum of the above quantum well with different degrees of lateral alloy mixing as functions of $`y^{}`$ (dashed lines) are shown in Fig. 5 for both clamped and unclamped cases. The strain Hamiltonian used here is the same as the one used in the self-assembled QWR at the corresponding $`y^{}`$. All material parameters are chosen the same as in Ref. for temperature at 77K, except that the deformation potential $`c_1`$ of GaAs is slightly modified from $`6.8eV`$ to $`7.1eV`$ so that $`a_1+a_2+c_1=9.8eV`$ is in agreement with the experimental measurements and the valence-band offset between GaAS and InAs used here is 0.26 eV according to the model solid theory. To correct the band gap of SPS and alloy due to the bowing effect, we add a correction term $`0.4x(1x)eV`$ to the diagonal element for the s-like bond orbital in the Hamiltonian, where $`x`$ is the effective alloy composition of the SPS or alloy at a given $`y^{}`$. With the correction of bowing effect, our model gives band gaps for the (GaAs)<sub>2</sub>(InAs)<sub>2</sub> SPS, Ga<sub>1-x</sub>In<sub>x</sub>As alloy, and the superlattice made of (GaAs)<sub>2</sub>(InAs)<sub>2</sub> SPS grown on InP substrate all in very good agreement with available experimental observations. The photoluminescence (PL) measurements indicates that the (GaAs)<sub>2</sub>(InAs)<sub>2</sub> SPS grown on InP substrate has a gap around 0.75 eV. Here we obtain a band gap of 0.74 eV for the (GaAs)<sub>2</sub>(InAs)<sub>2</sub> SPS (with a strain distribution obtained again via the VFF model) and 0.79 eV for 8 pairs (or 100 Å) of (GaAs)<sub>2</sub>(InAs)<sub>2</sub> SPS sandwiched between Al<sub>0.24</sub>Ga<sub>0.24</sub>In<sub>0.52</sub>As confining barriers. This is also consistent with the PL measurements on the (GaAs)<sub>2</sub>(InAs)<sub>2</sub>/InP multiple quantum wells. The band-edge profiles shown in this figure suggest that for the clamped case with alloy mixing, both electrons and holes are confined in the Ga-rich region with an effective barrier height around $`0.13eV`$ for the electron and $`0.1eV`$ for the hole. Both offsets are large enough to give rise to strong lateral confinement for electrons and holes in the Ga-rich region. Without alloy mixing, the electron is not well confined, since the average energy between the Ga-rich and In-rich regions are nearly the same. For the unclamped case, both electrons and holes are confined in the In-rich region with an effective barrier height around $`0.20.3eV`$ for the electron and $`0.10.2eV`$ for the hole. We note that the band gap in the In-rich region is rather insensitive to the alloy mixing, while the effective barrier height can change substantially due to different alloy mixing. Fig. 6 shows the band-edge profiles for 8 pairs of (2/2.25) SPS \[as illustrated in Fig. 1(b)\] sandwiched between 60ÅAl<sub>0.24</sub>Ga<sub>0.24</sub>In<sub>0.52</sub>As barriers with and without alloy mixing for both clamped and unclamped cases. For the clamped case, the electron is confined in the In-rich region, while the hole is confined in the Ga-rich region. Thus, we have a spatially indirect QWR, which will have very weak inter-band optical transition. For the unclamped case, the band-edge profile is similar to the unclamped case of structure 1 with both the electron and hole confined in the In-rich region. The band gaps for all cases of structure 2 are consistently smaller than their counterparts in structure 1. Fig. 7 shows the near zone-center valence subband structures and squared optical matrix elements ($`P^2`$) for the QWR model structure 1 \[as depicted in Fig. 1(a)\] with thick capping layer (clamped case) and without alloy mixing. Here $`P^2`$ is defined as $`\frac{2}{m_0}_{s,s^{}}\|<\psi _v\|\widehat{ϵ}𝐩\|\psi _c>\|^2`$, where $`\psi _v`$ ($`\psi _c`$) denotes a valance (conduction) subband state. The symbol $`_{s,s^{}}`$ means a sum over two near degenerate pair of subbands in the initial and final states. $`\widehat{ϵ}`$ denotes the polarization of light. Here we consider only the $`x^{}`$ (along the QWR axis) and $`y^{}`$ (perpendicular to the QWR axis) components. Only the dispersion along the $`k_1`$ ($`[1\overline{1}0]`$) direction is shown, since the dispersion along the $`k_2`$ direction for the confined levels is is rather small due to strong lateral confinement. All subbands are two-fold degenerate at the zone center due to the Kramer’s degeneracy and they split at finite wave vectors as a result of lack of inversion symmetry in the system. The first three pairs of subbands are labeled V1, V2, and V3. They have unusually large energy separations compared with other valence subbands. This is because the first three pairs of subbands represent QWR confined states, whereas the other valence subbands (with energies below -110 meV) are unconfined by the lateral composition modulation. To examine the effects caused by the shear strain, we have also calculated the band structures with the shear strain set to zero. We found that without the shear strain the first pairs of valence subbands (V1) have much larger dispersion along the $`k_1`$ direction with an effective masses along $`[1\overline{1}0]`$ about a factor five smaller than that with the shear strain. This indicates that the V1 subband states for the case without shear strain are derived mainly from bond orbitals with $`x^{}`$-character, which leads to stronger overlap between two bond orbitals along the $`x^{}`$ direction, hence larger dispersion along $`k_1`$. When the shear strain is present, the character of bond orbitals in the V1 subband states change from mainly $`x^{}`$-like to predominantly $`y^{}`$-like; thus, the dispersion along $`k_1`$ becomes much weaker. This explains the very flat V1 subbands as shown in Fig. 6. The switching of orbital character in the V1 subbands can be understood as follows. When the shear strain is absent, we have $`ϵ_{x^{}x^{}}=ϵ_{y^{}y^{}}`$, and the confinement effect in the $`y^{}`$ direction pushes down the states with $`y^{}`$ character (which has smaller effective mass in the $`y^{}`$ direction), thus leaving the top valence subband (V1) to have predominantly $`x^{}`$ character. On the other hand, with the presence of shear strain as shown in structure 1, we have $`ϵ_{y^{}y^{}}<ϵ_{x^{}x^{}}`$ in Ga-rich region, and the bi-axial term of the strain potential forces the $`y^{}`$-like states to move above the $`x^{}`$-like states, overcoming the confinement effect. As a result, the V1 subbands become $`y^{}`$-like. The conduction subbands are approximately parabolic as usual with a zone-center subband minimum equal to 688 meV This gives an energy gap of 763 meV for QWR model structure 1. As seen in this figure, the C1-V1 transition for $`y^{}`$-polarization (dashed line) is more than twice that for the $`x^{}`$-polarization (solid line) with $`P_{[1\overline{1}0]}^2/P_{[110]}^20.36`$. This is consistent with the discussion for valence subband structures, where we concluded that the bond orbitals involved in the V1 subbands are predominantly $`y^{}`$-like. Next, we study the case with alloy mixing, in which a gradual lateral modulation in In composition substantially reduces the shear strain. Using a lateral alloy mixing described in Eq. (1) with $`x_m=0.8`$ and $`b=7a_{[110]}`$, we found significant change in band structures and optical properties compared to the case without alloy mixing. With alloy mixing we found that the first pairs of valence subbands (V1) have much larger dispersion than their counterparts in Fig. 6, and the atomic character in the $`V1`$ states change from predominantly $`y^{}`$-like (without alloy mixing) to $`x^{}`$-like (with alloy mixing). The change of atomic character in the V1 states is a consequence of the competition between the QWR confinement (which suppresses the $`y^{}`$ character) and the $`e_{xy}`$ shear strain (which suppresses the $`x^{}`$ character). As a result, the ratio $`P_{[1\overline{1}0]}^2/P_{[110]}^2`$ changes from 0.36 (without alloy mixing) to 3.5 (with alloy mixing), indicating a reversal of optical anisotropy. The above studies lead to the conclusion that for self-assembled QWRs with thick capping layer (clamped case), the optical anisotropy is sensitive to the degree of lateral alloy mixing. Next, we study the properties of self-assembled QWRs with thin capping layer (unclamped case). We found that for the unclamped case, the electron and hole are always confined in the In-rich region, regardless of the degree of alloy mixing and whether the QWR structure consists of (2/2) SPS or (2/2.25) SPS. Furthermore, since the strain in In-rich region has opposite sign compared to that in the Ga-rich region, the effect of shear strain on the VB states confined in In-rich region is also reversed. Namely, it pushes the energy of the $`y^{}`$-like VB states down relative to the $`x^{}`$-like VB states. Thus, both the QWR confinement and the shear strain effects are in favor of having a predominant $`x^{}`$ character in the V1 subband. Therefore, the inter-band optical transition always has a ratio $`P_{[1\overline{1}0]}^2/P_{[110]}^2`$ larger than 1. To illustrate this, we show in Fig. 8 the near zone-center valence subband structures and squared optical matrix elements ($`P^2`$) for the QWR model structure 1 \[as depicted in Fig. 1(a)\] with thin capping layer (unclamped case) and with alloy mixing described by $`x_m=0.8`$ and $`b=7a_{[110]}`$. Comparing the band structure in Fig. 8 with that in Fig. 7, we see that the V1 subband now has much larger dispersion (i.e. much smaller effective mass) than its counterpart in Fig. 6. There are two reasons for this. First, the hole in the unclamped case is now confined in the In-rich region rather than in the Ga-rich region as in the clamped case, thus having much smaller effective mass. Second, both the QWR confinement effect and the shear strain effect lead to a predominantly $`x^{}`$ character in the V1 states, which tend to yield a smaller effective mass in the $`k_1`$ direction. The optical properties shown in Fig. 8 are also distinctly different from that shown in Fig. 7. The inter-band optical transition (C1-V1) now has a much stronger polarization along \[1$`\overline{1}0`$\] than . The quick drop of the the $`y^{}`$ polarization in C1-V1 transition at $`k_10.035\frac{2\pi }{a}`$ is caused by the band mixing of V1 and V2 states as can be seen in the VB band structures in this figure. We have also calculated the near zone-center valence subband structures and squared optical matrix elements ($`P^2`$) for the QWR model structure 2 with thin capping layer (not shown). We found that the band structures and optical matrix elements are similar to Fig. 8, except that the energy separation between V1 and V2 subbands are smaller (by $`5mV`$) and the polarization ratio ($`P_{[1\overline{1}0]}^2/P_{[110]}^2`$) is larger. Finally, we list in Table I the band gaps and the polarization ratio ($`P_{[1\overline{1}0]}^2/P_{[110]}^2`$) for various QWR structures with different degrees of lateral alloy mixing (as indicated by the parameters $`x_m`$ and $`b`$). To calculate the polarization ratio in the photoluminescence (PL) spectra, we integrate the squared optical matrix element over the range of $`k_1`$ corresponding to the spread of exciton envelope function in the $`k_1`$ space. The exciton envelope function is obtained by solving the 1D Schrödinger equation for the exciton in the effective-mass approximation similar to what we did in Ref. 9. For unclamped (2/2) SPS structure, the band gap varies between 0.77 and 0.8 eV and the polarization ratio changes from 3.1 to 1.5 as the degree of alloy mixing is increased. For unclamped (2/2.25) SPS structure, both the band gap (around 0.74 eV) and the polarization ratio (around 2) are insensitive to the degree of alloy mixing. Experimentally, the PL peak for self-assembled InGaAs QWRs with similar dimensions as considered in our model structures are around 735 meV and the polarization ratio for different samples are found to be between 2 and 4.\[4-6\] Taking into account an exciton binding energy between 10 and 20 meV, we find that our theoretical predictions for both unclamped and clamped structures with lateral alloy mixing are consistent with the experimental finding. To sort out which one is closer to the experimental structure would require more detailed comparison between the theory and experiment for optical spectra involving the excited hole states such as the PL excitation spectra. ## V Conclusion We have calculated the band structures and optical matrix elements for the self-assembled GaInAs QWR grown by the SILO method. The actual SPS structure and the microscopic strain distribution has been taken into account. The effects of microscopic strain distribution on the valence subband structures and optical matrix elements are studied for two model self-assembled QWR structures, one with (2/2) SPS structure and the other with (2/2.25) SPS structure, and with varying degrees of lateral alloy mixing. The clamping effect due to thick capping layer is also examined. The valence force field (VFF) model is used to calculate the equilibrium atomic positions in the model QWR structures. This allows the calculation of the strain distribution at the atomistic level. We found that in model structures with thick capping layer (clamped case), the magnitude of shear strain is quite large (around 2%) for abrupt lateral composition modulation and it reduces substantially when the lateral alloy mixing is introduced. We found that the shear strain effect can alter the valence subband structures substantially when the hole is confined in the Ga-rich region, and it gives rise to a reversed optical anisotropy compared with QWR structures with negligible shear strain. This points to a possibility of ”shear strain-engineering” to obtain QWR laser structures of desired optical anisotropy. For model structures with thin capping layer (unclamped case), the effect of shear strain is less dramatic, since the hole is confined in the In-rich region, where both the QWR confinement and the shear strain have similar effect. We found that polarization ratio ($`P_{[1\overline{1}0]}^2/P_{[110]}^2`$) for the unclamped case is always larger than 1 for both (2/2) and (2/2.25) SPS structures, regardless of the degree of lateral alloy mixing. The band gap and polarization ratio obtained for both clamped and unclamped (2/2) SPS structures with lateral alloy mixing ($`x_m0.8`$) and for unclamped (2/2.25) SPS structures are consistent with experimental observations on most self-assembled InGaAs QWRs with similar specifications. For clamped (2/2) SPS structure without alloy mixing, we predicted a reversed optical anisotropy compared with the PL measurements for most InAs/GaAs self-assembled QWRs. However, there exist a few InAs/GaAs self-assembled QWR samples which display the reversed optical anisotropy at 77 K, and recent experimental studies also indicated a reversed optical anisotropy in many samples at 300 K. Our studies demonstrated one possibility for the reversed optical anisotropy, the shear strain effect. There are other possibilities for the reversed optical anisotropy at 300K, such as the thermal population of the excited hole states, which can have different atomic characters as discussed in the previous section. The temperature dependence of the PL peak and polarization in self-assembled InGaAs QWRs remains an intriguing question to be addressed in the future. ACKNOWLEDGEMENTS This work was supported in part by the National Science Foundation (NSF) under Grant No. NSF-ECS96-17153. S. Sun was also supported by a subcontract from the University of Southern California under the MURI program, AFOSR, Contract No. F49620-98-1-0474. We would like to thank K. Y. Cheng and D. E. Wohlert for fruitful discussions and for providing us with the detailed experimental data of the QWR structures considered here. Table I. List of band gaps and polarization ratios for all structures studied \| $`x_m`$ \| b/$`a_{110}`$ \| alloy mixing \| clamped/unclamped \| SPS structure \| gap(eV) \| $`P_{}/P_{}`$ \| \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| \| 1.0 \| 0 \| no \| clamped \| (2/2) \| 0.765 \| 0.36 \| \| 0.8 \| 7 \| yes \| clamped \| (2/2) \| 0.763 \| 3.5 \| \| 1.0 \| 0 \| no \| unclamped \| (2/2) \| 0.791 \| 3.1 \| \| 0.7 \| 7 \| yes \| unclamped \| (2/2) \| 0.766 \| 1.8 \| \| 0.6 \| 7 \| yes \| unclamped \| (2/2) \| 0.80 \| 1.82 \| \| 0.6 \| 15 \| yes \| unclamped \| (2/2) \| 0.798 \| 1.5 \| \| 1.0 \| 0 \| no \| unclamped \| (2/2.25) \| 0.731 \| 2.2 \| \| 1.0 \| 7 \| yes \| unclamped \| (2/2.25) \| 0.74 \| 2.1 \| Figure Captions Fig. 1. Schematic sketch of the unit cell of the SILO quantum wire for two model structures considered. Each unit cell consists of 8 pairs of (2/2) or (2/2.25) GaAs/InAs short-period superlattices (SPS). In structure 1, the (2/2) SPS contains four diatomic layers (only the cation planes are illustrated) and the SPS is repeated eight times. In structure 2, four pairs of (2/2.25) SPS (or 17 diatomic layers) form a period, and the period is repeated twice in the unit cell. Filled and open circles indicate Ga and In rows (each row extends infinitely along the $`[1\overline{1}0]`$ direction). Fig. 2. Diagonal strain distribution of model structure 1 \[corresponding to Fig. 1(a)\] with thick capping layer (clamped case). Solid: $`ϵ_{x^{}x^{}}`$. Dashed: $`ϵ_{y^{}y^{}}`$. Dash-dotted: $`ϵ_{zz}`$. Parameters for alloy mixing: $`x_m=0.8`$ and $`b=7a_{[110]}`$. Fig. 3. Diagonal strain distribution of model structure 1 with thin capping layer (unclamped case). Solid: $`ϵ_{x^{}x^{}}`$. Dashed: $`ϵ_{y^{}y^{}}`$. Dash-dotted: $`ϵ_{z^{}z^{}}`$. Parameters for alloy mixing: $`x_m=0.8`$ and $`b=7a_{[110]}`$. Fig. 4. Diagonal strain distribution of model structure 2 \[corresponding to Fig. 1(b)\] with thin capping layer (unclamped case). Solid: $`ϵ_{x^{}x^{}}`$. Dashed: $`ϵ_{y^{}y^{}}`$. Dash-dotted: $`ϵ_{z^{}z^{}}`$. Parameters for alloy mixing: $`x_m=0.8`$ and $`b=7a_{[110]}`$. Fig. 5. Conduction and valence band edges for 8 pairs of (2/2) SPS structures sandwiched between 60ÅAl<sub>0.24</sub>Ga<sub>0.24</sub>In<sub>0.52</sub>As barriers with different degrees of lateral alloy mixing as functions of the coordinate $`y^{}`$ for both clamped (thick capping layer) and unclamped (thin capping layer) case. In upper pane, the dashed line is for abrupt modulation (without alloy mixing) and the solid line is for lateral alloy mixing with $`x_m=0.8`$ and $`b=7a_{[110]}`$. In lower panel, dash-dotted line: without alloy mixing; solid line: $`x_m=0.7`$ and $`b=7a_{[110]}`$; dotted lines: $`x_m=0.6`$ and $`b=7a_{[110]}`$; dashed lines: $`x_m=0.6`$ and $`b=15a_{[110]}`$. Fig. 6. Conduction and valence band edges for 8 pairs of (2/2.25) SPS structures sandwiched between 60ÅAl<sub>0.24</sub>Ga<sub>0.24</sub>In<sub>0.52</sub>As barriers with different degrees of lateral alloy mixing as functions of the coordinate $`y^{}`$ for both clamped (thick capping layer) and unclamped (thin capping layer) case. Dashed lines: without alloy mixing. Solid lines: with alloy mixing described by $`x_m=1.0`$ and $`b=7a_{[110]}`$. Fig. 7. Valance subband structures and squared optical matrix elements for self-assembled QWR made of (2/2) SPS structure with thick capping layer (clamped case) and without alloy mixing. Fig. 8. Valance subband structures and squared optical matrix elements for self-assembled QWR made of (2/2) SPS structure with thin capping layer (unclamped case) and with alloy mixing described by $`x_m=0.8`$ and $`b=7a_{[110]}`$.
warning/0005/cond-mat0005096.html	ar5iv	text	# Optical Properties of the Spin-Ladder Compound Sr14Cu24O41 ## I Introduction The Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub> compound is one of the three stable phases in the Sr-Cu-O system which can be synthesized under ambient pressure. The other two stable phases are Sr<sub>2</sub>CuO<sub>3</sub>, which has simple chains of Cu ions, and SrCuO<sub>2</sub> with zigzag chains of Cu ions. This oxide has unique crystal structure based on two sublattices; one of them consists of Cu<sub>2</sub>O<sub>3</sub> two-leg ladders and the second one is formed by CuO<sub>2</sub> chains. These two sublattices are incompatible along one crystallographic direction, thus resulting in an 1D incommensurate structure a1 . According to the structural analysis a1 , the ladder sublattice is face-centered-orthorhombic (space group $`Fmmm`$) with a lattice parameters a=1.1459 nm, b=1.3368 nm and c<sub>Ladder</sub>=0.3931 nm. There are two ladder layers with two ladders per unit cell, Fig. 1. The chain sublattice is A-centered orthorhombic (space group $`Amma`$), with nearly the same a and b axes but different c axis, c<sub>Chain</sub>=0.2749 nm. However, Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub> can be considered as nearly commensurate structure at 7xc<sub>Ladder</sub>=2.7372 nm and 10xc<sub>Chain</sub>=2.7534 nm. The schematic illustration of Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub> crystal structure is given in Fig. 1. The physical properties of Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub> have attracted a lot of attention recently a2 ; a3 ; a4 ; a5 ; a6 ; a7 ; a8 ; a9 ; a10 ; a11 ; a12 , in connection with the rich physics associated with the S=1/2 Heisenberg antiferromagnetic quasi one-dimensional (1D) structures and the discovery of superconductivity in Sr<sub>0.4</sub>Ca<sub>13.6</sub>Cu<sub>24</sub>O<sub>41</sub> under high pressure a2 . The progress in this field has been summarized in Ref. a13 . Various magnetic a3 ; a4 ; a5 , NMR a6 ; a7 and neutron scattering a8 ; a9 measurements, showed that Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub> has two kinds of magnetic gaps. The first one, $`\mathrm{\Delta }_L`$ = 32.5 meV, is attributed to the singlet-triplet excitation in the Cu<sub>2</sub>O<sub>3</sub> spin-ladders, with the exchange energies along the legs (J=130 meV) and the rungs (J<sub>0</sub>=72 meV) a8 . The second one, $`\mathrm{\Delta }_C`$ = 11.5 meV, is argued to arise from the spin dimer formation in the CuO<sub>2</sub> chains a8 , with an antiferromagnetic intradimer coupling J<sub>1</sub>=11.2 meV. Similar ratio of the superexchange interaction energies, J<sub>0</sub>/J $``$ 0.5, as well as the magnitude of J<sub>0</sub>=(950$`\pm `$300) K, is found in the <sup>17</sup>O and <sup>63</sup>Cu NMR measurements a6 . The Raman spectra of Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub> were measured previously a14 ; a15 . From the comparison between Raman spectra of the various layered tetragonal cuprates, Abrashev et al.. a14 concluded that the main contribution to the spectra comes from the Raman forbidden-infrared active LO phonons and the two-magnon scattering. Furthermore, Sugai et al. a15 argued that, besides strong two-magnon features, some low-frequency modes in the Raman spectra are also magnetic in origin, since they have similar energies to those found in the neutron scattering experiments a8 . Still, for the proper identification of the magnetic modes, the temperature dependent Raman spectra, as well as the spectra in magnetic field, are required. Also, detailed analysis of the lattice dynamics and comparison between Raman and infrared (IR) spectra are indispensable due to the incommensurability of the structure. Therefore, we present here the Raman and IR spectra at various temperatures between 5 and 300 K in order to make more complete assignment of the vibrational modes in Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub>. The Raman spectra are also measured under resonant conditions, with a laser light energy close to gap values. The correlated electron tight-binding model of electronic structure is used to estimate the hopping parameters, in fact adjusted to the measured gaps and exchange energies. ## II Experimental details The present work was performed on (010) oriented single crystal plates with dimensions typically about 5 x 1 x 6 mm<sup>3</sup> in the a, b and c axes, respectively. The infrared measurements were carried out with a BOMEM DA-8 FIR spectrometer. A DTGS pyroelectric detector was used to cover the wave number region from 100 to 700 cm<sup>-1</sup>; a liquid nitrogen cooled HgCdTe detector was used from 500 to 1500 cm<sup>-1</sup>. The spectra were collected with the 2 cm<sup>-1</sup> resolution. The low temperature reflectivity spectra in the range from 30 to 5000 cm<sup>-1</sup> were measured using Bruker IFS 133v FIR-spectrometer with Oxford-Cryostat. The Raman spectra were recorded in the backscattering configuration using micro- and macro- Raman systems with Dilor triple monochromator including liquid nitrogen cooled CCD-detector. An Ar- and Kr -ion lasers were used as an excitation source. We measured the pseudodielectric function with a help of a rotating-analyzer ellipsometer. We used a Xe-lamp as a light source, a double monochromator with 1200 lines/mm gratings and an S20 photomultiplier tube as a detector. The polarizer and analyzer were Rochon prisms. The measurements were performed in the 1.6-5.6 eV energy range. Optical reflectivity spectra were measured at room temperature in the 200-2000 nm spectral range using Perkin-Elmer Lambda 19 spectrophotometer. ## III Electronic structure The electronic structure of Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub> is calculated using the tight-binding method for correlated electrons a16 . Recent exact diagonalization and a variational Monte Carlo simulations revealed that electronic structure of Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub> is well described by single ladder energy spectrum a17 . It means that electron energy dispersions are governed mainly by electrons in the ladder. According to the ARPES measurements $`\left[10\right]`$, the chains contribute to the electronic structure of Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub> with a dispersionless band. Without entering any lengthy discussions about the substance nonstoichiometry and the carrier transport between chains and ladders, we assume further on that the ladder unit Cu$`{}_{}{}^{1+n}{}_{2}{}^{}`$O$`{}_{}{}^{2}{}_{3}{}^{}`$ has total charge $`2`$. In other words, there is one hole, $`n=1`$, per copper ion in the ladder for negligibly small hybridization of its $`d_{x^2y^2}`$orbitals with the $`p_y`$orbital of intermediate oxygens. The angle between Cu atoms of neighboring ladders is near right angle $`\left(88.7^{}\right)`$, Fig. 1. In our consideration of the electronic structure we assumed that the directions between the nearest neighbor Cu ions form an ideal right angle. The Hamiltonian for the correlated copper holes in the ladder with two rungs, $`ab`$ and $`cd`$, per a unit cell can be written as $`H`$ $`=`$ $`2t{\displaystyle \underset{p,\sigma }{}}\mathrm{cos}p_y\left[a_\sigma ^+\left(p\right)a_\sigma \left(p\right)+b_\sigma ^+\left(p\right)b_\sigma \left(p\right)+c_\sigma ^+\left(p\right)c_\sigma \left(p\right)+d_\sigma ^+\left(p\right)d_\sigma \left(p\right)\right]`$ $`t_0{\displaystyle \underset{p,\sigma }{}}[a_\sigma ^+\left(p\right)b_\sigma \left(p\right)+c_\sigma ^+\left(p\right)d_\sigma \left(p\right)+H.c.]`$ $`t_{xy}{\displaystyle \underset{p,\sigma }{}}[e^{i\sqrt{2}p_x}+e^{i(\sqrt{2}p_x+p_y)}][a_\sigma ^+\left(p\right)d_\sigma \left(p\right)+H.c.]`$ $`t_{xy}{\displaystyle \underset{p,\sigma }{}}(1+e^{ip_y})[b_\sigma ^+\left(p\right)c_\sigma \left(p\right)+H.c.]+U{\displaystyle \underset{i=a,b,c,d}{}}n_i^jn_i^j\mu {\displaystyle \underset{i=a,b,c,d}{}}n_i^j,`$ where $`a,b,c,d`$ represent chains, (see Fig.1.(a)), $`t`$ $`(t_0)`$ are the values of the carrier hopping along legs (rungs), $`t_{xy\text{ }}`$is a hopping amplitude between ladders, $`U`$ is the Anderson-Hubbard repulsion, and $`\mu `$ is the chemical potential and other notations as usual. The $`x`$ and $`y`$ axis of the reference system are taken along the a and the c crystallographic directions, respectively. Applying the $`X`$operator machinery a18 , the correlation split energy bands are governed by zeros of the inverse Green’s function in the first perturbation order with respect to tunneling matrix: $$\widehat{D}_p^1\left(\omega \right)=\left(\begin{array}{cc}\text{A}_{ab}\hfill & \text{B}\hfill \\ \text{B}^{}\hfill & \text{A}_{cd}\hfill \end{array}\right),$$ (2) where $`\widehat{A}_{ab}`$ $`=`$ $`\begin{array}{c}a\text{ }\{\begin{array}{c}\text{0+}\hfill \\ \text{-2}\hfill \end{array}\hfill \\ b\text{ }\{\begin{array}{c}\text{0+}\hfill \\ \text{-2}\hfill \end{array}\hfill \end{array}\left(\begin{array}{cccc}\frac{i\omega _n\mu }{f_{0+}}+r\hfill & r\hfill & t_0\hfill & t_0\hfill \\ r\hfill & \frac{i\omega _n\mu +U}{f_2}+r\hfill & t_0\hfill & t_0\hfill \\ t_0\hfill & t_0\hfill & \frac{i\omega _n\mu }{f_{0+}}+r\hfill & r\hfill \\ t_0\hfill & t_0\hfill & r\hfill & \frac{i\omega _n\mu +U}{f_2}+r\hfill \end{array}\right),`$ $`\widehat{B}`$ $`=`$ $`\left(\begin{array}{cccc}0\hfill & 0\hfill & D\hfill & D\hfill \\ 0\hfill & 0\hfill & D\hfill & D\hfill \\ C\hfill & C\hfill & 0\hfill & 0\hfill \\ C\hfill & C\hfill & 0\hfill & 0\hfill \end{array}\right)`$ $`\text{with }r`$ $`=`$ $`2t\mathrm{cos}p_y,C=t_{xy}\left(1+e^{ip_y}\right),D=t_{xy}\left[e^{i\sqrt{2}p_x}+e^{i(\sqrt{2}p_x+p_y)}\right]\text{.}`$ Here the correlation factors $`f_{0+}`$, $`f_2`$ are determined by fermion occupation $`n`$ per site. Namely, for the considered paramagnetic phase they are $`f_{0+}=1n/2,`$ $`f_2=n/2`$ and all equal to $`1/2`$ ($`n=1`$). After an analytical continuation, $`i\omega _n\xi +i\delta `$, in Eq. $`\left\|\widehat{D}_p^1\left(\omega \right)\right\|=0`$ one can find the bonding /antibonding correlation energy dispersions as follows: $$\xi _B^\pm \left(p\right)=\frac{1}{2}\left[\epsilon _p^{1,2}+\sqrt{\left(\epsilon _p^{1,2}\right)^2+U^2}\right],$$ (7) $$\xi _A^\pm \left(p\right)=\frac{1}{2}\left[\epsilon _p^{1,2}\sqrt{\left(\epsilon _p^{1,2}\right)^2+U^2}\right],$$ (8) where $`\epsilon _p^{1,2}=2t\mathrm{cos}p_y\pm \sqrt{t_0^2+\left(2t_{xy}\mathrm{cos}\frac{p_y}{2}\right)^2\pm 4t_0t_{xy}\mathrm{cos}\frac{p_y}{2}\mathrm{cos}\frac{\sqrt{2}p_x}{2}}`$. For derivation of these energy dispersions from the eight-by-eight fold secular equation (see Eq.(2)) it was useful to apply the theorem about the decomposition of determinant with respect to diagonal elements (see Appendix A in Ref. a19 ). The subbands $`\xi _B^{}`$ and $`\xi _A^{}`$ are completely occupied by carriers with concentration $`n=1`$ per copper site of ladder for the chosen chemical potential $`\mu =U/2`$. The nearest unoccupied energy band is $`\xi _B^+`$ and the correlation gap in electronic structure can be estimated as $`\mathrm{\Delta }_{corr}`$ $`=`$ $`\mathrm{min}\xi _B^+\left(p\right)\mathrm{max}\xi _A^{}\left(p\right)=`$ $`{\displaystyle \frac{1}{2}}\left[\sqrt{\left(t_0+2t2t_{xy}\right)^2+U^2}+\sqrt{\left(t_0+2t\right)^2+U^2}\right]\left(t_0+2t+t_{xy}\right).`$ For the dimensionless energies, $`\omega _p^{+,}=\epsilon _p^{1,2}/2t`$, $`\tau _0=t_0/2t,\tau =t_{xy}/2t,`$ the non-correlated electron density of states per spin, $`\rho `$($`\epsilon `$)$`=\underset{p_{x,}p_y}{}`$\[$`\delta `$($`\epsilon \omega _p^+`$)$`+\delta `$($`\epsilon \omega _p^{}`$)\], for the unit cell volume, is defined analytically as follows $`\rho _0^\alpha \left(12\tau +\alpha \tau _0\epsilon 1+2\tau \tau _0\right)`$ $`=`$ $`{\displaystyle \frac{4}{\pi ^2\sqrt{k_\alpha \tau }}}K\left(q_\alpha \right),`$ $`\rho _0^\alpha \left(1+2\tau +\alpha \tau _0\epsilon {\displaystyle \frac{1}{2}}+\tau +\alpha \tau _0\right)`$ $`=`$ $`{\displaystyle \frac{4}{\pi ^2q_\alpha \sqrt{k_\alpha \tau }}}F(\mathrm{arcsin}a_\alpha ;{\displaystyle \frac{1}{q_\alpha }}),`$ (10) $`\rho _0^\alpha \left(1+2\tau +\alpha \tau _0\epsilon 1+\alpha \tau _0\right)`$ $`=`$ $`{\displaystyle \frac{4}{\pi ^2q_\alpha \sqrt{k_\alpha \tau }}}K\left({\displaystyle \frac{1}{q_\alpha }}\right).`$ Eqs.(6) represent the analytical expressions for the electron density of states via elliptic integrals $`F`$ and $`K`$ of the 1-st kind in the Legendre normal form with modulus $`q_\alpha =\sqrt{\left[2\tau \left(\tau +k_\alpha \right)+1\left(\epsilon +\alpha \tau _0\right)^2\right]/k_\alpha \tau }/2`$, argument $`a_\alpha =\sqrt{\left[2\tau \left(\tau +k_\alpha \right)+1\left(\epsilon +\alpha \tau _0\right)^2\right]/\left[\left(1+\epsilon +\alpha \tau _0\right)\left(\tau +k_\alpha \right)k_\alpha \right]}`$, where $`k_\alpha =\sqrt{\tau ^2+2\left(1\epsilon \alpha \tau _0\right)}`$ and $`\alpha =\pm `$. The electron-electron repulsion, $`U`$, splits the density of the non-correlated electronic states $`\rho _0`$. The correlated electron density of states is $`\rho \left(\epsilon \right)`$ $`=`$ $`{\displaystyle \frac{2}{1+\frac{\xi _+^{^{}}}{\sqrt{\left(\xi _+^{^{}}\right)^2+\left(\frac{U}{2t}\right)^2}}}}\rho _0\left(\xi _+^{^{}}\right)+{\displaystyle \frac{2}{1\frac{\xi _{}^{^{}}}{\sqrt{\left(\xi _{}^{^{}}\right)^2+\left(\frac{U}{2t}\right)^2}}}}\rho _0\left(\xi _{}^{^{}}\right),\text{where}`$ (11) $`\xi _\pm ^{^{}}`$ $`=`$ $`{\displaystyle \frac{\xi _\pm ^^2\left(\frac{U}{2t}\right)^2}{\xi _\pm }}+S`$ $`\text{and }S`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left[\sqrt{\left(1+2\tau +\tau _0\right)^2+\left({\displaystyle \frac{U}{t}}\right)^2}\sqrt{\left(1+\tau _0\right)^2+\left({\displaystyle \frac{U}{t}}\right)^2}2\tau \right],`$ are expressed via dimensionless correlated energies $`\xi _\pm \xi ^\pm \left(p\right)/2t`$. With the help of Eq.(7) and Eqs.(6) one can calculate the correlated electron density of states for corresponding energy ranges. Its explicit form naturally includes the first kind elliptic integrals. The results of the calculations are plotted in Fig. 2. The overlap of the energy ranges for the electronic dispersions, Eq.(3), leads to the special features of the correlated electronic structure at $`L_5=1/2\left(12\tau +\tau _0\right)S1/2\sqrt{\left(1+2\tau +\tau _0\right)^2+\left(U/t\right)^2}`$and $`L_7=1/2\left(1/2+\tau +\tau _0\right)S1/2\sqrt{\left(1/2+\tau +\tau _0\right)^2+\left(U/t\right)^2}`$ in the lower correlated band. Logarithmic divergencies inside the band at $`L_2=1/2\left(1+2\tau \tau _0\right)S1/2\sqrt{\left(1+2\tau \tau _0\right)^2+\left(U/t\right)^2}`$, $`L_4=1/2\left(1\tau _0\right)S1/2\sqrt{\left(1\tau _0\right)^2+\left(U/t\right)^2}`$ , $`L_6=1/2\left(1+2\tau +\tau _0\right)S1/2\sqrt{\left(1+2\tau +\tau _0\right)^2+\left(U/t\right)^2}`$ and at the correlated band edge $`L_8=1/2\left(1+\tau _0\right)S1/2\sqrt{\left(1+\tau _0\right)^2+\left(U/t\right)^2}`$ are clear manifestations of the 2D electronic structure of Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub> compound. We would like to emphasize that in the one-dimensional case $`\left(t_{xy}0\right)`$ the electron density of states is taking features of a single spin-ladder without any logarithmic peaks, then the divergencies become square-root like and they are located at the band edges, $`\epsilon =\pm 1`$: $`\rho \left(\epsilon \right)=8K\left(0\right)/\pi ^2\sqrt{1\epsilon ^2}=4/\pi \sqrt{1\epsilon ^2}`$. ## IV Experimental results The dielectric function $`ϵ_2`$ of Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub> is shown in Fig.3 in the spectral region from 1.6 eV to 5.5 eV. These spectra were computed from the measured Fourier coefficients using the equations for an isotropic case. Consequently, $`ϵ_2`$ represents a complicated average of the projections of the dielectric tensor on the sample surface. We presented the spectra of the (010) surface taken with the a-axis, Fig. 3(a) and c axis, Fig. 3(b), in the plane of incidence (PI). According to the prescription given by Aspnes a20 , we attribute these components to the components of the dielectric tensor $`ϵ_2^{aa}`$ and $`ϵ_2^{cc}`$. The bands with the energies of 2.4, 4.1, and 4.7 eV for the a-axis and at about 1.86, 2.34, 2.5, and 4.3 eV are found for the c axis in the plane of incidence, respectively. Inset (a) in Fig. 3 shows reflectivity spectra of Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub>. These spectra are calculated from measured dielectric functions $`ϵ_1`$ and $`ϵ_2`$. Inset (b) in Fig. 3 represents the unpolarized optical reflectivity of Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub> measured at room temperature. In addition to peaks, previously observed in ellipsometric measurements, a new peak at about 1.4 eV appears. The room temperature polarized far-infrared reflectivity spectra of Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub> are given in Fig. 4. The open circles are the experimental data and the solid lines represent the spectra computed using a four-parameter model for the dielectric constant: $$ϵ(\omega )=ϵ_{\mathrm{}}(\underset{j=1}{\overset{n}{}}\frac{\omega _{LO,j}^2\omega ^2+ı\gamma _{LO,j}\omega }{\omega _{TO,j}^2\omega ^2+ı\gamma _{TO,j}\omega }\frac{\omega _p^2}{\omega (\omega ı\tau ^1)}),$$ (12) where $`\omega _{LO,j}`$ and $`\omega _{TO,j}`$ are longitudinal and transverse frequencies of the j<sup>th</sup> oscillator, $`\gamma _{LO,j}`$ and $`\gamma _{TO,j}`$ are their corresponding dampings, $`\omega _p`$ is the plasma frequency, $`\tau `$ is the free-carrier relaxation time and $`ϵ_{\mathrm{}}`$ is the high-frequency dielectric constant. The best fit parameters are given in Table I. The agreement between observed and calculated reflectivity spectra is rather good. For the E$`\|\|`$a polarization, eight oscillators with TO frequencies at about 164, 194, 219, 249, 283.5, 310.4, 554 and 623 cm<sup>-1</sup> are clearly seen. In the E$`\|\|`$c polarization, Fig. 4(b), nine oscillators at 135, 148, 253, 293, 345, 486, 540, 596 and 620 cm<sup>-1</sup> are observed. Besides phonons, our model includes the Drude expression for light scattering on free carriers. We obtained the plasma frequency at about 4000 cm<sup>-1</sup> (1000 cm<sup>-1</sup>) for the E$`\|\|`$c (E$`\|\|`$a) polarizations. The room temperature Raman spectra of Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub>, for (aa) and (cc) polarized configurations are presented in Figs. 5(a) and 5(b). These spectra consist of only A<sub>g</sub> symmetry modes. Four modes at 246, 302, 548, and 582 cm<sup>-1</sup> are clearly seen. The low temperature (cc) Raman spectra are given in Figs. 5(c)-(g). By lowering temperature below 200 K, the modes narrow and in addition the new modes appear. We will discuss them later on. The Raman spectra of Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub>, excited by different lines of Ar and Kr lasers at 8 K, are shown in Fig. 6 for the (cc) and (aa) polarized configurations in the spectral ranges: (a) from 700 to 1400 cm<sup>-1</sup>, (b) from 1675 to 1975 cm<sup>-1</sup> and (c) from 2600 to 3300 cm<sup>-1</sup>. Anticipating our conclusions, we divide the Raman spectra in three different energy regions: one phonon (0-700 cm<sup>-1</sup>), two-phonon (700-1400 cm<sup>-1</sup>) and two-magnon region (above 1500 cm<sup>-1</sup>). The mode at about 2900 cm<sup>-1</sup> and a broad structure at about 1900 cm<sup>-1</sup> are magnetic in origin, according to their intensity and frequency dependence as a function of the temperature, see Fig. 7. ## V Discussion The average unit cell of Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub> consists of four formula units with 316 atoms in all. Since there is a large number of atoms in the unit cell, we can expect a very large number of optically active modes. Consequently, the lattice dynamical calculation is practically impossible. All atoms have 4(e) position symmetry of $`Pcc2`$ $`(C_{2v}^3)`$ space group a1 . Factor-group-analysis (FGA) yields the following distribution of vibrational modes: $$\mathrm{\Gamma }_{Sr_{14}Cu_{24}O_{41}}=237A_1(𝐄\|\|𝐜,aa,bb,cc)+237A_2(ab)+237B_1(𝐄\|\|𝐚,ac)+237B_2(𝐄\|\|𝐛,bc)$$ (13) According to this representation one can expect 948 modes which are both Raman and infrared active. Experimentally, the number of observed modes is less then ten for each polarization. Because of that, we consider separately the contribution of each sublattice unit. As mentioned earlier, the space group of ladder sublattice is $`Fmmn`$ $`(D_{2h}^{23})`$. The site symmetries of Sr, Cu, O<sub>1</sub> and O<sub>2</sub> atoms are (8h), (8g), (8g) and (4b), respectively. The FGA for the ladder structure (Sr<sub>2</sub>Cu<sub>2</sub>O<sub>3</sub>) yields a21 : $`\mathrm{\Gamma }_{Ladder}`$ = 3$`A_g+3B_{1g}+2B_{2g}+B_{3g}+4B_{1u}+4B_{2u}+4B_{3u}`$ The space group of a chain sublattice is $`Amma`$ $`(D_{2h}^{17})`$. The site symmetries of Cu and O atoms are (4c) and (8f). The FGA gives for the chain structure: $`\mathrm{\Gamma }_{Chain}`$ = $`3A_g+3B_{1g}+2B_{2g}+B_{3g}+A_u+2B_{1u}+3B_{2u}+3B_{3u}`$, Since $`Amma`$ is not a standard setting for $`D_{2h}^{17}`$ space group ($`Cmcm`$) we use $`C_s^{xy}`$ site symmetry for oxygen atoms and $`C_{2v}^x`$ symmetry for Cu atoms in above representations. Thus, the total number of vibrational modes from both sub-units is: $$\mathrm{\Gamma }=6A_g(aa,bb,cc)+6B_{1g}(ab)+4B_{2g}(ac)+2B_{3g}(bc)+A_u+6B_{1u}(𝐄\|\|𝐜)+7B_{2u}(𝐄\|\|𝐛)+7B_{3u}(𝐄\|\|𝐚)$$ (14) According to this analysis we should expect 6A<sub>g</sub> modes; one mode from vibrations of the Sr atoms, two modes which originate from vibrations of Cu atoms and other three A<sub>g</sub> modes are due to oxygen vibrations. In order to assign the observed A<sub>g</sub> modes we compare our spectra with the corresponding spectra of the Cu-O based materials with similar structural units as in Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub>. For example, in SrCuO<sub>2</sub> a22 and YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub> a23 the Cu-O double layers exist and resemble the one leg of the ladder structure in Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub>. The Cu-O chains, formed from copper oxide squares with the common edges, as in Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub>, are also present in CuO a24 and CuGeO<sub>3</sub> a25 . Thus, the lowest frequency mode in Fig. 5(a) at 246 cm<sup>-1</sup> can be assigned to vibrations of the Cu ladder atoms (see Fig. 1). The corresponding A<sub>g</sub> mode of copper vibrations in YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub> (SrCuO<sub>2</sub>) appears at 250 (263) cm<sup>-1</sup>. The next mode is found at 302 cm<sup>-1</sup>. This mode represents the vibrations of the chain oxygen atoms along the a-axis and appears in CuO at the same frequency a24 . The mode at 548 cm<sup>-1</sup> is breathing mode of O<sub>1</sub> oxygen ladder atoms. Similar mode appears in SrCuO<sub>2</sub> at 543 cm<sup>-1</sup>. The second oxygen ladder atom (O<sub>2</sub>), Fig. 1, is situated in the center of inversion ($`D_{2h}`$ site symmetry) and has no Raman activity. The highest frequency Raman mode in Figs. 5(a)-(b), at about 582 cm<sup>-1</sup>, is caused by the chain oxygen vibrations along the c-axis. Corresponding mode appears in CuO at 633 cm<sup>-1</sup> a24 and in CuGeO<sub>3</sub> at 594 cm<sup>-1</sup> a25 . The vibrations of Sr atoms, with frequency of 188 cm<sup>-1</sup> as in SrCuO<sub>2</sub>, are not observed in the spectra. Finally, the A<sub>g</sub> mode at 153 cm<sup>-1</sup>, see Fig. 8, originates from the vibrations of the Cu atom in chains. Again, similar mode is found in YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub> at 153 cm<sup>-1</sup>. By lowering temperature Raman peaks narrow and at about T=150 K the new modes appear (Fig. 5 and Table II). Similar effects are found in the IR spectra as well. This temperature coincides with the charge ordering temperature established in the NMR and neutron scattering experiments a7 ; a9 . NMR study a7 showed the splitting of signal from Cu<sup>3+</sup> ions into two peaks below 200K suggesting the occurrence of the charge ordering. This effect is confirmed by Cox et al. a11 . They measured synchrotron x-ray scattering on Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub> single crystals and showed the appearance of the satellite peaks at (00l) positions. The results are interpreted in terms of a charge-ordered model involving both dimerization between two-nearest-neighbors of Cu<sup>2+</sup> ions surrounding a Cu<sup>3+</sup> ion on a Zhang-Rice singlet site, and dimerization between nearest-neighbors of Cu<sup>2+</sup> ions. Fig. 8(a) shows the (cc) polarized low temperature (T=10 K) Raman spectra of Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub> in the 125 - 750 cm<sup>-1</sup> spectral region, excited with 647.1 nm (1.91 eV) and 488 nm (2.54 eV) energies. There are a remarkable difference between Raman spectra for these excitation lines which appears due to resonance effects. Namely, both lines are very close to gap energies for polarization along c-axis, see Fig. 3. The reflectivity spectra measured at 10 K for the E$`\|\|`$c and the E$`\|\|`$a polarizations are given in Figs. 8(b) and 8(c), respectively. In order to compare Raman with IR data we shown in the same figure the $`ϵ_2(\omega )`$ and the $`Im[1/ϵ(\omega )]`$ spectra. These spectra are obtained using Kramers-Kronig analysis of reflectivity data. The TO and LO mode frequencies are obtained as peak positions of the $`ϵ_2(\omega )`$ and the $`Im[1/ϵ(\omega )]`$, respectively. For most of all Raman active modes we found theirs infrared counterparts (some of them are denoted by vertical lines) either for the E$`\|\|`$a or the E$`\|\|`$c polarizations. The appearance of similar lines in IR and Raman spectra, if not being the consequence of the symmetry, may also be attributed to the resonant conditions. It is well documented a26 ; a27 that Raman forbidden IR active LO modes appear in the Raman spectra of the insulating Cu-O based materials for the laser line energies close to gap values. The appearance of these modes in the Raman spectra is explained by Fröhlich interaction a26 . Here, since the 647.1 nm line is very close to the gap values (1.86 eV, see Fig. 3) one can expected that the IR LO modes appear in the Raman spectra of Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub> as well. Such effect is usually accompanied by the observation of the strong phonon overtones, as we shown in Fig. 6(a). All modes with energies higher then 700 cm<sup>-1</sup> are in fact the second order combinations (overtones) of the low-energy modes. The assignment and the frequencies of these modes are given in Table III. Therefore, the properties of the modes observed in our spectra may be understood in terms of the available symmetry combined with resonant effects, thus making proper identification of the low-temperature phonons practically impossible without detailed structural analysis. Now we focus on magnetic properties. As it is mentioned earlier, the neutron scattering and NMR measurements estimated the spin-ladder gap value to be at $`\mathrm{\Delta }_L`$ = 32.5 meV (268 cm<sup>-1</sup>) a8 or at 40.5 meV (326 cm<sup>-1</sup>) a28 , respectively. Because of that, we paid special attention to the 200-350 cm<sup>-1</sup> spectral range, Fig. 5 (left panel). By lowering temperature we observe the appearance of the new modes at 262, 293, and 317 cm<sup>-1</sup>. At the same time we find the modes with the same energy in the low-temperature IR spectra, Figs. 8(b)-(c). The 262 cm<sup>-1</sup> mode is very close in energy to the magnetic gap, thus possibly one-magnon excitation, as proposed by Sugai et al.a15 . However, the origin of the one-magnon excitation in the light scattering process usually comes from the spin-orbit interaction, which is found to be very small in transition metal oxides due to quenched orbital momentum of the transition metal ions. Moreover, below 100 K, this mode has nearly the same temperature dependency of the frequency and intensity, like all other low-temperature modes and we identify them as zone edge phonons, which become Raman active due to the zone folding effect caused by the charge ordering transition a7 . Yet another type of magnetic excitations is expected to appear in the Raman spectra of two-leg-ladders at energies close to $`2\mathrm{\Delta }_L`$ 534 cm<sup>-1</sup> a29 . The mode at 498 cm<sup>-1</sup> (see Fig. 5 (g), right panel) shows a typical asymmetric shape with a tail towards high-frequencies, as expected for the onset of the two-magnon continuum. Its energy difference from $`2\mathrm{\Delta }_L`$ could be the magnon binding energy. However, as it is shown in Fig. 8(b), this mode is positioned between TO and LO frequencies of very strong IR active mode in the E$`\|\|`$c spectra. This mode also shows a strong resonant enhancement, see Fig. 8 (a). Therefore, at this stage, it is hard to make definite conclusions about the origin of this mode and further experiments are needed to clarify this issue. Similar discussion holds for the spectral range around the twice the spin-gap value associated with a chains, $`2\mathrm{\Delta }_C`$180 cm<sup>-1</sup> a8 , where continuum-like feature is also found in the Raman spectra. Finally, we discuss the modes in the spectral range above 1500 cm<sup>-1</sup>. The strongest mode in the spectra is centered at about 2840 cm<sup>-1</sup> for the 488 nm excitation and at about 2880 cm<sup>-1</sup> for the 514.5 nm excitation line. This feature decrease in intensity and shifts to lower energies at higher temperatures, Fig. 7. The same structure is already observed in many copper oxides at similar frequencies a27 ; a29 . Thus, all observed effects indicate the two-magnon origin of this mode. The energy of the two-magnon mode, associated with a top of the magnon brunch, in copper oxide insulators is about 3J, where J represents the exchange interaction. In the case of the two-leg-ladders, its energy position for different polarized configurations in the Raman spectra may be used to estimate the exchange parameters parallel ($`J`$) and perpendicular ($`J_0`$) to the ladders a30 . Since the energy position of the two-magnon peak is the same for (cc) and (aa) polarizations, Fig. 6(c), such an analysis suggests that $`J=J_0`$. This conclusion is fully in agreement with Raman scattering data of Ref.a15 and recent high-energy neutron scattering measurementsa31 , but inconsistent with previous neutron a8 , NMRa6 and magnetization measurements, which estimated $`J_0/J`$ ratio to be between 0.5 and 0.8. The discrepancy may be related to the fact that in the previous measurements, the high-energy magnetic excitations were not observed directly, but estimated indirectly from the low-energy spin gap measurements assuming an ideal model a32 . Thus, the Raman scattering is more direct method to obtain $`J_0/J`$ ratio. The two-magnon mode at 2880 cm<sup>-1</sup> gives exchange energy $`J`$=119 meV, very close to the neutron scattering value $`J`$=130 meV a8 . It is also interesting to note that the two-magnon mode is asymmetric with the spectral weight shifted to higher frequencies. Such a spectral shape of the two-magnon mode can be a consequence of the resonance a30 , or it can be related to the bound-hole-pair effects a15 . Still, further experiments on the hole doped crystals are necessary to clarify this point. Let us consider Fig. 6(b), where a weak structure appears in (cc) spectra at about 1920 cm<sup>-1</sup>. Its energy is exactly equal to 2J and varies with a laser line frequency in a similar way as the two-magnon mode at 2880 cm<sup>-1</sup>. The energy shift of this mode as a function of temperature was not seen because of its very low intensity, thus leaving the origin of this mode as an open question. In addition to the 1920 cm<sup>-1</sup> mode, we found a weak structure at about 1700 cm<sup>-1</sup>, as well. This mode does not possess any noticeable temperature dependencies of energy and intensity. We concluded that this mode is an overtone phonon mode. Its frequency can be represented as the third order of the 568 cm<sup>-1</sup> mode ($`3c_3`$, see Table II). Finally, we discuss the electron energy dispersions and density of states, which are calculated using the model described in Sect. III. First, we analyze the influence of the energy transfer (hopping) to the correlation gap. Fig. 9 shows the correlation gap $`\mathrm{\Delta }_{corr}`$ vs. the Anderson-Hubbard parameter $`U`$. The plotted curves are calculated using Eq.(5) and $`J=4t^2/U`$=0.13 eV for different hopping energy ratios $`t_{xy}/t`$ and $`t_0/t`$. From Fig. 9 we conclude that the main influence to the correlation gap value comes from the Anderson-Hubbard parameter U. Because of that, by knowing the electronic gap from the ellipsometric or optical absorption measurements and the exchange energy J from the Raman spectroscopy we can determine the hopping parameters and the onsite electron-electron repulsion $`U`$. However, relative ratio of the hopping energies perpendicular and parallel to the legs as well as the interladder hopping, does not influence much the correlation gap . Namely, a decrease of the transfer energy along the rungs from $`t_0/t`$=1 to $`t_0/t`$=0.5 increases $`\mathrm{\Delta }_{corr}`$ for about 6%. Also, an increase of the interladder hopping $`t_{xy}`$ from 0 to 10% of the hoping value $`t`$ along the legs, produces a decrease of $`\mathrm{\Delta }_{corr}`$ for about 2.5%. Therefore, the interladder effects on electronic structure of Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub> are found to be negligibly small even though the distance between the neighboring ladders is short (see Fig. 1). The correlation gap is observed at 1.4 eV. This value is determined as a maximum of dielectric function $`ϵ_2(\omega )`$ obtained from KKA of the unpolarized reflectivity data, see Inset (b) of Fig. 3(b). Using this value and the fact that $`t_0/t=1`$ (comes form $`J_0/J=1`$), $`t_{xy}=0`$ we obtained U=2.1 eV. Similar values for the correlation gap and U has been also found in $`SrCuO_2`$ a34 . Energy dispersion, shown in Inset of Fig. 2, allows us to assign the 1.86 eV peak in Fig. 3(a) to the gap value ($`\mathrm{\Delta }_1`$=1.87 eV) between bonding and antibonding bands at the Z-point. Also, 2.4 eV peak from the ellipsometric measurements corresponds to the gap from the lowest occupied band ($`L_6`$, Fig. 2) to the highest empty band at Z-point of Brillouine zone ($`\mathrm{\Delta }_2`$=2.4 eV). By comparison of our measured and calculated gap values, with previously published results a17 ; a35 , we found that our $`t_0/t`$ ratio is close to the ratio determined in a17 . Mizuno et al., a35 , calculated the optical conductivity for small clusters, simulating the ladders and the chains. They obtained the gap for the ladder at about 1.7 eV while the contribution from the chains mainly emerges at a higher energy showing the large spectral weight at around 2.6 eV. These values are very close to our experimental results. Therefore, the peaks at 1.86, 2.4 and 2.5 eV, see Fig. 3(b), may correspond to the ladders and the chains, respectively. In conclusion, we studied the optical properties of the Sr<sub>14</sub>Cu<sub>24</sub>O<sub>41</sub> single crystal. The lattice vibrations are analyzed using the far-infrared reflectivity and Raman scattering measurements in the wide frequency and temperature range. At temperatures below 150 K the new IR and Raman modes appear, presumably due to the charge ordering. The two-magnon excitations are found in the Raman spectra which could be related to minimal ($`2\mathrm{\Delta }`$) and maximal (twice the top of the magnon brunch) magnon energy. The exchange constants along the legs and rungs of the ladders are found to be the same, $`J=J_0`$ 120 meV. The optical reflectivity and the ellipsometric measurements are used to study the charge dynamics. The gap values of 1.4 eV, 1.86 eV (2.5 eV) for the ladders (chains) along the c-axis and 2.4 eV along the a-axis are found. These results are analyzed using tight-binding approach for the correlated electrons. The correlation gap value of 1.4 eV is calculated with the transfer energy (hopping) parameters $`t`$=t<sub>0</sub>=0.26 eV, along and perpendicular to the legs, and U=2.1 eV, as a Coulomb repulsion. ## VI Acknowledgment We thank W. König for low-temperature infrared measurements. Z.V.P., V.A.I and O.P.K acknowledge support from the Research Council of the K.U. Leuven and DWTC. The work at the K.U. Leuven is supported by the Belgian IUAP and Flemish FWO and GOA Programs. M.J.K thanks Roman Herzog - AvH for partial financial support.
warning/0005/hep-th0005067.html	ar5iv	text	# 1 Introduction ## 1 Introduction One of the important motivations in favour of string theory in the mid-eighties was the fact that it seemed to include in principle all the ingredients required to embed the observed standard model (SM) physics inside a fully unified theory with gravity. The standard approach when trying to embed the standard model into string theory has traditionally been an top-down approach. One starts from a string theory like e.g. the $`E_8\times E_8`$ heterotic and reduces the number of dimensions, supersymmetries and the gauge group by an appropriate compactification leading to a massless spectrum as similar as possible to the SM. The paradigm of this approach has been the compactification of the $`E_8\times E_8`$ heterotic on a CY manifold with Euler characteristic $`\chi =\pm 6`$, leading to a three-generation $`E_6`$ model. Further gauge symmetry breaking may be achieved e.g. by the addition of Wilson lines and a final breakdown of $`D=4`$, $`𝒩=1`$ supersymmetry is assumed to take place due to some field-theoretical non-perturbative effects . Other constructions using compact orbifolds or fermionic string models follow essentially the same philosophy . Although since 1995 our view of string theory has substantially changed, the concrete attempts to embed the SM into string theory have essentially followed the same traditional approach. This is the case for instance in the construction of a M-theory compactifications on CY$`\times S^1/Z_2`$ , or of F-theory compactifications on Calabi-Yau (CY) four-folds , leading to new non-perturbative heterotic compactifications. This is still a top-down approach in which matching of the observed low-energy physics is expected to be achieved by searching among the myriads of CY three- or four-folds till we find the correct vacuum <sup>1</sup><sup>1</sup>1For recent attempts at semirealistic model building based on Type IIB orientifolds see refs. .. The traditional top-down approach is in principle a reasonable possibility but it does not exploit fully some of the lessons we have learnt about string theory in recent years, most prominently the fundamental role played by different classes of p-branes (e.g. D-branes) in the structure of the full theory, and the important fact that they localize gauge interactions on their worldvolume without any need for compactification at this level. It also requires an exact knowledge of the complete geography of the compact extra dimensions (e.g. the internal CY space) in order to obtain the final effective action for massless modes. We know that, for example, Type IIB D3-branes have gauge theories with matter fields living in their worldvolume. These fields are localized in the four-dimensional world-volume, and their nature and behaviour depends only on the local structure of the string configuration in the vicinity of that four-dimensional subspace. Thus, as far as gauge interactions are concerned, it seems that the most sensible approach should be to look for D-brane configurations with world volume field theories resembling as much as possible the SM field theory, even before any compactification of the six transverse dimensions. Following this idea, in the present article we propose a bottom-up approach to the embedding of the SM physics into string theory. Instead of looking for particular CY compactifications of a $`D=10,11,12`$ dimensional structure down to four-dimensions we propose to proceed in two steps: i) Look for local configurations of D-branes with worldvolume theories resembling the SM as much as possible. In particular we should search for a gauge group $`SU(3)\times SU(2)\times U(1)`$ but also for the presence of three chiral quark-lepton generations. Asking also for $`D=4`$ $`𝒩=1`$ unbroken supersymmetry may be optional, depending on what our assumptions about what solves the hierarchy problem are. At this level the theory needs no compactification and the D-branes may be embedded in the full 10-dimensional Minkowski space. On the other hand, gravity still remains ten-dimensional, and hence this cannot be the whole story. ii) The above local D-brane configuration may in general be part of a larger global model. In particular, if the six transverse dimensions are now compactified, one can in general obtain appropriate four-dimensional gravity with Planck mass proportional to the compactification radii. This two-step process is illustrated in Figure 1. An important point to realize is that, although taking the first step i.e. finding a ‘SM brane configuration’ may be relatively very restricted, the second step may be done possibly in myriads of manners. Some properties of the effective Lagrangian (e.g. the gauge group, the number of generations, the normalization of coupling constants) will depend only on the local structure of the D-brane configuration. Hence many phenomenological questions can be addressed already at the level of step i). On the other hand, other properties, like some Yukawa couplings and Kähler metrics, will be dependent on the structure of the full global model. In this paper we present the first specific realizations of this bottom-up approach. We believe that the results are very promising and lead to new avenues for the understanding of particle theory applications of string theory. In particular, and concerning the two step bottom-up approach described above we find that: i) One can obtain simple configurations of Type IIB D3, D7 branes with world-volume theories remarkably close to the SM (or some left-right symmetric generalizations). They correspond to collections of D3/D7 branes located at orbifold singularities. The presence of additional branes beyond D3-branes (i.e. D7-branes) is dictated by tadpole cancellation conditions. Finding three quark-lepton generations and $`𝒩=1`$ SUSY turns out to be quite restrictive leading essentially to $`\mathrm{IC}^3/ZZ_3`$ singularities or some variations (including related models with discrete torsion, $`ZZ_3`$ orbifolds of the conifold singularity, or a non-abelian orbifold singularity based on the discrete group $`\mathrm{\Delta }_{27}`$). In these models extra gauged $`U(1)`$’s with triangle anomalies (cured by a generalized Green-Schwarz mechanism) are generically present but decouple at low energies. The appearance of the weak hypercharge $`U(1)_Y`$ in these models is particularly elegant, corresponding to a unique universal linear combination which is always anomaly-free and yields the SM hypercharge assignments automatically. In the case of non-supersymmetric $`ZZ_N`$ orbifold singularities, three quark-lepton generations may be obtained for any $`N>4`$, but the resulting weak angle tends to be too small. ii) These local ‘SM configurations’ may be embedded into compact models yielding correct D=4 gravity. We construct specific compact Type IIB orbifold and orientifold models which contain subsectors given by the realistic D-brane configurations discussed above. In order to cancel global tadpoles (i.e, the total untwisted RR D-brane charges) one can add anti-D3 and/or anti-D7 branes. These antibranes are stuck at orbifold fixed points (in order to ensure stability against brane-antibrane annihilation) and lead to models with hidden sector gravity-induced supersymmetry breaking. Other compact models with unbroken $`D=4`$ $`𝒩=1`$ supersymmetry may be easily obtained in the more general framework of F-theory compactifications, and we construct an specific example of this type. In this approach the embedding of SM physics into F-theory is quite different from those followed previously: the interesting physics resides on D3-branes, rather than D7-branes. As we comment above, some properties of the low-energy physics of the compact models will only depend on the local D3-brane configuration. That is for example the case of hypercharge normalization. The D-brane configurations leading to unbroken $`𝒩=1`$ SUSY predict a tree-level value for the weak angle $`\mathrm{sin}^2\theta _W=3/14=0.215`$, different from the $`SU(5)`$ standard result 3/8. This is compatible with standard logarithmic gauge coupling unification if the string scale is of order $`M_s10^{11}`$ GeV, as we discuss in the text. Other phenomenological aspects like Yukawa couplings depend not only on the singularity structure, but also on the particular form of the compactification. Notice in this respect that although the physics of the D3-branes will be dominated by the presence of the singularity, the D7-branes wrap subspaces in the compact space and are therefore more sensitive to its global structure. The structure of this paper is as follows. In the following chapter we present some general results which will be needed in the remaining sections. We describe the general massless spectrum and couplings of Type IIB D3- and D7-branes on Abelian orbifold singularities, both for the $`𝒩=1`$ and $`𝒩=0`$ cases. We also discuss the consistency conditions of these configurations (cancellation of RR twisted tadpoles), as well as the appearance of non-anomalous $`U(1)`$ gauge symmetries in this class of theories. In chapter 3 we apply the formalism discussed in chapter 2 to the search of realistic three-generation D-brane configurations sitting on $`\mathrm{IR}^6/ZZ_N`$ singularities. We present specific simple $`𝒩=1`$ models leading to the SM gauge group with three quark-lepton generations. We also present an alternative $`SU(3)\times SU(2)_L\times SU(2)_R\times U(1)_{BL}`$ three generation model. We also discuss the case of non-SUSY $`ZZ_N`$ singularities and present an specific three generation non-SUSY model based on a $`ZZ_5`$ singularity. Finally, we discuss different generalizations yielding also three quark-lepton generations. In particular we discuss orbifold singularities with discrete torsion, models based on non-abelian orbifolds, as well as some models based on certain non-orbifold singularities. Although the massless spectrum of these new possibilities is very similar to the models based on the $`ZZ_3`$ singularity, some aspects like Yukawa couplings get modified, which may be interesting phenomenologically. We argue that locating the D3-branes on an orientifold (rather than orbifold) point does not lead to standard model configurations and hence is not very promising. Finally we discuss non-supersymmetric models constructed using branes and antibranes. In chapter 4 we proceed to the second step in our approach and embed the realistic D3/D7 configurations found in the previous chapters into a compact space. We present examples based on type IIB orbifolds and orientifolds. As we mentioned above, in this models the global RR charges may be canceled by the addition of anti-D-branes which are trapped at the fixed points. Some of these models are T-dual to the models recently studied in . We also discuss the construction of models with unbroken $`D=4`$, $`𝒩=1`$ SUSY by considering F-theory compactifications with the SM embedded on D3-branes, and construct an specific example. In chapter 5 we briefly discuss some general phenomenological questions, like gauge coupling unification and Yukawa couplings. We leave our final comments and outlook for chapter 6. In order not to obstruct continuity in the reading with many details, we have five appendices. The first four give the details of each of the generalizations mentioned in section 3.6. The last appendix deals with the issue of $`T`$-duality on some of the compact models in Section 4. ## 2 Three-Branes and Seven-Branes at Abelian Orbifold Singularities In this section we introduce the basic formalism to compute the spectrum and interactions on the world-volume of D3-branes at $`\mathrm{IR}^6/ZZ_N`$ singularities <sup>2</sup><sup>2</sup>2Other cases, like $`ZZ_N\times ZZ_M`$ orbifolds can be studied analogously, and we will skip their discussion.. These have been discussed in , but our treatment is more general in that we allow the presence of D7-branes. We also discuss several aspects of these field theories, to be used in the remaining sections. Before entering the construction, we would like to explain our interest in placing the D3-branes on top of singularities. The reason is that D3-branes sitting at smooth points in the transverse dimensions lead to $`𝒩=4`$ supersymmetric field theories. The only known way to achieve chirality in this framework is to locate the D3-branes at singularities, the simplest examples being $`\mathrm{IR}^6/ZZ_N`$ orbifold singularities. Another important point in our approach is that we embed all gauge interactions in D3-branes. Our motivation for this is the fact that the appearance of SM fermions in three copies is difficult to achieve if color and weak interactions live in e.g. D3- and D7- branes, respectively (as will be manifest from the general spectra below). ### 2.1 Brane Spectrum We start by considering the case of a generic, not necessarily supersymmetric, singularity. Later on we discuss the specific $`𝒩=0`$ non supersymmetric and $`𝒩=1`$ supersymmetric realizations. Consider a set of $`n`$ D3-branes at a $`\mathrm{IR}^6/\mathrm{\Gamma }`$ singularity with $`\mathrm{\Gamma }SU(4)`$ where, for simplicity, we take $`\mathrm{\Gamma }=ZZ_N`$. Before the projection, the world-volume field theory on the D3-branes is a $`𝒩=4`$ supersymmetric $`U(n)`$ gauge theory. In $`𝒩=1`$ language, it contains $`U(n)`$ vector multiplets, and three adjoint chiral multiplets $`\mathrm{\Phi }^r`$, $`r=1,2,3`$, with interactions determined by the superpotential $`W`$ $`=`$ $`{\displaystyle \underset{r,s,t}{}}ϵ_{rst}\mathrm{Tr}(\mathrm{\Phi }^r\mathrm{\Phi }^s\mathrm{\Phi }^t)`$ (2.1) In terms of component fields, the theory contains $`U(n)`$ gauge bosons, four adjoint fermions transforming in the $`\mathrm{𝟒}`$ of the $`SU(4)_R`$ $`𝒩=4`$ R-symmetry group, and six adjoint real scalar fields transforming in the $`\mathrm{𝟔}`$. The $`ZZ_N`$ action on fermions is given by a matrix $$𝐑_\mathrm{𝟒}=\mathrm{diag}(e^{2\pi ia_1/N},e^{2\pi ia_2/N},e^{2\pi ia_3/N},e^{2\pi ia_4/N})$$ (2.2) with $`a_1+a_2+a_3+a_4=0\mathrm{mod}N`$. The action of $`ZZ_N`$ on scalars can be obtained from the definition of the action on the $`\mathrm{𝟒}`$, and it is given by the matrix $$𝐑_\mathrm{𝟔}=\mathrm{diag}(e^{2\pi ib_1/N},e^{2\pi ib_1/N},e^{2\pi ib_2/N},e^{2\pi ib_2/N},e^{2\pi ib_3/N},e^{2\pi ib_3/N})$$ (2.3) with $`b_1=a_2+a_3`$, $`b_2=a_1+a_3`$, $`b_3=a_1+a_2`$. Scalars can be complexified, the action on them being then given by $`𝐑_{esc}=\mathrm{diag}(e^{2\pi ib_1/N},e^{2\pi ib_2/N},e^{2\pi ib_3/N})`$. Notice that, since scalars have the interpretation of brane coordinates in the transverse space, eq. (2.3) defines the action of $`ZZ_N`$ on $`\mathrm{IR}^6`$ required to form the quotient $`\mathrm{IR}^6/ZZ_N`$. The action of the $`ZZ_N`$ generator $`\theta `$ must be embedded on the Chan-Paton indices. In order to be more specific we consider the general embedding given by the matrix $`\gamma _{\theta ,3}=\mathrm{diag}(\mathrm{I}_{n_0},e^{2\pi i/N}\mathrm{I}_{n_1},\mathrm{},e^{2\pi i(N1)/N}\mathrm{I}_{n_{N1}})`$ (2.4) where $`\mathrm{I}_{n_i}`$ is the $`n_i\times n_i`$ unit matrix, and $`_in_i=n`$. The theory on D3-branes at the $`\mathrm{IR}^6/ZZ_N`$ singularity is obtained by keeping the states invariant under the combined (geometrical plus Chan-Paton) $`ZZ_N`$ action . World-volume gauge bosons correspond to open string states in the NS sector, of the form $`\lambda \psi _{\frac{1}{2}}^\mu \|0`$, with $`\mu `$ along the D3-brane world-volume, and $`\lambda `$ the Chan-Paton wavefunction. The projection for gauge bosons is then given by $`\lambda =\gamma _{\theta ,3}\lambda \gamma _{\theta ,3}^1`$ (2.5) The projection for each of the three complex scalars, $`\lambda \mathrm{\Psi }_{\frac{1}{2}}^r\|0`$ (with $`r=1,2,3`$ labeling a complex plane transverse to the D3-brane) is $`\lambda =e^{2\pi ib_r/N}\gamma _{\theta ,3}\lambda \gamma _{\theta ,3}^1`$ (2.6) The four fermions in the D3-brane world-volume, labeled by $`\alpha =1,\mathrm{},4`$ are described by string states in the R sector, of the form $`\lambda \|s_1,s_2,s_3,s_4`$, with $`s_i=\pm \frac{1}{2}`$ and $`_is_i=\mathrm{odd}`$. The projection for left-handed fermions, $`s_4=\frac{1}{2}`$, leads to $$\lambda =e^{2\pi ia_\alpha /N}\gamma _{\theta ,3}\lambda \gamma _{\theta ,3}^1$$ (2.7) The final spectrum in the $`33`$ sector is $`\mathrm{Vectors}`$ $`{\displaystyle \underset{i=0}{\overset{N1}{}}}U(n_i)`$ $`\mathrm{Complex}\mathrm{Scalars}`$ $`{\displaystyle \underset{r=1}{\overset{3}{}}}{\displaystyle \underset{i=0}{\overset{N1}{}}}(n_i,\overline{n}_{ib_r})`$ $`\mathrm{Fermions}`$ $`{\displaystyle \underset{\alpha =1}{\overset{4}{}}}{\displaystyle \underset{i=0}{\overset{N1}{}}}(n_i,\overline{n}_{i+a_\alpha })`$ (2.8) where subindices will be understood modulo $`N`$ throughout the paper. The interactions are obtained by keeping the surviving fields in the interactions of the original $`𝒩=4`$ theory. Notice that the spectrum is, generically, non supersymmetric. Instead, when $`b_1+b_2+b_3=0`$, we have $`a_4=0`$ and the $`ZZ_N`$ action is in $`SU(3)`$. This case corresponds to a supersymmetric singularity. The fermions with $`\alpha =4`$ transforming in the adjoint representation of $`U(n_i)`$ become gauginos, while the other fermions transform in the same bifundamental representations as the complex scalars. The different fields fill out complete vector and chiral multiplets of $`𝒩=1`$ supersymmetry. In general, we would also like to include D7-branes in the configuration. Let us center on D7-branes transverse to the third complex plane $`Y_3`$, denoted D7<sub>3</sub>-branes in what follows. Open strings in the $`37_3`$ and $`7_33`$ sectors contribute new fields in the D3-brane world-volume. In the R sector, there are fermion zero modes in the NN and DD directions $`Y_4`$, $`Y_3`$. Such states are labeled by $`\lambda \|s_3;s_4`$, with $`s_3=s_4=\pm \frac{1}{2}`$, where $`s_4`$ defines the spacetime chirality. The projection for left-handed fermions $`s_4=\frac{1}{2}`$ are $`\lambda _{37_3}=e^{i\pi b_3/N}\gamma _{\theta ,3}\lambda \gamma _{\theta ,7_3}^1,\lambda _{7_33}=e^{i\pi b_3/N}\gamma _{\theta ,7_3}\lambda \gamma _{\theta ,3}^1`$ (2.9) Scalars arise from the NS sector, which contains fermion zero modes in the DN directions ($`Y_1`$, $`Y_2`$). States are labeled as $`\lambda \|s_1,s_2`$, with $`s_1=s_2=\pm 1/2`$. The projections for $`\lambda \|\frac{1}{2},\frac{1}{2}`$ are $`\lambda _{37_3}=e^{i\pi (b_1+b_2)/N}\gamma _{\theta ,3}\lambda \gamma _{\theta ,7_3}^1,\lambda _{7_33}=e^{i\pi (b_1+b_2)/N}\gamma _{\theta ,7_3}\lambda \gamma _{\theta ,3}^1`$ (2.10) We can give the resulting spectrum quite explicitly. Let us consider the Chan-Paton embedding $`\gamma _{\theta ,7_3}`$ $`=`$ $`\mathrm{diag}(\mathrm{I}_{u_0},e^{2\pi i/N}\mathrm{I}_{u_1},\mathrm{},e^{2\pi i(N1)/N}\mathrm{I}_{u_{N1}})\mathrm{for}b_3=\mathrm{even}`$ (2.11) $`\gamma _{\theta ,7_3}`$ $`=`$ $`\mathrm{diag}(e^{\pi i\frac{1}{N}}\mathrm{I}_{u_0},e^{2\pi i\frac{3}{N}}\mathrm{I}_{u_1},\mathrm{},e^{2\pi i\frac{2N1}{N}}\mathrm{I}_{u_{N1}})\mathrm{for}b_3=\mathrm{odd}`$ The resulting spectrum is $`\begin{array}{cccc}b_3=\mathrm{even}\hfill & \hfill & \mathrm{Fermions}\hfill & _{i=0}^{N1}[(n_i,\overline{u}_{i+\frac{1}{2}b_3})+(u_i,\overline{n}_{i+\frac{1}{2}b_3})]\hfill \\ & & \mathrm{Complex}\mathrm{Scalars}\hfill & _{i=0}^{N1}[(n_i,\overline{u}_{i\frac{1}{2}(b_1+b_2)})+(u_i,\overline{n}_{i\frac{1}{2}(b_1+b_2)})]\hfill \\ b_3=\mathrm{odd}\hfill & \hfill & \mathrm{Fermions}\hfill & _{i=0}^{N1}[(n_i,\overline{u}_{i+\frac{1}{2}(b_31)})+(u_i,\overline{n}_{i+\frac{1}{2}(b_3+1)})]\hfill \\ & & \mathrm{Complex}\mathrm{Scalars}\hfill & _{i=0}^{N1}[(n_i,\overline{u}_{i\frac{1}{2}(b_1+b_2+1)})+(u_i,\overline{n}_{i\frac{1}{2}(b_1+b_21)})]\hfill \end{array}`$ (2.16) The computation is identical for other D7<sub>r</sub>-branes, transverse to the $`r^{th}`$ complex plane, i.e. with world-volume defined by the equation $`Y_r=0`$. Notice that for a general twist, D7-branes with world-volume $`_{r=1}^3\beta _rY_r=0`$, with arbitrary complex coefficients $`\beta _r`$, are not consistent with the orbifold action (more precisely, they are not invariant under the orbifold action, and suitable $`ZZ_N`$ images should be included). For twists with several equal eigenvalues, such D7-branes are possible (see footnote 5). Notice that in the non-compact setting, fields in the 77 sector are non-dynamical from the viewpoint of the D3-brane world-volume field theory. For instance, 77 gauge groups correspond to global symmetries, and 77 scalars act as parameters of the D3-brane field theory. Only after compactification of the transverse space, as in the models discussed in Section 4, 77 fields become four-dimensional, and should be treated on an equal footing with 33 and 37, 73 fields. We conclude this section by restricting the above results to the case of singularities $`\mathrm{IC}^3/ZZ_N`$ preserving $`𝒩=1`$ supersymmetry on the D3-brane world-volume. That is, for $`a_4=0`$, and hence $`a_1+a_2+a_3=0\mathrm{mod}N`$. The spectrum is given by $`\begin{array}{cccc}\mathrm{𝟑𝟑}& \mathrm{Vector}\mathrm{mult}.& _{i=0}^{N1}U(n_i)& \\ & \mathrm{Chiral}\mathrm{mult}.& _{i=0}^{N1}_{r=1}^3(n_i,\overline{n}_{i+a_r})& \\ \mathrm{𝟑𝟕}_\mathrm{𝟑},\mathrm{𝟕}_\mathrm{𝟑}\mathrm{𝟑}& \mathrm{Chiral}\mathrm{mult}.& _{i=0}^{N1}[(n_i,\overline{u}_{i\frac{1}{2}a_3})+(u_i,\overline{n}_{i\frac{1}{2}a_3})]& a_3\mathrm{even}\\ & & _{i=0}^{N1}[(n_i,\overline{u}_{i\frac{1}{2}(a_3+1)})+(u_i,\overline{n}_{i\frac{1}{2}(a_31)})]& a_3\mathrm{odd}\end{array}`$ (2.21) We will denote $`\mathrm{\Phi }_{i,i+a_r}^r`$ the $`33`$ chiral multiplet in the representation $`(n_i,\overline{n}_{i+a_r})`$. We also denote (assuming $`a_3=\mathrm{even}`$ for concreteness) $`\mathrm{\Phi }_{i,i\frac{1}{2}a_3}^{(37_3)}`$, $`\mathrm{\Phi }_{i,i\frac{1}{2}a_3}^{(7_33)}`$ the $`37_3`$ and $`7_33`$ chiral multiplets in the $`(n_i,\overline{u}_{i\frac{1}{2}a_3})`$, $`(u_i,\overline{n}_{i\frac{1}{2}a_3})`$. With this notation, the interactions are encoded in the superpotential $`W`$ $`=`$ $`{\displaystyle \underset{r,s,t=1}{\overset{3}{}}}ϵ_{rst}\mathrm{Tr}(\mathrm{\Phi }_{i,i+a_r}^r\mathrm{\Phi }_{i+a_r,i+a_r+a_s}^s\mathrm{\Phi }_{i+a_r+a_s,i}^t)+{\displaystyle \underset{i=0}{\overset{N1}{}}}\mathrm{Tr}(\mathrm{\Phi }_{i,i+a_3}^3\mathrm{\Phi }_{i+a_3,i+\frac{1}{2}a_3}^{(37_3)}\mathrm{\Phi }_{i+\frac{1}{2}a_3,i}^{(7_33)})`$ ### 2.2 Anomaly and Tadpole Cancellation With the fermionic spectrum at hand we can proceed to compute non-abelian anomalies and establish the constraints for consistent anomaly-free theories. Moreover, such constraints can be rephrased in terms of twisted tadpole cancellation conditions . Let us address the computation of the non-abelian anomaly for $`SU(n_i)`$, in a case with D7<sub>r</sub>-branes, $`r=1,2,3`$. Let us assume $`b_r=\mathrm{even}`$ for concreteness, and denote by $`u_j^r`$ the number of entries with phase $`e^{2\pi ij/N}`$ in $`\gamma _{\theta ,7_r}`$. The $`SU(n_j)`$ non-abelian anomaly cancellation conditions are $`{\displaystyle \underset{\alpha =1}{\overset{4}{}}}(n_{i+a_\alpha }n_{ia_\alpha })+{\displaystyle \underset{r=1}{\overset{3}{}}}(u_{i+\frac{1}{2}b_r}^ru_{i\frac{1}{2}b_r}^r)=0`$ (2.22) These conditions are equivalent to the consistency conditions of the string theory configuration, namely cancellation of RR twisted tadpoles. To make this explicit, we use $`\begin{array}{ccc}n_j=\frac{1}{N}_{k=1}^Ne^{2\pi ikj/N}\mathrm{Tr}\gamma _{\theta ^k,3}\hfill & ;\hfill & u_j^r=\frac{1}{N}_{k=1}^Ne^{2\pi ikj/N}\mathrm{Tr}\gamma _{\theta ^k,7_r}\hfill \end{array}`$ (2.24) and substitute in (2.22). We obtain $`{\displaystyle \frac{2i}{N}}{\displaystyle \underset{k=1}{\overset{N}{}}}e^{2\pi ijk/N}[{\displaystyle \underset{\alpha =1}{\overset{4}{}}}\mathrm{sin}(2\pi ka_\alpha /N)\mathrm{Tr}\gamma _{\theta ^k,3}+{\displaystyle \underset{r=1}{\overset{3}{}}}\mathrm{sin}(\pi kb_r/2)\mathrm{Tr}\gamma _{\theta ^k,7_r}]=\mathrm{\hspace{0.17em}0}`$ (2.25) Using the identity $`{\displaystyle \underset{\alpha =1}{\overset{4}{}}}\mathrm{sin}(2\pi ka_\alpha /N)=4{\displaystyle \underset{r=1}{\overset{3}{}}}\mathrm{sin}(\pi kb_r/N)`$ (2.26) the Fourier-transformed anomaly cancellation condition is recast as $`[{\displaystyle \underset{r=1}{\overset{3}{}}}2\mathrm{sin}(\pi kb_r/N)]\mathrm{Tr}\gamma _{\theta ^k,3}+{\displaystyle \underset{r=1}{\overset{3}{}}}2\mathrm{sin}(\pi kb_r/N)\mathrm{Tr}\gamma _{\theta ^k,7_r}=\mathrm{\hspace{0.17em}0}`$ (2.27) These are in fact the twisted tadpole cancellation conditions. Notice that the contributions from the different disk diagrams are weighted by suitable sine factors, which arise from integration over the string center of mass in NN directions. ### 2.3 Structure of $`U(1)`$ Anomalies and Non-anomalous $`U(1)`$’s An important property of systems of D3-branes at singularities leading to chiral world-volume theories is the existence of mixed $`U(1)`$-nonabelian gauge anomalies. These field theory anomalies are canceled by a generalized Green-Schwarz mechanism mediated by closed string twisted modes (see for a similar mechanism in six dimensions). Consider the generically non-supersymmetric field theory constructed from D3- and D7<sub>r</sub>-branes at a $`ZZ_N`$ singularity, with spectrum given in (2.8), (2.16). Assuming $`b_r=\mathrm{even}`$ for concreteness, the mixed anomaly between the $`j^{th}`$ $`U(1)`$ (that within $`U(n_j)`$) and $`SU(n_l)`$ is $`A_{jl}={\displaystyle \frac{1}{2}}n_j{\displaystyle \underset{\alpha =1}{\overset{4}{}}}(\delta _{l,j+a_\alpha }\delta _{l,ja_\alpha })+{\displaystyle \frac{1}{2}}\delta _{j,l}[{\displaystyle \underset{\alpha =1}{\overset{4}{}}}(n_{j+a_\alpha }n_{ja_\alpha })+{\displaystyle \underset{r=1}{\overset{3}{}}}(u_{j\frac{1}{2}b_r}^ru_{j+\frac{1}{2}b_r}^r)]`$ The anomaly is not present if $`n_j`$ or $`n_l`$ vanish. After using the cancellation of cubic non-abelian anomalies (2.22) (i.e. the tadpole cancellation conditions), the remaining piece is given by $$A_{jl}=\frac{1}{2}n_j\underset{\alpha =1}{\overset{4}{}}(\delta _{l,j+a_\alpha }\delta _{l,ja_\alpha })$$ (2.28) This anomaly can be rewritten as $`A_{jl}`$ $`=`$ $`{\displaystyle \frac{i}{2N}}{\displaystyle \underset{k=1}{\overset{N1}{}}}[n_j\mathrm{exp}(2i\pi kj/N)\mathrm{exp}(2i\pi kl/N){\displaystyle \underset{r=1}{\overset{3}{}}}2\mathrm{sin}(\pi kb_r/N)]`$ (2.29) which makes the factorized structure of the anomaly explicit. The anomaly is canceled by exchange of closed string twisted modes , which have suitable couplings to the gauge fields on the D-brane world-volume . An important fact is that anomalous $`U(1)`$’s get a tree-level mass of the order of the string scale , and therefore do not appear in the low-energy field theory dynamics on the D3-brane world-volume. It is therefore interesting to discuss the existence of non-anomalous $`U(1)`$’s and their structure. Concretely, we consider linear combinations of the $`U(1)`$ generators $$Q_c=\underset{j=0}{\overset{N1}{}}c_j\frac{Q_{n_j}}{n_j}$$ (2.30) (we take $`c_j=0`$ if $`n_j=0`$). The condition for $`U(1)`$’s free of $`Q_cSU(n_l)^2`$ anomalies reads $$\frac{1}{2}\underset{\alpha =1}{\overset{4}{}}\underset{j=0}{\overset{N1}{}}c_j(\delta _{l,j+a_\alpha }\delta _{l,ja_\alpha })=\underset{\alpha =1}{\overset{4}{}}(c_{la_\alpha }c_{l+a_\alpha })=0$$ (2.31) for all $`l=0,\mathrm{}N1`$. It is clear that $`c_j=const.`$ lead to an anomaly free combination $`Q_{diag.}`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{N1}{}}}{\displaystyle \frac{Q_{n_i}}{n_i}}`$ (2.32) This generically non-anomalous $`U(1)`$ plays a prominent role in the realistic models of Section 3. In general (2.31) gives $`N1`$ independent conditions for the $`N`$ unknowns $`c_j`$, and (2.32) is the only non-trivial solution (as long as no $`n_j`$ vanishes). In certain cases, however, the number of independent equations may be smaller, and additional non-anomalous $`U(1)`$’s appear. In order to compute the number of independent equations we rewrite the unknowns $`c_j`$ in terms of the new variables $`r_k=_{j=0}^{N1}e^{2i\pi kj/N}c_j`$. The original variables are given by $$c_j=\frac{1}{N}\underset{k=0}{\overset{N1}{}}e^{2\pi ikj/N}r_k.$$ (2.33) By replacing in (2.31) above we find $`_{\alpha =1}^4\mathrm{sin}(2\pi ka_\alpha /N)r_k=0`$ and, by using the identity (2.26) we obtain $$[\underset{r=1}{\overset{3}{}}\mathrm{sin}(\pi kb_r/N)]r_k=0$$ (2.34) Thus, we have managed to diagonalize the matrix corresponding to eq. (2.31). We see that, besides the diagonal solution (corresponding to $`r_i=0`$ for $`i=0.\mathrm{},N1`$) other non trivial solutions are possible whenever there exist twists $`\theta ^k`$ that leave, at least, one unrotated direction. In other words, the rank of the matrix is given by the number of twists rotating all complex planes. It is possible to describe explicitly the non-anomalous $`U(1)`$’s, as follows. Let us consider a $`ZZ_M`$ subgroup, generated by a certain twist $`\theta ^p`$, leaving e.g. the third complex plane fixed, hence $`pa_1=pa_2`$ and $`pa_3=pa_4`$. For each value of $`I=1,\mathrm{},M1`$, there is one non-anomalous $`U(1)`$, defined by $`c_j=\delta _{pj,I}`$. This satisfies the condition (2.31) by a cancellation of contributions from different values of $`\alpha `$. Hence we obtain one additional non-anomalous $`U(1)`$ per twist leaving some complex plane fixed <sup>3</sup><sup>3</sup>3There are arguments suggesting these additional non-anomalous $`U(1)`$’s are nevertheless massive due to their mixing with closed string twisted modes . This observation however, would not change our analysis in the following sections, since our model building involves only the diagonal combination (2.32), for which such mixing vanishes.. Let us consider some explicit examples. For instance, for $`ZZ_3`$ with twist $`v=\frac{1}{3}(1,1,2)`$, equation (2.31) leads to $`b_0=b_1=b_2`$ and the only anomaly free combination is (2.32). Consequently, it is not possible to have an anomaly free $`U(1)`$, unless all $`n_i0`$. This will always be the case for $`ZZ_N`$ orbifold actions (supersymmetric or not) without twists with fixed planes. Consequently, and as we discuss further in section 3, if we are interested in gauge groups similar to the Standard Model one, with an anomaly free (hypercharge) $`U(1)`$, then all $`n_i`$’s in (2.4) should be non vanishing. The situation is different when subgroups with fixed planes exist. For instance, consider the $`ZZ_6\frac{1}{6}(1,1,2)`$ example. The $`\mathrm{𝟑𝟑}`$ spectrum reads $`U(n_0)\times U(n_1)\times U(n_2)\times U(n_3)\times U(n_4)\times U(n_5)\times U(n_6)`$ $`2[(n_0,\overline{n}_1)+(n_1,\overline{n}_2)+(n_2,\overline{n}_3)+(n_3,\overline{n}_4)+(n_4,\overline{n}_5)+(n_5,\overline{n}_0)+(n_3,\overline{n}_0)]`$ (2.35) $`+(n_0,\overline{n}_4)+(n_2,\overline{n}_0)+(n_4,\overline{n}_2)+(n_1,\overline{n}_5)+(n_3,\overline{n}_1)+(n_5,\overline{n}_3)`$ (2.36) and the anomaly matrix (2.28) is given by $$2T_{ij}^{\alpha \beta }=\left(\begin{array}{cccccc}0& 2n_0& n_0& 0& n_0& 2n_0\\ 2n_1& 0& 2n_1& n_1& 0& n_1\\ n_2& 2n_2& 0& 2n_2& n_2& 0\\ 0& n_3& 2n_3& 0& 2n_3& n_3\\ n_4& 0& n_4& 2n_4& 0& 2n_4\\ 2n_5& n_5& 0& n_5& 2n_5& 0\end{array}\right)$$ (2.37) The search for anomaly-free combinations leads to eq.(2.31) above, which has two non-trivial independent solutions. In fact, (2.33) becomes $`c_i=r_0+(1)^ir_3`$, namely $`c_0=c_2=c_4`$ and $`c_1=c_3=c_5`$ indicating that two anomaly-free abelian factors can be present. In particular, by choosing $`n_1=n_3=n_5=0`$ and $`n_4=3,n_2=2,n_0=1`$ we obtain a field theory with Standard Model group $`SU(3)\times SU(2)\times U(1)`$ and one generation $`(\mathrm{𝟑},\overline{\mathrm{𝟐}})_{1/6}+(\mathrm{𝟏},\mathrm{𝟐})_{1/2}+(\overline{\mathrm{𝟑}},\mathrm{𝟏})_{2/3}`$, with subscripts giving $`U(1)`$ charges (fields in a suitable 37 sector should complete this to a full SM generation). ## 3 Particle Models from Branes at Singularities In this Section we discuss the embedding of the Standard Model (and related Left-Right symmetric extensions) in systems of D3-branes at $`\mathrm{IC}^3/ZZ_N`$ singularities, with $`ZZ_NSU(3)`$. We also discuss possible extensions to more general cases. ### 3.1 Number of Generations and Hypercharge Let us start by recalling the structure of the field theory on D3-branes at a $`\mathrm{IC}^3/ZZ_N`$ singularity, defined by the twist $`v=(a_1,a_2,a_3)/N`$. It has $`𝒩=1`$ supersymmetry, and the following field content: there are vector multiplets with gauge group $`_{i=0}^{N1}U(n_i)`$, and $`3N`$ $`𝒩=1`$ chiral multiplets $`\mathrm{\Phi }_{i,i+a_r}^r`$, $`i=0,\mathrm{},N1`$, $`r=1,2,3`$, transforming in the representation $`(n_i,\overline{n}_{i+a_r})`$. The interactions are encoded in the superpotential $`W=ϵ_{rst}\mathrm{Tr}(\mathrm{\Phi }_{i,i+a_r}^r\mathrm{\Phi }_{i+a_r,i+a_r+a_s}^s\mathrm{\Phi }_{i+a_r+a_s,i}^t).`$ (3.1) In general, the configurations will also include D7-branes, but for the moment we center of general features in the 33 sector. We are interested in constructing theories similar to the standard model or some simple extension thereof. In particular, we will be interested in constructing models which explicitly contain an $`SU(3)\times SU(2)`$ factor, to account for color and weak interactions. Besides simplicity, this choice has the additional advantage (discussed in detail below) that a non-anomalous $`U(1)`$ leading to correct hypercharge assignments arises naturally. Hence we consider models in which two factors, which without loss of generality we take $`U(n_0)`$ and $`U(n_j)`$, are actually $`U(3)`$, $`U(2)`$. Since the matter content contains only bi-fundamental representations, the number of generations is given by the number of left-handed quarks, i.e. fields in the representation $`(\mathrm{𝟑},\mathrm{𝟐})`$. This is given by the number of twist eigenvalues $`a_r`$ equal to $`j`$ (minus the number of $`a_r`$ equal to $`j`$), which is obviously at most three, this maximum value occurring only for the $`\mathrm{IC}^3/ZZ_3`$ singularity, with twist $`v=(1,1,2)/3`$. For this reason, this singularity will play a prominent role in our forthcoming models. Before centering on that concrete case, it will be useful to analyze the issue of hypercharge. In order to obtain a theory with standard model gauge group, one needs the presence of at least one non-anomalous $`U(1)`$ to play the role of hypercharge. Happily, our analysis in section 2.3 has shown that, as long as no $`n_i`$ vanishes, D3-branes at singularities always have at least one non-anomalous $`U(1)`$ generated by $`Q_{diag}`$ in (2.32). Therefore, a possibility to obtain the standard model gauge group, with no additional non-abelian factors would be to consider models with group <sup>4</sup><sup>4</sup>4Such choice for the Chan-Paton embedding for D3-branes leads in general to non-vanishing tadpoles, which can be canceled by the introduction of D7-branes. We leave their discussion for the explicit examples below. $`U(3)\times U(2)\times U(1)^{N2}`$. In the generic case, only the diagonal combination $`Q_{diag}`$ $`=`$ $`\left({\displaystyle \frac{1}{3}}Q_3+{\displaystyle \frac{1}{2}}Q_2+{\displaystyle \underset{s=1}{\overset{N2}{}}}Q_1^{(s)}\right)`$ (3.2) will be non-anomalous (the overall minus sign is included for later convenience). In a generic orbifold all other $`N1`$ additional $`U(1)`$ factors will be anomalous and therefore massive, with mass of the order of the string scale. Of course, the fact that we have a non-anomalous $`U(1)`$ does not guarantee it has the right properties of hypercharge. Quite surprisingly this is precisely the case for (3.2). For instance, fields transforming in the $`(\mathrm{𝟑},\mathrm{𝟐})`$ representation have $`Q_{diag}`$ charge $`\frac{1}{3}+\frac{1}{2}=\frac{1}{6}`$, as corresponds to left-handed quarks. Fields transforming in the $`(\mathrm{𝟑},\mathrm{𝟏})`$ (necessarily with charge $`1`$ under one of the $`Q_1^{(s)}`$ generators) have a $`Q_{diag.}`$ charge $`\frac{1}{3}+1=\frac{2}{3}`$, as corresponds to right-handed U quarks, etc. Analysis of the complete spectrum requires information about the D7-brane sector, and is postponed until the construction of explicit examples in section 3.3. Notice that correct hypercharge assignments would not be obtained had our starting point been e.g. $`SU(4)\times SU(2)`$, hence our interest in the $`SU(3)\times SU(2)`$ structure. It is worth noticing that normalization of this hypercharge $`U(1)`$ depends on $`N`$. In fact, by normalizing $`U(n)`$ generators such that $`\mathrm{Tr}T_a^2=\frac{1}{2}`$ the normalization of $`Y`$ generator is fixed to be $$k_1=5/3+2(N2)$$ (3.3) This amounts to a dependence on $`N`$ in the Weinberg angle, namely $$\mathrm{sin}^2\theta _W=\frac{1}{k_1+1}=\frac{3}{6N4}$$ (3.4) Thus the weak angle decreases as $`N`$ increases. Notice that the $`SU(5)`$ result $`3/8`$ is only obtained for a $`Z_2`$ singularity. However in that case the (33) spectrum is necessarily vector-like and hence one cannot reproduce the SM spectrum. We will also be interested in constructing left-right symmetric extensions of the standard model. In particular, we consider gauge groups with a factor $`SU(3)\times SU(2)_L\times SU(2)_R`$, which are obtained by choosing suitable values for three of the entries $`n_i`$ in the Chan-Paton embedding. As above, the corresponding tadpole must be canceled by additional D7-branes, whose details we postpone for the moment. The number of generations is again given by the number of representations $`(\mathrm{𝟑},\mathrm{𝟐},\mathrm{𝟏})`$, and is equal to three only for the $`ZZ_3`$ orbifold. To reproduce hypercharge after the breaking of the right-handed $`SU(2)`$ factor, an essential ingredient is the existence of a non-anomalous $`(BL)`$ $`U(1)`$ in the theory. In order to obtain at least one non-anomalous $`U(1)`$ in the D3-branes, we are led to consider models with group $`U(3)\times U(2)\times U(2)\times U(1)^{N3}`$. Generically, only the diagonal combination $`Q_{diag}`$ $`=`$ $`2\left({\displaystyle \frac{1}{3}}Q^{(3)}+{\displaystyle \frac{1}{2}}Q^{(2_L)}+{\displaystyle \frac{1}{2}}Q^{(2_R)}+{\displaystyle \underset{s=1}{\overset{N3}{}}}Q_s^{(1)}\right)`$ (3.5) (the overall factor is included for convenience) is non-anomalous. Interestingly, the charges under this non-anomalous $`U(1)`$ turn out to have the correct $`BL`$ structure. For instance, fields transforming in the $`(\mathrm{𝟑},\mathrm{𝟐},\mathrm{𝟏})`$ or $`(\mathrm{𝟑},\mathrm{𝟏},\mathrm{𝟐})`$ representations have $`Q_{diag}`$ charge is $`2(\frac{1}{3}\frac{1}{2})=\frac{1}{3}`$, correct for quark fields. Again, the discussion for the complete spectrum is postponed to the explicit examples in section 3.4. Notice, as above, that normalization of $`BL`$ generator is $`N`$ dependent and leads to $`k_{BL}=8/3+8(N2)`$. Since hypercharge is given by $`Y=T_R^3+Q_{BL}`$ (with $`T_R^3`$ the diagonal generator of $`SU(2)_R`$) the values (3.3) and, thus, (3.4) are reobtained for hypercharge normalization and Weinberg angle. We find it is quite remarkable that the seemingly complicated hypercharge structure in the standard model is easily accomplished by the structure of the diagonal $`U(1)`$ in this class of orbifold models. We conclude by remarking that in cases with additional non-anomalous $`U(1)`$’s $`Q_c`$ (2.30), they could be used as hypercharge or $`BL`$ generators, as long as the $`U(1)`$ factors in $`U(3)`$ and $`U(2)`$ belong to the corresponding linear combination in $`Q_c`$ (as in the $`ZZ_6`$ example in section 2.3). However, since the presence of these $`U(1)`$’s is not generic, we will not analyze this possibility in detail. Moreover, they are not present in the case of $`ZZ_3`$ singularity, which is the only candidate to produce three-generation models. ### 3.2 Generalities for $`\mathrm{IC}^3/ZZ_3`$ In the following we construct some explicit examples of standard model or left-right symmetric theories based on the $`\mathrm{IC}^3/ZZ_3`$ singularity. This is the most attractive case, since it leads naturally to three-family models. It also illustrates the general technology involved in model building using branes at singularities. Consider a set of D3-branes and D7<sub>r</sub>-branes at a $`\mathrm{IC}^3/ZZ_3`$ orbifold singularity, generated by the twist $`v=\frac{1}{3}(1,1,2)`$. Its action on the Chan-Paton factors is in general given by the matrices $`\begin{array}{ccc}\gamma _{\theta ,3}=\mathrm{diag}(\mathrm{I}_{n_0},\alpha \mathrm{I}_{n_2},\alpha ^2\mathrm{I}_{n_3})\hfill & ;& \gamma _{\theta ,7_1}=\mathrm{diag}(\mathrm{I}_{u_0^1},\alpha \mathrm{I}_{u_2^1},\alpha ^2\mathrm{I}_{u_3^1})\hfill \\ \gamma _{\theta ,7_3}=\mathrm{diag}(\mathrm{I}_{u_0^3},\alpha \mathrm{I}_{u_1^3},\alpha ^2\mathrm{I}_{u_2^3})\hfill & ;& \gamma _{\theta ,7_2}=\mathrm{diag}(\mathrm{I}_{u_0^2},\alpha \mathrm{I}_{u_2^2},\alpha ^2\mathrm{I}_{u_3^2})\hfill \end{array}`$ (3.8) with $`\alpha =e^{2\pi i/3}`$. The notation, slightly different from that in section 2.1, is more convenient for $`\mathrm{IC}^3/ZZ_3`$. The full spectrum is given by $`\begin{array}{cc}\mathrm{𝟑𝟑}& U(n_0)\times U(n_1)\times U(n_2)\\ & 3[(n_0,\overline{n}_1)+(n_1,\overline{n}_2)+(n_2,\overline{n}_0)]\\ \mathrm{𝟑𝟕}_𝐫,\mathrm{𝟕}_𝐫\mathrm{𝟑}& (n_0,\overline{u}_{}^{r}{}_{1}{}^{})+(n_1,\overline{u}_{}^{r}{}_{2}{}^{})+(n_2,\overline{u}_{}^{r}{}_{0}{}^{})+\\ & +(u_0^r,\overline{n}_1)+(u_1^r,\overline{n}_2)+(u_2^r,\overline{n}_0)\end{array}`$ (3.13) The superpotential <sup>5</sup><sup>5</sup>5It is possible to consider the generic case of D7<sup>β</sup>-branes, with world-volume defined by $`_r\beta _rY_r=0`$, which preserve the $`𝒩=1`$ supersymmetry of the configuration for arbitrary complex $`\beta _r`$ . The 37<sup>β</sup>, 7<sup>β</sup>3 spectra are as above, but the superpotential is $`W=_i_r\beta _r\mathrm{Tr}(\mathrm{\Phi }^r\mathrm{\Phi }^{37^\beta }\mathrm{\Phi }^{7^\beta 3})`$, with fields from a single mixed sector coupling to 33 fields from all complex planes. terms are $`W={\displaystyle \underset{i=0}{\overset{2}{}}}{\displaystyle \underset{r,s,t=1}{\overset{3}{}}}ϵ_{rst}\mathrm{Tr}(\mathrm{\Phi }_{i,i+1}^r\mathrm{\Phi }_{i+1,i+2}^s\mathrm{\Phi }_{i+2,i}^t)+{\displaystyle \underset{i=0}{\overset{2}{}}}{\displaystyle \underset{r=1}{\overset{3}{}}}\mathrm{Tr}(\mathrm{\Phi }_{i,i+1}^r\mathrm{\Phi }_{i+1,i+2}^{37_r}\mathrm{\Phi }_{i+2,i}^{7_r3})`$ (3.14) The twisted tadpole cancellation conditions are $`\mathrm{Tr}\gamma _{\theta ,7_3}\mathrm{Tr}\gamma _{\theta ,7_1}\mathrm{Tr}\gamma _{\theta ,7_2}+3\mathrm{T}\mathrm{r}\gamma _{\theta ,3}=0`$ (3.15) where the relative signs, coming from the sine prefactors, cancel those in the definitions in (3.8), hence it is consistent to ignore both. Eqs (3.15) are equivalent to the non-abelian anomaly cancellation conditions. ### 3.3 Standard Model and Branes at $`\mathrm{IC}^3/ZZ_3`$ Singularity Following the general arguments in section 3.1, the strategy to obtain a field theory with standard model gauge group from the $`ZZ_3`$ singularity is to choose a D3-brane Chan-Paton embedding $`\gamma _{\theta ,3}`$ $`=`$ $`\mathrm{diag}(\mathrm{I}_3,\alpha \mathrm{I}_2,\alpha ^2\mathrm{I}_1)`$ (3.16) The simplest way to satisfy the tadpole conditions (3.15) is to introduce only one set of D7-branes, e.g. D7<sub>3</sub>-branes, with Chan-Paton embedding $`u_0^3=0`$, $`u_0^1=3`$, $`u_0^2=6`$. The gauge group on the D3-branes is $`U(3)\times U(2)\times U(1)`$, whereas in the D7<sub>3</sub>-branes is $`U(3)\times U(6)`$ on each. Note that, before compactification, the latter behave as global symmetries in the worldvolume of the D3-branes. The D7<sub>3</sub>-branes group can be further broken by global effects, since the corresponding branes are extended along some internal dimensions. An alternative procedure to obtain a smaller group on the D7-branes is to use all three kinds of D7-branes, as depicted in Figure 2. For instance, a very symmetrical choice consistent with (3.16) is $`u_0^r=0`$, $`u_1^r=1`$, $`u_1^r=2`$, for $`r=1,2,3`$. Each kind of D7-brane then carries a $`U(1)\times U(2)`$ group. The spectrum for this latter model is given in table 1 (for later convenience we have also included states in the $`7_r7_r`$ sectors; their computation is analogous to the computation of the 33 sector). In the last column we give the charges under the anomaly-free combination (3.2): $$Y=\left(\frac{1}{3}Q_3+\frac{1}{2}Q_2+Q_1\right)$$ (3.17) As promised, it gives the correct hypercharge assignments for standard model fields. A pictorial representation of this type of models is given in Figure 3. We find it remarkable that such a simple configuration produces a spectrum so close to that of the standard model. In particular, we find encouraging the elegant appearance of hypercharge within this framework, as the only linear combination (2.32) of $`U(1)`$ generators which is naturally free of anomalies in systems of D3-branes at orbifold singularities. There is another interesting advantage in the fact that hypercharge arises exclusively from the 33 sector. Notice that it allows for the fields in the 77 sector to acquire nonvanishing vev’s <sup>6</sup><sup>6</sup>6Strictly speaking, 77 fields are not dynamical from the point of view of the D3-brane field theory in the non-compact context. This comment should be regarded as applied to compact models with a local behaviour given by the above system of D3- and D7<sub>r</sub>-branes at a $`ZZ_3`$ singularity, like those in Section 5. without breaking hypercharge. These vevs can be used to further break the $`\mathrm{𝟕𝟕}`$ gauge groups, and produce masses for the extra triplets in the 37 sectors. Then the only remaining light triplets are those of the standard model. Notice also that the same argument does not apply to the doublets in 37 sectors, which remain massless even after the 77 fields acquire vevs. Hence the three families of leptons remain light. This behavior is reminiscent of the models in , basically because (see Appendix E) their realistic sector has (in a T-dual version) a local structure very similar to our non-compact $`ZZ_3`$ singularity. The model constructed above, once embedded in a global context, may provide the simplest semirealistic string compactifications ever built. Indeed, in Section 4 we will provide explicit compact examples of this kind. Let us once again emphasize that, however, many properties of the resulting theory will be independent of the particular global structure used to achieve the compactification, and can be studied in the non-compact version presented above, as we do in Section 5. ### 3.4 Left-Right Symmetric Models and the $`\mathrm{IC}^3/ZZ_3`$ Singularity One may use a similar approach to construct three-generation models with left-right symmetric gauge group. Following the arguments in section 3.1, we consider the D3-brane Chan-Paton embedding $`\gamma _{\theta ,3}`$ $`=`$ $`\mathrm{diag}(\mathrm{I}_3,\alpha \mathrm{I}_2,\alpha ^2\mathrm{I}_2)`$ (3.18) The corresponding tadpoles can be canceled for instance by D7<sub>r</sub>-branes, $`r=1,2,3`$ with the symmetric choice $`u_0^r=0`$, $`u_1^r=u_2^r=1`$. The gauge group on D3-branes is $`U(3)\times U(2)_L\times U(2)_R`$, while each set of D7<sub>r</sub>-branes contains $`U(1)^2`$. As explained above, the combination (3.5) $$Q_{BL}=2\left(\frac{1}{3}Q_3+\frac{1}{2}Q_L+\frac{1}{2}Q_R\right)$$ (3.19) is non-anomalous, and in fact behaves as $`BL`$. The spectrum for this model, with the relevant $`U(1)`$ quantum numbers is given in table 2. A pictorial representation of this type of models is given in Figure 4. We can see that the triplets from the 37<sub>r</sub> sectors can become massive after the singlets of the $`\mathrm{𝟕}_𝐫\mathrm{𝟕}_𝐫`$ sector acquire a nonvanishing vev, leaving a light spectrum really close to left-right theories considered in phenomenological model-building. Left-right symmetric models have several interesting properties that were recently emphasized in , (besides the original discussions in ). As a practical advantage, it is the simplest to construct, and it allows to distinguish easily the Higgs fields from the leptons, since they transform differently under the gauge group. Phenomenologically, it allows to have gauge unification at the intermediate scale. Most of the properties studied in , regarding gauge unification, proton stability, fermion masses, etc, are inherited by theories based on the above construction. We will leave the phenomenological discussion to Section 5. We conclude this section with a comment that applies to models both of SM and LR type. A familiar property of D3-branes at singularities is that sets of D3-branes with traceless Chan-Paton embedding can combine and move off the singularity into the bulk. From the viewpoint of the world-volume field theory this appears as a flat direction along which the gauge group is partially broken, and which is parametrized by a modulus associated to the bulk position of the set of branes. This process is possible in the models discussed above, due to the existence of sets of D3-branes with Chan-Paton embedding $`\mathrm{diag}(1,\alpha ,\alpha ^2)`$ (one such set for SM type configurations, and two for LR type models). In fact, it is easy to find the corresponding flat directions in the field theory we computed. This fact will be important in the discussion of some compactified models in Section 5. ### 3.5 Non-Supersymmetric Singularities In the supersymmetric case the $`\mathrm{IC}^3/ZZ_3`$ singularity is singled out. It is the only case leading to three-generation models, due to the fact that three complex directions are equally twisted. The situation is quite different for non-supersymmetric singularities. In fact, by looking at the 33 fermionic spectrum (2.8), namely $`_{\alpha =1}^4_{i=0}^{N1}(n_i,\overline{n}_{i+a_\alpha })`$, we notice that for $`a_1=a_2=a_3=a`$ (then $`a_4=3a`$) we have $$3(n_i,\overline{n}_{i+a})+(n_i,\overline{n}_{i3a})$$ (3.20) and therefore a potential triplication. This singularities are therefore well-suited for model building of non-supersymmetric realistic spectra. We would like to point out that, despite the lack of supersymmetry, these models do not contain tachyons, neither in open nor in closed (untwisted or twisted) string sectors. Without loss of generality we can choose $`a=1`$ <sup>7</sup><sup>7</sup>7For a $`ZZ_N`$ singularity, we must have $`\mathrm{gcd}(a,N)=1`$, hence there exists $`p`$ such that $`pa=1\mathrm{mod}N`$, and we may choose $`\theta ^p`$ to generate $`ZZ_N`$. This only implies a harmless redefinition on the $`n_i`$’s in $`\gamma _{\theta ,3}`$. corresponding to $`ZZ_N`$ twist given by $`(1,1,1,3)/N`$. Thus, we observe that $`N=2`$ leads to non-chiral theories, $`N=4`$ to four-generation models, while each model with $`N5`$ leads to three-generations models. For instance, by choosing $`n_0=3,n_1=2`$ and $`n_i=1`$ for $`i=3,\mathrm{},N1`$, a Standard Model gauge group $`SU(3)\times SU(2)\times U(1)_Y\times U(1)^{N1}`$ is found with $`\mathrm{𝟑𝟑}`$ fermions transforming as $$3[(\mathrm{𝟑},\mathrm{𝟐})_{1/6}+(\overline{\mathrm{𝟑}},\mathrm{𝟏})_{2/3}+(\mathrm{𝟏},\mathrm{𝟐})_{1/2}]+(N3)\mathrm{singlets}.$$ (3.21) Thus, we obtain three generations of left-handed quarks and right-handed U type quarks. Since such a matter content is anomalous, extra contributions coming from D7-branes are expected to complete the spectrum to the three generations of SM quarks and leptons. Notice that presence of D7-branes only produce fundamentals of $`\mathrm{𝟑𝟑}`$ groups in $`\mathrm{𝟑𝟕}+\mathrm{𝟕𝟑}`$ sectors and therefore the number of generations is not altered. Interestingly enough, the correct SM hypercharge assignments above correspond to the anomaly-free diagonal $`U(1)`$ combination (2.32) $$Y=Q_{diag}=\left(\frac{1}{3}Q_{n_0}+\frac{1}{2}Q_{n_1}+\underset{j=3}{\overset{N1}{}}Q_{n_j}\right)$$ (3.22) Hence hypercharge can arise by the same mechanism discussed in section 3.1 for supersymmetric models. It is interesting to study the Weinberg angle prediction for these models. The computation follows that in section 3.1, leading to the result (3.4), which gives $`\mathrm{sin}^2\theta _W=0.214,0.115\mathrm{}`$ for $`N=3,N=5,\mathrm{}`$ respectively. Hence we see that, even though many non-supersymmetric singularities lead to three-generation models, in general they yield too low values of the Weinberg angle. It is interesting that the value gets worse as the order of the singularity increases, suggesting that simple configurations are better suited to reproduce realistic particle models. A $`ZZ_5`$ example As a concrete example of the above discussion let us consider the $`ZZ_5`$ singularity acting on the $`\mathrm{𝟒}`$ with twist $`(a_1,a_2,a_3,a_4)=(1,1,1,3)`$. Hence the action on the $`\mathrm{𝟔}`$ is given by $`(b_1,b_2,b_3)=(2,2,2)`$. The general D3-brane Chan-Paton matrix has the form $$\gamma _{\theta ,3}=\mathrm{diag}(\mathrm{I}_{n_0},\alpha \mathrm{I}_{n_1},\alpha ^2\mathrm{I}_{n_2},\alpha ^3\mathrm{I}_{n_3},\alpha ^4\mathrm{I}_{n_4})$$ (3.23) with $`\alpha =e^{2\pi i/5}`$. Anomaly/tadpole cancellation conditions require $`4(\alpha \alpha ^4)\mathrm{Tr}\gamma _{\theta ,3}+{\displaystyle \underset{r=1}{\overset{3}{}}}\gamma _{\theta ,7_r}=0`$ (3.24) and, therefore, D7-branes must be added. Let us consider the case with only D7<sub>3</sub>-branes, with Chan-Paton action $$\gamma _{\theta ,7_3}=\mathrm{diag}(\mathrm{I}_{u_0},\alpha \mathrm{I}_{u_1},\alpha ^2\mathrm{I}_{u_2},\alpha ^3\mathrm{I}_{u_3},\alpha ^4\mathrm{I}_{u_4})$$ (3.25) The massless spectrum reads $`\mathrm{𝟑𝟑}`$ $`\mathrm{Vectors}`$ $`U(n_0)\times U(n_1)\times U(n_2)\times U(n_3)\times U(n_4)`$ $`\mathrm{Fermions}`$ $`3[(n_0,\overline{n}_1)+(n_1,\overline{n}_2)+(n_2,\overline{n}_3)+(n_3,\overline{n}_4)+(n_4,\overline{n}_0)]+`$ $`(n_0,\overline{n}_2)+(n_1,\overline{n}_3)+(n_2,\overline{n}_4)+(n_3,\overline{n}_0)+(n_4,\overline{n}_1)]`$ $`\mathrm{Cmplx}.\mathrm{Sc}.`$ $`3[(n_0,\overline{n}_3)+(n_1,\overline{n}_4)+(n_2,\overline{n}_0)+(n_3,\overline{n}_1)+(n_4,\overline{n}_2),]`$ $`\mathrm{𝟑𝟕}_\mathrm{𝟑}`$ $`\mathrm{Fermions}`$ $`(n_0,\overline{u}_1)+(n_1,\overline{u}_2)+(n_2,\overline{u}_3)+(n_3,\overline{u}_4)+(n_4,\overline{u}_0)`$ $`\mathrm{Cmplx}.\mathrm{Sc}.`$ $`(n_0,\overline{u}_3)+(n_1,\overline{u}_4)+(n_2,\overline{u}_0)+(n_3,\overline{u}_1)+(n_4,\overline{u}_2)`$ $`\mathrm{𝟕}_\mathrm{𝟑}\mathrm{𝟑}`$ $`\mathrm{Fermions}`$ $`(u_0,\overline{n}_1)+(u_1,\overline{n}_2)+(u_2,\overline{n}_3)+(u_3,\overline{n}_4)+(u_4,\overline{n}_0)`$ $`\mathrm{Cmplx}.\mathrm{Sc}.`$ $`(u_0,\overline{n}_3)+(u_1,\overline{n}_4)+(u_2,\overline{n}_0)+(u_3,\overline{n}_1)+(u_4,\overline{n}_2)`$ where $`n_i`$, $`u_i`$’s are constrained by (3.24). To be more specific, let us build up a Left-Right model, by choosing $`n_0=3,n_1=n_4=2`$ and $`n_2=n_3=1`$, which leads to a gauge group $`SU(3)\times SU(2)\times SU(2)\times U(1)_{BL}\times [U(1)^5]`$. The choice $`u_0=0,u_1=3,u_2=u_3=7,u_4=3`$ ensures tadpole cancellation. The $`BL`$ charge is provided by the anomaly-free diagonal combination (2.32). From the generic spectrum above we find 33 fermions transform, under LR group as $`3[(3,2,1)_{\frac{1}{3}}+(1,2,1)_1+(1,1,1)_0+(1,1,2)_1+(\overline{3},1,2)_{\frac{1}{3}}]`$ $`(3,1,1)_{\frac{4}{3}}+(\overline{3},1,1)_{\frac{4}{3}}+(1,2,1)_{+1}+(1,1,2)_1+(1,2,2)_0`$ (3.26) while scalars live in $`3[(3,1,1)_{\frac{4}{3}}+(\overline{3},1,1)_{\frac{4}{3}}+(1,2,1)_1+(1,1,2)_1+(1,2,2)_0]`$ Fermions in 37 +73 sectors transform as $`(3,1,1;\overline{3},1,1,1)_{\frac{2}{3}}+(\overline{3},1,1;1,1,1,3)_{\frac{2}{3}}+`$ $`+(1,2,1;1,\overline{7},1,1)_1+(1,1,2;1,1,7,1)_1+\mathrm{LR}\mathrm{singlets}`$ (3.27) while scalars do as $`(3,1,1;1,1,\overline{7},1)_{\frac{2}{3}}+(1,2,1;1,1,1,\overline{3})_1+(1,1,2;1,\overline{7},1,1)_1+`$ $`+(1,1,2;3,1,1,1)_1+(\overline{3},1,1;1,7,1,1)_{\frac{2}{3}}+(1,2,1;1,1,7,1)_1+\mathrm{LR}\mathrm{singlets}`$ (3.28) We obtain three quark-lepton families from the 33 sector. The remaining fields, vector-like with respect to the LR group could acquire masses by breaking of the 77 groups. ### 3.6 Other Possibilities It is interesting to compare the kind of models we have constructed, using abelian orbifold singularities, with the field theories arising from D3-branes at more complicated singularities <sup>8</sup><sup>8</sup>8We have not discussed $`\mathrm{IC}^3/(ZZ_N\times ZZ_M)`$ singularities. However, it is a simple exercise to check they lead to models similar to those of $`\mathrm{IC}^3/ZZ_N`$ singularities. In particular, realistic SM and LR gauge groups are easily achieved by choosing suitable Chan-Paton actions, and hypercharge (or the $`BL`$ $`U(1)`$ in LR models) arises as a diagonal combination, straightforward generalization of (2.32). However, no three-generation models exist within this class. Also, the Weinberg angle, given by (3.4) by replacing $`N`$ by $`NM`$, is too small for any singularity of this type.. In this section we discuss some relevant cases. The details for the construction of the corresponding field theories can be found in appendices A, B, C and D. #### 3.6.1 Orbifold Singularities with Discrete Torsion Orbifold singularities $`\mathrm{IC}^3/(ZZ_{M_1}\times ZZ_{M_2})`$ lead to different models depending on their discrete torsion. In appendix A we review the field theory on singularities $`\mathrm{IC}^3/(ZZ_N\times ZZ_M\times ZZ_M)`$, with $`ZZ_N`$ twist $`(a_1,a_2,a_3)/N`$, and discrete torsion $`e^{2\pi i/M}`$ between the $`ZZ_M`$ twists. The final spectrum on the D3-brane world-volume is identical to that of the $`\mathrm{IC}^3/ZZ_N`$ singularity (2.21), but the superpotential is modified to (A.9). Hence it follows that the phenomenologically most interesting models in this class are those obtained from the $`ZZ_3`$ orbifold by a further $`ZZ_M\times ZZ_M`$ projection with discrete torsion. The resulting spectrum coincides with (3.13), and can lead to three-generation SM or LR models. The superpotential in the 33 sector is modified to $`W`$ $`=`$ $`\mathrm{Tr}[\mathrm{\Phi }_{i,i+1}^1\mathrm{\Phi }_{i+1,i+2}^2\mathrm{\Phi }_{i+2,i}^3]e^{2\pi i\frac{1}{M}}\mathrm{Tr}[\mathrm{\Phi }_{i,i+1}^1\mathrm{\Phi }_{i+1,i+2}^3\mathrm{\Phi }_{i+2,i}^2]`$ (3.29) Since the spectra we obtain from such singularities are identical to those in the simpler case of $`\mathrm{IC}^3/ZZ_3`$, we may question the interest of these models. They have two possible applications we would like to mention. The first, explored in section 5, is that discrete torsion enters as a new parameter that modifies the superpotential of the theory, and therefore the pattern of Yukawa couplings in phenomenological models. The second application is related to the process of moving branes off the singularity into the bulk, a phenomenon that, as discussed in Section 3 can take place in the realistic models there constructed. In $`ZZ_N\times ZZ_M\times ZZ_M`$ singularities with discrete torsion, the process can also occur, but the minimum number of branes allowed to move into the bulk is $`NM^2`$. In particular, in SM theories constructed from the $`\mathrm{IC}^3/(ZZ_3\times ZZ_M\times ZZ_M)`$ singularity, there are only six D3-branes with traceless Chan-Paton embedding, so motion into the bulk is forbidden for $`M2`$. For LR models, motion into the bulk of the twelve traceless D3-branes is forbidden for $`M3`$. In fact, the trapping is only partial, since branes are still allowed to move along planes fixed under some $`ZZ_M`$ twist . This partial trapping will be exploited in section 4.1 in the construction of certain compact models. #### 3.6.2 Non-Abelian Orbifold Singularities We may also consider the field theories on D3-branes at non-abelian orbifold singularities. The rules to compute the spectrum, along with the relevant notation, are reviewed in Appendix B. They allow to search for phenomenologically interesting spectra. The classification of non-abelian discrete subgroups of $`SU(3)`$ and several aspects of the resulting field theories have been explored in . An important feature in trying to embed the standard model in such field theories is that of triplication of families. Seemingly there is no non-abelian singularity where the spectrum appears in three identical copies. A milder requirement with a chance of leading to phenomenological models would be the appearance of three copies of at least one representation, i.e. $`a_{ij}^\mathrm{𝟑}=3`$ for suitable $`ij`$. Going through the explicit tables in there is one group with this property, $`\mathrm{\Delta }_{3n^2}`$ for $`n=3`$, on which we center in what follows. The 33 spectrum <sup>9</sup><sup>9</sup>9The gauge group in page 16 of corresponds to the particular case of Chan-Paton embedding given by the regular representation. Here we consider a general choice. for this case is $`{\displaystyle \underset{i=1}{\overset{9}{}}}U(n_i)\times U(n_{10})\times U(n_{11})`$ $`{\displaystyle \underset{i=1}{\overset{9}{}}}(n_i,\overline{n}_{10})+3(n_{10},\overline{n}_{11})+{\displaystyle \underset{i=1}{\overset{9}{}}}(n_{11},\overline{n}_i)`$ (3.30) If the triplicated representation is chosen to give left-handed quarks, a potentially interesting choice is given by $`n_i=1`$, $`n_{10}=3`$, $`n_{11}=2`$. However, and due to the additional factors $`r_i`$ (in our case $`r_i=1`$ for $`i=1,\mathrm{},9`$, $`r_{10}=r_{11}=3`$) in (B.4), the diagonal combination does not lead to correct hypercharge assignments. Besides (B.4) there are eight non-anomalous $`U(1)`$’s given by combinations $`Q_c=_{i=1}^9c_iQ_{n_i}+Q_{n_{10}}+\frac{3}{2}Q_{n_{11}}`$, with $`_{i=1}^9c_i=9`$. The charge structure under the combination $`c_1=c_2=c_3=3`$, $`c_i=0`$ for $`i=4,\mathrm{},9`$ is particularly interesting. We obtain $`SU(3)\times SU(2)\times U(1)(\times U(1)^8)`$ $`3(3,2)_{1/6}+3(1,2)_{1/2}+6(1,2)_{1/2}+3(\overline{3},1)_{2/3}+6(\overline{3},1)_{1/3}`$ (3.31) Hence, the spectrum under this $`U(1)`$ contains fields present in the standard model, and with the possibility of leading to three net copies (if suitable 37, 73 sectors are considered), even though they would still be distinguished by their charges under the additional $`U(1)`$’s. It is conceivable that this model leads to phenomenologically interesting field theories by breaking the additional $`U(1)`$ symmetries at a large enough scale. However, it is easy to see that the Weinberg angle, which, taking into account (B.3), is still given by (3.4) where now $`N=11`$, is exceedingly too small $`\mathrm{sin}^2\theta _W=3/62`$. Nevertheless we hope the model is illustrative on the type of field theory spectra one can achieve using non-abelian singularities. #### 3.6.3 Non-Orbifold Singularities In appendix C we discuss the construction of field theories on D3-branes at some non-orbifold singularities, which is in general rather involved. We also discuss that a promising spectrum is obtained from a partial blow-up of a $`ZZ_3`$ quotient of the conifold. The general spectrum is given in (C.12), (C.19). There are several possibilities to construct phenomenologically interesting spectra from this field theory. For the purpose of illustration, let us consider one example with $`n_2=3`$, $`n_0=2`$, $`n_1^{}=2`$, $`n_1^{}=1`$ and $`v_2=1`$, $`x_2=3`$, $`x_3=5`$ (all others vanishing). This choice leads to the spectrum $`SU(3)\times SU(2)\times SU(2)\times U(1)_{diag}\times U(1)^{}`$ $`\mathrm{𝟑𝟑}`$ $`3(3,2,1)_{\frac{1}{3}}+(\overline{3},1,1)_{\frac{4}{3}}+2(\overline{3},1,2)_{\frac{1}{3}}+(1,2,2)_0+2(1,2,1)_1+(1,1,2)_1`$ $`\mathrm{𝟑𝟕},\mathrm{𝟕𝟑}`$ $`(3,1,1;1,1)_{\frac{2}{3}}+(1,1,1;1,1)_2+(3,1,1;\overline{3},1)_{\frac{2}{3}}+`$ (3.32) $`(1,1,2;3,1)_1+(1,2,1;1,\overline{5})_1+(\overline{3},1,1;1,5)_{\frac{2}{3}}`$ with subindices giving charges under $`U(1)_{diag}`$. We obtain three net generations, but the right-handed quarks have a different embedding into the left-right symmetry: whereas two generations of right-handed quarks are standard $`SU(2)_R`$ doublets the other generation are $`SU(2)_R`$ singlets. Notice that an interesting property this model illustrates is that one can achieve three quark-lepton generations without triplication of the $`(1,2,2)`$ Higgs multiplets. This could be useful in order to suppress flavour changing neutral currents for these models. #### 3.6.4 Orientifold Singularities There is another kind of singularities that arises naturally in string theory, which we refer to as orientifold singularities. They arise from usual geometric singularities which are also fixed under an action $`\mathrm{\Omega }g`$, where $`\mathrm{\Omega }`$ reverses the world-sheet orientation, and $`g`$ is an order two geometric action. Orientifolds were initially considered in and have recently received further attention . Starting with the field theory on a set of D3-branes at a usual singularity, the main effect of the orientifold projection is to impose a $`ZZ_2`$ identification on the fields. Most examples in the literature deal with orientifolds of orbifolds singularities , on which we center in the following (see for orientifold of some simple non-orbifols singularities). To be concrete, we consider orientifolds of $`\mathrm{IC}^3/ZZ_N`$ singularities with odd $`N`$ . The spectrum before the orientifold projection is given in (2.21). This configuration can be modded out by $`\mathrm{\Omega }R_1R_2R_3(1)^{F_L}`$ (where $`R_r`$ acts as $`Y_rY_r`$, and $`F_L`$ is left-handed world-sheet fermion number), preserving $`𝒩=1`$ supersymmetry. The action of the orientifold projection on the 33 spectrum amounts to identifying the gauge groups $`U(n_i)`$ and $`U(n_i)`$, in such a way (due to world-sheet orientation reversal) that the fundamental representation $`n_i`$ is identified with the anti-fundamental $`\overline{n}_i`$. Consequently, there is an identification of the chiral multiplet $`\mathrm{\Phi }_{i,i+a_r}^r`$ and $`\mathrm{\Phi }_{ia_r,i}^r`$. Finally, since $`\mathrm{\Omega }`$ exchanges the open string endpoints, 37<sub>r</sub> fields $`\mathrm{\Phi }_{i,i\frac{1}{2}b_r}^{(37_r)}`$ map to 7<sub>r</sub>3 fields $`\mathrm{\Phi }_{Ni+\frac{1}{2}b_r,Ni}^{7_r3}`$. Notice that the gauge group $`U(n_0)`$ is mapped to itself and, similarly, chiral multiplets $`\mathrm{\Phi }_{i,i+a_r}^r`$ are mapped to themselves if $`i+a_r=Ni`$. There exist two possible orientifold projections, denoted ‘SO’ and ‘Sp’, differing in the prescription of these cases. The SO projection projects the $`U(n_0)`$ factor down to $`SO(n_0)`$, and the bi-fundamental $`\mathrm{\Phi }_{i,i+a_r}^r(=\mathrm{\Phi }_{i,i}^r)`$ to the two-index antisymmetric representation of the final $`U(n_i)`$. The Sp projection chooses instead $`USp(n_0)`$, and two-index symmetric representations. It is easy to realize that D3-branes at orientifold singularities will suffer from a generic difficulty in yielding realistic spectra. The problem lies in the fact that the orientifold projection removes from the spectrum the diagonal $`U(1)`$ (2.32) (as is obvious, since the $`U(1)`$ in $`U(n_0)`$ is automatically lost), which was crucial in obtaining correct hypercharge in our models in section 3. For illustration, we can check explicitly that, in the only candidate to yield three-family models, the orientifold of the $`\mathrm{IC}^3/ZZ_3`$ singularity, no realistic spectra arise. The general 33 spectrum for this orientifold (choosing, say the $`SO`$ projection) is $`SO(n_0)\times U(n_1)`$ $`3[(n_0,\overline{n}_1)+(1,{\displaystyle \frac{1}{2}}n_1(n_11))]`$ (3.33) The $`U(1)`$ factor is anomalous, with anomaly canceled by a GS mechanism . A seemingly interesting possibility would be $`SO(3)\times U(3)SU(3)\times SU(2)\times U(1)`$. However, the $`U(1)`$ factor does not provide correct hypercharge assignments. Moreover, it is anomalous and therefore not even present in the low-energy theory. A second possibility, yielding a Pati-Salam model $`SO(4)\times U(4)SU(2)\times SU(2)\times SU(4)\times U(1)`$ is unfortunately vector-like. In fact, the most realistic spectrum one can construct is obtained for $`n_0=1`$, $`n_1=5`$, yielding a $`U(5)`$ (actually $`SU(5)`$) gauge theory with chiral multiplets in three copies of $`\overline{\mathrm{𝟓}}+\mathrm{𝟏𝟎}`$. This GUT-like theory, which constitutes a subsector in a compact model considered in , does not however contain Higgs fields to trigger breaking to the standard model group. We hope this brief discussion suffices to support our general impression that orientifold singularities yield, in general, field theories relatively less promising than orbifold singularities. #### 3.6.5 Non-Supersymmetric Models from Antibranes We would like to conclude this section by considering a further set of field theories, obtained by considering branes and antibranes at singularities. The rules to compute the spectrum are reviewed in Appendix C. For simplicity we center on systems of branes and antibranes at $`\mathrm{IC}^3/ZZ_N`$ singularities, with $`ZZ_NSU(3)`$, even though the class of models is clearly more general. Notice that, due to the presence of branes and antibranes, the resulting field theories will be non-supersymmetric. Let us turn to the discussion of the generic features to be expected in embedding the standard model in this type of D3/$`\overline{\mathrm{D3}}`$ systems. The first possibility we would like to consider is the case with the standard model embedded on D3- and $`\overline{\mathrm{D3}}`$-branes. From (D.1), tachyons arise whenever the Chan-Paton embedding matrices $`\gamma _{\theta ,3}`$, $`\gamma _{\theta ,\overline{3}}`$ have some common eigenvalue. Denoting $`n_j`$, $`m_j`$ the number of eigenvalues $`e^{2\pi ij/N}`$ in $`\gamma _{\theta ,3}`$, $`\gamma _{\theta ,\overline{3}}`$, tachyons are avoided only if one considers models where $`n_i`$ vanishes when $`m_i`$ is non-zero, and vice-versa. This corresponds, for instance, to embedding $`SU(3)`$ in the D3-branes and $`SU(2)`$ in the $`\overline{D}3`$-branes. One immediate difficulty is manifest already at this level, regarding the appearance of hypercharge. As we remark in Appendix D, it is a simple exercise to extend the analysis of $`U(1)`$ anomalies of section 2.3 to the case with antibranes, with the result that (generically) there are two non-anomalous diagonal combinations (2.32) if no $`n_i`$, $`m_i`$ vanish. But this case is excluded by the requirement of absence of tachyons. Hence tachyon-free models lack the diagonal linear combination of $`U(1)`$’s and in general fail to produce correct hypercharge. It is still possible that in certain models, suitable tachyon-free choices of $`n_i`$, $`m_i`$ may produce hypercharge out of non-diagonal additional $`U(1)`$’s. Instead of exploring this direction (which in any case would not be available in the only three-family case of the $`\mathrm{IC}^3/ZZ_3`$ singularity), we turn to a different possibility. An alternative consists in embedding the standard model in, say D3-branes, but to satisfy the tadpole cancellation conditions using $`\overline{\mathrm{D7}}`$-branes (and possibly D7-branes as well). The resulting models are closely related to those in Section 3, differing from them only in the existence of $`3\overline{7}`$, $`\overline{7}3`$ sectors. Let us consider a particular example, with a set of D3-, D7<sub>3</sub>\- and $`\overline{\mathrm{D7}}_3`$-branes at a $`\mathrm{IC}^3/ZZ_3`$ singularity with twist $`v=(1,1,2)/3`$, with Chan-Paton embeddings $`\gamma _{\theta ,3}=\mathrm{diag}(\mathrm{I}_3,\alpha \mathrm{I}_2,\alpha ^2\mathrm{I}_1),\gamma _{\theta ,7_3}=\mathrm{diag}(\alpha ^2\mathrm{I}_3),\gamma _{\theta ,\overline{7}_3}=\mathrm{diag}(\mathrm{I}_3)`$ (3.34) which satisfy the tadpole condition $`\mathrm{Tr}\gamma _{\theta ,7_3}\mathrm{Tr}\gamma _{\theta ,\overline{7}_3}+3\mathrm{T}\mathrm{r}\gamma _{\theta ,3}=0`$ (3.35) The resulting spectrum on the D3-branes is $`SU(3)\times SU(2)\times U(1)_Y`$ $`\mathrm{𝟑𝟑}`$ $`3(3,2)_{\frac{1}{6}}+3(1,2)_{\frac{1}{2}}+3(\overline{3},1)_{\frac{2}{3}}`$ $`\mathrm{𝟑𝟕}_\mathrm{𝟑},\mathrm{𝟕}_\mathrm{𝟑}\mathrm{𝟑}`$ $`(1,2;\overline{3},1)_{\frac{1}{2}}+(\overline{3},1;3,1)_{\frac{1}{3}}`$ $`\mathrm{𝟑}\overline{\mathrm{𝟕}}_\mathrm{𝟑},\overline{\mathrm{𝟕}}_\mathrm{𝟑}\mathrm{𝟑}`$ $`(1,1;1,3)_1+(1,2;1,\overline{3})_{\frac{1}{2}}`$ $`(\mathrm{left}\mathrm{handed}\mathrm{fermions})`$ (3.36) $`(3,1;1,\overline{3})_{\frac{1}{3}}+(\overline{3},1;1,3)_{\frac{1}{3}}`$ $`(\mathrm{complex}\mathrm{scalars})`$ In the supersymmetric sectors, 33, 37 and 73, the above representations correspond to chiral multiplets, while in the non-supersymmetric 3$`\overline{7}`$, $`\overline{7}`$3 the spectrum for fermions and scalars is different. The complete spectrum corresponds to a three-generation model. Notice that models of this type can be obtained from our constructions in section 3 simply by adding an arbitrary number of $`\overline{\mathrm{D7}}`$-branes with traceless Chan-Paton factors, and annihilating fractional D7- and $`\overline{\mathrm{D7}}`$-branes with identical Chan-Paton phase (this corresponds to the condensation of the corresponding tachyons). Since many of the relevant properties of the models in section 3 are inherited by the non-supersymmetric theories with branes and antibranes, we do not pursue their detailed discussion here. ## 4 Embedding into a Compact Space As discussed in the introduction, even though SM gauge interactions propagate only within the D3-branes even in the non-compact setup, gravity remains ten-dimensional. In order to reproduce correct four-dimensional gravity the transverse space must be compact. In this section we present several simple string compactifications which include the local structures studied in section 3, as particular subsectors. Since these local structures contain D3- and D7-branes, the corresponding RR charges must be cancelled in the compact space. A simple possibility is to cancel these charges by including antibranes (denoted $`\overline{\mathrm{D3}}`$-, $`\overline{\mathrm{D7}}`$-branes) in the configuration. General rules to avoid tachyons and instabilities against brane-antibrane annihilation have been provided in (see for related models), which we exploit in section 4.1 to construct compact models based on toroidal orbifolds. A second possibility is to use orientifold planes. In section 4.2 we provide toroidal orientifolds without $`\overline{\mathrm{D3}}`$-branes, the D3-brane charge being cancelled by orientifold 3-planes (O3-planes). All of the above models contain some kind of antibranes, breaking supersymmetry in a hidden sector lying within the compactification space. It may be possible to find completely supersymmetric embeddings of our local structure by using toroidal orientifolds including O3- and O7-planes, even though we have not found such examples within the class of toroidal abelian orientifolds. This is hardly surprising, since this class is rather restricted. We argue that suitable $`𝒩=1`$ supersymmetric compactifications of our local structures exist in more general frameworks, like compactification on curved spaces. In fact, in section 4.3 we present an explicit F-theory compactification of this type. ### 4.1 Inside a Compact Type IIB Orbifold A simple family of Calabi-Yau compactifications with singularities is given by toroidal orbifolds $`T^6/ZZ_N`$. The twist $`ZZ_N`$ is constrained to act cristalographycally, and a list of the possibilities preserving supersymmetry is given in . Several such orbifolds include $`\mathrm{IC}^3/ZZ_3`$ singularities, the simplest being $`T^6/ZZ_3`$, with 27 singularities arising from fixed points of $`ZZ_3`$. Parameterizing the three two-tori in $`T^6`$ by $`z_r`$, $`r=1,2,3`$, with $`z_rz_r+1z_r+e^{2\pi i/3}`$, the 27 fixed points are located at $`(z_1,z_2,z_3)`$ with $`z_r=0`$, $`\frac{1}{\sqrt{3}}e^{\pi i/6}`$, $`\frac{1}{\sqrt{3}}e^{\pi i/6}`$. For simplicity, we denote these values by $`0`$, $`+1`$, $`1`$ in what follows. We now present several compact versions of the models in section 3, based on this orbifold. #### 4.1.1 Simple Examples The simplest possibility, in order to embed the local structures of section 3 in a compact orbifold, is to consider models with only one type of D7-brane. We will also consider the addition of Wilson lines, a modification which, being a global rather than a local phenomenon, was not available in the non-compact case. We first construct an example of a left-right symmetric model, which illustrates the general technique, and then briefly mention the construction of a model with standard model group. A Left-Right Symmetric Model Consider a compact $`T^6/ZZ_3`$ orbifold, and locate seven D3-branes with $`\gamma _{\theta ,3}=\mathrm{diag}(\mathrm{I}_3,\alpha \mathrm{I}_2,\alpha ^2\mathrm{I}_2)`$ (4.1) at the origin. To cancel twisted tadpoles at the origin we will add six D7<sub>1</sub>-branes at the origin in the first complex plane, with Chan-Paton embedding $$\gamma _{\theta ,7}=\mathrm{diag}(\alpha \mathrm{I}_3,\alpha ^2\mathrm{I}_3)$$ (4.2) The local structure around the origin is exactly as in the LR models studied in section 3.4, hence the final model will automatically contain a $`SU(3)\times SU(2)_L\times SU(2)_R\times U(1)_{BL}`$ gauge group, with three quark-lepton generations. That one can make this statement at such an early stage in the construction of the model is the main virtue of the bottom-up approach we are presenting. In the following, we simply discuss how to complete the model satisfying the string consistency conditions, and will not bother computing the detailed spectrum from the additional sectors. The choice (4.2) cancels the twisted tadpoles at the origin. But the worldvolume of these D7<sub>1</sub>-branes also overlap with other eight fixed points $`(0,n,p)`$, where $`n,p=0,\pm 1`$, so we will also have to ensure twisted tadpole cancellation at them. A simple option would be to add one single D3-brane with $`\gamma _{\theta ,3}=1`$ at each. However we consider the following more interesting possibility, which in addition breaks the (too large) D7<sub>1</sub>-brane $`U(3)\times U(3)`$ gauge group. Consider adding a Wilson line e.g. along the second complex plane, and acting on seven-branes as $$\gamma _{W,7}=\mathrm{diag}(\mathrm{I}_1,\alpha \mathrm{I}_1,\alpha ^2\mathrm{I}_1,\mathrm{I}_1,\alpha \mathrm{I}_1,\alpha ^2\mathrm{I}_1)$$ (4.3) The D7<sub>1</sub>-brane gauge group is broken down to $`U(1)^6`$ (some of which will actually be anomalous and therefore not really present). The fixed points $`(0,\pm 1,p)`$, $`p=0,\pm 1`$ feel the presence of the Wilson line , hence the D7<sub>1</sub> contribution to twisted tadpoles is given by $`\mathrm{tr}\gamma _{\theta ,7}\gamma _{W,7}`$ or $`\mathrm{tr}\gamma _{\theta ,7}\gamma _{W,7}^2`$, which vanish. Hence twisted tadpole cancellation does not require any additional branes at those fixed points. On the other hand the points $`(0,0,\pm 1)`$ do not feel the presence of the Wilson line and require e.g. the presence of one D3-brane with $`\gamma _{\theta ,3}=1`$ on each of them to achieve tadpole cancellation <sup>10</sup><sup>10</sup>10An amusing possibility is to add, instead of these single D3-branes, ‘new brane-worlds’ of seven 3-branes with (4.1), so that the model contains three different but analogous universes.. The above configuration has a total of nine D3-branes and six D7<sub>1</sub>-branes. The simplest possibility to cancel untwisted tadpoles is to locate six $`\overline{\mathrm{D7}}_1`$-branes at a different point, say $`Y_1=1`$, in the first complex dimension, with the same Chan-Paton twist matrix (4.2). Twisted and untwisted tadpoles are cancelled by adding one $`\overline{\mathrm{D3}}`$-brane at each of the nine fixed points $`(1,m,p)`$, $`m,p=0,\pm 1`$, with Chan-Paton embedding $`\gamma _{\theta ,\overline{3}}=\mathrm{diag}(\mathrm{I}_1)`$. The number of branes minus that of antibranes vanishes. The complete configuration is schematically depicted in Figure 5 A comment on the stability of this brane configuration is in order. In this compact configuration the D7-branes and $`\overline{\mathrm{D7}}`$-branes are trapped at the singular points, which they cannot leave. The flat directions associated to motions into the bulk are not present in their world-volume field theory. Moreover such motions would violate the twisted tadpole cancellation conditions. Hence, these objects cannot leave the fixed points and annihilate each other into the vacuum. On the other hand, groups of three D3-branes in the LR subsector may in principle abandon the origin and travel into the bulk in $`ZZ_3`$-invariant configurations, attracted by the $`\overline{\mathrm{D3}}`$-branes living at the fixed points $`(1,m,p)`$. Eventually the D3-branes could reach the $`\overline{\mathrm{D3}}`$-branes, and partial brane-antibrane annihilation would follow. This annihilation cannot be complete, since again the tadpole cancellation conditions would be violated. Since the complete configuration is non-supersymmetric, a more detailed analysis of the forces involved would be needed to find out whether the emission of D3-branes from the origin to the bulk is actually dynamically preferred or not, but we expect this specific model to suffer from such instability. On the other hand, as we will argue below, there are other ways to embed the set of SM or LR D3-branes into a compact space in which brane-antibrane annihilation is not possible. One Standard-like Model From the above construction, it is clear how to construct a compact orbifold model with an SM subsector, as we sketch in the following. We place the six D3-branes at the origin with $`\gamma _{\theta ,3}=\mathrm{diag}(\mathrm{I}_3,\alpha \mathrm{I}_2,\alpha ^2\mathrm{I}_1)`$ (4.4) and nine D7<sub>1</sub>-branes with $$\gamma _{\theta ,7}=\mathrm{diag}(\alpha \mathrm{I}_3,\alpha ^2\mathrm{I}_6)$$ (4.5) The local structure at the origin is of the type considered in section 3.3, and will lead to a three-generation $`SU(3)\times SU(2)_L\times U(1)_Y`$ subsector. In order to break the $`U(3)\times U(6)`$ D7<sub>1</sub>-brane gauge group, we add a Wilson line along the second complex plane, with $$\gamma _{W,7}=\mathrm{diag}(\mathrm{I}_1,\alpha \mathrm{I}_1,\alpha ^2\mathrm{I}_1,\mathrm{I}_2,\alpha \mathrm{I}_2,\alpha ^2\mathrm{I}_2)$$ (4.6) which break the D7<sub>1</sub>-brane gauge group to $`U(2)^3\times U(1)^3`$. Again, the fixed points $`(0,\pm 1,p)`$, $`p=0,\pm 1`$, feel the Wilson lines and receive no contribution to twisted tadpoles, hence do not require the presence of D3-branes. On the other hand the fixed points $`(0,0,\pm 1)`$ require additional D3-branes for tadpole cancellation. The simplest option is to add three 3-branes at each of these with $`\gamma _{\theta ,3}=\mathrm{diag}(\mathrm{I}_2,\alpha \mathrm{I}_1)`$. Let us add nine $`\overline{\mathrm{D7}}_1`$-branes at e.g. $`Y_1=+1`$, with $`\gamma _{\theta ,\overline{7}}=\mathrm{diag}(\mathrm{I}_3,\alpha \mathrm{I}_3,\alpha ^2\mathrm{I}_3).`$ (4.7) We also add a Wilson line $`W^{}`$ along the second complex plane, with embedding $`\gamma _{W^{},\overline{7}}=\mathrm{diag}(\mathrm{I}_3,1,\alpha ,\alpha ^2,1,\alpha ,\alpha ^2).`$ (4.8) Twisted tadpoles are generated only at points of the form $`(+1,\pm 1,p)`$, $`p=0,\pm 1`$. They can be cancelled by adding two $`\overline{\mathrm{D3}}`$-branes with $`\gamma _{\theta ,\overline{3}}=\mathrm{diag}(\alpha ,\alpha ^2)`$ at each. This completes the model, which satisfies the requirements of twisted and untwisted tadpole cancellation. #### 4.1.2 Trapping of Branes through Discrete Torsion As noted in section 3.6, orbifold models with discrete torsion may lead to partial trapping of D3-branes along certain directions. In this section we present a model with D3-branes and $`\overline{\mathrm{D3}}`$-branes where this partial trapping is used to avoid brane-antibrane annihilation. Consider the compact orbifold $`T^6/(ZZ_3\times ZZ_2\times ZZ_2)`$, with discrete torsion in the $`ZZ_2\times ZZ_2`$ factor. We place a set of twelve D3-branes at the origin $`(0,0,0)`$, with Chan-Paton embedding <sup>11</sup><sup>11</sup>11Actually, one should specify the embeddings of the $`ZZ_2`$ twists. They are given, up to phases, by the choice (A.2), and we do not give them explicitly to simplify the discussion. $`\gamma _{\theta ,3}=\mathrm{diag}(\mathrm{I}_6,\alpha \mathrm{I}_4,\alpha \mathrm{I}_2)`$ (4.9) As discussed in appendix A, the additional $`ZZ_2\times ZZ_2`$ twist breaks the gauge group to $`U(3)\times U(2)\times U(1)`$. In order to ensure cancellation of twisted tadpoles at the origin, we introduce D7<sub>1</sub>-branes at $`Y_1=0`$, with Chan-Paton embedding $`\gamma _{\theta ,7_1}=\mathrm{diag}(\alpha \mathrm{I}_6,\alpha ^2\mathrm{I}_{12})`$ (4.10) The local structure around the origin is precisely that of the SM theories in section 3.6.1, hence the final model will contain a three-generation $`SU(3)\times SU(2)_L\times U(1)_Y`$ sector. Due to the $`ZZ_2\times ZZ_2`$ twist, the D7<sub>1</sub>-brane gauge group would be $`U(3)\times U(6)`$. To reduce it further, we introduce a Wilson line $`W_2`$ along the second complex plane, embedded as $`\gamma _{W_2,7}=\mathrm{diag}(\mathrm{I}_2,\alpha \mathrm{I}_2,\alpha ^2\mathrm{I}_2;\mathrm{I}_4,\alpha \mathrm{I}_4,\alpha ^2\mathrm{I}_4)`$ (4.11) The surviving D7<sub>1</sub>-brane gauge group is $`U(1)^3\times U(2)^3`$. The D7<sub>1</sub>-brane contribution to twisted tadpoles at fixed points $`(0,\pm 1,p)`$ for $`p=0,\pm 1`$, vanishes, hence no D3-branes are required at them. At fixed points of the form $`(0,0,\pm 1)`$, tadpoles generated by the D7<sub>1</sub>-branes can be cancelled by placing D3-branes with $`\gamma _{\theta ,3}=\mathrm{diag}(\mathrm{I}_4,\alpha \mathrm{I}_2)`$. Notice that these D3-branes are stuck at the fixed points. So far we have placed twelve D3-branes at $`(0,0,0)`$ and six D3-branes at each of the points $`(0,0,\pm 1)`$. Also, there are eighteen D7<sub>1</sub>-branes at $`Y_1=0`$. In order to cancel the D7<sub>1</sub>-brane untwisted tadpole, we introduce nine anti-D7<sub>1</sub>-branes at $`Y_1=\pm 1`$, with $`\gamma _{\theta ,\overline{7}}=\mathrm{diag}(\mathrm{I}_3,\alpha \mathrm{I}_3,\alpha ^2\mathrm{I}_3).`$ (4.12) We also add a Wilson line $`W^{}`$ along the second complex plane, with embedding $`\gamma _{W_2^{},\overline{7}}=\mathrm{diag}(\mathrm{I}_3,1,\alpha ,\alpha ^2,1,\alpha ,\alpha ^2).`$ (4.13) They induce no tadpoles at the points of the form $`(+1,0,p)`$, $`p=0,\pm 1`$. The non-vanishing tadpoles at the points of the form $`(+1,\pm 1,p)`$, $`p=0,\pm 1`$, can be cancelled by adding two anti-D3-branes at each, with $`\gamma _{\theta ,\overline{3}}=\mathrm{diag}(\alpha ,\alpha ^2)`$. The $`\overline{\mathrm{D7}}_1`$-branes at $`(1,n,p)`$, with $`n,p=0,\pm 1`$, and the corresponding $`\overline{\mathrm{D3}}`$-branes are mapped to the above by the $`ZZ_2\times ZZ_2`$ operations, hence their embedding is determined by the above. The final configuration contains 18 D7<sub>1</sub>-branes and 18 $`\overline{\mathrm{D7}}_1`$-branes, as well as 24 D3-branes and 24 $`\overline{\mathrm{D3}}`$-branes, and is shown in figure 6. All RR charges are suitably cancelled. One important virtue of this model (and others of its kind) is that all $`\overline{\mathrm{D3}}`$-branes are stuck at fixed points, and so are most D3-branes. In fact, the only D3-branes not completely trapped are those in the SM sector at the origin $`(0,0,0)`$. As explained in section 3.6.1, they are however partially trapped, being able to move away at most in only one complex direction. Hence, they can never reach the points $`(\pm 1,\pm 1,p)`$ where anti-D3-branes sit, and brane-antibrane annihilation is therefore not possible. This example illustrates how to employ partial trapping by discrete torsion to improve the stability properties of the models. ### 4.2 Inside a Compact Type IIB Orientifold In this section we will present an explicit type IIB orientifold model with the following properties: 1. The standard model is embedded on a D3-brane at a $`\mathrm{IC}^3/ZZ_3`$ singularity. 2. There are no anti-D3-branes that could destabilize the system. 3. There are equal number of D7-branes and anti-D7-branes, breaking supersymmetry, trapped at different $`ZZ_3`$ fixed points, providing a clean example of gravity mediated supersymmetry breaking. Models of this type are similar to the ones constructed in (based on the general framework in , see also ), but in those cases the standard model was inside a higher dimensional brane (D7 or D9-brane) and supersymmetry was broken by the presence of lower dimensional anti-branes ($`\overline{\mathrm{D3}}`$\- or $`\overline{\mathrm{D5}}`$-branes). Furthermore those examples did not include models where all the anti-branes were completely trapped, and it was left as a dynamical question if the configuration was stable. The class of models we are going to present succeeds in trapping all antibranes. Nevertheless, it is well known that $`T`$-dualizing these models with respect to the different extra dimensions we should get to models where the standard model is embedded into D7 or D9-branes and therefore models like those considered in . We will then find the $`T`$ duals of these models which happen to be complicated versions of the models in in terms of several Wilson lines. We will submit the reader to the appendices where we discuss in detail the $`T`$ duality between the models of this section and those of . #### 4.2.1 Explicit Models As an illustration to this class of models let us consider in detail the following orientifold of the $`T^6/ZZ_3`$ orbifold. The orientifold action is $`\mathrm{\Omega }(1)^{F_L}R_1R_2R_3`$ with $`R_i`$ reflection on the $`i^{th}`$ plane. There are 64 orientifold three-planes (O3-planes), which are localized at points in the internal space. To cancel their RR charge we need a total of 32 D3 branes. We will distribute them among the 27 orbifold fixed points $`(m,n,p)`$, $`m,n,p=0,\pm 1`$. We will also introduce an equal number of D7 and $`\overline{\mathrm{D7}}`$-branes. Among these 27 points, only the origin $`(0,0,0)`$ is fixed under the orientifold action, hence it is an orientifold singularity, of the type mentioned in section 3.6.4. The cancellation of tadpoles at this point requires $$3\mathrm{Tr}\gamma _{\theta ,3}+(\mathrm{Tr}\gamma _{\theta ,7}\mathrm{Tr}\gamma _{\theta ,\overline{7}})=12$$ (4.14) As shown in appendix C, this orientifold singularity is not suitable to incorporate the standard model. Fortunately, $`ZZ_3`$ fixed points other than the origin are not fixed under the orientifold action, and therefore are $`\mathrm{IC}^3/ZZ_3`$ orbifold singularities, which we may employ to embed local configurations as those in section 3.3. The strategy is then to concentrate at an orbifold point different from the origin, say $`(0,1,0)`$ where we want to have the standard model (its $`ZZ_2`$ mirror image $`(0,1,0)`$ will have identical spectrum). We choose a suitable twist matrix for the D3 branes such as that in section 3.3. At these points the tadpole condition is simply: $$\mathrm{Tr}\gamma _{\theta ,7}\mathrm{Tr}\gamma _{\theta ,\overline{7}}+3\mathrm{T}\mathrm{r}\gamma _{\theta ,3}=0$$ (4.15) The challenge to construct a model is to be able to accommodate the 32 D3-branes among the different orbifold/orientifold points such that the corresponding twist matrices satisfy the tadpole conditions (4.14), (4.15) and of course having the same number of D7 and $`\overline{\mathrm{D7}}`$-branes. The other degree of freedom we have is the possibility of adding Wilson lines. They play an important role because at a given plane, a Wilson line differentiates among the different orbifold fixed points. A Left-Right Symmetric Model We can satisfy the tadpole condition (4.14) in many ways, but for future convenience we choose having eight D3-branes at the origin with $$\gamma _{\theta ,3}=\mathrm{diag}(\alpha \mathrm{I}_4,\alpha ^2\mathrm{I}_4)$$ (4.16) Let us now discuss the rest of the points $`(0,n,p)`$. Since these are only orbifold points the condition for tadpole cancellation is (4.15). If we want to have a standard-like model at a D3-brane on one of these points we must also have to introduce the D7 or $`\overline{\mathrm{D7}}`$-branes. Let us introduce D7<sub>1</sub>-branes at $`Y_1=0`$. In order not to alter the tadpole cancellation at the point $`(0,0,0)`$ the twist matrix for the 7-branes has to be traceless, so we take six D7<sub>1</sub>-branes with: $$\gamma _{\theta ,7}=\mathrm{diag}(\mathrm{I}_2,\alpha \mathrm{I}_2,\alpha ^2\mathrm{I}_2)$$ (4.17) In order to have nontrivial impact in the other points $`(0,n,p)`$ we add a Wilson line $`\gamma _W`$ on the second complex plane. We choose $$\gamma _{W,7}=\mathrm{diag}(\alpha ,\alpha ^2,\mathrm{I}_2,\mathrm{I}_2)$$ (4.18) The addition of this Wilson line implies that $`\mathrm{Tr}\gamma _{\theta ,7}\gamma _W=\mathrm{Tr}\gamma _{\theta ,7}\gamma _W^2=3`$ so we need to have D3-branes only at the points $`(0,n,p)`$ with $`n0`$ with twist matrix satisfying $`\mathrm{Tr}\gamma _{\theta ,3}=1`$. We achieve this at each of the points $`(0,\pm 1,0)`$ by adding seven D3 branes with twist matrix $$\gamma _{\theta ,3}=\mathrm{diag}(\mathrm{I}_3,\alpha \mathrm{I}_2,\alpha ^2\mathrm{I}_2)$$ (4.19) The local structure at these points is like that of left-right models in section 3.3. Hence they will lead to a $`SU(3)\times SU(2)_L\times SU(2)_R`$ gauge group with three quark-lepton generations. Next consider the points $`(0,\pm 1,\pm 1)`$. We can easily cancel tadpoles by introducing one single D3 brane at each of the four points, with $`\gamma _{\theta ,3}=1`$ Finally we have to consider the points $`(m,n,p)`$ with $`m0`$. So far we have six D7<sub>1</sub>-branes and $`8+7+7+4=26`$ D3-branes. Therefore we need to introduce six more D3 branes to complete 32, and six $`\overline{\mathrm{D7}}`$-branes. We can easily achieve this by locating three $`\overline{\mathrm{D7}}_1`$-branes at $`Y_1=\pm 1`$ with twist matrix and Wilson lines given by $$\gamma _{\theta ,\overline{7}}=\mathrm{I}_3\gamma _{\overline{W}\overline{7}}=(1,\alpha ,\alpha ^2)$$ (4.20) We then have $`\mathrm{Tr}\gamma _{\theta ,\overline{7}}\gamma _{\overline{W},\overline{7}}=\mathrm{Tr}\gamma _{\theta ,\overline{7}}\gamma _{\overline{W},\overline{7}}^2=0`$ and do not need to add D3 branes at the points $`(1,\pm 1,p)`$. At the points $`(1,0,p)`$ we can put a single D3-brane with $`\gamma _{\theta ,3}=1`$, canceling the tadpoles from $`\gamma _{\theta ,\overline{7}}`$. At the points $`(1,n,p)`$ we have a similar structure, consistently with the orientifold projection. In total, we have added precisely the six $`\overline{\mathrm{D7}}`$-branes and six D3-branes, as needed for consistency. Since this model provides an interesting example of a realistic compactification with gravity mediated supersymmetry breaking, we give some details on the computation of the spectrum in the visible sector. The D7-branes, with Chan-Paton embedding (4.17) would (before the orientifold projection) lead to a group $`U(2)^3`$ with matter in three copies of $`(2,2,1)+(1,2,2)+(2,1,2)`$. The Wilson line (4.18) breaks the gauge symmetry to $`U(1)^2\times U(2)\times U(2)`$ with matter in three copies of $`(2,\overline{2})`$ with zero charge under the $`U(1)^2`$. A simple way to read the spectrum is by using the shift vectors defined in the Cartan-Weyl basis (in analogy with ) for the twist matrix and Wilson line which in this case are: $`V_{7_0}`$ $`=`$ $`{\displaystyle \frac{1}{3}}(0,0,1,1,1,1)`$ $`W_{7_0}`$ $`=`$ $`{\displaystyle \frac{1}{3}}(1,1,0,0,0,0)`$ (4.21) We keep states of the form $`P=(1,1,0,\mathrm{})`$ (and permutations) satisfying $`PV=1/3`$ and $`PW=0`$. The orientifold projection identifies two pairs of groups and further breaks the symmetry to $`U(1)\times U(2)`$ with the simple matter consisting of three singlets (antisymmetrics of $`SU(2)`$) with charge +2 under the $`U(1)`$ inside $`U(2)`$. The D3-branes at the origin, with Chan-Paton embedding (4.16) lead, before the orientifold projection, to an $`U(4)\times U(4)`$ gauge group with matter on three copies of $`(4,\overline{4})`$. After the orientifold projection this becomes a single $`U(4)`$ group with matter on three copies of a 6. For states in the 37 sector for D3-branes at the origin, we use the extended shift vector $$\stackrel{~}{V}=\frac{1}{3}(1,1,1,1,1,1,1,1)V_{7_0}$$ (4.22) with the first eight entries corresponding to the original $`U(4)^2`$ on the D3-branes. We then keep the vectors <sup>12</sup><sup>12</sup>12The 73 sector would have the opposite relative signs between the D7 and D3 branes; but it is identified to the 37 sector by the orientifold projection and need not be considered independently. $`P=(1,0,\mathrm{};1,0,\mathrm{})`$ satisfying $`P\stackrel{~}{V}=1/3`$. We obtain the matter fields transforming as $`(4,2)_0+(\overline{4},1)_1+(\overline{4},1)_1`$ under the final $`U(4)\times U(2)\times U(1)`$. For D3-branes at $`(0,\pm 1,0)`$ the Chan-Paton twist is (4.19), we obtain the spectrum of the LR model. The gauge group is $`U(3)\times U(2)\times U(2)`$, with two anomalous $`U(1)`$’s being actually massive, and the diagonal combination giving $`BL`$. In the 33 sector, we obtain matter fields 33 sector: $$3\left[(3,\overline{2},1)+(1,2,\overline{2})+(\overline{3},1,2)\right]$$ (4.23) These can be explicitly seen by using the effective shift vector $`V_3=\frac{1}{3}(0,0,0,1,1,1,1)`$. For matter in 37, 73 sectors, the gauge quantum numbers are easily computed by considering the shift vectors $`V_3(V_7\pm W_7)`$, and keeping the state vectors $`P=(1,0,\mathrm{}0)(1,0\mathrm{},0)`$ and permutations for the 37 sector, and the opposite sign for the 73 sector. We get: 37 sector: $$(3,1,1;1)_1+(3,1,1;2)_0+(1,2,1;1)_1+(1,2,1;2)_0$$ (4.24) 73 sector: $$(\overline{3},1,1;1)_1+(\overline{3},1,1;2)_0+(1,1,2;1)_1+(1,1,2;2)_0$$ (4.25) The orientifold projection map the sets of branes at $`(0,1,0)`$ and $`(0,1,0)`$ to each other, so only one copy of the LR model is obtained. Notice that these sector contain some extra vector-like chiral fields beyond those obtained in section 3.4. Six pairs of colour triplets will in general become massive once $`(7_i7_i)`$ states get vevs. The remaining extra vector-like fields beyond the leptons transform like $`(3,1,1)+(1,2,1)+(1,1,2)+h.c.`$. As discussed in chapter 6, the presence of extra $`SU(2)_R`$ doublets is in fact welcome in order to give rise of the required $`SU(2)_R`$ gauge symmetry breaking. For D3-branes at fixed points $`(0,\pm 1,\pm 1)`$, with $`\gamma _{\theta ,3}=1`$, we have one $`U(1)`$ gauge group, with no 33 matter, at each brane. The orientifold projection relates two pairs of these points reducing the group to $`U(1)^2`$. The matter on each of these two 37 sectors can be easily computed as above (shift vector $`0V_7`$) and transforms as $`1_{1,1}+2_{0,1}`$ under the $`\left(U(2)\times U(1)\right)\times U(1)`$ (where the first $`U(1)`$ arises from D7-branes and the second one from D3-branes). The 73 sector gives multiplets $`1_{1,1}+2_{0,1}`$. A Standard-like Model Let us now briefly present a model with a SM subsector. We put twelve D7-branes at $`Y_1=0`$ and four D3-branes at the origin with $`\gamma _{\theta ,3}=\mathrm{diag}(\alpha \mathrm{I}_2,\alpha ^2\mathrm{I}_2),\gamma _{\theta ,7}=\mathrm{diag}(\alpha \mathrm{I}_6,\alpha ^2\mathrm{I}_6)`$ (4.26) satisfying the tadpole condition (4.14) at the origin. We also introduce a Wilson line along the second plane, with action $`\gamma _{W,7}=\mathrm{diag}(1,\alpha ,\alpha ^2,\mathrm{I}_9)`$ (4.27) The contribution to twisted tadpoles for the orbifold fixed points $`(0,\pm 1,0)`$ is given by $`\mathrm{tr}\gamma _7\gamma _W=\mathrm{tr}\gamma _7\gamma _W^2=6\alpha ^2+3\alpha `$. It allows us to put six D3-branes at each of the two points with $`\gamma _{\theta ,3}=\mathrm{diag}(\mathrm{I}_3,\alpha \mathrm{I}_2,\alpha ^2\mathrm{I}_1)`$ (4.28) The structure at these points is of the SM type considered in section 3.3 (it differs from it by an irrelevant overall phase in the Chan-Paton embeddings). Hence we obtain one three-generation SM sector. At each of the points $`(0,\pm 1,\pm 1)`$, we cancel the tadpoles by placing three D3-branes with $`\gamma _3=(\mathrm{I}_2,\alpha I_1)`$ Finally, we put two D3-branes with $`\gamma _3=\mathrm{I}_2`$ at each of the points $`(0,0,\pm 1)`$. The total number of D3-branes we have introduced is $`4+6+6+12+4=32`$, and there are also twelve D7-branes. The twelve $`\overline{\mathrm{D7}}`$-branes required to cancel the untwisted tadpole introduced by the D7-branes, could be evenly distributed at the positions $`Y_1=\pm 1`$, with $`\gamma _{\theta ,\overline{7}}=(\mathrm{I}_2,\alpha \mathrm{I}_2,\alpha ^2\mathrm{I}_2)`$. Being traceless, this embedding does not require to introduce more D3-branes nor Wilson lines. Unfortunately, these $`\overline{\mathrm{D7}}`$-branes can move off the fixed points into the bulk, and annihilate partially against the D7-branes at the origin. The result of such process is not necessarily disastrous for the SM sector we had constructed. In particular, it leads to a local structure of the type studied in section 3.6.5, leading to a non-supersymmetric SM with three generations. Hence the model would provide a realistic non-supersymmetric model, but necessarily requiring a TeV string scale. Alternatively, in order to avoid the partial annihilation of branes and anti-branes and obtain a gravity mediated supersymmetry breaking scenario we can modify the model in a simple way. Let us add, for instance, a second Wilson line, now in the third plane, acting on the D7-branes with $$\gamma _{W^{},7}=(\mathrm{I}_3;1,\alpha ,\alpha ^2,1,\alpha ,\alpha ^2,1,\alpha ,\alpha ^2)$$ (4.29) This Wilson line has no effect on the points $`(0,\pm 1,0)`$. Since $`\mathrm{Tr}\gamma _{\theta ,7}\gamma _{W^{},7}^n=3\alpha `$ for $`n=1,2`$, we need to add two D3 branes at the points $`(0,0,\pm 1)`$ with $`\gamma _{\theta ,3}=\mathrm{diag}(1,\alpha ^2)`$ to cancel tadpoles. However, since $`\mathrm{Tr}\gamma _{\theta ,7}\gamma _{W,7}^m\gamma _{W^{},7}^n=0`$ for $`m,n=1,2`$, we do not need to add any D3-brane at the points $`(0,\pm 1,\pm 1)`$. Therefore the remaining 12 D3-branes must sit at some of the fixed points $`(\pm 1,m,n)`$. We can achieve that having at each of the two planes $`Y_1=\pm 1`$ six $`\overline{\mathrm{D7}}_1`$-branes with $`\gamma _{\theta ,\overline{7}}=(\alpha \mathrm{I}_3,\alpha ^2\mathrm{I}_3)`$ with one Wilson line in the second plane with $`\gamma _{W,\overline{7}}=(1,\alpha ,\alpha ^2,1,\alpha ,\alpha ^2)`$ which precisely requires two D3 branes at each of the three fixed points points $`(1,0,m)`$ with $`\gamma _3=(\alpha ,\alpha ^2)`$ to cancel tadpoles. Putting a similar distribution at the plane $`Y_1=1`$, we end up with a model in which the anti-branes are also trapped, they cannot leave the $`Y_1=\pm 1`$ planes. The modulus parameterizing this motion has been removed by the Wilson line We have then succeeded in constructing a standard-like model with antibranes trapped at a hidden sector located at $`Y_1=\pm 1`$ without any danger to annihilate with the D7 branes of the visible sector, so it is a genuine gravity-mediated scenario. ### 4.3 Inside a Calabi-Yau Since our local singularities reproducing the semirealistic spectra in Section 3 contain D7-branes, the general framework to discuss compact models is F-theory . F-theory describes compactification of type IIB string theory on curved manifolds in the presence of seven-branes. Cancellation of magnetic charge under the RR axion implies the compact models should in general include a set of $`(p,q)`$ seven-branes, not mutually local, and around which the type IIB complex coupling constant $`\tau `$ suffers non-trivial $`SL(2,ZZ)`$ monodromies. Compactification to four dimensions on a three-complex dimensional manifold $`B_3`$, with a set of $`(p,q)`$ seven-branes wrapped on two-complex dimensional hypersurfaces in $`B_3`$, can be encoded in the geometry of a four-fold $`X_4`$, elliptically fibered over $`B_3`$, with $`(p,q)`$ seven-branes represented as the two-complex dimensional loci in $`B_3`$ over which the elliptic fiber degenerates by shrinking a $`(p,q)`$ one-cycle. In order to preserve $`𝒩=1`$ supersymmetry in four dimensions, $`X_4`$ must be Calabi-Yau, while $`B_3`$ in general is not. An important feature of four-dimensional compactifications of F-theory is that they generate a non-zero tadpole for the type IIB 4-form $`A_4^+`$ , its value (in the absence of 3-form fluxes ) being given, in units of D3-brane charge, by $`\frac{1}{24}\chi (X_4)`$, where $`\chi (X_4)`$ is the Euler characteristic of $`X_4`$. This tadpole may be cancelled by the addition of D3-branes (or instantons on the 7-brane gauge bundles, see below). Hence, F-theory compactifications contain the basic ingredients employed in our local structures (namely D7-branes, D3-branes, and a non-trivial geometry $`B_3`$) in a completely supersymmetric fashion. Clearly, local structures of the type studied in Section 3 will appear embedded in F-theory compactification where the base $`B_3`$ of the fourfold contains $`\mathrm{IC}^3/ZZ_3`$ singularities. In the following we construct a particular simple model containing a sector of D7-branes and D3-branes located at a $`\mathrm{IC}^3/ZZ_3`$ singularity, and reproducing the spectrum of a LR model of the type studied in Section 3.4. In order to keep the model tractable, we will take a simple fourfold as our starting point, namely $`X_4=X^{(1)}\times X^{(2)}`$, where $`X^{(1)}`$ is an elliptically fibered K3, and $`X^{(2)}`$ is also a K3, which we will take to be a $`T^4/ZZ_3`$ orbifold. Unfortunately K3$`\times `$K3 cannot contain $`\mathrm{IC}^3/ZZ_3`$ singularities, hence our final model will be a $`ZZ_3`$ quotient of $`X_4`$. The following construction of the model is pictured in figure 8 Let us start by describing $`X_4`$. The K3 manifold $`X^{(1)}`$ is a fibration of $`T^2`$ over a $`\mathrm{IP}_1`$, which can be defined in the Weierstrass form $`y^2`$ $`=`$ $`x^3+f_8(z,w)x+g_{12}(z,w)`$ (4.30) where $`x`$, $`y`$ parameterize $`T^2`$, and $`[z,w]`$ are projective coordinates in $`\mathrm{IP}_1`$. The functions $`f_8`$, $`g_{12}`$ are of degree $`8`$ and $`12`$ respectively in their arguments. The elliptic fibration degenerates at the 24 points $`[z,w]`$ given by zeroes of the discriminant $`\delta _{24}`$ $`=`$ $`4f_8(z,w)^3+27g_{12}(z,w)^2`$ (4.31) These degenerations signal the presence of 24 $`(p,q)`$ seven-branes sitting at a point in $`\mathrm{IP}_1`$ and wrapping $`X^{(2)}`$ completely. We will be interested in a particular family of K3 manifolds, containing six coincident D7-branes sitting at $`z=0`$. Geometrically, we require the fibration (4.30) to have an $`I_6`$ Kodaira type degeneration at $`z=0`$, meaning $`f_8`$ and $`g_{12}`$ are non-zero at $`z=0`$, but $`\delta _{24}`$ vanishes as $`z^6`$. We will be more explicit below. For $`X_4=K3\times K3`$, $`\chi (X_4)=24\times 24`$, and the compactification leads to an $`A_4^+`$-tadpole of $`24`$ units. It will be useful to rederive this result from a different point of view. On the world-volume $`M_7`$ of each seven-brane (of any $`(p,q)`$ type) there exist couplings $`{\displaystyle \frac{1}{24}}{\displaystyle _{M_7}}RRA_4^++{\displaystyle _{M_7}}FFA_4^+`$ (4.32) where $`R`$ is the Ricci tensor of the induced curvature, and $`F`$ is the field strength tensor for the world-volume gauge fields. In the absence of world-volume gauge instantons the second term drops, and the contribution of the first for each seven-brane wrapped on the K3 $`X^{(2)}`$ (recall $`\chi (K3)=_{K3}RR=24`$) is $`_{M_4}A_4^+`$, where $`M_4`$ is four-dimensional spacetime. Namely, we get an $`A_4^+`$-tadpole of $`1`$ for each of the 24 seven-branes in the model, hence a total tadpole of $`24`$. This tadpole can be cancelled by adding 24 D3-branes sitting at a point in $`B_3=\mathrm{IP}_1\times X^{(2)}`$, or by turning on instantons on the seven-brane gauge bundles. We will be interested in considering $`X^{(2)}`$ in the $`T^4/ZZ_3`$ orbifold limit. We start with $`T^4`$ parametrized by two coordinates $`z_1`$, $`z_2`$ with identifications $`z_iz_i+1`$, $`z_iz_i+e^{2\pi i/3}`$ (we set the radii to unity for the sake of simplicity), and mod out by the $`ZZ_3`$ action $`\theta :z_1e^{2\pi i/3}z_1,z_2e^{2\pi i/3}z_2`$ (4.33) The quotient is a K3 manifold, flat everywhere except at the nine fixed points of $`\theta `$ (located at points $`(z_1,z_2)`$ with $`z_1`$, $`z_2`$ = $`0`$, $`\frac{1}{3}(1+e^{\pi i/3})`$ or $`\frac{1}{3}(e^{\pi i/3}+e^{2\pi i/3})`$), which descend to singularities $`\mathrm{IC}^2/ZZ_3`$ in the quotient. The curvature is concentrated at those points, and gives a contribution of $`_{C^2/Z_3}RR=8/3`$ at each. Hence the $`1`$ $`A_4^+`$-tadpole at each seven-brane is split in nine $`1/9`$ $`A_4^+`$-tadpole, arising from the nine points $`\mathrm{IC}^2/ZZ_3`$ wrapped by the seven-brane world-volume. There are no contributions from world-volume gauge instantons if the gauge bundle over these $`\mathrm{IC}^2/ZZ_3`$ is trivial; in more familiar terms, if the $`ZZ_3`$ action is embedded trivially on the seven-brane indices. Given this $`\theta `$-embedding, one may worry about twisted tadpoles. However, since the $`ZZ_3`$ twist leaves fixed the complete $`\mathrm{IP}_1`$, such tadpoles receive contributions from all seven-branes, and they cancel by the same reason the overall axion RR-charge cancels, namely by cancellation among contributions of different $`(p,q)`$ seven-branes. As mentioned above, in order to obtain $`\mathrm{IC}^3/ZZ_3`$ singularities we are forced to consider a quotient of the four-fold $`X_4`$ considered above, by $`ZZ_3`$, with generator $`\omega `$ acting simultaneously on $`\mathrm{IP}_1`$ and $`X^{(2)}`$. We choose the action of $`\omega `$ on $`\mathrm{IP}_1`$ to be given by $`\omega _1:ze^{2\pi i/3}z,ww`$ (4.34) For the configuration to be invariant under this action, the K3 $`X^{(1)}`$ must be given by (4.30) with the functions $`f_8`$, $`g_{12}`$ depending only on $`z^3`$ $`f_8(z,w)`$ $`=`$ $`Aw^8+Bz^3w^5+Cz^6w^2`$ $`g_{12}(z,w)`$ $`=`$ $`Dw^{12}+Ez^3w^9+Fz^6w^6+Gz^9w^3+Hz^{12}`$ (4.35) The requirement that there are six D7-branes at $`z=0`$ amounts to the vanishing of the coefficients of order $`1`$ and $`z^3`$ in $`\delta _{24}`$ $`4A^3+27D^2=0,12A^2B+54DE=0`$ (4.36) It will be convenient to require no seven-branes to be present at $`[z,w]=[1,0]`$, hence we require $`H0`$. The configuration we have constructed therefore contains six D7-branes at $`z=0`$, and the remaining 18 seven-branes distributed on $`\mathrm{IP}_1`$ in a $`ZZ_3`$ symmetric fashion. Notice that not only the locations, but the $`(p,q)`$ types of these seven-branes are consistent with the symmetry, since the F-theory K3 $`X^{(1)}`$ is $`ZZ_3`$ invariant. Also, notice that no $`(p,q)`$ seven-branes (other than the six at $`z=0`$) are fixed under $`\omega `$. We have to specify also the action of $`\omega `$ on $`X^2`$, which we take to be $`\omega _2:z_1e^{2\pi i/3}z_1+{\displaystyle \frac{1}{3}}(1+e^{\pi i/3}),z_2`$ $`e^{2\pi i/3}z_2+{\displaystyle \frac{1}{3}}(1+e^{\pi i/3})`$ (4.37) (The shifts have been included to avoid $`\omega `$-fixed points to coincide with $`\theta `$-fixed points). That this is a symmetry of $`X^{(2)}=T^4/ZZ_3`$ follows from the fact that $`\omega _2`$ is a symmetry of $`T^4`$, and that it commutes with the action $`\theta `$ (4.33) required to form $`T^4/ZZ_3`$. The action $`\omega _2`$ has nine fixed points in $`T^4`$, given by points $`(z_1,z_2)`$ with $`z_i=\frac{1}{3}e^{\pi i/3}`$, $`\frac{1}{3}(1+2e^{\pi i/3})`$, $`\frac{1}{3}(e^{2\pi i/3}+2e^{\pi i/3})`$. The action of $`\theta `$ maps them to each other so that in the quotient they yield just three singular points. Similarly, $`\omega _2`$ maps the nine fixed points of $`\theta `$ to each other, so that in the quotient they give rise to just three singularities. Notice that in the final model the $`T^4`$ is modded out by $`\theta `$ and $`\omega _2`$, and consequently by all other twists they generate. In this respect it is important to notice that $`\theta \omega _2`$ and $`\theta ^2\omega _2`$ (and their inverses) are freely acting due to the shifts in (4.37), and do not generate new singularities. The final model is obtained by quotienting $`\mathrm{IP}_1\times X^{(2)}`$ by $`\omega `$, the combined action of (4.34), (4.37). One can see that it preserves $`𝒩=1`$ supersymmetry in the quotient by noticing that the holomorphic 4-form in $`X_4`$, given by $`(dx/y)(dz/w)dz_1dz_2`$ is invariant, hence the quotient is still a four-fold. Notice that the fixed points of $`\omega `$ correspond to $`[z,w]=[0,1],[1,0]`$ in $`\mathrm{IP}_1`$ times the three fixed points in $`X^{(2)}`$. We must also specify the action of $`\omega `$ on the seven-branes. For the $`(p,q)`$ seven-branes at $`z0`$, the action is fully specified by the geometrical action in $`\omega _1`$. The six D7-branes at $`z=0`$, sit at a fixed point of $`\omega `$, and hence may suffer a non-trivial action on their Chan-Paton indices, which we choose to be given by $`\gamma _{\omega ,7}=\mathrm{diag}(e^{2\pi i\frac{1}{3}}\mathrm{I}_3,e^{2\pi i\frac{2}{3}}\mathrm{I}_3)`$ (4.38) This amounts to choosing a non-trivial $`U(6)`$ bundle on D7-brane world-volume. This will lead to a non-zero contribution to world-volume instanton number, and hence influence the computation of the $`A_4^+`$-tadpole in the quotient. Moreover, the embedding (4.38) also leads to non-zero $`\omega `$-twisted tadpoles which must be properly cancelled. Let us turn to the computation of the $`A_4^+`$ tadpole. Since the $`\omega `$ action has fixed points, the Euler characteristic of the quotient four-fold is not simply $`\chi (X_4)/3`$. Direct computation seems rather involved, hence we prefer to compute the $`A_4^+`$ tadpole by using the seven-brane world-volume couplings (4.32). The 18 $`(p,q)`$ seven-branes at $`z0`$ are not invariant under $`\omega `$, hence their world-volume $`RR`$ contribution comes only from the nine $`\theta `$-fixed points, yielding a total $`A_4^+`$ tadpole of $`18`$. This is reduced to $`6`$ by the $`\omega `$ action (equivalently, the 18 seven-branes in $`X_4`$ descend to just six in the quotient by $`\omega `$). The six D7-branes at $`z=0`$ sit at a $`\omega `$ fixed point, hence their $`RR`$ coupling receives contributions from three $`\mathrm{IC}^2/ZZ_3`$ (from nine $`\theta `$-fixed $`\omega `$-identified points), and three $`\mathrm{IC}^2/ZZ_3`$ (from nine $`\omega `$-fixed $`\theta `$-identified points). The total contribution is $`\frac{1}{24}\times \frac{8}{3}\times 3\times 2\times 6=4`$. There is a further contribution, from the instanton number of the nontrivial bundle implied by (4.38). This can be computed (see for a discussion of instanton numbers for orbifold spaces) to be $`+1`$ per $`\mathrm{IC}^3/ZZ_3`$ generated by $`\omega `$, yielding an overall contribution $`+3`$ to the $`A_4^+`$ tadpole. Hence, the total $`A_4^+`$ tadpole in the model is $`7`$, and one must introduce $`7`$ D3-branes (as counted in the quotient) to achieve a consistent model. Recall that the choice of the twist (4.38) generates non-zero $`\omega `$-twisted tadpoles at the nine fixed points of $`\omega _2`$ (three in the quotient), at $`z=0`$. Happily, these are $`\mathrm{IC}^3/ZZ_3`$ fixed points of the type studied in section 3.2, for which the twisted tadpole condition reads $`\mathrm{Tr}\gamma _{\omega ,7}+3\mathrm{T}\mathrm{r}\gamma _{\omega ,3}`$ $`=`$ $`0`$ (4.39) The tadpoles are easily cancelled by choosing $`\gamma _{\omega ,3}=(I_3,e^{2\pi i\frac{1}{3}}I_2,e^{2\pi i\frac{2}{3}}I_2)`$ at three of them (one after the $`\theta `$-identification) and $`\gamma _{\theta ,3}=I_1`$ at six (two in the quotient by $`\theta `$). This employs 3 D3-branes (as counted in the quotient) out of the 7 we had available. The remaining 4 can be placed at a generic, smooth, point in the quotient (so that they would be described by $`4\times 3\times 3`$ in the covering space of the $`\theta `$ and $`\omega `$ actions). This concludes the construction of the model. We have succeeded in constructing a $`𝒩=1`$ supersymmetric compact model with a local singularity of the form $`\mathrm{IC}^3/ZZ_3`$, on which D3-branes and D7-branes sit. This local structure is of the kind considered in section 3.4, and leads to a left-right symmetric model $`SU(3)\times SU(2)\times SU(2)\times U(1)_{BL}`$ with three generations of standard model particles. Notice that in this concrete example the D7-brane gauge group is relatively large, but it should not be difficult to modify the construction (for instance, by adding Wilson lines on $`T^4`$) to reduce it. We prefer not to complicate the discussion for the moment. We would like to conclude with a remark. F-theory model building (see e.g. ) has centered on embedding the standard model gauge group on the seven-branes, while D3-branes are basically a useless sector. In our framework, the situation is inverted, with the D3-branes playing the main role, even if seven-branes also contribute to the spectrum. One particular advantage of our alternative embedding, which we have stressed throughout the paper, is that the relevant sector can be determined by using just the local behaviour of the compactification manifold. In this respect it is clear that many other consistent supersymmetric embeddings of our local structures may be achieved by considering more general four-folds with $`\mathrm{IC}^3/ZZ_3`$ singularities on the base, on which D3-branes and D7-branes (with gauge bundles locally reproducing the required monodromy) are located. ## 5 Some Phenomenological Aspects Obtaining $`SU(3)\times SU(2)_L\times U(1)_Y`$ (or its left-right symmetric extension $`SU(3)\times SU(2)_L\times SU(2)_R\times U(1)_{BL}`$) and three quark-lepton generations in a simple manner is already remarkable. However we would like also to know if the class of models discussed in the previous sections is sufficiently rich to accommodate other important phenomenological features like gauge coupling unification, a hierarchical structure of quark/lepton masses and proton stability. In this chapter we would like to address these issues. Let us make clear from the beginning that we do not intend here to give a full account of these questions but rather to evaluate the phenomenological potential of these models without compromising ourselves with a particular one. We will limit ourselves to the case of models obtained from D-branes at $`𝒩=1`$ supersymmetric singularities. As already remarked above, some properties like number of generations, gauge group and gauge coupling normalization will mostly be controlled by the local D-brane configurations introduced in section 3. On the other hand other aspects like some Yukawa couplings or the $`𝒩=1`$ Kahler metrics will be in general dependent on the particular compact space chosen in order to embed the realistic D-brane set. Thus, for example, the Kahler metric of the matter fields will in general be different for an embedding inside a compact $`ZZ_3`$ orbifold compared to an embedding inside an F-theory model. Thus, the complete effective action will certainly depend on the specific model. We will try to discuss certain features which seem generic among the class of models considered in the previous chapters. In particular, that is the case of the question of gauge coupling unification. We will also study some aspects of the Yukawa couplings both in the SM and LR type of models. As a general conclusion we believe that the structure of the models seems to be sufficiently rich to be able to accommodate the main phenomenological patterns of the fermion mass spectrum. Finally, we comment on the structure of mass scales and in particular the value of the string scale $`M_s`$ that should be considered in these constructions. Since this depends on the compactification scale, the choice of $`M_s`$ is also model dependent. i) Gauge Coupling Unification This question obviously depends on whether we are considering a SM like the one considered in section 3.3 or a left-right symmetric model as in section 3.4. As we will see, from the gauge coupling unification point of view the left-right models look more interesting. Let us consider first the case of the SM. As mentioned in section 3, the tree-level gauge coupling constants at the string scale $`M_s`$ are in the ratios $`g_3^2:g_2^2:g_1^2=1:1:11/3`$. The one-loop running between the weak scale $`M_Z`$ and the string scale $`M_s`$ is governed by the equations $`\mathrm{sin}^2\theta _W(M_Z)=`$ $`{\displaystyle \frac{3}{14}}\left(1+{\displaystyle \frac{11}{6\pi }}\alpha (M_Z)(b_2{\displaystyle \frac{3}{11}}b_1)\mathrm{log}({\displaystyle \frac{M_s}{M_Z}})\right)`$ $`{\displaystyle \frac{1}{\alpha _3(M_Z)}}=`$ $`{\displaystyle \frac{3}{14}}\left({\displaystyle \frac{1}{\alpha (M_Z)}}{\displaystyle \frac{1}{2\pi }}\left(b_1+b_2{\displaystyle \frac{14}{3}}b_3\right)\mathrm{log}\left({\displaystyle \frac{M_s}{M_Z}}\right)\right)`$ (5.1) where $`b_i`$ are the $`\beta `$-function one-loop coefficients. Notice that the tree level result for $`\mathrm{sin}^2\theta _W`$$`=3/14=0.215`$ is only slightly below the measured value at $`M_Z`$, $`\mathrm{sin}^2\theta _W(M_Z)=0.231`$. Thus in order to get the correct sign for the one-loop correction we need to have $`b_2>3b_1/11`$. With the massless spectrum in Table 1 we have $`b_3=0`$ , $`b_2=3`$, $`b_1=15`$, although as mentioned in section 3.3, generically some three pairs of colour triplets will be heavy, leading to $`b_3=3`$, $`b_2=3`$, $`b_1=13`$. In both cases the one-loop correction has the wrong sign and one gets $`\mathrm{sin}^2\theta _W(M_Z)=0.18`$, $`M_s=2.2\times 10^{15}`$ GeV ($`\mathrm{sin}^2\theta _W(M_Z)=0.204`$, $`M_s=10^{10}`$ GeV) respectively. Thus standard logarithmic gauge coupling unification within the particular SM configurations of section 3.3 does not seem to work. We cannot exclude, however that a more sophisticated configuration of D-branes (yielding, in particular, some extra massless doublets) could work . Let us now move to the left-right model case. In this case the gauge coupling unification question works remarkably well, as already pointed out in ref.. Indeed, the massless spectrum of the $`SU(3)\times SU(2)_L\times SU(2)_R\times U(1)_{BL}`$ model of section 3.4 is essentially identical to the model considered in that reference. There are now two regions for the running, $`M_R<Q<M_s`$, where the gauge group is the left-right symmetric one and $`M_Z<Q<M_R`$ where the gauge group is the SM one. Thus $`M_R`$ is the scale of breaking of the left-right symmetry. If the $`\beta `$-function coefficients of the left-right gauge group are denoted by $`B_3,B_L,B_R`$ and $`B_{BL}`$, the one-loop running yields $`\mathrm{sin}^2\theta _W(M_Z)=`$ $`{\displaystyle \frac{3}{14}}(1+{\displaystyle \frac{11\alpha _e(M_Z)}{6\pi }}[(B_L{\displaystyle \frac{3}{11}}B_1^{})\mathrm{log}\left({\displaystyle \frac{M_s}{M_R}}\right)`$ (5.3) $`+(b_2{\displaystyle \frac{3}{11}}b_1)\mathrm{log}\left({\displaystyle \frac{M_R}{M_Z}}\right)])`$ $`{\displaystyle \frac{1}{\alpha _e(M_Z)}}{\displaystyle \frac{14}{3\alpha _3(M_Z)}}=`$ $`{\displaystyle \frac{1}{2\pi }}[(b_1+b_2{\displaystyle \frac{14}{3}}b_3)\mathrm{log}\left({\displaystyle \frac{M_R}{M_Z}}\right)`$ $`+(B_1^{}+B_L{\displaystyle \frac{14}{3}}B_3)\mathrm{log}\left({\displaystyle \frac{M_s}{M_R}}\right)]`$ where one defines $$B_1^{}=B_R+\frac{1}{4}B_{BL}$$ (5.4) and $`b_i`$ are the $`\beta `$-function coefficients with respect to the SM group. For the generic (i.e., no extra triplets) massless spectrum found in section 3.4 one has $`B_3=3`$, $`B_L=B_R=+3`$ and $`B_{BL}=16`$. To get a numerical idea let us assume that the left-right symmetry is broken not far away from the weak scale <sup>13</sup><sup>13</sup>13This is in fact the most natural assumption if both electroweak and $`SU(2)_R`$ symmetry breaking are triggered by soft SUSY-breaking soft terms. , e.g. at $`M_R`$ 1 TeV, and that the spectrum below that scale is given by that of the MSSM (so that $`b_3=3`$, $`b_2=1`$ and $`b_1=11`$). <sup>14</sup><sup>14</sup>14A more general quantitative analysis for different values of $`M_R`$ and SUSY-breaking soft terms may be found in ref.. The general agreement is in general found as long as $`M_R<1`$ TeV or so. Then one finds the results: $$\mathrm{sin}^2\theta _W(M_Z)=0.231;M_s=9\times 10^{11}\mathrm{GeV}$$ (5.5) in remarkable agreement with low-energy data. This agreement requires $`SU(2)_R`$ to be broken not much above 1 TeV. In this connection, notice that the massless spectrum of the left-right model in section 3.4 does not have the required fields in order to do this breaking. However, slight variations like the explicit model in section 4.2.1 do have extra $`SU(2)_R`$ doublets which can produce the breaking. On the other hand it is easy to check that the presence of the additional vector-like chiral fields in that model does not affect the one-loop coupling unification and hence the agreement remains. In summary, although logarithmic gauge coupling unification in the SM would require some modification of the models, couplings unify nicely (with equal precision than in the MSSM) in the case of the $`SU(3)\times SU(2)_L\times SU(2)_R\times U(1)_{BL}`$ models. In this case the string scale coincides with the unification scale and should be of order $`10^{12}`$ GeV. ii) Yukawa Couplings Many phenomenological issues depend on the Yukawa coupling structure of the models. At the renormalizable level there are essentially three type of superpotential couplings involving physical chiral fields : a) $`(33)^3`$ , b) $`(33)(73)(37)`$ and c) $`(73)(37)(77)`$. Whereas the couplings of types b) and c) depend on the structure of D7-branes and hence are more sensitive to global effects, that is not the case of couplings of type a) which involve couplings among different $`(33)`$ chiral matter fields. Those are expected to depend mostly in the structure of the orbifold singularity on top of which the D3-branes reside. Indeed, for the more general case of a $`ZZ_3\times ZZ_M\times ZZ_M`$ singularity in the presence of discrete torsion (as in appendix A) the structure of those superpotentials has the form: $`W`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{2}{}}}ϵ^{abc}[\mathrm{\Phi }_{i,i+1}^a\mathrm{\Phi }_{i+1,i+2}^b\mathrm{\Phi }_{i+2,i}^c]+(1e^{2\pi i\frac{1}{M}})s^{abc}[\mathrm{\Phi }_{i,i+1}^a\mathrm{\Phi }_{i+1,i+2}^b\mathrm{\Phi }_{i+2,i}^c]`$ (5.6) where (see equation (A.9)) $`s^{abc}`$ is such that $`s^{321}=s^{132}=s^{213}=1`$ and all other components vanish. The gauge group is $`U(n_0)\times U(n_1)\times U(n_2)`$ and the $`\mathrm{\Phi }_{i,j}^a`$ are bifundamental representations $`(n_i,\overline{n_j})`$. In the absence of discrete torsion the second piece vanishes and the superpotentials are purely antisymmetric. In the particular case of the SM, these include superpotential couplings corresponding to Yukawa couplings giving masses to U-type quarks, i.e., couplings of type $`h_{abc}Q_L^aU_R^bH^c`$. The corresponding Yukawa couplings have thus the form: $$h_{abc}=ϵ_{abc}+(1e^{2\pi i\frac{1}{M}})s_{abc}$$ (5.7) with $`a,b,c=1,2,3`$ are family indices. However, to make contact with the low energy Yukawa couplings we have to recall that in a general $`𝒩=1`$ supergravity theory the kinetic terms of these chiral fields will not be canonical but will be given by a a Kahler metric $`K_{ab}`$ which will in general be a function of the compactification moduli. Consider, to be specific, the case in which the compact case is a standard $`ZZ_3`$ toroidal orbifold or orientifold, like the examples shown in chapter 4. In this case one has $`K_{ab}=(T_{a\overline{b}}+T_{a\overline{b}}^{})^1`$, where $`T_{a\overline{b}}`$ are the nine untwisted Kahler moduli of the $`ZZ_3`$ orbifold. Now, one can easily check that the low-energy physical Yukawa couplings corresponding to canonically normalized standard model fields $`h_{abc}^0`$ will be related with the supergravity base ones $`h_{abc}`$ as : $$h_{abc}^0=h_{lmn}\mathrm{exp}(K/2)(K_{al}K_{bm}K_{cn})^{1/2}$$ (5.8) where we have neglected a complex phase irrelevant for our purposes and $`K`$ denotes the full Kahler potential. The latter is given in the weak coupling limit by the expression $`K=\mathrm{log}(S+S^{})\mathrm{log}det(T_{a\overline{b}}+T_{a\overline{b}}^{})`$, where $`S`$ is the complex dilaton field. Thus altogether, the low energy physical Yukawa couplings of U-type quarks in these models will have the general form: $$h_{abc}^0=(S+S^{})^{1/2}[ϵ_{abc}+(1e^{2\pi i\frac{1}{M}})s^{lmn}det(t_{jk})^{1/2}(t_{al}t_{bm}t_{cn})^{1/2}]$$ (5.9) where we define $`t_{ab}=T_{a\overline{b}}+T_{a\overline{b}}^{}`$. Notice that the Yukawa couplings are proportional to the gauge coupling constant $`g`$, since $`ReS=2/\lambda `$, $`\lambda `$ being the Type IIB dilaton. Also, in the absence of discrete torsion ($`s^{abc}=0`$) these Yukawa couplings do not depend on the moduli. To see whether this kind of structure has a chance to describe the observed pattern of U-quark masses, let us consider the case in which only one of the three Higgs fields, say $`H_1`$, eventually gets a vev of order the weak scale, $`H_1=v`$. The U-quark mass matrix will be thus given by $`M_{bc}^U=vh_{1bc}^0`$. Let us assume also that the moduli with largest vevs are $`t_{11},t_{33}`$ and the off-diagonal ones $`t_{23},t_{32}`$. Then the U-quark mass matrix would have the general form in leading order: $$M_{bc}^U=(S+S^{})^{1/2}v\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& h_{123}^0\\ 0& h_{132}^0& h_{133}^0\end{array}\right)$$ (5.10) Plugging the above expressions for the Yukawa couplings one finds for $`\|t_{33}\|>>\|t_{32}\|`$ three U-quark eigenstates with masses of order $`gv(0,1/\delta ,\delta )`$, with $`\delta =\|t_{33}/t_{32}\|^{1/2}`$, yielding a hierarchical structure for the U-quarks. Thus the structure of the U-quarks Yukawa couplings for this class of models is in principle sufficiently rich to accommodate the observed hierarchy, at least as long as off-diagonal moduli have modest hierarchical vevs. In the particular setting of toroidal $`Z_3`$ orbifolds discrete torsion seems also to be required. The same discussion applies directly to the case of the left-right symmetric model of section 3.4, the only difference being that in this case the Yukawa couplings $`h_{abc}Q_L^aQ_R^bH^c`$ give masses both to U-quarks and D-quarks. For both a hierarchical structure should be possible to be accommodated. The rest of the phenomenologically relevant renormalizable Yukawa couplings are of type b) and c), i.e., they involve directly the D7-brane sectors and are thus more compactification dependent. Let us describe some qualitative features for the SM and the LR models in turn. In the case of the SM of section 3.3 (and its compact versions), there are renormalizable couplings of type $`(33)_i(7_i3)(37_i)`$, where $`(33)_i`$ denotes the chiral fields in the $`(33)`$ sector associated to the i-th complex plane. Looking at the quantum numbers in Table 1 one sees that in particular there are couplings of type $`(3,2)(\overline{3},1;2^{})(1,2;2^{})`$ which provide Yukawa couplings for the D-quarks. One of the doublets inside $`(1,2;2^{})`$ should be identified with the normal left-handed leptons and the other with Higgs fields. By definition, the doublet appearing in the D-quark Yukawa coupling will be identified with the Higgs field. As pointed out in the footnote at the end of section 3.2, in generic D7-brane configurations fields from a given $`37`$ sector can couple to $`33`$ fields from all three complex planes. Thus this flexibility will also in general allow for a hierarchy of D-quark masses. Concerning lepton masses, they do not appear at the renormalizable level since all left- and right-handed leptons as well as the Higgs field belong to $`(37_i)`$ sectors and there are no $`(37)^3`$ type of couplings on the disk. They may however appear from non-renormalizable couplings of type $`(37)^n`$ in compact models if some standard model singlets in some other $`(37)`$ sectors get a vev. This will be very model dependent. In the case of the left-right symmetric model of section 3.4 (and its compact versions), the $`(33)^3`$ type of couplings already give masses to both U- and D-quarks. Concerning leptons, all anomaly free gauge symmetries allow for couplings of type $`(33)(37)(73)`$ including terms transforming like $`(1,2,2)(1,2,1)(1,1,2)`$ which would give rise to standard Dirac masses for leptons. Looking at Table 2 one observes that those couplings are in fact forbidden by some anomalous $`U(1)`$ charges. Nevertheless one expects those couplings to be present once appropriate twisted closed string insertions (which are charged under anomalous $`U(1)`$’s) are made . Again, for different geometrical configurations of D7-branes there will be enough flexibility to obtain a hierarchy of Yukawa couplings also for leptons. In the case of the left-right model we also have to worry about neutrino masses. Generically they will get Dirac masses of the order of the charged lepton masses. However, once $`SU(2)_R`$ symmetry-breaking takes place through the vevs of the extra $`SU(2)_R`$ doublets discussed above, the right-handed neutrinos may combine with some singlets in $`(7_i7_j)`$ sectors and get masses of order $`M_R`$. We will not give details here which would be very model dependent. Let us note however that the effective left-handed neutrino masses would thus be of order $`m_l^2/M_R`$, where $`m_l`$ are the neutrino Dirac masses. For values of $`m_l`$ of order the electron mass and $`M_R`$ of order 1 TeV, they would have masses of order 1 eV, compatible with experimental bounds. Let us finally comment about proton stability in these schemes. This is an important question if, as seems to be indicated by gauge coupling unification, the string scale is chosen to be well below the Planck mass (e.g. at scales of order $`10^{12}`$ GeV or below). Since for this question higher dimensional operators are in principle important, it is not possible to make model independent statements about sufficient proton stability. We would like to underline however, that from this point of view the left-right symmetric models seem more promising since the symmetry $`BL`$ is gauged and hence dimension four baryon or lepton violating operators are forbidden to start with. Furthermore, B/L-violating dimension 5 operators like F-terms of the form $`Q_LQ_LQ_LL`$ (which respect $`BL`$) are of type $`(33)^3(73)`$ and hence are forbidden on the disk. In fact, it was shown in ref. , that in a left-right model analogous to the compact orientifold of chapter 5, the proton is absolutely stable due to the presence of induced baryon and lepton parity $`ZZ_2`$ symmetries. Thus we believe that, at least in the left-right class of models, a sufficiently stable proton may naturally be achieved. iii) String Scale versus Planck Mass In our approach, four-dimensional gravity arises once we embed our SM or LR D-brane systems into a finite compact space. The precise relationship between the string scale $`M_s`$ and four-dimensional Planck scale $`M_p`$ will depend on the compactification. In toroidal orbifold compactifications with factorized two-tori they are related by $$M_p=\frac{2\sqrt{2}M_s^4}{\lambda M_1M_2M_3}$$ (5.11) where $`M_i`$ are the compactification mass scales of the three tori. Thus, as is well known , values for $`M_s`$ much below the Planck scale may be possible if some of the compact dimensions are very large. Until we get a handle on the dynamics fixing the radii, the value of the string scale is only constrained by accelerator limits to be $`M_s>1`$ TeV or so. In particular, one can accommodate intermediate scale values of order $`M_s10^{10}10^{12}`$ GeV, as suggested by the logarithmic gauge coupling unification arguments above. Similar qualitative statements can be made concerning the Planck versus string scales in F-theory models with the SM sector living on the world-volume of D3-branes, like those discussed in chapter 5. Thus, e.g. one can match the observed size of the Planck scale by making very large the volume of the base $`P_1`$ in the first K3 $`X^{(1)}`$, while maintaining the volume of the second K3, $`X^{(2)}`$ of order the string scale. It is quite interesting to remark that such intermediate values for the string scale turn out to be also suggested in a completely independent manner in models containing anti-branes like the orientifolds in chapter 5. Indeed, in those models supersymmetry is broken at the string scale $`M_s`$ due to the presence of anti-D7-branes. However the physical SM or LR sectors are away from these antibranes in the $`Y_1`$ transverse dimension. Let us suppose that e.g. one has $`M_2,M_3=M_s`$ and $`M_1=M_s^2/(\lambda M_p)`$ so that one obtains the correct Planck scale. This means taking very large $`Y_1`$ radii which will also mean that the susy-breaking effects from the anti-D7-branes will be felt in a Planck mass suppressed manner in the visible SM world. Thus one expects effective SUSY-breaking contributions of order $`M_s^2/M_p`$ in the SM sector. If $`M_s`$ is of order $`10^{10}10^{12}`$ GeV, then these SUSY-breaking contributions would be of order the weak scale, as required. The structure of these models would be very similar to hidden-sector supersymmetry-breaking models, the main difference being that in the present case the string and SUSY-breaking scales do coincide. Let us also comment that the approach here presented is also consistent with low (i.e. 1-10 TeV) string scale scenarios , just by choosing larger compact radii. In these cases an alternative understanding of the gauge coupling unification problem should be found. Finally the $`F`$-theory constructions of the previous section allows for the possibility of a more standard scenario. Being supersymmetric, we do not have to identify the string scale with the supersymmetry breaking scale. We may in principle conceive models of this type for which the string scale is closer to the Planck or GUT scales. The left-right symmetric model presented in that section still requires unification at the intermediate scale but in general, achieving or not gauge coupling unification will clearly depend on the spectrum of each model. ## 6 Final Comments and Outlook We have presented what we consider to be the first steps on a new way to deal with string model building. The fact that D3-brane worlds can appear naturally on type IIB string theory suggests that looking directly to the physics on these D3-branes may tell us many interesting phenomenological aspects that depend very weakly on the details of the compactified space. This is so since only gravitational strength interactions feel the extra dimensions. It would be interesting to extend this bottom-up construction of realistic models to other string theories, for instance using type 0B D3-branes on singularities , or to M-theory. The latter framework was explored in , where an explicit M5-brane configuration led to a $`N=2`$ supersymmetric (and therefore non-chiral) toy model with SM group and correct gauge non-abelian and $`U(1)`$ quantum numbers for some of the standard model fields. It is remarkable how easy it is to obtain realistic models from this approach and how powerful the structure of singularities turn out to be in order to achieve interesting phenomenological properties, such as the number of families, the ubiquity of hypercharge and the value of Weinberg’s angle. For instance a concrete prediction of the class of models considered is that the number of families does not exceed 3 for supersymmetric models and 4 for non-supersymmetric ones. The reason being essentially the $`𝒩=4`$ symmetry underlying the 33 spectrum, which limits the number of replicated fermions after the orbifold projection. This is a very small (and realistic) number compared with more standard compactifications in which the number of families depends on the topology of the compact manifold and is naturally very large. It is worth remarking that the consideration of the non-compact models constructed here may have an importance per se. At present there is increasing interest on non-compact extra dimensions raised after the work of Randall and Sundrum . It would be tempting to speculate that a similar mechanism to localize gravity on the D3-branes may be present in models similar to ours. However it is not clear if this could be achievable. Furthermore, during the past few years, non-compact brane models were used in order to obtain information on the dynamics of gauge theories from the structure of branes. Much work has been done in this direction, , although most of the activity was based on extended supersymmetric models and non-chiral $`𝒩=1`$ models. Chiral theories were constructed in , and seen to be related to theories on D3-branes at singularities . Our models extend this program by addressing (and actually succeeding in) the construction of realistic models using brane configurations. Independent of these issues, the fact that in compact models we have been able to obtain realistic $`𝒩=0`$ and $`𝒩=1`$ models from type IIB strings at singularities may be the most important concrete result of this work. We have explored several mechanisms and generalizations to current model building, each of them deserving further study. The new approach, while sharing some of the good properties of the constructions of (gauge unification, supersymmetry breaking, proton stability), has several clear advantages over previous discussions, besides the fact that it is more general: from the stability of the models, by trapping anti-branes avoiding brane/anti-brane annihilation, to the inclusion of new degrees of freedom, such as discrete torsion, which allow more flexibility on the structure of fermion masses and helps to stabilize the models. A scenario with intermediate fundamental scale is clearly favoured by the non supersymmetric models which can also have gauge coupling unification at the same scale, especially in the case of the left-right symmetric models. On the other hand, the fact that we have also obtained quasi-realistic $`𝒩=1`$ supersymmetric models allows for the possibility of realistic scenarios with fundamental scales as large as the Planck or GUT scales. It is worth pointing out however that, as usual in string theory, grand unified gauge models do not seem to appear naturally in this scheme, whereas the standard model and its left-right symmetric extension are very easy to obtain. There are clearly many avenues that deserve to be explored in the future, starting from our approach. Probably the most pressing one being the understanding of moduli and dilaton stabilization after supersymmetry breaking and the corresponding value of the vacuum energy. Acknowledgements We thank M. Cvetic, S. Franco, M. Klein, P. Langacker, J. Maldacena, B. Ovrut, M. Plümacher and R. Rabadán for useful discussions. A. M. U. thanks the Department of Applied Mathematics and Theoretical Physics of Cambridge University, and the Department of Physics and Astronomy of University of Pennsylvania for hospitality, and M. González for encouragement and support. G. A. thanks the High Energy Theory Group of Harvard University for hospitality. This work has been partially supported by CICYT (Spain), the European Commission (grant ERBFMRX-CT96-0045) and PPARC. G.A work is partially supported by ANPCyT grant 03-03403. Appendices ## Appendix A Orbifold Singularities with Discrete Torsion Consider a $`\mathrm{IC}^3/(ZZ_{M_1}\times ZZ_{M_2})`$ orbifold singularity, with the two factors of the orbifold group generated by twists $`g_1`$, $`g_2`$. The effect of discrete torsion in the closed string sector was analyzed in (see also ). It amounts to the introduction of an additional phase $`ϵ(g_1,g_2)`$ in the action of $`g_1`$ on $`g_2`$-twisted sector states. This phase must be an $`s^{th}`$ root of unity, where $`s=\mathrm{gcd}(M_1,M_2)`$. Hence discrete torsion modifies the twisted sector spectrum, and introduces relative phases among the different contributions to the torus partition function. The effect of discrete torsion on open string sector, and consequently on systems of D-branes at singularities, has been analyzed only recently (see for orientifold examples). In the presence of discrete torsion the embedding of the orbifold twists in the Chan-Paton indices forms a projective representation of the orbifold group , i.e. the Chan-Paton embedding matrices $`\gamma _g`$ obey the group law up to phases, concretely $`\gamma _{g_1}\gamma _{g_2}=ϵ(g_1,g_2)\gamma _{g_2}\gamma _{g_1}`$ (A.1) The field theory on a set of D3-branes at such an orbifold singularity with discrete torsion has been discussed in detail in . In the following we center on a particular case which illustrates the main features. Consider a $`\mathrm{IC}^3/(ZZ_N\times ZZ_M\times ZZ_M)`$ singularity, where the generator $`\theta `$ of $`ZZ_N`$ acts on $`\mathrm{IC}^3`$ through the twist $`(a_1,a_2,a_3)/N`$, and the generators $`\omega _1`$, $`\omega _2`$ of the $`ZZ_M`$’s act through the twists $`v_1=(1,0,1)/M`$, $`v_2=(0,1,1)/M`$ <sup>15</sup><sup>15</sup>15In certain cases, for instance when $`N`$ and $`M`$ are coprime, the group is actually $`ZZ_{NM}\times ZZ_M`$. However, it will be convenient to make the $`ZZ_M\times ZZ_M`$ action more manifest.. Let us consider the case of discrete torsion $`ϵ(\omega _1,\omega _2)=e^{2\pi i/M}`$ in the $`ZZ_M\times ZZ_M`$ part. Consequently we choose the following embedding in the Chan-Paton indices $`\gamma _{\theta ,3}`$ $`=`$ $`\mathrm{diag}(I_{Mn_0},e^{2\pi i\frac{1}{N}}I_{Mn_1},\mathrm{},e^{2\pi i\frac{N1}{N}}I_{Mn_{N1}})`$ $`\gamma _{\omega _1,3}`$ $`=`$ $`\mathrm{diag}(1,e^{2\pi i\frac{1}{M}},\mathrm{},e^{2\pi i\frac{M1}{M}})\mathrm{diag}(I_{n_0},I_{n_1},\mathrm{},I_{n_{N1}})`$ $`\gamma _{\omega _2,3}`$ $`=`$ $`\left(\begin{array}{cccc}& 1& & \\ & & \mathrm{}& \\ & & & 1\\ 1& & & \end{array}\right)\mathrm{diag}(I_{n_0},I_{n_1},\mathrm{},I_{n_{N1}})`$ (A.2) Notice that $`\gamma _{\omega _1,3}\gamma _{\omega _2,3}=e^{2\pi i\frac{1}{M}}\gamma _{\omega _2,3}\gamma _{\omega _1,3}`$, in agreement with (A.1). The spectrum is obtained after imposing the projections $`\begin{array}{cccc}\mathrm{Vector}\hfill & \lambda =\gamma _{\theta ,3}\lambda \gamma _{\theta ,3}^1\hfill & \lambda =\gamma _{\omega _1,3}\lambda \gamma _{\omega _1,3}^1\hfill & \lambda =\gamma _{\omega _2,3}\lambda \gamma _{\omega _2,3}^1\hfill \\ \mathrm{Chiral}_1\hfill & \lambda =e^{2\pi i\frac{a_1}{N}}\gamma _{\theta ,3}\lambda \gamma _{\theta ,3}^1\hfill & \lambda =e^{2\pi i\frac{1}{M}}\gamma _{\omega _1,3}\lambda \gamma _{\omega _1,3}^1\hfill & \lambda =\gamma _{\omega _2,3}\lambda \gamma _{\omega _2,3}^1\hfill \\ \mathrm{Chiral}_2\hfill & \lambda =e^{2\pi i\frac{a_2}{N}}\gamma _{\theta ,3}\lambda \gamma _{\theta ,3}^1\hfill & \lambda =\gamma _{\omega _1,3}\lambda \gamma _{\omega _1,3}^1\hfill & \lambda =e^{2\pi i\frac{1}{M}}\gamma _{\omega _2,3}\lambda \gamma _{\omega _2,3}^1\hfill \\ \mathrm{Chiral}_3\hfill & \lambda =e^{2\pi i\frac{a_3}{N}}\gamma _{\theta ,3}\lambda \gamma _{\theta ,3}^1\hfill & \lambda =e^{2\pi i\frac{1}{M}}\gamma _{\omega _1,3}\lambda \gamma _{\omega _1,3}^1\hfill & \lambda =e^{2\pi i\frac{1}{M}}\gamma _{\omega _2,3}\lambda \gamma _{\omega _2,3}^1\hfill \end{array}`$ (A.7) The projection process is relatively simple. For vector multiplets, the $`\theta `$ projection yields a gauge group $`_{i=0}^{N1}U(Mn_i)`$. The $`\omega _1`$ projection splits it into $`_{i=0}^{N1}U(n_1)^N`$. Finally, the $`\omega _2`$ action identifies the different gauge factors with equal rank, leaving a remaining gauge group $`_{i=0}^{N1}U(n_i)`$. Proceeding analogously, the complete spectrum is $`\mathrm{Vector}`$ $`{\displaystyle \underset{i=0}{\overset{N1}{}}}U(n_i)`$ $`\mathrm{Chiral}`$ $`{\displaystyle \underset{r=1}{\overset{3}{}}}{\displaystyle \underset{i=0}{\overset{N1}{}}}(n_i,\overline{n}_{i+a_r})`$ (A.8) Notice this is identical to the spectrum of the $`\mathrm{IC}^3/ZZ_N`$ singularity with twist $`(a_1,a_2,a_3)/N`$, in section 2.1 or 3.1. However, the effect of the additional $`ZZ_M\times ZZ_M`$ twist and of the discrete torsion are manifest in the superpotential, which reads $`W`$ $`=`$ $`Tr[\mathrm{\Phi }_{i,i+a_1}^1\mathrm{\Phi }_{i+a_1,i+a_1+a_2}^2\mathrm{\Phi }_{i+a_1+a_2,i}^3e^{2\pi i\frac{1}{M}}\mathrm{\Phi }_{i,i+a_1}^1\mathrm{\Phi }_{i+a_1,i+a_1+a_3}^3\mathrm{\Phi }_{i+a_1+a_3,i}^2]`$ (A.9) which differs from the 33 piece in (3.1). We spare the reader the detailed discussion of the introduction of D7-branes in the configuration. The resulting spectrum is identical to that obtained for $`ZZ_N`$ singularities in Section 3.1. Notice that, since the twists $`\omega _1`$, $`\omega _2`$ have traceless Chan-Paton embeddings, the corresponding disk tadpoles vanish. The only constraints on the integers $`n_i`$ arise from cancellation of tadpoles in $`ZZ_N`$ twisted sectors. The corresponding conditions are those for $`ZZ_N`$ singularities (as expected, since they are basically the anomaly cancellation conditions). ## Appendix B Non-Abelian Orbifolds The rules to compute the spectrum of D3-branes at non-abelian orbifold singularities have been discussed in for $`\mathrm{\Gamma }SU(2)`$, and generalized in . In the following we center on $`𝒩=1`$ supersymmetric field theories in the case $`\mathrm{\Gamma }SU(3)`$, studied in detail in (see for the general case $`\mathrm{\Gamma }SU(4)`$). The computation of the spectrum is just a natural group-theoretical generalization of the abelian case. Let $`\{_i\}_{i=1,\mathrm{},N}`$ denote the irreducible representations of $`\mathrm{\Gamma }`$, with $`N`$ the number of conjugacy classes of $`\mathrm{\Gamma }`$. We are interested in the field theory on a set of D3-branes at a $`\mathrm{IC}^3/\mathrm{\Gamma }`$ singularity, where the action of $`\mathrm{\Gamma }`$ on $`\mathrm{IC}^3`$ is specified by a three-dimensional representation $`^{(\mathrm{𝟑})}`$, and its action on the D3-brane Chan-Paton indices is given by a representation $`^{CP}`$, which decomposes as $`^{CP}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}n_i_i`$ (B.1) The field theory on D3-branes at $`\mathrm{IC}^3/\mathrm{\Gamma }`$ is obtained by projecting the field theory on D3-branes on $`\mathrm{IC}^3`$ onto states invariant under the combined (geometric plus Chan-Paton) action of $`\mathrm{\Gamma }`$. The resulting gauge group <sup>16</sup><sup>16</sup>16In the literature on D-branes on non-abelian orbifold, the representation $`^{CP}`$ is usually taken to be $`k`$ copies of the regular representation, $`n_i=k\mathrm{dim}_i`$, which satisfies tadpole cancellation conditions without the need of additional branes. We will consider more general choices of $`n_i`$, with the understanding that additional D7-branes will cancel the corresponding tadpoles. Unfortunately, the inclusion of D7-branes has not been discussed in the literature, hence we will not be explicit about the 37 sector in our models, and simply assume it provides additional chiral multiplets to cancel gauge theory anomalies from the 33 sector. is $`_{i=1}^NU(n_i)`$. There are also $`a_{ij}^\mathrm{𝟑}`$ chiral multiplets in the $`(n_i,\overline{n}_j)`$ representation, where the adjacency matrix $`a_{ij}^\mathrm{𝟑}`$ is defined by the decomposition $`^{(\mathrm{𝟑})}_i`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{N}{}}}a_{ij}^\mathrm{𝟑}_j`$ (B.2) The superpotential is obtained by substituting the surviving chiral multiplets $`\mathrm{\Phi }_{ij}`$ in the $`𝒩=4`$ superpotential (2.1), but we will not need it explicitly. An important observation is that the gauge couplings for the different factors are not equal in field theories from non-abelian orbifolds. The gauge coupling for the $`i^{th}`$ group is given by $`\tau _i={\displaystyle \frac{r_i\tau }{\|\mathrm{\Gamma }\|}}`$ (B.3) with $`r_i=\mathrm{dim}_i`$, and $`\tau `$ an overall value independent of $`i`$. The analysis of $`U(1)`$ anomalies has not been performed in the literature, but can be easily generalized from the abelian case <sup>17</sup><sup>17</sup>17We thank A. Hanany for conversations on this point.. Assuming the 37, 73 sectors contribute a set of bifundamental multiplets which cancel the non-abelian anomaly in the 33 sector, the mixed $`Q_{n_i}`$-$`SU(n_j)^2`$ anomaly reduces to $`A_{ij}=\frac{1}{2}n_i(a_{ij}^{(\mathrm{𝟑})}a_{ji}^{(\mathrm{𝟑})})`$. It is possible to check that the diagonal combination $`Q_{diag}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{r_i}{n_i}}Q_{n_i}`$ (B.4) is automatically non-anomalous. Also, as long as no $`n_i`$ vanishes, there is one additional non-anomalous $`U(1)`$’s for each non-trivial conjugacy classes leaving some complex plane invariant. An important difference with respect to the abelian case is that, due to the presence of the factors $`r_i`$ above, the structure $`SU(3)\times SU(2)`$ does not in general guarantee the correct hypercharge assignments for the non-anomalous $`U(1)`$’s. ## Appendix C Non-Orbifold Singularities In this section we briefly discuss systems of D3-branes at non-orbifold singularities. The only known systematic approach to construct the world-volume field theory is to realize the non-orbifold singularity as a partial resolution of a suitable orbifold singularity <sup>18</sup><sup>18</sup>18Alternative approaches include the direct construction of field theories with the right symmetries and the correct moduli space , and the use of T-duality to configurations of NS-branes and D4-branes .. However, the identification of the field-theory interpretation of the partial blowing-up is very involved beyond the simplest examples . Consequently, it is difficult to make general statements about the phenomenological potential of these systems. We would like to center on a particular family of singularities, which can be constructed as orbifolds of the conifold singularity (see also ), and their partial blow-ups. They have potential interest since generically they lead to chiral $`𝒩=1`$ supersymmetric field theories on the world-volume of D3-branes, and include abelian orbifold singularities as particular cases. Moreover, they have a T-dual representation in terms of configurations of NS-fivebranes and D5-branes, known as ‘brane diamonds’ , generalizing the brane box constructions in , which allow a systematic search of interesting field theories within this class. Instead of extending on such analysis, we point out that the only singularity in this family, leading to full triplication of chiral multiplet content in the 33 sector is in fact the $`\mathrm{IC}^3/ZZ_3`$ singularity. A milder requirement would be to have triplication not for the full 33 sector, but at least for some representation. Besides the $`ZZ_3`$ orbifold, there is only one singularity in this family fulfilling this milder requirement. Let us briefly discuss its construction. Consider the conifold, given by the hypersurface $`xy=zw`$ in $`\mathrm{IC}^4`$. The field theory on D3-branes at a conifold was determined in . The gauge group is $`U(N_1)\times U(N_2)`$ and there are chiral multiplets $`A_1`$, $`A_2`$ in the $`(\text{ }\text{ }\text{ }\text{ }\text{ },\overline{\text{ }\text{ }\text{ }\text{ }\text{ }})`$, and $`B_1`$, $`B_2`$ in the $`(\overline{\text{ }\text{ }\text{ }\text{ }\text{ }},\text{ }\text{ }\text{ }\text{ }\text{ })`$. There is also a superpotential $`W=\mathrm{Tr}(A_1B_1A_2B_2)\mathrm{Tr}(A_1B_2A_2B_1)`$ (C.1) We are going to consider a quotient of this variety by a $`ZZ_3`$ twist with generator $`\theta `$ acting as $`xe^{2\pi i/3}x;ze^{2\pi i/3}z;ye^{2\pi i/3}y;we^{2\pi i/3}w`$ (C.2) which preserves $`𝒩=1`$ supersymmetry on the D3-brane world-volume. The geometric action is reflected in the field theory as $`A_1e^{\pi i}A_1;B_1e^{\pi i/3}B_1;A_2e^{\pi i/3}A_2;B_2e^{\pi i/3}B_2`$ (C.3) Finally, the action of $`\theta `$ may be embedded on the $`U(N_1)`$, $`U(N_2)`$ gauge degrees of freedom, through the matrices $`\gamma _{\theta ,3}^{(1)}=\mathrm{diag}(I_{n_0},e^{2\pi i\frac{1}{3}}I_{n_1},e^{2\pi i\frac{2}{3}}I_{n_2})`$ ; $`\gamma _{\theta ,3}^{(2)}=\mathrm{diag}(e^{\pi i}I_{n_0^{}},e^{\pi i\frac{5}{3}}I_{n_1^{}},e^{\pi i\frac{1}{3}}I_{n_2^{}})`$ (C.4) The field theory on D3-branes at the $`ZZ_3`$ orbifold of the conifold is obtained by projecting the conifold field theory onto states invariant under the combined (geometric plus Chan-Paton) action of $`\theta `$. The resulting spectrum is $`U(n_0)\times U(n_0^{})\times U(n_1)\times U(n_1^{})\times U(n_2)\times U(n_2^{})`$ $`\begin{array}{cccc}a_0:(n_0,\overline{n}_0^{}),& b_2:(n_2,\overline{n}_0^{}),& c_0:(\overline{n}_1,n_0^{}),& d_0:(\overline{n}_1,n_0^{})\\ a_1:(n_1,\overline{n}_1^{}),& b_0:(n_0,\overline{n}_1^{}),& c_1:(\overline{n}_2,n_1^{}),& d_1:(\overline{n}_2,n_1^{})\\ a_2:(n_1,\overline{n}_1^{}),& b_1:(n_1,\overline{n}_2^{}),& c_2:(\overline{n}_0,n_2^{}),& d_2:(\overline{n}_0,n_2^{})\end{array}`$ (C.8) with superpotential $`W`$ $`=`$ $`\mathrm{Tr}[a_0c_0b_1d_2a_0d_0b_1c_2+a_1c_1b_2d_0a_1d_1b_2c_0+a_2c_2b_0d_1a_2d_2b_0c_1]`$ (C.9) In general the gauge anomalies in this theory are non-vanishing, but a suitable set of D7-branes can be introduced to yield the configuration consistent. Assuming a 37 sector with an arbitrary set of bifundamental representations, it is possible to perform the analysis of $`U(1)`$ anomalies. In general, mixed anomalies do not vanish, but have a nice factorized structure which suggests they are cancelled through a Green-Schwarz mechanism mediated by closed string twisted modes, generalizing the observation in to these non-orbifold singularities. We also conclude that the only non-anomalous $`U(1)`$ in this example is the diagonal combination $`Q_{diag}=_{i=0}^2(Q_{n_i}/n_i+Q_{n_i^{}}/n_i^{})`$. The singularity leading to field theories with triplication of at least some chiral multiplet is obtained as a partial blow-up of this $`ZZ_3`$ quotient of the conifold. The partial resolution is manifested in the field theory as an F-flat direction, which is not D-flat by itself, but requires turning on a Fayet-Illiopoulos term, i.e. turning on a vev for blow-up modes. For instance, when $`n_0=n_0^{}`$, $`n_2=n_2^{}`$ there is one such direction corresponding to diagonal vevs $`a_0=v_0`$, $`a_2=v_2`$, along which the field theory is reduced to $`U(n_0)\times U(n_1)\times U(n_1^{})\times U(n_2)`$ $`\begin{array}{ccc}a_1:(n_1,\overline{n}_1^{});& b_0:(n_0,\overline{n}_1^{});& c_0,d_0:(\overline{n}_1,n_0)\\ b_1:(n_1,\overline{n}_2);& c_1,d_1:(\overline{n}_2,n_1^{});& b_2,c_2,d_2:(n_2,\overline{n}_0)\end{array}`$ (C.12) with superpotential $`W=v_0\mathrm{Tr}(c_0b_1d_2d_0b_1c_2)+\mathrm{Tr}(a_1c_1b_2d_0a_1d_1b_2c_0)+v_2\mathrm{Tr}(c_2b_0d_1d_2b_0c_1)`$ (C.13) This corresponds to a partial blow-up to a certain singular variety $`X`$, whose form could be determined from the above choice of expectation values. It is possible to follow the above construction in the presence of D7-branes, but since the computation is rather involved, we just quote the result. One may introduce four kinds of D7-branes, containing a total of ten unitary group factors, with ranks denoted $`u_1`$, and $`v_i`$, $`w_i`$, $`x_i`$. The 37, 73 sectors are $`\begin{array}{ccccccccccc}(n_1^{},\overline{u}_1)& ,& (u_1,\overline{n}_1)& & a_1& ;& (n_1^{},\overline{v}_1)& ,& (v_1,\overline{n}_0)& & b_0\\ (n_2,\overline{v}_2)& ,& (v_2,\overline{n}_1)& & b_1& ;& (n_0,\overline{v}_3)& ,& (v_3,\overline{n}_2)& & b_2\\ (n_1,\overline{w}_1)& ,& (w_1,\overline{n}_0)& & c_0& ;& (n_2,\overline{w}_2)& ,& (w_2,\overline{n}_1^{})& & c_1\\ (n_0,\overline{w}_3)& ,& (w_3,\overline{n}_2)& & c_2& ;& (n_1,\overline{x}_1)& ,& (x_1,\overline{n}_0)& & d_0\\ (n_2,\overline{x}_2)& ,& (x_2,\overline{n}_1^{})& & d_1& ;& (n_0,\overline{x}_3)& ,& (x_3,\overline{n}_2)& & d_2\end{array}`$ (C.19) we have also indicated by an arrow the 33 field to which the corresponding 37, 73 fields couple in the superpotential. Tadpole conditions cannot be computed directly, but we assume they amount to cancellation of non-abelian anomalies. This should also guarantee that $`U(1)`$ anomalies cancel by the Green-Scharz mechanism. In this case there are two non-anomalous $`U(1)`$’s generated by $`Q_{diag}={\displaystyle \frac{Q_{n_0}}{n_0}}+{\displaystyle \frac{Q_{n_1}}{n_1}}+{\displaystyle \frac{Q_{n_1^{}}}{n_1^{}}}+{\displaystyle \frac{Q_{n2}}{n_2}};Q^{}={\displaystyle \frac{2}{n_1}}Q_{n_1}{\displaystyle \frac{Q_{n_1^{}}}{n_1^{}}}{\displaystyle \frac{Q_{n_2}}{n_2}}`$ (C.20) When $`n_1=n_1^{}`$, a further diagonal vev $`a_1=v_1`$ corresponds to a blow-up of $`X`$ to the $`\mathrm{IC}^3/ZZ_3`$ singularity. Indeed, the resulting field theory agrees with (3.13), after the replacements $`b_i,c_i,c_i\mathrm{\Phi }_{i,i+1}^r`$, $`v_i,w_i,x_iu_i^r`$. The 33 piece of the superpotential $`W=v_0\mathrm{Tr}(b_1d_2c_0b_1c_2d_0)+v_1\mathrm{Tr}(b_2d_0c_1b_2c_0d_1)+v_2\mathrm{Tr}(b_0d_1c_2b_0c_1d_2)`$ (C.21) is more flexible than the 33 piece in (3.14), a fact that may have phenomenological applications. ## Appendix D Non-Supersymmetric Models from Antibranes The computation of the spectrum in a set of D3/$`\overline{\mathrm{D3}}`$-branes at orbifold singularities in the presence of D7/$`\overline{\mathrm{D7}}`$-branes can be extracted from (for a recent review on non-supersymmetric brane configurations with references to previous material see ). The only differences with respect to the case without antibranes is the opposite GSO projections in the brane-antibrane sectors. For instance, in $`\overline{\mathrm{𝟑}}\overline{\mathrm{𝟑}}`$ sectors, GSO projection is as usual and the spectrum is analogous to that in 33 sectors. We consider a generically non-susy $`ZZ_N`$ singularity. Let us introduce D3, $`\overline{\mathrm{D3}}`$ Chan-Paton embeddings $`\gamma _{\theta ,3}`$, $`\gamma _{\theta ,\overline{3}}`$ with $`n_j`$, $`m_j`$ eigenvalues $`e^{2\pi ij/N}`$. Consider the $`\mathrm{𝟑}\overline{\mathrm{𝟑}}+\overline{\mathrm{𝟑}}\mathrm{𝟑}`$ sector. In the NS sector the complex tachyon, $`\lambda \|0`$ survives the GSO, whereas the massless states $`\mathrm{\Psi }_{\frac{1}{2}}\|0`$ do not. The orbifold projection on the Chan-Paton factors for the tachyons is $`\lambda =\gamma _{\theta ,3}\lambda \gamma _{\theta ,\overline{3}}^1`$ $`\lambda =\gamma _{\theta ,\overline{3}}\lambda \gamma _{\theta ,3}^1`$ (D.1) leading to complex tachyons in the representation $`_{i=0}^{N1}[(n_i,\overline{m}_i)+(m_i,\overline{n}_i)]`$. In the R sector, we obtain states $`\lambda \|s_1,s_2,s_3,s_4`$, with $`s_i=\pm \frac{1}{2}`$ and $`_is_i=\mathrm{even}`$. We get right-handed spacetime fermions $`s_4=\frac{1}{2}`$ with the projection $`\begin{array}{ccc}\lambda _{3\overline{3}}=e^{2\pi ia_\alpha /N}\gamma _{\theta ,3}\lambda _{3\overline{3}}\gamma _{\theta ,\overline{3}}^1\hfill & & \lambda _{\overline{3}3}=e^{2\pi ia_\alpha /N}\gamma _{\theta ,\overline{3}}\lambda _{\overline{3}3}^{(a)}\gamma _{\theta ,3}^1\hfill \end{array}`$ (D.3) They transform in the representation $`_{i=0}^{N1}[(n_i,\overline{m}_{i+a_\alpha })+(m_i,\overline{n}_{i+a_\alpha })]`$. The $`\mathrm{𝟑}\overline{\mathrm{𝟕}_𝐫}+\overline{\mathrm{𝟕}}_𝐫\mathrm{𝟑}`$ and $`\overline{\mathrm{𝟑}}\mathrm{𝟕}_𝐫+\mathrm{𝟕}_𝐫\overline{\mathrm{𝟑}}`$ sectors are analogous to the $`\mathrm{𝟑𝟕}_𝐢+\mathrm{𝟕}_𝐢\mathrm{𝟑}`$ sector, with a few modifications due only to the opposite GSO projections. Consider introducing $`\overline{\mathrm{D7}}_3`$-branes, with $`\gamma _{\theta ,7_3}`$ having $`v_j`$ eigenvalues $`e^{2\pi ij/N}`$, and let us center on the case $`b_3=\mathrm{even}`$. Scalars arise from the NS sector, which has fermion zero modes along the DN directions, $`Y_1`$, $`Y_2`$. Massless states have the form $`\lambda \|s_1,s_2`$, with $`s_1=s_2=\pm \frac{1}{2}`$. The orbifold projection for $`\|\frac{1}{2},\frac{1}{2}`$ gives $`\lambda _{3\overline{7}_3}=e^{2\pi i(b_1b_2)/2}\gamma _{\theta ,3}\lambda _{3\overline{7}_3}\gamma _{\theta ,\overline{7}_3}^1`$ $`\lambda _{\overline{7}_33}=e^{2\pi i(b_1b_2)/2}\gamma _{\theta ,\overline{7}_3}\lambda _{\overline{7}_33}\gamma _{\theta ,3}^1`$ (D.4) We obtain complex scalars in the representation $`_{i=0}^{N1}[(n_i,\overline{v}_{i+\frac{1}{2}(b_1b_2)})+(v_i,\overline{n}_{i+\frac{1}{2}(b_1b_2)})]`$. Spacetime fermions arise from the R sector, which has fermion zero modes along the DD and NN directions, $`Y_3`$, $`Y_4`$. Massless states are $`\lambda \|s_3;s_4`$, with $`s_3=s_4=\pm \frac{1}{2}`$. The orbifold projection for $`s_4=\frac{1}{2}`$ gives $`\lambda _{3\overline{7}_3}=e^{2\pi ib_3/2}\gamma _{\theta ,3}\lambda _{3\overline{7}_3}\gamma _{\theta ,\overline{7}_3}^1`$ $`\lambda _{\overline{7}_33}=e^{2\pi ib_3/2}\gamma _{\theta ,\overline{7}_3}\lambda _{\overline{7}_33}\gamma _{\theta ,3}^1`$ (D.5) We obtain right-handed fermions in the representation $`_{i=0}^{N1}[(n_i,\overline{v}_{i\frac{1}{2}b_3})+(v_i,\overline{n}_{i\frac{1}{2}b_3})]`$. The spectrum in the $`\overline{3}7_3`$, $`7_3\overline{3}`$ is computed analogously. The analysis of gauge anomalies in the presence of antibranes can be performed in analogy with that in sections 2.2 and 2.3 . The non-abelian gauge anomalies are equivalent to twisted RR tadpole cancellation conditions, which have the form (2.27) with the replacement $`\mathrm{Tr}\gamma _{\theta ^k,p}\mathrm{Tr}\gamma _{\theta ^k,p}\mathrm{Tr}\gamma _{\theta ^k,\overline{p}}`$. This reflects the opposite RR charges of branes and antibranes. Concerning $`U(1)`$ anomalies, one can show the existence of two non-anomalous $`U(1)`$ per twist with fixed planes, associated to combinations (2.30) for branes and antibranes, respectively. ## Appendix E T-Duality of the Orientifold Models It is well known that $`T`$-duality changes Dirichlet and Neumann boundary conditions, therefore $`T`$-duality with respect to the three complex planes should map our models, where the standard model is embedded in D3-branes, to models where the standard model arises from D5, D7 or D9 branes. Since some standard-like models from D7 and D9 branes have recently been build , it is worth seeing the explicit correspondence between those models and our present constructions. In particular we will show that, in a suitable T-dual version, the realistic sectors in the models in correspond to local structures of D3- and D7-branes at $`\mathrm{IC}^3/ZZ_3`$ singular points, of the type studied in section 3. Hence, our results ‘explain’ the natural appearance of certain features (like hypercharge or $`BL`$ $`U(1)`$ gauge factors, or three generations) in the models in . Dual of the Left-Right Model of Let us start with the simplest model in , a left-right symmetric model constructed from a type IIB orientifold with D7<sub>3</sub>-branes at $`Y_3=0`$, and D3, $`\overline{\mathrm{D3}}`$-branes at some of the fixed points. Since the D7-branes sit at the orientifold plane $`O7`$, the twist matrices $`\gamma _{\theta ,7}`$ and $`\gamma _{\theta ,3}`$ have to satisfy the tadpole condition $$\mathrm{Tr}\gamma _7+3\left(\mathrm{Tr}\gamma _3\mathrm{Tr}\gamma _{\overline{3}}\right)=4$$ (E.1) Taking $`\gamma _7=(\stackrel{~}{\gamma }_7,\stackrel{~}{\gamma }_7^{})`$ with $$\stackrel{~}{\gamma }_7=\mathrm{diag}(\alpha \mathrm{I}_3,\alpha ^2\mathrm{I}_2,\mathrm{I}_2;\mathrm{I}_2,\alpha \mathrm{I}_7)$$ (E.2) and Wilson line in the first complex plane $`\stackrel{~}{\gamma }_W=\mathrm{diag}(\alpha \mathrm{I}_7,\mathrm{I}_9)`$, the gauge group is $`U(3)\times U(2)_L\times U(2)_R\times \left[SO(4)\times U(7)\right]`$. To satisfy the tadpole constraint, we just need to add D3-branes at the three points $`(1,m,0)`$ with $`\gamma _3=\mathrm{diag}(\alpha ,\alpha ^2)`$ (since $`\mathrm{Tr}\gamma _7\gamma _W^2=1`$). As shown in , this is a three-family model. Additional $`\overline{\mathrm{D3}}`$-branes must be introduced to cancel the untwisted tadpole, but we will not discuss them in detail here. To find the dual of this model in the approach of the present paper we have to dualize with respect to the first two complex planes so that the D7<sub>3</sub>-branes become D3-branes and viceversa. Also under $`T`$-duality Wilson lines become displacement of the branes, although the explicit mapping is rather involved. The best strategy in the present case is to construct the model explicitly. Since the Wilson lines above only affect the 14 entries leading to the LR model, these correspond to the D3-branes that sit outside the origin in the T-dual version. Therefore we consider 18 D3-branes at the origin, with twist matrix $`\gamma _{\theta ,3}=(\stackrel{~}{\gamma }_3,\stackrel{~}{\gamma }_3^{})`$ and $`\stackrel{~}{\gamma }_3=\mathrm{diag}(\mathrm{I}_2,\alpha \mathrm{I}_7)`$ (E.3) This should correspond to the ‘hidden’ part of the original model (last entries in (E.2), and indeed it leads to an $`SO(4)\times U(7)`$ gauge group. Since the origin is an $`O3`$-plane, the tadpole condition is (4.14). Since $`\mathrm{Tr}\gamma _3=3`$ we can cancel tadpoles introducing six D7-branes at $`Y_3=0`$ with $`\gamma _{\theta ,7}=\mathrm{diag}(\alpha \mathrm{I}_3,\alpha ^2\mathrm{I}_3)`$. After introducing these D7-branes, we have to worry about tadpole cancellation at the other fixed points in the $`Y_1,Y_2`$ planes, given by (4.15). At the points $`(\pm 1,0,0)`$ (mapped into each other under the orientifold action $`(1)^F\mathrm{\Omega }R_1R_2R_3`$) we can put seven D3-branes on each, with twist matrix $$\gamma _{\theta ,3}=\mathrm{diag}(\mathrm{I}_3,\alpha \mathrm{I}_2,\alpha ^2\mathrm{I}_2)$$ (E.4) giving rise to a three-generation $`SU(3)\times SU(2)_L\times SU(2)_R\times U(1)_{BL}`$ LR model of the type studied in section 3.4. Notice that we have introduced $`18`$ D3-branes at the origin and $`7+7=14`$ at the $`(\pm 1,0,0)`$ points, already saturating the total number of allowed D3-branes ($`32`$). Therefore, twisted tadpoles at the remaining fixed points should cancel without additional D3-branes. This can be achieved by introducing a Wilson line in the $`Y_2`$ direction (hence not affecting the points $`(m,0,0)`$) , such that $`\mathrm{Tr}\gamma _7\gamma _W=\mathrm{Tr}\gamma _7\gamma _W^2=0`$. The right choice is $$\gamma _W=\mathrm{diag}(1,\alpha ,\alpha ^2;1,\alpha ^2,\alpha )$$ (E.5) Therefore we do not have to introduce D3-branes at any of the six fixed points $`(m,\pm 1,0)`$. Notice also that this Wilson line breaks the D7-brane gauge group $`U(3)`$ to $`U(1)^3`$, as in the original model. The appearance of this Wilson line in the dual model reflects the fact that in the original model the six D3-branes were distributed among three different fixed points. The last consistency requirement is to have an equal number of D7- and $`\overline{\mathrm{D7}}`$-branes. To satisfy it, we put three anti-D7-branes with $`\gamma _{\overline{7}}=(1,\alpha ,\alpha ^2)`$ at the two mirror locations $`Y_3=\pm 1`$. Since this matrix is traceless there is no need to add more D3-branes to cancel tadpoles. These anti-branes can move freely in the $`Y_3`$ direction. Moreover, it is possible to turn on continuous Wilson lines along their world-volume, i.e. the directions $`Y_1`$, $`Y_2`$. This corresponds to possibility of moving the $`\overline{\mathrm{D3}}`$-branes in the original model in the three complex dimensions of the bulk. They are naturally dual to the anti D3-branes of the original model. A Variant Left-Right Orientifold Model and its Dual Let us consider a variant left-right orientifold model slightly different from the one in section 4. We satisfy the tadpole equations (4.14) at the origin by having 12 D7-branes and 4 D3-branes with: $$\gamma _7=(\alpha \mathrm{I}_6,\alpha ^2\mathrm{I}_6),\gamma _3=(\alpha ,\alpha ^2,\alpha ,\alpha ^2)$$ (E.6) In the second plane, we add a Wilson line: $$\gamma _W=(1,\alpha ,\alpha ^2,\mathrm{I}_3,1,\alpha ,\alpha ^2,\mathrm{I}_3)$$ (E.7) Since $`\mathrm{Tr}\gamma _7\gamma _W=\mathrm{Tr}\gamma _7\gamma _W^2=3`$ we can put a left-right model at each of the points $`(0,\pm 1,0)`$ in terms of 7 D3-branes with $$\gamma _3=(\mathrm{I}_3,\alpha \mathrm{I}_2,\alpha ^2\mathrm{I}_2)$$ (E.8) satisfying the tadpole condition (4.15). In order to minimize the number of D3-branes needed to cancel tadpoles (so we do not exceed 32) it is convenient to add a second Wilson line now in the third plane: $$\gamma _W^{}=(\mathrm{I}_3,1,\alpha ,\alpha ^2,\mathrm{I}_3,1,\alpha ,\alpha ^2)$$ (E.9) Since $`\mathrm{Tr}\gamma _7\gamma _W^{}=\mathrm{Tr}\gamma _7\gamma _W^2=3`$ we can cancel tadpoles at the points $`(0,0,\pm 1)`$ by putting just one single D3-brane on each with $`\gamma _3=1`$. The advantage of adding the second Wilson line is that at the points $`(0,\pm 1,\pm 1)`$ there is no need to add any D3-brane because $`\mathrm{Tr}\gamma _7\gamma _W^m\gamma _W^n=0`$ for $`m,n=1,2`$. Therefore we have at the plane $`Y_1=0`$ 12 D7-branes and a total of $`4+7+7+1+1=20`$ D3-branes. At each of the planes $`Y_1=\pm 1`$ we have to put 6 $`\overline{\mathrm{D7}}`$-branes and 6 D3-branes distributed on the fixed points $`(\pm 1,m,n)`$. We achieve tadpole cancellation by having: $$\gamma _{\overline{7}}=(\alpha \mathrm{I}_3,\alpha ^2\mathrm{I}_3),\gamma _W=(1,\alpha ,\alpha ^2,1,\alpha ,\alpha ^2)$$ (E.10) So we have $`\mathrm{Tr}\gamma _{\overline{7}}=3`$ and can cancel tadpoles at the points $`(\pm 1,0,m)`$ adding 2 D3-branes at each point with $`\gamma _3=(\alpha ,\alpha ^2)`$. This saturates the number of D3-branes ($`20+2\times 6=32`$) which is OK since at the points $`(\pm 1,\pm 1,m)`$ there is no need to add D3-branes since $`\mathrm{Tr}\gamma _{\overline{7}}\gamma _W=\mathrm{Tr}\gamma _{\overline{7}}\gamma _W^2=0`$. We can easily see that the gauge group coming from the D3-branes is $`[U(3)\times U(2)_L\times U(2)_R]\times U(2)\times U(1)^7`$ whereas from the D7-branes is simply $`U(1)^{12}`$. The full spectrum may be computed following the same lines as the left-right model of section 4. In order to find the $`T`$ dual of this model, in the formalism of we know that at the plane $`Y_1=0`$ we should have 20 D7 branes and 12 D3 branes. To obtain the same gauge group as the model above we choose: $$\stackrel{~}{\gamma }_7=(\alpha \mathrm{I}_3,\alpha ^2\mathrm{I}_2,\mathrm{I}_2,\alpha \mathrm{I}_2,\alpha ^2),\gamma _W=(\alpha \mathrm{I}_7,\mathrm{I}_2,\alpha ).$$ (E.11) Since $`\mathrm{Tr}\gamma _7=4`$ and $`\mathrm{Tr}\gamma _7\gamma _W=\mathrm{Tr}\gamma _7\gamma _W^2=1`$ we cancel tadpoles by having 2 D3-branes at each of the 6 points $`(0,\pm 1,m)`$ with $`\gamma _3=(\alpha ,\alpha ^2)`$. At each of the planes $`Y_1=\pm 1`$ we put 6 D7-branes with: $$\stackrel{~}{\gamma }_7=\alpha \mathrm{I}_3,\stackrel{~}{\gamma }_W=(1,\alpha ,\alpha ^2)$$ (E.12) Since $`\mathrm{Tr}\gamma _7=3`$ and $`\mathrm{Tr}\gamma _7\gamma _W=\mathrm{Tr}\gamma _7\gamma _W^2=0`$ we cancel tadpoles by adding 2 anti D3 branes at each of the points $`(\pm 1,0,n)`$. We can then see that there is an explicit D3/D7 duality among the two models indicating that they are $`T`$ dual of each other. In both cases anti-branes are trapped, as in the orientifold models of section 4.2 providing good examples of gravity mediated supersymmetry breaking. One of the positive points about this left-right model is that the extra gauge symmetries are essentially abelian, and most of the $`U(1)`$ symmetries get broken by the Green-Schwarz mechanism.
warning/0005/cond-mat0005274.html	ar5iv	text	# Tiles and colors ## 1 Lattices and tilings In many models for solid state physics the lattice as a simple idealization of the crystal is a given ingredient. Of course in reality the crystal is the result of the interactions and statistics of the particles in the system. The ideal lattice can be viewed as a periodic repetition of the unit cell. This makes it periodic and it is typically also symmetric under a specific discrete rotational symmetry group. The compatibility of the translational and rotational symmetry gives the well known crystallographic restrictions on the possible rotational symmetry groups. The quasicrystal is a state of matter which does not observe this restriction. This is possible only because it is not strictly periodic. Instead it has a property called quasiperiodicity. A quasiperiodic function in $`d`$ dimensions is a restriction of a periodic function in a dimension higher than $`d`$ to a $`d`$-dimensional hyperplane. The analogue for quasicrystals of the ideal lattice is the quasilattice or perfect quasiperiodic tiling. It is built as a repetition of more than one ‘unit cell’ also called tile. Examples of these structures are the famous Penrose tilings with rhombi as tiles, or their three-dimensional analogues with Ammann rhombohedra. In contrast to lattices, these quasilattices may have non-crystallographic rotational symmetries, reflecting the genuinely crystallographic symmetries of the periodic functions in higher dimensions. However, these non-crystallographic symmetries are not true symmetries of the function itself, but only of the absolute value of its Fourier transform. Because it is precisely this quantity that is measured in diffraction experiments, the symmetry is nevetheless quite real. ## 2 Random tilings as statistical models It was proposed by Elser that these noncrystallographic symmetries can be achieved not only by a quasiperiodic arrangement, but also by an ensemble of arbitrary packings of the tiles. This happens in the same way that the high temperature phase of a statistical model reflects the symmetry of the hamiltonian, not in any particular configuration, but only in the complete ensemble and in statistical averages. For instance, in the case of the Penrose rhombi, the angles between the edges are all multiples of $`\pi /5`$. Therefore if the plane is covered by copies of these rhombi, then also all the edges in the whole tiling have angles with one another which are multiples of $`\pi /5`$. As a consequence the continuous rotational symmetry of the plane, is reduced to at most a ten-fold discrete symmetry. The complete ensemble of tilings with these tiles will naturally have this tenfold rotational symmetry, unless the symmetry is further reduced by possible symmetry breaking schemes. This observation has led to the study of what is now known as random tilings. Despite their name these models have no intrinsic randomness. These are discrete statistical models of which the configurations are tilings of space by means of a limited set of tiles. They are called random only to emphasize the difference with perfect quasiperiodic tilings. In principle one may introduce an interaction energy between adjacent tiles, but in most cases studied there is no other interaction than the full packing constraint, i.e. that the tiles cover space without holes or overlaps. Though the phrase of random tilings came up in the study of quasicrystals, many random tilings actually live on the lattice . The dimer problem is a typical example, but also e.g. the hard hexagon model can be viewed as a tiling with hexagons and triangles. This paper, however, is concerned with random tilings of which the tiles do not fit together on a regular lattice. ## 3 Solvable random tilings On first sight the lack of an underlying lattice is a great difficulty in the analysis. However this difficulty is not essential as the configurations of the tiling can be mapped on a lattice by means of a geometric deformation. With this approach Widom studied the tilings of the plane with squares and equilateral triangles. Due to the fact that all angles of these objects are multiples of $`\pi /6`$ radians, this model exhibits twelvefold rotational symmetry. Widom showed that the transfer matrix of the tiling can be diagonalized by means of a Bethe Ansatz (BA). The corresponding Bethe-Ansatz equations were subsequently solved by Kalugin in the thermodynamic limit. This solution led to an exact expression of the extensive part of the entropy of the tiling. Later De Gier and the present author have found two other cases of quasicrystalline random tilings which can be solved by similar techniques. The tiles are rectangles and isosceles triangles, immediate generalizations of the squares and equilateral triangles. The rotational symmetry of the maximally symmetric phase is different, namely octagonal and decagonal respectively. Some example configuration of these three solved tilings is shown in Figure 1. The triangular tile has top angle $`\alpha =2\pi /n`$, where the integer $`n`$ takes the values 4, 5 and 6. The rectangle is simply defined to have sides matching in length with the legs (length 1) and the base (length $`2\mathrm{sin}\alpha /2`$) of the triangle. An obvious question at this point is in what way these tilings are different from other rectangle-triangle tilings, in which the angle $`\alpha `$ is another rational or even irrational fraction of $`2\pi `$. To see this it is necessary to consider in more detail the geometry of the tiling. Consider a vertex of the tiling, where $`i`$ triangles meet with their top angle $`\alpha `$, $`j`$ triangles with their base angle, and $`k`$ rectangles. Then obviously $$i\alpha +j\frac{\pi \alpha }{2}+k\frac{\pi }{2}=2\pi ,$$ (1) and the sum $`j+k`$ is even. Irrespective of the value of $`\alpha `$ three solutions of these conditions always exist, namely $`(i,j,k)=(0,0,4)`$, $`(1,2,2)`$ and $`(2,4,0)`$. Tilings with only vertices of these types we will call generic tilings. In all configurations of these tilings it is possible to vary continuously the angle $`\alpha `$, while the base angle and the length of the edges vary accordingly. Therefore the angle $`\alpha `$ plays no role in counting the number of ways the plane can be tiled. However, the restriction to this limited set of vertex types is so severe that these tilings have zero entropy: any finite region can be tiled in at most one way. For this reason these models are not of great interest from the view point of statistical physics. The tilings are possibly more interesting when $`\alpha `$ is chosen such that other combinations for $`(i,j,k)`$ are possible. Though there are many ways to allow for other solutions of (1), there are other conditions to cope with. The total number of base angles of an entire configuration is always twice the number of top angles. Therefore when there is a vertex in which $`2i<j`$ there must in the same configuration also be a vertex with $`2i>j`$. Therefore the only values of $`\alpha `$ giving rise to other than generic tiling configurations, are those which admit solutions of (1) both with $`2i<j`$ and $`2i>j`$. The task to find these values of $`\alpha `$ is elementary, but tedious. The result is that $`\alpha =2\pi /n`$ with $`n=3`$, 4, 5 or 6. All other values of $`\alpha `$ give only generic tilings. The cases $`n=4`$, 5 and 6 are mentioned above and are precisely those tilings that have been solved by means of the Bethe Ansatz. The case $`n=3`$ has not been discussed in the literature before. The tiles, a triangle with sides 1, 1 and $`\sqrt{3}`$, and a rectangle with sides 1 and $`\sqrt{3}`$ are such that the vertices of this tiling together with the mid-points of the rectangles form precisely the sites of a triangular lattice. It is tempting to believe that this tiling is also solvable by BA, but this does not appear to be the case. It may still be true that the rectangle-triangle tilings with $`n=4`$, 5 and 6 are members of an infinite sequence of solvable models. It may be that for higher values of $`n`$ the two tiles, the triangle and the rectangle do not suffice, and one may need to introduce more tiles as $`n`$ increases. On the same token it may be that for $`n=3`$ the two tiles are already too many, and one should work only with the triangle. This indeed gives a solvable tiling with finite entropy: when there are no rectangles pairs of triangles always share their long side, and thus form a rhombus. This rhombus tiling has been studied in many guises. ## 4 Integrability of the square-triangle tiling In almost all cases where the BA approach to the diagonalization of a transfer matrix or quantum Hamiltonian is effective, these operators are members of a commuting family The commutativity of this family is proven by the fact that the local Boltzmann weights satisfy the Yang-Baxter (YB) equation. Such a connection to a YB structure is not apparant for the solvable random tilings. However, in the case of the square-triangle tiling, a connection has been found. It turns out that the square-triangle model is equivalent to a limit of a known vertex model (associated with the affine Lie algebra A$`{}_{2}{}^{}{}_{}{}^{(1)}`$), provided with fields of the Perk-Schultz type. This vertex model does solve the YB equation. Without repeating the complete argument I will summarize the connection, and make some comments for later reference. For the derivation the reader is referred to. The Boltzmann weights denoted as $`W(\alpha ,\beta ;\gamma ,\delta )`$, with as successive arguments the states of the (left, bottom; top, right) legs of the vertex. $`W(\mu ,\mu ;\mu ,\mu )`$ $`=`$ $`X_\mu ^2\mathrm{sinh}(u+\lambda )`$ $`W(\mu ,\nu ;\mu ,\nu )`$ $`=`$ $`X_\mu X_\nu \mathrm{e}^{u\mathrm{sgn}(\nu \mu )}\mathrm{sinh}\lambda `$ (2) $`W(\mu ,\nu ;\nu ,\mu )`$ $`=`$ $`X_\mu ^2(Y_{6\mu \nu })^{2\mathrm{sgn}(\mu \nu )}\mathrm{sinh}u`$ where $`\mu `$ and $`\nu \mu `$ take the values 1, 2 and 3. The limit involves both the spectral parameter $`u`$ and the field parameters $`X_\mu `$ and $`Y_\mu `$. The square-triangle model is recovered when the limit is taken in two steps. In the first step the spectral parameter is taken to $`u=\lambda `$, at which some of the vertex weights vanish. This is a point of special symmetry where the vertices of the square lattice can be factorized in vertices of the hexagonal lattice. In fact both the spatial symmetry of the hexagonal lattice and the internal permutation symmetry of the three vertex states are fully observed. At this point the model is equivalent to the three-coloring problem of the edges of the honeycomb lattice such that equally colored edges never meet in a vertex. This problem was solved by Baxter. In the second step a combination of the Perk-Schultz fields is taken to an extreme limit. $`X_1^1=X_2=X_3=Y_1^1=Y_2^1=Y_3^1=x`$ and take $`x0`$. In this limit a few of the remaining Boltzmann weights vanish as well. Because in the first limit the spectral parameter as a free variable is lost, so is also the YB structure. In order to retain the spectral parameter one might wish to take the field limit first. However in that procedure one of the vertex weights would diverge, or after suitable normalization, be the only one to survive. Thus we see the two limits do not commute. As a consequence we come to the counter-intuitive conclusion that the Boltzmann weights of the square-triangle tiling are a limit of a solution to the YB equation, but these weights themselves are not a proper solution. It thus appears that the solvability of square-triangle model is related to the YB equation, albeit somewhat remotely. In the remainder of this paper we investigate in more detail the solvable rectangle-triangle tiling with octagonal symmetry, and in particular discuss the question if it also can be viewed as a limit of a YB-solvable model. For this purpose we review the solution, though in an approach different from that published before. Here we try to keep the symmetry of the tiling as much as possible. ## 5 The octagonal rectangle-triangle tiling This model was first considered by Cockayne in a slightly different language. The plane (or a finite periodic section of it) is tiled with rectangular, isosceles triangles and rectangles of which the sides, length 1 and $`\sqrt{2}`$, match those of the triangle. One may, as in a canonical ensemble fix the relative density of each of the tiles, or as in the grand canonical ensemble control the relative density by means of an activity variable. The tiles are not permitted to overlap, and may leave no space uncovered. Since all the angles in the tiles are multiples of $`\pi /4`$, the angle between any two edges in a tiling configuration must also be such multiple. Therefore the rectangles may occur in four, and the triangles in eight different orientations, shown in figure 2. As a first step in the analysis the tiles are deformed, such that they are commensurable with a regular square lattice but continue to cover the plane without holes or overlap. The deformation is illustrated in figure 2 and may be described as follows. Let all the edges in the original tiling make an angle with the horizontal equal to an odd multiple of $`\pi /8`$ radians. The short edges are rotated over $`\pi /8`$ to the left or to the right, such that they end up horizontally or vertically. The long edges are also rotated over $`\pm \pi /8`$ but such that they end up in one of the diagonal directions. As a result the triangles are rotated rigidly and the rectangles are deformed into parallelograms. The deformed tiles fit precisely on the square lattice, and any tiling configuration can be viewed as the state of a lattice model. The tile edges in the resulting lattice are marked by either solid or dotted lines to distinguish the original orientations. This is necessary because two differently oriented edges are mapped onto the same orientation. The fact that the mapping, while it deforms the entire configuration, does not create any holes or overlaps, follows directly from the observation that it is defined as a mapping of the edges, i.e. the shared boundary between adjacent tiles. The configurations of the original tiling are mapped one-to-one on configuratings of the lattice tiling. Therefore the combinatorial problem of counting the number of possible tilings seems unaltered. However, because the overall shape and the area of a tiling is altered by the mapping, the problem of tiling a given section of the plane with a given number of tiles is not the same. In the thermodynamic limit this distinction will be insignificant. The change of area can be accounted for, and the shape does not matter for the thermodynamic functions. The lattice representation of the tiling can be compactly described by listing the possible states of each elementary square of the lattice, as shown in figure 3. The sides of the elementary squares can be in one of four states, and represent either a short edge of the tiling, or a diagonal of a rectangle. The local configurations of the elementary faces of the lattice can be represented as fourteen different states of a vertex model which has four possible states for each edge. Figure 3 shows a representation in terms of two types of arrows, which will be called open and filled. This notation makes immediately evident two conservation laws, the net flux of open and that of filled arrows. We will refer to this lattice model as fourteen-vertex model. ## 6 Bethe Ansatz equations A coordinate BA formalism can be set to diagonalize the vertex model shown in figure 3. For the reference state we choose the open arrow up. The quasi-excitations are the locations with a vertical filled arrow, up or down. The open arrow down serves as a bound state between these two different excitations. Besides the conservation of open and filled flux, there is an additional conservation law, which is less evident. It turns out that at every application of the transfer matrix the excitations move one step either to the left or to the right. Therefore the lattice can be divided into two sublattices such that the number of excitations on each sublattice remains constant from row to row. As a result there are in total three conserved number and correspondingly three families of BA variables. In order to find the most general solution the vertex weights of the model are kept general. As the BA eigenstate takes form, however, several restrictions on the weights are necessary. In the first place it turns out that the weights of the vertices almost factorize into weights associated with the tiles. Complete factorization would imply completely non-interacting tiles (apart from the ban on holes and overlaps). The one exception is that where two triangles form a square by sharing their long edge, this configurations has an extra weight $`1/2`$. The tiles may have different weights in the different orientation, but there are restrictions on these weights. We denote the weights of the oriented triangles as $`t_j`$, with $`j=1,\mathrm{},8`$ and of the rectangles as $`r_j`$, with $`j=1,\mathrm{},4`$, both in the order in which they are shown in figure 2. Then the restrictions are: $`t_1t_5`$ $`=`$ $`t_3t_7`$ $`t_2t_6`$ $`=`$ $`t_4t_8`$ (3) $`r_1r_3`$ $`=`$ $`\pm r_2r_4`$ Because we favor a statistical interpretation of the tiling we choose the positive sign in the last line. The resulting BA equations are, for a lattice with even horizontal size $`L`$, and periodic boundary conditions. $`x_p^L`$ $`=`$ $`()^{n_x+1}A{\displaystyle \underset{k=1}{\overset{n_\lambda }{}}}\left(\lambda _kx_p^2\right)`$ $`y_q^L`$ $`=`$ $`()^{n_y+1}A{\displaystyle \underset{k=1}{\overset{n_\lambda }{}}}\left(\lambda _ky_q^2\right)`$ (4) $`1`$ $`=`$ $`()^{n_\lambda +1}B{\displaystyle \underset{p=1}{\overset{n_x}{}}}\left(\lambda _kx_p^2\right){\displaystyle \underset{q=1}{\overset{n_y}{}}}\left(\lambda _ky_q^2\right)`$ In terms of the solutions $`x_p`$, $`y_q`$ and $`\lambda _k`$ of these equations the eigenvalue of the transfer matrix reads $$\mathrm{\Lambda }=C\underset{k=1}{\overset{n_\lambda }{}}\lambda _k\underset{p=1}{\overset{n_x}{}}x_p^1\underset{q=1}{\overset{n_y}{}}y_q^1$$ (5) The constants are expressed in the tile weights: $`A`$ $`=`$ $`\left({\displaystyle \frac{t_1t_5r_3}{t_2t_6r_4}}\right)^{n_\lambda }`$ $`B`$ $`=`$ $`\left({\displaystyle \frac{t_1t_5}{r_4}}\right)^L\left({\displaystyle \frac{t_1t_5r_3}{t_2t_6r_4}}\right)^{n_x+n_y}`$ (6) $`C`$ $`=`$ $`(t_1t_5)^L\left({\displaystyle \frac{r_1}{t_1t_5}}\right)^{n_x+n_yn_\lambda }\left({\displaystyle \frac{r_2}{r_1}}\right)^{n_\lambda }.`$ The form of these BA equations is typical for tiling models. The eigenvalue expression is a simple product rather than a sum of, say four, products, according to the number of edge states. The factors in the BA equations are binomial rather than rational or trigonometric. ## 7 Search for a Yang-Baxter structure Even though the tiling model has a number of free parameters in the solution, none of those plays the role of a spectral parameter. The only way that the parameters feature in the eigenvalue is via the combination $`C`$ as an overall factor. As the BA equations are the most general given the set of vertex configurations with non-zero weight, the only way to introduce a spectral parameter is to include other vertex configurations. The natural choice is to include only configurations that satisfy the flux conservation of both types of arrows. Our attempts in this direction have all failed, as they led only to null solutions. Incapable of an exhaustive search, we can not make strong statements about this possibility. However, the situation may be similar to the case of square-triangle model, described above. If that is so it would be futile to seek solutions of the YB equations which include the fourteen-vertex model. Instead this model would be a singular limit of such a solution, and not a true member. It is on this assumption that we proceed. Guided by the results concerning the integrability of the square-triangle model, we seek a vertex model of which the vertex configurations include those of the fourteen-vertex model, but have a higher symmetry. This would still have a fixed value of the spectral parameter, and permit the introduction of field parameters, such that in a special limit the fourteen-vertex model is recovered. A natural choice of a symmetrized version of the fourteen-vertex model is extending the permitted vertex states with those obtained by arrow inversion from the original fourteen. The resulting set of vertices are twenty-four states in which one open and one filled arrow enter and leave the vertex. We propose that it is possible to assign weights to these twenty-four vertex states such that (i) they form a member of a solution of the YB equation, and (ii) permit a limit in which they reduce to the fourteen vertices of the octagonal tiling problem. To be more precise we propose that there exists a solution of the YB equation $$\underset{\mu ^{\prime \prime },\nu ^{\prime \prime },\gamma }{}W^{\prime \prime }(\mu ,\nu ;\nu ^{\prime \prime },\mu ^{\prime \prime })W(\mu ^{\prime \prime },\beta ;\gamma ,\mu ^{})W^{}(\nu ^{\prime \prime },\gamma ;\alpha ,\nu ^{})=$$ $$\underset{\mu ^{\prime \prime },\nu ^{\prime \prime },\gamma }{}W^{\prime \prime }(\mu ^{\prime \prime },\nu ^{\prime \prime };\nu ^{},\mu ^{})W^{}(\nu ,\beta ;\gamma ,\nu ^{\prime \prime })W(\mu ,\gamma ;\alpha ,\mu ^{\prime \prime }),$$ (7) such that the Greek symbols take the four arrow states as values. Part (i) of the proposal states that the symbol $`W`$ is non-zero only for the twenty-four cases of its arguments when one open and one filled arrow enter and leave the corresponding vertex. Guided by the YB structure of the square-triangle model, we expect that the symbols $`W^{}`$ and $`W^{\prime \prime }`$ may have nonzero entries besides these ones. Part (ii) of the proposal states that the solution manifold permits a limit to be taken in which the symbol $`W`$ is reduced to weights of the fourteen vertices corresponding to the octagonal tiling model (figure 3). In this limit some elements of the symbols $`W^{}`$ or $`W^{\prime \prime }`$ may well diverge. Before discussing tests of this proposal we make a few remarks concerning the symmetry of the model. Consider a configuration of these vertices on a square lattice. The edge states are denoted by the kind of arrow, open or filled, and by the direction, up or down for vertical edges and left or right for horizontal ones. If the sites of the lattice are divided into two square sublattices A and B, the direction of the arrows can be denoted as A to B or vice versa. With this labelling, the four edges incident in the same vertex are always in different states. In other words the configurations of the 24-vertex model are in one-to-one correspondence to the colorings of the edges of the square lattice with four colors, such that in no vertex two edges of the same color meet. Clearly in this coloring problem the four colors can be freely permuted without altering the ensemble of coloring configurations. Since any of the twenty-four vertex configurations can be turned into any other one by a suitable permutation of the four edge states, the only weight assignment invariant for these permutations gives all vertices equal weight. It is likely that if there is an integrable manifold in this 24-vertex model, it includes the symmetric point where all weights are equal, equivalent to the four-coloring problem of the edges of the square lattice. In order to test the above proposal we set up and attempt to solve the YB equations (7). The most restrictive and simple approach is to insert for all three symbols, $`W`$, $`W^{}`$ and $`W^{\prime \prime }`$ just the twenty-four weights described above, setting all other weights equal to zero. This turns out to lead only to a trivial null solution. Thus, the full solution, if it exists reuires more than the twenty-four vertices. A natural extension is with those vertices in which the total flux of each types of arrows is still conserved. Those vertex configuration are the well-known six-vertex configuration, now in two types, those with only open arrows and those with only filled arrows, twelve configurations in total. Together with the original twenty-four vertices this makes thirty-six. We have not succeeded in solving the YB equations for these thirty-six weights for each of the symbols $`W`$, $`W^{}`$ and $`W^{\prime \prime }`$, due to the complexity of the problem. However, by allowing for $`W`$ only the original twenty-four weights while retaining the full thirty-six weights for $`W^{}`$ and $`W^{\prime \prime }`$, the the problem is considerably simpler and could in fact be solved. Unfortunately, the resulting weights for $`W`$ can not be made all positive. Therefore, contrary to our expectation the four-coloring problem is not included in the solution. Even so, part (i) of the proposal is verified, albeit with some negative weights. We note that the twenty-four weights can be made equal in absolute value, even though they can not all be made positive. The solution of this reduced set of YB equation has still a great number of free parameters. This freedom permits the reduction of the weights $`W`$ to only the fourteen non-zero weights of the octagonal tiling model. Interestingly, in that limit the same restriction on the vertex weights is found as those following from the BA, with the sign in the third line of (3) no longer free, but negative. In other words, one, or three of the weights of the oriented rectangles must be negative. Thus we have not found the YB structure behind the integrability of the octagonal tiling problem, but something intriguingly close to it. We can not recover the tiling problem from a YB structure, but then the tile weights can not be all positive. Thus part (ii) of the proposal is almost verified: the solution to the YB equation has a limit which up to signs reproduces the weights of the tiling model. ## 8 The four-coloring problem In the previous section we encountered the combinatorial problem of coloring the edges of the square lattice in four colors, such that two equally colored edges never meet in the same vertex. This four-coloring problem has been studied also in its own right. It is related to fully packed loop models on the square lattice which has application to physics of polymers in the melt. It is also related to the Hamiltonian walk problem . As described above, we have not succeeded in finding a solution to the YB equation which includes this model, at least not the coloring problem in which all configurations have positive weight. This does not imply that the model is not integrable. To yet have an indication of its integrability we have attempted to construct BA eigenectors to the transfer matrix. We used the formulation of the 24-vertex model. Thus the edges of the square lattice wrapped on a cylinder of even circumference $`L`$ all carry an open or filled arrow, pointing along the edge. The configurations in which one arrow of each type points into each vertex have weight one, all other configurations have weight zero. The coordinate (nested) BA approach indeed yields eigenvectors to the transfer matrix, rather similar in structure to those of the octagonal tiling model (4). With respect to a reference state, all open arrow up, the excitations are the filled arrows. The open arrow down serves again as a bound state of two opposite filled arrows. With every application of the transfer matrix the excitations that do not form a bound state, move a one step, to the right or to the left. As a consequence not only the number of excitations is conserved, but also their distribution over two sublattices. Precisely as in the fourteen-vertex model, two families of variables give the momenta of the excitations on the even and odd sublattice respectively, and a third set of variables is associated with the distribution of the filled down arrows among all of the filled arrows. The resulting BA equations read in suitable variables $$\left(\frac{1+u_j}{1u_j}\right)^{L/2}=\left(\underset{m=1}{\overset{n_w}{}}\frac{w_mu_j+1}{u_jw_m+1}\right)\left(\underset{k=1}{\overset{n_u}{}}\frac{u_ju_k+2}{u_ju_k2}\right)$$ $$\left(\frac{1+v_j}{1v_j}\right)^{L/2}=\left(\underset{m=1}{\overset{n_w}{}}\frac{w_mv_j+1}{v_jw_m+1}\right)\left(\underset{k=1}{\overset{n_v}{}}\frac{v_jv_k+2}{v_jv_k2}\right)$$ (8) $$1=\left(\underset{j=1}{\overset{n_u}{}}\frac{w_mu_j+1}{u_jw_m+1}\right)\left(\underset{k=1}{\overset{n_v}{}}\frac{w_mv_k+1}{v_kw_m+1}\right)\left(\underset{l=1}{\overset{n_w}{}}\frac{w_lw_m+2}{w_lw_m2}\right)$$ The eigenvalue can be written in terms of the solutions of these equations $`\mathrm{\Lambda }`$ $`=`$ $`\left({\displaystyle \underset{j=1}{\overset{n_u}{}}}{\displaystyle \frac{1+u_j}{1u_j}}\right)^{1/2}\left({\displaystyle \underset{k=1}{\overset{n_v}{}}}{\displaystyle \frac{1+v_k}{1v_k}}\right)^{1/2}\left({\displaystyle \underset{m=1}{\overset{n_w}{}}}{\displaystyle \frac{2w_m}{w_m}}\right)`$ (9) $`+`$ $`\left({\displaystyle \underset{j=1}{\overset{n_u}{}}}{\displaystyle \frac{1u_j}{1+u_j}}\right)^{1/2}\left({\displaystyle \underset{k=1}{\overset{n_v}{}}}{\displaystyle \frac{1v_k}{1+v_k}}\right)^{1/2}\left({\displaystyle \underset{m=1}{\overset{n_w}{}}}{\displaystyle \frac{2w_m}{w_m}}\right)`$ These equations have been derived for the sectors where the excitations are relatively sparse. We have no proofs for these equations in full generality. We have numerically verified them for arbitrary sectors of lattices up to twelve sites in the horizontal direction. When the roots of the equation include a $`w_m=0`$ the eigenvalue is undetermined from these expressions. This ambiguity could be resolved if we could find the eigenvalue expression with a spectral parameter. Altogether the BA approach shows that indeed this four-coloring problem is integrable. These BA equations look very much like those derived by Martins for mixed SU($`N`$) vertex models for the case $`N=4`$. However, the sign of some of the factors is different (those involving both $`u`$ and $`w`$ or both $`v`$ and $`w`$). The eigenvalue (9), aside from the absence of a spectral parameter here, differs from also from that of Martins in an overall factor, which depends on the BA roots. In fact the comparison inspires to introduce a spectral parameter $`u`$ in our expression for the eigenvalue. This process may seem arbitrary, but it is highly constrained for the following consideration. It is known that the BA equations follow directly from the expression for the eigenvalue by the requirement that the eigenvalue be an entire function of the spectral parameter. Here we follow the reverse argument, we know the BA equations and the eigenvalue for one value of the spectral parameter, and propose a specific dependence on the spectral parameter such that the poles from the different terms in the expression cancel against each other as a consequence of the BA equations. $$\mathrm{\Lambda }(u)=\left(\underset{j=1}{\overset{n_u}{}}\frac{1+u_j}{1u_j}\right)^{1/2}\left(\underset{j=1}{\overset{n_v}{}}\frac{1v_j}{1+v_j}\right)^{1/2}\times $$ $$[(1u^2)^{L/2}\left(\underset{j=1}{\overset{n_v}{}}\frac{1+v_j2u}{1v_j+2u}\right)\left(\underset{j=1}{\overset{n_w}{}}\frac{2w_j+2u}{w_j2u}\right)$$ $$+(1u^2)^{L/2}\left(\underset{j=1}{\overset{n_u}{}}\frac{1u_j+2u}{1+u_j2u}\right)\left(\underset{j=1}{\overset{n_w}{}}\frac{2w_j+2u}{w_j2u}\right)$$ $$+(u+u^2)^{L/2}\left(\underset{j=1}{\overset{n_u}{}}\frac{3u_j+2u}{1+u_j2u}\right)$$ $$+(u^2u)^{L/2}\left(\underset{j=1}{\overset{n_v}{}}\frac{3+v_j2u}{1v_j+2u}\right)]$$ (10) The value $`u=0`$ corresponds to the original expression (9). One check of this proposal is that it now allows us to calculate the eigenvalue unambiguously also in those cases where one of the $`w_m`$ roots vanishes. Indeed it turns out that these cases the limit $`u0`$ of the expression correspond to eigenvalues of the transfer matrix that were undetermied before. The eigenvalue expression (10) together with the BA eigenvectors of the transfer matrix of the 24-vertex model unambiguously defines a matrix. If this matrix can be built up from local Boltzmann weights is yet to be verified. The whole expression has some unusual features, such as the obvious overall factor, independent of the spectral parameter, and involving a square root. ## 9 Summary Quasicrystals show rotational symmetries which are forbidden by the rules of crystallography. This is possible because quasicrystals are aperiodic. From the fact that quasicrystals exist we must conclude their structure is a minimum of the free energy. It can be argued that in quasicrystals, more than in crystals, the entropy is, compared with the energy, plays a significant role in determining this minimum. For these materials random tilings serve as model systems. This paper reviews some results of two-dimensional random tilings of which the entropy has been calculated exactly. These models have phases with octagonal, decagonal and dodecagonal spatial rotational symmetry. They appear as three member of an infinite sequence, but there properties, among which solvability, uniquely set them apart from the rest of the sequence. We have discussed the question why this is the case. While these models have been solved by means of the Bethe Ansatz, an underlying Yang-Baxter structure is not apparent. Only in the dodecagonal case the relation with a solution to the YB equation has been known for a few years. Here the connection involves a model family which includes the coloring problem of the edges of the hexagonal lattice with three colors, such that the edges meeting in a vertex have different colors. The tiling model is a singular limit of this family. This paper takes a few steps at finding a similar connection for the solved tiling model with octagonal symmetry. Its transfer matrix also permits a BA solution. The search for the YB structure behind this integrability is not completed. However, the trace seems to lead, just as in the dodecagonal case, by a coloring problem. In this case it is the problem of coloring the edges of the square lattice in four colors such that nowhere two edges with the same color meet in a vertex. We give the BA equations of this problem, thus showing that the model is integrable. ## Acknowledgments The author wishes to thank the organizers of “Baxter’s revolution in mathematical physics”, and R. Tateo, S.O. Warnaar, J.C. de Gier and D. dei Cont for many valuable discussions.
warning/0005/cond-mat0005478.html	ar5iv	text	# The High Magnetic Field Phase Diagram of a Quasi-One Dimensional Metal ## Abstract We present a unique high magnetic field phase of the quasi-one dimensional organic conductor (TMTSF)<sub>2</sub>ClO<sub>4</sub>. This phase, termed ”Q-ClO<sub>4</sub>”, is obtained by rapid thermal quenching to avoid ordering of the ClO<sub>4</sub> anion. The magnetic field dependent phase of Q-ClO<sub>4</sub> is distinctly different from that in the extensively studied annealed material. Q-ClO<sub>4</sub> exhibits a spin density wave (SDW) transition at $``$ 5 K which is strongly magnetic field dependent. This dependence is well described by the theoretical treatment of Bjelis and Maki. We show that Q-ClO<sub>4</sub> provides a new B-T phase diagram in the hierarchy of low-dimensional organic metals (one-dimensional towards two-dimensional), and describe the temperature dependence of the of the quantum oscillations observed in the SDW phase. Quasi-one dimensional organic metals have the general character of a large bandwidth along the molecular stacking (chain) direction, followed by significantly smaller band widths in the inter-chain and inter-plane directions. For the Bechgaard salts these ratios are ($`t_a`$, $`t_b`$, $`t_c`$: 1000, 250, 3 meV) respectively. For sufficiently small transverse bandwidths, a one-dimensional conductor will undergo an instability, at a critical temperature, to an insulating ground state. Following Yamaji, this temperature has an anisotropic bandwidth dependence in terms of the so called ”imperfect nesting parameter”, $`\epsilon _0`$ = $`t_b^2/t_a`$. Hence the more two-dimensional the material is, the larger $`\epsilon _0`$ will be, and the lower the temperature where the instability, or ”nesting” will occur. In the case of the Bechgaard salts, a spin density wave (SDW) ground state is formed. For sufficiently large $`\epsilon _0`$, the low temperature ground state remains metallic, but high magnetic fields can effectively reduce $`\epsilon _0`$ (i.e. drive the system more one-dimensional), and a field induced spin density wave state (FISDW) can be stabilized. This latter phenomena has been the subject of extensive experimental and theoretical work. In this communication we consider a system where $`\epsilon _0`$ is large, but where an SDW state still forms at a temperature T<sub>SDW</sub> = 5 K. This allows a unique situation where T<sub>SDW</sub> increases by a factor of two in high magnetic fields due to the close competition between $`\epsilon _0`$ and the magnetic energy. To accomplish this we employ the well-studied organic conductor (TMTSF)<sub>2</sub>ClO<sub>4</sub>. What is particular to our approach is that the material has been prepared in a very rapidly (30K/sec) thermally quenched state (i.e. Q-ClO<sub>4</sub>) to preserve the high temperature electronic structure, which is comprised of two, open orbit, warped Fermi surface sheets. Otherwise, if the material is slowly cooled, the tetrahedral anion ClO<sub>4</sub> undergoes an ordering transition around 24 K, the unit cell doubles in the b direction, and the resulting Fermi surface becomes more complex. Generally, it is this relaxed state of the material (hereafter R-ClO<sub>4</sub>), which has been most extensively studied. In contrast, Q-ClO<sub>4</sub> is in the class of Bechgaard salts (TMTSF)<sub>2</sub>X where X = PF<sub>6</sub>, AsF<sub>6</sub>, NO<sub>3</sub> which form an SDW state at ambient pressure below a transition temperature $`T_{SDW}`$ (12 K for X = PF<sub>6</sub> and AsF<sub>6</sub>, 10K for X = NO<sub>3</sub>, and 5 K for X = Q-ClO<sub>4</sub>). In these cases a high magnetic field improves the nesting condition, and $`T_{SDW}`$ increases with magnetic field. We have determined that not only does this happen for the case of Q-ClO<sub>4</sub> , but the effect is the most dramatic, and leads to a new phase diagram in the hierarchy of nesting parameters ( 7 K $``$ $`\epsilon _0`$ $``$ 22.5 K ) in the Bechgaard salts. The measurements reported here were carried out on two samples in 50 T (sample $`\mathrm{\#}1`$) and 60 T (sample $`\mathrm{\#}2`$) pulsed-field magnets at the Los Alamos National Laboratory. Electrical transport contacts were made via graphite paint and 12$`\mu `$$`m`$ gold wires, with a dc four-terminal technique with a current of 50 $`\mu `$$`A`$. The current was applied transverse to the layers along the c-axis, as were the voltage contacts. The magnetic field is also along the c-axis. To ensure that the samples fully quenched, the samples were put in direct contact with liquid helium from room temperature as rapidly as possible. Estimated cooling rates were of order 30 K/sec or greater. The use of graphite paint appeared to greatly enhance the reliability of the contacts and to reduce degradation (cracking) of the samples during the rapid cool-downs. Systematic temperature measurements for each run were performed with a single quench to preserve the anion disorder in the samples. In Fig. 1a we show a summary of the magnetoresistance measurements for sample $`\mathrm{\#}1`$ above 5 K, and in Fig. 1b similar measurements at lower temperatures for both samples $`\mathrm{\#}1`$ and $`\mathrm{\#}2`$ are presented. In the inset of Fig. 1a the temperature dependence of the c-axis resistivity is shown at zero magnetic field. Here the upturn is the onset of the spin density wave transition, as established by, for instance NMR studies . In Fig. 2 we show an expanded view of the behavior of the magnetic field dependent resistance, where the data has been plotted vs. the square of the magnetic field. All data shown is for temperatures higher than $`T_{SDW}`$ (B=0). Here the magnetoresistance remains small, and quadratic in field, until a critical field (hereafter $`B_{SDW}`$) is reached. We have fit the 10.6 K resistance data in Fig. 2 to the quadratic function (expected for the case of an open orbit metal) below the point where the slope changes. Aside from a weak temperature dependence, all of the data follow this general functional dependence until the point $`B_{SDW}`$ is reached. $`B_{SDW}`$ is temperature dependent, and is manifested as a change in the field dependence of the magnetoresistance. This defines the threshold field between the metallic and SDW ground states. Above $`B_{SDW}`$, additional structure appears (hereafter $`B_{SDW}^{}`$) where the field dependence of the magnetoresistance changes again. At higher fields quantum oscillations of frequency F=190 T become evident. We will return to these last two points in the discussion below. In Fig. 3 we summarize the dependence of $`B_{SDW}`$ and $`B`$ in terms of the corresponding temperatures $`T_{SDW}`$ and $`T`$ based on the analysis of Fig. 2 and the zero field value from the inset of Fig. 1a. The new phase diagram, as defined by Fig. 3, is the main result of the present work. To put the new data for Q-ClO<sub>4</sub> in perspective, we have included previous results for AsF<sub>6</sub>, PF<sub>6</sub>, and NO<sub>3</sub>, (Refs. ) for the field dependence of $`T_{SDW}`$ in terms of the theoretical framework, which describes the magnetic field dependence of $`T_{SDW}`$ as given by Bjelis and Maki in the form $$\mathrm{ln}\left[\frac{T_{SDW}}{T_{SDW0}}\right]\underset{l=\mathrm{}}{\overset{\mathrm{}}{}}J_l^2(e_0)\left\{Re\mathrm{\Psi }(\frac{1}{2}+2ilx_1)\mathrm{\Psi }(\frac{1}{2})\right\}.$$ (1) Here $`e_0`$=$`\epsilon _0`$/$`\omega _b`$, where $`\epsilon _0`$ is the imperfect nesting parameter described above, and $`\omega _b`$=$`ev_FbB`$. The latter is the effective cyclotron frequency along the b-axis, where $`v_F`$ is the Fermi velocity. $`J_l`$ and $`\mathrm{\Psi }`$ are Bessel and digamma functions respectively, and $`x`$<sub>1</sub> = $`\omega `$<sub>b</sub>/4$`\pi `$$`T_{SDW}`$. $`T_{SDW0}`$ corresponds to the transition temperature for perfect nesting. This expression, in an asymptotic form, successfully describes the field dependence of $`T_{SDW}`$ for (TMTSF)<sub>2</sub>NO<sub>3</sub> Ref. and (TMTSF)<sub>2</sub>PF<sub>6</sub> Ref.. For Q-ClO<sub>4</sub>, Eq. 1 is also applied, but its full (non-asymptotic) form must be used to obtain proper convergence at low magnetic fields due to the large imperfect nesting parameter $`\epsilon _0`$ needed to describe the data. (The fitting parameters for all three materials are listed in the caption of Fig. 3.) It is clear that $`T_{SDW}`$(B =0) decreases with increasing $`\epsilon _0`$. However, since the general Fermi surface topology is very similar, all three materials approach a common transition temperature in the high field limit. For completeness, the phase diagram of R-ClO<sub>4</sub> is also shown in Fig. 3. A point that must be clearly made is that the Q-ClO<sub>4</sub> phase diagram is very different from that studied in the R-ClO<sub>4</sub> case. For R-ClO<sub>4</sub> the Fermi surface topology involves double open orbit sheets due to the anion ordering, the imperfect nesting is so large that at zero field the system is metallic, even superconducting (T<sub>c</sub> = 1.2 K), and the maximum field induced $`T_{SDW}`$ is less than 6 K. Another factor which distinguishes Q-ClO<sub>4</sub> from R-ClO<sub>4</sub> is that the quantum oscillation frequency in the magnetoresistance is 190 T, as confirmed by this (see below) and previous independent studies. In contrast, for R-ClO<sub>4</sub> the frequency is 250 T. This confirms that the Fermi surface topologies are fundamentally different. Returning for a moment to the observation of the $`B_{SDW}^{}`$ feature (Figs. 2 and 3), it is not clear what assignment to make to it (i.e. a SDW sub-phase for instance). It falls outside the prediction of Eq. 1 since there only one transition ($`B_{SDW}`$) is predicted. Hence further experimental and theoretical work will be needed to fully understand the significance of $`B_{SDW}^{}`$. We next turn to the nature of the oscillations in the magnetoresistance which appear above $`B_{SDW}`$ (see Fig. 1). The oscillatory component of the magnetoresistance, plotted as $`\mathrm{\Delta }R/R`$ vs. inverse magnetic field is shown in Fig. 4a and 4b for representative temperatures. (Here R is the non-oscillatory background magnetoresistance.) The oscillation amplitude increases until about 5 K, but then rapidly vanishes at lower temperatures. In Fig. 4c we show the amplitude of the oscillations (normalized by R ) as a function of temperature. Such oscillations, commonly called ”rapid oscillations” or ”RO” appear in a number of the quasi-one dimensional organic salts, including the Bechgaard salts and the (DMET-TSeF)<sub>2</sub>AuCl<sub>2</sub> salts. These oscillations are periodic in inverse field, with a frequency of between 190 and 250 T. If considered as closed orbits, they represent about 3% of the first Brillouin zone. A general observation is that they only occur when the original open orbit Fermi surface has undergone a reconstruction. This can happen by an anion ordering transition and/or nesting of the Fermi surface. In many cases where a spin density wave ground state is stabilized, the amplitude of the rapid oscillations first increases with decreasing temperature, but below a characteristic temperature $`T^{}`$ (typically between 2 to 4K), the amplitude attenuates exponentially for $`T0`$. These oscillations have been described as a magnetic breakdown phenomena in an imperfectly nested Fermi surface topology. However, below $`T^{}`$ there is an improvement of the nesting condition, such that more of the reconstructed Fermi surface becomes gapped, and the magnetic breakdown becomes less probable. As shown in Fig. 4c, and following Ref., this behavior may be modeled by the standard Lifshitz - Kosevich description for quantum oscillations above $`T^{}`$, and an exponential attenuation below $`T^{}`$. Here the parameters of the model involve an effective mass m\* = 1.2 m<sub>0</sub>, a Dingle temperature $`T_D`$ = 7 K, a $`T^{}`$ = 3.75 K, and a low temperature magnetic breakdown gap of 10 K. The Dingle temperature associated with the quantum oscillations discussed above gives insight into the nature of carrier scattering, and therefore disorder, in Q-ClO<sub>4</sub>. Typically for ”clean” organic metals where strong quantum oscillations are observed, $`T_D`$ is of order 1 K. For Q-ClO<sub>4</sub> it is considerably higher ( 7 K). Since anion disorder is necessary to stabilize the Q-ClO<sub>4</sub> state, this could be a possible source for the enhanced scattering. However, similar studies of the quantum oscillations in AsF<sub>6</sub> and PF<sub>6</sub> anion systems, where there is no disorder involved, show a similar value of $`T_D`$. Furthermore, given that the SDW state, which should be sensitive to disorder, readily develops in Q-ClO<sub>4</sub>, we speculate that in spite of the anion disorder, the material is in some sense still relatively clean, and the large value of $`T_D`$ may be characteristic of the SDW (antiferromagnetic ) nature of these systems. One other feature of the low temperature data is shown in Fig. 1b. Below about 3 K there are two distinct changes (shown by the arrows) in the magnetoresistance, one at around 12 T, and the other around 30 T. At present we have no explanation for these features. One possibility is that this behavior is related to the $`T^{}`$ mechanism in some manner. Another is that these changes are some vestige of the anion-ordered state that only shows up when the resistivity of the predominantly quenched state is sufficiently high. (The characteristic fields are close to some of the major FISDW phase boundaries in the R-ClO<sub>4</sub> compound.) In summary, the present work shows the evolution of the metal-to-spin density wave transition for open orbit, quasi-one dimensional metals with changes in their nesting parameters. By mapping the magnetic field dependence of the spin-density wave transition, we provide a description of a new field dependent ground state of the material (TMTSF)<sub>2</sub>ClO<sub>4</sub> where its low temperature anion ordering transition has been suppressed. The theory of Bjelis and Maki provide an excellent framework to describe the general field dependent features of these systems. New aspects of the Q-ClO<sub>4</sub> ground state include the full description of the temperature dependence of the quantum oscillations in the spin density wave phase, which attenuate exponentially below 5 K, and a nearly quadratic field dependence of the resistance in the lower field, high temperature metallic phase of the system. Some evidence is also present for additional sub-phase structure in the B-T phase diagram, but further work will be needed to make accurate assignments to these features. This work should serve as a guide to future investigations of the magnetic field dependent mechanisms and properties of quasi-one dimensional metals. Acknowledgments: We are indebted to K. Maki for testing the convergence of Eq. 1 at low fields, and to D. Agterberg who gave us valuable advice on the computation involved. Support from NSF-DMR 95-10427 and 99-71474 (J.S.B.) and a cooperative agreement between NSF-DMR95-27035 and the State of Florida (NHMFL) is acknowledged. FIGURE CAPTIONS Figure 1. a) Magnetoresistance $`\mathrm{\Delta }`$R/R(B=0) of thermally quenched (TMTSF)<sub>2</sub>ClO<sub>4</sub> vs. pulsed magnetic field above 5 K (sample #1). Inset: zero field c-axis resistivity vs. temperature. b) Low temperature magnetoresistance vs. pulsed magnetic field for sample $`\mathrm{\#}1`$ and $`\mathrm{\#}2`$. Arrows indicate the position of structure in the lowest temperature data (see text). Figure 2. Detail of the resistance of sample $`\mathrm{\#}1`$ vs. the square of the magnetic field for different temperatures. The solid line is a fit of the high temperature (10.4 K) resistance to a quadratic field dependence. $`B_{SDW}`$ ($`T_{SDW}`$) is defined as the point of deviation of the resistance from a quadratic behavior in field (open circles), and $`B_{SDW}^{}`$ ($`T_{SDW}^{}`$) is a second change in the field dependence (see arrows) of the resistance at higher fields. (Complete temperature labels, left to right, are 5.5, 6.45, 6.5, 6.54, 7.25, 7.52, 7.90, 8.5, 8.98,9.0,10.0, and 10.6 K.) Figure 3. Temperature vs. magnetic field phase diagram of Q-ClO<sub>4</sub> based on the values of $`T_{SDW}`$ vs. $`B_{SDW}`$ ( also $`T_{SDW}^{}`$ vs. $`B_{SDW}^{}`$) obtained from Fig. 2. The solid lines are fits of the theoretical expression of Bjelis and Maki from Eq.1 with the parameters $`v_F`$ , $`T_0`$, and $`\epsilon _0`$ for X = AsF<sub>6</sub> ($`v_F`$ = 2.4 x 10<sup>5</sup> $`m/s`$; $`\epsilon _0`$ = 7 K; $`T_0`$ = 11.5 K), X = NO<sub>3</sub> ($`v_F`$ = 2.4 x 10<sup>5</sup> m/s; $`\epsilon _0`$ = 13 K; $`T_0`$ = 11.0 K), and X = Q-ClO<sub>4</sub> ($`v_F`$ = 1.1 x 10<sup>5</sup> m/s; $`\epsilon _0`$ = 22.5 K; $`T_0`$ = 13.0 K). The dashed line is a polynomial fit to $`T_{SDW}^{}`$ vs. $`B_{SDW}^{}`$. The lower phase diagram is for R-ClO<sub>4</sub> after Ref. . The finely dashed line is the main second order phase boundary between the metallic and FISDW phases. The other two phase lines T<sub>H</sub> and T<sub>R</sub> delineate sub-phases (See Ref. .) Figure 4. Temperature dependence of the quantum oscillations (normalized by the background magnetoresistance). The frequency is 190 $`\pm `$ 5 T. a) Low temperature behavior (sample #2) below 5 K. b) High temperature behavior (sample #1) above 5 K. c) Plot of the amplitudes vs. temeperature in the range 30 to 50 T for both samples measured. The solid line is a fit to a modified version of the standard Lifshitz-Kosevich formula for quantum oscillation amplitudes in metals (see text).
warning/0005/hep-ph0005103.html	ar5iv	text	# Experimental Bounds on Masses and Fluxes of Nontopological Solitons ## I Introduction In quantum field theory there exist “nontopological solitons” such as Q-balls, Fermi-balls, and neutrino-balls, the stabilities of which are based on conservation of global U(1) charges, rather than on topological quantum numbers. For example, Q-balls are stabilised by the conservation of a U(1) charge of scalar fields, while Fermi-balls and neutrino-balls are stabilised by the conservation of the number of such fermions that have Yukawa couplings with scalar fields. If such nontopological solitons exist, they may solve or at least be closely related to important problems in cosmology such as dark matter, the baryon number asymmetry of the universe, and gamma-ray bursts. Although the ideas of such nontopological solitons are very attractive, the qualitative properties of them, mass scale, charge size, typical energy scale, and cosmic abundance, are so ambiguous that there are many orders of magnitude in the parameter space that must be considered. It is desirable to make the parameter regions of nontopological solitons as narrow as possible through the use of currently available observational data. This will clarify the regions that are to be searched through experimentation in the near future. Although there are reports of a number of useful observational or phenomenological analyses to examine the allowed parameter regions of Q-balls, there does not yet exist an available analysis. In fact, Kusenko et al. discussed the experimental signatures of Q-balls and pointed out powerful detection methods for neutral Q-balls. The Gyrlyanda experiments at Lake Baikal reported the flux limit of neutral Q-balls applying these methods to their monopole search experiments. The report also gave rough estimates of bounds to be obtained with other monopole search experiments, ”Baksan” scintillators and old aged mica. Bakari et al. calculated the energy losses of Q-balls in matter and concluded the various MACRO detectors can search for charged Q-balls <sup></sup><sup></sup>* The explicit values of the flux upper limits on Q-balls are not available yet. . When it comes to Fermi-balls, no experimental limits on the flux have been reported to date. In this paper, we comprehensively discuss the flux limits of Q-balls and Fermi-balls. As their mass could be very large, we use various kinds of results from direct searches for supermassive magnetic monopoles, nuclearites, and/or heavy primary cosmic rays. We also estimate the sensitivity of present and future experiments for the purpose of searching for nontopological solitons. While these experiments were mainly designed for different purposes, we stress their importance in the specific analysis of nontopological solitons. In the following discussion, we analyze the flux bounds taking the charge $`Z_Q`$ for the case of $`Z_Q=0,1,2,3,10,\mathrm{and}137`$, as typical values. If the charge is larger than 137, the Q-ball will have the electromagnetic properties at low energy similar to the case of $`Z_Q=137`$, having a geometrical cross section $`\pi R_Q^2`$ with the effective radius of at least $`R_Q1\text{Å}`$ . For the case of Fermi-balls, we assume that the electric charge to be large enough in order to assure its stability against deformation and fragmentation, and that the electromagnetic properties at low energy is the same as the case of $`Z_F=137`$ unless the radius exceeds $`1\mathrm{\AA }`$. ## II Experimental Bounds on Q-ball ### A Q-ball Properties In the following our arguments will be restricted to the case of thick wall Q-balls, since most attractive Q-balls in SUSY are of this case. We briefly review the properties of Q-balls in order to clarify our notations and assumptions in estimating the flux bounds. Q-balls consist of a complex scalar field with a conserved global U(1) charge. The scalar field, $`\phi `$ vanishes outside the Q-ball, while it forms a coherent state inside with a time-dependent phase, $$\phi (t,\stackrel{}{x})=\phi (\stackrel{}{x})e^{i\omega t}.$$ (1) In this circumstance, the Q-ball is at the lowest energy state with a fixed global U(1) charge. This is derived by minimizing the Q-ball mass $`M_Q`$, $$M_Q=𝑑x\left[\left\|_t\phi \right\|^2+\left\|\stackrel{}{}\phi \right\|^2+V(\left\|\phi \right\|^2)\right],$$ (2) under the constraint on the charge $`Q`$, $$Q=i𝑑x\left[\phi ^{}_t\phi \phi _t\phi ^{}\right].$$ (3) Here we assume We stress that our result does not lose generality by making this assumption. The experimental flux limits, which we obtain later as a function of mass of the Q-ball, are independent of this assumption in case where Q-balls are charged. This is due to the fact that it solely depends on the mass of the Q-ball. However, in the case of neutral Q-balls, they are dependent on this assumption, since they rely only on their cross section rather than mass. Of course, the flux limitations of neutral Q-balls were independent of the assumption if we expressed them as a function of the cross section. Readers, who are interested in neutral Q-balls with another type of potential, can easily accommodate our results by matching the mass of the Q-balls with such of ours, that has the same cross section as them. that for large $`\phi `$ the scalar potential $`V(\left\|\phi \right\|^2)`$ is almost “flat”, i.e., $`V(\left\|\phi \right\|^2)M_S^4`$, where $`M_S`$ is a constant with a mass dimension. (The potential of this type is known to be present in supersymmetric theories, in which $`M_S`$ is a SUSY breaking scale). In this case, the Q-ball mass is obtained as $$M_Q=\frac{4\pi \sqrt{2}}{3}M_SQ^{3/4},$$ (4) with its radius, $$R_Q=\frac{1}{\sqrt{2}}M_S^1Q^{1/4}.$$ (5) The proportionality, $`M_QQ^{3/4}`$, in Eq.(4) leads to the stability of the Q-ball for $`Q`$ large enough. If we consider a baryonic Q-ball, i.e., a B-ball, A B-ball is the Q-ball the conserved U(1) charge of which is baryon number. for example, the condition of stability against nucleon emission is $`m_{\mathrm{proton}}>\frac{}{Q}M_Q`$, i.e., $$Q>5.0\times 10^{14}\left(\frac{M_S}{\text{TeV}}\right)^4.$$ (6) It is pointed out that such large Q-balls may have been created in the early universe. If this is the case and such large stable Q-balls have survived until present day, they would contribute toward the dark matter in the Galaxy. Their flux should then satisfy $$FF_{DM}\frac{\rho _{DM}v}{4\pi M_Q}7.2\times 10^5\left(\frac{\text{GeV}}{M_Q}\right)\mathrm{cm}^2\mathrm{sec}^1\mathrm{sr}^1,$$ (7) where $`\rho _{DM}`$ is the energy density of the dark matter in the Galaxy, $`\rho _{DM}0.3\mathrm{GeV}/\mathrm{cm}^3`$, and $`v`$ is the Virial velocity of the Q-ball, $`v3\times 10^7\mathrm{cm}/\mathrm{sec}`$. In the following analysis for Q-balls and Fermi-balls, we assume for simplicity the velocity of Q-balls to be $`v=10^3c`$, where $`c`$ is the light velocity. In order to estimate efficiencies in the detection of Q-balls with various detectors, Kusenko et al. classified relic solitons into two groups according to the properties of their interaction with matter: Supersymmetric Electrically Neutral Solitons(SENS) and Supersymmetric Electrically Charged Solitons(SECS). In this paper we use the terms “SENS” and “SECS” only for Q-balls. In case of SENS, a process similar to proton decay may occur in the thin layer of the Q-ball’s surface, when the Q-ball collides with nuclei. The energy release of roughly 1 GeV per nucleon is carried away by pions through this process. We call this the “KKST process”. In case of SECS, with a positive charge, the KKST process in collision with nuclei will be strongly suppressed by Coulomb repulsion. Only electromagnetic processes will take place in this circumstance. A charged Q-ball with a small velocity is accompanied by an atomic-like cloud of electrons and interacts with matter similarly to an atom with a heavy mass. If the charge is very large, the effective interaction radius is approximately $`1\text{Å}`$, which is similar to the case of ”nuclearites”. Negatively charged Q-balls are of little interest to us at present. If negatively charged Q-balls happen to exist, the detection of them should be much easier than that of positively charged Q-balls. This is due to both the electromagnetic process and the KKST process occurring during collision. As a result of this, the excluded parameter region of the positively charged Q-ball is also excluded for the negatively charged ones. For the purpose of this paper, we will only discuss the neutral and positively charged Q-balls. ### B Bounds on Flux and Mass of Neutral Q-ball (SENS) As pointed out by Kusenko et al., a neutral Q-ball would produce the KKST process, when it collides with a nucleus, absorbing it and emitting pions with the total energy of approximately 1 GeV per nucleon. If the cross section of such a process is large and successive events of this type are detected along a single trajectory, it will be a signal detecting a Q-ball. We note that such a Q-ball process is different from the Rubakov effect even though they are similar to each other. The cross section for the former process depends little on whether the target nucleus possesses a magnetic moment or not, while the cross section for the latter process strongly depends on it. <sup>§</sup><sup>§</sup>§ If the nucleus has a finite magnetic moment, the cross section of the Rubakov process should have an enhancement factor of $`1/\beta ^210^6`$ with $`\beta `$ being the relative velocity (in unit of the light velocity) between the monopole and the nucleus. Since the interaction cross section of the KKST process is expected of a geometrical size, the energy loss of SENS is determined by $$\frac{dE}{dx}=\pi R_Q^2v^2\rho ,$$ (8) similar to the case of nuclearite collision with a nucleus. Here $`R_Q`$ and $`\rho `$ are the radius of the Q-ball and the density of the target matter, respectively. In order for the Q-balls to reach the detector site and penetrate the detector, the following condition should be satisfied as: $$M_Q\frac{1}{7.3}\pi R_Q^2\rho 𝑑x$$ (9) with integration over the trajectory of Q-balls. This leads to $$M_Q4.2\times 10^{39}\left(\frac{\rho L}{\mathrm{gr}\mathrm{cm}^2}\right)^3\left(\frac{\mathrm{TeV}}{M_S}\right)^8\mathrm{GeV},$$ (10) where $`L`$ is the path length for the Q-ball to traverse matter (air, water, and/or rock) and to penetrate the detector. Since we are interested in the parameter regions of $`M_S>100\mathrm{G}\mathrm{e}\mathrm{V}`$, $`\rho 10\mathrm{gr}\mathrm{cm}^3`$, $`L2R_{\mathrm{Earth}}1.3\times 10^{10}\mathrm{cm}\mathrm{gr}\mathrm{cm}^2`$ and $`M_QM_S`$, this condition is always satisfied in any experiments on the earth. Most experimental searches for monopole-catalyzed proton decay (the Rubakov effect) are also sensitive to the KKST process, and are able to give stringent bounds on SENS flux. Let us review the available data, not only the already reported results of the Gyrlyanda experiments at Lake Baikal on SENS flux but also other typical experiments reinterpreting them to get SENS flux bounds. The Gyrlyanda experiments reported that the flux of SENS has the bound : $`F<3.9\times 10^{16}\mathrm{cm}^2\mathrm{sec}^1\mathrm{sr}^1,`$ (11) if the cross section for the KKST process is $`\sigma >1.9\times 10^{22}\text{cm}^2`$. This corresponds to the lower limit of the SENS mass, $`M_Q>1.0\times 10^{21}\left({\displaystyle \frac{M_S}{\text{TeV}}}\right)^4\text{GeV},`$ (12) since Eq.(4) and Eq.(5) relate $`\sigma `$ to $`M_Q`$ as $$\sigma =1.9\times 10^{36}\left(\frac{\text{TeV}}{M_S}\right)^{8/3}\left(\frac{M_Q}{\text{GeV}}\right)^{2/3}\text{cm}^2.$$ (13) The authors of Ref. also estimated the SENS flux limit which will be given in the BAKSAN experiments: $`F<3.0\times 10^{16}\mathrm{cm}^2\mathrm{sec}^1\mathrm{sr}^1,`$ (14) if the cross section is $`\sigma >5.0\times 10^{26}\text{cm}^2`$, which corresponds to $`M_Q>4.2\times 10^{15}\left({\displaystyle \frac{M_S}{\text{TeV}}}\right)^4\text{GeV}.`$ (15) Although the above two experiments give considerably stringent bounds to SENS, we must examine how stringent limits will be given by other monopole searches and future cosmic ray experiments. The Kamiokande Cherenkov detector, which has 3 000 tons of water, approximately 1 000 photo-tubes and is located at approximately 1 000 meters underground, gave the limits of monopole with 335 days of live time. We reinterpret the results to obtain SENS flux limits approximately: $`F<3\times 10^{12},3\times 10^{14},\mathrm{and}3\times 10^{15}\mathrm{cm}^2\mathrm{sec}^1\mathrm{sr}^1,`$ (16) for SENS with $`\sigma =0.1,1,\mathrm{and}10\text{mb}`$, which correspond to $`M_Q=4.0\times 10^{11},1.2\times 10^{13},\mathrm{and}5.6\times 10^{13}\left({\displaystyle \frac{M_S}{\text{TeV}}}\right)^4\text{GeV},`$ (17) respectively. The Super-Kamiokande experiments with 50 000 tons of water, the Cherenkov detector would obtain more stringent flux limits than those of Kamiokande by almost two orders of magnitude and with three years of observation time. The MACRO, large underground detectors consisting of three kinds of subdetectors ( i.e. scintillation counters, streamer tubes, and nuclear track detectors(CR-39)) could search for SENS. Although the flux limits of SENS are reported to be obtained as $`F10^{16}\mathrm{cm}^2\mathrm{sec}^1\mathrm{sr}^1`$ in Ref., we do not give this flux limit here since the detection efficiency is too difficult for us to estimate. The AMANDA Cherenkov detectors located at the South Pole under ice, are designed mainly to detect relativistic particles. If they could also detect slow particles, the large area of the array would be of great help. For a few years to observe SENS, they could obtain their flux limits as $`F10^{17}\mathrm{cm}^2\mathrm{sec}^1\mathrm{sr}^1/ϵ`$ with the detection efficiency, $`ϵ`$. The flux limits, however, are not estimated here since the detection efficiencies are also difficult for us to evaluate. The TA(Telescope Array Project) with effective aperture of $`S\mathrm{\Omega }ϵ=6\times 10^3\mathrm{km}^2\mathrm{sr}`$ ($`S\mathrm{\Omega }=6\times 10^4\mathrm{km}^2\mathrm{sr}`$ and the duty factor $`ϵ=0.1`$) is planned to detect cosmic rays with energy, $`10^{16}10^{20}\text{eV}`$, by detecting the air fluorescence. Using a special trigger, it may detect slow particles with large energy loss. With such a trigger, it may be possible to search for SENS flux at the level of $`F<1\times 10^{21}\mathrm{cm}^2\mathrm{sec}^1\mathrm{sr}^1`$, provided that the fluorescence light yield of the event is equivalent to that of air showers with the minimum energy $`E_{min}=10^{16}\text{eV}`$. Considering energy release of roughly 1 GeV per absorbed nucleon in the KKST process, one obtains the condition of detectability in TA: $$\sigma \rho 𝑑xE_{min}/\xi _{SENS},$$ (18) where $`\xi _{SENS}`$ is the ratio of the efficiency of fluorescence light yield per total energy loss of SENS to that of extensive air showers. In the KKST process we estimate $`\xi _{SENS}1,`$ since emitted pions are relativistic. Assuming the average density of air $`\rho 0.5\times 10^3\mathrm{g}/\mathrm{cm}^3`$ and the path length of SENS in the atmosphere $`L20\text{km}`$, we obtain the lower limit of SENS mass from Eq.(18), $`M_Q>0.7\times 10^{24}\left({\displaystyle \frac{M_S}{\text{TeV}}}\right)^4\text{GeV}.`$ (19) The OA(OWL-AIRWATCH) detector located on a satellite with effective aperture $`5\times 10^5\mathrm{km}^2\mathrm{sr}`$ (with duty cycle of $`0.2`$ taken into account), is used to observe the atmospheric fluorescence of air showers with energy $`10^{19}10^{20}`$ eV or more. If it could also detect slow particles with high efficiency, it would be possible to search for SENS during a one year time constraint at the flux level of $`F<10^{23}\mathrm{cm}^2\mathrm{sec}^1\mathrm{sr}^1`$. For SENS to yield the amount of fluorescence light comparable to that of air showers with $`10^{19}`$eV, their mass should satisfy $`M_Q>2\times 10^{28}\left(\frac{M_S}{\text{TeV}}\right)^4\text{GeV}`$. The bounds given in the various experiments mentioned above are summarized in Fig.1(a). This figure also gives the flux to be expected if the dark matter of the Galaxy consists mainly of SENS(see Eq.(7)). Note that these bounds of excluded regions depend on $`M_S`$. ( The lower mass bound in each experiment is calculated taking $`M_S=1\text{TeV}`$ in this figure. In case of TA and OA, we added those for $`M_S=100\mathrm{GeV}`$. ) The region in which it is expected that future experiments are able to be conducted, are also given there. If we assume that the dark matter of the Galaxy consists mainly of SENS ( i.e. if the flux is just on the DM-limit line of Fig.1(a) ), we obtain the excluded regions of SENS mass as functions of the parameter $`M_S`$ as in Fig.1(b). This figure shows that SENS with the U(1) charge $`Q10^{22}`$ are cosmologically interesting in the region of $`M_S6\times 10^3\mathrm{GeV},`$ while that TA and OA could search for SENS of $`10^{25}Q10^{35}`$ in the region of $`M_S6\times 10^3\mathrm{GeV}.`$ ### C Bounds on Flux and Mass of Charged Q-ball (SECS) Charged Q-balls (SECS) interact with matter in a similar way to nuclearites. The rate of energy loss of SECS in matter is $$\frac{dE}{dx}=\sigma \rho v^2,$$ (20) where $`\sigma `$ is the cross section of SECS collision with matter, $`\rho `$ is the density of the matter, and $`v`$ is the relative velocity of SECS and the matter. Substituting $`E=\frac{1}{2}M_Qv^2`$ in the above equation and integrating with respect to $`x`$, one obtains $$\rho L=\frac{M_Q}{\sigma }\mathrm{ln}\frac{v_0}{v_c},$$ (21) where $`L`$ is the range of SECS in the medium, $`v_0`$ is an initial velocity of SECS($`3\times 10^7\mathrm{cm}/\mathrm{sec}`$), and $`v_c`$ is a final velocity of SECS. We use the estimation in the case of nuclearites, $`v_c1.7\times 10^4\mathrm{cm}/\mathrm{sec}`$ in rock. SECS with a velocity below this value are quickly brought to rest. In the case of SECS with a radius larger than $`1\text{Å}`$ ( i.e., $`M_Q>2.1\times 10^{30}\left(M_S/\text{TeV}\right)^4\mathrm{GeV}`$ ), the cross section is simply given by $`\sigma =\pi R_Q^2`$. From Eqs.(4), (5) and (21), we obtain $$\rho L=6.2\times 10^{12}\left(\frac{M_S}{\text{TeV}}\right)^{8/3}\left(\frac{M_Q}{\text{GeV}}\right)^{1/3}\mathrm{gr}/\mathrm{cm}^2.$$ (22) Such heavy SECS should, however, have too small a flux to be detected according to Eq.(7), and thus, are not interesting to us. We have considered only SECS with a radius smaller than $`1\text{Å}`$ in the following. If such SECS have a significantly large electric charge ($`Z_Q\alpha ^1137`$ where $`\alpha `$ is the fine structure constant), the cross section of their collision with matter is not controlled by their intrinsic radius $`R_Q`$, but by the size of the surrounding electron cloud. This cloud is never smaller than $`1\text{Å}`$. We note that this situation is similar to the case of nuclearites. We have $`\sigma =\pi R_{eff}^2`$ with $`R_{eff}=1\text{Å}`$, which leads to the relation $$\rho L=4.0\times 10^8\left(\frac{M_Q}{\text{GeV}}\right)\mathrm{gr}/\mathrm{cm}^2.$$ (23) Equation (23) shows that SECS with $`Z_Q137`$ should be as heavy as $$M_Q>2.5\times 10^7\left(\frac{\rho L}{\mathrm{gr}/\mathrm{cm}^2}\right)\text{GeV}$$ (24) in order to penetrate the medium with density $`\rho `$ and length $`L`$. If the charge of SECS is small, i.e. $`Z_Q137`$ (and the intrinsic SECS radius $`R_Q`$ is also smaller than $`1\text{Å}`$) , the effective cross section should be smaller than $`\pi (1\mathrm{\AA })^2`$. We then need more delicate treatments to estimate the rate of the energy loss. It is known that there are two kinds of interaction that contribute to the energy loss of SECS at low velocity $`\beta 10^3`$ ( i.e. interaction with electrons and interaction with nuclei ), $`dE/dx=(dE/dx)_{\mathrm{electrons}}+(dE/dx)_{\mathrm{nuclei}}`$. The rate of electronic energy loss is estimated as Ref., $$\left(\frac{dE}{dx}\right)_{\mathrm{electrons}}=8\pi \alpha a_0\frac{v}{v_0}N_e\frac{Z_Q^{7/6}}{\left(Z_Q^{2/3}+Z^{2/3}\right)^{3/2}},$$ (25) where $`v_0`$ is given by $`\alpha c`$, $`N_e`$ is the number density of electrons in the medium, $`Z`$ is the atomic number of the medium, and $`a_0`$ is the Bohr radius. The rate of energy loss caused by interaction with the nuclei of the medium is given in Ref. $$\left(\frac{dE}{dx}\right)_{\mathrm{nuclei}}=4\pi \alpha aN_ZZ_QZ\frac{M_Q}{M_Q+M}\frac{A\mathrm{ln}Bϵ}{Bϵ(Bϵ)^C},$$ (26) with $`A=0.56258`$, $`B=1.1776`$, $`C=0.62680`$, $`M`$ is the mass of the target nucleus, $`N_Z`$ is the number density of target nuclei and $$ϵ=\frac{aMM_Q\beta ^2}{2\alpha Z_QZ(M_Q+M)},a=\frac{0.8853a_0}{\left(\sqrt{Z_Q}+\sqrt{Z}\right)^{2/3}}.$$ (27) In case of $`M_QM`$ we calculated $`(dE/dx)_\rho (dE/dx)/\rho `$ (which depends only loosely on the density of the medium) for $`\mathrm{SiO}_2`$ and obtained $`(dE/dx)_\rho =0.16,0.41,0.71`$, and $`2.9\mathrm{GeV}/\mathrm{gr}/\mathrm{cm}^2`$ for SECS with typical charge $`Z_Q=1,2,3,\mathrm{and}10,`$ respectively. We also find the energy loss rate $`(dE/dx)_\rho `$ depends little on matter medium(air, water, and/or rock). We use the above values for the following estimation of energy loss. The condition for SECS to penetrate the medium with density $`\rho `$ and length $`L`$ is roughly given by $$\frac{1}{2}M_Q\beta ^2>\left(\frac{dE}{dx}\right)_\rho \rho L.$$ (28) Since SECS have a large cross section when they collide with matter, they should have kinetic energy large enough to reach a detector and to penetrate it. This situation is unlike that of SENS. Let us take the following experiments as typical in detecting SECS. The MACRO with its three subdetectors are sensitive to SECS for any values of $`Z_Q`$; in order for SECS to reach the underground detector site from above with $`\rho L3.7\times 10^5\mathrm{gr}\mathrm{cm}^2`$, the lower limit of the SECS mass $`M_Q`$ is given by Eq.(24) and Eq.(28), as $`M_Q>1.2\times 10^{11},3.0\times 10^{11},5.3\times 10^{11},2.1\times 10^{12},\mathrm{and}9.3\times 10^{12}`$ GeV for $`Z_Q=1,2,3,10,\mathrm{and}137,`$ respectively. The flux limit suggested by the results of no events of monopole search experiments is roughly $`F10^{16}\mathrm{cm}^2\mathrm{sec}^1\mathrm{sr}^1`$. It is obvious that the inclusion of searches for Q-balls from below, decreases the flux upper limit by a factor of two, although the lower limits for the mass to reach the detector increase by a factor of $`(\rho L)_{\mathrm{from}\mathrm{below}}/(\rho L)_{\mathrm{from}\mathrm{above}}=(6.6\times 10^9)/(3.7\times 10^5)1.8\times 10^4`$ . This feature for SECS is common to other experiments such as OYA, NORIKURA, KITAMI, AKENO, UCSD$`II`$, KEK, and MICA. The OYA experiments with CR-39 plastic track detectors (located at the depth of $`\rho L=10^4\mathrm{cm}^2\mathrm{sec}^1\mathrm{sr}^1`$), searched for monopoles and nuclearites for 2.1 years. The flux limit of these experiments corresponds to that of SECS, $`F<3.2\times 10^{16}\mathrm{cm}^2\mathrm{sec}^1\mathrm{sr}^1`$. Since the detectors are sensitive to the restricted energy loss larger than $`0.14\mathrm{GeV}\mathrm{gr}/\mathrm{cm}^2`$, it is possible that SECS with a charge of $`Z_Q2`$ may be detected. The condition for SECS to reach the detector is $`M_Q>8.2\times 10^9\mathrm{GeV},1.4\times 10^{10}\mathrm{GeV},5.8\times 10^{10}\mathrm{GeV},\mathrm{and}2.5\times 10^{11}\mathrm{GeV}`$ for $`Z_Q=2,3,10,\mathrm{and}137,`$ respectively. The KEK experiments with scintillation counters placed at ground level should be sensitive to SECS at any value of $`Z_Q`$. As the detection threshold is $`0.01I_{min}`$ where $`I_{min}`$ is the minimum energy loss of the ionizing particle with $`Z=1`$. We interpret that their flux upper limit for strange quark matter (nuclearites) holds for the SECS: $`F<3.2\times 10^{11}\mathrm{cm}^2\mathrm{sec}^1\mathrm{sr}^1`$. The mass bounds for SECS to reach detectors located at ground level are given by $`M_Q>3.1\times 10^8\mathrm{GeV},8.2\times 10^8\mathrm{GeV},1.4\times 10^9\mathrm{GeV},5.8\times 10^9\mathrm{GeV},\mathrm{and}2.5\times 10^{10}\mathrm{GeV}`$ for $`Z_Q=1,2,3,10,\mathrm{and}137`$, respectively. The KITAMI experiments with CN (cellulose nitrate) nuclear track detectors installed in houses at sea level, give the flux limit $`F<5.2\times 10^{15}\mathrm{cm}^2\mathrm{sec}^1\mathrm{sr}^1`$. We reinterpret this as the SECS flux limit. Since the ionization energy loss of SECS should be larger than $`1.3\mathrm{GeV}/\mathrm{gr}/\mathrm{cm}^2`$ for a chemically etchable track to be made in CN, we estimate that SECS with $`Z_Q10`$ can be detected with the CN detectors. At ground level with $`\rho L=10^3\mathrm{gr}\mathrm{cm}^2`$, SECS heavier than $`5.8\times 10^9\mathrm{GeV}\mathrm{and}2.5\times 10^{10}\mathrm{GeV}`$ can penetrate the atmosphere and reach the detector from above. This is the case for $`Z_Q=10\mathrm{and}137`$, respectively. The AKENO experiments with Helium-gas counters located at ground level should be sensitive to SECS at any value of $`Z_Q`$. The detector is composed of many layers of proportional counters and concrete shields. The detector is sensitive to tracks with an ionization level larger than $`10I_{\mathrm{mim}}`$. This should therefore assure the detectability for SECS with any value of $`Z_Q`$. SECS with masses of $`M_Q>5.8\times 10^8\mathrm{GeV},1.5\times 10^9\mathrm{GeV},`$ $`2.6\times 10^9\mathrm{GeV}`$ $`,1.1\times 10^{10}\mathrm{GeV},\mathrm{and}4.7\times 10^{10}\mathrm{GeV}`$, can penetrate the earth and the concrete shields of the detector from above for $`Z_Q=1,2,3,10,\mathrm{and}137`$, respectively. Due to the fact that no events have been detected, we estimate the flux limit as $`F<1.8\times 10^{14}\mathrm{cm}^2\mathrm{sec}^1\mathrm{sr}^1`$. The UCSD$`II`$ experiments with He-CH<sub>4</sub> proportional tubes placed at ground level, should be sensitive to SECS with any value of $`Z_Q`$, as the detection threshold is $`0.09I_{mim}`$. The flux upper limit for monopoles $`F=1.8\times 10^{14}\mathrm{cm}^2\mathrm{sec}^1\mathrm{sr}^1`$ should be the same as that for SECS. The mass bounds are the same as in the above case of the KEK experiments. The NORIKURA CR-39 experiments were conducted at the top of Mt.Norikura in the search for monopoles and strange quark matter. They are sensitive to tracks with ionization larger than $`0.35\mathrm{GeV}/\mathrm{gr}/\mathrm{cm}^2`$. This assures that SECS with $`Z_Q3`$ will be detectable. Since the detector is installed at a high altitude, it is more sensitive to lighter SECS compared to detectors such as those used by KITAMI, AKENO, UCSD$`II`$, and KEK, which are installed at ground level. The flux limit is $`F=2.2\times 10^{14}\mathrm{cm}^2\mathrm{sec}^1\mathrm{sr}^1`$. The MICA analysis with ancient mica crystals of $`0.60.9\times 10^9`$ years old, was made to search for monopoles. It should also be sensitive to SECS with $`Z_Q10`$ as the detection threshold is $`2.4\mathrm{GeV}/\mathrm{gr}/\mathrm{cm}^2`$ . In order for SECS coming from above to reach the mica crystals located at 3 km underground with $`\rho L=7.5\times 10^5\mathrm{gr}/\mathrm{cm}^2`$, SECS should be heavier than $`4.4\times 10^{12}\mathrm{GeV}\mathrm{and}1.9\times 10^{13}\mathrm{GeV}`$ for $`Z_Q=10\mathrm{and}137`$, respectively. In case of the monopole search, the capture of an aluminium atom by a monopole was taken into account. The detection efficiency was estimated at $`0.15`$ for monopoles coming from above and 0 for those from below, due to this consideration. Therefore, this decrease of efficiency caused the flux limits to be less stringent. In case of SECS search, however, SECS need not capture an aluminum atom and can be heavy enough to penetrate the earth. With these considerations we estimate the flux limit for SECS, which is better than that for monopoles by a factor of 6, to have $`F=2.3\times 10^{20}\mathrm{cm}^2\mathrm{sec}^1\mathrm{sr}^1`$ for SECS from above and a value improved by a factor of twelve for SECS from any direction. The SKYLAB experiments with Lexan track detectors (installed in the SKYLAB workshop located in space), reported that no superheavy relativistic nuclei ($`Z110`$) had been detected, which correspond to $`Z_Q10`$ in our case of slow particles. The Lexan detector and its container have a thickness $`\rho L=2\mathrm{g}\mathrm{r}\mathrm{cm}^2`$ in total, or effectively $`\rho L3\mathrm{g}\mathrm{r}\mathrm{cm}^2`$, giving the mass a lower limit of SECS $`M_Q>1.7\times 10^7\mathrm{GeV}\mathrm{and}7.5\times 10^7\mathrm{GeV}`$ for $`Z_Q=10\mathrm{and}137`$, respectively. This experiment gives the flux limit of SECS as $`F<3.8\times 10^{12}\mathrm{cm}^2\mathrm{sec}^1\mathrm{sr}^1`$. The AMS experiments may detect SECS in the future if a special trigger designed to detect slow particles is used. The magnetic spectrometer on the space station would have a large area, allowing us to record the flux upper limit as $`F\times 10^{11}\mathrm{cm}^2\mathrm{sec}^1\mathrm{sr}^1`$ given a one year observation time frame. The detection with thickness of $`\rho L10\mathrm{g}\mathrm{r}\mathrm{cm}^2`$ needs mass $`M_Q>3.1\times 10^6\mathrm{GeV},8.2\times 10^6\mathrm{GeV}`$, $`1.4\times 10^7\mathrm{GeV},5.8\times 10^7\mathrm{GeV},\mathrm{and}2.5\times 10^8\mathrm{GeV}`$ for $`Z_Q=1,2,3,10\mathrm{and}137`$, respectively. The experimental data and also the future possibilities mentioned above are summarized in Figures2-6. Figures 2(a)-6(a) represent limits of the flux and the mass for SECS. Figures2(a)-4(a) show that future AMS experiments and current MACRO experiments are very sensitive to a wide parameter region of mass, charge and flux of SECS. On the other hand, Figure5(a) and Figure6(a) show that the SKYLAB and MICA experiments give us a stringent exclusion of the lower and upper region of mass respectively. Figures2(b)-6(b) represent bounds on the U(1) charge $`Q`$ of SECS, versus the symmetry breaking parameter $`M_S`$ in the case that SECS are mainly contributing to dark matter in the Galaxy. If we assume that SECS are B-balls and impose their condition of stability, these figures show that only SECS with $`Q10^{2226}`$ remain to be considered. ## III Experimental Bounds on Fermi-ball ### A Fermi-ball Properties We have investigated which regions are allowed for the flux of Fermi-balls and discuss our findings in this section. The Fermi-ball, which is another kind of nontopological soliton, was first proposed by Macpherson and Campbell as a possibility for explaining the make-up of the dark matter in Our Galaxy. Fermi-balls with a large radius are found to be unstable against deformation and are therfore expected to fragment into very small Fermi-balls. The Fermi-ball then interacts with matter in a weak manner and seems too difficult to be detected. It was then proposed by Morris that Fermi-balls be electrically charged in order to improve its stability against deformation and fragmentation. Although a Fermi-ball with large fermion number is energetically unstable and could fragment into many smaller Fermi-balls, the Coulomb force is expected to suppress the process of fragmentation, increasing the life time of the large Fermi-ball so that it would become more stable. Owing to the electromagnetic interaction, various detectors become sensitive to Fermi-balls. In this paper we have focused on the experimental bounds for charged Fermi-balls rather than on neutral Fermi-balls. Let us first briefly review the properties of Fermi-balls in order to make our assumptions and terminologies clear. The charged Fermi-ball consists of three components: a scalar field, a large number of electrically charged heavy fermions, and also a large number of electrons (or positrons) which partly compensate for the electric charge of the heavy fermions. (When we consider the stability of Fermi-balls in the following, we assume that the electric charge of the heavy fermions is positive without loss of generality.) Inside the Fermi-ball, the scaler field has the value of a false vacuum (which is almost degenerate to the true vacuum The approximate degeneracy of energy densities of a true and false vacuum is necessary for the stability of large Fermi-balls. ), while outside the Fermi-ball, it has the value of a true vacuum. On the boundary wall, the energy density of the scalar field is higher than those in the two vacua. Since the heavy fermions have a large mass in both the true and false vacuum and have a vanishing mass on the wall, many fermions are strongly trapped on the wall. Since the electric field of the Fermi-ball is strong, electron-positron pairs may be created by the effect of the quantum field theory. This effect decreases the electric field strength by leaving the electrons on the surface and emitting the positrons to infinity. Assuming the difference of the energy density between two vacua is small enough to be neglected, we obtain the energy of the Fermi-ball with the number of heavy fermions $`N_F`$ and the number of electrons $`N_e`$, $`E_F=4\pi \mathrm{\Sigma }R^2+{\displaystyle \frac{2(N_F^{\frac{3}{2}}+N_e^{\frac{3}{2}})}{3R}}+{\displaystyle \frac{\alpha (N_FN_e)^2}{2R}}.`$ (29) Here, $`\mathrm{\Sigma }`$ is the surface tension of the Fermi-ball and $`R`$ is the radius of it. The radius of the stable Fermi-ball, $`R_F`$ is determined by balance We have two ways in differentiating Eq.(29) with respect to $`R`$. : either one fixes $`N_F`$ and $`N_e`$ , or one fixes $`N_F`$ and the electric field $``$ on the surface of the Fermi-ball has the critical value $`=m_e^2/e`$. Here we take the former and follow Morris. For the other possibility, see Ref.. of the surface tension energy (the first term) proportional to $`R^2`$, and the Fermi energy (the second term) or the Coulomb energy (the third term) which is proportional to $`R^1`$ as $`R_F=\left\{{\displaystyle \frac{1}{8\pi \mathrm{\Sigma }}}\left[{\displaystyle \frac{2(N_F^{\frac{3}{2}}+N_e^{\frac{3}{2}})}{3}}+{\displaystyle \frac{\alpha (N_FN_e)^2}{2R}}\right]\right\}^{\frac{1}{3}}.`$ (30) The energy of the Fermi-ball is then, $`E_F(=M_F)=12\pi \mathrm{\Sigma }R_F^2\kappa ^3R_F^2,`$ (31) which is the common relationship when the volume energy is neglected. Our following analyses of the Fermi-ball’s flux are based only on the above relation. Therefore, the experimental limits from these analyses do not depend on the details of Fermi-ball models. ### B Bounds on Flux and Mass of Fermi-ball Throughout this section, we discuss the physical parameter regions that have been excluded in experiments. If the electric charge of the heavy fermions is positive, the observational situation is almost the same as for the case where SECS has $`Z_Q137`$. This means that the effective radius can be taken as $`R_{eff}1\text{Å}`$ when the intrinsic radius is smaller than this value. In the case of the Fermi ball (not as is the case with the Q-ball), however, the radius can be larger than $`1\text{Å}`$ without becoming too heavy, as its mass is proportional to $`R_F^2`$ rather than $`R_F^3`$. The effective radius thus becomes, $`R_{eff}=\{\begin{array}{cc}\hfill 1\text{Å}& (\mathrm{for}R_F<1\text{Å})\hfill \\ \hfill R_F& (\mathrm{for}R_F1\text{Å}).\hfill \end{array}`$ (34) In the case where $`R_F`$ is large enough, Fermi-balls can be detected with not only detectors which are sensitive to SECS with $`Z_Q137`$ , but also with future detectors used for extensive air showers such as TA and OA. This detectability within future experiments is the main difference in the case of SECS with $`Z_Q137`$. We examine the bounds on Fermi-balls <sup>\</sup><sup>\</sup>\In the following, we assume that the electric charge of Fermi-balls is positive. In case Fermi-balls have a negative electric charge, it is much easier for us to detect them than those with a positive charge. This is followed by the fact that the negatively charged Fermi-balls may trap nuclei in collision with matter and emit mesons or photons with total energy of order 1GeV per nucleon. We note that all regions excluded for the positively charged Fermi-balls should also be excluded for the negatively charged Fermi-balls. to be given by AMS, SKYLAB, UCSD$`II`$, AKENO, KEK, NORIKURA, KITAMI, OYA, MACRO, MICA, TA, and OA in the following. First, we discuss the conditions in which the Fermi-ball reaches and penetrates the detectors. For Fermi-balls with $`R_{eff}=1\text{Å}`$ ($`R_F<1\AA `$), the condition is the same as that for SECS with $`Z_Q137`$, i.e. Eq.(24), $$M_F>2.5\times 10^7\left(\frac{\rho L}{\mathrm{gr}/\mathrm{cm}^2}\right)\mathrm{GeV}.$$ (35) For Fermi-balls with $`R_{eff}=R_F`$ ($`R_F>1\AA `$ ), the condition is different from that of SECS (Eq.(22)), since the relation between the mass and the radius of Fermi-balls is different from that of Q-balls. Using Eq.(21) and Eq.(31) with $`\sigma =\pi R_F^2`$, we obtain the condition $`\kappa 4.7\times 10^2\left({\displaystyle \frac{\rho L}{\mathrm{gr}/\mathrm{cm}^2}}\right)^{\frac{1}{3}}\mathrm{GeV}.`$ (36) This condition is independent of the mass of Fermi-balls in the case of $`R_F1\text{Å}`$. From Eq.(35) and Eq.(36) we see that all the experiments (except TA and OA) give the same bounds on mass and flux of Fermi-balls as on those of SECS with $`Z_Q137`$, as the detection efficiencies for these two kinds of solitons are the same for these detectors. We next discuss the conditions for detecting Fermi-balls with TA and OA. This requires a fluorescence light yield which corresponds to the air shower energy of $`E_{min}=10^{16}\text{eV}`$ for TA and $`E_{min}=10^{19}\text{eV}`$ for OA. This condition is satisfied if the energy loss of the Fermi-ball is $$\pi R_F^2\rho v^2LE_{min}/\xi _F.$$ (37) Here, the parameter $`\xi _F`$ is the ratio of the efficiency of fluorescence light yield per total energy loss for slow Fermi-balls to that for high energy cosmic ray protons. We estimated $`\xi _F`$ from the measurements of the efficiency of ionization for slow ions as $`\xi _F1/5`$. By taking the average density of air as $`\rho =5.0\times 10^4\mathrm{g}/\mathrm{cm}^3`$, the velocity of the Fermi-ball as $`v=10^3c`$ and the length of the trajectory as $`L=20\mathrm{km}`$, one obtains the required condition to observe Fermi-balls from Eq.(31) and Eq.(37) <sup>††</sup><sup>††</sup>††Here we did not assume the black body radiation from the Fermi-ball trajectory, since it is effective only for dense medium. , $`M_F6.5\times 10^{22}\left({\displaystyle \frac{\kappa }{10^3\mathrm{GeV}}}\right)^3\left({\displaystyle \frac{E_{min}}{10^{16}\mathrm{eV}}}\right)\mathrm{GeV}.`$ (38) If TA can observe slow particles <sup>‡‡</sup><sup>‡‡</sup>‡‡The TA experiments may be available for slow particle search with a special trigger, it will be able to search for Fermi-balls on order to give the same flux limit as that of SENS, $`F<1\times 10^{21}\mathrm{cm}^2\mathrm{sec}^1\mathrm{sr}^1`$. Equation(38) gives the mass bounds required to observe Fermi-balls, $`M_F6.5\times 10^{22}\left(\frac{\kappa }{10^3\mathrm{GeV}}\right)^3\mathrm{GeV}`$. If OA can observe slow particles as well, the flux bounds to be obtained are improved by two orders of magnitude $`F<1\times 10^{23}\mathrm{cm}^2\mathrm{sec}^1\mathrm{sr}^1`$, for $`M_F6.5\times 10^{25}\left(\frac{\kappa }{10^3\mathrm{GeV}}\right)^3\mathrm{GeV}`$. The flux limit with $`\kappa =10^3`$ GeV is shown in Figure7(a). This figure shows that quite broad regions are already excluded by available experimental data. Future experiments, TA and OA, may find it possible to search large regions which have not been previously accessible through current experiments. Figure7(b) shows the region of $`M_F\kappa `$ plane (hatched with solid lines and with half-tone dot meshing) to be excluded when we assume that the dark matter of the Galaxy consists mainly of Fermi-balls. Here we discuss how these results constrain the Morris’s simple Fermi-ball model. The electric field becomes the critical value ($`=m_e^2/e`$), near the surface due to the surrounding electrons in this model. The third term of Eq.(29) is $`m_e^4R^3/2\alpha `$ in this case. When this Coulomb energy is related to the surface energy as $`4\pi \mathrm{\Sigma }R^2=Cm_e^4R^3/2\alpha `$ <sup>\</sup><sup>\</sup>\ Since the surface energy must be the same order as the total energy, $`C`$ must not be much smaller than unity. Investigation into what range of $`C`$ is allowed to stabilize the Fermi-ball is underway. we obtain, $`M_F3(4\pi \mathrm{\Sigma })^3\left({\displaystyle \frac{2\alpha }{Cm_e^4}}\right)^2=6.0\times 10^{21}C^2\left({\displaystyle \frac{\kappa }{\mathrm{GeV}}}\right)^9\mathrm{GeV}.`$ (39) In order to prevent the Fermi-ball from forming a black hole, its radius should be larger than the Schwarzschild radius, $`M_F/M_{pl}^2`$, where $`M_{pl}`$ is the Planck mass. From this condition, we obtain the following constraint of $`\kappa `$ using Eq.(39), $`\kappa 3.3\times 10^4C^{\frac{1}{6}}\mathrm{GeV}.`$ (40) We draw lines which represent Eq.(39) and Eq.(40) in Figure7(b). We find that the TA and OA experiments are powerful to search the new region of approximately $`10^{17}\mathrm{GeV}M_F10^{29}\mathrm{GeV}`$ in Morris’s Model (taking $`C=1`$). ## IV Conclusion Quantum field theory allows the existence of such nontopological solitons as Q-balls and Fermi-balls, the stability of which is supported by the conservation of a global U(1) charge. These solitons may play important roles in cosmology in solving the problems of dark matter, baryogenesis, and gamma ray bursts. In this paper, we have considered Q-balls and Fermi-balls, both of which are typical nontopological solitons. We have examined the parameter regions, masses, fluxes, charges, and energy scales of Q-balls and Fermi-balls that are to be excluded. These exclusions were derived by analyzing the various existing and future searches for monopoles, nuclearites, and cosmic rays. We also analysed the existing results and analyses of the Q-ball searches of Gyrlyanda and MACRO. The experiments considered here include: (1) underground searches with Gyrlyanda, BAKSAN, Kamiokande, Super-Kamiokande, MACRO, OYA, MICA, and AMANDA; (2) searches on the earth’s surface with NORIKURA, KITAMI, KEK, AKENO, UCSD$`II`$, and TA; (3) space experiments with SKYLAB, AMS, and OA. We summarized these experimental data and obtained bounds on the mass and the flux of Q-balls and Fermi-balls. Of course, the precise estimation of bounds should be more carefully made by those who did, or will do, the experiments. We believe, however, that our estimations will give useful information in the research of Q-balls and Fermi-balls. We first investigated Q-balls with an electric charge $`Z_Q=0`$ (SENS) which can be detected through proton-decay like process as proposed in Ref.. We found that a considerably large region $`Q10^{25}`$ has already been excluded in $`QM_S`$ plane for $`M_S<100\mathrm{GeV}`$ only by existing experiments (see Figure1(b)), and that a wider region $`Q10^{35}`$ could be searched in future experiments of TA and OA. We also found that the region $`Q=B10^{22}`$ has been excluded for B-balls for any value of $`M_S`$. We next investigated Q-balls with $`Z_Q=1,2,3,10,\mathrm{and}137`$ (SECS) which interact with matter, mainly by electromagnetic force, just like nuclearites (though the relations between the radius and the mass are quite different). We saw that larger $`Z_Q`$ gives more stringent bounds on flux and mass of SECS. We found that for the value of $`M_S10^2\mathrm{GeV}`$, experimental data gives more stringent limitations on SECS global U(1) charge than the stability condition of B-balls. Q-balls with $`Q10^{2226}`$ still remain to be examined (see Figure2-6). We finally investigated Fermi-balls with electric charge $`Z_F137`$, which are expected to be rather stable against perturbative deformation and fragmentation. We obtained bounds on mass of Fermi-balls, $`M_F10^8`$ GeV and $`M_F10^{29}`$ GeV for $`\kappa 0.1`$ GeV in $`M_F\kappa `$ plane. If we further assume the Morris Model, we obtained more stringent constraints. We also noted the importance of future TA and OA experiments in the search for Fermi-balls as seen in Figure7(b). We lastly note that the charged solitons with a relatively light mass, $`M10^8`$GeV, on which we have not focussed our attention, have astrophysical difficulties. If we consider that the dark matter consists mainly of such solitons with large cross section with matter and radiation, we will be faced by difficulties including ‘the Galaxy-halo in-fall’, ‘too much heating of disk molecules’, and ‘too small density fluctuations in the early universe’. Therefore, We must show an interest in the window for the heavier charged solitons. From these experimental bounds, we comprehensively obtained the stringent limitations on the properties of Q-balls and the Fermi-balls and then noted the possible importance of future experiments of TA and OA. These bounds will help us study the unsolved cosmological problems that we have mentioned in this paper, by developing more realistic scenarios. In these cases, Q-balls and/or Fermi-balls play an essential role. Acknowledgement: We thank H.Tawara, M.Fukugita, M.Fukushima, and M.Sasaki for useful comments.
warning/0005/cond-mat0005463.html	ar5iv	text	# Trace and antitrace maps for aperiodic sequences, their extensions and applications ## I Introduction The trace-map technique, first introduced in 1983, has proven to be a powerful tool to investigate the electronic spectrum of various aperiodic systems, such as the Fibonacci sequence (FS), the Thue-Morse sequence (TMS), and the period-doubling sequence. It has also been applied to investigate other physical systems, for instance, kicked two-level systems, and classical and quantum spin systems. The technique was extended to study aperiodic systems in combination with the real-space renormalization-group technique. Recently, trace maps have been used to evaluate localization properties in a FS tight-binding model. This technique was transferred to the field of optics in order to see the scaling of the light transmission coefficient through a Fibonacci dielectric multilayer. Now, we consider light that is vertically transmitted through a Fibonacci multilayer of two materials $`a`$ and $`b`$ which is sandwiched by two media of type $`a`$. The FS is constructed by the substitution rule $`ba`$, $`aab`$. The corresponding transfer matrices $`A_l`$ are written as $`A_1`$ $`=`$ $`P_{ab}P_bP_{ba},`$ (1) $`A_2`$ $`=`$ $`P_a,`$ (2) $`A_{l+1}`$ $`=`$ $`A_lA_{l1},`$ (3) where $`P_{ab}(P_{ba})`$ stands for the propagation matrix from layer $`a`$ to $`b`$ ($`b`$ to $`a`$) and $`P_a`$ is the propagation matrix through the single layer $`a`$. They are given by $`P_{ab}`$ $`=`$ $`P_{ba}^1=\left(\begin{array}{cc}1& 0\\ 0& n_a/n_b\end{array}\right),`$ (4) $`P_a`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{cos}\delta _a& \mathrm{sin}\delta _a\\ \mathrm{sin}\delta _a& \mathrm{cos}\delta _a\end{array}\right),`$ (5) where $`\delta _a=kn_ad_a`$, $`n_a`$ is the refraction index of material $`a`$, $`d_a`$ denotes the thickness of the layers, and $`k`$ is the wave number in vacuum. The quantity $`\delta _a`$ is the phase difference between the ends of a layer. For material $`b`$, the quantities $`P_b`$, $`\delta _b`$, $`n_b`$, and $`d_b`$ are defined analogously. The transmission coefficient is given by $$t_l=\frac{4}{\|A_l\|^2+2},$$ (6) where $`\|A_l\|^2`$ is the sum of squares of the four elements of $`A_l`$. Since the transfer matrix is unimodular, we can express the transmission coefficient as $$t_l=\frac{4}{x_l^2+y_l^2},$$ (7) where $`x_l`$ and $`y_l`$ denote the trace and antitrace of the transfer matrix $`A_l`$, respectively. Here, the so-called “antitrace” of a $`2\times 2`$ matrix $$A=\left(\begin{array}{cc}A_{11}& A_{12}\\ A_{21}& A_{22}\end{array}\right)$$ (8) is defined as $`y_A=A_{21}A_{12}`$, which follows the notion of Ref.. From Eq. (7), we see that the transmission coefficient is completely determined by the trace and the antitrace, i.e., a complete description of the light transmission through general aperiodic multilayers requires both the trace and the antitrace map. Now, we consider a different system, namely a harmonic chain composed of two kinds of masses, $`m_a`$ and $`m_b`$, which are arranged according to the FS, and are coupled by two kinds of springs, $`K_{aa}`$ and $`K_{ab}=K_{ba}`$. Making use of the transfer-matrix formalism, the equation of motion is $`\left({\displaystyle \genfrac{}{}{0pt}{}{u_{n+1}}{u_n}}\right)`$ $`=`$ $`\left(\begin{array}{cc}\frac{a_n}{K_{n,n+1}}& \frac{K_{n,n1}}{K_{n,n+1}}\\ 1& 0\end{array}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{u_n}{u_{n1}}}\right)`$ (9) $`=`$ $`P_n\left({\displaystyle \genfrac{}{}{0pt}{}{u_n}{u_{n1}}}\right),`$ (10) where $`u_n`$ is the displacement of the $`n`$th atom from its equilibrium position, $`K_{n,n\pm 1}`$ denotes the strength of the harmonic coupling between neighboring atoms, $`a_n=K_{n,n1}+K_{n,n+1}m_n\omega ^2`$, and $`\omega `$ is the vibration frequency. From the corresponding global transfer matrix $`A_l=_{n=N_l}^1P_n`$, the transmission coefficient $`t_l`$ and the Lyapunov exponent $`\mathrm{\Gamma }_l`$ are given by $`t_l`$ $`=`$ $`{\displaystyle \frac{4\mathrm{sin}^2k}{(z_l\mathrm{cos}ky_l)^2+x_l^2\mathrm{sin}^2k}},`$ (11) $`\mathrm{\Gamma }_l`$ $`=`$ $`{\displaystyle \frac{1}{N_l}}\mathrm{ln}(\|A_l\|^2)={\displaystyle \frac{1}{N_l}}\mathrm{ln}(x_l^2+y_l^22),`$ (12) where $`z_l=(A_l)_{11}(A_l)_{22}`$, $`N_l`$ denotes the number of atoms in the chain, and $`\mathrm{cos}k=m_a/(2K_{aa})`$. As our third example of physical systems based on aperiodic substitution sequences, we consider the transmission in electronic systems. This is closely related to the harmonic chain considered above. The Schrödinger equation for a one-dimensional tight-binding model with nearest-neighbor hopping can be written in matrix form as follows $`\left({\displaystyle \genfrac{}{}{0pt}{}{\varphi _{n+1}}{\varphi _n}}\right)`$ $`=`$ $`\left(\begin{array}{cc}\frac{Eϵ_n}{t_{n,n+1}}& \frac{t_{n,n1}}{t_{n,n+1}}\\ 1& 0\end{array}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{\varphi _n}{\varphi _{n1}}}\right)`$ (13) $`=`$ $`M_n\left({\displaystyle \genfrac{}{}{0pt}{}{\varphi _n}{\varphi _{n1}}}\right),`$ (14) where $`\varphi _n`$ denotes the amplitude of the wave function in the Wannier representation, $`E`$ the corresponding energy, $`ϵ_n`$ the on-site energy at site $`n`$, and $`t_{n,n\pm 1}`$ the hopping matrix element between two neighboring sites. $`M_n`$ is the local transfer matrix associated with site $`n`$. The transmission coefficient is given by $`t_l`$ $`=`$ $`{\displaystyle \frac{4E^2}{(z_lE/2y_l)^2+x_l^2(1E^2/4)}},`$ (15) where the quantities $`x_l`$, $`y_l`$, and $`z_l`$ are again related to the global transfer matrix of the chain, i.e., the product $`A_l=_{n=N_l}^1M_n`$ of the local transfer matrices along the chain. Note the similarity to Eq. (11). The corresponding Lyapunov exponent $`\mathrm{\Gamma }_l`$ is given by the same expression (12). For the latter two systems, the Lyapunov exponent is completely determined by the trace and the antitrace; however, we need to know $`z_l`$ to calculate the transmission coefficient. Fortunately, it turns out that the maps for $`z_l`$ and $`y_l`$ are the same, as will be shown in Sec. IV. Therefore, the trace and antitrace map are sufficient to determine the transmission coefficient and the Lyapunov exponent. Thus it is desirable to construct antitrace maps for various aperiodic sequences, which is the motivation of the work presented here. The paper is organized as follows. In Sec. II, we give the antitrace maps for various classes of aperiodic sequences, including the FS, the TMS, the periodic-doubling sequence, and certain generalizations. The extension to arbitrary substitution rules and to maps for matrix elements are discussed in Sec. III and in Sec. IV, respectively. It is shown that the antitrace maps and the maps for matrix elements exist for arbitrary substitution rules and that the maps for non-diagonal elements and for the difference of the diagonal elements coincide with the antitrace maps. Applications to the computation of transmission coefficients and Lyapunov exponents in different aperiodic systems are investigated in Sec. V. Finally, in Sec. VI, we conclude. ## II Trace and antitrace maps for two-letter sequences We now proceed with the derivation of the antitrace maps of various classes of aperiodic sequences. We also include the corresponding trace maps for reasons which will become clear later. In this part, we make ample use of several relations for unimodular matrices. Therefore, we append a compilation of these relations in Appendix A. ### A Generalized Fibonacci sequences There are many kinds of generalized FSs. Here, we study two-letter sequences $`\text{FS}(m,n)`$ that can be generated by the inflation scheme $$S_0=b,S_1=a,S_{l+1}=S_l^mS_{l1}^n$$ (16) with arbitrary positive integers $`m`$ and $`n`$, where $`\text{FS}(1,1)`$ corresponds to the well-known standard FS. Equivalently, they can also be generated by the substitution rule $$ba,aa^mb^n.$$ (17) The total number of letters $`a`$ and $`b`$ in the word $`S_l`$ is denoted by $`F_l`$ and satisfies the recursion relation $$F_{l+1}=mF_l+nF_{l1},F_0=F_1=1.$$ (18) In the limit of an infinite sequence, the ratio of word lengths for subsequent inflation steps is given by $$\sigma =\underset{l\mathrm{}}{lim}\frac{F_{l+1}}{F_l}=\frac{m+\sqrt{m^2+4n}}{2}.$$ (19) Some values of $`\sigma `$ and commonly used terms for special cases of so-called “metallic means” are $`\text{FS}(1,1):`$ $`\sigma _g={\displaystyle \frac{1+\sqrt{5}}{2}}`$golden mean, $`\text{FS}(2,1):`$ $`\sigma _s=1+\sqrt{2}`$silver mean, $`\text{FS}(3,1):`$ $`\sigma _b={\displaystyle \frac{3+\sqrt{13}}{2}}`$bronze mean, $`\text{FS}(1,2):`$ $`\sigma _c=2`$copper mean, $`\text{FS}(1,3):`$ $`\sigma _n={\displaystyle \frac{1+\sqrt{13}}{2}}\text{nickel mean}.`$ It is known that the sequences $`\text{FS}(m,n)`$ with $`n=1`$ are quasiperiodic and those with $`n2`$ are always aperiodic. It is interesting to consider two further classes of generalized FSs $`b`$ $``$ $`b^{m1}a,ab^{m1}ab,`$ (20) $`b`$ $``$ $`b^{m2}a,ab^{m2}ab^{m2}ab.`$ (21) The first class (20) consists of the so-called Fibonacci-class sequences $`\text{FC}(m)`$, the second (21) occurs in the renormalization-group analysis of the energy spectrum of $`\text{FC}(m)`$ chains. It is easy to check that the inflation schemes of these two generalized FSs are the same as those for $`\text{FS}(m,1)`$, but they differ in the initial words. A natural further generalization of these sequences is given by $`bb^{mk}a,a(b^{mk}a)^kb,`$ (22) which we denote as $`\text{FC}(m,k)`$. Here, $`\text{FC}(m,1)`$ and $`\text{FC}(m,2)`$ correspond to the cases (20) and (21). The corresponding inflation scheme is $`S_0=b,S_1=b^{mk}a,S_{l+1}=S_l^mS_{l1},`$ (23) which is the same as that of $`\text{FS}(m,1)`$ apart from the different second initial word. #### 1 The Fibonacci sequence Let us commence with the simplest example $`\text{FS}(1,1)`$. We consider the case that the two letters $`a`$ and $`b`$ correspond to two basic unimodular transfer matrices $`A`$ and $`B`$, respectively. Denoting by $`A_l`$ the total transfer matrix corresponding to a word $`S_l`$, the matrix equivalent of Eq. (16) for FS(1,1) is $$A_{l+1}=A_{l1}A_l,$$ (24) where $`A_1=A`$ and $`A_0=B`$ are the transfer matrices of the two building blocks $`a`$ and $`b`$. Note the reversed order of matrix multiplication as compared to the concatenation of letters in Eq. (16), which occurs in the related tight-binding model that is usually considered, whereas the order of matrix multiplications is not reversed in the optical problem, compare Eq. (3). The well-known trace map reads $$x_{l+1}=x_{l1}x_lx_{l2}.$$ (25) Note that in part of the literature a factor $`1/2`$ is introduced in the definition of $`x_l`$. Here, we omitted this factor to keep symmetry between trace and antitrace. From Eq. (A.9), we obtain the antitrace map $`y_{l+1}=x_ly_{l1}+y_{l2}.`$ (26) The coefficients of the trace map are constants; however, those of the antitrace map include the traces. So, if we want to derive the antitrace map, the trace map must also be known. This is why we have to consider trace and antitrace maps at the same time. #### 2 Generalized Fibonacci sequences $`\mathrm{FS}(m,n)`$ For FS$`(m,n)`$, Eq. (16), the recursion relation for the transfer matrix is given by $`A_{l+1}`$ $`=`$ $`A_{l1}^nA_l^m`$ (27) $`=`$ $`\left(U_n^{(l1)}A_{l1}U_{n1}^{(l1)}I\right)`$ (29) $`\times \left(U_m^{(l1)}A_lU_{m1}^{(l1)}I\right).`$ Here, we used Eq. (A.1) and the definition of the functions $`U_n(x)=C_{n1}(x/2)`$ given in Appendix A in terms of the Chebyshev polynomials of the second kind $`C_n(x)`$. Furthermore, we introduced the notation $$U_n^{(l)}=U_n(x_{A_l}).$$ (30) From Eqs. (29), (A.8), and (A.9), the trace and the antitrace maps are obtained as $`x_{l+1}`$ $`=`$ $`U_n^{(l1)}U_m^{(l)}v_lU_{n1}^{(l1)}U_{m+1}^{(l)}`$ (32) $`U_{n+1}^{(l1)}U_{m1}^{(l)},`$ $`v_{l+1}`$ $`=`$ $`U_n^{(l1)}U_{m+1}^{(l)}v_lU_{n1}^{(l1)}U_{m+2}^{(l)}`$ (34) $`U_{n+1}^{(l1)}U_m^{(l)},`$ $`y_{l+1}`$ $`=`$ $`U_n^{(l1)}\left(U_m^{(l)}w_lU_{m1}^{(l)}y_{l1}\right)`$ (36) $`U_{n1}^{(l1)}U_m^{(l)}y_l,`$ $`w_{l+1}`$ $`=`$ $`U_n^{(l1)}\left(U_{m1}^{(l)}w_lU_{m2}^{(l)}y_{l1}\right)`$ (38) $`+\left(x_{l+1}U_{n1}^{(l1)}U_{m1}^{(l)}\right)y_l,`$ where $`v_l=x_{A_{l1}A_l}`$ and $`w_l=y_{A_{l1}A_l}`$. Note that the roles of $`v_l`$ and $`w_l`$ are subsidiary. Eqs. (32) and (34) constitute the trace map, Eqs. (36) and (38) give the corresponding antitrace map. For special cases, these expressions simplify considerably. For $`\text{FS}(1,n)`$, we obtain, using the properties (A.8) of the functions $`U_n(x)`$, $`x_{l+1}`$ $`=`$ $`U_n^{(l1)}v_lU_{n1}^{(l1)}x_l,`$ (39) $`v_{l+1}`$ $`=`$ $`U_n^{(l1)}x_lv_lU_{n1}^{(l1)}\left(x_l^21\right)`$ (41) $`U_{n+1}^{(l1)},`$ $`y_{l+1}`$ $`=`$ $`U_n^{(l1)}w_lU_{n1}^{(l1)}y_l,`$ (42) $`w_{l+1}`$ $`=`$ $`x_{l+1}y_l+U_n^{(l1)}y_{l1}.`$ (43) Similarly, for $`\text{FS}(m,1)`$, we find $`x_{l+1}`$ $`=`$ $`U_m^{(l)}v_lU_{m1}^{(l)}x_{l1},`$ (44) $`v_{l+1}`$ $`=`$ $`U_{m+1}^{(l)}v_lU_m^{(l)}x_{l1},`$ (45) $`y_{l+1}`$ $`=`$ $`U_m^{(l)}w_lU_{m1}^{(l)}y_{l1},`$ (46) $`w_{l+1}`$ $`=`$ $`x_{l+1}y_l+U_{m1}^{(l)}w_lU_{m2}^{(l)}y_{l1}.`$ (47) Eqs. (44)–(47) are quite different from Eqs. (39)–(43) above. The corresponding aperiodic sequences show rather different physical properties. We also point out that the trace and antitrace maps for the sequences $`\text{FC}(m,k)`$ are given by Eqs. (44)–(47) since they have the same inflation scheme. Eliminating the subsidiary variables $`v_l`$ and $`w_l`$ in Eqs. (32)–(38) for the general case $`\text{FS}(m,n)`$, we obtain $`x_{l+1}`$ $`=`$ $`{\displaystyle \frac{U_m^{(l)}U_n^{(l1)}}{U_m^{(l1)}}}\left(U_{m+1}^{(l1)}x_lU_{n+1}^{(l2)}+U_{n1}^{(l2)}\right)`$ (49) $`U_{m+1}^{(l)}U_{n1}^{(l1)}U_{m1}^{(l)}U_{n+1}^{(l1)},`$ $`y_{l+1}`$ $`=`$ $`{\displaystyle \frac{U_m^{(l)}U_n^{(l1)}}{U_m^{(l1)}}}\left(U_n^{(l2)}y_{l2}+U_{m1}^{(l1)}y_l\right)`$ (51) $`+U_{m+1}^{(l)}U_n^{(l1)}y_{l1}U_m^{(l)}U_{n1}^{(l1)}y_l.`$ Here, we used Eq. (A.8) to simplify the result. The above two equations are alternative forms of the trace and the antitrace map. Again, for the special cases $`m=1`$ or $`n=1`$, these equations simplify. For $`\text{FS}(1,n)`$, we find $`x_{l+1}`$ $`=`$ $`U_n^{(l1)}\left(U_{n1}^{(l2)}U_{n+1}^{(l2)}\right)+U_{n+1}^{(l1)}x_l,`$ (52) $`y_{l+1}`$ $`=`$ $`U_n^{(l1)}\left(U_n^{(l2)}y_{l2}+x_ly_{l1}\right)U_{n1}^{(l1)}y_l.`$ (54) The result for $`\text{FS}(m,1)`$ reads $`x_{l+1}`$ $`=`$ $`{\displaystyle \frac{U_m^{(l)}}{U_m^{(l1)}}}\left(U_{m+1}^{(l1)}x_lx_{l2}\right)U_{m1}^{(l)}x_{l1},`$ (56) $`y_{l+1}`$ $`=`$ $`{\displaystyle \frac{U_m^{(l)}}{U_m^{(l1)}}}\left(y_{l2}+U_{m1}^{(l1)}y_l\right)+U_{m+1}^{(l)}y_{l1}.`$ (58) For $`\text{FS}(1,1)`$, Eqs. (49)–(LABEL:eq:yyln1) reduce to Eqs. (25) and (26), as expected. For some other special values of $`m`$ and $`n`$, the trace and antitrace maps of $`\text{FS}(m,n)`$ are given in Appendix B. ### B Generalized Thue-Morse sequences Another type of aperiodic sequence is the celebrated TMS and its generalizations. Here, we consider generalized sequences $`\text{TMS}(m,n)`$ with inflation scheme $$bb^ma^n,aa^nb^m.$$ (60) Equivalently, $`\text{TMS}(m,n)`$ can be constructed as $`S_0`$ $`=`$ $`b,\stackrel{~}{S}_0=a,`$ (61) $`S_{l+1}`$ $`=`$ $`S_l^m\stackrel{~}{S}_l^n,\stackrel{~}{S}_{l+1}=\stackrel{~}{S}_l^nS_l^m.`$ (62) For $`m=n=1`$, this reduces to the standard TMS. The recursion relation for transfer matrices of $`\text{TMS}(m,n)`$ reads $$A_{l+1}=B_l^nA_l^m,B_{l+1}=A_l^mB_l^n$$ (63) where $`A_0`$ is the matrix corresponding to the building block $`b`$, and $`B_0`$ corresponds to $`a`$, respectively. Using the same method as above, we get $`x_{l+1}`$ $`=`$ $`U_n^{(l)}U_m^{(l)}v_lU_{n1}^{(l)}U_{m+1}^{(l)}U_{n+1}^{(l)}U_{m1}^{(l)},`$ (64) $`v_{l+1}`$ $`=`$ $`U_{2n}^{(l)}U_{2m}^{(l)}v_lU_{2n1}^{(l)}U_{2m+1}^{(l)}U_{2n+1}^{(l)}U_{2m1}^{(l)},`$ (66) where $`v_l=x_{B_lA_l}`$. These two equations determine the trace map. It is somewhat more complicated to derive the antitrace map because $`y_{A_l}y_{B_l}`$. We define $`y_l=y_{A_l}`$ and $`\stackrel{~}{y}_l=y_{B_l}`$. Then, from Eqs. (63),(A.1), (A.21), and (A.25), we have $`y_{l+1}`$ $`=`$ $`U_m^{(l)}\left(U_n^{(l)}w_lU_{n1}^{(l)}y_l\right)U_n^{(l)}U_{m1}^{(l)}\stackrel{~}{y}_l,`$ (68) $`\stackrel{~}{y}_{l+1}`$ $`=`$ $`U_m^{(l)}\left(U_n^{(l)}\stackrel{~}{w}_lU_{n1}^{(l)}y_l\right)U_n^{(l)}U_{m1}^{(l)}\stackrel{~}{y}_l,`$ (70) $`w_{l+1}`$ $`=`$ $`(U_{2n}^{(l)}U_m^{(l)}v_lU_{2n1}^{(l)}U_{m+1}^{(l)}`$ (73) $`U_{2n+1}^{(l)}U_{m1}^{(l)})U_m^{(l)}y_l+U_{2n}^{(l)}\stackrel{~}{y}_l,`$ $`\stackrel{~}{w}_{l+1}`$ $`=`$ $`(U_n^{(l)}U_{2m}^{(l)}v_lU_{n1}^{(l)}U_{2m+1}^{(l)}`$ (75) $`U_{n+1}^{(l)}U_{2m1}^{(l)})U_n^{(l)}\stackrel{~}{y}_l+U_{2m}^{(l)}y_l.`$ Here, $`w_l=y_{B_lA_l}`$ and $`\stackrel{~}{w}_l=y_{A_lB_l}`$. The antitrace map is completely determined by Eqs. (LABEL:eq:mnx)–(75). For $`n=1`$ and $`m=1`$, Eqs. (LABEL:eq:mnx)–(75) reduce to $`x_{l+1}`$ $`=`$ $`v_l,`$ (76) $`v_{l+1}`$ $`=`$ $`x_l^2(v_l2)+2,`$ (77) $`y_{l+1}`$ $`=`$ $`w_l,`$ (78) $`\stackrel{~}{y}_{l+1}`$ $`=`$ $`\stackrel{~}{w}_l,`$ (79) $`w_{l+1}`$ $`=`$ $`x_l[(x_{l+1}1)y_l+\stackrel{~}{y}_l],`$ (80) $`\stackrel{~}{w}_{l+1}`$ $`=`$ $`x_l[(x_{l+1}1)\stackrel{~}{y}_l+y_l].`$ (81) This yields the well-known trace map of the TMS $$x_{l+1}=x_{l1}^2(x_l2)+2,$$ (82) and the antitrace map $`y_{l+1}`$ $`=`$ $`x_{l1}\left[(x_l1)y_{l1}+\stackrel{~}{y}_{l1}\right],`$ (83) $`\stackrel{~}{y}_{l+1}`$ $`=`$ $`x_{l1}\left[(x_l1)\stackrel{~}{y}_{l1}+y_{l1}\right].`$ (84) The above two equations give $`y_{l+1}`$ $`=`$ $`x_{l1}[(x_l+x_{l2}2)y_{l1}`$ (86) $`+x_{l3}x_{l2}(2x_{l2})y_{l3}]`$ which is an alternative form of the antitrace map. From Eqs. (LABEL:eq:mnx)–(75), we can solve for the subsidiary quantities $`v_l`$, $`w_l`$, and $`\stackrel{~}{w}_l`$, for instance, $`v_{l+1}`$ $`=`$ $`{\displaystyle \frac{U_{2n}^{(l)}U_{2m}^{(l)}}{U_n^{(l)}U_m^{(l)}}}[x_l+U_{n1}^{(l)}(U_m^{(l)}x_lU_{m1}^{(l)})`$ (90) $`+U_{m1}^{(l)}(U_n^{(l)}x_lU_{n1}^{(l)})]`$ $`U_{2n1}^{(l)}\left(U_{2m}^{(l)}x_lU_{2m1}^{(l)}\right)`$ $`U_{2m1}^{(l)}\left(U_{2n}^{(l)}x_lU_{2n1}^{(l)}\right).`$ The combination of Eqs. (LABEL:eq:mnx) and (90) gives an alternative form of the trace map of TMS($`m,n`$). ### C Period-doubling sequence The period-doubling sequence can be generated by the substitution rule $`b`$ $``$ $`ba,ab^2,`$ (91) or the inflation scheme $`S_0=b,S_1=ba,S_{l+1}=S_lS_{l1}^2.`$ (92) The inflation scheme is the same as that of $`\text{FS}(1,2)`$. Therefore, from Eqs. (49)–(51), the trace and the antitrace maps are obtained as $`x_{l+1}`$ $`=`$ $`x_{l1}(x_lx_{l1}x_{l2}^2+2)x_l,`$ (93) $`y_{l+1}`$ $`=`$ $`x_{l1}x_{l2}y_{l2}+x_lx_{l1}y_{l1}y_l.`$ (94) This yields also the trace and antitrace map for the copper mean sequence $`\text{FS}(1,2)`$. ## III Arbitrary substitution sequences As there exist trace maps for arbitrary substitution sequences, one natural question is whether antitrace maps also exist for arbitrary sequences. The answer is affirmative. We commence our argument in analogy with the discussion in Ref. and restrict ourselves to the case of unimodular matrices. Let $`A_1,A_2,\mathrm{},A_r`$ be $`2\times 2`$ matrices and define the following $`2^r`$ matrices $$B_{ϵ_1ϵ_2\mathrm{}ϵ_r}=A_1^{ϵ_1}A_2^{ϵ_2}\mathrm{}A_r^{ϵ_r}$$ (95) where $`ϵ_j\{0,1\}`$ for $`1jr`$. Then, from Eq. (A.14), any monomial $`A_{j_1}A_{j_2}\mathrm{}A_{j_s}`$, with $`1j_ir`$ and $`1is`$, can be written as a linear combination of the matrices $`B_{ϵ_1ϵ_2\mathrm{}ϵ_r}`$, namely, $`A_{j_1}A_{j_2}\mathrm{}A_{j_s}`$ (96) $`=`$ $`{\displaystyle \underset{ϵ_1=0}{\overset{1}{}}}{\displaystyle \underset{ϵ_2=0}{\overset{1}{}}}\mathrm{}{\displaystyle \underset{ϵ_r=0}{\overset{1}{}}}c_{ϵ_1ϵ_2\mathrm{}ϵ_r}B_{ϵ_1ϵ_2\mathrm{}ϵ_r},`$ (97) where each coefficient is a polynomial in the traces $`x_{A_j}`$, $`1jr`$, and the traces $`x_{A_jA_k}`$, $`1j<kr`$. This result not only yields the trace map, but also gives the antitrace map for any substitution sequence. We define $$B_{ϵ_1ϵ_2\mathrm{}ϵ_r,l}=A_{1l}^{ϵ_1}A_{2l}^{ϵ_2}\mathrm{}A_{rl}^{ϵ_r}$$ (98) with $`ϵ_j\{0,1\}`$, $`1jr`$, and $`l0`$, where $`A_{jl}`$ is the unimodular $`2\times 2`$ matrix associated to the $`l`$-th iterate of the $`j`$-th letter. Since each matrix in $`B_{ϵ_1ϵ_2\mathrm{}ϵ_r,l+1}`$ is, by definition, a monomial in the matrices $`A_{jl}`$, they can be expanded in terms of the matrices $`B_{ϵ_1ϵ_2\mathrm{}ϵ_r,l}`$ according to (97). Then the trace of $`B_{ϵ_1ϵ_2\mathrm{}ϵ_r,l+1}`$ is a polynomial in the $`2^r1`$ traces of $`B_{ϵ_1ϵ_2\mathrm{}ϵ_r,l}`$; and the antitrace of $`B_{ϵ_1ϵ_2\mathrm{}ϵ_r,l+1}`$ is a polynomial in the $`2^r1`$ antitraces of $`B_{ϵ_1ϵ_2\mathrm{}ϵ_r,l}`$. Therefore, we conclude that both the trace and antitrace maps exist for arbitrary substitution sequences, and the dimension of the antitrace map is $`2^r1`$. Next, we present a concrete example to illustrate this conclusion. The Rudin-Shapiro sequence can be defined by means of a substitution rule on four letters. The substitution rule is $$aac,bdc,cab,ddb,$$ (99) and the corresponding matrix recursion relations are $`A_{l+1}`$ $`=`$ $`C_lA_l,B_{l+1}=C_lD_l,`$ (100) $`C_{l+1}`$ $`=`$ $`B_lA_l,D_{l+1}=B_lD_l.`$ (101) We have the useful relation $$D_l=C_lA_l^1B_l,$$ (102) which effectively reduces the sequence to three basic letters. Now, we choose the seven matrices $`A_l`$, $`B_l`$, $`C_l`$, $`D_l`$, $`A_lC_l`$, $`A_lB_l`$, and $`B_lC_l`$ as our basic set of matrices. In what follows, we denote the traces and antitraces by $`a_l=x_{A_l},b_l=x_{B_l},c_l=x_{C_l},d_l=x_{D_l},`$ (103) $`e_l=x_{A_lC_l},f_l=x_{A_lB_l},g_l=x_{B_lC_l},`$ (104) $`\stackrel{~}{a}_l=y_{A_l},\stackrel{~}{b}_l=y_{B_l},\stackrel{~}{c}_l=y_{C_l},\stackrel{~}{d}_l=y_{D_l},`$ (105) $`\stackrel{~}{e}_l=y_{A_lC_l},\stackrel{~}{f}_l=y_{A_lB_l},\stackrel{~}{g}_l=y_{B_lC_l}.`$ (106) By using Eqs. (A.14) and (A.21), we obtain $`A_{l+1}`$ $`=`$ $`C_lA_l,`$ (107) $`B_{l+1}`$ $`=`$ $`c_lD_la_lB_l+A_lB_l,`$ (108) $`C_{l+1}`$ $`=`$ $`B_lA_l,`$ (109) $`D_{l+1}`$ $`=`$ $`(a_lg_lc_lf_l)B_lc_lA_l+b_lD_l`$ (112) $`(g_lb_lc_l)A_lB_l`$ $`+(f_la_lb_l)C_lB_l+C_lA_l,`$ $`A_{l+1}C_{l+1}`$ $`=`$ $`f_lC_lA_lb_lC_l+C_lB_l,`$ (113) $`A_{l+1}B_{l+1}`$ $`=`$ $`c_l[(1a_l^2)B_l+e_lD_l+a_lA_lB_l]`$ (115) $`C_lB_l,`$ $`B_{l+1}C_{l+1}`$ $`=`$ $`b_lc_l[(f_la_la_l^2b_l+b_l)C_l`$ (120) $`+a_lD_lC_lB_l(f_la_lb_l)C_lA_l]`$ $`b_l[(f_la_lb_l)(a_lIA_l)`$ $`+b_lI+(a_l^21)B_la_lA_lB_l]`$ $`c_lC_l+I.`$ Note that the order of multiplication of two matrices on the right-hand sides of these equations may differ from the order of our basic matrix products $`A_lC_l`$, $`A_lB_l`$, or $`B_lC_l`$. We can use Eq. (A.14) to reverse the order to obtain a systems of equations that closes with our seven basic matrices. Now, from Eq. (120) and (A.9), the trace and antitrace maps are obtained as $`a_{l+1}`$ $`=`$ $`e_l,`$ (121) $`b_{l+1}`$ $`=`$ $`c_ld_la_lb_l+f_l,`$ (122) $`c_{l+1}`$ $`=`$ $`f_l,`$ (123) $`d_{l+1}`$ $`=`$ $`b_ld_la_lc_l+e_l,`$ (124) $`e_{l+1}`$ $`=`$ $`e_lf_lb_lc_l+g_l,`$ (125) $`f_{l+1}`$ $`=`$ $`c_l(d_le_la_l^2b_l+a_lf_l+b_l)g_l,`$ (126) $`g_{l+1}`$ $`=`$ $`b_lc_l(a_lc_lf_le_lf_la_l^2b_lc_l+a_lb_le_l+b_lc_l`$ (128) $`+a_ld_l+g_l)b_l^2c_l^2+2,`$ $`\stackrel{~}{a}_{l+1}`$ $`=`$ $`c_l\stackrel{~}{a}_l+a_l\stackrel{~}{c}_l\stackrel{~}{e}_l,`$ (129) $`\stackrel{~}{b}_{l+1}`$ $`=`$ $`c_l\stackrel{~}{d}_la_l\stackrel{~}{b}_l+\stackrel{~}{f}_l,`$ (130) $`\stackrel{~}{c}_{l+1}`$ $`=`$ $`b_l\stackrel{~}{a}_l+a_l\stackrel{~}{b}_l\stackrel{~}{f}_l,`$ (131) $`\stackrel{~}{d}_{l+1}`$ $`=`$ $`(a_lg_la_lb_lc_l)\stackrel{~}{b}_l+b_l\stackrel{~}{d}_l(g_lb_lc_l)\stackrel{~}{f}_l`$ (133) $`+(f_la_lb_l)(b_l\stackrel{~}{c}_l\stackrel{~}{g}_l)+a_l\stackrel{~}{c}_l\stackrel{~}{e}_l,`$ $`\stackrel{~}{e}_{l+1}`$ $`=`$ $`f_l(c_l\stackrel{~}{a}_l+a_l\stackrel{~}{c}_l\stackrel{~}{e}_l)+c_l\stackrel{~}{b}_l\stackrel{~}{g}_l,`$ (134) $`\stackrel{~}{f}_{l+1}`$ $`=`$ $`c_l(a_l^2\stackrel{~}{b}_l+e\stackrel{~}{d}_l+a_l\stackrel{~}{f}_l)b_l\stackrel{~}{c}_l+\stackrel{~}{g}_l,`$ (135) $`\stackrel{~}{g}_{l+1}`$ $`=`$ $`b_l(1c_l^2)(f_la_lb_l)\stackrel{~}{a}_l+b_l(1a_l^2c_l^2)\stackrel{~}{b}_l`$ (138) $`c_l\stackrel{~}{c}_l+a_lb_lc_l\stackrel{~}{d}_l+b_lc_l(f_la_lb_l)\stackrel{~}{e}_l+a_lb_l\stackrel{~}{f}_l`$ $`+b_lc_l\stackrel{~}{g}_l.`$ Thus, we derived the trace and antitrace maps of the Rudin-Shapiro sequence. Now, we discuss the dimension of the antitrace map. Let $`A`$, $`B`$, and $`C`$ be $`2\times 2`$ matrices. Then $`ABC`$ $`=`$ $`[(x_{ABC}x_{AB}x_Cx_Ax_{BC}+x_Ax_Bx_C)I`$ (142) $`+(x_{BC}x_Bx_C)Ax_{AC}B`$ $`+(x_{AB}x_Ax_B)C`$ $`+x_CAB+x_BAC+x_ABC]/2.`$ Taking the trace on both sides of Eq. (142), we are led to a trivial identity. However, if we take the antitrace, we obtain $`y_{ABC}`$ $`=`$ $`[(x_{BC}x_Bx_C)y_Ax_{AC}y_B`$ (145) $`+(x_{AB}x_Ax_B)y_C`$ $`+x_Cy_{AB}+x_By_{AC}+x_Ay_{BC}]/2.`$ The antitrace of any monomial can be written as a linear combination of a polynomial in the antitraces $`y_{A_j}`$, $`1jr`$, and the antitraces $`y_{A_jA_k}`$, $`1j<kr`$. Each coefficient is a polynomial in the traces $`x_{A_j}`$, $`1jr`$ and the traces $`x_{A_jA_k}`$, $`1j<kr`$. From this observation we conclude that the dimension of our antitrace map is $`r(1+r)/2`$, i.e., the dimension is reduced from $`2^r1`$. Here, for the dimension of the antitrace map, we do not take into account the dimension of the trace map, which enters the coefficients of the antitrace map. Thus, the full dimension of the trace and antitrace map is given by the sum of their respective dimensions. Let us consider two ternary sequences as examples. Our first example of a three-letter substitution rule and the corresponding recursion relation for the transfer matrices is $`ab,bc,cca,`$ (146) $`A_{l+1}=B_l,B_{l+1}=C_l,C_{l+1}=A_lC_l.`$ (147) Using Eqs. (A.15) and (A.21), we obtain $`A_{l+1}`$ $`=`$ $`B_l,`$ (148) $`B_{l+1}`$ $`=`$ $`C_l,`$ (149) $`C_{l+1}`$ $`=`$ $`A_lC_l,`$ (150) $`B_{l+1}A_{l+1}`$ $`=`$ $`C_lB_l,`$ (151) $`C_{l+1}B_{l+1}`$ $`=`$ $`x_{C_l}A_lC_lA_l,`$ (152) $`A_{l+1}C_{l+1}`$ $`=`$ $`B_lA_lC_l,`$ (153) $`B_{l+1}A_{l+1}C_{l+1}`$ $`=`$ $`B_lA_l+x_{B_lA_lC_l}C_l`$ (155) $`x_{B_lA_l}I.`$ Taking the trace of the above equation, we obtain the trace map. The dimension of the trace map is $`2^31=7`$. Taking the antitrace of the sixth line of the above equation, we can expand it according to Eq. (145). Therefore, the last equality in Eq. (155) is not necessary, and the dimension of the antitrace map is $`3(3+1)/2=6`$. Our second example is the three-component FS generated by $`ab,bc,cabc,`$ (156) $`A_{l+1}=B_l,B_{l+1}=C_l,C_{l+1}=C_lB_lA_l.`$ (157) The corresponding maps for the matrices are $`A_{l+1}`$ $`=`$ $`B_l,`$ (158) $`B_{l+1}`$ $`=`$ $`C_l,`$ (159) $`C_{l+1}`$ $`=`$ $`C_lB_lA_l,`$ (160) $`B_{l+1}A_{l+1}`$ $`=`$ $`C_lB_l,`$ (161) $`C_{l+1}B_{l+1}`$ $`=`$ $`B_lA_lx_{B_lA_l}I+x_{C_lB_lA_l}C_l,`$ (162) $`C_{l+1}B_{l+1}A_{l+1}`$ $`=`$ $`A_lx_{A_l}I+x_{C_lB_lA_l}C_lB_l.`$ (163) We see that, for this particular sequence, both the trace and antitrace maps are six-dimensional. ## IV Maps for matrix elements As discussed in Sec. I, we need to know all elements of the global transfer matrix in order to compute certain physical quantities. Thus, the trace and antitrace maps may not be sufficient, and one would like to determine analogous maps for each of the matrix elements. Actually, from Eq. (97), we know that such matrix element maps exist for any substitution rule, and Eqs. (120), (155), and (163) already contain examples of matrix element maps. Now, we investigate the maps for the matrix elements of the FS($`m`$,$`n`$) and TMS($`m`$,$`n`$). Using Eqs. (A.14), (A.21), and (A.25), we obtain the matrix maps of FS($`m,n`$) as (29) and $`M_{l+1}`$ $`=`$ $`U_n^{(l1)}U_{m1}^{(l)}M_l`$ (167) $`+\left(x_{l+1}U_{n1}^{(l1)}U_{m1}^{(l)}\right)A_l`$ $`U_n^{(l1)}U_{m2}^{(l)}A_{l1}`$ $`+\left(U_{n1}^{(l1)}U_{m2}^{(l)}+v_{l+1}x_lx_{l+1}\right)I,`$ where $`M_l=A_{l1}A_l`$. The traces $`x_{l+1}`$ and $`v_{l+1}`$ appearing on the right-hand side of Eq. (167) are given in terms of $`v_l`$ via Eqs. (32) and (34), respectively. Similarly, the matrix map of TMS($`m`$,$`n`$) is obtained as $`A_{l+1}`$ $`=`$ $`U_n^{(l)}\left(U_m^{(l)}N_lU_{m1}^{(l)}B_l\right)`$ (169) $`U_{n1}^{(l)}\left(U_m^{(l)}A_lU_{m1}^{(l)}I\right),`$ $`B_{l+1}`$ $`=`$ $`U_n^{(l)}\left(U_m^{(l)}\stackrel{~}{N}_lU_{m1}^{(l)}B_l\right)`$ (171) $`U_{n1}^{(l)}\left(U_m^{(l)}A_lU_{m1}^{(l)}I\right),`$ $`N_{l+1}`$ $`=`$ $`(U_{2n}^{(l)}U_m^{(l)}v_lU_{2n1}^{(l)}U_{m+1}^{(l)}`$ (174) $`U_{2n+1}^{(l)}U_{m1}^{(l)}\left)\right(U_m^{(l)}A_lU_{m1}^{(l)}I)`$ $`+U_{2n}^{(l)}B_lU_{2n+1}^{(l)}I,`$ $`\stackrel{~}{N}_{l+1}`$ $`=`$ $`(U_n^{(l)}U_{2m}^{(l)}v_lU_{n1}^{(l)}U_{2m+1}^{(l)}`$ (177) $`U_{n+1}^{(l)}U_{2m1}^{(l)}\left)\right(U_n^{(l)}B_lU_{n1}^{(l)}I)`$ $`+U_{2m}^{(l)}A_lU_{2m+1}^{(l)}I,`$ where $`N_l=B_lA_l`$ and $`\stackrel{~}{N}_l=A_lB_l`$. We can eliminate the subsidiary matrices $`M_l`$, $`N_l`$, and $`\stackrel{~}{N}_l`$ from Eqs. (29) and (167)–(177). For example, Eq. (167) becomes $`M_l`$ $`=`$ $`{\displaystyle \frac{1}{U_m^{(l1)}}}\left(U_{m1}^{(l1)}A_l+U_n^{(l2)}A_{l2}U_{n1}^{(l2)}I\right)`$ (179) $`+v_lIx_lx_{l1}I+x_lA_{l1}.`$ Thus, we obtain another form of the matrix map of FS($`m,n`$) given by Eqs. (29) and (179). For $`m=n=1`$, Eqs. (29) and (179) reduce to $`A_{l+1}=(x_{l+1}x_lx_{l1})I+x_lA_{l1}+A_{l2}.`$ (180) This is the matrix map of the FS. For the TMS, we find from Eqs. (169)–(177) for $`m=n=1`$ $`A_{l+1}`$ $`=`$ $`x_{l1}[(x_l1)A_{l1}+B_{l1}x_{l1}I]+I,`$ (181) $`B_{l+1}`$ $`=`$ $`x_{l1}[(x_l1)B_{l1}+A_{l1}x_{l1}I]+I.`$ (183) The maps for the matrix elements are easily obtained from the matrix map, thus we do not give them explicitly. Specifically, we consider the FS. From Eq. (180), it is interesting to find that the maps for the non-diagonal elements, and for the difference of the diagonal elements, coincide with the antitrace map, Eq. (26). From Eqs. (LABEL:eq:matmap40) and (LABEL:eq:matmap4), this fact also holds for the TMS. Actually, as again follows from Eq. (97), we arrive at the important conclusion that the maps for the antitrace, the non-diagonal elements, and the difference of the diagonal elements are all the same for arbitrary substitution rules. This means that the knowledge of the trace and antitrace maps suffices to compute any physical quantities related to the global transfer matrix. ## V Applications We now turn our attention to applications of the dynamical map method developed in this paper. In what follows, we are going to consider three examples. ### A Optical Multilayers As our first example, we show how to use the antitrace map to calculate light transmission coefficients. The transmission of light through aperiodic multilayers arranged according to the FS, the “non-Fibonacci” sequence, the TMS, and the generalized TMSs was studied in the literature. Possible applications of quasiperiodic multilayers as optical switches and memories have been suggested by Schwartz. Huang et al and Yang et al found an interesting switch-like property in the light transmission through a FC($`m`$) multilayer. Using the antitrace map, we re-investigate the light transmission through FC($`m`$) which is sandwiched by two media of type $`b`$. In analogy with the discussion of Ref. , we write the corresponding transfer matrices as $`A_1`$ $`=`$ $`P_b,`$ (185) $`A_2`$ $`=`$ $`P_b^{m1}P_{ba}P_aP_{ab},`$ (186) $`A_{l+1}`$ $`=`$ $`A_l^mA_{l1}.`$ (187) The recursion relation for the transfer matrix (187) is a little different from Eq. (16) for FS($`m`$,1). It can easily be seen that the trace map is the same, but that the antitrace map differs slightly. The antitrace map is given by $`y_{l+1}`$ $`=`$ $`U_m^{(l)}\overline{w}_lU_{m1}^{(l)}y_{l1},`$ (188) $`\overline{w}_{l+1}`$ $`=`$ $`x_{l+1}y_l+U_{m1}^{(l)}\overline{w}_lU_{m2}^{(l)}y_{l1}`$ (189) where $`\overline{w}=y_{A_lA_{l1}}`$. We consider the case that the light vertically transmits the multilayer and choose the thicknesses of the layers $`d_a`$ and $`d_b`$ appropriately in order to make $`n_ad_a=n_bd_b`$. Then, we have phase differences $`\delta _a=\delta _b=\delta `$, compare Eq. (5). For $`\delta =(n+1/2)\pi `$, the propagation matrices become $`P_a`$ $`=`$ $`P_b=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).`$ (190) From the above equation and Eqs. (187), (A.1), and (A.8), we can obtain the initial conditions for the trace and antitrace maps as $`x_1`$ $`=`$ $`0,`$ (191) $`x_2`$ $`=`$ $`U_{m1}(0)\varrho _1,`$ (192) $`v_2`$ $`=`$ $`U_{m2}(0)\varrho _1,`$ (193) $`y_1`$ $`=`$ $`2,`$ (194) $`y_2`$ $`=`$ $`U_{m2}(0)\varrho _1,`$ (195) $`\overline{w}_2`$ $`=`$ $`U_{m1}(0)\varrho _1,`$ (196) where $$\varrho _m=R^m+R^m,R=n_a/n_b.$$ (197) The initial conditions depend on the parameters $`\varrho _1`$ and $`m`$, while the recursion relations only depend on $`m`$. From Eq. (A.8), we know that $`U_m(0)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[1(1)^m]i^{m1}`$ (198) $`=`$ $`\{\begin{array}{cc}0\hfill & \text{for }m=2k\hfill \\ (1)^k\hfill & \text{for }m=2k+1\hfill \end{array}.`$ (199) Since $`U_{m+4}(0)=U_m(0)`$, i.e., the function $`U_m(0)`$ is periodic in $`m`$ with period four, the initial conditions, given in Table I, also show this periodicity. The initial conditions for FC($`2q`$), $`q=1,2,3,\mathrm{}`$, or for FC($`2q+1`$), only differ by the sign of the parameter $`R`$. Thus, it is natural to divide the FC($`m`$) into two classes, FC($`2q`$) and FC($`2q+1`$). From the initial conditions and recursion equations, we can directly obtain the trace, the antitrace and the transmission coefficients of FC($`2q`$), which are given in Table II. It can be seen that the trace and the antitrace vanish alternately. The trace shows periodicity with period four for odd values of $`q`$, and period two for even $`q`$, but the antitrace shows no periodicity. Thus, the transmission coefficient also is not periodic in $`l`$. For even $`l`$, the transmission coefficient does not depend on $`m`$. However, for odd $`l`$, the transmission coefficient depends on $`m`$ and $`l`$, see Table II. Table III shows the results for FC($`2q+1`$). In this case, the trace, the antitrace and the transmission coefficient are periodic in $`l`$ with period six. The transmission coefficients are the same for $`l=2`$, $`l=3`$, and $`l=6`$ and do not depend on $`m`$. We find that the multilayer is transparent for $`l=6i+1`$, $`i=0,1,2\mathrm{}`$. Here, we not only recover the recent results of Ref. , but also give a natural classification of FC($`m`$) and derive the periodicities of the trace and antitrace maps. ### B Harmonic Chains As our second example, we show how to apply the map for the matrix elements to calculate some physical quantities for a harmonically coupled Fibonacci chain. The transmission coefficient and the Lyapunov exponent were already given in Eqs. (11) and (12). We know the trace (25) and antitrace maps (26) for this system. In order to determine the transmission coefficient, we additionally need to know the map for the difference $`z_l`$ of the diagonal elements in Eq. (11). As discussed in Sec. IV, the map for $`z_l`$ is the same as the map for the antitrace $`y_l`$. Now, this leaves us with the problem to determine the initial conditions. By a so-called transfer matrix “renormalization”, the transfer matrix product can be re-written in terms of “renormalized” transfer matrices such that these are arranged according to the FS. Following the discussion in Ref. , we choose a special value of parameters $$\mathrm{\Omega }=\frac{\alpha 2\beta +1}{\alpha (1\beta )}=\frac{m_a\omega ^2}{K_{ab}}$$ (200) where $`\alpha =m_b/m_a`$ and $`\beta =K_{aa}/K_{ab}`$. The first two renormalized transfer matrices are $`A_1`$ $`=`$ $`\left(\begin{array}{cc}1& 0\\ \eta _1& 1\end{array}\right),`$ (201) $`A_2`$ $`=`$ $`\left(\begin{array}{cc}1& 0\\ \eta _2& 1\end{array}\right),`$ (202) where $`\eta _1=2(\alpha 2)`$ and $`\eta _2=2(1\alpha )`$. Note that these two matrices commute with each other for arbitrary values of $`\eta _1`$ and $`\eta _2`$. From this equation, we obtain $$A_3=A_1A_2=\left(\begin{array}{cc}1& 0\\ \eta _2\eta _1& 1\end{array}\right).$$ (203) Thus, the initial conditions are given by $`x_1=2,x_2=2,x_3=2,`$ (204) $`y_1=\eta _1,y_2=\eta _2,y_3=\eta _2\eta _1,`$ (205) $`z_1=z_2=z_3=0.`$ (206) From the antitrace map (26), we find that $`z_l=0`$ for all $`l`$. Using the trace map (25), we easily obtain $`x_{3i+1}=2`$, $`x_{3i+2}=2`$, and $`x_{3i}=2`$. That is, the trace map is periodic in $`l`$ with period three. Then, from Eqs. (11) and (12) the transmission coefficient and the Lyapunov exponent have the simple forms $`t_l^1`$ $`=`$ $`1+{\displaystyle \frac{y_l^2}{4\mathrm{sin}^2k}},`$ (207) $`\mathrm{\Gamma }_l`$ $`=`$ $`N^1\mathrm{ln}(y_l^2+2).`$ (208) From the initial conditions for $`y_l`$ and the antitrace map (26), we easily find that the modulus of $`y_l`$ is $$\|y_l\|=\|F_l\eta _2F_{l1}\eta _1\|,l3,$$ (209) where $`F_l`$ denotes the Fibonacci number defined by the recursion $`F_l=F_{l1}+F_{l2}`$ with $`F_0=F_1=1`$. Finally, the transmission coefficient and the Lyapunov exponent are obtained as $`t_l^1`$ $`=`$ $`1+{\displaystyle \frac{(F_l\eta _2F_{l1}\eta _1)^2}{4\mathrm{sin}^2k}},`$ (210) $`\mathrm{\Gamma }_l`$ $`=`$ $`N^1\mathrm{ln}[(F_l\eta _2F_{l1}\eta _1)^2+2].`$ (211) Thus, using our matrix element maps, we have re-derived the result of Ref. . ### C Electronic systems We now apply the trace and antitrace method to the transmission problem in electronic systems for the examples of the FS and the TMS. In what follows, we choose the parameters as $`ϵ_a=ϵ_b=ϵ`$, $`t_{ab}=1`$, and $`t_{aa}=t_{bb}=t`$. #### 1 Fibonacci sequence For the FS, there are actually four different local transfer matrices $`M_n`$ (14), because the hopping matrix elements depend on three subsequent letters in the FS. Nevertheless, the transfer matrix product can be re-written in terms of two matrices $`M_b`$ $`=`$ $`\left(\begin{array}{cc}Eϵ& 1\\ 1& 0\end{array}\right)\left(\begin{array}{cc}E+ϵ& 1\\ 1& 0\end{array}\right),`$ (212) $`M_a`$ $`=`$ $`\left(\begin{array}{cc}Eϵ& t\\ 1& 0\end{array}\right)\left(\begin{array}{cc}(Eϵ)/t& 1/t\\ 1& 0\end{array}\right)`$ (214) $`\times \left(\begin{array}{cc}E+ϵ& 1\\ 1& 0\end{array}\right),`$ such that the resulting transfer matrix product is again arranged according to the Fibonacci sequence. For the trace and antitrace maps, we only need to know the first three matrices $`A_1=M_a`$, $`A_2=M_bM_a`$, and $`A_3=M_aM_bM_a`$. From Eq. (214), these matrices and thus the initial conditions are easily obtained. In order to obtain an analytical result, we restrict ourselves to the case $`E=ϵ=0`$. For this particular choice of parameters, Eq. (15) simplifies to $$t_l=\frac{4}{x_l^2+y_l^2},$$ (215) which is formally the same as Eq. (7). The initial conditions become $`x_1=0,x_2=0,x_3=2,`$ (216) $`y_1=t1/t,y_2=t+1/t,y_3=0.`$ (217) From the trace and antitrace map equations (25)–(26) for the FS, we can easily find that both the trace $`x_l`$ and the antitrace $`y_l`$ are periodic in $`l`$ with period six. In one period the traces are $`0`$, $`0`$, $`2`$, $`0`$, $`0`$, $`2`$, and the antitraces are $`t1/t`$, $`t+1/t`$, $`0`$, $`t+1/t`$, $`t+1/t`$, $`0`$. From Eq. (215), we deduce that the transmission coefficient $`t_l`$ is periodic in $`l`$ with period three. For one period, the transmission coefficients are given by $`4/(t+1/t)^2`$, $`4/(t+1/t)^2`$, and $`1`$. If the hopping parameter $`t=1`$, the transmission coefficient $`t_l=1`$ for all values of $`l`$, which is the trivial (periodic) case. Next we consider the electronic transmission for the TMS. #### 2 Thue-Morse sequence We consider the on-site model for the TMS, i.e., the hopping parameter $`t=1`$. So there are only two kinds of transfer matrices $$B_0=\left(\begin{array}{cc}E+ϵ& 1\\ 1& 0\end{array}\right),A_0=\left(\begin{array}{cc}Eϵ& 1\\ 1& 0\end{array}\right).$$ (218) From these, we can calculate the matrices $`A_1`$, $`B_1`$, $`A_2`$, and $`B_2`$, and thus the initial conditions for the trace and antitrace map. Again, in order to obtain an analytical result, we limit ourselves to the case where the parameter $`ϵ`$ and the energy $`E`$ fulfill a particular relation, $`E=\sqrt{2+ϵ^2}`$. In this case, the initial conditions become $`x_0=\sqrt{2+ϵ^2}ϵ,x_1=0,x_2=24ϵ^2,`$ (219) $`y_0=\stackrel{~}{y}_0=2,y_1=\stackrel{~}{y}_1=2\sqrt{2+ϵ^2},y_2=\stackrel{~}{y}_2=4ϵ,`$ (220) $`z_0=\sqrt{2+ϵ^2}ϵ,\stackrel{~}{z}_0=\sqrt{2+ϵ^2}+ϵ,`$ (221) $`z_1=\stackrel{~}{z}_1=2,z_2=\stackrel{~}{z}_2=4ϵ\sqrt{2+ϵ^2},`$ (222) where $`z_l=(A_l)_{11}(A_l)_{22}`$ and $`\stackrel{~}{z}_l=(B_l)_{11}(B_l)_{22}`$. From Eq. (82), we deduce that the traces $`x_l=2`$ for all $`l3`$. From the antitrace map equations (83)–(84) and the above initial conditions, we easily find that $`y_l=z_l=0`$ for $`l3`$. Thus, we obtain the result that the transmission coefficient $`t_l=1`$ for $`l3`$. For $`l=1`$ and $`l=2`$, the transmission coefficients are given by $`t_1=(2ϵ^2)/(2+ϵ^2)`$ and $`t_2=(2ϵ^2)/(2+7ϵ^2+4ϵ^4)`$, respectively. The examples considered here show that trace and antitrace maps provide a convenient tool for the computation of physical quantities related to the global transfer matrices of aperiodic substitution systems. In the applications presented above, we mainly concentrated on obtaining analytical results, and therefore had to restrict the discussion to specific values of the parameters. The trace and antitrace map equations, of course, are not restricted to these cases, but there will be no simple closed-form solutions to the recursion relations in general. The particular parameter values considered above correspond to periodic orbits of the associated dynamical systems. These cases, and probably all examples where simple solutions exist, share the property that, at a certain stage, different transfer matrices commute with each other, and thus are simultaneously diagonalizable. This also explains why these systems turn out to be transparent, because it does not matter in which order one multiplies matrices that commute with each other. In spite of these comments, the method presented here is expedient and useful for the investigation of physical systems built on aperiodic substitution sequences, because the trace and antitrace map equations can very efficiently be used in numerical investigations of large, but finite, systems. ## VI Conclusions In conclusion, we have extended the well-studied trace-map method for the investigation of aperiodic substitution systems by considering corresponding maps for the antitrace and the matrix elements of the transfer matrices. Our main results are the following. Firstly, we obtained the trace and antitrace maps for various aperiodic sequences, such as generalized FSs and TMSs, the periodic-doubling sequence, examples of ternary sequences, and the four-letter Rudin-Shapiro sequence. The dimension of the dynamical systems defined by the trace map and our antitrace maps is $`r(r+1)/2`$ plus the dimension of the trace map itself, where $`r`$ denotes the number of basic letters in the aperiodic sequence. Secondly, we showed that trace and antitrace maps can be constructed for arbitrary substitution rules. Thirdly, we introduced analogous maps for specific matrix elements of the transfer matrix, but it turns out that the maps for the off-diagonal elements and those for the difference of the diagonal elements coincide with the antitrace map. Thus, from the trace and antitrace map, we can determine any physical quantity related to the global transfer matrix of the system. Finally, as examples of applications of the trace and antitrace map method, we investigated the transmission problem for optical multilayers, harmonic chains, and electronic systems arranged according to the FS or the TMS. The trace and antitrace map method developed here can be expected to have many applications in the study of one-dimensional aperiodic systems. ###### Acknowledgements. XW thanks Shaohua Pan, Guozhen Yang, Changpu Sun, and Hongchen Fu for their encouragement and kind help. This work has been supported by Deutsche Forschungsgemeinschaft. ## A Relations for unimodular matrices For convenience, we present a collection of relevant identities, which are used in the construction of the trace and antitrace maps in Secs. II, III, and IV. The $`n`$th power of a unimodular $`2\times 2`$ matrix $`A`$ can be written as $$A^n=U_n(x_A)AU_{n1}(x_A)I,$$ (A.1) where $`I`$ is the unit matrix and $`U_n(x_A)`$ $`=`$ $`{\displaystyle \frac{\lambda _+^n\lambda _{}^n}{\lambda _+\lambda _{}}},`$ (A.2) $`\lambda _\pm `$ $`=`$ $`{\displaystyle \frac{x_A\pm \sqrt{x_A^24}}{2}}.`$ (A.3) Here $`x_A`$ and $`\lambda _\pm `$ denote the trace and the two eigenvalues of $`A`$, respectively, and $`\lambda _+\lambda _{}=detA=1`$. The functions $`U_n(x)`$ are related to the Chebyshev polynomials of the second kind $`C_n(x)`$ by $`U_n(x)=C_{n1}(x/2)`$. From the definition of the functions $`U_n(x)`$, it follows $`U_1(x)`$ $`=`$ $`1,U_0(x)=0,`$ (A.4) $`U_1(x)`$ $`=`$ $`1,U_2(x)=x,`$ (A.5) $`U_3(x)`$ $`=`$ $`x^21,U_4(x)=x^32x,`$ (A.6) $`U_{n+1}(x)`$ $`=`$ $`xU_n(x)U_{n1}(x),`$ (A.7) $`U_n^2(x)`$ $`=`$ $`U_{n+1}(x)U_{n1}(x)+1.`$ (A.8) In order to study the antitrace maps, we need the following identity $$y_{AB}=x_By_A+x_Ay_By_{BA}$$ (A.9) for the antitraces of two unimodular $`2\times 2`$ matrices $`A`$ and $`B`$. Now, we briefly prove this identity by introducing an auxiliary matrix $$\gamma =\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\gamma ^2=I,det(\gamma )=1.$$ (A.10) For the matrix $`A`$, we have $$y_A=x_{A\gamma }.$$ (A.11) Then the antitrace of $`AB`$ is given by $$y_{AB}=x_{AB\gamma }=x_{\gamma A\gamma \gamma B}.$$ (A.12) Let $`A`$, $`B`$, and $`C`$ be unimodular matrices. Then $`x_{ABAC}`$ $`=`$ $`x_{AB}x_{AC}+x_{BC}x_Bx_C.`$ (A.13) Applying the above identity to Eq. (A.12) and using Eq. (A.11) again, we obtain Eq. (A.9). It should be pointed out that Eq. (A.9) is valid for any pair of $`2\times 2`$ matrices, and it follows directly from the identity $$AB=(x_{AB}x_Ax_B)I+x_AB+x_BABA$$ (A.14) which holds for any pair of $`2\times 2`$ matrices. The detailed proof of this identity can be found in Ref. . Here, we only need to consider unimodular matrices. For $`n=2`$, Eq. (A.1) becomes $$A^2=x_AAI.$$ (A.15) This is the well-known Cayley-Hamilton theorem. From the theorem, we have $`A+A^1`$ $`=`$ $`x_AI,`$ (A.16) $`x_{A^1}`$ $`=`$ $`x_A,`$ (A.17) $`x_{BA}+x_{BA^1}`$ $`=`$ $`x_Ax_B,`$ (A.18) $`y_{BA}+y_{BA^1}`$ $`=`$ $`x_Ay_B,`$ (A.19) From Eqs. (A.14) and (A.19), we can prove the following useful relations $`BAB`$ $`=`$ $`Ax_AI+x_{AB}B=x_{AB}BA^1,`$ (A.20) $`BA^1B`$ $`=`$ $`(x_Ax_Bx_{AB})BA.`$ (A.21) Finally, from Eqs. (A.1), (A.8), and (A.15), we obtain the following relations $`x_{A^2}`$ $`=`$ $`x_A^22,`$ (A.22) $`x_{A^n}`$ $`=`$ $`U_{n+1}(x_A)U_{n1}(x_A),`$ (A.23) $`y_{A^2}`$ $`=`$ $`x_Ay_A,`$ (A.24) $`y_{A^n}`$ $`=`$ $`U_n(x_A)y_A.`$ (A.25) This completes our collection of identities. ## B Antitrace maps for some metallic mean sequences The trace and antitrace maps for the golden mean and the copper mean sequences were discussed explicitly in the main part of this paper. Here, we give the trace and antitrace maps for some other prominent examples of metallic mean sequences. From Eqs. (49), (51), and (A.8), the trace and antitrace maps for the silver mean case ($`m=2,n=1`$) are obtained as $`x_{l+1}`$ $`=`$ $`{\displaystyle \frac{x_l}{x_{l1}}}[x_l(x_{l1}^21)x_{l2}]x_{l1},`$ (B.1) $`y_{l+1}`$ $`=`$ $`{\displaystyle \frac{x_l}{x_{l1}}}(y_{l2}+y_l)+(x_l^21)y_{l1}.`$ (B.2) For the bronze mean sequence ($`m=3,n=1`$), we find $`x_{l+1}`$ $`=`$ $`{\displaystyle \frac{x_l^21}{x_{l1}^21}}[x_l(x_{l1}^32x_{l1})x_{l2}]`$ (B.4) $`x_lx_{l1},`$ $`y_{l+1}`$ $`=`$ $`{\displaystyle \frac{x_l^21}{x_{l1}^21}}(y_{l2}+x_{l1}y_l)+(x_l^32x_l)y_{l1}.`$ (B.5) Finally, for the nickel mean case ($`m=1,n=3`$), the result reads $`x_{l+1}`$ $`=`$ $`(x_{l1}^21)(x_lx_{l1}x_{l2}^3+3x_{l2})`$ (B.8) $`x_lx_{l1},`$ $`y_{l+1}`$ $`=`$ $`(x_{l1}^21)(x_{l2}^21)y_{l2}`$ (B.10) $`+x_l(x_{l1}^21)y_{l1}x_{l1}y_l.`$
warning/0005/astro-ph0005141.html	ar5iv	text	# NLS1s and Sy1s: A comparison of ionized X-ray absorber properties ## Abstract The first results from a systematic study of warm absorbers in NLS1s using BeppoSAX public archive data are presented here. We confirm ASCA results showing that a warm absorber, as modeled by two oxygen K-shell absorption edges, is less frequent in NLS1s than in broad line (BL) Sy1s ($``$20% versus $``$50%). However, our study suggests that the ionization state of NLS1s is not lower than that of BLS1s, as opposed to the ASCA-based results. The soft excess temperatures of our sample, when fitted with blackbody emission models, lie within a small range of values ($``$0.02–0.15 keV in the rest frame) with no marked dependence on source luminosity. This is in agreement with ASCA-based findings for NLS1s and early results from IUE-ROSAT BL Sy1 observations. thanks: A.O. acknowledges a fellowship of the Swiss National Foundation for Scientific Research The BeppoSAX NLS1 sample The spectra of our sample were obtained with the LECS p97 (9), MECS b97 (1), HPGSPC m97 (8) and PDS f97 (4) instruments on-board BeppoSAX. Because of its good low-energy resolution and effective area down to 0.1 keV, BeppoSAX is very efficient for the study of the complex spectral features in the soft X-rays. The broad-band sensitivity of BeppoSAX also allows a much more accurate measurement of the underlying continuum than is possible with other X-ray observatories. Our sample is composed of all the publicly available BeppoSAX NLS1 spectra before December 1999, plus the proprietary data for one source (Mrk 335). Table 1 lists the sources with their X-ray luminosities and redshifts. The X-ray luminosities span more than 4 orders of magnitude, from the narrow-line QSO PKS 0558$``$504 down to NGC 4051 in its ultra-low state. The data were processed using the SAXDAS 2.0.0 data analysis package and the spectra were rebinned in order to allow the use of the $`\chi ^2`$ statistic. Where available, data were selected in the energy ranges 0.1–4.0 keV (LECS), 1.8–10 keV (MECS), 8.0–34 keV (HPGSPC), and 15–100 keV (PDS). Because the spectra of IRAS 13224$``$3809 and the low state of NGC 4051 have a relatively low signal to noise ratio, only data below 10 keV were considered for these two objects. Fit results for the warm absorber All our fit models include Galactic absorption as well as factors to allow for flux normalization uncertainties between the instruments. These factors were constrained within their standard ranges f:99 (3). Uncertainties are quoted as 90% confidence intervals for one parameter of interest. The single power-law model does not give a good fit to any of the spectra except PKS 0558$``$504, where the fit is marginally acceptable ($`\chi ^2`$= 1.19, 109 dof). A soft excess, in the form of either power-law or blackbody emission, gives significant improvement to the fit statistics of all spectra. In 4 spectra (IRAS 13224-3809, Mrk 335, Mrk 766, NGC 4051-high), the soft excess is slightly better fit by a blackbody component than by a power law. In the remaining 5 spectra, the opposite is true. We then tried adding a Fe K$`\alpha `$ line and 2 K-shell absorption edges (Ovii and Oviii) to the power-law plus blackbody continuum. The edge energies are fixed at $`\mathrm{E}=0.74`$ and 0.87 keV in the rest frame and the line at E$`{}_{\mathrm{rest}}{}^{}=6.4`$ keV with $`\sigma =0.1`$ keV. The fit results are listed in Table 2. Only two spectra of our sample (Mrk 335 and Mrk 766) show significantly improved fits (confidence level $`>`$ 95 %) with the addition of these three spectral features. Three others show a marginal improvement (TON S 180, NGC 4051-high and Ark 564). The remaining four spectra show no improvement in the fit statistics. In all 9 spectra, we obtain only upper limits for $`\tau _{\mathrm{O7}}`$. Likewise, in 7 spectra there are upper limits for $`\tau _{\mathrm{O8}}`$. Only two spectra (Mrk 335 and Mrk 766) have definite, albeit low, values for $`\tau _{\mathrm{O8}}`$. We have also tried fitting a warm absorber model (absori in Xspec) instead of the two O edges. In this case, 5 spectra out of 9 show a significant improvement in quality of fit over the 2-edge fits: Mrk 335, TON S 180, RE J 1034+39, Ark 563 and NGC 4051-low. For the first three in this list, the improvement in $`\chi ^2`$ has an F-statistic $`>15`$. Among the remaining spectra, it is evident that PKS 0558$``$504 has no warm absorber, and that absori brings no significant improvement over the 2-edge fit for NGC 4051-low or Mrk 766. However, Mrk 766 most likely does have two O absorption edges. In the cases of NGC 4051-low and IRAS 13224$``$3809, the poor quality of the spectra does not allow one to conclude anything about the presence of soft X-ray features. Comparisons with ASCA NLS1 results and with BLS1s Recent studies based on ASCA data have shown that warm absorbers are a common spectral component of BL type 1 AGNs re:97 (10) ge:98 (5). Reynolds re:97 (10) detected Ovii and Oviii K-shell absorption edges in 50% of his sample of 24, with $`>`$37% having optical depths greater than 0.2 for one or both edges. The average optical depths for Ovii and Oviii in the Reynolds sample are $`0.29\pm 0.07`$ and $`0.18\pm 0.06`$, respectively le:99 (7). These type 1 AGN results are in sharp contrast with the values from our BeppoSAX NLS1 sample (see Table 2), that has only 2 clear detections of O absorption edges in a total of 9 spectra, corresponding to $``$20%. A similar low proportion of $``$26% was found by Leighly le:99 (7) in her study of NLS1 ASCA spectra, where 6 spectra from 5 sources (out of 23) have significant O absorption edge detections. We note that all the sources of our sample are included in the larger ASCA sample le:99 (7), but only one of our two sources with significant edge detections (Mrk 766) is among Leighly’s five “warm-absorbed” NLS1. The average optical depths for Ovii and Oviii in the ASCA sample le:99 (7) are $`0.19\pm 0.14`$ and $`0.053\pm 0.020`$, respectively<sup>1</sup><sup>1</sup>1The soft excess is modeled by a blackbody, as with our BeppoSAX sample.. The low value of $`\tau _{\mathrm{O8}}`$ suggests that the ionization state of NLS1s is lower than that of BLS1s le:99 (7). However, looking at Table 2 we see that the BeppoSAX results do not support this conclusion. The ASCA sample includes objects which may have dusty warm absorbers (e.g. IRAS 13349+2438, IRAS 17020+4454 and IRAS 20181$``$2244). If the effects of the dust are not taken in to account correctly one can derive warm absorber ionization states which are too low ko:98 (6). Furthermore, it is a common assumption that the central black holes of NLS1s accrete at very high rates. In this case, a high power output is expected that could cause the ionization state of the warm absorber to be correspondingly high. ASCA-based studies of NLS1s le:99 (7) v:99 (11) have shown that certain sources (e.g Ark 564) have broad absorption features in the 1.1-1.4 keV range. We find that we can reproduce the BeppoSAX results on Ark 564 published by Comastri et al. co:98 (2), but not those of Vaughan et al. v:99 (11) based on ASCA data. In other words, a weak emission line or absorption edge may be present, but probably no absorption line as in v:99 (11). And when a warm absorber (absori) is included, the need for such a line is no longer very strong. We have not yet tested the other sources in our sample for such features, but it is clear that no large residuals are present between 1.1-1.4 keV in our absori fits. The EWs of the Fe K$`\alpha `$ line in our sample are low and, within the uncertainties, compatible with values expected in typical Sy 1 BLR conditions, i.e. column densities of 10<sup>22-23</sup> cm<sup>-2</sup>, leading to EWs of $``$50–150 keV. The soft excess temperatures of our sample, when fitted by blackbody emission, lie within a small range of values ($``$0.02–0.15 keV in the rest frame) with no marked dependency on source luminosity. This is in agreement with ASCA-based findings for NLS1s le:99 (7) v:99 (11) and also IUE-ROSAT BL Sy1 observations wa:94 (12).
warning/0005/math0005122.html	ar5iv	text	# References INDUCED REPRESENTATIONS OF THE TWO PARAMETRIC QUANTUM DEFORMATION U<sub>pq</sub>\[gl(2/2)\] Nguyen Anh Ky Department of Physics, Chuo University, Kasuga, Bunkyo–ku Tokyo 112-8551, Japan and Institute of Physics, P.O. Box 429, Bo Ho, Hanoi 10000, Vietnam Abstract The two–parametric quantum superalgebra $`U_{p,q}[gl(2/2)]`$ and its induced representations are considered. A method for constructing all finite–dimensional irreducible representations of this quantum superalgebra is also described in detail. It turns out that finite–dimensional representations of the two–parametric $`U_{p,q}[gl(2/2)]`$, even at generic deformation parameters, are not simply trivial deformations from those of the classical superalgebra $`gl(2/2)`$, unlike the one–parametric cases. Running title : Representations of $`U_{p,q}[gl(2/2)]`$. PACS numbers : 02.20Tw, 11.30Pb. MSC–class. : 81R50; 17A70. I. Introduction Introduced in 80’s as a result of the study on quantum integrable systems and Yang–Baxter equations the quantum groups have been intensively investigated in different aspects. Since then many (algebraic and geometric) structures and various applications of quantum (super) groups have been found (see in this context, for example, Refs. ). It turns out that quantum groups are related to unrelated, at first sight, areas of both physics and mathematics (Refs. and references therein). For applications of quantum groups, as in the non–deformed cases, we often need their explicit representations. However, despite of remarkable results in this direction the problem of investigating and constructing explicit representations of quantum groups, especially those for quantum superalgebras, is still far from being satisfactorily solved. Even in the case of one–parametric quantum superalgebras, explicit representations are mainly known for quantum Lie superalgebras of lower ranks and of particular types like $`U_q[osp(1/2)]`$, $`U_q[gl(1/n)]`$ (Refs. ), while for higher rank quantum Lie superalgebras , besides some q–oscillator representations which are most popular among those constructed, we do not know so much about other representations, in particular, the finite–dimensional ones which in many cases are related to trigonometric solutions of the quantum Yang–Baxter equations . Some general aspects and the module structure finite–dimensional representations of the quantum superalgebras $`U_q[gl(m/n)]`$ are considered in Ref. but, unfortunately, their explicit construction is still absent. So far explicit finite–dimensional irreducible representations are all known and classified only for those $`U_q[gl(m/n)]`$ with both $`m`$ and $`n2`$ (see Refs. ). What about multi–parametric deformations (first considered in ), this area is even less covered and results are much poorer. Some kinds of two–parametric deformations have been considered by several authors from different points of view (see Refs. and references therein) but, to our knowledge, explicit representations are known and/or classified in a few lower rank cases such as $`U_{p,q}[sl(2/1)]`$ and $`U_{p,q}[gl(2/1)]`$ only . The latter two–parametric quantum superalgebra $`U_{p,q}[gl(2/1)]`$ was consistently defined and investigated in where all its finite–dimensional irreducible representations were explicitly constructed and classified at generic deformation parameters. This $`U_{p,q}[gl(2/1)]`$, however, is still a small quantum superalgebra which can be defined without the so–called extra–Serre defining relations representing additional constraints on odd Chevalley generators in higher rank cases. Now, in order to include the extra–Serre relations on examination we consider a bigger two–parametric quantum superalgebra, namely $`U_{p,q}[gl(2/2)]`$, and its representations. This quantum superalgebra $`U_{p,q}[gl(2/2)]`$ resembles to the one–parametric quantum superalgebra $`U_{\sqrt{pq}}[gl(2/2)]`$ but can not be identified with the latter. Here we suppose $`pq`$, otherwise we should return to the case of $`U_q[gl(2/2)]`$ investigated already in . Another our motivation is that already in the non–deformed case, the superalgebras $`gl(n/n)`$, especially, their subalgebras $`sl(n/n)`$ and $`psl(n/n)`$, have special properties (in comparison with other $`gl(m/n)`$, $`mn`$) and, therefore, attract interest . Additionally, structures of two–parameter deformations considered in and here are, of course, richer than those of one–parameter deformations. Every deformation parameter can be independently chosen to take a separate generic value (including zero) or to be a root of unity. Combining the advantages of the previously developed methods for $`U_q[gl(2/2)]`$ and $`U_{p,q}[gl(1/2)]`$ (see Refs. ) we can construct all finite–dimensional representations of the two–parametric quantum Lie superalgebra $`U_{p,q}[gl(2/2)]`$. In the frame–work of this paper we consider representations at generic $`p`$ and $`q`$ only (i.e., $`p`$ and $`q`$ are not roots of unity), while representations at roots of unity are a subject of later separate investigations. In comparison with previous papers , the approach here is somewhat modified because of some specific features arising in the present case but the main steps in the construction procedure remain the same. Following this approach we can directly construct explicit representations of the quantum superalgebra $`U_{p,q}[gl(2/2)]`$ induced from some (usually, irreducible) finite–dimensional representations of the even subalgebra $`U_{p,q}[gl(2)gl(2)]`$ which itself is a quantum algebra. Since the latter is a stability subalgebra of $`U_{p,q}[gl(2/2)]`$ we expect the representations of $`U_{p,q}[gl(2/2)]`$ constructed are decomposed into finite–dimensional irreducible representations of $`U_{p,q}[gl(2)gl(2)]`$. For a clear description of this decomposition we shall introduce a $`U_{p,q}[gl(2/2)]`$–basis (i.e., a basis within a $`U_{p,q}[gl(2/2)]`$–module or briefly a basis of $`U_{p,q}[gl(2/2)]`$) which will be convenient for us in investigating the module structure. This basis (see (4.26)) can be expressed in terms of some basis of the even subalgebra $`U_{p,q}[gl(2)gl(2)]`$ which in turn represents a (tensor) product between two $`U_{p,q}[gl(2)]`$–bases referred to as the left and the right ones. As is shown in , the Gel’fand–Zetlin (GZ) patterns can serve again as a basis of finite–dimensional representations of $`U_{p,q}[gl(2)]`$. Thus, finite–dimensional representations of $`U_{p,q}[gl(2)gl(2)]`$ are realized in tensor products of two such $`U_{p,q}[gl(2)]`$ GZ bases. For generic $`p`$ and $`q`$, the finite–dimensional $`U_{p,q}[gl(2/2)]`$–modules constructed have similar structures to those of $`U_q[gl(2/2)]`$ investigated in and to those of $`gl(2/2)`$ investigated in . However, finite–dimensional representations of $`U_{p,q}[gl(2/2)]`$ at generic deformation parameters are not simply trivial deformations from those of $`gl(2/2)`$ that is the former can not be obtained from the latter by putting quantum deformation brackets in appropriate places, unlike many cases of one–parametric deformations. When one or both of $`p`$ and $`q`$ are roots of unity the structures of $`U_{p,q}[gl(2/2)]`$–modules are drastically different but we hope that the present method for construction of finite–dimensional representations of $`U_{p,q}[gl(2/2)]`$ at generic deformation parameters can be extended on its finite–dimensional representations at roots of unity. This paper is organized as follows. The two–parametric quantum superalgebra $`U_{p,q}[gl(2/2)]`$ is consistently defined in section 2 where we also describe how to construct its representations induced from representations of the stability subalgebra $`U_{p,q}[gl(2)gl(2)]`$. Section 3 is devoted to constructing finite–dimensional representations of $`U_{p,q}[gl(2/2)]`$. Finally, some comments and conclusions are made in section 4. Let us list some abbreviations and notations used throughout the paper: \| fidirmod(s) : finite-dimensional irreducible module(s), \| \| --- \| \| GZ basis : Gel’fand–Zetlin basis, \| \| QGZ basis : Quasi–Gel’fand–Zetlin basis, \| \| lin.env.{X} : linear envelope of X, \| \| $`p,q`$ : the deformation parameters, \| \| $`[x][x]_{p,q}=\frac{q^xp^x}{qp^1}`$, : a $`pq`$–deformation of a number or an operator $`x`$, \| \| $`V_l^{p,q}V_r^{p,q}`$ : a tensor product between two linear spaces $`V_l^{p,q}`$ and $`V_r^{p,q}`$ \| \| \| or a tensor product between a $`U_{p,q}[gl(2)_l]`$–module $`V_l^{p,q}`$ and \| \| \| a $`U_{p,q}[gl(2)_r]`$–module $`V_r^{p,q}`$, \| \| $`T^{p,q}V_0^{p,q}`$ : a tensor product between two $`U_{p,q}[gl(2)gl(2)]`$–modules \| \| \| $`T^{p,q}`$ and $`V_0^{p,q}`$, \| \| $`[E,F\}`$ : supercommutator between $`E`$ and $`F`$, \| \| $`[E,F]_rEFrFE`$ : an r–deformed commutator between $`E`$ and $`F`$, \| We hope the notations $`[x][x]_{p,q}`$ for quantum deformations, $`[m]`$ for highest weights (signatures) in (quasi–) GZ bases $`(m)`$ and \[ , \] for commutators do not confuse the reader. II. The quantum superalgebra $`U_{p,q}[gl(2/2)]`$ The quantum superalgebra $`U_{p,q}U_{p,q}[gl(2/2)]`$ as a two–parametric deformation of the universal enveloping algebra $`U[gl(2/2)]`$ of the Lie superalgebra $`gl(2/2)`$ is generated by the operators $`L_k`$, $`E_{12}`$, $`E_{23}`$, $`E_{34}`$, $`E_{21}`$, $`E_{32}`$, $`E_{43}`$ and $`E_{ii}`$ ($`1i4`$) called again Cartan–Chevalley generators and satisfying \| a) the super–commutation relations ($`1i,i+1,j,j+14`$): \| \| --- \| \| \| $`[E_{ii},E_{jj}]`$ \| = \| 0, \| (2.1a) \| \| \| $`[E_{ii},E_{j,j+1}]`$ \| = \| $`(\delta _{ij}\delta _{i,j+1})E_{j,j+1},`$ \| (2.1b) \| \| \| $`[E_{ii},E_{j+1,j}]`$ \| = \| $`(\delta _{i,j+1}\delta _{ij})E_{j+1,j}`$, \| (2.1c) \| \| \| \[even generator, $`L_k`$\] \| = \| 0, $`k=1,2,3`$, \| (2.1d) \| \| \| $`[E_{i,i+1},E_{j+1,j}\}`$ \| = \| $`\delta _{ij}\left(\frac{q}{p}\right)^{L_ih_i(1+\delta _{i2})/2}[h_i]`$, \| (2.1e) \| \| with $`h_i=(E_{ii}\frac{d_{i+1}}{d_i}E_{i+1,i+1})`$, $`L_1=L_l,L_2=0,L_3=L_r`$ and $`d_1=d_2=d_3=d_4`$ \| \| $`=1`$, \| \| b) the Serre–relations: \| \| --- \| \| \| \| $`[E_{12},E_{34}]`$ \| = \| $`[E_{21},E_{43}]`$ \| = 0, \| (2.2a) \| \| \| \| $`E_{23}^2`$ \| = \| $`E_{32}^2`$ \| = 0, \| (2.2b) \| \| \| $`[E_{12},E_{13}]_p`$ = \| $`[E_{21},E_{31}]_q`$ \| = \| $`[E_{24},E_{34}]_q`$ \| = $`[E_{42},E_{43}]_p=0`$, \| (2.2c) \| and c) the extra–Serre relations: $$\{E_{13},E_{24}\}=0,$$ $`(2.3a)`$ $$\{E_{31},E_{42}\}=0.$$ $`(2.3b)`$ Here, the operators $$E_{13}:=[E_{12},E_{23}]_{q^1},$$ $`(2.4a)`$ $$E_{24}:=[E_{23},E_{34}]_{p^1},$$ $`(2.4b)`$ $$E_{31}:=[E_{21},E_{32}]_{p^1},$$ $`(2.4c)`$ $$E_{42}:=[E_{32},E_{43}]_{q^1}$$ $`(2.4d)`$ and the operators composed in the following way \| $`E_{14}`$ \| := \| $`[E_{12},[E_{23},E_{34}]_{p^1}]_{q^1}`$ \| $`[E_{12},E_{24}]_{q^1}`$, \| (2.5a) \| \| --- \| --- \| --- \| --- \| --- \| \| $`E_{41}`$ \| := \| $`[E_{21},[E_{32},E_{43}]_{q^1}]_{p^1}`$ \| $`[E_{21},E_{42}]_{p^1}`$ \| (2.5b) \| are defined as new generators which, like $`E_{23}`$ and $`E_{32}`$, are all odd and have vanishing squares. These generators $`E_{ij}`$, $`1i,j4`$, are two–parametric deformation analogues ($`pq`$–analogues) of the Weyl generators $`e_{ij}`$, $`1i,j4`$, of the superalgebra $`gl(2/2)`$ whose universal enveloping algebra $`U[gl(2/2)]`$ is a classical limit of $`U_{p,q}[gl(2/2)]`$ when $`p,q1`$. The so–called maximal–spin operators $`L_k`$ are constants within a $`U_{p,q}[gl(2)]`$–fidirmod and are different for different $`U_{p,q}[gl(2)]`$–fidirmods. Therefore, commutators between these operators with the odd generators intertwining $`U_{p,q}[gl(2)]`$–fidirmods take concrete forms on concrete basis vectors. Other commutation relations between $`E_{ij}`$ follow from the relations (2.1)–(2.3) and the definitions (2.4) and (2.5). The subalgebra $`U_{p,q}[gl(2/2)_0](U_{p,q}[gl(2/2)]_0U_{p,q}[gl(2/2)])`$ is even and isomorphic to $`U_{p,q}[gl(2)gl(2)]U_{p,q}[gl(2)]U_{p,q}[gl(2)]`$ which is completely defined by $`L_1`$, $`L_3`$, $`E_{12}`$, $`E_{34}`$, $`E_{21}`$, $`E_{43}`$ and $`E_{ii}`$, $`1i4`$, $$U_q[gl(2/2)_0]=lin.env.\{L_1,L_3,E_{ij}i,j=1,2andi,j=3,4\}.$$ $`(2.6)`$ In order to distinguish two components $`U_{p,q}[gl(2)]`$ of $`U_{p,q}[gl(2/2)_0]`$ we set $$leftU_{p,q}[gl(2)]U_{p,q}[gl(2)_l]:=lin.env.\{L_1,E_{ij}i,j=1,2\},$$ $`(2.7)`$ $$rightU_{p,q}[gl(2)]U_{p,q}[gl(2)_r]:=lin.env.\{L_3,E_{ij}i,j=3,4\},$$ $`(2.8)`$ that is $$U_{p,q}[gl(2/2)_0]=U_{p,q}[gl(2)_lgl(2)_r].$$ $`(2.9)`$ Looking at the relations (2.1)–(2.3) we see that every of the odd spaces $`A_+`$ and $`A_{}`$ spanned on the positive and negative odd roots (generators) $`E_{ij}`$ and $`E_{ji}`$, $`1i2<j4`$, respectively $$A_+=lin.env.\{E_{14},E_{13},E_{24},E_{23}\},$$ $`(2.10)`$ $$A_{}=lin.env.\{E_{41},E_{31},E_{42},E_{32}\},$$ $`(2.11)`$ is a representation space of the even subalgebra $`U_{p,q}[gl(2/2)_0]`$ which, as seen from (2.1)–(2.2), is a stability subalgebra of $`U_{p,q}[gl(2/2)]`$. Therefore, we can construct representations of $`U_{p,q}[gl(2/2)]`$ induced from some (finite–dimensional irreducible, for example) representations of $`U_{p,q}[gl(2/2)_0]`$ which are realized in some representation spaces (modules) $`V_0^{p,q}`$ representing tensor products of $`U_{p,q}[gl(2)_l]`$–modules $`V_{0,l}^{p,q}`$ and $`U_{p,q}[gl(2)_r]`$–modules $`V_{0,r}^{p,q}`$ $$V_0^{p,q}(\mathrm{\Lambda })=V_{0,l}^{p,q}(\mathrm{\Lambda }_l)V_{0,r}^{p,q}(\mathrm{\Lambda }_r),$$ $`(2.12)`$ where $`\mathrm{\Lambda }`$’s are some signatures (such as highest weights, respectively) characterizing the modules (highest weight modules, respectively). Here $`\mathrm{\Lambda }_l`$ and $`\mathrm{\Lambda }_r`$ are referred to as the left and the right components of $`\mathrm{\Lambda }`$, respectively $$\mathrm{\Lambda }=[\mathrm{\Lambda }_l,\mathrm{\Lambda }_r].$$ $`(2.13)`$ If we demand $$E_{23}V_0^{p,q}(\mathrm{\Lambda })=0$$ $`(2.14)`$ hence $$U_{p,q}(A_+)V_0^{p,q}=0,$$ $`(2.15)`$ we turn the $`U_{p,q}[gl(2/2)_0]`$–module $`V_0^{p,q}`$ into a $`U_{p,q}(B)`$–module where $$B=A_+gl(2)gl(2).$$ $`(2.16)`$ The $`U_{p,q}[gl(2/2)]`$–module $`W^{p,q}`$ induced from the $`U_{p,q}[gl(2/2)_0]`$–module $`V_0^{p,q}`$ is the factor–space $$W^{p,q}=W^{p,q}(\mathrm{\Lambda })=[U_{p,q}V_0^{p,q}(\mathrm{\Lambda })]/I^{p,q}(\mathrm{\Lambda })$$ $`(2.17)`$ which, of course, depends on $`\mathrm{\Lambda }`$, where $$U_{p,q}U_{p,q}[gl(2/2)],$$ $`(2.18)`$ while $`I^{p,q}`$ is the subspace $$I^{p,q}=lin.env.\{ubvubvuU_{p,q},bU_{p,q}(B)U_{p,q},vV_0^{p,q}\}.$$ $`(2.19)`$ Using the commutation relations (2.1)–(2.3) and the definitions (2.4) and (2.5) we can prove the following analogue of the Poincaré–Birkhoff–Witt theorem Proposition 1: The quantum deformation $`U_{p,q}:=U_{p,q}[gl(2/2)]`$ is spanned on all possible linear combinations of the elements $$g=(E_{23})^{\eta _1}(E_{24})^{\eta _2}(E_{13})^{\eta _3}(E_{14})^{\eta _4}(E_{41})^{\theta _1}(E_{31})^{\theta _2}(E_{42})^{\theta _3}(E_{32})^{\theta _4}g_0,$$ $`(2.20)`$ or equivalently $$g=(E_{41})^{\theta _1}(E_{31})^{\theta _2}(E_{42})^{\theta _3}(E_{32})^{\theta _4}b,$$ $`(2.21)`$ where $`g_0U_{p,q}[gl(2/2)_0]`$, $`bU_{p,q}(B)`$ and $`\eta _i`$, $`\theta _i=0,1`$. Any vector $`w`$ from the module $`W^{p,q}`$ can be represented as $$w=uv,uU_{p,q},vV_0^{p,q}.$$ $`(2.22)`$ Then $`W^{p,q}`$ is a $`U_{p,q}[gl(2/2)]`$–module in the sense $$gwg(uv)=guvW^{p,q},$$ $`(2.23)`$ for $`g,uU_{p,q}`$, $`wW^{p,q}`$ and $`vV_0^{p,q}`$. Taking into account the fact that $`V_0^{p,q}(\mathrm{\Lambda })`$ is a $`U_{p,q}(B)`$–module we have $$W^q(\mathrm{\Lambda })=lin.env.\{(E_{41})^{\theta _1}(E_{31})^{\theta _2}(E_{42})^{\theta _3}(E_{32})^{\theta _4}vvV_0^{p,q},\theta _1,\mathrm{},\theta _4=0,1\}.$$ $`(2.24)`$ Consequently, a basis of $`W^{p,q}`$ can be constituted by taking all the vectors of the form $$\|\theta _1,\theta _2,\theta _3,\theta _4;(\lambda ):=(E_{41})^{\theta _1}(E_{31})^{\theta _2}(E_{42})^{\theta _3}(E_{32})^{\theta _4}(\lambda ),\theta _i=0,1,$$ $`(2.25)`$ where $`(\lambda )`$ is a (GZ, for example) basis of $`V_0^{p,q}V_0^{p,q}(\mathrm{\Lambda })`$. We refer to this basis of $`W^{p,q}`$ as the induced $`U_{p,q}[gl(2/2)]`$–basis (or simply, the induced basis) in order to distinguish it from another $`U_{p,q}[gl(2/2)]`$–basis introduced later and called a reduced basis which is more convenient for investigating the module structure of $`W^{p,q}`$. It is obvious that if the module $`V_0^{p,q}`$ is finite–dimensional so is the module $`W^{p,q}`$. Finite–dimensional representations of $`U_{p,q}[gl(2/2)]`$ are namely the subject of the next section. III. Finite–dimensional representations of $`U_{p,q}[gl(2/2)]`$ In this section we consider finite–dimensional representations of $`U_{p,q}[gl(2/2)]`$ induced from irreducible finite–dimensional representations of $`U_{p,q}[gl(2/2)_0]`$. We firstly construct the bases of the module $`W^q`$ and then find the explicit matrix elements for the finite–dimensional representations of $`U_{p,q}[gl(2/2)]`$. We can shown that the GZ patterns $$\left[\begin{array}{c}m_{12}m_{22}\\ m_{11}\end{array}\right]\left[\begin{array}{c}[m]\\ m_{11}\end{array}\right]$$ $`(3.1)`$ where $`m_{ij}`$ are complex numbers such that $`m_{12}m_{11}𝐙_+`$ and $`m_{11}m_{22}𝐙_+`$, can serve as a basis of a $`U_{p,q}[gl(2)]`$–fidirmod. Indeed, finite–dimesional representations of $`U_{p,q}[gl(2)]`$ are high weight and if the operators $`L`$ and $`E_{ij}`$, $`i,j=1,2,`$ are difined on the basis (3.1) as follows $`L\left[\begin{array}{c}m_{12}m_{22}\\ m_{11}\end{array}\right]`$ $`=`$ $`{\displaystyle \frac{1}{2}}(l_{12}l_{22}1)\left[\begin{array}{c}m_{12}m_{22}\\ m_{11}\end{array}\right],`$ $`E_{11}\left[\begin{array}{c}m_{12}m_{22}\\ m_{11}\end{array}\right]`$ $`=`$ $`(l_{11}+1)\left[\begin{array}{c}m_{12}m_{22}\\ m_{11}\end{array}\right],`$ $`E_{22}\left[\begin{array}{c}m_{12}m_{22}\\ m_{11}\end{array}\right]`$ $`=`$ $`(l_{12}+l_{22}l_{11}+2)\left[\begin{array}{c}m_{12}m_{22}\\ m_{11}\end{array}\right],`$ $`E_{12}\left[\begin{array}{c}m_{12}m_{22}\\ m_{11}\end{array}\right]`$ $`=`$ $`\left([l_{12}l_{11}][l_{11}l_{22}]\right)^{1/2}\left[\begin{array}{c}m_{12}m_{22}\\ m_{11}+1\end{array}\right],`$ $`E_{21}\left[\begin{array}{c}m_{12}m_{22}\\ m_{11}\end{array}\right]`$ $`=`$ $`\left([l_{12}l_{11}+1][l_{11}l_{22}1]\right)^{1/2}\left[\begin{array}{c}m_{12}m_{22}\\ m_{11}1\end{array}\right],(3.2)`$ they really satisfy commutation relations of $`U_{p,q}[gl(2)]`$ given in (2.1). Here the notation $$l_{ij}=m_{ij}i,i=1,2$$ $`(3.3a)`$ and later also the notation $$l_{ij}^{}=m_{ij}i+2,i=3,4,$$ $`(3.3b)`$ are used. Since the $`U_{p,q}[gl(2/2)_0]`$–fidirmod $`V_0^{p,q}`$ is decomposed into a $`U_{p,q}[gl(2)_l]`$–fidirmod $`V_{0,l}^{p,q}`$ and a $`U_{p,q}[gl(2)_r]`$–fidirmod $`V_{0,r}^{p,q}`$ via the tensor product $$V_0^{p,q}=V_{0,l}^{p,q}V_{0,r}^{p,q},$$ $`(3.4)`$ its basis, therefore, is a tensor product $$\left[\begin{array}{c}m_{13}m_{23}\\ m_{11}\end{array}\right]\left[\begin{array}{c}m_{33}m_{43}\\ m_{31}\end{array}\right]\left[\begin{array}{c}[m]_l\\ m_{11}\end{array}\right]\left[\begin{array}{c}[m]_r\\ m_{31}\end{array}\right](m)_l(m)_r(m)$$ $`(3.5)`$ between a GZ basis of $`V_{0,l}^{p,q}`$ spanned on the vectors $`(m)_l`$ and a GZ basis of $`V_{0,r}^{p,q}`$ spanned on the vectors $`(m)_r`$. Following the approach of Ref. (see aslo Ref.) and keeping the notations used there, we can represent the basis (3.5) of $`V_0^{p,q}`$ in the form $$\left[\begin{array}{ccc}\begin{array}{c}m_{13}m_{23}\\ m_{11}\end{array};\begin{array}{c}m_{33}m_{43}\\ m_{31}\end{array}\hfill & & \end{array}\right]\left[\begin{array}{ccc}\begin{array}{c}[m]_l\\ m_{11}\end{array};\begin{array}{c}[m]_r\\ m_{31}\end{array}\hfill & & \end{array}\right](m)$$ $`(3.6)`$ Then, the signature $`\mathrm{\Lambda }`$, which now is the highest weight, is given by the first row $`[m_{13},m_{23},m_{33},m_{43}][[m]_l,[m]_r][m]`$ common for all the basis vectors (3.6) of $`V_0^{p,q}`$: $$V_0^{p,q}V_0^{p,q}(\mathrm{\Lambda })=V_0^{p,q}([m])=V_{0,l}^{p,q}([m]_l)V_{0,r}^{p,q}([m]_r).$$ $`(3.7)`$ The explicit action of $`U_{p,q}[gl(2/2)_0]`$ on $`V_0^{p,q}([m])`$ follows directly from (3.2) and : $$g_0(m)=g_{0,l}(m)_l(m)_r+(m)_lg_{0,r}(m)_r$$ $`(3.8)`$ for $`g_0g_{0,l}g_{0,r}U_q[gl(2/2)_0]`$ and $`(m)V_0^q([m])`$. The basis vector with $`m_{11}=m_{13}`$ and $`m_{31}=m_{33}`$ $$[\begin{array}{c}m_{13}m_{23}\\ m_{13}\end{array};\begin{array}{c}m_{33}m_{43}\\ m_{33}\end{array}][\begin{array}{c}[m]_l\\ m_{13}\end{array};\begin{array}{c}[m]_r\\ m_{33}\end{array}](M)$$ $`(3.9)`$ satisfying the conditions \| $`E_{ii}(M)`$ \| = \| $`m_{i3}(M)`$, $`i=1,2,3,4`$, \| \| --- \| --- \| --- \| \| $`E_{12}(M)`$ \| = \| $`E_{34}(M)=0`$ \| (3.10) \| is the highest weight vector in $`V_0^{p,q}([m])`$. Therefore, as in the classical case ($`p=q=1`$) and in the case of one–parametric deformation ($`p=q`$) the highest weight $`[m]`$ is nothing but an ordered set of the eigen values of the Cartan generators $`E_{ii}`$ on the highest weight vector $`(M)`$. The latter is also highest weight vector in $`W^{p,q}([m])`$ because of the condition (2.14). All other, i.e. lower weight, basis vectors of $`V_0^{p,q}`$ can be obtained from the highest weight vector $`(M)`$ through acting on the latter by monomials of the lowering generators $`E_{21}`$ and $`E_{43}`$ in definite powers: $`(m)`$ $`=`$ $`\left({\displaystyle \frac{[m_{11}m_{23}]![m_{31}m_{43}]!}{[m_{13}m_{23}]![m_{13}m_{11}]![m_{33}m_{43}]![m_{33}m_{31}]!}}\right)^{1/2}`$ $`\times (E_{21})^{m_{13}m_{11}}(E_{43})^{m_{33}m_{31}}(M),(3.11)`$ where $`[n]`$’s stand for $$\frac{q^np^n}{qp^1}[n]_{p,q}[n],$$ $`(3.12)`$ while $$[n]!=[1][2]\mathrm{}[n1][n].$$ $`(3.13)`$ Therefore, the induced basis (2.25) of $`W^{p,q}(\mathrm{\Lambda })=W^{p,q}([m])`$ now takes the form $$\|\theta _1\theta _2,\theta _3,\theta _4;(m):=(E_{41})^{\theta _1}(E_{31})^{\theta _2}(E_{42})^{\theta _3}(E_{32})^{\theta _4}(m).$$ $`(3.14)`$ The subspace $`T^{p,q}`$ consisting of $$\|\theta _1,\theta _2,\theta _3,\theta _4:=(E_{41})^{\theta _1}(E_{31})^{\theta _2}(E_{42})^{\theta _3}(E_{32})^{\theta _4}$$ $`(3.15)`$ can be considered as a $`U_{p,q}[gl(2/2)_0]`$–adjoint module which is 16–dimensional when all the $`\theta _i`$ ($`i=1,2,3,4`$) take all two possible values 0 and 1, that is $`_{i=1}^4\theta _i`$ runs all over the range from 0 to 4. Thus $`W^{p,q}([m])`$ being a tensor product between two $`U_{p,q}[gl(2/2)_0]`$–modules: $$W^{p,q}([m])=T^{p,q}V_0^{p,q}([m]),$$ $`(2.24^{})`$ is, in general, a reducible $`U_{p,q}[gl(2/2)_0]`$–module and are decomposed into irreducible $`U_{p,q}[gl(2/2)_0]`$–submodules. We arrive at the next assertion Proposition 2: The induced $`U_{p,q}[gl(2/2)]`$–module $`W^{p,q}`$ is the linear span $$W^{p,q}([m])=lin.env.\{(E_{41})^{\theta _1}(E_{31})^{\theta _2}(E_{42})^{\theta _3}(E_{32})^{\theta _4}vvV_0^{p,q}([m]),\theta _i=0,1\},$$ $`(2.24^{\prime \prime })`$ which is decomposed into a direct sum of (sixteen, at most) $`U_{p,q}[gl(2/2)_0]`$–fidirmods $`V_k^{p,q}([m]_k)`$: $$W^{p,q}([m])=\underset{k=0}{\overset{15}{}}V_k^{p,q}([m]_k).$$ $`(3.16)`$ where $`[m]_k`$ are signatures of $`V_k^{p,q}V_k^{p,q}([m]_k)`$. Here, we call $`[m]_k[m_{12},m_{22},m_{32},m_{42}]_k`$ the local highest weights of the submodules $`V_k^{p,q}`$ in their GZ bases denoted now as $$\left[\begin{array}{ccc}\begin{array}{c}m_{12}m_{22}\\ m_{11}\end{array};\begin{array}{c}m_{32}m_{42}\\ m_{31}\end{array}\hfill & & \end{array}\right]_k(m)_k.$$ $`(3.17)`$ The highest weight $`[m]_0[m]`$ of $`V_0^{p,q}`$ being also the highest weight of $`W^{p,q}`$ is referred to as the global highest weight. We call $`[m]_k`$, $`k0`$, the local highest weights in the sense that they characterize the submodules $`V_k^{p,q}W^{p,q}`$ as $`U_{p,q}[gl(2/2)_0]`$–fidirmods only, while the global highest weight $`[m]`$ characterizes the $`U_{p,q}[gl(2/2)]`$–module $`W^{p,q}`$ as the whole. In the same way we define the local highest weight vectors $`(M)_k`$ in $`V_k^{p,q}`$ as those $`(m)_k`$ satisfying the conditions (cf. (3.10)) $`E_{ii}(M)_k`$ $`=`$ $`m_{i2}(M)_k,i=1,2,3,4,`$ $`E_{12}(M)_k`$ $`=`$ $`E_{34}(M)_k=0.(3.18)`$ The highest weight vector $`(M)`$ of $`V_0^{p,q}`$ is also the global highest weight vector in $`W^{p,q}`$ for which the condition (see (2.14)) $$E_{23}(M)=0$$ $`(3.19)`$ and the conditions (3.18) simultaneously hold. Let us denote by $`\mathrm{\Gamma }_k^{p,q}`$ the basis system spanned on the basis vectors $`(m)_k`$ (3.17) in each $`V_k^{p,q}([m])`$. For a basis of $`W^{p,q}`$ we can choose the union $`\mathrm{\Gamma }^{p,q}=_{k=0}^{15}\mathrm{\Gamma }_k^{p,q}`$ of all the bases $`\mathrm{\Gamma }_k^{p,q}`$, namely, a basis vector of $`W^{p,q}`$ has to be identified with one of the vectors $`(m)_k`$, $`0k15`$. The basis $`\mathrm{\Gamma }^{p,q}`$ is referred to as the $`U_{p,q}[gl(2/2)]`$–reduced basis or simply, the reduced basis. It is clear that every basis $`\mathrm{\Gamma }_k^{p,q}=\mathrm{\Gamma }_k([m]_k)^{p,q}`$ is labelled by a local highest weight $`[m]_k`$, while the basis $`\mathrm{\Gamma }^{p,q}=\mathrm{\Gamma }^{p,q}([m])`$ is labelled by the global highest weight $`[m]`$. Going ahead, we modify the notation (3.17) for the basis vectors in $`\mathrm{\Gamma }^{p,q}`$ as follows (cf. (4.26) in Ref. ) $$\left[\begin{array}{cccc}m_{13}\hfill & m_{23}& m_{33}& m_{43}\\ m_{12}\hfill & m_{22}& m_{32}& m_{42}\\ m_{11}\hfill & 0& m_{31}& 0\end{array}\right]_k\left[\begin{array}{ccc}\begin{array}{c}m_{12}m_{22}\\ m_{11}\end{array};\begin{array}{c}m_{32}m_{42}\\ m_{31}\end{array}\hfill & & \end{array}\right]_k(m)_k,$$ $`(3.20)`$ with $`k`$ running from 0 to 15 as for $`k=0`$ we must take into account $`m_{i2}=m_{i3}`$, $`i=1,2,3,4`$, i.e., $$(m)_0(m)=\left[\begin{array}{cccc}m_{13}\hfill & m_{23}& m_{33}& m_{43}\\ m_{13}\hfill & m_{23}& m_{33}& m_{43}\\ m_{11}\hfill & 0& m_{31}& 0\end{array}\right].$$ $`(3.21)`$ In (3.20) the first row $`[m]=[m_{13},m_{23},m_{33},m_{43}]`$ being the (global) highest weight of $`W^{p,q}`$ is fixed for all the vectors in the whole $`W^{p,q}`$ and characterizes this module itself, while the second row is a (local) highest weight of some submodule $`V_k^{p,q}`$ and tells us that the considered basis vector $`(m)_k`$ of $`W^{p,q}`$ belongs to this submodule in the decomposition (3.16) corresponding to the branching rule $`U_{p,q}[gl(2/2)]U_{p,q}[gl(2/2)_0]U_{p,q}[gl(1)gl(1)]`$. We refer to (3.20) as the quasi–Gel’fand–Zetlin (QGZ) basis. It is easy to see that the highest vectors $`(M)_k`$ in the notation (3.20) are $$(M)_k=\left[\begin{array}{cccc}m_{13}\hfill & m_{23}& m_{33}& m_{43}\\ m_{12}\hfill & m_{22}& m_{32}& m_{42}\\ m_{12}\hfill & 0& m_{32}& 0\end{array}\right]_k,k=0,1,\mathrm{}15.$$ $`(3.22)`$ The (global) highest weight vector $`(M)`$ (3.9) is given now by $$(M)=\left[\begin{array}{cccc}m_{13}\hfill & m_{23}& m_{33}& m_{43}\\ m_{13}\hfill & m_{23}& m_{33}& m_{43}\\ m_{13}\hfill & 0& m_{33}& 0\end{array}\right].$$ $`(3.23)`$ A highest weight vector $`(M)_k`$ expressed in terms of the induced basis (3.14) has the form of a homogeneous polynomial of a definite degree $`\eta `$ in negative odd generators ($`E_{ij}`$, $`1j2<i4`$) acting on $`(m)V_0^{p,q}([m])`$: $$(M)_k(M)_{\eta ,h}=\underset{\theta _i=0,1}{}C_{\eta ,h}(\theta _1,\theta _2,\theta _3,\theta _4)\|\theta _1,\theta _2,\theta _3,\theta _4;(m)$$ $`(3.24)`$ with $`\eta =_{i=1}^4\theta _i`$ fixed for every $`(M)_{\eta ,h}`$, and the coefficients $`C_{\eta ,h}`$ determined by solving Eqs. (3.18). Applying (3.11) to any $`(M)_{\eta ,h}`$ we find all the basis vectors $`(m)_{\eta ,h}`$ of the corresponding fidirmod $`V_{\eta ,h}^{p,q}`$ which is a linear space spanned on homogeneous polynomials of the negative odd generators of the same degree $`\eta `$ since (3.11) does not change $`\eta `$. Here we call $`\eta `$ the level of $`V_{\eta ,h}^{p,q}`$. It is easy to see that on the level $`\eta =0`$ there is only one fidirmod, namely $`V_0^{p,q}V_{0,1}^{p,q}`$, while on the next level $`\eta =1`$ there are four fidirmods, say, $`V_{1,h}^{p,q},h=1,2,3,4`$. On the level $`\eta =2`$ we can find six fidirmods $`V_{2,h}^{p,q},1h6`$, which are divided into two groups ($`h=1,2,3`$ and $`h=4,5,6`$) expressed in terms of two indenpendent groups of second order monomials of odd generators (3.15) acting on $`(m)`$. For $`\eta =3`$ the number of fidirmods is four, $`V_{3,h}^{p,q},h=1,2,3,4`$, and finally on $`\eta =4`$ we find again only one fidirmod $`V_{4,1}^{p,q}`$. However, this form (3.24) which was used in the one–parametric case is now inconvenient for us here to apply formula (3.11) in order to find all other (i.e., lower weight) reduced basis vectors. It is so because of the presence of the maximal–spin operators $`L_i`$ which are not diagonalized in the induced basis but in the reduced basis (since an eigenvalue of any $`L_i`$ is a fixed constant only within a $`U_{p,q}[gl(2)]`$–fidirmod (or fidirmod, for short) and changes from fidirmod to fidirmod). Applying (3.11) we have to push generators $`E_{21}`$ and $`E_{43}`$ to the right side until reaching $`V_0^{p,q}`$ by using commutation ralations (2.1)–(2.2) which give rise to $`L_i`$’s acting on the induced basis vectors. But it is extremtly difficult to get explicit actions of $`L_i`$ on the latter vectors before knowing how they are projected on the reduced basis which, however, we are now looking for. Instead, we will write down $`(M)_k(M)_{\eta ,h}`$ in a form convenient for applying (3.11) which leaves the $`\eta `$’s unchanged: $`(M)_0`$ $``$ $`(M)_{0,1}=a_0\|0,0,0,0;(M)(M),a_01,`$ $`(M)_1`$ $``$ $`(M)_{1,1}=a_1\|0,0,0,1;(M)a_1E_{32}(M),`$ $`(M)_2`$ $``$ $`(M)_{1,2}=a_2\left\{{\displaystyle \frac{1}{a_1}}E_{21}(M)_1{\displaystyle \frac{[2l+1]}{[2l]}}E_{32}E_{21}(M)\right\},`$ $`(M)_3`$ $``$ $`(M)_{1,3}=a_3\left\{{\displaystyle \frac{1}{a_1}}E_{43}(M)_1{\displaystyle \frac{[2l^{}+1]}{[2l^{}]}}E_{32}E_{43}(M)\right\},`$ $`(M)_4`$ $``$ $`(M)_{1,4}=a_4\{{\displaystyle \frac{1}{a_1}}E_{21}E_{43}(M)_1{\displaystyle \frac{1}{a_2}}E_{43}(M)_2{\displaystyle \frac{1}{a_3}}E_{21}(M)_3`$ $`{\displaystyle \frac{[2l+1][2l^{}+1]}{[2l][2l^{}]}}E_{32}E_{21}E_{43}(M)\},`$ $`(M)_5`$ $``$ $`(M)_{2,1}=a_5\|0,0,1,1;(M)a_5E_{42}E_{32}(M),`$ $`(M)_6`$ $``$ $`(M)_{2,2}=a_6\left\{{\displaystyle \frac{1}{a_5}}E_{21}(M)_5{\displaystyle \frac{[2l+2]}{[2l]}}E_{42}E_{32}E_{21}(M)\right\},`$ $`(M)_7`$ $``$ $`(M)_{2,3}=a_7\{{\displaystyle \frac{1}{a_5}}E_{21}^2(M)_5{\displaystyle \frac{1}{a_6}}{\displaystyle \frac{[2][2l+1]}{[2l]}}E_{21}(M)_6`$ $`{\displaystyle \frac{[2l+1][2l+2]}{[2l][2l1]}}E_{42}E_{32}E_{21}^2(M)\},`$ $`(M)_8`$ $``$ $`(M)_{2,4}=a_8\|0,1,0,1;(M)a_8E_{31}E_{32}(M),`$ $`(M)_9`$ $``$ $`(M)_{2,5}=a_9\left\{{\displaystyle \frac{1}{a_8}}E_{43}(M)_8{\displaystyle \frac{[2l^{}+2]}{[2l^{}]}}E_{31}E_{32}E_{43}(M)\right\},`$ $`(M)_{10}`$ $``$ $`(M)_{2,6}=a_{10}\{{\displaystyle \frac{1}{a_8}}E_{43}^2(M)_8{\displaystyle \frac{1}{a_9}}{\displaystyle \frac{[2][2l^{}+1]}{[2l^{}]}}E_{43}(M)_9`$ $`{\displaystyle \frac{[2l^{}+1][2l^{}+2]}{[2l^{}][2l^{}1]}}E_{31}E_{32}E_{43}^2(M)\},`$ $`(M)_{11}`$ $``$ $`(M)_{3,1}=a_{11}\|0,1,1,1;(M)a_{11}E_{31}E_{42}E_{32}(M),`$ $`(M)_{12}`$ $``$ $`(M)_{3,2}=a_{12}\left\{{\displaystyle \frac{1}{a_{11}}}E_{21}(M)_{11}{\displaystyle \frac{[2l+1]}{[2l]}}E_{31}E_{42}E_{32}E_{21}(M)\right\},`$ $`(M)_{13}`$ $``$ $`(M)_{3,3}=a_{13}\left\{{\displaystyle \frac{1}{a_{11}}}E_{43}(M)_{11}{\displaystyle \frac{[2l^{}+1]}{[2l^{}]}}E_{31}E_{42}E_{32}E_{43}(M)\right\},`$ $`(M)_{14}`$ $``$ $`(M)_{3,4}=a_{14}\{{\displaystyle \frac{1}{a_{11}}}E_{21}E_{43}(M)_{11}{\displaystyle \frac{1}{a_{12}}}E_{43}(M)_{12}{\displaystyle \frac{1}{a_{13}}}E_{21}(M)_{13}`$ $`{\displaystyle \frac{[2l+1][2l^{}+1]}{[2l][2l^{}]}}E_{31}E_{42}E_{32}E_{21}E_{43}(M)\},`$ $`(M)_{15}`$ $``$ $`(M)_{4,1}=a_{15}\|1,1,1,1;(M)a_{15}E_{41}E_{31}E_{42}E_{32}(M),(3.25)`$ where $`l=\frac{1}{2}(m_{13}m_{23})`$ and $`l^{}=\frac{1}{2}(m_{33}m_{43})`$, while $`a_k=a_k(p,q)`$ are coefficients depending, in general, on $`p`$ and $`q`$. Indeed, $`(M)_k`$ given in (3.25) form a set of all linear independent vectors satisfying the conditions (3.18). Looking at (3.25) we easily identify the highest weights $`[m]_k`$ \| $`[m]_0`$ \| = \| $`[m_{13},m_{23},m_{33},m_{43}]`$, \| \| --- \| --- \| --- \| \| $`[m]_1`$ \| = \| $`[m_{13},m_{23}1,m_{33}+1,m_{43}]`$, \| \| $`[m]_2`$ \| = \| $`[m_{13}1,m_{23},m_{33}+1,m_{43}]`$, \| \| $`[m]_3`$ \| = \| $`[m_{13},m_{23}1,m_{33},m_{43}+1]`$, \| \| $`[m]_4`$ \| = \| $`[m_{13}1,m_{23},m_{33},m_{43}+1]`$, \| \| $`[m]_5`$ \| = \| $`[m_{13},m_{23}2,m_{33}+1,m_{43}+1]`$, \| \| $`[m]_6`$ \| = \| $`[m_{13}1,m_{23}1,m_{33}+1,m_{43}+1]_6`$, \| \| $`[m]_7`$ \| = \| $`[m_{13}2,m_{23},m_{33}+1,m_{43}+1]`$, \| \| $`[m]_8`$ \| = \| $`[m_{13}1,m_{23}1,m_{33}+2,m_{43}]`$, \| \| $`[m]_9`$ \| = \| $`[m_{13}1,m_{23}1,m_{33}+1,m_{43}+1]_9`$, \| \| $`[m]_{10}`$ \| = \| $`[m_{13}1,m_{23}1,m_{33},m_{43}+2]`$, \| \| $`[m]_{11}`$ \| = \| $`[m_{13}1,m_{23}2,m_{33}+2,m_{43}+1]`$, \| \| $`[m]_{12}`$ \| = \| $`[m_{13}2,m_{23}1,m_{33}+2,m_{43}+1]`$, \| \| $`[m]_{13}`$ \| = \| $`[m_{13}1,m_{23}2,m_{33}+1,m_{43}+2]`$, \| \| $`[m]_{14}`$ \| = \| $`[m_{13}2,m_{23}1,m_{33}+1,m_{43}+2]`$, \| \| $`[m]_{15}`$ \| = \| $`[m_{13}2,m_{23}2,m_{33}+2,m_{43}+2]`$. (3.26) \| In the latest formula (3.26), with the exception of $`[m]_6`$ and $`[m]_9`$ where a degeneration is present, we skip the subscript $`k`$ in the r.h.s.. The proofs of (3.25) and (3.26) follow from direct computations. Using the rule (3.11) which now reads $`(m)_k`$ $`=`$ $`\left({\displaystyle \frac{[m_{11}m_{22}]![m_{31}m_{42}]!}{[m_{12}m_{22}]![m_{12}m_{11}]![m_{32}m_{42}]![m_{32}m_{31}]!}}\right)^{1/2}`$ $`\times (E_{21})^{m_{12}m_{11}}(E_{43})^{m_{32}m_{31}}(M)_k(3.11^{})`$ we can find all the basis vectors $`(m)_k`$ : $`(m)_0`$ $`=`$ $`\|0,0,0,0;(m),`$ $`(m)_1`$ $`=`$ $`a_1q\{\left({\displaystyle \frac{[l_{13}l_{11}][l_{33}l_{31}+1]}{[2l+1][2l^{}+1]}}\right)^{1/2}\|1,0,0,0;(m)^{+11}`$ $`q^{l^{}s^{}}\left({\displaystyle \frac{[l_{13}l_{11}][l_{31}l_{43}1]}{[2l+1][2l^{}+1]}}\right)^{1/2}\|0,1,0,0;(m)^{+1131}`$ $`+p^{l+s}\left({\displaystyle \frac{[l_{11}l_{23}][l_{33}l_{31}+1]}{[2l+1][2l^{}+1]}}\right)^{1/2}\|0,0,1,0;(m)`$ $`+p^{l+s}q^{l^{}s^{}}\left({\displaystyle \frac{[l_{11}l_{23}][l_{31}l_{43}1]}{[2l+1][2l^{}+1]}}\right)^{1/2}\|0,0,0,1;(m)^{31}\},`$ $`(m)_2`$ $`=`$ $`a_2q\left({\displaystyle \frac{q}{p}}\right)^{ls1}\left\{\left({\displaystyle \frac{[l_{11}l_{23}][l_{33}l_{31}+1]}{[2l][2l^{}+1]}}\right)^{1/2}\right\|1,0,0,0;(m)^{+11}`$ $`+q^{l^{}s^{}}\left({\displaystyle \frac{[l_{11}l_{23}][l_{31}l_{43}1]}{[2l][2l^{}+1]}}\right)^{1/2}\|0,1,0,0;(m)^{+1131}`$ $`+q^{l+s+1}\left({\displaystyle \frac{[l_{13}l_{11}][l_{33}l_{31}+1]}{[2l][2l^{}+1]}}\right)^{1/2}\|0,0,1,0;(m)`$ $`+q^{l+s+l^{}s^{}+1}\left({\displaystyle \frac{[l_{13}l_{11}][l_{31}l_{43}1]}{[2l][2l^{}+1]}}\right)^{1/2}\|0,0,0,1;(m)^{31}\},`$ $`(m)_3`$ $`=`$ $`a_3\{q\left({\displaystyle \frac{q}{p}}\right)^{l^{}s^{}}\left({\displaystyle \frac{[l_{13}l_{11}][l_{31}l_{43}1]}{[2l+1][2l^{}]}}\right)^{1/2}\|1,0,0,0;(m)^{+11}`$ $`+\left({\displaystyle \frac{q}{p}}\right)^{2l^{}}\left({\displaystyle \frac{[l_{13}l_{11}][l_{33}l_{31}+1]}{[2l+1][2l^{}]}}\right)^{1/2}\|0,1,0,0;(m)^{+1131}`$ $`+qp^{l+s}\left({\displaystyle \frac{q}{p}}\right)^{l^{}s^{}}\left({\displaystyle \frac{[l_{11}l_{23}][l_{31}l_{43}1]}{[2l+1][2l^{}]}}\right)^{1/2}\|0,0,1,0;(m)`$ $`p^{l+s}\left({\displaystyle \frac{q}{p}}\right)^{2l^{}}\left({\displaystyle \frac{[l_{11}l_{23}][l_{33}l_{31}+1]}{[2l+1][2l^{}]}}\right)^{1/2}\|0,0,0,1;(m)^{31}\},`$ $`(m)_4`$ $`=`$ $`a_4\left({\displaystyle \frac{q}{p}}\right)^{ls+l^{}s^{}1}\left\{q\left({\displaystyle \frac{[l_{11}l_{23}][l_{31}l_{43}1]}{[2l][2l^{}]}}\right)^{1/2}\right\|1,0,0,0;(m)^{+11}`$ $`p^{l^{}s^{}}\left({\displaystyle \frac{[l_{11}l_{23}][l_{33}l_{31}+1]}{[2l][2l^{}]}}\right)^{1/2}\|0,1,0,0;(m)^{+1131}`$ $`+q^{l+s+2}\left({\displaystyle \frac{[l_{13}l_{11}][l_{31}l_{43}1]}{[2l][2l^{}]}}\right)^{1/2}\|0,0,1,0;(m)`$ $`q^{l+s+1}p^{l^{}s^{}}\left({\displaystyle \frac{[l_{13}l_{11}][l_{33}l_{31}+1]}{[2l][2l^{}]}}\right)^{1/2}\|0,0,0,1;(m)^{31}\},`$ $`(m)_5`$ $`=`$ $`a_5\left({\displaystyle \frac{q}{p}}\right)^{l^{}s^{}+1}\left\{\left({\displaystyle \frac{[l_{13}l_{11}][l_{13}l_{11}1]}{[2l+1][2l+2]}}\right)^{1/2}\right\|1,1,0,0;(m)^{+11+1131}`$ $`+p^{l+s}\left({\displaystyle \frac{[l_{13}l_{11}][l_{11}l_{23}+1]}{[2l+1][2l+2]}}\right)^{1/2}\|0,1,1,0;(m)^{+1131}`$ $`p^{l+s+1}\left({\displaystyle \frac{[l_{13}l_{11}][l_{11}l_{23}+1]}{[2l+1][2l+2]}}\right)^{1/2}\|1,0,0,1;(m)^{+1131}`$ $`+p^{2(l+s)}\left({\displaystyle \frac{[l_{11}l_{23}][l_{11}l_{23}+1]}{[2l+1][2l+2]}}\right)^{1/2}\|0,0,1,1;(m)^{31}\},`$ $`(m)_6`$ $`=`$ $`{\displaystyle \frac{a_6}{a_5}}\left({\displaystyle \frac{[2l+1][2l+2][l_{13}l_{11}]}{[l_{11}l_{23}+1]}}\right)^{1/2}(m)_5a_6\left({\displaystyle \frac{q}{p}}\right)^{l^{}s^{}+1}{\displaystyle \frac{[2l+2]}{[2l]}}\times `$ $`\times \left\{[l_{13}l_{11}2]\left({\displaystyle \frac{[l_{13}l_{11}1]}{[l_{11}l_{23}+1]}}\right)^{1/2}\right\|1,1,0,0;(m)^{+11+1131}`$ $`+p^{l+s+1}[l_{13}l_{11}1]\|0,1,1,0;(m)^{+1131}`$ $`p^{l+s+2}[l_{13}l_{11}1]\|1,0,0,1;(m)^{+1131}`$ $`+p^{2(ls1)}\left([l_{11}l_{23}][l_{13}l_{11}]\right)^{1/2}\|0,0,1,1;(m)^{31}\},`$ $`(m)_7`$ $`=`$ $`{\displaystyle \frac{a_7}{a_5}}\left({\displaystyle \frac{[2l1][2l][2l+1][2l+2][l_{13}l_{11}1][l_{13}l_{11}]}{[l_{11}l_{23}][l_{11}l_{23}+1]}}\right)^{1/2}(m)_5`$ $`{\displaystyle \frac{a_7}{a_6}}[2][2l+1]\left({\displaystyle \frac{[2l1][l_{13}l_{11}1]}{[2l][l_{11}l_{23}]}}\right)^{1/2}(m)_6`$ $`a_7{\displaystyle \frac{[2l+1][2l+2]}{([2l1][2l])^{1/2}}}\left({\displaystyle \frac{q}{p}}\right)^{l^{}s^{}+1}\left\{{\displaystyle \frac{[l_{13}l_{11}2][l_{13}l_{11}3]}{([l_{11}l_{23}][l_{11}l_{23}+1])^{1/2}}}\right\|1,1,0,0;(m)^{+11+1131}`$ $`+p^{l+s+2}[l_{13}l_{11}2]\left({\displaystyle \frac{[l_{13}l_{11}1]}{[l_{11}l_{23}]}}\right)^{1/2}\|0,1,1,0;(m)^{+1131}`$ $`p^{l+s+3}[l_{13}l_{11}2]\left({\displaystyle \frac{[l_{13}l_{11}1]}{[l_{11}l_{23}]}}\right)^{1/2}\|1,0,0,1;(m)^{+1131}`$ $`+p^{2(ls2)}\left([l_{13}l_{11}1][l_{13}l_{11}]\right)^{1/2}\|0,0,1,1;(m)^{31}\},`$ $`(m)_8`$ $`=`$ $`a_8q^2\left({\displaystyle \frac{q}{p}}\right)^{ls1}\left\{\left({\displaystyle \frac{[l_{33}l_{31}+1][l_{33}l_{31}+2]}{[2l^{}+1][2l^{}+2]}}\right)^{1/2}\right\|1,0,1,0;(m)^{+11}`$ $`+q^{2l^{}}\left({\displaystyle \frac{[l_{33}l_{31}+2][l_{31}l_{43}1]}{[2l^{}+1][2l^{}+2]}}\right)^{1/2}\|1,0,0,1;(m)^{+1131}`$ $`+q^{2l^{}+1}\left({\displaystyle \frac{[l_{33}l_{31}+2][l_{31}l_{43}1]}{[2l^{}+1][2l^{}+2]}}\right)^{1/2}\|0,1,1,0;(m)^{+1131}`$ $`+q^{2(2l^{}+1)}\left({\displaystyle \frac{[l_{31}l_{43}2][l_{31}l_{43}1]}{[2l^{}+1][2l^{}+2]}}\right)^{1/2}\|0,1,0,1;(m)^{+113131}\},`$ $`(m)_9`$ $`=`$ $`{\displaystyle \frac{a_9}{a_8}}\left({\displaystyle \frac{[2l^{}+1][2l^{}+2][l_{33}l_{31}+2]}{[l_{31}l_{43}1]}}\right)^{1/2}(m)_8`$ $`a_9q^2\left({\displaystyle \frac{q}{p}}\right)^{ls1}{\displaystyle \frac{[2l^{}+2]}{[2l^{}]}}\left\{[l_{33}l_{31}]\left({\displaystyle \frac{[l_{33}l_{31}+1]}{[l_{31}l_{43}1]}}\right)^{1/2}\right\|1,0,1,0;(m)^{+11}`$ $`+q^{l^{}s^{}1}[l_{33}l_{31}+1]\|1,0,0,1;(m)^{+1131}`$ $`+q^{l^{}s^{}}[l_{33}l_{31}+1]\|0,1,1,0;(m)^{+1131}`$ $`+q^{2(l^{}s^{})}\left([l_{33}l_{31}+2][l_{31}l_{43}2]\right)^{1/2}\|0,1,0,1;(m)^{+113131}\},`$ $`(m)_{10}`$ $`=`$ $`{\displaystyle \frac{a_{10}}{a_8}}\left({\displaystyle \frac{[2l^{}1][2l^{}][2l^{}+1][2l^{}+2][l_{33}l_{31}+1][l_{33}l_{31}+2]}{[l_{31}l_{43}2][l_{31}l_{43}1]}}\right)^{1/2}(m)_8`$ $`{\displaystyle \frac{a_{10}}{a_9}}[2][2l^{}+1]\left({\displaystyle \frac{[2l^{}1][l_{33}l_{31}+1]}{[2l^{}][l_{31}l_{43}2]}}\right)^{1/2}(m)_9`$ $`a_{10}\left({\displaystyle \frac{q}{p}}\right)^{ls1}{\displaystyle \frac{[2l^{}+1][2l^{}+2]}{[2l^{}1][2l^{}]}}\{q^2[l_{33}l_{31}1][l_{33}l_{31}]\times `$ $`\times \left({\displaystyle \frac{[2l^{}1][2l^{}]}{[l_{31}l_{43}2][l_{31}l_{43}1]}}\right)^{1/2}\|1,0,1,0;(m)^{+11}`$ $`+q^{l^{}s^{}}[l_{33}l_{31}]\left({\displaystyle \frac{[2l^{}1][2l^{}][l_{33}l_{31}+1]}{[l_{31}l_{43}2]}}\right)^{1/2}\|1,0,0,1;(m)^{+1131}`$ $`+q^{l^{}s^{}+1}[l_{33}l_{31}]\left({\displaystyle \frac{[2l^{}1][2l^{}][l_{33}l_{31}+1]}{[l_{31}l_{43}2]}}\right)^{1/2}\|0,1,1,0;(m)^{+1131}`$ $`+q^{2(l^{}s^{})}\left([2l^{}1][2l^{}][l_{33}l_{31}+1][l_{33}l_{31}+2]\right)^{1/2}\|0,1,0,1;(m)^{+113131}\},`$ $`(m)_{11}`$ $`=`$ $`a_{11}\left({\displaystyle \frac{q}{p}}\right)^{ls+l^{}s^{}+1}\left\{p\left({\displaystyle \frac{[l_{13}l_{11}1][l_{33}l_{31}+2]}{[2l+1][2l^{}+1]}}\right)^{1/2}\right\|1,1,1,0;(m)^{+11+1131}`$ $`+p^{l+s+2}\left({\displaystyle \frac{[l_{11}l_{23}+1][l_{33}l_{31}+2]}{[2l+1][2l^{}+1]}}\right)^{1/2}\|1,0,1,1;(m)^{+1131}`$ $`+q^{l^{}s^{}+1}p\left({\displaystyle \frac{[l_{13}l_{11}1][l_{31}l_{43}2]}{[2l+1][2l^{}+1]}}\right)^{1/2}\|1,1,0,1;(m)^{+11+113131}`$ $`+q^{l^{}s^{}+2}p^{l+s+1}\left({\displaystyle \frac{[l_{11}l_{23}+1][l_{31}l_{43}2]}{[2l+1][2l^{}+1]}}\right)^{1/2}\|0,1,1,1;(m)^{+113131}\},`$ $`(m)_{12}`$ $`=`$ $`{\displaystyle \frac{a_{12}}{a_{11}}}\left({\displaystyle \frac{[2l][2l+1][l_{13}l_{11}1]}{[l_{11}l_{23}+1]}}\right)^{1/2}(m)_{11}`$ $`a_{12}p\left({\displaystyle \frac{q}{p}}\right)^{ls+l^{}s^{}}{\displaystyle \frac{[2l+1]}{[2l]}}\{[l_{13}l_{11}2]\left({\displaystyle \frac{[2l][l_{33}l_{31}+2]}{[l_{11}l_{23}+1][2l^{}+1]}}\right)^{1/2}\times `$ $`\times \|1,1,1,0;(m)^{+11+1131}`$ $`+p^{l+s+2}\left({\displaystyle \frac{[2l][l_{13}l_{11}1][l_{33}l_{31}+2]}{[2l^{}+1]}}\right)^{1/2}\|1,0,1,1;(m)^{+1131}`$ $`+q^{l^{}s^{}+1}[l_{13}l_{11}2]\left({\displaystyle \frac{[2l][l_{31}l_{43}2]}{[[l_{11}l_{23}+1]][2l^{}+1]}}\right)^{1/2}\|1,1,0,1;(m)^{+11+113131}`$ $`+q^{l^{}s^{}+2}p^{l+s+1}\left({\displaystyle \frac{[2l][l_{13}l_{11}1][l_{31}l_{43}2]}{[2l^{}+1]}}\right)^{1/2}\|0,1,1,1;(m)^{+113131}\},`$ $`(m)_{13}`$ $`=`$ $`{\displaystyle \frac{a_{13}}{a_{11}}}\left({\displaystyle \frac{[2l^{}][2l^{}+1][l_{33}l_{31}+2]}{[l_{31}l_{43}2]}}\right)^{1/2}(m)_{11}`$ $`a_{13}p\left({\displaystyle \frac{q}{p}}\right)^{ls+l^{}s^{}}{\displaystyle \frac{[2l^{}+1]}{[2l^{}]}}\{[l_{33}l_{31}+1]\left({\displaystyle \frac{[2l^{}][l_{13}l_{11}1]}{[2l+1][l_{31}l_{43}2]}}\right)^{1/2}\times `$ $`\times \|1,1,1,0;(m)^{+11+1131}`$ $`+p^{l+s+1}[l_{33}l_{31}+1]\left({\displaystyle \frac{[l_{11}l_{23}+1][2l^{}]}{[2l+1][l_{31}l_{43}2]}}\right)^{1/2}\|1,0,1,1;(m)^{+1131}`$ $`+q^{l^{}s^{}}\left({\displaystyle \frac{[l_{13}l_{11}1][l_{33}l_{31}+2][2l^{}]}{[2l+1]}}\right)^{1/2}\|1,1,0,1;(m)^{+11+113131}`$ $`+q^{l^{}s^{}+1}p^{l+s}\left({\displaystyle \frac{[l_{11}l_{23}+1][l_{33}l_{31}+2][2l^{}]}{[2l+1]}}\right)^{1/2}\|0,1,1,1;(m)^{+113131}\},`$ $`(m)_{14}`$ $`=`$ $`{\displaystyle \frac{a_{14}}{a_{11}}}\left({\displaystyle \frac{[2l][2l+1][l_{13}l_{11}1][2l^{}][2l^{}+1][l_{33}l_{31}+2]}{[l_{11}l_{23}+1][l_{31}l_{43}2]}}\right)^{1/2}(m)_{11}`$ $`{\displaystyle \frac{a_{14}}{a_{12}}}\left({\displaystyle \frac{[2l^{}][2l^{}+1][l_{33}l_{31}+2]}{[l_{31}l_{43}2]}}\right)^{1/2}(m)_{12}`$ $`{\displaystyle \frac{a_{14}}{a_{13}}}\left({\displaystyle \frac{[2l][2l+1][l_{13}l_{11}1]}{[l_{11}l_{23}+1]}}\right)^{1/2}(m)_{13}`$ $`a_{14}p\left({\displaystyle \frac{q}{p}}\right)^{ls+l^{}s^{}1}{\displaystyle \frac{[2l+1][2l^{}+1]}{[2l][2l^{}]}}\{[l_{13}l_{11}2][l_{33}l_{31}+1]\times `$ $`\times \left({\displaystyle \frac{[2l][2l^{}]}{[l_{11}l_{23}+1][l_{31}l_{43}2]}}\right)^{1/2}\|1,1,1,0;(m)^{+11+1131}`$ $`+p^{l+s+2}[l_{33}l_{31}+1]\left({\displaystyle \frac{[l_{13}l_{11}1][2l][2l^{}]}{[l_{31}l_{43}2]}}\right)^{1/2}\|1,0,1,1;(m)^{+1131}`$ $`+q^{l^{}s^{}}[l_{13}l_{11}2]\left({\displaystyle \frac{[2l][2l^{}][l_{33}l_{31}+2]}{[l_{11}l_{23}+1]}}\right)^{1/2}\|1,1,0,1;(m)^{+11+113131}`$ $`+q^{l^{}s^{}+1}p^{l+s+1}\left([2l][l_{13}l_{11}1][l_{33}l_{31}+2][2l^{}]\right)^{1/2}\|0,1,1,1;(m)^{+113131}\},`$ $`(m)_{15}`$ $`=`$ $`a_{15}(m)(3.27)`$ where $`l=\frac{1}{2}(m_{13}m_{23})`$, $`s=m_{11}\frac{1}{2}(m_{13}+m_{23})`$, $`l^{}=\frac{1}{2}(m_{33}m_{43})`$ and $`s^{}=m_{31}\frac{1}{2}(m_{33}+m_{43})`$, while $`(m)_k^{\pm ij}`$ is a GZ basis vector obtained from $`(m)_k`$ with replacing the element $`m_{ij}`$ by $`m_{ij}\pm 1`$. We can write down the coefficients in (3.27) all in terms of $`l`$, $`s`$, $`l^{}`$ and $`s^{}`$ only but here we leave them partially expressed in terms of $`l_{ij}`$ and $`l_{ij}^{}`$. From (3.27) we can immediately find all the (local) lowest weight vectors $`(M)_k^V`$ which, by definition, are annihilated by $`E_{21}`$ and $`E_{43}`$. Let us remind again that every firdirmod $`V_k^{p,q}`$ on a level $`\eta `$, spanned on linear combinations of $`\|\theta _1,\theta _2,\theta _3,\theta _4;(m)`$ in (3.14) with a fixed $`_{i=1}^4\theta _i\eta `$ is a linear space of homogeneous polynomials of a definite power $`\eta `$ in the negative odd generators $`E_{ij}`$ ($`1j<3i4`$) acting on $`(m)V_0^{p,q}([m])`$. Taking into account all results obtained above we have proved the following assertion Proposition 3 : Every $`U_{p,q}[gl(2/2)_0]`$–fidirmod $`V_k^{p,q}`$ in decomposition (3.16) is characterized by a highest weight $`[m]_k`$ given in (3.25) and is spanned by a GZ basis $`(m)_k`$ given in (3.27). The latest formula (3.27), in fact, represents a way in which the reduced basis is expressed in terms of the induced basis and vas versa it is not a problem for us to find the inverse relation between these bases (see the Appendix). For further convenience the vectors $`(m)_{\stackrel{~}{k}}(m)_k`$ (for $`k=6,7,9,10,12,13`$ and 14) are partially given via other $`(m)_k`$ which are completely expressed in terms of $`\|\theta _1,\theta _2,\theta _3,\theta _4;(m)`$. It is not difficult to write down the explicit decompositions of these $`(m)_{\stackrel{~}{k}}`$ in the induced basis. But here we prefer the expressions in (3.27) which are more compact and more convenient for finding the inverse relation between two bases and matrix elements of odd generators. Now we are ready to canculate the matrix elements of the generators $`E_{ij}`$. It is sufficient to canculate the matrix elements of the Cartan–Chevalley generators only, since $`U_{p,q}[gl(2/2)]`$ can be generated by these generators and any its representation in some basis is completely defined by their actions on the same basis. For the even generators which do not shift the $`\eta `$’s we readily have $`E_{11}(m)_k`$ $`=`$ $`(l_{11}+1)(m)_k,`$ $`E_{22}(m)_k`$ $`=`$ $`(l_{12}+l_{22}l_{11}+2)(m)_k,`$ $`E_{12}(m)_k`$ $`=`$ $`\left([l_{12}l_{11}][l_{11}l_{22}]\right)^{1/2}(m)_k^{+11},`$ $`E_{21}(m)_k`$ $`=`$ $`\left([l_{12}l_{11}+1][l_{11}l_{22}1]\right)^{1/2}(m)_k^{11},`$ $`L_1(m)_k`$ $`=`$ $`{\displaystyle \frac{1}{2}}(l_{12}l_{22}1)(m)_k,`$ $`E_{33}(m)_k`$ $`=`$ $`(l_{31}+1)(m)_k,`$ $`E_{44}(m)_k`$ $`=`$ $`(l_{32}+l_{42}l_{31}+2)(m)_k,`$ $`E_{34}(m)_k`$ $`=`$ $`\left([l_{32}l_{31}][l_{31}l_{42}]\right)^{1/2}(m)_k^{+31},`$ $`E_{43}(m)_k`$ $`=`$ $`\left([l_{32}l_{31}+1][l_{31}l_{42}1]\right)^{1/2}(m)_k^{31},`$ $`L_3(m)_k`$ $`=`$ $`{\displaystyle \frac{1}{2}}(l_{32}l_{42}1)(m)_k.(3.28)`$ As the matrix elements of $`E_{23}`$ and $`E_{32}`$ are very long expressions we only explain here how to find them. By construction a reduced basis vector $`(m)_k`$ in (3.27) belonging to a fidirmod $`V_k^{p,q}`$ on a level $`\eta `$ is a homogeneous polynomial of a power $`\eta `$ in odd generators $`E_{ij},1j<3i4`$ acting on $`(m)V_0^{p,q}`$. Under the action of $`E_{23}`$ (or $`E_{32}`$, respectively) this vector $`(m)_k`$ is shifted to other fidirmods $`V_k^{}^{p,q}`$ on the previous level $`\eta 1`$ (or on the next level $`\eta +1`$, respectively), i.e., we get on the r.h.s of $`E_{23}(m)_k`$ (or $`E_{32}(m)_k`$, respectively) a homogeneous polynomial of a degree $`\eta 1`$ (or $`\eta +1`$, respectively). Using the inverse relations (A.1) we can express the latter polynomials obtained in terms of the reduced basis, that is we get matrix elements of $`E_{23}`$ and $`E_{32}`$ in this basis. It is a standard way to find matrix elements but in practice we can use a trick making calculations simpler. Since $`E_{23}`$ commutes with $`E_{21}`$ and $`E_{43}`$ we first calculate the action of $`E_{23}`$ on the highest vectors only and then apply (3.11) to find all matrix elements of this generator on arbitrary $`(m)_k`$. It is less complicated to compute matrix elements of $`E_{32}`$ in the standard way but we can apply a similar trick, namely, we first calculate the action of $`E_{23}`$ (which commutes with $`E_{12}`$ and $`E_{34}`$) on the lowest weight vectors and then apply the rule inverse to (3.11). It can be shown that the representations constructed contain all finite–dimensional irreducible representations of $`U_{p,q}[gl(2/2)]`$ classified as typical or nontypical representations which are subjects of next investigations. IV. Conclusion We have considered the two–parametric quantum deformations $`U_{p,q}[gl(2/2)]`$ and described in detail a method for constructing its finite–dimensional representations. The representations constructed can be decomposed into finite–dimensional irreducible representations of the even subalgebra $`U_{p,q}[gl(2/2)_0]`$ and therefore can be given in bases of the latter. Using Poincaré-Birkhoff-Witt theorem and the induced representation method we constructed the induced basis of the induced module $`W^{p,q}`$. This basis, however, does not allow a clear description of a decomposition of $`W^{p,q}`$ into $`U_{p,q}[gl(2/2)_0]`$–fidirmods. It was the reason the reduced basis was introduced. The latter basis representing a union of GZ bases of the even subalgebra $`U_{p,q}[gl(2/2)_0]`$ according to the branching rule $`U_{p,q}[gl(2/2)]U_{p,q}[gl(2/2)_0]gl(1)gl(1)`$ is refered to as Quasi–GZ basis. This step is intermediate but of independent interest. Having these two bases, the induced and the reduced ones, and the relations between them we can find all matrix elements of finite–dimensional representations of $`U_{p,q}[gl(2/2)]`$. It turn out that the representations constructed contain all finite–dimensional irreducible representations of $`U_{p,q}[gl(2/2)]`$ and can be classified into typical and nontypical representations which are subjects of later papers. Looking at the basis transformations and the matrix elements we observe, even at generic deformation parameters, some ”anomalies” which are canceled out at $`p=q`$. It means that the finite–dimensional representations of the two-parametric quantum superalgebra $`U_q[gl(2/2)]`$ are not simply trivial deformations from those of the classical Lie superalgebra $`gl(2/2)`$ in the sense that they can not be found from classical analogues by putting quantum deformation brackets in appropriate places unlike many cases of one–parametric deformations. For example, the expressions $$\frac{q}{p}[2l][l_{13}l_{11}][2l+1][l_{13}l_{11}1]$$ $`(5.1)`$ and $$\frac{q}{p}[2l^{}][l_{33}l_{31}][2l^{}+1][l_{33}l_{31}1]$$ $`(5.2)`$ appearing in the basis transformations and matrix elements can be written in the forms: $$(\frac{q}{p}1)[2l][l_{13}l_{11}]+\left(\frac{q}{p}\right)^{l_{13}l_{11}1}[l_{11}l_{23}]$$ $`(5.1^{})`$ and $$(\frac{q}{p}1)[2l^{}][l_{33}l_{31}]+\left(\frac{q}{p}\right)^{l_{33}l_{31}1}[l_{31}l_{43}],$$ $`(5.2^{})`$ respectively. At $`p=q`$ the latest expressions become $`[l_{11}l_{23}]`$ and $`[l_{31}l_{43}]`$, respectively, exactly as in the one–parametric case . We hope that it is not very difficult to extend the present method to the case of one or both deformation parameters being roots of unity. For conclusion, let us emphasize that our method has an advantage that it avoids the use of the Clebsch–Gordan coefficients which are not always known, especially for higher rank (classical and quantum) groups and multi–parametric deformations. Acknowledgements I would like to thank the Nishina memorial foundation for financial support and the Department of Physics, Chuo University, Tokyo, Japan for warm hospitality. Fruitful discussions with K. Furuta, T. Inami and other members of the Theory Group of the Department of Physics, Chuo University, are also hereby acknowledged. This work was supported in part by the National Research Programme for Natural Sciences of Vietnam under grant number KT – 04.1.1. Appendix The induced basis (4.20) is expressed in terms of the reduced basis through the following inverse relation $`\|1,0,0,0;(m)`$ $`=`$ $`{\displaystyle \frac{1}{a_1}}q^{l+s1}p^{l^{}s^{}}\left({\displaystyle \frac{[l_{13}l_{11}+1][l_{33}l_{31}+1]}{[2l+1][2l^{}+1]}}\right)^{1/2}(m)_1^{11}`$ $`{\displaystyle \frac{1}{a_2}}{\displaystyle \frac{q^{l+s1}p^{l^{}s^{}1}}{[2l+1]}}\left({\displaystyle \frac{[2l][l_{11}l_{23}1][l_{33}l_{31}+1]}{[2l^{}+1]}}\right)^{1/2}(m)_2^{11}`$ $`{\displaystyle \frac{1}{a_3}}{\displaystyle \frac{q^{l+s}p^{l^{}s^{}}}{[2l^{}+1]}}\left({\displaystyle \frac{[l_{13}l_{11}+1][l_{31}l_{43}1][2l^{}]}{[2l+1]}}\right)^{1/2}(m)_3^{11}`$ $`+{\displaystyle \frac{1}{a_4}}{\displaystyle \frac{q^{l+s}p^{l^{}s^{}1}}{[2l+1][2l^{}+1]}}\left([2l][l_{11}l_{23}1][2l^{}][l_{31}l_{43}1]\right)^{1/2}(m)_4^{11},`$ $`\|0,1,0,0;(m)`$ $`=`$ $`{\displaystyle \frac{1}{a_1}}q^{l+s}\left({\displaystyle \frac{[l_{13}l_{11}+1][l_{31}l_{43}]}{[2l+1][2l^{}+1]}}\right)^{1/2}(m)_1^{11+31}`$ $`{\displaystyle \frac{1}{a_2}}{\displaystyle \frac{q^{l+s}}{p[2l+1]}}\left({\displaystyle \frac{[2l][l_{11}l_{23}1][l_{31}l_{43}]}{[2l^{}+1]}}\right)^{1/2}(m)_2^{11+31}`$ $`+{\displaystyle \frac{1}{a_3}}\left({\displaystyle \frac{p}{q}}\right)^{l^{}s^{}1}{\displaystyle \frac{q^{l+s}}{[2l^{}+1]}}\left({\displaystyle \frac{[l_{13}l_{11}+1][2l^{}][l_{33}l_{31}]}{[2l+1]}}\right)^{1/2}(m)_3^{11+31}`$ $`{\displaystyle \frac{1}{a_4}}\left({\displaystyle \frac{p}{q}}\right)^{l^{}s^{}2}{\displaystyle \frac{q^{l+s1}}{[2l+1][2l^{}+1]}}\left({\displaystyle \frac{[2l][l_{11}l_{23}1][2l^{}]}{[l_{33}l_{31}}}\right)^{1/2}(m)_4^{11+31},`$ $`\|0,0,1,0;(m)`$ $`=`$ $`{\displaystyle \frac{1}{a_1}}q^1p^{l^{}s^{}}\left({\displaystyle \frac{[l_{11}l_{23}][l_{33}l_{31}+1]}{[2l+1][2l^{}+1]}}\right)^{1/2}(m)_1`$ $`+{\displaystyle \frac{1}{a_2}}\left({\displaystyle \frac{p}{q}}\right)^{ls}{\displaystyle \frac{p^{l^{}s^{}1}}{[2l+1]}}\left({\displaystyle \frac{[2l][l_{13}l_{11}][l_{33}l_{31}+1]}{[2l^{}+1]}}\right)^{1/2}(m)_2`$ $`+{\displaystyle \frac{1}{a_3}}{\displaystyle \frac{p^{l^{}s^{}}}{[2l^{}+1]}}\left({\displaystyle \frac{[l_{11}l_{23}][2l^{}][l_{31}l_{43}1]}{[2l+1]}}\right)^{1/2}(m)_3`$ $`+{\displaystyle \frac{1}{a_4}}\left({\displaystyle \frac{p}{q}}\right)^{ls1}{\displaystyle \frac{p^{l^{}s^{}}}{[2l+1][2l^{}+1]}}\times `$ $`\times \left([2l][l_{13}l_{11}][2l^{}][l_{31}l_{43}1]\right)^{1/2}(m)_4,`$ $`\|0,0,0,1;(m)`$ $`=`$ $`{\displaystyle \frac{1}{a_1}}\left({\displaystyle \frac{[l_{11}l_{23}][l_{31}l_{43}+1]}{[2l+1][2l^{}+1]}}\right)^{1/2}(m)_1^{+31}`$ $`{\displaystyle \frac{1}{a_2}}\left({\displaystyle \frac{p}{q}}\right)^{ls1}{\displaystyle \frac{1}{[2l+1]}}\left({\displaystyle \frac{[2l][l_{13}l_{11}][l_{31}l_{43}]}{[2l^{}+1]}}\right)^{1/2}(m)_2^{+31}`$ $`{\displaystyle \frac{1}{a_3}}\left({\displaystyle \frac{p}{q}}\right)^{l^{}s^{}1}{\displaystyle \frac{1}{[2l^{}+1]}}\left({\displaystyle \frac{[l_{11}l_{23}][2l^{}][l_{33}l_{31}]}{[2l+1]}}\right)^{1/2}(m)_3^{+31}`$ $`{\displaystyle \frac{1}{a_4}}\left({\displaystyle \frac{p}{q}}\right)^{ls+l^{}s^{}2}{\displaystyle \frac{1}{[2l+1][2l^{}+1]}}\times `$ $`\times \left({\displaystyle \frac{[2l][2l^{}]}{[l_{13}l_{11}][l_{31}l_{43}]}}\right)^{1/2}(m)_4^{+31},`$ $`\|1,1,0,0;(m)`$ $`=`$ $`{\displaystyle \frac{1}{a_5}}q^{2(l+s)}\left({\displaystyle \frac{p}{q}}\right)^{l^{}s^{}}\left({\displaystyle \frac{[l_{13}l_{11}+1][l_{13}l_{11}+2]}{[2l+1][2l+2]}}\right)^{1/2}(m)_5^{1111+31}`$ $`+{\displaystyle \frac{1}{a_6}}\left({\displaystyle \frac{p}{q}}\right)^{l^{}s^{}3}{\displaystyle \frac{q^{2(ls)3}}{[2l+2]}}(p^2q^{2l1}+(pq)[2l1])\times `$ $`\times \left([l_{11}l_{23}1][l_{13}l_{11}1]\right)^{1/2}(m)_6^{1111+31}`$ $`{\displaystyle \frac{1}{a_7}}\left({\displaystyle \frac{p}{q}}\right)^{l^{}s^{}3}{\displaystyle \frac{q^{2(ls+1)}}{[2][2l+1][2l+2]}}\times `$ $`\times \left([2l][2l1][l_{11}l_{23}2][l_{11}l_{23}1]\right)^{1/2}(m)_7^{1111+31},`$ $`E^{(2)}`$ $``$ $`\|0,1,1,0;(m)p\|1,0,0,1;(m)`$ $`=`$ $`{\displaystyle \frac{1}{a_5}}q^{l+s+1}\left({\displaystyle \frac{p}{q}}\right)^{l^{}s^{}+1}[2]\left({\displaystyle \frac{[l_{13}l_{11}+1][l_{11}l_{23}]}{[2l+1][2l+2]}}\right)^{1/2}(m)_5^{11+31}`$ $`+{\displaystyle \frac{1}{a_6}}\left({\displaystyle \frac{p}{q}}\right)^{ls+l^{}s^{}2}{\displaystyle \frac{q^{l+s}}{[2l+2]}}\{[2][2l1][l_{13}l_{11}]`$ $`[2l](p^2[l_{13}l_{11}]+q^1[l_{13}l_{11}1])\}(m)_6^{11+31}`$ $`+{\displaystyle \frac{1}{a_7}}\left({\displaystyle \frac{p}{q}}\right)^{ls+l^{}s^{}2}{\displaystyle \frac{q^{l+s}}{[2l+1][2l+2]}}\times `$ $`\times \left([2l][2l1][l_{13}l_{11}][l_{11}l_{23}1]\right)^{1/2}(m)_7^{11+31},`$ $`\|0,0,1,1;(m)`$ $`=`$ $`{\displaystyle \frac{1}{a_5}}\left({\displaystyle \frac{p}{q}}\right)^{l^{}s^{}}\left({\displaystyle \frac{[l_{11}l_{23}][l_{11}l_{23}+1]}{[2l+1][2l+2]}}\right)^{1/2}(m)_5^{+31}`$ $`+{\displaystyle \frac{1}{a_6}}\left({\displaystyle \frac{p}{q}}\right)^{2(ls)+l^{}s^{}4}\{{\displaystyle \frac{p}{q}}[2l][l_{13}l_{11}2][2l1][l_{13}l_{11}1]\}\times `$ $`\times {\displaystyle \frac{1}{[2l+2]}}\left({\displaystyle \frac{[l_{13}l_{11}]}{[l_{11}l_{23}]}}\right)^{1/2}(m)_6^{+31}`$ $`{\displaystyle \frac{1}{a_7}}\left({\displaystyle \frac{p}{q}}\right)^{2(ls)+l^{}s^{}4}{\displaystyle \frac{\left([2l][2l1][l_{13}l_{11}1][l_{13}l_{11}]\right)^{1/2}}{[2][2l+1][2l+2]}}(m)_7^{+31},`$ $`\|1,0,1,0;(m)`$ $`=`$ $`{\displaystyle \frac{1}{a_8}}q^1p^{2(l^{}+s^{})}\left({\displaystyle \frac{q}{p}}\right)^{l+s}\left({\displaystyle \frac{[l_{33}l_{31}+1][l_{33}l_{31}+2]}{[2l^{}+1][2l^{}+2]}}\right)^{1/2}(m)_8^{11}`$ $`+{\displaystyle \frac{1}{a_9}}q^1p^{2(l^{}s^{})1}\left({\displaystyle \frac{p}{q}}\right)^{ls}{\displaystyle \frac{\left([l_{33}l_{31}+1][l_{31}l_{43}1]\right)^{1/2}}{[2l^{}+2]}}(m)_9^{11}`$ $`{\displaystyle \frac{1}{a_{10}}}\left({\displaystyle \frac{p}{q}}\right)^{ls}{\displaystyle \frac{qp^{2(l^{}s^{})1}}{[2][2l^{}+1][2l^{}+2]}}\times `$ $`\times \left([2l^{}][2l^{}1][l_{31}l_{43}2][l_{31}l_{43}1]\right)^{1/2}(m)_{10}^{11},`$ $`E^{(+2)}`$ $``$ $`\|0,1,1,0;(m)+q^1\|1,0,0,1;(m)`$ $`=`$ $`{\displaystyle \frac{1}{a_8}}q^1p^{l^{}s^{}1}\left({\displaystyle \frac{p}{q}}\right)^{ls+1}[2]\left({\displaystyle \frac{[l_{33}l_{31}+1][l_{31}l_{43}]}{[2l^{}+1][2l^{}+2]}}\right)^{1/2}(m)_8^{11+31}`$ $`+{\displaystyle \frac{1}{a_9}}\left({\displaystyle \frac{p}{q}}\right)^{ls+l^{}s^{}2}{\displaystyle \frac{p^{l^{}s^{}}}{q[2l^{}+2]}}\times `$ $`\times \left\{[2][2l1][l_{33}l_{31}][2l^{}]\left(q^2[l_{33}l_{31}]+p[l_{33}l_{31}1]\right)\right\}(m)_9^{11+31}`$ $`+{\displaystyle \frac{1}{a_{10}}}\left({\displaystyle \frac{p}{q}}\right)^{ls+l^{}s^{}2}{\displaystyle \frac{p^{l^{}s^{}}}{q[2l^{}+1][2l^{}+2]}}\times `$ $`\times \left([2l][2l1][l_{33}l_{31}][l_{31}l_{43}1]\right)^{1/2}(m)_{10}^{11+31},`$ $`\|0,1,0,1;(m)`$ $`=`$ $`{\displaystyle \frac{1}{a_8}}\left({\displaystyle \frac{p}{q}}\right)^{ls}\left({\displaystyle \frac{[l_{31}l_{43}][l_{31}l_{43}+1]}{[2l^{}+1][2l^{}+2]}}\right)^{1/2}(m)_8^{11+31+31}`$ $`+{\displaystyle \frac{1}{a_9}}\left({\displaystyle \frac{p}{q}}\right)^{ls+2(l^{}s^{})4}{\displaystyle \frac{1}{[2l^{}+2]}}\left({\displaystyle \frac{[l_{33}l_{31}]}{[2l^{}+1]}}\right)^{1/2}\times `$ $`\times \left({\displaystyle \frac{p}{q}}[2l^{}][l_{33}l_{31}2][2l^{}1][l_{33}l_{31}1]\right)(m)_9^{11+31+31}`$ $`{\displaystyle \frac{1}{a_{10}}}\left({\displaystyle \frac{p}{q}}\right)^{ls+2(l^{}s^{})4}{\displaystyle \frac{1}{[2][2l^{}+1][2l^{}+2]}}\times `$ $`\times \left([2l^{}][2l^{}1][l_{33}l_{31}1][l_{33}l_{31}]\right)^{1/2}(m)_{10}^{11+31+31},`$ $`\|1,1,1,0;(m)`$ $`=`$ $`q^{l+s2}p^{l^{}s^{}2}\left({\displaystyle \frac{p}{q}}\right)^{ls+l^{}s^{}}\times `$ $`\times \{{\displaystyle \frac{1}{a_{11}}}\left({\displaystyle \frac{[l_{13}l_{11}+1][l_{33}l_{31}+1]}{[2l+1][2l^{}+1]}}\right)^{1/2}(m)_{11}^{1111+31}`$ $`+{\displaystyle \frac{1}{a_{12}}}{\displaystyle \frac{q}{[2l+1]}}\left({\displaystyle \frac{[2l][l_{11}l_{23}1][l_{33}l_{31}+1]}{[2l^{}+1]}}\right)^{1/2}(m)_{12}^{1111+31}`$ $`+{\displaystyle \frac{1}{a_{13}}}{\displaystyle \frac{q^2}{[2l^{}+1]}}\left({\displaystyle \frac{[l_{13}l_{11}+1][2l^{}][l_{31}l_{43}1]}{[2l+1]}}\right)^{1/2}(m)_{13}^{1111+31}`$ $`+{\displaystyle \frac{1}{a_{14}}}{\displaystyle \frac{q^3}{[2l+1][2l^{}+1]}}\left([2l][l_{11}l_{23}1][2l^{}][l_{31}l_{43}1]\right)^{1/2}(m)_{14}^{1111+31}\},`$ $`\|1,0,1,1;(m)`$ $`=`$ $`p^{l^{}s^{}}\left({\displaystyle \frac{p}{q}}\right)^{2(ls)+l^{}s^{}3}\{{\displaystyle \frac{1}{a_{11}}}q^3(q[2l][l_{13}l_{11}]p[2l+1][l_{13}l_{11}1])\times `$ $`\times \left(p[2l^{}+1]q^2[2l^{}]\right)\left({\displaystyle \frac{[l_{33}l_{31}+1]}{[2l+1][2l^{}+1][l_{11}l_{23}]}}\right)^{1/2}(m)_{11}^{11+31}`$ $`{\displaystyle \frac{1}{a_{12}}}{\displaystyle \frac{q^2}{[2l+1]}}(p[2l^{}+1]q^2[2l^{}])\times `$ $`\times \left({\displaystyle \frac{[2l][l_{13}l_{11}][l_{33}l_{31}+1]}{[2l^{}+1]}}\right)^{1/2}(m)_{12}^{11+31}`$ $`{\displaystyle \frac{1}{a_{13}}}{\displaystyle \frac{q^1}{[2l^{}+1]}}(q^1[2l][l_{13}l_{11}]p[2l+1][l_{13}l_{11}1])\times `$ $`\times \left({\displaystyle \frac{[2l^{}][l_{31}l_{43}1]}{[2l+1][l_{11}l_{23}]}}\right)^{1/2}(m)_{13}^{11+31}`$ $`+{\displaystyle \frac{1}{a_{14}}}{\displaystyle \frac{1}{[2l+1][2l^{}+1]}}\left([2l][l_{13}l_{11}][2l^{}][l_{31}l_{43}1]\right)^{1/2}(m)_{14}^{11+31}\},`$ $`\|1,1,0,1;(m)`$ $`=`$ $`q^{l+s3}\left({\displaystyle \frac{p}{q}}\right)^{ls+2(l^{}s^{})4}\{{\displaystyle \frac{1}{a_{11}}}(q[2l]p^2[2l+1])\times `$ $`\times \left(p[2l^{}+1][l_{33}l_{31}1]q[2l^{}][l_{33}l_{31}]\right)`$ $`\times \left({\displaystyle \frac{[l_{13}l_{11}+1]}{[2l+1][2l^{}+1][l_{31}l_{43}]}}\right)^{1/2}(m)_{11}^{1111+31+31}`$ $`{\displaystyle \frac{1}{a_{12}}}{\displaystyle \frac{q}{[2l+1]}}(p[2l^{}+1][l_{33}l_{31}1]q[2l^{}][l_{33}l_{31}]\times )`$ $`\times \left({\displaystyle \frac{[2l][l_{11}l_{23}1]}{[2l^{}+1][l_{31}l_{43}]}}\right)^{1/2}(m)_{12}^{1111+31+31}`$ $`+{\displaystyle \frac{1}{a_{13}}}{\displaystyle \frac{q}{[2l^{}+1]}}(q[2l]p^2[2l+1])\times `$ $`\times \left({\displaystyle \frac{[l_{13}l_{11}+1][2l^{}][l_{33}l_{31}]}{[2l+1]}}\right)^{1/2}(m)_{13}^{1111+31+31}`$ $`+{\displaystyle \frac{1}{a_{14}}}{\displaystyle \frac{q^2}{[2l+1][2l^{}+1]}}\left([2l][l_{11}l_{23}1][2l^{}][l_{33}l_{31}]\right)^{1/2}(m)_{14}^{1111+31+31}\},`$ $`\|0,1,1,1;(m)`$ $`=`$ $`\left({\displaystyle \frac{p}{q}}\right)^{2(ls+l^{}s^{}2)}\{{\displaystyle \frac{1}{a_{11}}}q^2(p[2l+1][l_{13}l_{11}1]q[2l][l_{13}l_{11}])\times `$ $`\times (p[2l^{}+1][l_{33}l_{31}1]q[2l^{}][l_{33}l_{31}])\times `$ $`\times \left({\displaystyle \frac{1}{[2l+1][2l^{}+1][l_{11}l_{23}][l_{31}l_{43}]}}\right)^{1/2}(m)_{11}^{11+31+31}`$ $`+{\displaystyle \frac{1}{a_{12}}}{\displaystyle \frac{q^1}{[2l+1]}}(p[2l^{}+1][l_{33}l_{31}1]q[2l^{}][l_{33}l_{31}])\times `$ $`\times \left({\displaystyle \frac{[2l][l_{13}l_{11}]}{[2l^{}+1][l_{31}l_{43}]}}\right)^{1/2}(m)_{12}^{11+31+31}`$ $`+{\displaystyle \frac{1}{a_{13}}}{\displaystyle \frac{q^1}{[2l^{}+1]}}(p[2l+1][l_{13}l_{11}1]q[2l][l_{13}l_{11}])\times `$ $`\times \left({\displaystyle \frac{[2l^{}][l_{33}l_{31}]}{[2l+1][l_{11}l_{23}]}}\right)^{1/2}(m)_{13}^{11+31+31}`$ $`{\displaystyle \frac{1}{a_{14}}}{\displaystyle \frac{1}{[2l+1][2l^{}+1]}}\left([2l][l_{13}l_{11}][2l^{}][l_{33}l_{31}]\right)^{1/2}(m)_{14}^{11+31+31}\},`$ $`\|1,1,1,1;(m)`$ $`=`$ $`{\displaystyle \frac{1}{a_{15}}}(m)^{1111+31+31}.`$ $`(A.1)`$
warning/0005/hep-ph0005007.html	ar5iv	text	# 1 Introduction ## 1 Introduction The firmly established disappearance of muon neutrinos of cosmic ray origin strongly points toward the existence of neutrino oscillations . The first generation long baseline (LBL) experiments — K2K , MINOS , OPERA and ICANOE — will give a conclusive and unambiguous signature of the oscillation mechanism and will provide the first precise measurements of the parameters governing the oscillation mechanism. MiniBOONE and the LBL programs will test the LSND signal. A neutrino “factory” is based on the decay of muons circulating in a storage ring. Neutrino factories raised the interest of the physics community, since they appear natural follow-ups to the current experimental LBL program and could open the way to future muon colliders. As many studies have shown , the physics potential of such facilities are indeed very vast. An entry-level neutrino factory could test the LSND signal in a background free environment. More importantly, a neutrino factory source would be of sufficiently high intensity to perform very long baseline (transcontinental) experiments. It could also bring the neutrino sector into the realm of precision measurements. The neutrino oscillation phenomenology may be complicated and involve a combination of transitions to $`\nu _e`$, $`\nu _\mu `$ and $`\nu _\tau `$. It is quite evident that future neutrino factories will provide ideal conditions for the neutrino oscillation physics. The neutrino flavor phenomenology could be completely explored: a precise measurement of the mass difference and mixing matrix elements is achievable, a test of the unitarity of the mixing matrix can be performed, a direct detection of Earth matter effects is feasible and CP violation effects could be studied on the leptonic sector . The combination of data from atmospheric neutrino and first generation LBL experiments will provide some preliminary information on the possible sub-leading electron mixing. A neutrino factory can largely improve the sensitivity on this mixing angle. Neutrino sources from muon decays provide clear advantages over neutrino beams from pion decays. The exact neutrino helicity composition is a fundamental tool to study neutrino oscillations. It can be easily selected, since $`\mu ^+e^+\nu _e\overline{\nu _\mu }`$ and $`\mu ^{}e^{}\overline{\nu _e}\nu _\mu `$ can be separately obtained. At a neutrino factory, one could independently study the following flavor transitions: $`\mu ^{}e^{}`$ $`\overline{\nu _e}`$ $`\nu _\mu `$ (3) $`\nu _ee^{}\mathrm{appearance}`$ $`\nu _\mu \mathrm{disappearance},\mathrm{same}\mathrm{sign}\mathrm{muons}`$ $`\nu _\tau \tau ^{}\mathrm{appearance},\mathrm{high}\mathrm{energy}\mathrm{nu}^{}\mathrm{s}`$ $``$ $`\overline{\nu _e}\mathrm{disappearance}`$ (4) $``$ $`\overline{\nu _\mu }\mu ^+\mathrm{appearance},\mathrm{wrong}\mathrm{sign}\mathrm{muons}`$ (5) $``$ $`\overline{\nu _\tau }\tau ^+\mathrm{appearance},\mathrm{high}\mathrm{energy}\mathrm{nu}^{}\mathrm{s}`$ (6) plus 6 other charge conjugate processes initiated from $`\mu ^+`$ decays. The other main advantages over traditional pion beams are (1) the beam is free of systematics and the composition is well known, therefore ideal for disappearance studies; (2) the two neutrinos in the beam have opposite helicities, therefore one can envisage oscillation appearance searches without intrinsic beam backgrounds (3) the muon energy is monochromatic and in principle adaptable (4) muon storage rings allow for multiple baselines, and hence a complete exploration of the $`L/E`$ domains of oscillations (5) because very high intensity will be needed for muon colliders, very intense muon sources will produce very intense neutrino sources, at least a factor 100 more intense than existing high energy facilities. While physics motivations are well understood, it is not yet clear which design of detector would best allow to take full advantage of the neutrino factory beams. We think that, in order to fully explore the neutrino oscillation processes, the detector should be capable of: 1. measuring and identifying all three lepton flavors: electron, muon and tau; 2. measuring the sign of the lepton charge; 3. separating between charged and neutral current interactions. Experimentally, it is a very challenging task to build detectors with (1) mass scales of the order of tens of ktons required for long-baseline experiments, (2) with sufficient granularity to cleanly identify electron and tau leptons and (3) which measure the charge of these leptons. Various solutions have been explored recently. One based on nuclear emulsions and magnetized iron has been discussed in . The main challenge there is to reach the required mass. Large magnetized calorimeters have been discussed in . Such high-density detectors, while “easily” conceived as massive objects, have intrinsically very coarse granularity and only allow the clean measurement of muons. They certainly do not have sufficient power to adequately identify and measure electron or tau charged current states. In this paper, we are motivated by the recent progress made in the direction of the design of the multikton ICANOE detector: accordingly, we consider a 10 kton (fiducial) high granularity low density liquid argon imaging target, complemented with a high-acceptance external muon spectrometer. Thanks to its extremely high granularity target and its excellent calorimetric properties, this design provides the clean identification and measurement of all three neutrino flavors: electron, muons and taus. However, only the sign of the muons reaching the muon spectrometer can be determined<sup>1</sup><sup>1</sup>1The possibility of the measurement of the electron charge will be addressed in a future work.. The aim of this paper is to understand the potentials of a non-magnetized high-granularity target detector which, compared to traditional high density iron calorimeters, brings the measurement of electrons and taus. We study the physics potentials of such a detector configuration for three-family neutrino mixing and for three different baselines (732, 2900 and 7400 km). In Section 2 we summarize three-family neutrino mixing framework, including the treatment of propagation through matter. In Section 3, we explain in details how the event distributions and rates are obtained from a detailed simulation of neutrino interactions and detector effects. In Section 4, we construct for given oscillation scenarios, event variable distributions for various event classes. All events can be subdivided into the electron, the same sign muons, the wrong sign muons and the no lepton samples. In addition, in Section 5 we discuss the possibility to further discriminate final states between $`\nu _e`$ and $`\nu _\mu `$ origins from $`\nu _\tau `$ by means of kinematical analysis of the events. We also address in Section 6 the possibility to tag quasi-elastic events which provide indirect neutrino helicity discrimination however at a large statistical price. Section 7 is devoted to describing our fits of the oscillation parameters. The fits are expected to give back the input reference oscillation parameters and are used to estimate the precision with which we can estimate these parameters. Section 8 presents the results for the important case in which the oscillation effects can be approximated by one mass scale. Since the information about the oscillation parameter is redundantly available in the visible energy distributions of the various event classes, we address in Section 9 the question of the consistency between the different observed oscillations processes. The ability to treat the appearance of electron or tau neutrinos gives good over-constraints on the mixing matrix. Finally, in Section 10, we analyze the three-family scenario including possible $`CP`$-violation. ## 2 Three-family neutrino oscillation framework ### 2.1 Mixing matrix parameterization We consider neutrino oscillations in a three-family scenario: the flavor eigenstates $`\nu _\alpha (\alpha =e,\mu ,\tau )`$ are related to the mass eigenstates $`\nu _i^{}(i=1,2,3)`$ by the mixing matrix U $$\nu _\alpha =U_{\alpha i}\nu _i^{}$$ (7) and we parameterize it as: $$U=\left(\begin{array}{ccc}c_{12}c_{13}& s_{12}c_{13}& s_{13}e^{i\delta }\\ s_{12}c_{23}c_{12}s_{13}s_{23}e^{i\delta }& c_{12}c_{23}s_{12}s_{13}s_{23}e^{i\delta }& c_{13}s_{23}\\ s_{12}s_{23}c_{12}s_{13}c_{23}e^{i\delta }& c_{12}s_{23}s_{12}s_{13}c_{23}e^{i\delta }& c_{13}c_{23}\end{array}\right)$$ (8) with $`s_{ij}=\mathrm{sin}\theta _{ij}`$ and $`c_{ij}=\mathrm{cos}\theta _{ij}`$. We confine without loss of generality the mixing angle $`\theta _{13}`$ to values in the interval $`[0,\pi /4]`$, and $`\theta _{12}`$, $`\theta _{23}`$, $`\delta `$ to the interval $`[0,\pi /2]`$. We present results in terms of $`\mathrm{sin}^22\theta _{13}`$, $`\mathrm{sin}^2\theta _{23}`$ and $`\mathrm{sin}^2\theta _{12}`$, all running in the interval $`[0,1]`$. The reason for this choice can be for example seen in Appendix A, where we recall oscillation probabilities in the one mass scale approximation. The oscillation probabilities for $`\nu _e\nu _\mu `$ and $`\nu _e\nu _\tau `$ depend on $`\mathrm{sin}^22\theta _{13}`$ and on the sin and cosine of $`\theta _{23}`$. The angle $`\theta _{23}`$ must span the interval $`[0,\pi /2]`$, however, the $`\theta _{13}`$ can vary within interval $`[0,\pi /4]`$. Note that we consider small $`\theta _{13}`$ angles, so the $`\mathrm{cos}^4\theta _{13}`$ dependence in $`\nu _\mu \nu _\tau `$ is very mild. For $`\delta =0`$ (i.e. $`U`$ is real), the general expression for the three-family neutrino oscillation probability is: $$P(\nu _\alpha \nu _\beta ;E,L)=P(\overline{\nu }_\alpha \overline{\nu }_\beta ;E,L)=\delta _{\alpha \beta }4\underset{j>k}{}J_{\alpha \beta jk}\mathrm{sin}^2\left(\mathrm{\Delta }_{jk}\right)$$ (9) where in natural units $`\mathrm{\Delta }_{jk}\mathrm{\Delta }m_{jk}^2L/4E=(m_j^2m_k^2)L/4E`$, $`E`$ is the neutrino energy, $`L`$ is the neutrino path-length, and the Jarlskog term is $`J_{\alpha \beta jk}=U_{\beta j}U_{\beta k}U_{\alpha j}U_{\alpha k}`$. We naturally assign the mass difference squared $`\mathrm{\Delta }m_{12}^2`$ to explain the solar neutrino deficit and the mass difference squared $`\mathrm{\Delta }m_{32}^2`$ to describe the atmospheric neutrino anomaly. We will take as reference value $`\mathrm{\Delta }m_{32}^2=3.5\times 10^3\mathrm{eV}^2`$. To cover possible ranges of this value, we will also consider two other values $`\mathrm{\Delta }m_{32}^2=5\times 10^3\mathrm{eV}^2`$ and $`\mathrm{\Delta }m_{32}^2=7\times 10^3\mathrm{eV}^2`$. We will always assume maximal (2-3)-mixing $`\mathrm{sin}^22\theta _{23}=1`$. In case of the solar neutrino deficit solution, the values for $`\mathrm{\Delta }m_{12}^2`$ and mixing $`\mathrm{sin}^22\theta _{12}`$ are not uniquely defined by experiments. We will limit ourself to the LMA-MSW solution with parameters $`\mathrm{\Delta }m_{12}^2=1\times 10^4\mathrm{eV}^2`$ and hypothesize a maximal (1-2)-mixing $`\mathrm{sin}^2\theta _{12}=0.5`$. In this paper, we do not consider the LSND result which would force us to include more states with new parameters beyond three-family mixing. This choice implies that we will always work in a situation where $`\|\mathrm{\Delta }m_{21}^2\|<\|\mathrm{\Delta }m_{32}^2\|\|\mathrm{\Delta }m_{31}^2\|`$. In first approximation, the oscillation phenomena governed by the two mass differences decouple and the effects produced by $`\mathrm{\Delta }m_{12}^2`$ are small at high energy for the considered baselines. In the first part of this paper, we will neglect $`\mathrm{\Delta }m_{12}^2`$ effects and work in the so-called “one mass scale approximation”. In a second phase, we will be concerned with $`CP`$-violation effects and will have to include $`\mathrm{\Delta }m_{12}^2`$. For simplicity, we will take $`m_1<m_2<m_3`$ which implies $`\mathrm{\Delta }m_{32}^2>0`$. We recall that neutrino oscillations through matter can be used to distinguish $`\mathrm{\Delta }m_{32}^2>0`$ from $`\mathrm{\Delta }m_{32}^2<0`$. In vacuum, we will express the oscillation probability as a function of the seven following parameters: (a) the three mixing angles $`\theta _{12}`$,$`\theta _{13}`$,$`\theta _{23}`$; (b) the two mass differences squared $`\mathrm{\Delta }m_{12}^2`$, $`\mathrm{\Delta }m_{32}^2`$; (c) the baseline L; (d) the neutrino energy $`E`$. ### 2.2 Matter effects Since we will consider very long distances between neutrino production and detection, this will only be possible in practice for neutrinos traveling inside the Earth. In this case, the neutrino oscillation probabilities will be modified by an additional diagram due to the interaction of electron neutrino with the electrons in the matter. One can maintain the neutrino oscillation formalism derived in vacuum but define effective masses and mixing angles valid in matter. For example, the effective masses will result from the diagonalization of the Hamiltonian: $$U\left(\begin{array}{ccc}m_1^2& 0& 0\\ 0& m_2^2& 0\\ 0& 0& m_3^2\end{array}\right)U^{}+\left(\begin{array}{ccc}D& 0& 0\\ 0& 0& 0\\ 0& 0& 0\end{array}\right)$$ (10) where $$D=2\sqrt{2}G_Fn_eE=7.56\times 10^5eV^2(\frac{\rho }{gcm^3})(\frac{E}{GeV})$$ (11) Here, $`n_e`$ is the electron density and $`\rho `$ the matter density. For anti-neutrinos, we must replace $`D`$ by $`D`$. For $`\rho (gcm^3)E(GeV)40`$, the effective mass parameter $`D`$ is of the order of the mass splitting that is derived from the atmospheric neutrino anomaly. We then expect matter effects to be important. Within the two-family mixing scheme, the modification of the flavor transition in matter is taken into account by the mixing angle in matter $`\theta _m`$, which is: $$\mathrm{sin}^22\theta _m(D)=\frac{\mathrm{sin}^22\theta }{\mathrm{sin}^22\theta +(\frac{D}{\mathrm{\Delta }m^2}\mathrm{cos}2\theta )^2}$$ (12) For neutrinos, a resonance condition will be met when $`D\mathrm{\Delta }m^2\mathrm{cos}2\theta `$ and the oscillation amplitude will reach a maximum. The resonant neutrino energy $`E^{res}`$ is $$E^{res}\frac{1.32\times 10^4\mathrm{cos}2\theta \mathrm{\Delta }m^2(eV^2)}{\rho (g/cm^3)}$$ (13) Rather than two-family mixing, we have adopted throughout this study three-family framework. In this context, we use the analytic expressions for the matter mixing angles and mass eigenvalues calculated in (see Appendix A). The mass eigenvalues in matter $`M_1`$, $`M_2`$ and $`M_3`$ are: $`M_1^2`$ $`=`$ $`m_1^2+{\displaystyle \frac{A}{3}}{\displaystyle \frac{1}{3}}\sqrt{A^23B}S{\displaystyle \frac{\sqrt{3}}{3}}\sqrt{A^23B}\sqrt{1S^2}`$ (14) $`M_2^2`$ $`=`$ $`m_1^2+{\displaystyle \frac{A}{3}}{\displaystyle \frac{1}{3}}\sqrt{A^23B}S+{\displaystyle \frac{\sqrt{3}}{3}}\sqrt{A^23B}\sqrt{1S^2}`$ (15) $`M_3^2`$ $`=`$ $`m_1^2+{\displaystyle \frac{A}{3}}+{\displaystyle \frac{2}{3}}\sqrt{A^23B}S`$ (16) where $`A`$, $`B`$ and $`S`$ are given in the Appendix. For the mixing angles in matter the analytical expressions read: $`\mathrm{sin}^2\theta _{12}^m`$ $`=`$ $`{\displaystyle \frac{(M_2^4\alpha M_2^2+\beta )\mathrm{\Delta }M_{31}^2}{\mathrm{\Delta }M_{32}^2(M_1^4\alpha M_1^2+\beta )\mathrm{\Delta }M_{31}^2(M_2^4\alpha M_2^2+\beta )}}`$ (17) $`\mathrm{sin}^2\theta _{13}^m`$ $`=`$ $`{\displaystyle \frac{M_3^4\alpha M_3^2+\beta }{\mathrm{\Delta }M_{31}^2\mathrm{\Delta }M_{32}^2}}`$ (18) $`\mathrm{sin}^2\theta _{23}^m`$ $`=`$ $`{\displaystyle \frac{G^2s_{23}^2+F^2c_{23}^2+2GFc_{23}s_{23}c_\delta }{G^2+F^2}}`$ (19) where $`\alpha `$, $`\beta `$, $`G`$ and $`F`$ are found in the Appendix. To illustrate matter effects in three-neutrino mixing framework, we show in Figure 1 the values of the mixing angles in matter, plotted as a function of $`D/\mathrm{\Delta }m_{32}^2`$, or equivalently of $`\rho \times E`$. The parameter values in vacuum correspond to our reference values for atmospheric and LMA-MSW solar experiments. The resonant behavior of $`\theta _{13}^m`$ is clearly visible. It gives maximum oscillation at a neutrino energy of about 12 GeV for a density of 3.7 $`gcm^3`$. There is a similar resonant behavior for $`\theta _{12}^m`$ but it occurs at low energy since it is driven by $`\mathrm{\Delta }m_{21}^2`$. For $`D>\mathrm{\Delta }m_{32}^2`$, the angles $`\mathrm{sin}^22\theta _{12}^m`$ and $`\mathrm{sin}^2\theta _{23}^m`$ tend to rise slightly, because the non-vanishing $`\mathrm{\Delta }m_{21}^2`$ splitting removes the degeneracy between muon and tau flavors. Our results have been computed assuming a constant density along the whole neutrino path, and equal to the mean density, obtained integrating over the earth profile . This approximation yields, as shown in e.g. Ref. , similar results to those obtained by numerical integration using the actual Earth’s density profile. The oscillation probability through matter will be a function of eight parameters: (a) the vacuum three mixing angles $`\theta _{12}`$,$`\theta _{13}`$,$`\theta _{23}`$; (b) the vacuum two mass differences squared $`\mathrm{\Delta }m_{12}^2`$, $`\mathrm{\Delta }m_{23}^2`$; (c) the average earth density $`\rho `$; (d) the baseline L; (e) the neutrino energy $`E`$. ## 3 Choice of baseline and event rates The exact parameters of a neutrino factory are not yet completely fixed and realistic scenarios are in the process to be defined. However, based on , we can assume that the muons in the storage ring have an energy $`E_\mu =30`$ GeV and that after one year of operation, the factory should deliver about $`10^{20}`$ “useful” muons decays of both polarities in the straight section pointing towards the far detector location. We base our ultimate reach on $`10^{21}`$ “useful” muons decays. Even an integrated intensity of $`10^{22}`$ might be eventually reachable. We compute the fluxes assuming unpolarized muons and disregarding muon beams divergences within the storage ring. We integrate the expected event rates using a neutrino-nucleon Monte-Carlo generator . The total charged current (CC) cross section is technically subdivided into three parts: the exclusive quasi-elastic scattering channel $`\sigma _{QE}`$ and the inelastic cross section $`\sigma _{inelasic}`$ which includes all other processes except charm production which is included separately. Table 1 summarizes the expected rates for the 10 kton fiducial mass and $`10^{20}`$ muon decays (expected 1 year of operation). $`N_{tot}`$ is the total number of events and $`N_{qe}`$ is the number of quasi-elastic events. Even though our study is site non-specific, the chosen baselines could correspond to the distances between the Laboratori Nazionale del Gran Sasso (LNGS) and neutrino factories at (1) CERN ($`L=732`$ km, $`<\rho _{Earth}>=2.8`$ g/cm<sup>3</sup>), (2) Canary Islands ($`L=2900`$ km, $`<\rho _{Earth}>=3.2`$ g/cm<sup>3</sup>) and (3) Fermilab ($`L=7400`$ km, $`<\rho _{Earth}>=3.7`$ g/cm<sup>3</sup>). ## 4 The four main classes of events Muon identification, charge and momentum measurement provide discrimination between $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ charged current (CC) events. Good $`\nu _e`$ CC versus $`\nu `$ NC discrimination relies on the fine granularity of the target. Finally, the identification of $`\nu _\tau `$ CC events requires a precise measurement of all final state particles. It is natural to classify the events in four classes. We illustrate them for the case of $`\mu ^{}`$ stored in the ring. 1. Right sign muons ($`rs\mu `$): the leading muon has the same charge as those circulating inside the ring. Their origin is from 1. non-oscillated $`\nu _\mu `$ CC 2. $`\nu _\mu \nu _\tau `$ CC, $`\tau ^{}\mu ^{}`$ decays 3. hadron decays in neutral currents. 2. Wrong sign muons ($`ws\mu `$): the leading muon has opposite charge to those circulating inside the ring. Opposite-sign leading muons can only be produced by neutrino oscillations, since there is no component in the beam that could account for them. 1. $`\overline{\nu }_e\overline{\nu }_\mu `$ oscillations 2. $`\overline{\nu }_e\overline{\nu }_\tau `$ oscillations, $`\tau ^+\mu ^+`$ decays 3. hadron decays in neutral currents. 3. Electrons ($`e`$): events with a prompt electron and no primary muon identified. Events with leading electron or positron are produced by the charged-current interactions of the following neutrinos: 1. non-oscillated $`\overline{\nu _e}`$ neutrinos 2. $`\nu _\mu \nu _e`$ oscillations 3. $`\overline{\nu }_e\overline{\nu }_\tau `$ or $`\nu _\mu \nu _\tau `$ oscillations with $`\tau e`$ decays 4. No Lepton ($`0\mathrm{}`$): events corresponding to NC interactions or $`\nu _\tau `$ CC events followed by a hadronic decay of the tau lepton. Events with no leading electrons or muons will be used to study the $`\nu _\mu \nu _\tau `$ oscillations. These events can be produced in 1. neutral current processes 2. $`\overline{\nu }_e\overline{\nu }_\tau `$ or $`\nu _\mu \nu _\tau `$ oscillations with $`\tau hadrons`$ decays The last two classes can only be cleanly studied in a fine granularity detector. The most effective way to fit the oscillation parameters is to study the visible energy distribution of the four classes of events defined above, since assuming the unoscillated spectra are known, they contain direct information on the oscillation probabilities. Of course, for electron or muon charged current events, the visible energy reconstructs the incoming neutrino energy. In the case of neutral currents or the charged current of tau neutrinos, the visible energy is less than the visible energy because of undetected neutrinos in the final state. The information is in this case degraded but can still be used. Our analyses are performed on samples of fully generated Monte-Carlo events, which include proper kinematics of the events, full hadronization of the recoiling jet and proper exclusive polarized tau decays when relevant<sup>2</sup><sup>2</sup>2Our simulation has been bench-marked on the comparison of the kinematic features of the lepton and hadronic jet of real neutrino data accumulated in the NOMAD experiment (see e.g. ).. Nuclear effects, which are taken into account by the FLUKA model , are included as they are important for a proper estimation of the tau kinematical identification. The detector response is included in our analyses using a fast simulation which parameterizes the momentum and angular resolution of the emerging particles, using essentially the following values: electromagnetic shower $`3\%/\sqrt{E}1\%`$, hadronic shower $`20\%/\sqrt{E}5\%`$, and magnetic muon momentum measurement $`20\%`$. Hadron decay background can be quite large in a low density target and could be quite dangerous. Fortunately, it can be easily suppressed by a cut on the muon candidate momentum, $`P_\mu >2`$ GeV, which reduces it to a tolerable level. Figure 2 illustrates the relative background expected for $`\nu _\mu `$ and $`\overline{\nu }_\mu `$ NC and CC processes as a function of the muon momentum. After the cut, the expected contamination for $`\nu _\mu `$ CC events is at the level of $`10^5`$. Real charged current events maintain an efficiency above $`95\%`$. The visible energy is computed as the modulus of the vector sum of the momenta of each visible particle in the event. Figures 666 and 6 show the reconstructed visible energy at the baseline $`L=7400`$km normalized to $`10^{20}\mu `$’s for each event class for a specific oscillation scenario with $`\mathrm{\Delta }m_{32}^2=3.5\times 10^3`$ eV<sup>2</sup>, $`\mathrm{sin}^2\theta _{23}=0.5`$ and $`\mathrm{sin}^22\theta _{13}=0.05`$. The different contributions including backgrounds for each event class have been evidenced in the plots. For example, in Figure 6, the different processes that contribute to the right-sign muon class are unoscillated muons, taus and background events. ## 5 Further classification based on kinematical analysis An efficient identification of $`\nu _\tau `$ induced charged current events requires a precise measurement of all final state particles. Excellent calorimetry allows to take full advantage of the special kinematic features of $`\nu _\tau `$ events. We independently search for the leptonic and hadronic tau decay modes. For the $`\tau l\nu \nu `$ decay mode, the main background comes from $`\nu _l`$ CC. To enhance the separation between $`\tau `$ and background events, we demand the event missing $`P_T`$ to be larger than 0.6 GeV and the transverse momentum of the lepton candidate,$`P_T^l`$, to be smaller than 0.5 GeV. This set of cuts is referred to as “loose cuts” and it will be used to perform a check on appearance/disappearance consistency (see subsection 9). In Table 2, we show that the overall $`\tau `$ efficiency for “loose cuts” is $`40\%`$ for a CC background level of $`5\times 10^3`$. Figure 7 shows the energy spectra for the four event classes after application of these cuts, for $`\nu _\tau `$ CC and other types of events. No energy cut has been applied, but the fact of using the energy spectra in the fit also exploits the difference in energy spectra of $`\nu _\tau `$ CC events. A set of “tight cuts” is also applied, aiming at having one expected background event at the farthest location in case $`10^{20}`$ “useful” muons decays are delivered. As we can see from Table 2, the overall tau efficiency in this case amounts up to $`20\%`$. For hadronic decays the most important source of background correspond to NC events. If we demand a $`P_T^{miss}`$ smaller than 1 GeV and a transverse momentum of the hadron candidate with respect to the total event momentum, $`Q_T`$, larger than 0.5 GeV (“loose cuts”) only $`2\%`$ of the initial background survives for an overall tau efficiency of $`30\%`$. If we require only one NC background event survivor at $`L=7400`$ km (“tight cuts”), the signal efficiency drops to $`6\%`$ due to the stringent $`Q_T`$ requirement imposed. ## 6 Quasi-elastic final states The quasi-elastic process, while rare, is a clean process that allows to separate neutrino from anti-neutrino events, in principle for all neutrino flavors, since $`\nu _{\mathrm{}}+n\mathrm{}^{}+p`$ and $`\overline{\nu }_{\mathrm{}}+p\mathrm{}^++n`$. The recoil proton is easily identifiable within the high-granularity target. This channel is particularly interesting to study oscillation in the electron channel. Starting from negative muons circulating in the storage ring, we look for exclusive electron-proton final states. These provide “background-free” oscillation signals: $`\mu ^{}e^{}`$ $`\overline{\nu _e}`$ $`\nu _\mu `$ (21) $`\nu _e+ne^{}+p`$ $`\nu _\tau +n\tau ^{}+pe^{}\nu \nu +p`$ since $`\overline{\nu _e}+pe^++n`$. The selection of quasi-elastic events is the only way to identify the helicity of neutrino electrons in absence of measurement of the electron charge. ## 7 Oscillation parameters fitting Given the adopted parameterization of the mixing matrix, we have a priori a total of 7 free parameters, which can be represented by the vector: $$\stackrel{}{P}=(\mathrm{\Delta }m_{21}^2,\mathrm{\Delta }m_{32}^2,\mathrm{sin}^2\theta _{12},\mathrm{sin}^22\theta _{13},\mathrm{sin}^2\theta _{23},\delta ,\rho )$$ (22) The values of the parameters governing the oscillations are extracted from a global fit of the visible energy distributions obtained for each event class. The fit is performed with the MINUIT package, and is expected to get back the same values of the parameters, starting from the reference distributions. At each iteration, a different set of parameters is probed, and with the same procedure used to get the reference histograms. For a given polarity $`\lambda `$ of the muons in the storage ring, we compute $`\chi ^2`$’s of the difference between the binned oscillated spectra, which will be function of the parameters, and the reference histograms. We define a $`\chi ^2`$ for each of the four classes of events, i.e. the electrons ($`e`$), the right-sign muon ($`rs\mu `$), the wrong sign muons ($`ws\mu `$) and the no lepton class ($`0\mathrm{}`$): $$\chi _{\lambda ,all}^2=\chi _{\lambda ,e}^2+\chi _{\lambda ,rs\mu }^2+\chi _{\lambda ,ws\mu }^2+\chi _{\lambda ,0\mathrm{}}^2$$ (23) where $$\chi _{\lambda ,c}^2(\stackrel{}{P})=\underset{i}{}\left(\frac{N_i^c(\stackrel{}{P})N_i^c(\stackrel{}{P}_{ref})}{\sigma _{i,c}^2}\right)^2.$$ (24) The sum runs over 25 equally spaced energy bins ($`E_i`$, $`i=1,25`$). Here, $`N_i^c(\stackrel{}{P}_{ref})`$ is the expectation for the reference values and $`N_i^c(\stackrel{}{P})`$ is the “data” obtained for a given set of the oscillation parameters $`\stackrel{}{P}`$. $`\sigma _{i,c}`$ contains both statistical and systematic contributions to the total error. In order to improve the statistical treatment of bins with low statistics, a bin that possesses less than 40 events is assumed to be Poisson distributed and therefore its contribution is computed as $`2(N_i^c(\stackrel{}{P})N_i^c(\stackrel{}{P}_{ref}))+2N_i\mathrm{ln}(N_i^c(\stackrel{}{P})/N_i^c(\stackrel{}{P}_{ref}))`$ (see Ref. ). The systematic error takes into account the uncertainties in the knowledge of the beam, neutrino cross sections and selection efficiencies and we assume it amounts up to $`2\%`$ uncorrelated from bin to bin. We assume that a neutrino factory will operate with alternate runs of opposite muon polarities, therefore eight energy distributions can be fitted simultaneously: $$\chi _{all}^2=\chi _{+,all}^2+\chi _{,all}^2$$ (25) The values of the fitted parameters $`\stackrel{}{P}`$ are obtained minimizing the $`\chi ^2(\stackrel{}{P})`$. It is in practice not always possible to fit all the free parameters, since for some parameter-space regions, the oscillation effects at the chosen baselines and energies can be negligible. In particular, this can be the case for $`\mathrm{\Delta }m_{21}^2`$ and $`\mathrm{sin}^2\theta _{12}`$ which drive the solar oscillations. For values $`\mathrm{\Delta }m_{21}^210^4\mathrm{eV}^2`$, we are insensitive to the “solar” sector. We therefore adopted successive fitting procedures with an increasing number of free parameters. At first, we can simplify the three-family oscillation picture if the oscillations produced by $`\mathrm{\Delta }m_{21}^2`$ can be neglected at the considered baselines and energies. In this case, only the mass difference squared $`\mathrm{\Delta }m_{32}^2`$ and the two mixing angles $`\theta _{13}`$ and $`\theta _{23}`$ are relevant. The oscillation probabilities are given in the Appendix. For example, for $`\nu _e\nu _\mu `$ oscillations, it is: $$P(\nu _e\nu _\mu ,E,L)=\mathrm{sin}^2(2\theta _{13}^m)\mathrm{sin}^2(\theta _{23}^m)\mathrm{\Delta }_{32}^2$$ (26) where $`\mathrm{\Delta }_{32}^2=\mathrm{sin}^2\left((M_3^2M_2^2)L/4E\right)`$. The fit has in this case 4 free parameters, which can be represented by the vector: $$\stackrel{}{P}_{1ms}=(\mathrm{\Delta }m_{32}^2,\mathrm{sin}^22\theta _{13},\mathrm{sin}^2\theta _{23},\rho )$$ (27) The $`\delta `$ phase has disappeared, since within this approximation it becomes unphysical. The results are presented in Section 8, assuming a $`\mathrm{\Delta }m_{32}^2`$ parameter varying in the range of values favoured by current atmospheric data ($`\mathrm{\Delta }m_{32}^2=3.5,5,7\times 10^3\mathrm{eV}^2`$), a maximal (2-3)-mixing $`\mathrm{sin}^2\theta _{23}=0.5`$ and a $`\theta _{13}`$ value compatible with CHOOZ results and recent fits to data ($`\mathrm{sin}^22\theta _{13}=0.05`$). We return to the general three-family scenario in Section 10 where we consider sensitivity to the $`CP`$ violation. In order to compute the precision of the determination of the parameters, we consider two methods: (1) a one-dimensional “scan” of a given parameter; the other variables are left free and mininized at each step; the resp. 1,2,3 sigmas are given by resp. $`\chi _{min}^2+1`$, $`+4`$ and $`+9`$. (2) a two-dimensional “scan” of a two-parameter plane; the other variables are left free and minimized at each point in the plane; the resp. 68%, 90%, 99% C.L. are given by resp. $`\chi _{min}^2+2.3`$, $`+6.0`$ and $`+9.2`$. ## 8 Results for one mass scale approximation ### 8.1 Case of two-family mixing : $`\theta _{13}=0`$ The detection of the dip in the energy distribution of the right-sign muon sample dominates the precision on the measurement of the mixing angle $`\theta _{23}`$ and the mass difference $`\mathrm{\Delta }m_{32}^2`$. This is true provided that the beam energy and baseline are chosen in such a way that the $`\nu _\mu `$ disappearance maximum is visible in the oscillated spectrum. Table 3 summarizes the expected accuracies. With $`10^{20}`$ muon decays of each polarity, precisions of $`12\%`$ are expected in the determination of $`\mathrm{\Delta }m_{32}^2`$ while for the precision on $`\mathrm{sin}^2\theta _{23}`$ is around 10$`\%`$, in agreement with results quoted in . We compare in figure 8 how the precision on the mass difference and mixing angle changes when experimental resolutions and backgrounds are disregarded. We see that our fits are barely affected and only for the cases where the dip is not seen, the instrumental effects and backgrounds spoil at the level of a few per cent the accuracy on the oscillation parameters. ### 8.2 Case of three-family mixing : $`\theta _{13}0`$ In order to obtain the best sensitivity on the mixing angle $`\theta _{13}`$, the search for wrong sign muons is ideal, since it will be a direct signature for $`\nu _e\nu _\mu `$ oscillations. For this kind of study, since we are dealing with the smallest number of signal events, the sensitivity does strongly depend on the ability of rejecting background. Figure 9 shows for $`L=7400`$ km, the sensitivity on $`\theta _{13}`$ for two different muon normalizations ($`10^{20}`$and $`10^{21}`$ muon decays of each polarity). For each pair of values $`(\mathrm{\Delta }m_{32}^2,\theta _{13})`$, the fit was performed leaving $`\theta _{23}`$ free. In the obtaining previous plot, we did not on purpose apply strong background cuts, since we believe that very high rejection powers obtained on paper may not stand the proof of real experimental conditions, with non-gaussian behaviours, tails of distributions etc. To also give the maximum of sensitivity that can be obtained in principle, we also illustrated in figure 9 the effect that a background free environment would have in the expected sensitivity. At 90$`\%`$ C.L., we obtain $`\mathrm{sin}^22\theta _{13}<10^310^4`$ depending on the number of muon decays. This represents two orders of magnitude improvement with respect to quoted sensitivities at CNGS. In case we consider $`10^{21}`$ muon decays and backgrounds are reduced to a negligible level, the obtained sensitivity is consistent with the one quoted in . For a background-free environment gives the best sensitivity, the amount of information added by other event classes (i.e. the electrons) is negligible. #### 8.2.1 Determination of $`\theta _{13}`$ The measurement of $`\theta _{13}`$ can profit from long baselines, since matter effects will enhance the oscillation signal. In presence of backgrounds, it is more favorable to enhance neutrinos signal even at the cost of the suppression of the anti-neutrino oscillations. Table 4 summarizes the expected precision on the measurement of the oscillation parameters. In this case, since for the chosen value of $`sin^22\theta _{13}=0.05`$ the number of signal events is quite large, there is not any more a strong need for a background-free environment. Therefore, the inclusion of other event classes, like the electrons, can help to constrain the oscillation parameters. ### 8.3 Sensitivity to $`\theta _{13}`$ with quasi-elastic events An exclusive way of detecting the effects of a non vanishing $`\theta _{13}`$ is through the appearance of wrong sign electrons. Although, with the assumed detector configuration, there is no ability to directly measure the charge of the leading electrons, there is the possibility of disentangling final state electrons from positrons through the use of quasi-elastic events. We look for electron-proton final states. For example, in the target of ICANOE, a proton can be resolved if its kinetic energy is larger than about 100 MeV corresponding to a range of more than 2 cm. Table 5 shows the expected rates contributing to the electron class before any cut is applied. The background is twofold: (a) quasi-elastic $`\nu _\tau `$ CC events followed by $`\tau e`$, (b) $`\overline{\nu }_e`$ CC with an extra proton from nuclear origin. This last background can be estimated from data themselves studying the reaction $`\overline{\nu }_\mu p\mu ^+n`$. Demanding a back to back electron-proton event topology with a proton kinetic energy in excess of 100 MeV, we estimate the expected background to be less than one event for an overall signal efficiency of $`50\%`$. The quasi-elastic channel provides, in a “background free” environment, 20 to 40 gold-plated events depending on the selected baseline. It is very clean channel but is limited by statistics. ### 8.4 Fit of the average Earth density parameter In matter, the amplitude of neutrino oscillation goes through a maximum for an energy given by equation 13. Since $`\theta _{13}`$ is small, the MSW resonance peak is only a function of $`\rho `$ and $`\mathrm{\Delta }m_{32}^2`$, This can be seen in Figure 11. Since the mass difference is constrained by the disappearance of right-sign muons, $`\rho `$ is well-determined by the energy distribution of wrong sign muons. We extract the density from the fit, leaving it as a free parameter, as well as $`\mathrm{\Delta }m_{32}^2`$, $`\theta _{23}`$ and $`\theta _{13}`$. The precision on the determination of $`\rho `$ depends on the baseline, as shown in table 6, obtained considering $`2\times 10^{21}`$ muon decays. For the longest baseline, where matter effect are large, a precision as good as 2% can be obtained. The influence of $`\rho `$ in our fits is addressed in figure 11. We see that for $`L=7400`$ km and three different muon normalizations, the fact that $`\rho `$ is either considered as a free parameter or fixed during the fit does not influence the accuracy in the determination of the mixing angles. ## 9 Over-constraining the oscillation parameters The information about the oscillation parameter is redundantly available in the visible energy distributions of the various event classes. This allows us to address the question of the consistency between the different observed oscillations processes. Given the high statistical accuracy of the measurement, the consistency can be tested with good accuracy. In the three active neutrino mixing scheme, this implies that $`\mathrm{\Sigma }_{y=e,\mu ,\tau }P(\nu _x\nu _y)`$ should be equal to one for $`x=e,\mu ,\tau `$ and the same holds for anti-neutrinos. Other models can predict different values (i.e., oscillations to sterile neutrinos exist, the sum would be smaller than one). Let us concentrate on the oscillations into $`\tau `$ neutrinos. In the case of negative muons in the ring, they can be originated from $`\nu _\mu \nu _\tau `$ or $`\overline{\nu }_e\overline{\nu }\tau `$ oscillations. The latter case is particularly interesting, since coupling a neutrino factory with a detector with $`\tau `$ identification capabilities is probably the only way to identify and measure such a process. The $`\overline{\nu }_e\overline{\nu }\tau `$ can be revealed experimentally from the presence of $`\tau `$ candidates in the wrong-sign muon sample, due to the opposite helicity of electron and muon neutrinos in the beam. To have a quantitative estimation of the consistency check of the various oscillation modes to $`\tau `$ neutrino, we assign two global normalization factors, $`\alpha `$ and $`\beta `$ to the oscillation probabilities $`P(\nu _\mu \nu _\tau )`$ and $`P(\nu _e\nu _\tau )`$. If no new phenomena occur, these parameters should be exactly one. From a global fit to the visible energy distributions it is possible to extract the values of these parameters, and the precision obtainable on their measurement. To select $`\tau `$ events from the background, still retaining a high efficiency on the signal, we apply a set of loose kinematic cuts (see table 2). Background levels of the order of one event can be reached applying tighter cuts, but in these cases the statistics is too small and the results obtained are slightly worse. Since the $`\tau `$ lepton decays into muons, electrons and hadrons, we expect that a fit to all event classes would result in a remarkable improvement on the precision for the $`\alpha `$ parameter. Figure 12 shows how $`\alpha `$ determination improves as the different event classes are included in the fit. Table 7 shows the expected precisions in the determination of $`\alpha `$ and $`\beta `$. We observe that for $`10^{20}`$ decays, $`\alpha `$ is better determined (accuracy around $`5\%`$) for $`L=2900`$km, however for $`10^{21}`$ muons, the accuracy is about $`1\%`$ regardless of the baseline and the mass difference and therefore $`\nu _\mu `$ oscillations into a sterile neutrino can be largely ruled out. The accuracy on $`\beta `$, and therefore the first experimental evidence for $`\nu _e\nu _\tau `$ is much worst, given the smaller statistics available, since this oscillation probability is smaller than the corresponding $`\nu _\mu \nu _\tau `$ by a factor $`\mathrm{sin}^22\theta _{13}`$, taken to be in this case 0.05. We observe that somewhat better determination exists at $`L=7400`$km since this oscillation mode is largely influenced by matter effects given its dependence on $`\theta _{13}`$. We conclude that a precision at the level of a few per cent on the observation of $`\nu _e\nu _\tau `$ oscillations would require $`O(10^{22})`$ useful muon decays. ## 10 Results for the general three-family scenario and CP-violation Let us consider now a more complex scenario. In this case, the value of the mass difference $`\mathrm{\Delta }m_{12}^2`$ is not any more negligible, and is actually assumed to be $`10^4eV^2`$, one of the highest values compatible with the large mixing angle MSW solution for solar neutrinos . The oscillation does not depend any more on only three parameters, but all four independent angles of the mixing matrix and the two mass differences become important. In the most general case, the phase $`\delta `$ can be different from zero, producing a complex mixing matrix, and thus generating CP violation. We recall that for neutrinos propagating in vacuum, the oscillation probability after a distance $`L`$ for neutrinos can be expressed as: $`P(\nu _\alpha \nu _\beta ;E,L)`$ $`=`$ $`P_{CPeven}(\alpha ,\beta ;E,L)+P_{CPodd}(\alpha ,\beta ;E,L)`$ (28) and for antineutrinos : $`P(\overline{\nu }_\alpha \overline{\nu }_\beta ;E,L)`$ $`=`$ $`P_{CPeven}(\alpha ,\beta ;E,L)P_{CPodd}(\alpha ,\beta ;E,L)`$ (29) As an illustration, the probability for $`\nu _\mu `$ to $`\nu _e`$ conversion in vacuum, assuming three family mixing and CP-violation, is given by: $`P_{CPodd}(\mu ,e)`$ $`=`$ $`2\mathrm{cos}(\theta _{13}^m)^2\mathrm{sin}\delta ^m\mathrm{sin}(2\theta _{12}^m)\mathrm{sin}(\theta _{13}^m)\mathrm{sin}(2\theta _{23}^m)\times `$ $`\mathrm{sin}({\displaystyle \frac{\mathrm{\Delta }M_{12}^2L}{4E}})\mathrm{sin}({\displaystyle \frac{\mathrm{\Delta }M_{13}^2L}{4E}})\mathrm{sin}({\displaystyle \frac{\mathrm{\Delta }M_{23}^2L}{4E}})`$ where the $`CP`$-phase in matter is: $$e^{i\delta _m}=\frac{(G^2e^{i\delta }F^2e^{i\delta })s_{23}c_{23}+GF(c_{23}^2s_{23}^2)}{\sqrt{(G^2s_{23}^2+F^2c_{23}^2+2GFc_{23}s_{23}c_\delta )(G^2c_{23}^2+F^2s_{23}^22GFc_{23}s_{23}c_\delta )}}$$ (31) Here, $`G`$ and $`F`$ are given in the Appendix. In case of degenerate mass differences, $`\|\mathrm{\Delta }m_{21}^2\|\|\mathrm{\Delta }m_{32}^2\|`$, the $`\delta `$ phase is significantly modified from its original value in vacuum for $`\delta >\pi /4`$ radians, while in the case $`\|\mathrm{\Delta }m_{32}^2\|>\|\mathrm{\Delta }m_{21}^2\|`$ the CP phase is almost unaffected by the presence of matter. In addition, neutrino propagation in a dense medium makes more difficult to experimentally extract a genuine CP violation signal, since the asymmetrical behavior of matter, with respect to neutrinos and anti-neutrinos, induces fake CP violation effects. ### 10.1 $`CP`$ statistical significance To evaluate the sensitivity to this kind of measurement as a function of the oscillation parameters and of the selected baseline, we compare the case where CP violation is maximal ($`\delta =\pi /2`$) and the case of no CP violation ($`\delta =0`$). Figure 13 shows the sensitivity $$S_{CP}(E_i)\frac{N(\delta =\pi /2,E_i)N(\delta =0,E_i)}{\sqrt{N(\delta =0,E_i)}}$$ (32) i.e. the difference between the two extreme cases, divided by the statistical error (being N the number of events in each energy bin, for a given value of $`\delta `$). We can see that for the baseline of 732 km, this quantity is positive for almost the full energy range, so there is no real shape variation in the spectrum, and the CP effect is similar to what would be obtained with a change in the angle $`\theta _{13}`$. On the other hand, for larger baselines this curve crosses the zero, and the CP violation produces a visible deformation of the energy spectrum. ### 10.2 Fitting of the $`\delta `$ parameter To evaluate the precision reachable on the measurement of these three oscillations parameters, we perform fits assuming the following reference values: $`\mathrm{\Delta }m_{32}^2`$ $`=`$ $`3.5,5,7\times 10^3\mathrm{eV}^2`$ $`\mathrm{\Delta }m_{12}^2`$ $`=`$ $`1\times 10^4\mathrm{eV}^2`$ $`\mathrm{sin}^2\theta _{23}`$ $`=`$ $`0.5\theta _{23}=45^o`$ $`\mathrm{sin}^22\theta _{13}`$ $`=`$ $`0.05\theta _{13}=6.5^o`$ $`\mathrm{sin}^2\theta _{12}`$ $`=`$ $`0.5`$ $`\delta 0`$ (33) Figure 14 shows for the baseline L=2900 km, the expected visible energy spectra for electron and wrong sing muon class for the two cases $`\delta =0`$ and $`\delta =\pi /2`$, and their difference in terms of number of events. As expected, in absolute values, the CP violation affects the two classes in a similar way, but the effect is much more visible and significant for the wrong-sign muons, since the total number of events is smaller. Figure 14 is normalized to $`10^{21}\mu ^+`$ decays, since this is the minimum amount of decays required to produce a statistically significant observation of CP violation. For the largest baseline $`L=7400`$ km, CP violation effects are almost undetectable even in case $`\delta =\pi /2`$ since, wrong sing muon appearance is largely dominated by matter effects. At $`L=732`$ km, the effect of a non-vanishing $`\delta `$ is more striking thanks to the higher event rates and the smaller neutrino path in matter. However, as we formerly pointed out, this effect is similar to the one produced by a smaller value of $`\theta _{13}`$ and therefore, both effects cannot be disentangled with a single measurement at this baseline as shown in figure 15, where the correlations between $`\delta `$ and $`\theta _{13}`$ prevent a precise determination of any of them. Nonetheless, a measurement at $`L=7400`$ km where CP violation effects are negligible and the most accurate determination of $`\theta _{13}`$ exits, can be combined with data collected at $`L=732`$ km to produce a precise determination of $`\delta `$. Another possibility to unveil the existence of CP violation is to perform a single measurement at $`L=2900`$ km. As shown in figure 14, at this distance the effect of $`\delta 0`$ is twofold: not only the event rate is modified but also the spectral shape. This last effect cannot be produced by a change on $`\theta _{13}`$. Figure 15 shows that at this distance the correlation between $`\delta `$ and $`\theta _{13}`$ has diminished and therefore a better determination of the parameters can be achieved with a single measurement. For a baseline of $`L=2900`$ km, matter effects are not so strong as for the longest baseline. The possibility of measuring the CP-violating phase $`\delta `$ is not spoiled by the fake asymmetries due to the matter interactions, but both effects, even if correlated, can be measured at the same time. This is shown in figure 16, where the result of a simultaneous fit to the average matter density $`\rho `$ and the phase $`\delta `$ is presented for the two cases where $`\rho `$ is left as a free parameter or it is known beforehand with a $`3\%`$ accuracy. In case no CP violation is observed, the allowed $`\delta `$ values at 90$`\%`$ C.L. as a function of $`\mathrm{\Delta }m_{12}^2`$ are shown in figure 17 for two different muon normalizations. On the other hand, if a significant effect is detected, the precision achievable on the measurement of the CP phase $`\delta `$ is shown in figure 18. We fit $`\delta `$ using all event classes, leaving the five parameters governing the oscillation free. Assuming a reference value of $`90^o`$ for $`\delta `$ and $`10^{21}`$ muon decays of each polarity, we get: $`90\pm 15^o`$. Therefore a precision around 20$`\%`$ is expected. Finally, we note that the change on the expected precision for $`\theta _{13}`$ is negligible when $`\mathrm{\Delta }m_{12}^2`$ and $`\theta _{12}`$ are included in the fit. ### 10.3 Use of quasi-elastic events Quasi-elastic events can also be useful to spot the presence of CP violation. Table 8 shows the expected number of QE electron events after kinematics cuts (a back to back electron-proton final state with proton kinetic energy in excess of 100 MeV). Three different assumptions for the total number of muon decays and the baseline of 2900 km have been assumed. In case CP is conserved, we expect 35 (96) quasi-elastic electron events for $`\mathrm{\Delta }m_{32}^2=3.5(7)\times 10^3`$ eV<sup>2</sup> and $`10^{21}`$ muon decays of each polarity. In case $`\delta =\pi /2`$, we expect 26 and 85 events for the two mass differences considered. The effect is at the one sigma level. To obtain a statistically conclusive signal for CP violation would require more than $`10^{21}`$ muon decays. ## 11 Conclusions In this paper, we tried to understand in deeper detail the capabilities an experiment at the Neutrino Factory, with the aim of exploiting as much as possible the possibility of studying several neutrino transitions at the same time. For this reason, we have considered as a baseline detector a large Liquid Argon TPC with external muon identifier, the most versatile design proposed so far for large neutrino experiments. Assuming an oscillation scenario favoured by present experimental results, the leading oscillation would be between the second and third neutrino family, which are maximally mixed. Therefore, a very precise determination of the parameters governing this transition, $`\theta _{23}`$ and $`\mathrm{\Delta }m_{23}^2`$ is essential also for the understanding of all other processes. This is mainly achieved using the information coming from the $`\nu _\mu `$ disappearance (right-sign muon class), provided that the baseline and beam energy are chosen in such a way that the first oscillation maximum is visible as a dip in the oscillated spectrum. The maximal sensitivity to $`\theta _{13}`$ is achieved for very small background levels, since we are looking in this case for small signals; most of the information is coming from the clean wrong-sign muon class, and from quasi-elastic events. On the other hand, if its value is not too small, for a measurement of $`\theta _{13}`$, the signal/background ratio could be not so crucial, and also the other event classes can contribute to this measurement. Like for a B-Factory, a $`\nu `$-Factory should have among its aims the overconstraining of the oscillation pattern, in order to look for unexpected new physics effects. This can be achieved in global fits of the parameters, where the unitarity of the mixing matrix is not strictly assumed. Using a detector able to identify the $`\tau `$ lepton production via kinematic means, it is possible to verify the unitarity in $`\nu _\mu \nu _\tau `$ and $`\nu _e\nu _\tau `$ transitions. For this latter, the possibility of a kinematical $`\tau `$ identification for wrong-sign muon events could allow for the first time a clear identification of this type of oscillations. The study of CP violation in the lepton system is a very fascinating subject, and probably the most ambitious goal of this kind of machines. It will only be possible for high beam intensities, and if the parameters governing the solar neutrino deficit are in the region usually indicated as large mixing angle MSW solution. Matter effect can mimic CP violation; however, a multiparameter fit at the right baseline can allow a simultaneous determination of matter and CP-violating parameters. Also for CP-violation measurements most of the information would come from the wrong-sign muon class, but since in this case the electron class would also be affected, the study of these events (and of the very clean quasi-elastic interactions) can provide essential cross-checks for these delicate measurements. ## Appendix A Oscillations in matter In reference , the authors compute analytic expressions for the mass eigenvalues, mixing angles and CP-violation phase in matter, assuming the mixing matrix $`U`$ is parametrized “à la CKM”. Since several missprints were observed, we reproduce here the corrected expressions: For neutrinos, the mass eigenvalues in matter $`M_1,M_2`$ and $`M_3`$ are: $`M_1^2`$ $`=`$ $`m_1^2+{\displaystyle \frac{A}{3}}{\displaystyle \frac{1}{3}}\sqrt{A^23B}S{\displaystyle \frac{\sqrt{3}}{3}}\sqrt{A^23B}\sqrt{1S^2}`$ (34) $`M_2^2`$ $`=`$ $`m_1^2+{\displaystyle \frac{A}{3}}{\displaystyle \frac{1}{3}}\sqrt{A^23B}S+{\displaystyle \frac{\sqrt{3}}{3}}\sqrt{A^23B}\sqrt{1S^2}`$ (35) $`M_3^2`$ $`=`$ $`m_1^2+{\displaystyle \frac{A}{3}}+{\displaystyle \frac{2}{3}}\sqrt{A^23B}S`$ (36) where $`A`$ $`=`$ $`\mathrm{\Delta }m_{21}^2+\mathrm{\Delta }m_{31}^2+D`$ (37) $`B`$ $`=`$ $`\mathrm{\Delta }m_{21}^2\mathrm{\Delta }m_{31}^2+D[\mathrm{\Delta }m_{31}^2c_{13}^2+\mathrm{\Delta }m_{21}^2(c_{13}^2c_{12}^2+s_{13}^2)]`$ (38) $`C`$ $`=`$ $`D\mathrm{\Delta }m_{21}^2\mathrm{\Delta }m_{31}^2c_{13}^2c_{12}^2`$ (39) $`D`$ $`=`$ $`2\sqrt{2}G_FN_eE\mathrm{for}\mathrm{neutrinos}`$ (40) $`S`$ $`=`$ $`\mathrm{cos}\left[{\displaystyle \frac{1}{3}}\mathrm{arccos}\left({\displaystyle \frac{2A^39AB+27C}{2\sqrt{(A^23B)^3}}}\right)\right]`$ (41) and $`c_{ij}=\mathrm{cos}\theta _{ij}`$ and $`s_{ij}=\mathrm{sin}\theta _{ij}`$. The mixing angles and CP phase in matter are: $`\mathrm{sin}^2\theta _{12}^m`$ $`=`$ $`{\displaystyle \frac{(M_2^4\alpha M_2^2+\beta )\mathrm{\Delta }M_{31}^2}{\mathrm{\Delta }M_{32}^2(M_1^4\alpha M_1^2+\beta )\mathrm{\Delta }M_{31}^2(M_2^4\alpha M_2^2+\beta )}}`$ (42) $`\mathrm{sin}^2\theta _{13}^m`$ $`=`$ $`{\displaystyle \frac{M_3^4\alpha M_3^2+\beta }{\mathrm{\Delta }M_{31}^2\mathrm{\Delta }M_{32}^2}}`$ (43) $`\mathrm{sin}^2\theta _{23}^m`$ $`=`$ $`{\displaystyle \frac{G^2s_{23}^2+F^2c_{23}^2+2GFc_{23}s_{23}c_\delta }{G^2+F^2}}`$ (44) $`e^{i\delta _m}`$ $`=`$ $`{\displaystyle \frac{(G^2e^{i\delta }F^2e^{i\delta })s_{23}c_{23}+GF(c_{23}^2s_{23}^2)}{\sqrt{(G^2s_{23}^2+F^2c_{23}^2+2GFc_{23}s_{23}c_\delta )(G^2c_{23}^2+F^2s_{23}^22GFc_{23}s_{23}c_\delta )}}}`$ (45) where $`\alpha `$ $`=`$ $`m_3^2c_{13}^2+m_2^2(c_{13}^2c_{12}^2+s_{13}^2)+m_1^2(c_{13}^2s_{12}^2+s_{13}^2)`$ (46) $`\beta `$ $`=`$ $`m_3^2c_{13}^2(m_2^2c_{12}^2+m_1^2s_{12}^2)+m_2^2m_1^2s_{13}^2`$ (47) $`G`$ $`=`$ $`[\mathrm{\Delta }m_{31}^2(M_3^2m_1^2\mathrm{\Delta }m_{21}^2)\mathrm{\Delta }m_{21}^2(M_3^2m_1^2\mathrm{\Delta }m_{31}^2)s_{12}^2]c_{13}s_{13}`$ (48) $`F`$ $`=`$ $`(M_3^2m_1^2\mathrm{\Delta }m_{31}^2)\mathrm{\Delta }m_{21}^2c_{12}s_{12}c_{13}`$ (49) For anti-neutrinos, we must replace $`D`$ by $`D`$. In the case of one-mass scale approximation, the oscillation probabilities are simply The oscillation probabilities are: $`P(\nu _e\nu _e,E,L)`$ $`=`$ $`1\mathrm{sin}^2(2\theta _{13}^m)\mathrm{\Delta }_{32}^2`$ (50) $`P(\nu _e\nu _\mu ,E,L)`$ $`=`$ $`\mathrm{sin}^2(2\theta _{13}^m)\mathrm{sin}^2(\theta _{23}^m)\mathrm{\Delta }_{32}^2`$ $`P(\nu _e\nu _\tau ,E,L)`$ $`=`$ $`\mathrm{sin}^2(2\theta _{13}^m)\mathrm{cos}^2(\theta _{23}^m)\mathrm{\Delta }_{32}^2`$ $`P(\nu _\mu \nu _\mu ,E,L)`$ $`=`$ $`14\mathrm{cos}^2(\theta _{13}^m)\mathrm{sin}^2(\theta _{23}^m)\left[1\mathrm{cos}^2(\theta _{13}^m)\mathrm{sin}^2(\theta _{23}^m)\right]\mathrm{\Delta }_{32}^2`$ $`P(\nu _\mu \nu _\tau ,E,L)`$ $`=`$ $`\mathrm{cos}^4(\theta _{13}^m)\mathrm{sin}^2(2\theta _{23}^m)\mathrm{\Delta }_{32}^2`$ where $`\mathrm{\Delta }_{32}^2=\mathrm{sin}^2\left((M_3^2M_2^2)L/4E\right)`$.
warning/0005/cond-mat0005362.html	ar5iv	text	# Shift in the velocity of a front due to a cut-off ## 1 Introduction Equations describing the propagation of a front between a stable and an unstable state appear in a large variety of situations in physics, chemistry and biology. One of the simplest equations of this kind is the Fisher-Kolmogorov equation $$\frac{h}{t}=\frac{^2h}{x^2}+hh^3,$$ (1) which describes the evolution of a space and time dependent concentration $`h(x,t)`$ in a reaction-diffusion system. This equation, originally introduced to study the spread of advantageous genes in a population , has been widely used in other contexts, in particular to describe the time dependence of the concentration of some species in a chemical reaction . For such an equation, the uniform solutions $`h=1`$ and $`h=0`$ are respectively stable and unstable and it is known that for initial conditions such that $`h(x,0)1`$ as $`x\mathrm{}`$ and $`h(x,0)0`$ as $`x+\mathrm{}`$ there exists a one parameter family $`F_v`$ of traveling wave solutions (indexed by their velocity $`v`$) of the form $$h(x,t)=F_v\left(xvt\right),$$ (2) with $`F_v`$ decreasing, $`F_v\left(z\right)1`$ as $`z\mathrm{}`$ and $`F_v\left(z\right)0`$ as $`z\mathrm{}`$. The analytic expression of the shape $`F_v`$ is in general not known but one can determine the range of velocities $`v`$ for which solutions of type (2) exist. If one assumes an exponential decay $$F_v\left(z\right)e^{\gamma z}\text{for large }z,$$ (3) it is easy to see by replacing (2) and (3) into (1) that the velocity $`v`$ is given by $$v\left(\gamma \right)=\gamma +\frac{1}{\gamma }.$$ (4) As $`\gamma `$ is arbitrary, this shows the well known fact that the range of possible velocities is $`v2`$. The minimal velocity $`v_0=2`$ is reached for $`\gamma _0=1`$ and for steep enough initial conditions $`h(x,0)`$ (which decay faster than $`e^{\gamma _0x}`$), the solution selected for large $`t`$ is the one corresponding to this minimal velocity $`v_0`$. Equations of type (1) are obtained either as the large scale limit or as the mean field limit of physical situations which are discrete at the microscopic level (particles, lattice models, etc.) As the number of particles is an integer, the concentration $`h(x,t)`$ could be thought as being larger than some $`\epsilon `$, which would correspond to the value of $`h(x,t)`$ when a single particle is present. Equations of type (1) appear then as the limit of the discrete model when $`\epsilon 0`$. Several authors have already noticed in their numerical works that the speed $`v_\epsilon `$ of the discrete model converges slowly, as $`\epsilon `$ tends to 0, towards the minimal velocity $`v_0`$. We believe that the main effect of having $`\epsilon 0`$ is to introduce a cut-off in the tail of the front, and that this changes noticeably the speed. The speed of the front is in general governed by its tail. In the present work, we consider equations similar to (1), which we modify in such a way that whenever $`h(x,t)`$ is much smaller than a cut-off $`\epsilon `$, it is replaced by 0. The cut-off $`\epsilon `$ can be introduced by replacing (1) by $$\frac{h}{t}=\frac{^2h}{x^2}+\left(hh^3\right)a\left(h\right),$$ (5) with $`a\left(h\right)`$ $`=1`$ $`\text{if }h>\epsilon ,`$ (6) $`a\left(h\right)`$ $`1`$ $`\text{if }h\epsilon .`$ For example, one could choose $`a\left(h\right)=1`$ for $`h\epsilon `$ and $`a\left(h\right)=h/\epsilon `$ for $`h\epsilon `$. Another choice that we will use in section 4 is simply $`a\left(h\right)=1`$ if $`h>\epsilon `$ and $`a\left(h\right)=0`$ if $`h\epsilon `$. The question we address here is the effect of the cut-off $`\epsilon `$ on the velocity $`v_\epsilon `$ of the front. We will show that the velocity $`v_\epsilon `$ converges, as $`\epsilon 0`$, to the minimal velocity $`v_0`$ of the original problem (without cut-off) and that the main correction to the velocity of the front is $$v_\epsilon v_0\frac{\pi ^2\gamma _0^2}{2}v^{\prime \prime }\left(\gamma _0\right)\frac{1}{\left(\mathrm{log}\epsilon \right)^2}$$ (7) for an equation of type (1) for which the velocity is related to the exponential decay $`\gamma `$ of the shape (2) by some relation $`v\left(\gamma \right)`$. (Everywhere we note by $`v_0`$ the minimal velocity and $`\gamma _0`$ the corresponding value of the decay $`\gamma `$.) In the particular case of equation (1), where $`v\left(\gamma \right)`$ is given by (4), this becomes $$v_\epsilon 2\frac{\pi ^2}{\left(\mathrm{log}\epsilon \right)^2}.$$ (8) In section 2 we describe an equation of type (1) where both space and time are discrete, so that simulations are much easier to perform. The results of the numerical simulations of this equation are described in section 3: as $`\epsilon 0`$, the velocity is seen to converge like $`\left(\mathrm{log}\epsilon \right)^2`$ to the minimal velocity $`v_0`$, and the shape of the front appears to take a scaling form. In section 4 we show that for equations of type (1) in presence of a small cut-off $`\epsilon `$ as in (5), one can calculate both the shape of the front and the shift in velocity. The results are in excellent agreement with the numerical data of section 3. In section 5 we consider a model defined, for a finite number $`N`$ of particles, by some microscopic stochastic dynamics which reduces to the front equation of sections 3 and 4 in the limit $`N\mathrm{}`$. Despite the presence of noise, our simulations indicate that in this case too, the velocity dependence of the front decays slowly (as $`\left(\mathrm{log}N\right)^2`$) to the minimal velocity $`v_0`$ of the front. ## 2 A discrete front equation To perform numerical simulations, it is much easier to study a case where both time and space are discrete variables. We consider here the equation $$h(x,t+1)=g(x,t)\mathrm{\Theta }\left[g(x,t)\epsilon \right],$$ (9a) where $$g(x,t)=1\left[1ph(x1,t)\left(1p\right)h(x,t)\right]^2.$$ (9b) Time is a discrete variable and if initially the concentration $`h(x,0)`$ is only defined when $`x`$ is an integer, $`h(x,t)`$ remains so at any later time. Because $`t`$ and $`x`$ are both integers, the cut-off $`\epsilon `$ can be introduced as in (9) in the crudest way using a Heaviside $`\mathrm{\Theta }`$ function. (We have checked however that other ways of introducing the cut-off $`\epsilon `$ as in (5, 6) do not change the results.) Equation (9) appears naturally (in the limit $`\epsilon =0`$) in the problem of directed polymers on disordered trees (where the energy of the bonds is either $`1`$ with probability $`p`$ or $`0`$ with probability $`1p`$). At this stage we will not give a justification for introducing the cut-off $`\epsilon `$. This will be discussed in section 5. We consider for the initial condition a step function $`h(x,0)`$ $`=0`$ if $`x0`$, (10) $`h(x,0)`$ $`=1`$ if $`x<0`$. Clearly for such an initial condition, $`h(x,t)=1`$ for $`x<0`$ at all times. As $`h(x,t)1`$ behind the front and $`h(x,t)0`$ ahead of the front, we define the position $`X_t`$ of the front at time $`t`$ by $$X_t=\underset{x=0}{\overset{+\mathrm{}}{}}h(x,t).$$ (11) The velocity of the front $`v_\epsilon `$ can then be calculated by $$v_\epsilon =\underset{t\mathrm{}}{lim}\frac{X_t}{t}=X_{t+1}X_t,$$ (12) where the average is taken over time. (Note that as $`h(x,t)`$ is only defined on integers, the difference $`X_{t+1}X_t`$ is time dependent and has to be averaged as in (12).) When $`\epsilon =0`$, the evolution equation (9) becomes $$h(x,t+1)=1\left[1ph(x1,t)\left(1p\right)h(x,t)\right]^2.$$ (13) As for (1), there is a one parameter family of solutions $`F_v`$ of the form (2) indexed by the velocity $`v`$ which is related (3) to the exponential decay $`\gamma `$ of the shape by $$v\left(\gamma \right)=\frac{1}{\gamma }\mathrm{log}\left(2pe^\gamma +2\left(1p\right)\right).$$ (14) (This relation is obtained as (4) by considering the tail of the front where $`h(x,t)`$ is small and where therefore (13) can be linearized.) One can show that for $`p<\frac{1}{2}`$, $`v\left(\gamma \right)`$ reaches a minimal value $`v_0`$ smaller than 1 for some $`\gamma _0`$, whereas for $`p\frac{1}{2}`$, $`v\left(\gamma \right)`$ is a strictly decreasing function of $`\gamma `$, implying that the minimal velocity is $`v_0=\underset{\gamma \mathrm{}}{lim}v\left(\gamma \right)=1`$. We will not discuss here this phase transition and we assume from now on that $`p<\frac{1}{2}`$. Table 1 gives some values of $`v_0`$ and $`\gamma _0`$ obtained from (14). It is important to notice that for $`p<\frac{1}{2}`$, the function $`v\left(\gamma \right)`$ has a single minimum at $`\gamma _0`$. Therefore, there are in general two choices $`\gamma _1`$ and $`\gamma _2`$ of $`\gamma `$ for each velocity $`v`$. For $`vv_0`$, the exponential decay of $`F_v\left(z\right)`$ is dominated by min$`(\gamma _1,\gamma _2)`$. As $`vv_0`$, the two roots $`\gamma _1`$ and $`\gamma _2`$ become equal, and the effect of this degeneracy gives (in a well chosen frame) $$F_{v_0}\left(z\right)Aze^{\gamma _0z}\text{for large }z\text{,}$$ (15) where $`A`$ is a constant. This large $`z`$ behavior can be recovered by looking at the general solution of the linearized form of equation (13) $$h(x,t+1)=2ph(x1,t)+2\left(1p\right)h(x,t).$$ (16) ## 3 Numerical determination of the velocity We iterated numerically (9) with the initial condition (10) for several choices of $`p<\frac{1}{2}`$ and for $`\epsilon `$ varying between $`0.03`$ and $`10^{17}`$. We observed that the speed is usually very easy to measure because, after a short transient time, the system reaches a periodic regime for which $$h(x,t+T)=h(xY,t)$$ (17) for some constants $`T`$ and $`Y`$. The speed $`v_\epsilon `$ of the front is then simply given by $$v_\epsilon =\frac{Y}{T}.$$ (18) For example, for $`p=0.25`$ and $`\epsilon =10^5`$, we find $`T=431`$ and $`Y=343`$ so that $`v_\epsilon =343/431`$. The emergence of this periodic behavior is due to the locking of the dynamical system of the $`h(x,t)`$ on a limit cycle. Because $`Y`$ and $`T`$ are integers, our numerical simulations give the speed with an infinite accuracy. For each choice of $`p`$ and $`\epsilon `$, we measured the speed of the front, as defined by (12) and its shape. Figure 1 is a log-log plot of the difference $`v_0v_\epsilon `$ versus $`\epsilon `$ (varying between $`0.03`$ to $`10^{17}`$) for three choices of the parameter $`p`$. The solid lines on the plot indicate the value predicted by the calculations of section 4. We see on this figure that the velocity $`v_\epsilon `$ converges slowly towards the minimal velocity $`v_0`$ as $`\epsilon 0`$. Our simulations, done over several orders of magnitude (here, fifteen), reveal that the convergence is logarithmic: $`v_0v_\epsilon \left(\mathrm{log}\epsilon \right)^2`$. As the front is moving, to measure its shape, we need to locate its position. Here we use expression (11) and we measure the shape $`s_\epsilon \left(z\right)`$ of the front at a given time $`t`$ relative to its position $`X_t`$ by $$s_\epsilon \left(z\right)=h(z+X_t,t)$$ (19) When the system reaches the limit cycle (17), the shape $`s_\epsilon \left(z\right)`$ becomes roughly independent of the time chosen. (In fact it becomes periodic of period $`T`$, but the shape $`s_\epsilon `$ has a smooth envelope.) We have measured this shape at some arbitrary large enough time to avoid transient effects. As we expect $`s_\epsilon \left(z\right)`$ to look more and more like $`F_{v_0}\left(z\right)`$ as $`\epsilon `$ tends to 0, we normalize this shape by dividing it by $`e^{\gamma _0z}`$. The result $`s_\epsilon \left(z\right)e^{\gamma _0z}`$ is plotted versus $`z`$ for $`p=0.25`$ and $`\epsilon =10^9`$, $`10^{11}`$, $`10^{13}`$, $`10^{15}`$ and $`10^{17}`$ in figure 2. On the left part of the graph, our data coincide over an increasing range as $`\epsilon `$ decreases, indicating that far from the cut-off, the shape converges to expression (15) of $`F_{v_0}\left(z\right)`$. On the right part, the curves increase up to a maximum before falling down to some small value which seems to be independent of $`\epsilon `$. When $`\epsilon `$ is multiplied by a constant factor (here $`10^2`$), the maximum as well as the right part of the curves are translated by a constant amount. This indicates that for $`\epsilon `$ small enough, the shape $`s_\epsilon \left(z\right)`$ in the tail (that is for $`z`$ large) takes the scaling form $$s_\epsilon \left(z\right)\left\|\mathrm{log}\epsilon \right\|G\left(\frac{z}{\left\|\mathrm{log}\epsilon \right\|}\right)e^{\gamma _0z}.$$ (20) We will see that our analysis of section 4 does predict this scaling form. As one expects this shape to coincide with the asymptotic form (15) of $`F_{v_0}\left(z\right)`$ for $`1z\left\|\mathrm{log}\epsilon \right\|`$, the scaling function $`G\left(y\right)`$ should be linear for small $`y`$. ## 4 Calculation of the velocity for a small cut-off The first remark we make is that as soon as we introduce a cut-off through a function $`a\left(h\right)`$ which is everywhere smaller than 1, the velocity $`v_\epsilon `$ of the front is lowered compared to the velocity obtained in the absence of a cut-off. This is easy to check by comparing a solution $`h_\epsilon (x,t)`$ of (5) where $`a\left(h\right)`$ is present and a solution $`h_0(x,t)`$ of (1). If initially $`h_\epsilon (x,0)<h_0(x,0)`$, the solution $`h_\epsilon `$ will never be able to take over the solution $`h_0`$. Indeed, would the two functions $`h_\epsilon (x,0)`$ and $`h_0(x,0)`$ coincide for the first time at some point $`x`$, we would have at that point $`^2h_\epsilon /x^2^2h_0/x^2`$ and together with the effect of $`a\left(h\right)`$ this would bring back the system in the situation where $`h_\epsilon (x,t)<h_0(x,t)`$ . This shows that $`v_\epsilon v_0`$. For the calculation of the velocity $`v_\epsilon `$, we will consider first the modified Fisher-Kolmogorov equation (5) when the cut-off function $`a\left(h\right)`$ is simply given by $$a\left(h\right)=\mathrm{\Theta }\left(h\epsilon \right).$$ (21) In this section we will calculate the leading correction to the velocity when $`\epsilon `$ is small and we will obtain the scaling function $`G`$ which appears in (20). Then we will discuss briefly how our analysis could be extended to more general forms of the cut-off function $`a\left(h\right)`$ or to other traveling wave equations such as (9). As $`v_\epsilon `$ is the velocity of the front, its shape $`s_\epsilon \left(z\right)=h(z+v_\epsilon t,t)`$ in the asymptotic regime satisfies $$v_\epsilon s_\epsilon ^{}+s_\epsilon ^{\prime \prime }+\left(s_\epsilon s_\epsilon ^2\right)a\left(s_\epsilon \right)=0.$$ When $`\epsilon `$ is small, with the choice (21) for $`a\left(h\right)`$, we can decompose the range of values of $`z`$ into three regions: > where $`s_\epsilon \left(z\right)`$ is not small compared to 1. > > where $`\epsilon <s_\epsilon \left(z\right)1`$. > > where $`s_\epsilon \left(z\right)<\epsilon `$. In region I, the shape of the front $`s_\epsilon `$ looks like $`F_{v_0}`$ whereas in regions II and III, as $`s_\epsilon `$ is small, it satisfies the linear equations $`v_\epsilon s_\epsilon ^{}+s_\epsilon ^{\prime \prime }+s_\epsilon =0`$ in region II, (22) $`v_\epsilon s_\epsilon ^{}+s_\epsilon ^{\prime \prime }=0`$ in region III. (23) These linear equations (22,23) can be solved easily. The only problem is to make sure that the solution in region II and its derivative coincides with $`F_{v_0}`$ at the boundary between I and II and with the solution valid in region III at the boundary between II and III. If we call $`\mathrm{\Delta }`$ the shift in the velocity $$\mathrm{\Delta }=v_0v_\epsilon ,$$ (24) and if we note $`\gamma _r\pm i\gamma _i`$ the two roots of the equation $`v\left(\gamma \right)=v_\epsilon `$, the shape $`s_\epsilon `$ is given in the three regions by $`s_\epsilon \left(z\right)`$ $`F_{v_0}\left(z\right)`$ in region I, $`s_\epsilon \left(z\right)`$ $`Ce^{\gamma _rz}\mathrm{sin}\left(\gamma _iz+D\right)`$ in region II, (25) $`s_\epsilon \left(z\right)`$ $`\epsilon e^{v_\epsilon \left(zz_0\right)}`$ in region III, and we can determine the unknown quantities $`C`$, $`D`$, $`z_0`$ and $`v_\epsilon `$ by using the boundary conditions. For large $`z`$ we know from (15) that $`F_{v_0}\left(z\right)Aze^{\gamma _0z}`$ for some $`A`$. Therefore, as $`\gamma _0\gamma _r\mathrm{\Delta }`$ and $`\gamma _i\mathrm{\Delta }^{1/2}`$, the boundary conditions between regions I and II impose, to leading order in $`\mathrm{\Delta }^{1/2}`$, that $`C=A/\gamma _i`$ and $`D=0`$. At the boundary between regions II and III, we have $`s_\epsilon \left(z\right)=\epsilon `$ and $`z=z_0`$. If we impose the continuity of $`s_\epsilon `$ and of its first derivative at this point, we get $$Ae^{\gamma _rz_0}\mathrm{sin}\left(\gamma _iz_0\right)=\epsilon \gamma _i,$$ (26a) and $$Ae^{\gamma _rz_0}\left[\gamma _r\mathrm{sin}\left(\gamma _iz_0\right)+\gamma _i\mathrm{cos}\left(\gamma _iz_0\right)\right]=v_\epsilon \epsilon \gamma _i.$$ (26b) Taking the ratio between these two relations leads to $$\gamma _r\frac{\gamma _i}{\mathrm{tan}\left(\gamma _iz_0\right)}=v_\epsilon .$$ (27) When $`\mathrm{\Delta }`$ is small, $`\gamma _r\gamma _0=1`$, $`v_\epsilon v_0=2`$ and $`\gamma _i\mathrm{\Delta }^{1/2}`$. Thus the only way to satisfy (27) is to set $`\gamma _iz_0\pi `$ and $`\pi \gamma _iz_0\gamma _i\mathrm{\Delta }^{1/2}`$. Therefore, (26) implies to leading order that $`z_0\left(\mathrm{log}\epsilon \right)/\gamma _0`$ and the condition $`\gamma _iz_0\pi `$ gives $$\gamma _i\frac{\pi }{z_0}\frac{\pi \gamma _0}{\left\|\mathrm{log}\epsilon \right\|}$$ (28) Then, as $`\gamma _i`$ is small, the difference $`\mathrm{\Delta }=v_0v_\epsilon `$ is given by $$v_0v_\epsilon \frac{1}{2}v^{\prime \prime }\left(\gamma _0\right)\gamma _i^2\frac{v^{\prime \prime }\left(\gamma _0\right)\pi ^2\gamma _0^2}{2\left(\mathrm{log}\epsilon \right)^2}$$ (29) which is the result announced in (7) and (8). A different cut-off function $`a\left(h\right)`$ should not affect the shape of $`s_\epsilon `$ in the region II or the size $`z_0`$ of region II. Only the precise matching between regions II and III might be modified and we do not think that this would change the leading dependency of $`z_0`$ in $`\epsilon `$ which controls everything. In fact there are other choices of the cut-off function $`a\left(h\right)`$ (piecewise constant) for which we could find the explicit solution in region III, confirming that the precise form of $`a\left(h\right)`$ does not change (28). The generalization of the above argument to equations other than (1) (and in particular to the case studied in sections 2 and 3) is straightforward. Only the form of the linear equation is changed and the only effect on the final result (7) is that one has to use a different function $`v\left(\gamma \right)`$. When expression (7) is compared in figure 1 to the results of the simulations, the agreement is excellent. Moreover, in region II, one sees from (4) and (28) that $$s_\epsilon \left(z\right)\frac{A}{\pi \gamma _0}\left\|\mathrm{log}\epsilon \right\|\mathrm{sin}\left(\frac{\pi \gamma _0z}{\left\|\mathrm{log}\epsilon \right\|}\right)e^{\gamma _0z},$$ (30) which also agrees with the scaling form (20). Recently, for a simple model of evolution governed by a linear equation, the velocity was found to be the logarithm of the cut-off to the power $`\frac{1}{3}`$. This result was obtained by an analysis which has some similarities to the one presented in this section. ## 5 A stochastic model Many models described by traveling wave equations originate from a large scale limit of microscopic stochastic models involving a finite number $`N`$ of particles . Here we study such a microscopic model, the limit of which reduces to (13) when $`N\mathrm{}`$. Our numerical results, presented below, indicate a large $`N`$ correction to the velocity of the form $`v_Nv_0a\left(\mathrm{log}N\right)^2`$ with a coefficient $`a`$ consistent with the one calculated in section 4 for $`\epsilon =\frac{1}{N}`$. The model we consider in this section appears in the study of directed polymers and is, up to minor changes, equivalent to a model describing the dynamics of hard spheres . It is a stochastic process discrete both in time and space with two parameters: $`N`$, the number of particles, and $`p`$, a real number between 0 and 1. At time $`t`$ ($`t`$ is an integer), we have $`N`$ particles on a line at integer positions $`x_1\left(t\right)`$, $`x_2\left(t\right)`$, …, $`x_N\left(t\right)`$. Several particles may occupy the same site. At each time-step, the $`N`$ positions evolve in the following way: for each $`i`$, we choose two particles $`j_i`$ and $`j_i^{}`$ at random among the $`N`$ particles. (These two particles do not need to be different.) Then we update $`x_i\left(t\right)`$ by $$x_i\left(t+1\right)=\mathrm{max}(x_{j_i}\left(t\right)+\alpha _i,x_{j_i^{}}\left(t\right)+\alpha _i^{}),$$ (31) where $`\alpha _i`$ and $`\alpha _i^{}`$ are two independent random numbers taking the value 1 with probability $`p`$ or 0 with probability $`1p`$. The numbers $`\alpha _i`$, $`\alpha _i^{}`$, $`j_i`$ and $`j_i^{}`$ change at each time-step. Initially ($`t=0`$), all particles are at the origin so that we have $`x_i\left(0\right)=0`$ for all $`i`$. At time $`t`$, the distribution of the $`x_i\left(t\right)`$ on the line can be represented by a function $`h(x,t)`$ which counts the fraction of particles strictly at the right of $`x`$. $$h(x,t)=\frac{1}{N}\underset{x_i\left(t\right)>x}{}1.$$ (32) Obviously $`h(x,t)`$ is always an integral multiple of $`\frac{1}{N}`$. At $`t=0`$, we have $`h(x,0)=1`$ if $`x<0`$ and $`h(x,0)=0`$ if $`x0`$. One can notice that the definition of the position $`X_t`$ of the front used in (11) coincides with the average position of the $`N`$ particles $$X_t=\underset{x=0}{\overset{+\mathrm{}}{}}h(x,t)=\frac{1}{N}\underset{i=1}{\overset{N}{}}x_i\left(t\right).$$ (33) Given the positions $`x_i\left(t\right)`$ of all the particles (or, equivalently given the function $`h(x,t)`$), the $`x_i\left(t+1\right)`$ become independent random variables. Therefore, given $`h(x,t)`$, the probability for each particle to have at time $`t+1`$ a position strictly larger than $`x`$ is given by $`h(x,t+1)\|h(x,t)`$ (34) $`=1\left[1ph(x1,t)\left(1p\right)h(x,t)\right]^2.`$ The difficulty of the problem comes from the fact that one can only average $`h(x,t+1)`$ over a single time-step. On the right hand side of (34) we see terms like $`h^2(x,t)`$ or $`h(x1,t)h(x,t)`$ and one has to calculate all the correlations of the $`h(x,t)`$ in order to find $`h(x,t+1)`$. This makes the problem very difficult for finite $`N`$. However, given $`h(x,t)`$, the $`x_i\left(t+1\right)`$ are independent and in the limit $`N\mathrm{}`$, the fluctuations of $`h(x,t+1)`$ are negligible. Therefore, when $`N\mathrm{}`$, $`h(x,t)`$ evolves according to the deterministic equation (13). As the initial condition is a step function, we expect the front to move, in the limit $`N\mathrm{}`$, with the minimal velocity $`v_0`$ of (14). For large but finite $`N`$, we expect the correction to the velocity to have two main origins. First, $`h(x,t)`$ takes only values which are integral multiples of $`\frac{1}{N}`$, so that $`\frac{1}{N}`$ plays a role similar to the cut-off $`\epsilon `$ of section 2. Second, $`h(x,t)`$ fluctuates around its average and this has the effect of adding noise to the evolution equation (13). In the rest of this section we present the results of simulations done for large but finite $`N`$ and we will see that the shift in the velocity seems to be very close to the expression of section 4 when $`\epsilon =\frac{1}{N}`$. With the most direct way of simulating the model for $`N`$ finite, it is difficult to study systems of size much larger than $`10^6`$. Here we use a more sophisticated method allowing $`N`$ to become huge. Our method, which handles many particles at the same time, consists in iterating directly $`h(x,t)`$. Knowing the function $`h(x,t)`$ at time $`t`$, we want to calculate $`h(x,t+1)`$. We call respectively $`x_{\text{min}}`$ and $`x_{\text{max}}`$ the positions of the leftmost and rightmost particles at time $`t`$ and $`l=x_{\text{max}}x_{\text{min}}+1`$. In terms of the function $`h(x,t)`$, one has $`0<h(x,t)<1`$ if and only if $`x_{\text{min}}x<x_{\text{max}}`$. Obviously, all the positions $`x_i\left(t+1\right)`$ will lie between $`x_{\text{min}}`$ and $`x_{\text{max}}+1`$. The probability $`p_k`$ that a given particle $`i`$ will be located at position $`x_{\text{min}}+k`$ at time $`t+1`$ is $$p_k=h(x_{\text{min}}+k1,t+1)h(x_{\text{min}}+k,t+1),$$ (35) with $`h(x,t+1)`$ given by (34). Obviously, $`p_k0`$ only for $`0kl`$. The probability to have, for every $`k`$, $`n_k`$ particles at location $`x_{\text{min}}+k`$ at time $`t+1`$ is given by $`P(n_0,n_1,\mathrm{},n_l)=`$ $`{\displaystyle \frac{N!}{n_0!n_1!\mathrm{}n_l!}}p_0^{n_0}p_1^{n_1}\mathrm{}p_l^{n_l}`$ $`\times \delta \left(Nn_0n_1\mathrm{}n_l\right).`$ (36) Using a random number generator for a binomial distribution, expression (36) allows to generate random $`n_k`$. This is done by calculating $`n_0`$ according to the distribution $$P\left(n_0\right)=\frac{N!}{n_0!\left(Nn_0\right)!}p_0^{n_0}\left(1p_0\right)^{Nn_0},$$ (37) then $`n_1`$ with $`P\left(n_1\|n_0\right)`$ $`=`$ $`{\displaystyle \frac{\left(Nn_0\right)!}{n_1!\left(Nn_0n_1\right)!}}\left({\displaystyle \frac{p_1}{1p_0}}\right)^{n_1}`$ (38) $`\times \left(1{\displaystyle \frac{p_1}{1p_0}}\right)^{Nn_0n_1},`$ and so on. This method can be iterated to produce the $`l+1`$ numbers $`n_0`$, $`n_1`$, …, $`n_l`$ distributed according to (36). Then we construct $`h(x,t+1)`$ by $`h(x,t+1)=`$ $`1`$ if $`x<x_{\text{min}}`$, $`h(x,t+1)=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=k+1}{\overset{l}{}}}n_i`$ $`h(x,t+1)=`$ $`0`$ if $`x>x_{\text{max}}+1`$. (39) As the width $`l`$ of the front is roughly of order $`\mathrm{log}N`$, this method allows $`N`$ to be very large. Using this method with the generator of random binomial numbers given in , we have measured the velocity $`v_N`$ of the front for several choices of $`p`$ (0.05, 0.25 and 0.45) and for $`N`$ ranging from 100 to $`10^{16}`$. We measured the velocities with the expression $$v_N=\frac{X_{10^6}X_{10^5}}{\mathrm{9\hspace{0.17em}10}^5}.$$ (40) Figure 3 is a log-log plot of the difference $`v_0v_N`$ versus $`\frac{1}{N}`$ compared to the prediction (7) for $`\epsilon =\frac{1}{N}`$. The variation of $`v_N`$ when using longer times or different random numbers were not larger than the size of the symbols. We see on figure 3 that the speed $`v_N`$ of the front seems to be given for large $`N`$ by $$v_Nv_0\frac{K}{\left(\mathrm{log}N\right)^2},$$ (41) where the coefficient $`K`$ is not too different from the prediction (7). The agreement is however not perfect. The shift $`v_0v_N`$ seems to be proportional to $`\left(\mathrm{log}N\right)^2`$, but the constant looks on figure 3 slightly different from the one predicted by (7). A possible reason for this difference could have been the discretization of the front: instead of only cutting off the tail as in sections 3 and 4, here the whole front $`h(x,t)`$ is constrained to take values multiple of $`\frac{1}{N}`$. One might think that this could explain this discrepancy. However, we have checked numerically (the results are not presented in this paper) that equation (13) with $`h(x,t)`$ constrained to be a multiple of a cut-off $`\epsilon `$ does not give results significantly different from the simpler model of sections 3 and 4 with only a single cut-off. So we think that the full discretization of the front can not be responsible for a different constant $`K`$. The discrepancy observed in figure 3 is more likely due to the effect of the randomness of the process. It is however not clear whether this mismatch would decrease for even larger $`N`$. It would be interesting to push further the numerical simulations and check the $`N`$-dependence of the front velocity for very large $`N`$. ## 6 Conclusion We have shown in the present work that a small cut-off $`\epsilon `$ in the tail of solutions of traveling wave equations has the effect of selecting a single velocity $`v_\epsilon `$ for the front. This velocity $`v_\epsilon `$ converges to the minimal velocity $`v_0`$ when $`\epsilon 0`$ and the shift $`v_0v_\epsilon `$ is surprisingly large (78). Very slow convergences to the minimal velocity have been observed in a number of cases as well as the example of section 5. As the effect of the cut-off on the velocity is large, it is reasonable to think that it would not be much affected by the presence of noise. The example of section 5 shows that the cut-off alone gives at least the right order of magnitude for the shift and it would certainly be interesting to push further the simulations for this particular model to see whether the analysis of section 4 should be modified by the noise. The numerical method used in section 5 to study a very large ($`N10^{16}`$) system was very helpful to observe a logarithmic behavior. We did not succeed to check in earlier works whether the correction was logarithmic, mostly because the published data were usually too noisy or obtained on a too small range of the parameters. Still even if the cut-off was giving the main contribution to the shift of the velocity, other properties would remain very specific to the presence of noise like the diffusion of the position of the front . Our approach of section 4 shows that the effect of a small cut-off is the existence of a scaling form (2030) which describes the change in the shape of the front in its steady state. The effect of initial conditions for usual traveling wave equations (with no cut-off) leads to a very similar scaling form for the change in the shape of the front in the transient regime. This is explained in the appendix where we show how the logarithmic shift of the position of a front due to initial conditions can be recovered. ## A Effect of initial conditions on the position and on the shape of the front In this appendix we show that ideas very similar to those developed in section 4 allow one to calculate the position and the shape at time $`t`$ of a front evolving according to (1), or to a similar equation, given its initial shape. The main idea is that in the long time limit, there is a region of size $`\sqrt{t}`$ ahead of the front which keeps the memory of the initial condition. We will recover in particular the logarithmic shift in the position of the front due to the initial condition , namely that if the initial shape is a step function $`h(x,0)=`$ $`0`$ if $`x>0`$, (A.1) $`h(x,0)=`$ $`1`$ if $`x<0`$, then the position $`X_t`$ of the front at time $`t`$ increases like $$X_t2t\frac{3}{2}\mathrm{log}t.$$ (A.2) More generally, if initially $`h(x,0)=`$ $`x^\nu e^{\gamma _0x}`$ if $`x>0`$, (A.3) $`h(x,0)=`$ $`1`$ if $`x<0`$, we will show that for $`\nu >2`$ $$X_t2t+\frac{\nu 1}{2}\mathrm{log}t,$$ (A.4) whereas the shift is given by (A.2) for $`\nu <2`$. Here, there is no cut-off but the transient behavior in the long time limit gives rise to a scaling function very similar to the one discussed in section 4. If we write the position of the front at time $`t`$ as $$X_t=v_0tc\left(t\right),$$ (A.5) we observed numerically (as in figure 2 of section 3) and we are going to see in the following that the shape of the front takes for large $`t`$ the scaling form $$h(x,t)=t^\alpha G\left(\frac{xX_t}{t^\alpha }\right)e^{\gamma _0\left(xX_t\right)},$$ (A.6) very similar to (2030). If we use (A.5) and (A.6) into the linearized form of equation (1), we get using, the fact that $`v_0=2`$ and $`\gamma _0=1`$, $$\frac{1}{t^\alpha }G^{\prime \prime }+\frac{1}{t^{1\alpha }}\left(\alpha zG^{}\alpha G\right)+t^\alpha \dot{c}G=\dot{c}G^{},$$ (A.7) where $`z=\left(xX_t\right)t^\alpha `$. By writing that the leading orders of the different terms of (A.7) are comparable, we see that we must have $`\alpha `$ $`=`$ $`{\displaystyle \frac{1}{2}},`$ (A.8) $`\dot{c}`$ $``$ $`{\displaystyle \frac{\beta }{t}},`$ (A.9) for some $`\beta `$, and that the right hand side of (A.7) is negligible. Therefore, the equation satisfied by $`G`$ is $$\frac{d^2}{dz^2}G+\frac{z}{2}\frac{d}{dz}G+\left(\beta \frac{1}{2}\right)G=0,$$ (A.10) and the position of the front is given by $$X_tv_0t\beta \mathrm{log}t.$$ (A.11) As in section 4, we expect that as $`t\mathrm{}`$, the front will approach its limiting form and therefore that for $`z`$ small, the shape will look like (15). Therefore we choose the solution $`G_\beta \left(z\right)`$ of (A.10) which is linear at $`z=0`$. This solution can be written as an infinite sum $`G_\beta \left(z\right)`$ $`=A{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\left(1\right)^n}{\left(2n+1\right)!}}z^{2n+1}{\displaystyle \underset{i=0}{\overset{n1}{}}}\left(\beta +i\right),`$ $`=A{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\left(1\right)^n}{\left(2n+1\right)!}}{\displaystyle \frac{\mathrm{\Gamma }\left(n+\beta \right)}{\mathrm{\Gamma }\left(\beta \right)}}z^{2n+1}.`$ (A.12) (The second expression is not valid when $`\beta `$ is a non-positive integer.) To determine $`\beta `$, one can notice that the scaling form (A.6) has to match the initial condition when $`x`$ is large and $`t`$ of order 1. We thus need to calculate the asymptotic behavior of $`G\left(z\right)`$ when $`z`$ is large. For certain values of $`\beta `$, there exist closed expressions of the sum (A.12). For instance, $`G_2\left(z\right)`$ $`=`$ $`A\left(z+{\displaystyle \frac{z^3}{3}}+{\displaystyle \frac{z^5}{60}}\right),`$ $`G_{\frac{7}{2}}\left(z\right)`$ $`=`$ $`A\left(z{\displaystyle \frac{z^3}{3}}+{\displaystyle \frac{z^5}{60}}\right)e^{\frac{z^2}{4}},`$ $`G_1\left(z\right)`$ $`=`$ $`A\left(z+{\displaystyle \frac{z^3}{6}}\right),`$ $`G_{\frac{5}{2}}\left(z\right)`$ $`=`$ $`A\left(z{\displaystyle \frac{z^3}{6}}\right)e^{\frac{z^2}{4}},`$ (A.13) $`G_0\left(z\right)`$ $`=`$ $`Az,`$ $`G_{\frac{3}{2}}\left(z\right)`$ $`=`$ $`Aze^{\frac{z^2}{4}},`$ $`G_1\left(z\right)`$ $`=`$ $`Ae^{\frac{z^2}{4}}{\displaystyle _0^z}e^{\frac{t^2}{4}}𝑑t,`$ $`G_{\frac{1}{2}}\left(z\right)`$ $`=`$ $`A{\displaystyle _0^z}e^{\frac{t^2}{4}}𝑑t,`$ One can check directly on (A.10) that $`G_\beta `$ has a symmetry $$G_\beta \left(z\right)=ie^{z^2/4}G_{\frac{3}{2}\beta }\left(iz\right).$$ (A.14) For any $`\beta `$, one can obtain the large $`z`$ behavior of $`G\left(z\right)`$. To do so, we note that for $`\beta >0`$, one can rewrite (A.12) as $`G_\beta \left(z\right)`$ $`=`$ $`{\displaystyle \frac{A}{\mathrm{\Gamma }\left(\beta \right)}}{\displaystyle _0^{\mathrm{}}}𝑑tt^{\beta \frac{3}{2}}\mathrm{sin}\left(\sqrt{t}z\right)e^t,`$ (A.15) $`=`$ $`{\displaystyle \frac{2A}{\mathrm{\Gamma }\left(\beta \right)}}z^{12\beta }{\displaystyle _0^{\mathrm{}}}𝑑tt^{2\beta 2}\mathrm{sin}\left(t\right)e^{\frac{t^2}{z^2}}.`$ For $`0<\beta <1`$, the second integral in (A.15) has a non zero limit and this gives the asymptotic behavior of $`G_\beta \left(z\right)`$ $$G_\beta \left(z\right)\frac{2A}{\mathrm{\Gamma }\left(\beta \right)}\mathrm{cos}\left(\pi \beta \right)\mathrm{\Gamma }\left(2\beta 1\right)z^{12\beta }.$$ (A.16) From (A.12), one can also show that $$G_\beta ^{\prime \prime }=\frac{\mathrm{\Gamma }\left(\beta +1\right)}{\mathrm{\Gamma }\left(\beta \right)}G_{\beta +1},$$ (A.17) implying that (A.16) remains valid for all $`\beta `$ except for $`\beta =\frac{3}{2}`$, $`\frac{5}{2}`$, $`\frac{7}{2}`$, etc., where the amplitude in (A.16) vanishes. For these values of $`\beta `$, $`G_\beta \left(z\right)`$ decreases faster than a power law (see (A)). The functions $`G_\beta `$ calculated so far are acceptable scaling functions for the shape of the front only for $`\beta \frac{3}{2}`$. Indeed, one can see in (A.16) that for $`\frac{3}{2}<\beta <\frac{5}{2}`$ the function $`G_\beta \left(z\right)`$ is negative for large $`z`$. In fact, for all $`\beta >\frac{3}{2}`$, this function changes its sign at least once, so that the scaling form (A.6) is not reachable for an initial $`h(x,0)`$ which is always positive. It is only for $`\beta \frac{3}{2}`$ that $`G_\beta `$ remains positive for all $`z>0`$. Looking at the asymptotic form (A.16), we see that if initially $`h(x,0)=x^\nu e^{\gamma _0x}`$, the only function $`G_\beta \left(z\right)`$ which has the right large $`z`$ behavior is such that $`12\beta =\nu `$, and this gives, together with (A.11), the expression (A.4) for the shift of the position. As the cases $`\beta >\frac{3}{2}`$ are not reachable, all initial conditions corresponding to $`\nu <2`$ or steeper (such as step functions) give rise to $`G_{\frac{3}{2}}`$ and the shift in position given by (A.2). All the analysis of this appendix can be extended to other traveling wave equations such as (13), with more general functions $`v\left(\gamma \right)`$ (having a non-degenerate minimum at $`\gamma _0`$) as in (14). Then the expressions (A.2A.4) of the shift become $$X_tv_0t\frac{3}{2\gamma _0}\mathrm{log}t$$ (A.18) and $$X_tv_0t\frac{1\nu }{2\gamma _0}\mathrm{log}t.$$ (A.19) We thank C. Appert, V. Hakim and J.L. Lebowitz for useful discussions.
warning/0005/astro-ph0005420.html	ar5iv	text	# Galactic Observations at 31, 42 and 144 GHz with the Mobile Anisotropy Telescope ## 1 Introduction Galactic radiation at centimeter and millimeter wavelengths is attributed to a combination of the following sources of emission: synchrotron, free-free, thermal interstellar dust and variations in the electric and magnetic dipoles of dust (i.e. spinning dust). Large-area, multi-frequency surveys help us to understand the objects and environments which emit in the far-infrared and assist in determining the mechanisms driving microwave emission. The atmosphere is relatively opaque at these frequencies making such surveys difficult to perform from low altitudes. At present, the only comprehensive microwave survey data available is from the Cosmic Background Explorer (COBE) Bennett et al. (1996). It will not be until the launch of the Microwave Anisotropy Probe (MAP) satellite in late 2000 that a full sky survey with pixel resolution smaller than $`1^{}`$ will be completed. Bright millimeter-wavelength sources are important for cosmology as calibrators and a possible Cosmic Microwave Background (CMB) foreground contaminate Toffolatti et al. (1999); Sokasian et al. (1998). There are several open questions about this class of object that include: what are their spectra; are they variable; are they polarized? We address these questions for six Galactic sources using observations made with the Mobile Anisotropy Telescope (MAT). We focus on characterizing emission from unresolved and localized millimeter-bright regions. These data will also help assess daytime atmospheric conditions relevant to observations from similar Chilean sites. We present maps of the Galactic Plane centered on a declination of $`60^{}`$ (approximately $`280^{}<l<305^{}`$ and $`8^{}<b<4^{}`$) covering several bright regions including the Carinae Nebula and RCW 57. All of the millimeter-bright objects we observe have IRAS counterparts or are known HII regions. The IRAS survey indicates that most of these contain substructure unresolved by the MAT observations. Since they cannot be uniquely labeled with previous surveys we denote each bright region with number from 1 to 6. Figure Galactic Observations at 31, 42 and 144 GHz with the Mobile Anisotropy Telescope shows the region covered by these observations. Table 1 lists some of the known radio and millimeter-bright sources observed in other surveys which comprise these regions. The list of associated names gives the brightest far-infrared and radio sources in the vicinity. ## 2 Instrument and Observations MAT is a 0.8 m diameter aperture telescope based on the design in Wollack et al. (1997). Three frequency bands (31, 42, 144 GHz) in eight channels are monitored. The telescope is steerable in azimuth but fixed in elevation. The focal plane contains five corrugated feeds. Observations are made in two polarizations in Ka-band (26-36 GHz) and Q-band (36-46 GHz) using total power high electron mobility transistors (HEMT) amplifiers. Observations are also made in D-band with two single polarization SIS receivers with dual sidebands ($``$ 138-140 GHz and $``$ 148-150 GHz with the local oscillator center at 144 GHz). The calibration and antenna patterns are measured in the field using multiple observations of Jupiter ($`\theta _{disk}`$ 30”-45”). Table 2 lists the MAT beam characteristics for each frequency. Antenna patterns are modeled as two-dimensional Gaussian profiles parameterized by an amplitude and two orthogonal widths<sup>1</sup><sup>1</sup>1FWHM is related to $`\sigma `$ by $`\sigma =\mathrm{\Theta }_{FWHM}/[2(2ln2)^{\frac{1}{2}}]`$. $`\sigma _H`$ and $`\sigma _V`$. The H and V denote the horizontal (chop direction) and vertical (meridian) line with respect to the local horizon. Except for one of the SIS channels (D2) not located in the center of the focal plane, all antenna patterns are well approximated by Gaussian profiles. The center of the focal plane is fixed at an elevation of $`40^{}.76`$. A flat mirror sweeps the beams with an amplitude of $`3^{}`$ on the sky in a sinusoidal pattern at 3.7 Hz. When the mirror is set to zero amplitude, the center of the focal plane is at an azimuth of $`207^{}.41`$ East of North. Each day, MAT views a band about the South Celestial Pole of constant mean declination. To view calibrators and other bright objects, the telescope base is slewed to a new azimuth and the source drifts through the field of view. This procedure typically requires 20 minutes. Observations were carried out at 5200 m from Cerro Toco<sup>2</sup><sup>2</sup>2The Cerro Toco site of the Universidad Católica de Chile was made available through the generosity of Prof. Hernán Quintana , Dept. of Astronomy and Astrophysics. It is near the proposed MMA site. at latitude 22$`\stackrel{}{\mathrm{.}}`$95 S, longitude 67$`\stackrel{}{\mathrm{.}}`$78 W. For this paper we present data obtained between August 30, 1998 and October 14, 1998. The primary goal of the experiment was to use nightly observations far from the Galactic Plane ($`b>15^{}`$) to measure the anisotropy in the CMB radiation Miller et al. (1999); Torbet et al. (1999); Dodelson et al. (1999). This observation scheme also provides a daily view of a large section of the Galactic Plane. ### 2.1 Calibration Jupiter’s emission spectrum is complicated by several molecular lines. However, for wideband observations such as those made by MAT, Jupiter is stable and its intrinsic brightness temperature has been measured to $`5\%`$ Ulich et al. (1981); Griffin et al. (1986); Goldin et al. (1997). We use the values 152 K, 160 K, and 170 K for the 31, 42 and 144 GHz bands respectively. Four terms contribute most to the uncertainty of the calibration of each channel: (1) the standard deviation in the measured solid angle; (2) the standard deviation in the measured temperature of Jupiter for each channel; (3) the intrinsic uncertainty of the brightness temperature of Jupiter; (4) the uncertainty due to calibration drift between Jupiter observations. For the SIS channel used in this analysis, the first three terms contribute uncertainties of 5.5%, 7%, and 5%, respectively. For the HEMT channels, the corresponding values are less than 7%, 10%, and 5%. These terms add in quadrature. The fourth term is primarily the result of diurnal temperature variations. To monitor calibration changes between observations of Jupiter, the HEMT-based channels are injected with a signal from a thermally-stabilized noise source for 40 ms twice every 100 seconds. The SIS-based channels are coupled to a 149 GHz tone at an effective temperature of $`1K`$. These data indicate a 5% random uncertainty in both the SIS and HEMT calibrations and a $`<2.5\%`$ gain change at the time of the Galactic observations. After accounting for the gain change, we estimate the total calibration uncertainty for each channel by adding the uncorrelated errors in quadrature giving a 11% error for the SIS channel and a 14% error for the HEMT channels. The uncertainty quoted in Miller et al. is smaller than that stated here because there was little temperature change, and consequently little gain drift, between the Jupiter observations and the CMB observations. ### 2.2 Pointing Pointing is established through both rising and setting observations of Jupiter. A nominal central azimuth and beam elevation are assumed initially. A Global Positioning System (GPS) time-stamp is stored with each data point and is used to fix the time and location of the centroid of the measured antenna pattern for each horn. Ephemeris data are then used to calculate a position error. The nominal central azimuth and beam elevation<sup>3</sup><sup>3</sup>3Measurements of Galactic sources indicated a tilt in the base of the telescope of $`0^{}.10`$. are then adjusted to minimize the variance between the expected and predicted location of Jupiter. This model allows an accurate reconstruction of pointing over the full azimuthal range of the telescope for the data considered here. After coadding the data for this analysis, the effective center of each of the 5 horns’ antenna patterns is determined with an absolute accuracy of $`0^{}.045`$ and a relative accuracy of $`<0^{}.01`$ degrees. ## 3 Data Reduction We consider a subset of daily observations of the Galactic Plane that contains several millimeter-bright sources. Each of the 8 amplifier channels is processed independently. We do not include the 144 GHz channel D2 since its larger beam greatly decreases its sensitivity to unresolved sources. Of the total collected Galactic data, 54% are rejected due to signal saturation in one or more channels and $`<1\%`$ are rejected based on telemetry drop-outs and noise-injector operation; this procedure leaves approximately 15 observations of 2.68 hours each. Data are coadded for each chopper cycle based on the position of the chopping mirror. The chopper position over a cycle never deviates from the average cycle for the entire observation region by more that 1%. The 144 GHz channels are binned into samples of $`0^{}.045`$ in azimuth for 40 chopper cycles (about 11 seconds). The Ka and Q channels are binned into samples three times larger. By fixing the number of chopper cycles coadded and knowing both the time associated with each data point and the location of each horn antenna pattern (from previous observations of Jupiter), we construct an unfiltered map of the sky. The quality of daytime observations at this time of year is degraded from those at night. This effect is well correlated to the rise in local temperature above $`0^{}`$C. We remove the contribution from atmospheric emission to these Galactic observations using a purely spatial filter. After performing a $`10\sigma `$ cut on isolated points (i.e. inconsistent with the point spread function) we remove a slowly varying polynomial from each coadded azimuthal strip. In this technique, a sliding box-car average is generated. Points greater than $`2.5\sigma `$ as well as the four nearest neighbors are flagged. The sliding box-car average is then regenerated with the flagged points removed from consideration. The results are insensitive to small changes in this flagging level. In all channels, the box-car is 2 times the size of the FWHM of the antenna pattern for that channel. This second average is subtracted from the unfiltered data. Monte Carlo analysis on simulated maps constructed from measured atmospheric data at this site show that this filter effects measurements of the amplitude and width of a Gaussian beam by $`<1\%`$ for the atmospheric conditions encountered in this data set. The largest adverse effects occur at the edge of the maps where incomplete coverage of a source can be confused for an atmospheric variation. This occurs with two sources for one feed horn (Q3&4) and, hence, the measured source amplitudes for these channels have been omitted from all averages. ## 4 Maps After filtering, each daily map is coadded with observations from other days to produce final images of the region. Each map is comprised of between 10 and 15 coadded viewings. No attempt has been made to coadd maps made from different amplifier channels. We show representative plots of the Galactic region at 31, 42 and 144 GHz in Figure Galactic Observations at 31, 42 and 144 GHz with the Mobile Anisotropy Telescope, Galactic Observations at 31, 42 and 144 GHz with the Mobile Anisotropy Telescope, and Galactic Observations at 31, 42 and 144 GHz with the Mobile Anisotropy Telescope respectively. Also shown are close-ups of each source region for the frequencies at which they were observed. Each source close-up is comprised of color contours and an intensity color-bar associated with the MAT data. IRAS<sup>4</sup><sup>4</sup>4Released data product at http://www.ipac.caltech.edu/ intensity contours have also been superimposed to demonstrate relative pointing and additional bright sources in the region not seen at all of the MAT frequencies. We are able to place constraints on both polarization and variability of those sources observed at 31 and 42 GHz. Though all detector channels are polarized, the antenna response of orthogonal polarizations in Ka and Q bands are matched. This allows the fractional polarization $`2(E_{\|\|}E_{})/(E_{\|\|}+E_{})`$ over the source to be constrained as well as test for systematic errors. Table 3 gives the measured fractional polarization limits for each source over a $`3\sigma `$ region. All sources are unpolarized to the level of our sensitivity. Nearly all compact radio sources demonstrate some variability on times scales ranging from a few days to years. The variability in emission for the MAT sources is constrained over the observing period in all frequency bands. We place limits on the daily and total variation of peak brightness temperatures for all sources observed at 31 and 42 GHz except MAT 4. Due to the lower signal-to-noise of MAT 4 and all sources at 144 GHz we group these observations into 3 subsets in time; the limits are only on the total variation. Daily variation is defined as the standard deviation of all observations about a line. Total variation is the standard deviation on the mean value of all observations. Table 4 shows the limits to which all sources are stable over the observing period. No statistically significant change in brightness is indicated over any time frame. ## 5 MAT Sources: Discussion The six MAT sources are extended at $`100\mu `$m in the IRAS maps and cannot be unambiguously resolved in any of the MAT beams. For this reason, the flux density from each identified bright region is determined within a circle whose radius is the average $`1\sigma `$ and $`2\sigma `$ beam width for that frequency. That is, the flux density at frequency $`\nu `$ is: $$S_\nu =_{n\sigma }\frac{2kT(x,y)\nu ^2}{c^210^{26}}𝑑x𝑑y\mathrm{Jy}$$ (1) where $`n\sigma `$ denotes the area (number of $`\sigma `$) over which to integrate and T(x,y) is the measured temperature profile over this region. We also quote the amplitude (i.e. peak brightness temperature) of the best-fit Gaussian profile (widths unconstrained) for each source. Table 5 details these results. We average the measurements from all channels of a given frequency. The results at 42 GHz are rescaled to a single beam size of $`\mathrm{\Omega }_{42GHz}=16.5\times 10^5`$str before averaging. Due to the low signal-to-noise, Q2 is not included in the 42 GHz average. The quoted uncertainty results from the quadrature sum of the uncorrelated components of the error divided by the number of channels coadded. In addition to calibration error, the sum includes an uncertainty in the baseline of the filtered maps of 4% for the SIS channel and 3% for the HEMT channels. MAT 1 is an extended HII region containing a single $`100\mu `$m-bright IRAS point source. MAT 2 is a compact HII region. Since they are observed only at 42 GHz we cannot constrain the spectral indices. MAT 4 is also an extended HII region containing two moderately bright $`100\mu `$m-bright IRAS point sources. It is seen at 31 and 42 GHz but is not measurable at 144 GHz. The implied spectral index based on the low frequency observations is $`\alpha =0.72`$ where $`S\nu ^\alpha `$. In this frequency range, the emission is dominated by synchrotron radiation. MAT 3 is dominated by emission from the Keyhole Nebula. This region contains the known millimeter-bright and variable remnant $`\eta `$Carinae. Cox et al. (1995) have stated that at wavelengths greater than 1 cm, $`\eta `$Carinae is dominated by free-free emission but that between 1 cm and 1 mm the spectral index rises to $`1`$. Cox et al. also give evidence of time variability in a one year period of $`30\%`$. Since our beam size at 144 GHz is significantly larger, we are not able to constrain the emission of $`\eta `$Carinae but rather the emission from the larger Keyhole Nebula in which it is contained. This region also contains a second source (IRAS 10414-5919) known to be bright at both 4.8 GHz Condon et al. (1995) and $`100\mu `$m. Therefore, there is no reason to expect that our map would be centered on $`\eta `$Carinae. Given the extended emission of this region we have not attempted to constrain the spectral index. MAT 5 and MAT 6 both contain two IRAS-bright sources which are clearly seen at 144 GHz. However, the 31 and 42 GHz observations indicate that only one source dominates emission at lower frequencies implying a highly inverted spectrum. For MAT 6a, if we use the flux limit at 42 GHz in conjunction with the measured flux at 144 GHz assuming it is an unresolved source we find a spectral index of $`\alpha 1.5`$ where $`S\nu ^\alpha `$. For MAT 5a, we find $`\alpha 2.0`$. Both are consistent with the classical spectra of compact HII regions where $`S\nu ^2`$. We would like to thank Suzanne Staggs, Dave Rusin, Simon Dicker, Jeff Klein, and Neill Reid for helpful conversations and assistance in preparing this manuscript. This work was supported by an NSF NYI award, a Cottrell Award from the Research Corporation, a David and Lucile Packard Fellowship (to LP), a NASA GSRP fellowship to AM, an NSF graduate fellowship to MN, a NSF Career award (AST-9732960, to MD), NSF grants PHY-9222952, PHY-9600015, AST-9732960, and the University of Pennsylvania. The data will be made public upon publication of this article.
warning/0005/math-ph0005003.html	ar5iv	text	# Untitled Document LAVAL-PHY-00-22 Generating-function method for tensor products. L. Bégin<sup>1</sup> Work supported by NSERC (Canada)., C. Cummins<sup>♯2</sup> and P. Mathieu<sup>2</sup> Work supported by NSERC (Canada) and FCAR (Québec). Département de Physique, Université Laval, Québec, Canada G1K 7P4 CICMA, Department of Mathematics and Statistics, Concordia University, Montréal, Québec, Canada H3G 1M8 Abstract: This is the first of two articles devoted to a exposition of the generating-function method for computing fusion rules in affine Lie algebras. The present paper is entirely devoted to the study of the tensor-product (infinite-level) limit of fusions rules. We start by reviewing Sharp’s character method. An alternative approach to the construction of tensor-product generating functions is then presented which overcomes most of the technical difficulties associated with the character method. It is based on the reformulation of the problem of calculating tensor products in terms of the solution of a set of linear and homogeneous Diophantine equations whose elementary solutions represent “elementary couplings”. Grobner bases provide a tool for generating the complete set of relations between elementary couplings and, most importantly, as an algorithm for specifying a complete, compatible set of “forbidden couplings”. 11/98 (revised 06/99,01/00) 1. Introduction 1.1. Orientation Fusion rules yield the number of independent couplings between three given primary fields in conformal field theories. We are interested in fusion rules in unitary conformal field theories that have a Lie group symmetry, that is, those whose generating spectrum algebra is an affine Lie algebra at integer level. These are the Wess-Zumino-Witten models . Primary fields in these cases are in 1-1 correspondence with the integrable representations of the appropriate affine Lie algebra at level $`k`$. Denote this set by $`P_+^{(k)}`$ and a primary field by the corresponding affine weight $`\widehat{\lambda }`$. Fusion coefficients $`𝒩_{\widehat{\lambda }\widehat{\mu }}^{(k)}^{\widehat{\nu }}`$ are defined by the product $$\widehat{\lambda }\times \widehat{\mu }=\underset{\nu P_+^{(k)}}{}𝒩_{\widehat{\lambda }\widehat{\mu }}^{(k)}{}_{}{}^{\widehat{\nu }}\widehat{\nu }$$ $`(1.1)`$ (For a review of conformal field theory and in particular fusion rules, see ; to a large extend we follow the notation of this reference.) In the infinite-level limit and for fields with finite conformal dimensions, the purely affine condition on weight integrability is relaxed and the primary fields are solely characterised by their finite part, required to be an integrable weight of the corresponding finite Lie algebra. Recall that a finite weight $`\lambda `$ is characterised by its expansion coefficients in terms of the fundamental weights $`\omega _i`$ $$\lambda =\underset{i=1}{\overset{r}{}}\lambda _i\omega _i=(\lambda _1,\mathrm{},\lambda _r)$$ $`(1.2)`$ where $`r`$ is the rank of the algebra. The numbers $`\lambda _i`$’s are the Dynkin labels. The set of weights with non-negative Dynkin labels (the integrable weights) is denoted by $`P_+`$. In the infinite-level limit the fusion coefficients reduce to tensor-product coefficients: $$\underset{k\mathrm{}}{lim}𝒩_{\widehat{\lambda }\widehat{\mu }}^{(k)}{}_{}{}^{\widehat{\nu }}=𝒩_{\lambda \mu }^{}{}_{}{}^{\nu }.$$ $`(1.3)`$ where $`𝒩_{\lambda \mu }^{}{}_{}{}^{\nu }`$ is defined by $$\lambda \mu =\underset{\nu P_+}{}𝒩_{\lambda \mu }^{}{}_{}{}^{\nu }\nu $$ $`(1.4)`$ By abuse of notation, we use the same symbol for the highest weight and the highest-weight representation. Notice that $$𝒩_{\lambda \mu }{}_{}{}^{\nu }=𝒩_{\lambda \mu \nu ^{}}$$ $`(1.5)`$ where $`\nu ^{}`$ denotes the highest weight of the representation conjugate to that of $`\nu `$. Equivalently, $`𝒩_{\lambda \mu \nu ^{}}`$ gives the multiplicity of the scalar representation in the triple product $`\lambda \mu \nu ^{}`$. A tensor-product generating function codes the information for all the tensor products of a given algebra in a single function defined by $$G(L,M,N)=\underset{\lambda ,\mu ,\nu P_+}{}𝒩_{\lambda \mu }^{}{}_{}{}^{\nu }L^\lambda M^\mu N^\nu $$ $`(1.6)`$ where $`L^\lambda =L_1^{\lambda _1}\mathrm{}L_r^{\lambda _r}`$ and similarly for $`M^\mu `$ and $`N^\nu `$. $`G`$ can generally be expressed as a simple closed function of its variables. For instance, for $`su(2)`$, it reads $$G(L,M,N)=\frac{1}{(1LM)(1LN)(1MN)}$$ $`(1.7)`$ An example of basic global information that can be deduced from a generating function is the integrality as well as the positivity of the tensor-product coefficients. More importantly, from our point of view, is that in the context of fusion rules, the construction of the simplest generating functions led to the discovery of the notion of threshold levels . Moreover, as shown in the sequel paper, setting up a fusion generating function is a way to obtain explicit expressions for these threshold levels. Our new approach to fusion-rule generating functions, which originates from the generalisation of techniques developed in the present paper on tensor products, leads to a further new concept, that of a fusion basis. 1.2. Overview of the paper The present article is organised as follows. We start by explaining in detail the construction of tensor-product generating functions for finite Lie algebras. The first construction which is presented is the character method developed by Sharp and his collaborators (section 2). Although it is conceptually very simple, the character method is limited by its inherent computational difficulties: the disproportion between the simplicity of the resulting form of the generating function and the intermediate calculations is enormous. This motivates our alternative approach to the construction of tensor-product generating function. It is based on the reformulation of the problem of calculating tensor products in terms of the solution of a set of linear and homogeneous Diophantine equations (cf. section 3). The elementary solutions of these Diophantine equations represent “elementary couplings”. For $`sp(4)`$, the use of the Berenstein-Zelevinsky inequalities to obtain the elementary couplings and their relations (cf. the analysis of section 6) is new. The key difficulty is finding the numerous relations that exist in general between the elementary solutions. From the Diophantine-equation point of view, the decomposition of a solution may not be unique because different sums of elementary solutions could yield the same result. To solve this problem we first “exponentiate” it: given a solution $`\alpha =(\alpha _1,\mathrm{},\alpha _k)`$ to our system of linear Diophantine equations, we introduce formal variables $`X_1,\mathrm{},X_k`$ and consider the monomial $`X_1^{\alpha _1}\mathrm{}X_k^{\alpha _k}`$. The linear span, $`R`$, of all such monomials is a “model” for the generating function for the solutions to the original set of linear Diophantine equations (see section 5), since the Poincaré series of $`R`$ is the required generating function. This series can be calculated using Grobner basis methods. For $`su(N)`$ there is a remarkable graphical construction for computing tensor product multiplicities, the famous Berenstein-Zelevinsky triangles. These are introduced in section 6. We also discuss the analogous construction for $`sp(4)`$, whose diagrammatic representation is new. But the main interest of these re-formulations is that it yields a simple and systematic way of obtaining the elementary couplings from the construction of a vector basis. Thus we get a new way of constructing the corresponding generating functions. 2. Generating-function for tensor products: the character method 2.1. The character method for the construction of the tensor-product generating function: the $`su(2)`$ case The method developed by Sharp and collaborators for constructing generating functions for tensor products is based on manipulations of the character generating functions . Although simple in principle, these manipulations become rather cumbersome as the rank of the algebra is increased. To illustrate the method, we will work in complete detail the simplest example, the $`su(2)`$ case. The first step is the derivation of the character generating function. The Weyl character formula for a general algebra of rank $`r`$ and a highest-weight representation $`\lambda `$ is $$\chi _\lambda =\frac{\xi _{\lambda +\rho }}{\xi _\rho }$$ $`(2.1)`$ where $`\rho `$ is the finite Weyl vector, $`\rho =_{i=1}^r\omega _i`$, and where the characteristic function $`\xi `$ is defined as $$\xi _{\lambda +\rho }=\underset{wW}{}ϵ(w)e^{w(\lambda +\rho )}$$ $`(2.2)`$ where $`ϵ(w)`$ is the signature of the Weyl reflection $`w`$ and $`W`$ is the Weyl group. For $`su(2)`$, $`W`$ contains two elements: $`1,s_1`$. With $$x=e^{\omega _1}$$ $`(2.3)`$ the $`su(2)`$ characteristic function $`\xi `$ for the representation of highest weight $`m\omega _1(m)`$ is $$x^{m+1}x^{m1}$$ $`(2.4)`$ The character reads then $$\chi _m=\frac{x^{m+1}x^{m1}}{xx^1}=\frac{x^mx^{m2}}{1x^2}=x^m+x^{m2}+\mathrm{}+x^m$$ $`(2.5)`$ The character generating function $`\chi _L`$ is obtained by multiplying the above expression by $`L^m`$ where $`L`$ is a dummy variable, and summing over all positive values of $`m`$: $$\begin{array}{cc}\hfill \chi _L(x)& =\underset{0}{\overset{\mathrm{}}{}}L^m\chi _m=\frac{1}{xx^1}\underset{0}{\overset{\mathrm{}}{}}L^m(x^{m+1}x^{m1})\hfill \\ & =\frac{1}{1x^2}\left(\frac{1}{1Lx}\frac{x^2}{1Lx^1}\right)=\frac{1}{(1Lx)(1Lx^1)}\hfill \end{array}$$ $`(2.6)`$ We should point out here that in all generating functions in this paper, expressions of the form $`1/(1a)`$ should be formally expanded in positive powers of $`a`$. So for example, $`1/(1Lx^1)=1+Lx^1+L^2x^2+\mathrm{}`$. By construction, the character of the highest weight $`(m)`$ can be recovered from the power expansion of $`\chi _L`$ as the coefficient of the term $`L^m`$. The characteristic generating function $`\xi _L`$ is defined by $$\chi _L(x)=\frac{\xi _L}{\xi _0}$$ $`(2.7)`$ and it reads $$\xi _L(x)=\frac{xx^1}{(1Lx)(1Lx^1)}=\frac{x}{1Lx}\frac{x^1}{1Lx^1}$$ $`(2.8)`$ the last form being the one that results directly from (2.5). The tensor product of two highest-weight representations can be obtained from the product of the corresponding characters: $$\chi _m\chi _n=\underset{\mathrm{}}{}𝒩_{mn}^{}{}_{}{}^{\mathrm{}}\chi _{\mathrm{}}$$ $`(2.9)`$ This information can be extracted from the product of the corresponding generating functions. We are thus led to consider the product $`\chi _L(x)\chi _M(x)`$. To simplify the analysis of the resulting expression, notice that the information concerning the representations occurring in the tensor product is coded in the leading term of the character, i.e., the term $`x^{m+1}`$. To insure that every positive power of $`x`$ singles out a highest-weight representation, we can multiply both sides by $`\xi _0`$. To read off these terms, we can focus on the terms with strictly positive powers of $`x`$ in the product $`\chi _L(x)\chi _M(x)\xi _0(x)`$. If we require the Dynkin label of the representations (and not their shifted value), it is more convenient to divide by $`x`$ before doing the projection, now restricted to the non-negative powers of $`x`$. The truncation of an expression by its negative powers of $`x`$ will be denoted by the MacMahon symbol $`\mathrm{\Omega }`$, defined by $$\underset{}{\overset{x}{\Omega }}\underset{\mathrm{}}{\overset{\mathrm{}}{}}c_nx^n=\underset{n0}{}c_nx^n$$ $`(2.10)`$ When there is no ambiguity concerning the variable in terms of which the projection is defined, it is omitted from the $`\mathrm{\Omega }`$ symbol. We are thus interested in the projection of the following expression $$\begin{array}{cc}\hfill \chi _L(x)\chi _M(x)\xi _0(x)x^1& =\chi _L(x)\xi _M(x)x^1\hfill \\ & =\frac{1}{(1Lx)(1Lx^1)}\left(\frac{1}{1Mx}\frac{x^2}{1Mx^1}\right)\hfill \end{array}$$ $`(2.11)`$ For these manipulations, we use systematically the following simple identities: $$\begin{array}{cc}\hfill \frac{1}{(1Ax)(1Bx^1)}& =\frac{1}{(1AB)}\left(\frac{1}{1Ax}+\frac{Bx^1}{1Bx^1}\right)\hfill \\ & =\frac{1}{(1AB)}\left(\frac{Ax}{1Ax}+\frac{1}{1Bx^1}\right)\hfill \\ & =\frac{1}{(1AB)}\left(\frac{1}{1Ax}+\frac{1}{1Bx^1}1\right)\hfill \end{array}$$ $`(2.12)`$ There are two terms to analyse. The first is $$\frac{1}{(1Lx)(1Lx^1)(1Mx)}=\frac{1}{(1Lx)(1LM)}\left(\frac{1}{1Mx}+\frac{Lx^1}{1Lx^1}\right)$$ $`(2.13)`$ The first part is not affected by the projection and the second can be written as $$\frac{Lx^1}{(1Lx)(1LM)(1Lx^1)}=\frac{Lx^1}{(1LM)(1L^2)}\left(\frac{Lx}{1Lx}+\frac{1}{1Lx^1}\right)$$ $`(2.14)`$ The second term of this expression contains only negative powers of $`x`$ and can thus be ignored and the first part is unaffected by the projection. We have thus, for the first term of (2.11) $$\underset{}{\Omega }\frac{1}{(1Lx)(1Lx^1)(1Mx)}=\frac{1}{(1Lx)(1LM)}\left(\frac{1}{1Mx}+\frac{L^2}{1L^2}\right)$$ $`(2.15)`$ The projection of the second term of (2.11) is: $$\begin{array}{cc}& \underset{}{\Omega }\frac{x^2}{(1Lx)(1Lx^1)(1Mx^1)}\hfill \\ & =\underset{}{\Omega }\frac{x^2}{(1Lx^1)(1LM)}\left(\frac{1}{1Lx}+\frac{Mx^1}{1Mx^1}\right)\hfill \\ & =\underset{}{\Omega }\frac{x^2}{(1Lx^1)(1LM)(1Lx)}\hfill \\ & =\underset{}{\Omega }\frac{x^2}{(1LM)(1L^2)}\left(\frac{Lx}{1Lx}+\frac{1}{1Lx^1}\right)\hfill \\ & =\underset{}{\Omega }\frac{Lx^1}{(1LM)(1L^2)(1Lx)}\hfill \\ & =\frac{Lx^1}{(1LM)(1L^2)}\left(\frac{1}{1Lx}1\right)\hfill \\ & =\frac{L^2}{(1LM)(1L^2)(1Lx)}\hfill \end{array}$$ $`(2.16)`$ Subtracting (2.16) from (2.15), we find that $$\underset{}{\Omega }\chi _L(x)\xi _M(x)x^1=\frac{1}{(1LM)(1Lx)(1Mx)}$$ $`(2.17)`$ Replacing $`x`$ by $`N`$, we thus get $$G^{su(2)}(L,M,N)=\frac{1}{(1LM)(1LN)(1MN)}$$ $`(2.18)`$ 2.2. The abstract setting: Poincaré series, elementary couplings and relations; defining a model As we shall see it is frequently useful have a model, $`R`$, for a generating function $`G(X_1,\mathrm{},X_k)`$ such as (2.18). By this we mean a commutative $`Q`$-algebra with an identity, graded by $`N^k`$, ( $`N=\{0,1,2,3,\mathrm{}\}`$) $$R=_{\alpha N^k}R_\alpha ,R_\alpha R_\beta R_{\alpha +\beta }$$ $`(2.19)`$ and such that its Poincaré series (also frequently called Hilbert series) $$F(R)=\underset{\alpha N^k}{}dim_Q(R_\alpha )X^\alpha $$ satisfies $$F(R)=G(X_1,\mathrm{},X_k).$$ $`(2.20)`$ For example, for (2.18), with $`X_1=L,X_2=M,X_3=N`$, we can take $`R=Q[E_1,E_2,E_3]`$, which is the polynomial ring generated by the formal variables $`E_1,E_2,E_3`$ (in fact all our examples $`R`$ is either a subring or quotient of a polynomial ring) with the grading of $`E_1,E_2`$ and $`E_3`$ being $`(1,1,0)`$, $`(1,0,1)`$ and $`(0,1,1)`$. The homogeneous subspaces are spanned by $`E_1^aE_2^bE_3^c`$, $`a,b,cN`$ with grade $`(a+b,a+c,b+c)`$ and so $$F(R)=\underset{(a,b,c)N^3}{}X_1^{a+b}X_2^{a+c}X_3^{b+c}=G^{su(2)}(X_1,X_2,X_3)$$ $`(2.21)`$ as required. If $`R`$ is generated by elements $`E_1,\mathrm{},E_s`$ and is a model for a generating function $`G`$ for tensor products (or fusion products) then we call $`E_1,\mathrm{},E_s`$ a set of “elementary couplings” for $`G`$. It should perhaps be stressed that a priori the variables $`X_1,\mathrm{},X_k`$ and $`E_1,\mathrm{},E_s`$ are unrelated. We shall refer to the $`E`$’s as model variables and the $`X`$’s as grading variables. If the grading vector of $`E_i`$ is $`\alpha ^i`$ , $`i=1,\mathrm{},s`$ then there is an associated monomial in the grading variables: $`X^{\alpha ^i}`$, for which we will use the notation $`g(E_i)`$. For example in the above example we have $`g(E_1)=X_1^1X_2^1X_3^0=LM`$. However, to avoid tedious repetition when writing down generating functions we shall often write, for example, $`1/(1E_1)`$ rather than $`1/(1g(E_1))`$. In all such cases where model variables appear in a generating function they should be replaced by the corresponding monomial in the grading variables. In the case of tensor products we use the notation “$`E:g(E):\text{product}`$” to denote a set of elementary couplings with their “exponentiated” grading and the corresponding term in the tensor product. So in the example above we would write: \| \| $`E_1:LM:`$ \| $`(1)(1)(0),`$ \| \| --- \| --- \| --- \| \| \| $`E_2:LN:`$ \| $`(1)(0)(1),`$ \| \| \| $`E_3:MN:`$ \| $`(0)(1)(1)`$ \| $`(2.22)`$ Having made the distinction between grading and model variables, it should be noted that there are cases where we can identify the model as a ring generated by monomials in the grading variables. So in the above example we could define $`E_1=LM`$, $`E_2=LN`$ and $`E_3=MN`$ and take the model for our generating function to be the subring of $`Q[L,M,N]`$ generated by $`E_1,E_2`$ and $`E_3`$. However, it is not always desirable, or even possible, to make this identification. We close this section with two examples of how models for the $`su(2)`$ character generating function can be constructed. The first method, which has been exploited by Sharp et al (see ) to construct character generating functions, amounts to finding an algebra $`R`$ which is a module for the Lie algebra $`su(2)`$ and such that, as an $`su(2)`$ module, $`R`$ is isomorphic to $`_{i1}V_i`$ where $`V_i`$ is the irreducible $`su(2)`$ module of dimension $`i`$. In this case we can take $`R=Q[p,q]`$ with the generators of $`su(2)`$ being given by differential operators: $$h=p\frac{}{p}q\frac{}{q},x_{}=q\frac{}{p},x_+=p\frac{}{q}$$ $`(2.23)`$ The $`su(2)`$ highest-weight vectors are $`p^i`$, $`i0`$ and a basis of the irreducible submodule of dimension $`i`$ is just given by the monomials of degree $`i`$ in $`p`$ and $`q`$. We can give $`R`$ an $`N^3`$ grading by taking the degree of $`p`$ to be $`(1,1,0)`$ and of $`q`$ to be $`(1,0,1)`$. Here the first grading index specifies the representation while the other two refer to a particular weight. As $`R=Q[p,q]`$ the Poincaré function for $`R`$ is, $$\frac{1}{(1p)(1q)}$$ with the understanding, as explained above, that $`p`$ and $`q`$ should be replaced by the corresponding expression in terms of the grading variables. Let us denote these grading variables here by $`L`$ (which exponentiates the representation index) and $`x,y`$ (exponentially related to the weights). The Poincaré function reads then $$\frac{1}{(1Lx)(1Ly)}$$ $`(2.24)`$ Another way of constructing a model for the weight generating function, which makes more natural the $`N^3`$ grading, is to observe that the complete set $`SU(2)`$ weight vectors of finite dimensional irreducible $`su(2)`$ modules are in 1-1 correspondence with one-rowed Young tableaux. If the Young tableau has $`c`$ boxes filled with $`a`$ 1’s and $`b`$ 2’s then there is a constraint $$a+bc=0,a,b,c0$$ $`(2.25)`$ and so the solutions to this linear Diophantine equation are in 1-1 correspondence with the complete set of $`SU(2)`$ weight vectors. Thus to find a model for the weight generating function it is sufficient to find a model for the solutions to (2.25). It is not difficult to see that every solution to this equation is a linear combination (with non-negative coefficients) of the two fundamental solutions: $`(a,b,c)=(1,0,1)`$ and $`(a,b,c)=(0,1,1)`$. Let $`R`$ be the subring of $`Q[A,B,C]`$ generated by the monomials $`E_1=AC,E_2=BC`$. Considering the exponents of the monomials $`E_1`$ and $`E_2`$, we see that the monomials in $`R`$ correspond to the solutions of (2.25) and hence taking the natural grading on $`R`$ ensures that the Poincaré series of $`R`$ is the generating function for the solutions to (2.25) and hence is the required generating function. In this example there are no relations between $`E_1`$ and $`E_2`$ and so $`R`$ is isomorphic to the polynomial ring in two variables (as expected) and so the Poincaré function is once again (with $`Ax,By,CL`$) given by (2.24). 2.3. Multiple $`su(2)`$ tensor products In order to illustrate the occurrence of linear relations between elementary couplings, consider the problem of finding the multiplicity of a given representation $`\zeta `$ in the triple product $`\lambda \mu \nu `$. In terms of character generating functions, this amounts to considering the product $`\chi _L(x)\chi _M(x)\chi _N(x)\chi _P(x)`$, or equivalently, $`\chi _L(x)\chi _M(x)\xi _N(x)x^1\xi _P(x)x^1`$. The left side is then projected onto positive powers of $`x`$. We are thus led to consider $$\underset{}{\Omega }\frac{1}{(1Lx)(1Lx^1)(1Mx)(1Mx^1)}\left(\frac{1}{1Nx}\frac{x^2}{1Nx^1}\right)$$ $`(2.26)`$ The projection of each term is worked out as previously and the resulting expression is found to be, with $`x`$ replaced by $`P`$: $$G(L,M,N,P)=\frac{1LMNP}{(1LP)(1MP)(1NP)(1LM)(1LN)(1MN)}$$ $`(2.27)`$ This is the sought for generating function. Here we would like to have a model with 6 elementary couplings corresponding to the terms in the denominator of the generating function: \| \| $`E_1:LM:`$ \| $`(1)(1)(0)(0)`$ \| \| --- \| --- \| --- \| \| \| $`E_2:LN:`$ \| $`(1)(0)(1)(0)`$ \| \| \| $`E_3:LP:`$ \| $`(1)(0)(0)(1)`$ \| \| \| $`E_4:MN:`$ \| $`(0)(1)(1)(0)`$ \| \| \| $`E_5:MP:`$ \| $`(0)(1)(0)(1)`$ \| \| \| $`E_6:NP:`$ \| $`(0)(0)(1)(1)`$ \| $`(2.28)`$ and there must be a linear relation (in this context, such a relation is often called a syzygy in the physics literature - see in particular and related works) between the following products (signalled by a term in the numerator) which has grading $`LMNP`$: $$E_1E_6,E_2E_5,E_3E_4$$ $`(2.29)`$ It is not difficult to see that a model is given by $`Q[e_1,e_2,e_3,e_4,e_5,e_6]/I`$ where $`E_i=e_i+I,i=1,\mathrm{},6`$ and $`I=(ae_1e_6+be_2e_5+ce_3e_4)`$ is the ideal generated by the polynomial $`ae_1e_6+be_2e_5+ce_3e_4`$ for any choice of $`a,b,cQ`$ not all zero. The elements of $`R`$ have the form $`m+I`$ with $`mQ[e_1,\mathrm{},e_6]`$. However there is no canonical way of choosing the representatives $`m`$. Take for example the case $`a=b=c=1`$. (Usually we will construct a model for our generating function as explained above and this construction will fix the values of $`a,b`$ and $`c`$). In $`R`$ we have $`E_1E_6=(E_2E_5+E_3E_4)`$ and so we can take as a basis for $`R`$ the set of (equivalences classes of ) monomials which do not contain the product $`E_1E_6`$. In this case we say that we have chosen to make $`E_1E_6`$ a ‘forbidden product’. Similarly we can forbid the products $`E_2E_5`$ or $`E_3E_4`$. As we shall see later, the choice of forbidden products corresponds to a choice of term ordering. Before leaving this example, we would like to rework it from a different point of view, as an illustration of the ‘composition’ technique of generating functions. Let $`G(L,M,R)`$ describe the tensor product corresponding to $`\chi _L\chi _M\chi _R`$ and similarly let $`G(Q,N,P)`$ correspond to $`\chi _Q\chi _N\chi _P`$. We are interested the product $`\chi _L(x)\chi _M(x)\chi _N(x)\chi _P(x)`$, but treated from the product of the two generating functions $`G`$. We thus want to enforce the constraint $`R=Q`$ in the product $`G(L,M,R)G(Q,N,P)`$. The idea – which is used in the references in mainly in relation with the construction of generating functions for branching functions – is to multiply this product by $`(1Q^1R^1)^1`$ and, in the expansion in powers of $`R`$ and $`Q`$, keep only terms of order zero in both variables: with an obvious notation we have $$\begin{array}{cc}\hfill \underset{=}{\overset{R}{\Omega }}\underset{=}{\overset{Q}{\Omega }}& G(L,M,R)G(Q,N,P)\frac{1}{1Q^1R^1}\hfill \\ & =\underset{=}{\overset{R}{\Omega }}\underset{=}{\overset{Q}{\Omega }}\underset{n}{}A_n(L,M)R^n\underset{m}{}B_m(N,P)Q^m\underset{\mathrm{}}{}R^{\mathrm{}}Q^{\mathrm{}}\hfill \\ & =\underset{p}{}A_p(L,M)B_p(L,M)\hfill \end{array}$$ $`(2.30)`$ which is manifestly equivalent to considering $$\underset{=}{\overset{x}{\Omega }}G(L,M,x)G(x^1,N,P)$$ $`(2.31)`$ With the explicit expressions for the generating functions, we have thus $$\underset{=}{\overset{x}{\Omega }}\frac{1}{(1Lx)(1Mx)(1LM)}\frac{1}{(1Px^1)(1Nx^1)(1NP)}$$ $`(2.32)`$ A brief and by now standard analysis yields directly the generating function (2.27). 2.4. The $`sp(4)`$ case As a final example, consider the $`sp(4)`$ case. With the $`x_i=e^{\omega _i},i=1,2`$, the characteristic function is found to be $$\begin{array}{cc}\hfill \xi _{(m,n)}& =x_1^{m+1}x_2^{n+1}x_1^{m1}x_2^{m+n+2}x_1^{n+m+5}x_2^{n1}+x_1^{m+2n+3}x_2^{mn2}\hfill \\ & +x_1^{m2n3}x_2^{n+m+2}x_1^{m+1}x_2^{mn2}x_1^{m2n3}x_2^n+x_1^{m1}x_2^{n1}\hfill \end{array}$$ $`(2.33)`$ and the characteristic generating function is $$\begin{array}{cc}\hfill \xi _{L_1,L_2}& =\frac{1}{(1L_1x_1)(1L_1x_1x_2^1)(1L_2x_2^1)(1L_2x_1^2x_2)}\hfill \\ & \times (\frac{1+L_2}{(1L_2x_1^2x_2^1)(1L_2x_2^1)}+\frac{(1+L_2)L_1x_1}{(1L_1x_1)(1L_2x_1^2x_2^1)}\hfill \\ & +\frac{L_1x_1^1x_2}{(1L_1x_1)(1L_1x_1^1x_2)})\hfill \end{array}$$ $`(2.34)`$ From this we construct the character generating function and then we can proceed to the tensor-product generating function. This is again extremely cumbersome. The result is $$\begin{array}{cc}& G^{sp(4)}(L_1,L_2,M_1,M_2,N_1,N_2)\hfill \\ & =[(1M_1N_1)(1L_1N_1)(1L_1M_1)(1M_2N_2)(1L_2N_2)(1L_2M_2)]^1\hfill \\ & \times (\frac{1}{(1L_2M_1N_1)(1L_2M_1^2N_2)}+\frac{L_2M_2N_1^2}{(1L_2M_1N_1)(1L_2M_2N_1^2)}\hfill \\ & +\frac{L_1^3M_2^2N_1N_2}{(1L_1M_2N_1)(1L_1^2M_2N_2)}+\frac{L_1M_2N_1}{(1L_1M_2N_1)(1L_2M_2N_1^2)}\hfill \\ & +\frac{L_1^2M_2N_2}{(1L_1M_1N_2)(1L_1^2M_2N_2)}+\frac{L_1M_1N_2}{(1L_1M_1N_2)(1L_2M_1^2N_2)})\hfill \end{array}$$ $`(2.35)`$ From this expression, we read off the following list of elementary couplings (recall that the first variable is a model variable and then we write the corresponding monomial in the grading variables): $$\begin{array}{cc}\hfill A_1:M_1N_1,& A_2:L_1N_1,A_3:L_1M_1\hfill \\ \hfill B_1:M_2N_2,& B_2:L_2N_2,B_3:L_2M_2\hfill \\ \hfill C_1:L_2M_1N_1,& C_2:L_1M_2N_1,C_3:L_1M_1N_2\hfill \\ \hfill D_1:L_1^2M_2N_2,& D_2:L_2M_1^2N_2,D_3:L_2M_2N_1^2.\hfill \end{array}$$ $`(2.36)`$ However, not all the products of the model variables can be linearly independent: there are linear relations between: \| \| $`C_iC_j,`$ \| $`A_kD_k,`$ \| $`\mathrm{and}A_iA_jB_k`$ \| \| --- \| --- \| --- \| --- \| \| \| $`D_iD_j,`$ \| $`A_k^2B_iB_j,`$ \| $`\mathrm{and}B_kC_k^2`$ \| \| \| $`C_iD_i,`$ \| $`A_jB_kC_k,`$ \| $`\mathrm{and}A_kB_jC_j`$ \| $`(2.37)`$ for $`i,j,k`$ a cyclic permutation of $`1,2,3`$ and repeated indices are not summed. (It is plain that the three sets of products found to be linearly related must have the same Dynkin labels.) A specific form of the generating function, as expressed in terms of the elementary couplings, amounts to a specific choice of a set of forbidden couplings among those that are related by a linear relation. 3. Tensor-product descriptions 3.1. The need for a tensor-product description It is clear that one major technical complication of the character method is that it starts at too fundamental a level, namely the character of the separate representations. One natural way to proceed is to start from a combinatorial description of the tensor-product rules. Such a description already takes into account the action of the Weyl group and encodes the various subtractions of the singular vectors. But how do we make the connection with the generating-function approach? The key is to find a combinatorial description which can be expressed as a set of linear Diophantine inequalities. Given this set of inequalities, there is an algorithm, again due to MacMahon, for constructing a generating function. (This is an adaptation of a method developed by Elliot for the analysis of linear Diophantine equalities and for this reason the algorithm is often referred to as the Elliot-MacMahon method. For a detailed discussion of the algorithm, see in particular vol. 2 section VIII of .) This method is conceptually similar to the character method, except that the starting point is substantially closer to the end result. See section 7.3 for a slight generalisation of this algorithm. Although the description of tensor products via linear Diophantine equations is a more efficient route to finding the generating function than the character one, complications associated to the $`\mathrm{\Omega }`$ projections remain a source of technical difficulty that severely limits the practical applicability of the method. A more powerful approach to our problem is to use the techniques of computational algebra. We start with a description of the tensor-product multiplicities as solutions to linear Diophantine inequalities. Efficient algorithms exists for finding the fundamental solutions to these inequalities . From these we find directly a model for the generating function using Grobner basis techniques. (This is roughly the inverse of MacMahon’s method which was originally conceived as a technique to generate the elementary couplings and their linear relations through the construction of the generating function. Here, the elementary couplings and their relations are first obtained and used as the input for the construction of the generating function.) 4. The LR rule ($`su(N)`$) For $`su(N)`$ tensor products there is a particularly convenient description based on Littlewood-Richardson tableaux supplemented by the stretched-product operation (defined below) . Integrable weights in $`su(N)`$ can be represented by tableaux: the weight $`(\lambda _1,\lambda _2,\mathrm{},\lambda _{N1})`$ is associated to a left justified tableau of $`N1`$ rows with $`\lambda _1+\lambda _2+\mathrm{}+\lambda _{N1}`$ boxes in the first row, $`\lambda _2+\lambda _3+\mathrm{}+\lambda _{N1}`$ boxes in the second row, etc. Equivalently, the tableau has $`\lambda _1`$ columns of 1 box, $`\lambda _2`$ columns of 2 boxes, etc. The scalar representation has no boxes, or equivalently, any number of columns of $`N`$ boxes. The Littlewood-Richardson rule is a simple combinatorial description of the tensor product of two $`su(N)`$ representations $`\lambda \mu `$. The second tableau ($`\mu `$) is filled with numbers as follows: the first row with $`1`$’s, the second row with $`2`$’s, etc. All the boxes with a $`1`$ are then added to the first tableau according to following restrictions: 1) the resulting tableau must be regular: the number of boxes in a given row must be smaller or equal to the number of boxes in the row immediately above; 2) the resulting tableau must not contain two boxes marked by $`1`$ in the same column. All the boxes marked by a $`2`$ are the added to the resulting tableaux according to the above two rules (with $`1`$ is replaced by $`2`$) and the further restriction: 3) in counting from right to left and top to bottom, the number of $`1`$’s must always be greater or equal to the number of $`2`$’s. The process is repeated with the boxes marked by a $`3,4,\mathrm{},N1`$, with the additional rule that the number of $`i`$’s must always be greater or equal to the number of $`i+1`$’s when counted from right to left and top to bottom. The resulting Littlewood-Richardson (LR) tableaux are the Young tableaux of the irreducible representations occurring in the decomposition. These rules can be rephrased in an algebraic way as follows . Define $`n_{ij}`$ to be the number of boxes $`i`$ that appear in the LR tableau in the row $`j`$. The LR conditions read $$\lambda _{j1}+\underset{i=1}{\overset{k1}{}}n_{i,j1}\underset{i=1}{\overset{k}{}}n_{ij}01k<jN$$ $`(4.1)`$ and $$\underset{j=i}{\overset{k}{}}n_{i1j1}\underset{j=i}{\overset{k}{}}n_{ij}02ikN\mathrm{and}iN1.$$ $`(4.2)`$ The weight $`\mu `$ of the second tableau and the weight $`\nu `$ of the resulting LR tableau are respectively given by $$\begin{array}{cc}\hfill \underset{j=i}{\overset{N}{}}n_{ij}& =\underset{j=i}{\overset{N1}{}}\mu _ji=1,2,\mathrm{},N1,\hfill \\ \hfill \nu _j\lambda _j+\underset{i=1}{\overset{N1}{}}n_{ij+1}& =\underset{i=1}{\overset{\mathrm{min}(j,N1)}{}}n_{ij}j=1,2,\mathrm{},N1.\hfill \end{array}$$ $`(4.3)`$ Hence, given three weights $`\lambda ,\mu `$ and $`\nu `$, the number of non-negative integers solutions $`\{n_{ij}\}`$ satisfying the above conditions gives the multiplicity $`𝒩_{\lambda \mu }^\nu `$ of $`\nu `$ in the tensor product $`\lambda \mu `$. The combined equations (4.1) and (4.2) constitute a set of linear and homogeneous inequalities. As described in , the Hilbert basis theorem guarantees that every solution can be expanded in terms of the elementary solutions of these inequalities. We can construct a model for the solutions of the equations (4.1) and (4.2) by introducing new formal variables $`A_i`$, $`1it`$ where $`t`$ is the total number of variables in (4.1) and (4.2). Then the subring of $`Q[A_i;1it]`$ generated by the monomials $`A^\alpha `$ with $`\alpha `$ a solution of (4.1) and (4.2) provides the required model. This ring $`R`$ will be generated by a finite set of monomials $`E_j`$ $`1js`$ which we call elementary couplings corresponding to the elementary solutions of (4.1) and (4.2). Thus $`R`$ is isomorphic to $`Q[e_1,\mathrm{},e_s]/I`$ under the mapping $`\varphi :e_iE_i`$ where $`I`$ is some ideal. Each element of $`I`$ corresponds, via the map $`\varphi `$, to a relation between the elementary couplings. In the case of LR tableaux, there is a nice pictorial representation of the model $`R`$. Consider the set of formal linear combinations of LR tableaux with rational coefficients. It is given a ring structure by defining the stretched product of two LR tableaux (denoted by $``$) to be the tableau obtained by fusing the two tableaux and reordering the numbers in each row in increasing order . More algebraically, if we denote the empty boxes of a LR tableau by a 0, so that $$n_{0j}=\underset{i=j}{\overset{N1}{}}\lambda _ij=1,2,\mathrm{},N1$$ $`(4.4)`$ we can characterise completely a tableau by the data $`\{n_{ij}\}`$ with now $`i0`$. It is clear the set of numbers $`\{n_{ij}\}`$ with $`i0`$, or equivalently, $`\{\lambda _i,n_{ij}\}`$ with $`i1`$, is a complete set of variables for the description of the tensor products. Then, the tableau obtained by the stretched product of the tableaux $`\{n_{ij}\}`$ and $`\{n_{ij}^{}\}`$ is simply described by the numbers $`\{n_{ij}+n_{ij}^{}\}`$. Here is a simple example: $$\begin{array}{c}\begin{array}{c}\text{ 1 1 1 1 1 2 2 3 4 }\end{array}\end{array}\begin{array}{c}\begin{array}{c}\text{ 1 1 1 1 1 2 1 2 }\end{array}\end{array}=\begin{array}{c}\begin{array}{c}\text{ 1 1 1 1 1 1 1 1 1 1 2 2 1 2 2 3 4 }\end{array}\end{array}$$ $`(4.5)`$ This ring of tableaux is isomorphic to the model $`R`$ constructed above and we do not distinguish between them. Thus we specify a set of elementary couplings (i.e. a set of generators of $`R`$) as a set of elementary LR Tableaux. 4.1. Example: the $`su(2)`$ case The complete set of inequalities for $`su(2)`$ variables $`\{\lambda _1,n_{11},n_{12}\}`$ is simply $$\lambda _1n_{12}n_{11}0n_{12}0$$ $`(4.6)`$ The other weights are fixed by $$\mu _1=n_{11}+n_{12}\nu _1=\lambda _1+n_{11}n_{12}$$ $`(4.7)`$ By inspection, the elementary solutions of this set of inequalities are $$(\lambda _1,n_{11},n_{12})=(1,0,1),(1,0,0),(0,1,0)$$ $`(4.8)`$ which correspond respectively to $`E_1,E_2,E_3`$ in (2.22). These correspond to the following LR tableaux: $$E_1:\begin{array}{c}\text{ 1 1 }\end{array},E_2:\begin{array}{c}\text{ 1 }\end{array},E_3:\begin{array}{c}\text{ 1 }\end{array}$$ $`(4.9)`$ It is also manifest that there are no linear relations between these couplings. The generating function is thus simply: $$G^{su(2)}=\frac{1}{(1E_1)(1E_2)(1E_3)}$$ $`(4.10)`$ 4.2. Example: multiple tensor products in the $`su(2)`$ case Consider the problem of finding the multiplicity of the representation $`\zeta `$ in the triple product $`\lambda \mu \nu \zeta `$. As a first step, the LR rule applies as before: with $`n_{11}+n_{12}=\mu _1`$, we have $`\lambda _1n_{12}`$. After the first product, we re-apply the LR rule with now $`\lambda _1`$ replaced by $`\lambda _1+n_{11}n_{12}`$ and $`n_{ij}`$ replaced by $`m_{ij}`$ with $`m_{11}+m_{12}=\nu _1`$. The LR gives $`\lambda _1+n_{11}n_{12}m_{12}`$. The two inequalities for the $`su(2)`$ quadruple product are then: $$\lambda _1n_{12}\lambda _1+n_{11}n_{12}m_{12}n_{ij}0m_{ij}0$$ $`(4.11)`$ The elementary solutions are then, in the order: name of the coupling, corresponding Dynkin labels and the 5-vector $`(\lambda _1,n_{11},n_{12},m_{11},m_{12})`$,: \| \| $`E_1:`$ \| $`(1)(1)(0)(0)`$ \| $`(1,0,1,0,0)`$ \| \| --- \| --- \| --- \| --- \| \| \| $`E_2:`$ \| $`(1)(0)(1)(0)`$ \| $`(1,0,0,0,1)`$ \| \| \| $`E_3:`$ \| $`(1)(0)(0)(1)`$ \| $`(1,0,0,0,0)`$ \| \| \| $`E_4:`$ \| $`(0)(1)(1)(0)`$ \| $`(0,1,0,0,1)`$ \| \| \| $`E_5:`$ \| $`(0)(1)(0)(1)`$ \| $`(0,1,0,0,0)`$ \| \| \| $`E_6:`$ \| $`(0)(0)(1)(1)`$ \| $`(0,0,0,1,0)`$ \| $`(4.12)`$ The linear relation, whose existence was signalled by the character method, is $$E_3E_4=E_2E_5:(1,1,0,0,1),E_1E_6:(1,0,1,1,0)$$ $`(4.13)`$ Choosing to forbid the product $`E_3E_4`$, the generating function can be written in the form $$\begin{array}{cc}\hfill G& =\frac{1E_3E_4}{(1E_1)(1E_2)(1E_3)(1E_4)(1E_5)(1E_6)}\hfill \\ & =\left(\underset{i=1,2,5,6}{}\frac{1}{1E_i}\right)\left(\frac{1}{1E_3}+\frac{E_4}{1E_4}\right)\hfill \end{array}$$ $`(4.14)`$ The latter form makes manifest the absence of $`E_3E_4`$. We could represent the elementary couplings in terms of tableaux, where the boxes with 1’s refers to the $`\mu `$ tableau and those with 2’s originate from the $`\nu `$ tableau. (Warning: the resulting tableaux describing the four-products are not necessarily LR tableaux.) Hence, $`n_{1j}`$ gives the number of 1’s in row $`j`$ of the composed tableau while $`m_{1k}`$ gives the number of 2’s in row $`k`$. The elementary tableaux are \| \| $`E_1:\begin{array}{c}\text{ 1 1 }\end{array}`$ \| $`E_2:\begin{array}{c}\text{ 1 2 }\end{array},`$ \| $`E_3:\begin{array}{c}\text{ }\text{ }\text{1}\text{ }\text{ }\end{array}`$ \| \| --- \| --- \| --- \| --- \| \| \| $`E_4:\begin{array}{c}\text{ 1 2 }\end{array}`$ \| $`E_5:\begin{array}{c}\text{ }\text{ }\text{1}\text{ }\text{ }\end{array},`$ \| $`E_6:\begin{array}{c}\text{ }\text{ }\text{2}\text{ }\text{ }\end{array}`$ \| $`(4.15)`$ From this representation, the relation reads $$E_3E_4=E_2E_5:\begin{array}{c}\text{ 1 1 2 }\end{array},E_1E_6:\begin{array}{c}\text{ 1 2 1 }\end{array}$$ $`(4.16)`$ 4.3. Example: the $`su(4)`$ case The $`su(4)`$ LR conditions are: \| \| $`\lambda _1n_{12}`$ \| $`n_{11}n_{22}`$ \| \| --- \| --- \| --- \| \| \| $`\lambda _2n_{13}`$ \| $`n_{11}+n_{12}n_{22}+n_{23}`$ \| \| \| $`\lambda _2+n_{12}n_{13}+n_{23}`$ \| $`n_{11}+n_{12}+n_{13}n_{22}+n_{23}+n_{24}`$ \| \| \| $`\lambda _3n_{14}`$ \| $`n_{22}n_{33}`$ \| \| \| $`\lambda _3+n_{13}n_{14}+n_{24}`$ \| $`n_{22}+n_{23}n_{33}+n_{34}`$ \| \| \| $`\lambda _3+n_{13}+n_{23}n_{14}+n_{24}+n_{34}`$ \| \| $`(4.17)`$ The tensor-product elementary couplings are: $$\begin{array}{cc}\hfill A_1& :\begin{array}{c}\text{ 1 2 3 }\end{array},A_2:\begin{array}{c}\text{ 1 1 1 1 }\end{array},A_3:\begin{array}{c}\text{ 1 }\end{array},B_1:\begin{array}{c}\text{ 1 2 }\end{array},B_2:\begin{array}{c}\text{ 1 1 1 2 }\end{array},B_3:\begin{array}{c}\text{ 1 1 }\end{array},\hfill \\ \hfill C_1& :\begin{array}{c}\text{ 1 }\end{array},C_2:\begin{array}{c}\text{ 1 1 2 3 }\end{array},C_3:\begin{array}{c}\text{ 1 1 1 }\end{array},D_1^{}:\begin{array}{c}\text{ 1 1 1 }\end{array},D_2^{}:\begin{array}{c}\text{ 1 1 }\end{array},D_3^{}:\begin{array}{c}\text{ 1 1 2 }\end{array},\hfill \end{array}$$ $`(4.18)`$ together with $$D_1:\begin{array}{c}\text{ 1 1 1 2 3 }\end{array},D_2:\begin{array}{c}\text{ 1 1 1 2 1 3 }\end{array},D_3:\begin{array}{c}\text{ 1 1 1 1 2 }\end{array},E_1:\begin{array}{c}\text{ 1 1 1 1 1 2 }\end{array},E_2:\begin{array}{c}\text{ 1 1 1 2 }\end{array},E_3:\begin{array}{c}\text{ 1 1 1 2 1 3 }\end{array}$$ $`(4.19)`$ The Dynkin-label transcription of the elementary couplings reads \| \| $`A_1:(0,0,0)(0,0,1)(0,0,1)`$ \| $`D_1^{}:(0,1,0)(1,0,0)(0,0,1)`$ \| \| --- \| --- \| --- \| \| \| $`A_2:(0,0,1)(1,0,0)(0,0,0)`$ \| $`D_2^{}:(1,0,0)(1,0,0)(0,1,0)`$ \| \| \| $`A_3:(1,0,0)(0,0,0)(1,0,0)`$ \| $`D_3^{}:(1,0,0)(0,1,0)(0,0,1)`$ \| \| \| $`B_1:(0,0,0)(0,1,0)(0,1,0)`$ \| $`D_1:(0,1,0)(0,0,1)(1,0,0)`$ \| \| \| $`B_2:(0,1,0)(0,1,0)(0,0,0)`$ \| $`D_2:(0,0,1)(0,0,1)(0,1,0)`$ \| \| \| $`B_3:(0,1,0)(0,0,0)(0,1,0)`$ \| $`D_3:(0,0,1)(0,1,0)(1,0,0)`$ \| \| \| $`C_1:(0,0,0)(1,0,0)(1,0,0)`$ \| $`E_1:(1,0,1)(0,1,0)(0,1,0)`$ \| \| \| $`C_2:(1,0,0)(0,0,1)(0,0,0)`$ \| $`E_2:(0,1,0)(0,1,0)(1,0,1)`$ \| \| \| $`C_3:(0,0,1)(0,0,0)(0,0,1)`$ \| $`E_3:(0,1,0)(1,0,1)(0,1,0)`$ \| $`(4.20)`$ For $`su(4)`$ there are 15 relations : \| \| $`D_j^{^{}}D_k=C_iE_i`$ \| $`D_jD_k^{^{}}=B_iC_jC_k`$ \| $`E_iE_j=B_kD_kD_k^{^{}}`$ \| \| --- \| --- \| --- \| --- \| \| \| $`D_iE_i=C_jB_kD_k`$ \| $`D_i^{^{}}E_i=B_jD_j^{^{}}C_k`$ \| \| $`(4.21)`$ with $`i,j,k`$ a cyclic permutation of $`1,2,3`$. To construct the generating function, we need to select forbidden couplings. It turns out that when there are more that one relation, complications may arise. We must ensure that the selected forbidden couplings are complete, which means that no further (usually higher-order) relations are required for a unique decomposition of a given coupling. A technique that is tailor-made for dealing with problems of that type is that of Grobner bases. This will be introduced in the next section. At this point, we simply indicate a complete choice of forbidden couplings, namely $`\{E_iE_j,D_i^{}E_i,D_iE_i,D_jD_i^{},D_j^{}D_i\}`$. This yields then a model for the generating function, which then reads : $$\begin{array}{cc}\hfill G^{su(4)}=& (\underset{i=1}{\overset{3}{}}\stackrel{~}{A}_i\stackrel{~}{B}_i\stackrel{~}{C}_i)(\stackrel{~}{D}_1^{}\stackrel{~}{D}_2^{}\stackrel{~}{D}_3^{}+E_1\stackrel{~}{E_1}\stackrel{~}{D}_2^{}\stackrel{~}{D}_3^{}+D_3\stackrel{~}{D}_3\stackrel{~}{D}_3^{}\stackrel{~}{E}_1\hfill \\ & +D_2\stackrel{~}{D}_2\stackrel{~}{D}_3\stackrel{~}{E}_1+D_1\stackrel{~}{D}_1\stackrel{~}{D}_2\stackrel{~}{D}_3+E_3\stackrel{~}{E}_3\stackrel{~}{D}_1\stackrel{~}{D}_2+D_1^{}\stackrel{~}{D}_1^{}\stackrel{~}{D}_1\stackrel{~}{E}_3\hfill \\ & +D_2^{}E_3\stackrel{~}{D}_2^{}\stackrel{~}{E}_3\stackrel{~}{D}_1^{}+E_2\stackrel{~}{E}_2\stackrel{~}{D}_1^{}\stackrel{~}{D}_3^{}+E_2D_1\stackrel{~}{E}_2\stackrel{~}{D}_1\stackrel{~}{D}_1^{}+E_2D_3\stackrel{~}{E}_2\stackrel{~}{D}_3\stackrel{~}{D}_3^{}\hfill \\ & +D_1D_3E_2\stackrel{~}{D}_1\stackrel{~}{D}_3\stackrel{~}{E}_2+D_2D_2^{}\stackrel{~}{D}_2\stackrel{~}{D}_2^{}\stackrel{~}{E}_1+D_2D_2^{}E_3\stackrel{~}{D}_2\stackrel{~}{D}_2^{}\stackrel{~}{E}_3).\hfill \end{array}$$ $`(4.22)`$ where $$\stackrel{~}{M}_i=(1M_i)^1.$$ $`(4.23)`$ 5. Diophantine inequalities: elementary couplings, relations and Grobner bases We introduce the idea of the Grobner basis via a simple example (see also ). Suppose $`R`$ is a model for a generating function, where $`R=Q[x,y,z,t]/I`$ and $`I=(xyt,zyt)`$ is the ideal generated by $`xyt`$ and $`zyt`$, with an $`N^2`$ grading given by $`(1,0),(0,1),(1,0)`$ and $`(1,1)`$ for $`x,y,z`$ and $`t`$. Writing $`\overline{x}=x+I`$ and similarly for the other variables, we have in $`R`$ that $`\overline{x}\overline{y}=\overline{t}`$ and $`\overline{z}\overline{y}=\overline{t}`$. These two expressions give two re-write rules : $`xyt`$ and $`zyt`$. These rules can be used to simplify any monomial. The aim is to find a re-write rule which, when iterated, produces unique representatives for the classes of $`I`$. If this is the case, then a vector space basis of $`R`$ would consist of terms of the form $`m+I`$ with $`m`$ a monomial which is not divisible by any of the left-hand sides of the rewrite rules. In the example above, if we had ‘good’ rewrite rules then a basis for $`R`$ would be represented by monomials not containing $`xy`$ or $`zy`$, i.e. monomials of the form either $`y^at^b`$ or $`x^az^bt^c`$. The generating function which counts these monomials is: $$\frac{1}{(1AB)}\left(\frac{B}{1B}+\frac{1}{(1A)^2}\right),$$ $`(5.1)`$ The exponent of $`A`$ carries the first grading index and $`B`$ the second. However this generating function is not correct. It contains the term $`2A^2B`$ corresponding to the monomials $`xt`$ and $`zt`$. But the polynomial $`z(xyt)x(zyt)=xtzt`$ is also in $`I`$ and hence in $`R`$ we have $`\overline{x}\overline{t}=\overline{z}\overline{t}`$ and so the space of grade $`(2,1)`$ has dimension 1 rather than 2. This problem can also be seen as a problem with the re-write rules. If we start with $`xyz`$ then we can use the first re-write rule: $`xyztz`$ or the second: $`xyzxt`$. We cannot apply any further re-write rules and so this set of re-write rules does not produce a unique representative. The solution is to include the rule $`xtzt`$. This gives a set of 3 rules: $`xyt`$, $`zyt`$ and $`xtzt`$. It turns out that this is a ‘good’ set and so a basis for $`R`$ is given by (the classes of) monomials of the form $`y^at^b`$, $`x^az^b`$ and $`z^at^b`$ which gives the generating function: $$\frac{1}{(1AB)(1B)}+\frac{A}{(1A)^2}+\frac{A}{(1A)(1AB)}$$ $`(5.2)`$ The set of ‘good’ generators, $`xyt,zyt,xtzt`$ we have found for $`I`$ is known as a Grobner basis . The general procedure for constructing a Grobner basis given a set of generating polynomials is as follows. First choose a term ordering, which is an ordering on monomials with the property that any chain $`m_1>m_2>\mathrm{}`$ has finite length. For example we can order the variables by $`x>y>z>t`$ and then order all monomials by the corresponding lexicographic (dictionary) order, for example: $`x^2y>xyz>y^3`$. For each generator of our ideal $`I`$, select the monomial which is highest with respect to the given term ordering. This is then the term which appears on the left of the re-write rule. The lexicographic ordering gives the first two re-write rules of our example: $`xyt`$ and $`zyt`$. Next, for each pair of leading terms find the lowest common multiple and simplify it in the two possible ways. In this case there is only one pair of leading terms and the lowest common multiple is $`xyz`$ which simplifies to $`xt`$ and $`yt`$. Continue to apply the re-write rules until the terms do not simplify further. If the resulting pair of terms are the same, then proceed to the next pair of leading terms, otherwise add a new re-write rule. In this case we add $`xtyt`$. Proceed until no pair of leading terms gives a new rule. This is the case for the rules we now have. For example the two rules $`xyt`$ and $`xtzt`$ appears to give a new rule by simplifying $`xyt`$ to both $`t^2`$ and $`yzt`$. However the second term can be further reduced to $`t^2`$ and so no new rule is required. This algorithm for computing Grobner bases is known as Buchberger’s algorithm. Improvements on this basic algorithm mean that it is now feasible to find Grobner bases for quite large sets of generating polynomials. (The web pages of the computer-algebra information network at the address http://cand.can.nl/CAIN contain information about many of the programs currently available.) Although it is not clear from this example, Grobner bases are a very versatile tool for performing explicit calculations. We end this section with an illustrative example relevant to our discussion of tensor-product generating functions. Consider a set of linear Diophantine equations: $$M\alpha =0,\alpha 0$$ $`(5.3)`$ with $`M`$ an integer matrix and $`\alpha `$ a vector of non-negative integers. We would like to construct a generating function for the solutions to this set of equations: $$\underset{\alpha }{}x^\alpha .$$ $`(5.4)`$ A non-trivial example is given by the Diophantine equations that describe a $`3\times 3`$ magic square: $$\left(\begin{array}{ccc}a& b& c\\ d& e& f\\ g& h& i\end{array}\right)$$ $`(5.5)`$ with non-negative entries and equal row and column sums. The magic square condition (the sum of each row and each column is the same, say equal to $`t`$) gives the following set of equations: \| \| $`a+b+c=t`$ \| $`a+d+g=t`$ \| \| --- \| --- \| --- \| \| \| $`d+e+f=t`$ \| $`b+e+h=t`$ \| \| \| $`g+h+i=t`$ \| $`c+f+i=t`$ \| $`(5.6)`$ With $`\alpha `$ standing for the column vector with entries $`(a,b,c,d,e,f,g,h,i,t)`$, the matrix $`M`$ reads $$M=\left(\begin{array}{cccccccccc}1& 1& 1& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 1& 1& 1& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 1& 1& 1\\ 1& 0& 0& 1& 0& 0& 1& 0& 0& 1\\ 0& 1& 0& 0& 1& 0& 0& 1& 0& 1\\ 0& 0& 1& 0& 0& 1& 0& 0& 1& 1\end{array}\right)$$ $`(5.7)`$ There is a straightforward algorithm for finding the basic set of solutions which yields: \| \| $`\alpha _1=(0,0,1,0,1,0,1,0,0,1)`$ \| $`\alpha _4=(1,0,0,0,0,1,0,1,0,1)`$ \| \| --- \| --- \| --- \| \| \| $`\alpha _2=(0,1,0,0,0,1,1,0,0,1)`$ \| $`\alpha _5=(0,1,0,1,0,0,0,0,1,1)`$ \| \| \| $`\alpha _3=(0,0,1,1,0,0,0,1,0,1)`$ \| $`\alpha _6=(1,0,0,0,1,0,0,0,1,1)`$ \| $`(5.8)`$ We shall use $`A,B,\mathrm{},T`$ to denote the “grading variables” of this example so that the exponent of $`A`$ carries the value of $`a`$ and so on. A model for the generating function is given by the subring $`S`$ of $`Q[A,B,C,D,E,F,G,H,I,T]`$ generated by monomials corresponding to the 6 elementary solutions, \| $`E_1`$ \| $`=CEGT,E_2`$ \| $`=BFGT,E_3`$ \| $`=CDHT,`$ \| \| --- \| --- \| --- \| --- \| \| $`E_4`$ \| $`=AFHT,E_5`$ \| $`=BDIT,E_6`$ \| $`=AEIT`$ \| $`(5.9)`$ The monomials in $`S`$ correspond to magic squares. For example $`E_1^2E_4E_6=A^2C^2E^3FG^2HIT^4S`$ corresponds to a square with row and column sums equal to 4: $$\left(\begin{array}{ccc}2& 0& 2\\ 0& 3& 1\\ 2& 1& 1\end{array}\right).$$ $`(5.10)`$ Note that in this example it is convenient to construct our model as a subring of the ring of grading variables. Thus each “elementary coupling” $`E_i`$ is actually equal to the corresponding monomial in the grading variables. However, there are relations between these generators and so it is not immediately clear how to construct the Poincaré series for $`S`$. What we require is an isomorphism of $`S`$ with $`R=Q[e_1,\mathrm{},e_6]/I`$ such that $`e_iE_i`$, $`i=1,\mathrm{},6`$ and such that we have a Grobner basis of the ideal $`I`$ (the ‘ideal of relations’). Fortunately, such an isomorphism is easily constructed using Grobner-basis methods. Introduce the ring $`Q[A,B,C,D,E,F,G,H,I,T,e_1,\mathrm{},e_6]`$ with the lexicographic ordering $$A>B>C>D>E>F>G>H>I>T>e_1>\mathrm{}>e_6$$ $`(5.11)`$ Let $`J`$ be the ideal generated by $`E_1e_1,\mathrm{},E_6e_6`$. This is not necessarily a Grobner basis with respect to this term ordering. Let $`G`$ be the Grobner basis for $`J`$ with the given ordering. Then it can be shown that $`GQ[e_1,\mathrm{},e_6]`$ is a Grobner basis for the ideal of relations $`I`$ which we require. In this case $`G`$ is quite large, but its intersection with $`Q[e_1,\mathrm{},e_6]`$ is $`e_1e_4e_5e_2e_3e_6`$. The corresponding relation in $`R`$ is $`E_1E_4E_5E_2E_3E_6`$ and these two terms do indeed give the same magic square, so that indeed we have found a relation between the generators of $`R`$. The Poincaré series for $`Q[e_1,\mathrm{},e_6]/I`$ is easily computed: $$\begin{array}{cc}\hfill \frac{1}{(1E_2)(1E_3)(1E_6)}& (\frac{1}{(1E_1)(1E_4)}\hfill \\ & +\frac{E_5}{(1E_1)(1E_5)}+\frac{E_4E_5}{(1E_4)(1E_5)})\hfill \end{array}$$ $`(5.12)`$ 6. Berenstein-Zelevinsky Triangles 6.1. Generalities The previous examples make clear the usefulness of a re-expression of the tensor-product calculation in terms of Diophantine inequalities. The Littlewood-Richardson algorithm yields a set of such inequalities only for $`su(N)`$. Fortunately, Berenstein and Zelevinsky have expressed the solution of the multiplicity of a given tensor product as a counting problem for the number of integral points in a convex polytope. For a given algebra, the polytope is formulated in terms of a characteristic set of inequalities. For $`su(N)`$, these reduce to the LR set of inequalities. For the other classical algebras, except $`sp(4)`$, the proposed set of inequalities is a conjecture. 6.2. BZ triangles for $`sp(4)`$ The combinatorial description of tensor products for $`sp(4)`$ is not as simple as in the $`su(N)`$ case: a standard LR product must be supplemented by a division operation and modification rules . Given the BZ set of inequalities, the natural way to proceed, as just mentioned, is to interpret these as the appropriate inequalities for the description of the tensor products. These inequalities are as follows: \| \| $`\lambda _1p`$ \| $`\mu _1q`$ \| \| --- \| --- \| --- \| \| \| $`\lambda _2r_1/2`$ \| $`\mu _1q+r_1r_2`$ \| \| \| $`\lambda _2r_1/2+qp`$ \| $`\mu _1p+r_1r_2`$ \| \| \| $`\lambda _2r_2/2+qp`$ \| $`\mu _2r_2/2`$ \| \| \| $`\nu _1=r_2r_12p+\lambda _1+\mu _1`$ \| $`\nu _2=pqr_2+\lambda _2+\mu _2`$ \| $`(6.1)`$ (Our notation is different from that used in : the relation is $`r_1=m_1,r_2=m_2,p=m_{12},q=m_{12}^{}`$.) The $`sp(4)`$ tensor product coefficient $`𝒩_{\lambda \mu \nu }`$ is thus given by the number of solutions of the above system with $`r_1,r_22N`$ et $`p,qN`$ ($`N`$ being the set of nonnegative integers). A proper set of variables for a complete description of a particular tensor-product coupling is thus $`\{\lambda _1,\lambda _2,\mu _1,\mu _2,r_1,r_2,p,q\}`$. We give the list of elementary couplings, adding to each coupling the corresponding four-vector $`[r_1,r_2,p,q]`$: \| \| $`A_1:`$ \| $`(0,0)(1,0)(1,0)[0,0,0,0]`$ \| \| --- \| --- \| --- \| \| \| $`A_2:`$ \| $`(1,0)(0,0)(1,0)[0,0,0,0]`$ \| \| \| $`A_3:`$ \| $`(1,0)(1,0)(0,0)[0,0,1,1]`$ \| \| \| $`B_1:`$ \| $`(0,0)(0,1)(0,1)[0,0,0,0]`$ \| \| \| $`B_2:`$ \| $`(0,1)(0,0)(0,1)[0,0,0,0]`$ \| \| \| $`B_3:`$ \| $`(0,1)(0,1)(0,0)[2,2,0,0]`$ \| \| \| $`C_1:`$ \| $`(0,1)(1,0)(1,0)[0,0,0,1]`$ \| \| \| $`C_2:`$ \| $`(1,0)(0,1)(1,0)[0,2,1,0]`$ \| \| \| $`C_3:`$ \| $`(1,0)(1,0)(0,1)[0,0,1,0]`$ \| \| \| $`D_1:`$ \| $`(2,0)(0,1)(0,1)[0,2,2,0]`$ \| \| \| $`D_2:`$ \| $`(0,1)(2,0)(0,1)[2,0,0,0]`$ \| \| \| $`D_3:`$ \| $`(0,1)(0,1)(2,0)[0,2,0,0]`$ \| $`(6.2)`$ The unspecified linear relations mentioned in (2.37) can now be obtained. To find those products that are equal in the current situation we need only compare their corresponding sets of four-vectors $`[r_1,r_2,p,q]`$ (which are additive in products of couplings). We thus find for instance that $$C_1C_2=A_3D_3:[0,2,1,1]A_1A_2B_3:[2,2,0,0]$$ $`(6.3)`$ Proceeding in this way for the other cases, we find the following complete list of relations: \| \| $`C_1C_2=A_3D_3,`$ \| $`C_2C_3=A_1D_1`$ \| $`C_3C_1=A_1A_3B_2`$ \| \| --- \| --- \| --- \| --- \| \| \| $`D_1D_2=B_3C_3^2`$ \| $`D_2D_3=A_1^2B_2B_3`$ \| $`D_1D_3=B_2C_2^2`$ \| \| \| $`C_1D_1=A_3B_2C_2`$ \| $`C_2D_2=A_1B_3C_3`$ \| $`C_3D_3=A_1B_2C_2`$ \| $`(6.4)`$ The use of the BZ inequalities to find the elementary couplings and their relations is novel. (An off-shoot of our construction is that it provides an indirect proof of the validity of the BZ inequalities since we recover from it the result of derived from the character method.) A possible choice of forbidden products is the one given in : $$\{C_iC_j,D_iD_j,C_iD_i\}$$ $`(6.5)`$ with $`i,j=1,2,3\mathrm{and}ij`$. It leads to the generating function: $$\begin{array}{cc}\hfill G^{sp(4)}=\left(\underset{i=1}{\overset{3}{}}\stackrel{~}{A}_i\stackrel{~}{B}_i\right)& (\stackrel{~}{C}_1\stackrel{~}{D}_2+D_3\stackrel{~}{C}_1\stackrel{~}{D}_3+C_2D_1\stackrel{~}{C}_2\stackrel{~}{D}_1\hfill \\ & +C_2\stackrel{~}{C}_2\stackrel{~}{D}_3+D_1\stackrel{~}{C}_3\stackrel{~}{D}_1+C_3\stackrel{~}{C}_3\stackrel{~}{D}_2)\hfill \end{array}$$ $`(6.6)`$ Of course, by modifying the ordering in the Grobner basis, we can get other choices of forbidden couplings. Here is another set of forbidden couplings that can be obtained: $`\{D_iD_j,C_iD_i,A_1D_1,A_3D_3,A_1A_3B_2\}`$. The corresponding generating function reads $$\begin{array}{cc}\hfill G^{sp(4)}& =\stackrel{~}{B}_1\stackrel{~}{B}_2\stackrel{~}{B}_3[\left(\underset{i=1}{\overset{3}{}}\stackrel{~}{A}_i\right)\stackrel{~}{C}_i(1A_1A_3B_2)+D_3\stackrel{~}{D}_3\stackrel{~}{A}_1\stackrel{~}{A}_2\stackrel{~}{C}_1\stackrel{~}{C}_2\hfill \\ & +D_1\stackrel{~}{D}_1\stackrel{~}{A}_2\stackrel{~}{A}_3\stackrel{~}{C}_2\stackrel{~}{C}_3+D_2\stackrel{~}{D}_2\stackrel{~}{A}_1\stackrel{~}{A}_2\stackrel{~}{A}_3\stackrel{~}{C}_1\stackrel{~}{C}_3(1A_1A_3B_2)].\hfill \end{array}$$ $`(6.7)`$ These two generating functions are equivalent when rewritten in terms of the grading variables, that is, in terms of Dynkin labels. However, they originate from two distinct models. The second one turns out to be well adapted to the fusion extension. 7. A vector basis approach to the construction of generating functions In this section, we present a simple and systematic way of generating by hand all the elementary solutions of a set of linear homogeneous inequalities starting from the well-known construction of a vector basis. The first step amounts to reformulate the system of inequalities in terms of equalities. We then look for the elementary independent solutions by relaxing the positivity requirement. In other words, we construct the vector basis. In a final step, we find the minimal linear combinations of these vector basis elements that yield positive solutions. This will also provide an illustration of MacMahon’s projection technique. The result of this projection is the desired tensor-product generating function. Hence, this approach turns out to be a new way of constructing the tensor-product generating functions. (This generic method, referred to as being novel for tensor products, is certainly well-known in general: it is discussed in the first reference of .) 7.1. Graphical representations as BZ triangles for $`su(N)`$ Consider the direct transformation of the LR inequalities to equalities by introducing an appropriate number of new non-negative integer variables. Consider first the $`su(2)`$ case, for which there is a single inequality: $`\lambda _1n_{12}`$. We transform this into an equality by introducing the non-negative integer $`a`$ defined by $$\lambda _1=n_{12}+a$$ $`(7.1)`$ The expression for $`\nu _1`$ becomes then $`\nu _1=\lambda _1+n_{11}n_{12}=a+n_{11}`$. Since $`\mu _1=n_{11}+n_{12}`$, we are led naturally to a triangle representation of the tensor product: $$\lambda \mu \nu \begin{array}{c}a\\ n_{12}n_{11}\end{array}$$ $`(7.2)`$ We read off the Dynkin label of the $`\lambda `$ representation from the sum of the two integers that form the left side of the triangle, that of the $`\mu `$ representation from the bottom of the triangle and the $`\nu _1`$ label is the sum of the two integers that form the right side. A more uniform notation amounts to setting $`a=m_{12}`$ and $`n_{11}=l_{12}`$, in terms of which the triangle looks quite symmetrical: $$\begin{array}{c}m_{12}\\ n_{12}l_{12}\end{array}$$ $`(7.3)`$ with $$\lambda _1=m_{12}+n_{12}\mu _1=n_{12}+l_{12}\nu _1=m_{12}+l_{12}$$ $`(7.4)`$ These numbers $`m_{12}`$ and $`l_{12}`$ plays the role of $`n_{12}`$ in the permuted versions of the tensor product. The triangle combinatorial reformulation of the tensor product problem is as follows: the number of triangles that can be formed from non-negative integers $`n_{12},m_{12}`$ and $`l_{12}`$ that add up to the Dynkin labels of the representations under study according to the above relations gives the multiplicity of the triple coupling $`\lambda \mu \nu `$, or equivalently, the multiplicity of the scalar representation in the product $`\lambda \mu \nu (0)`$ (since for $`su(2)`$, $`\nu ^{}=\nu `$). The situation for $`su(3)`$ is somewhat more complicated. The transformation of the LR inequalities (4.1, 4.2) into equalities in this case takes the form \| \| $`\lambda _1=n_{12}+a`$ \| $`n_{11}=n_{22}+d`$ \| \| --- \| --- \| --- \| \| \| $`\lambda _2=n_{13}+b`$ \| $`n_{11}+n_{12}=n_{22}+n_{23}+e`$ \| \| \| $`\lambda _2+n_{12}=n_{13}+n_{23}+c`$ \| \| $`(7.5)`$ The expression for the other weights becomes \| \| $`\mu _1=n_{13}+e`$ \| $`\mu _2=n_{22}+n_{23}`$ \| \| --- \| --- \| --- \| \| \| $`\nu _1=a+d`$ \| $`\nu _2=n_{22}+c`$ \| $`(7.6)`$ Since there are two expressions for both $`n_{11}`$ and $`\lambda _2`$, there follows the compatibility relations: $$n_{12}+d=n_{23}+en_{23}+c=b+n_{12}$$ $`(7.7)`$ By adding these two relations, we find: $$c+d=b+e$$ $`(7.8)`$ Again we are led naturally to a triangle representation: with $`\zeta =\nu ^{}`$ this reads $$\begin{array}{c}a_{}\\ n_{12}d_{}\\ b_{}c_{}\\ n_{13}en_{23}n_{22}\end{array}$$ $`(7.9)`$ We read the Dynkin labels from the sides of the triangles, from $`\lambda _1`$ to $`\zeta _2`$ in an anti-clockwise rotation starting from the top of the triangle, exactly as for $`su(2)`$, except that here there are two labels on each sides. Notice that the compatibility conditions amounts to the equality of the sums of the extremal points of the three pairs of opposite sides of the hexagon obtained by dropping the three corners of the triangle. Again a more symmetrical notation is: $$a=m_{13}b=m_{23}c=m_{12}d=l_{23}e=l_{12}n_{22}=l_{13}$$ $`(7.10)`$ in terms of which the triangle reads $$\begin{array}{c}m_{13}\\ n_{12}l_{23}\\ m_{23}m_{12}\\ n_{13}l_{12}n_{23}l_{13}\end{array}$$ $`(7.11)`$ with labels fixed by: \| \| $`\lambda _1=m_{13}+n_{12}`$ \| $`\lambda _2=m_{23}+n_{13}`$ \| \| --- \| --- \| --- \| \| \| $`\mu _1=n_{13}+l_{12}`$ \| $`\mu _2=n_{23}+l_{13}`$ \| \| \| $`\zeta _1=l_{13}+m_{12}`$ \| $`\zeta _2=l_{23}+m_{13}`$ \| $`(7.12)`$ The hexagon conditions read: $$\begin{array}{cc}\hfill n_{12}+m_{23}& =n_{23}+m_{12},\hfill \\ \hfill l_{12}+m_{23}& =l_{23}+m_{12},\hfill \\ \hfill l_{12}+n_{23}& =l_{23}+n_{12}.\hfill \end{array}$$ $`(7.13)`$ In terms of triangles, the problem of finding the multiplicity of the $`su(3)`$ tensor product $`\lambda \mu \zeta 0`$ boils down to enumerating the number of triangles made with non-negative integers that form a bipartition of the Dynkin labels and that satisfy the above three hexagon relations. (For the $`su(N)`$ generalisation, see ). Here is the rationale for the labelling $`n_{ij},m_{ij},l_{ij}`$ from the triangle point of view . If $`e_i`$ are orthonormal vectors in $`𝐑^N,`$ then the positive roots of $`su(N)`$ can be represented in the form $`e_ie_j,1i<jN.`$ The triangle encodes three sums of positive roots: $$\begin{array}{cc}\hfill \mu +\zeta \lambda ^{}& =\underset{i<j}{}l_{ij}(e_ie_j),\hfill \\ \hfill \zeta +\lambda \mu ^{}& =\underset{i<j}{}m_{ij}(e_ie_j),\hfill \\ \hfill \lambda +\mu \zeta ^{}& =\underset{i<j}{}n_{ij}(e_ie_j),\hfill \end{array}$$ $`(7.14)`$ The hexagon relations are simply the consistency conditions for these three expansions. Clearly, the variables $`n_{ij}`$ that appear in the above relations are exactly the $`n_{ij}`$ that appear in the LR tableaux for the product $`\lambda \mu \zeta ^{}=\nu `$. 7.2. From a vector basis to the generating function: the $`su(3)`$ case Given the transcription of inequalities into equalities, we can easily extract the corresponding basis vectors. This is the starting point of a new method for constructing the tensor-product generating functions. To keep things concrete, we focus on the $`su(3)`$ case. The goal is to first get a vector basis and then to project it to get the elementary couplings. The generating function is a direct result of this procedure. The equality version of the LR inequalities are (7.12) and (7.13); they underlie the construction of the BZ triangle (7.11). The last hexagon condition of (7.13) is the difference of the previous two so it is not an independent relations. We thus have a total of 15 variables: $`\lambda _1,\mathrm{},\zeta _2,l_{12},\mathrm{},n_{23}`$ and 8 equations. The number of independent variables is thus 7. These will be chosen to be $`m_{13},m_{23},l_{13},l_{23},n_{12},n_{13},n_{23}`$. The dependent variables are fixed as follows: \| \| $`\lambda _1=m_{13}+n_{12}`$ \| $`\lambda _2=m_{23}+n_{13}`$ \| \| --- \| --- \| --- \| \| \| $`\mu _1=n_{13}+n_{12}+l_{23}n_{23}`$ \| $`\mu _2=n_{23}+l_{13}`$ \| \| \| $`\zeta _1=n_{12}+m_{23}+l_{13}n_{23}`$ \| $`\zeta _2=l_{23}+m_{13}`$ \| \| \| $`l_{12}=n_{12}+l_{23}n_{23}`$ \| $`m_{12}=n_{12}+m_{23}n_{23}`$ \| $`(7.15)`$ We now look for the elementary solutions of this system (without invoking the constraint that all the above dependent variables should be necessarily positive). The sought basis vectors are obtained by setting one of the variable $`m_{13},\mathrm{},n_{23}`$ to 1 and all other set equal to zero. This produces (in order) the triangles $`E_2,E_5,E_6,E_3,E_7,E_4`$ and $`Z_1`$ displayed below: $$\begin{array}{c}E_2:(1,0)(0,0)(0,1)\\ \\ \begin{array}{c}1\\ 00\\ 00\\ 0000\end{array}\end{array}\begin{array}{c}E_3:(0,0)(1,0)(0,1)\\ \\ \begin{array}{c}0\\ 01\\ 00\\ 0100\end{array}\end{array}$$ $$\begin{array}{c}E_4:(0,1)(1,0)(0,0)\\ \\ \begin{array}{c}0\\ 00\\ 00\\ 1000\end{array}\end{array}\begin{array}{c}E_5:(0,1)(0,0)(1,0)\\ \\ \begin{array}{c}0\\ 00\\ 11\\ 0000\end{array}\end{array}\begin{array}{c}E_6:(0,0)(0,1)(1,0)\\ \\ \begin{array}{c}0\\ 00\\ 00\\ 0001\end{array}\end{array}$$ $$\begin{array}{c}E_7:(1,0)(1,0)(1,0)\\ \\ \begin{array}{c}0\\ 10\\ 01\\ 0100\end{array}\end{array}\begin{array}{c}Z_1:(0,0)(1,1)(1,0)\\ \\ \begin{array}{c}0\\ 00\\ 01\\ 0110\end{array}\end{array}$$ $`(7.16)`$ These are all genuine BZ triangles except for $`Z_1`$ which has some negative entries. However, at this level, there are no relations between these elementary solutions (the basis vectors are independent), hence the decomposition of any solution in terms of these 7 basic ones is unique. All solutions are then freely generated from the following function: $$G=\frac{1}{(1E_2)(1E_3)(1E_4)(1E_5)(1E_6)(1E_7)(1Z_1)}$$ $`(7.17)`$ To recover the generating function for all tensor products from the above expression, we need to project out terms that lead to triangles with negative entries. To achieve this, we introduce the grading variables associated to the above couplings (compare the above triangles with the general form given in (7.11)): \| \| $`E_2:M_{13},`$ \| $`E_3:L_{12}L_{23}`$ \| $`E_4:N_{13}`$ \| \| --- \| --- \| --- \| --- \| \| \| $`E_5:M_{12}M_{23}`$ \| $`E_6:L_{13}`$ \| $`E_7:L_{12}M_{12}N_{12}`$ \| \| \| \| $`Z_1:L_{12}^1M_{12}^1N_{23}`$ \| \| $`(7.18)`$ Our generating function follows from the projection of the above function $`G`$, re-expressed in terms of the grading variables, to positive powers of $`L_{12}`$ and $`M_{12}`$. Equivalently, one can re-scale $`L_{12}`$ by $`x`$ and $`M_{12}`$ by $`y`$ and project to positive powers of $`x`$ and $`y`$ and set $`x=y=1`$ in the result. This is equivalent to the rescaling $$E_3xE_3E_5yE_5E_7xyE_7Z_1x^1y^1Z_1$$ $`(7.19)`$ We are thus led to consider $$\underset{}{\overset{x}{\Omega }}\underset{}{\overset{y}{\Omega }}G(E_2,xE_3,\mathrm{},x^1y^1Z_1)$$ $`(7.20)`$ Keeping only those terms which depend explicitly upon $`x`$ or $`y`$, we have then $$\begin{array}{cc}\hfill \underset{}{\overset{x}{\Omega }}\underset{}{\overset{y}{\Omega }}& \frac{1}{(1xE_3)(1yE_5)(1xyE_7)(1x^1y^1Z_1)}\hfill \\ & =\frac{1}{(1xE_3)(1yE_5)(1E_7Z_1)}\left(\frac{1}{1xyE_7}+\frac{x^1y^1Z_1}{1x^1y^1Z_1}\right)\hfill \end{array}$$ $`(7.21)`$ No more work is needed for the first term. For the second one, we have $$\begin{array}{cc}\hfill \underset{}{\overset{x}{\Omega }}\underset{}{\overset{y}{\Omega }}& \frac{x^1y^1Z_1}{(1xE_3)(1E_7Z_1)(1x^1Z_1E_5)}\left(\frac{yE_5}{1yE_5}+\frac{1}{1x^1y^1Z_1}\right)\hfill \\ & =\underset{}{\overset{x}{\Omega }}\frac{x^1E_5Z_1}{(1E_5)(1E_7Z_1)(1xE_3)(1x^1Z_1E_5)}\hfill \\ & =\underset{}{\overset{x}{\Omega }}\frac{x^1E_5Z_1}{(1E_5)(1E_7Z_1)(1E_3E_5Z_1)}\left(\frac{xE_3}{1xE_3}+\frac{1}{1x^1Z_1E_5}\right)\hfill \\ & =\frac{E_3E_5Z_1}{(1E_5)(1E_7Z_1)(1E_3E_5Z_1)(1E_3)}\hfill \end{array}$$ $`(7.22)`$ We then introduce the following two new elementary couplings $$E_1=E_7Z_1,E_8=E_3E_5Z_1$$ $`(7.23)`$ Collecting the two terms resulting from the projection, we end up with $$G^{su(3)}=\left(\underset{i=1}{\overset{8}{}}\stackrel{~}{E}_i\right)(1E_7E_8)$$ $`(7.24)`$ which is indeed the $`su(3)`$ tensor-product generating function. It is worth pointing out that the Elliot-MacMahon algorithm that has been presented here as a method distinct from the vector basis, can be reinterpreted in a way that makes the two approaches equivalent. This is done in section 3 of the first reference of . There, the elementary solutions are not obtained as above by setting successively one dependent variables equal to 1 and the others equal to 0, but in reading them off directly from the columns of the $`8\times 7`$ matrix of the matrix version of the above equation: $$\left(\begin{array}{ccccccc}1& 1& 0& 0& 0& 0& 0\\ 0& 0& 1& 1& 0& 0& 0\\ 0& 1& 0& 1& 1& 0& 1\\ 0& 0& 0& 0& 1& 1& 0\\ 0& 1& 1& 0& 1& 1& 0\\ 1& 0& 0& 0& 0& 0& 1\\ 0& 1& 0& 0& 1& 0& 1\\ 0& 1& 1& 0& 1& 0& 0\end{array}\right)\left(\begin{array}{c}m_{13}\\ n_{12}\\ m_{23}\\ n_{13}\\ n_{23}\\ l_{13}\\ l_{23}\end{array}\right)=\left(\begin{array}{c}\lambda _1\\ \lambda _2\\ \mu _1\\ \mu _2\\ \zeta _1\\ \zeta _2\\ l_{12}\\ m_{12}\end{array}\right)$$ The exponentiated version of the columns gives the elementary solutions written below. This leads to the so-called ‘crude’ generating function that is then projected onto the positive solutions by the usual method. 7.3. General aspects of the vector basis construction In general, of course, the fundamental solutions to the linear system may have non-integral values of the variables. However the corresponding terms in the generating function can be eliminated by rationalising all the denominator terms and then keeping only those terms in the numerator that have integral exponents. This suggests the following modification of MacMahon’s algorithm. Consider the system of equations $$Mx=0,xN^k$$ $`(7.25)`$ where $`M`$ is a matrix of rank $`s`$. We thus have $`k`$ variables and $`s`$ relations between them. The dimension of the vector basis is thus $`ks`$. We will denote the independent (free) variables as $`x_i`$, $`i=1,\mathrm{},ks`$ and the remaining ones as $`\stackrel{~}{x}_j`$, $`j=1,\mathrm{},s`$. To find a generating function for the solutions of this system: 1. First construct a basis in $`Q^k`$ for the solutions of $`Mx=0`$ by setting $`x_i=1`$ with all other $`x_j`$ zero ($`j=1,\mathrm{},ks,ji`$). Denote by $`\stackrel{~}{x}_j^{(1)}`$ the value of the dependent variable $`\stackrel{~}{x}_j`$ evaluated at $`x_1=1`$ with all other $`x_i`$ zero. The basis then reads $$\begin{array}{cc}& ϵ_1=(1,0,0\mathrm{},0;\{\stackrel{~}{x}_j^{(1)}\}),\hfill \\ & ϵ_2=(0,1,0\mathrm{},0;\{\stackrel{~}{x}_j^{(2)}\}),\hfill \\ & \mathrm{}\hfill \\ & ϵ_{ks}=(0,0,0,\mathrm{},1;\{\stackrel{~}{x}_j^{(ks)}\})\hfill \end{array}$$ $`(7.26)`$ By construction, the $`ϵ_i`$’s are linearly independent. However notice that in general the $`\stackrel{~}{x}_j^{(i)}`$ might be rational. 2. From the form of the $`ϵ_i`$’s, it follows that any solution to (7.25) can be written as $`_ic_iϵ_i`$ with $`c_i`$ non-negative integers. In particular this means that every solution to (7.25) corresponds to a term in the generating function: $$G(X)=\frac{1}{(1X^{ϵ_1})(1X^{ϵ_2})\mathrm{}(1X^{ϵ_s})}$$ $`(7.27)`$ where $`X_1,\mathrm{},X_k`$ are grading variables. 3. $`G(x)`$ may contain negative or fractional exponents due to the occurrence of $`\stackrel{~}{x}_j^{(i)}`$ in the exponents. These are eliminated by first using MacMahon’s algorithm to eliminate any negative exponents and then rationalising denominators and keeping only terms with integral exponents in the numerators. The result is the generating function for the solutions to (7.25). This algorithm, however, does not seem to be optimal in all case. 7.4. Multiple $`su(2)`$ products from the vector basis construction A simple and different application of the formalism just developed is furnished by the analysis of $`su(2)`$ quadruple tensor products. This application is different in that it does not rely on the triangle description and as such, its formulation is less direct. This does not mean however that there are no diagrammatic representations for the quadruple product. In fact, having a set of inequalities, we can transform then into equalities, as it is done below, and from them set up a diagrammatic representation. In the present case, it would correspond to two adjacent $`su(2)`$ triangles, one upside down, with their adjacent sides forced to be equal. However, our analysis will not rely on such a description. It will serve as a preparation the somewhat more complicated $`sp(4)`$ example treated in the following section. The Diophantine description of this problem has been presented in section 4.2. It is based on the two inequalities (4.11) which are readily transformed into equalities by the introduction of two non-negative integers $`a_1,a_2`$: $$\lambda _1=n_{12}+a_1\lambda _1+n_{11}n_{12}=m_{12}+a_2$$ $`(7.28)`$ However this system does not contain any reference to the variable $`m_{11}`$ and for this reason we introduce the further constraint $`m_{11}0`$ which calls for a new non-negative integer variable: $$m_{11}=a_3$$ $`(7.29)`$ We have thus a total of 8 variables : $`\{\lambda _1,n_{11},n_{12},m_{11},m_{12},a_1,a_2,a_3\}`$ and 3 equations. There are thus 5 independent variables, chosen to be $`\{a_1,a_2,a_3,n_{12},m_{12}\}`$. The basis vectors, with components ordered as follows $$(a_1,a_2,a_3,n_{12},m_{12};\lambda _1,n_{11},m_{11})$$ $`(7.30)`$ are obtained by successively setting equal to 1 one of $`\{a_1,a_2,a_3,n_{12},m_{12}\}`$ and the others equal to 0. These basis vectors together with their exponentiated version written in terms of appropriate grading variables read: \| \| $`(1,0,0,0,0;1,1,0)`$ \| $`:L_1N_{11}^1𝒜_1`$ \| \| --- \| --- \| --- \| \| \| $`(0,1,0,0,0;0,1,0)`$ \| $`:N_{11}𝒜_2`$ \| \| \| $`(0,0,1,0,0;0,0,1)`$ \| $`:M_{11}𝒜_3`$ \| \| \| $`(0,0,0,1,0;1,0,0)`$ \| $`:L_1N_{12}`$ \| \| \| $`(0,0,0,0,1;0,1,0)`$ \| $`:N_{11}M_{12}`$ \| $`(7.31)`$ The desired generating function is obtained from the projection to positive powers of $`N_{11}`$ of the function $$\frac{1}{(1L_1N_{11}^1𝒜_1)(1N_{11}𝒜_2)(1L_1N_{12})(1N_{11}M_{12})(1M_{11}𝒜_3)}$$ $`(7.32)`$ The projection operation is done by the familiar method and the result, after setting all $`𝒜_i=1`$ is $$G=\frac{1L_1N_{11}M_{12}}{(1L_1N_{12})(1L_1M_{12})(1L_1)(1N_{11}M_{12})(1N_{11})(1M_{11})}$$ $`(7.33)`$ from which we read of the 6 elementary couplings $`E_1,\mathrm{},E_6`$ (in the order where they appear in the denominator) given in (4.12) and the relation $`E_3E_4=E_2E_5`$. The above function is exactly the one derived in section 4.2. 7.5. $`sp(4)`$ diamonds and the vector basis derivation of the generating function The system of inequalities (6.1) pertaining to $`sp(4)`$ can be transformed into a system of equations in the standard way: by setting $`r_1/2=s_1`$ and $`r_2/2=s_2`$ and introducing the non-negative integers $`a_i`$, we get : \| \| $`\lambda _1=p+a_1`$ \| $`\nu _2=a_4+a_8`$ \| \| --- \| --- \| --- \| \| \| $`\lambda _2=s_1+a_2`$ \| $`a_2+p=a_3+q`$ \| \| \| $`\mu _1=q+a_5`$ \| $`a_3+s_1=a_4+s_2`$ \| \| \| $`\mu _2=s_2+a_8`$ \| $`a_5+2s_2=a_6+2s_1`$ \| \| \| $`\nu _1=a_1+a_7`$ \| $`a_6+q=a_7+p`$ \| $`(7.34)`$ This leads to a diamond-type graphical representation of the tensor product that has the advantage over the one presented in of being linear in that the sum of two diamonds is also a diamond. This is illustrated in Fig. 1. In Fig. 1, all data pertaining to the first (second) Dynkin label appear at the left (right). Dotted lines relate those two points that compose the label indicated beside it. Opposite continuous lines are constrained to be equal, with the length of a line being defined as the sum of its extremal points except for the lines delimited by the points $`(a_6,s_1)`$ and $`(a_5,s_2)`$ where the point $`s_i`$ is counted twice (the little bar besides $`s_1`$ and $`s_2`$ being a reminder of this). Explicitly, for those lines, we have thus the constraint $`a_6+2s_1=a_5+2s_2`$. Given a triple $`sp(4)`$ product, the number of such diamonds that can be drawn with non-negative entries yields the multiplicity of the product. For instance, the two diamonds that describe the triple coupling $`(1,1)(1,1)(2,0)`$ are shown in Fig 2. The dimension of the vector basis is 8 (18 variables and 10 equations, the last four equations above being linearly independent). As our free variables we choose the set $`\{s_1,s_2,p,q,a_1,a_3,a_6,a_8\}`$. The 8 basis vectors in terms of grading variables are: \| \| $`_1:L_2M_1^2N_2𝒜_4𝒜_5^2S_1`$ \| $`_2:M_1^2M_2N_2^1𝒜_4^1𝒜_5^2S_2`$ \| \| --- \| --- \| --- \| \| \| $`_3:L_1L_2^1N_1^1𝒜_2^1𝒜_7^1P`$ \| $`_4:L_2M_1N_1𝒜_2𝒜_7Q`$ \| \| \| $`_5:L_1N_1𝒜_1`$ \| $`_6:L_2N_2𝒜_2𝒜_3𝒜_4`$ \| \| \| $`_7:M_1N_1𝒜_5𝒜_6𝒜_7`$ \| $`_8:M_2N_2𝒜_8`$ \| $`(7.35)`$ The generating function is obtained by first projecting of the function $`(1_i)^1`$ to positive powers for each grading variables and then by setting all grading variables equal to 1 except for $`L_i,M_i,N_i`$’s. The $`sp(4)`$ elementary couplings are simple products of the $`_i`$’s (the following $`A_{1,2,3}`$ should not be confused with the above grading variables): \| \| $`A_1=_7`$ \| $`A_2=_5`$ \| $`A_3=_3_4`$ \| \| --- \| --- \| --- \| --- \| \| \| $`B_1=_8`$ \| $`B_2=_6`$ \| $`B_3=_1_2`$ \| \| \| $`C_1=_4`$ \| $`C_2=_2_3_6_7^2`$ \| $`C_3=_1_3_7`$ \| \| \| $`D_1=_2_3^2_6^2_7^2`$ \| $`D_2=_1`$ \| $`D_3=_2_6_7^2`$ \| $`(7.36)`$ The complete list of $`sp(4)`$ elementary couplings (6.2) are thus recovered. 8. Conclusion As was stressed in the introduction, the main purpose of this work is to prepare the ground for the analysis of fusion rules, which is the subject of a sequel paper. In this paper, we have reviewed the existing techniques for computing tensor-product generating functions and presented a comparative assessment of their virtues and limitations. We also focused on a model formulation linking generating functions to Poincaré series, an idea first introduced in and extended in . Our contribution has been to rephrase this program more explicitly, clarify some issues and to exemplify the procedure with many examples, some of which are new. An extended version of this article is available on the Los Alamos server . Acknowledgement: We thank R.T Sharp, J. Patera and M. Walton for useful discussions. L. B. thanks S. Lantagne and H. Roussel for computing guidance. REFERENCES relax1.V.G. Knizhnik and A.B. Zamolodchikov, Nucl. Phys. B247 (1984) 83. relax2.D. Gepner and E. Witten, Nucl. Phys. B278 (1986) 493. relax3.P. Di Francesco, P. Mathieu, D. Sénéchal, Conformal Field Theory, Springer Verlag 1997. relax4.C.J. Cummins, P. Mathieu and M.A. Walton, Phys. Lett. B254 (1991) 390. relax5.R. Gaskell, A. Peccia and R.T. Sharp, J. Math. Phys. 19 (1978) 727; J. Patera and R.T. Sharp, in Recent advances in group theory and their applications to spectroscopy, ed. J. Domini, New-York, Plenum; J. Patera and R.T. Sharp, in Lecture Notes in Physics (New York, Springer Verlag, 1979) vol. 94, p. 175; M. Couture and R.T. Sharp, J. Phys. A13 (1980) 1925; R.T. Gaskell and R.T. Sharp, J. Math. Phys. 22 (1981) 2736; C. Bodine and R.T. Gaskell, J. Math. Phys. 23 (1982) 2217; R.V. Moody, J. Patera and R.T. Sharp, J. Math. Phys. 24 (1983) 2387; J. Patera and R.T. Sharp, J. Phys. A13 (1983) 397; Y. Giroux, M. Couture and R..T. Sharp, J. Phys. A17 (1984) 715. relax6.P. MacMahon, Combinatory analysis, 2 vols (1917,1918), reprinted by Chelsea, third edition, 1984. relax7.M. Hongoh, R.T. Sharp and D.E. Tilley, J.Math. Phys. 15 (1974) 782. relax8.E.B. Elliott, Quart. J. Math. 34 (1903) 388. relax9.G. Huet, An algorithm to generate the basis of solutions to homogeneous Diophantine equations, Information Processing Lett. 7 (1978) 144-7. relax10.M.Couture, C.J.Cummins and R.T.Sharp, J.Phys A23 (1990) 1929. relax11.R.P. Stanley, Duke Math. J. 40 (1973) 607; Combinatorics and Commutative Algebra, (Boston: Birkhauser) (1983). relax12.R.T Sharp and D. Lee, Revista Mexicana de Fisica 20 (1971) 203. relax13.R. Froberg, An introduction to Grobner bases, Wiley, New York 1997. relax14.B. Buchberger, Applications of Grobner basis in non-linear computational geometry in Trends in Computer Algebra, Lecture Notes in Computer Science 296 ed R Jansen (Berlin: Springer) (1989) 52-80. relax15. A.D. Berenstein and A.V. Zelevinsky, J. Geom. Phys. 5 (1989) 453. relax16.G.R.E. Black, R.C King and B.G. Wybourne, J. Phys. A: Math. Gen. 16 (1983) 1555. relax17.L. Bégin, A.N. Kirillov, P. Mathieu and M. Walton, Lett. Math. Phys. 28 (1993) 257. relax18.A.D. Berenstein and A.Z. Zelevinsky, J. Algebraic Combinat. 1 (1992) 7. relax19.The original idea of looking for a diagrammatic representation of $`sp(4)`$ tensor products along theses lines is due to M. Walton. relax20.L. Bégin, C. Cummins and P. Mathieu Generating functions for tensor products, hep-th/9811113. ......................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................$``$$``$$``$$``$$``$$``$$``$$``$$``$$``$$``$$``$$`s_1`$$`s_2`$$`q`$$`p`$$`\lambda _1`$$`\mu _2`$$`\nu _1`$$`\nu _2`$$`\mu _1`$$`\lambda _2`$$`a_6`$$`a_5`$$`a_2`$$`a_3`$$`a_4`$$`a_7`$$`a_8`$$`a_1`$............................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................................ 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warning/0005/hep-ph0005271.html	ar5iv	text	# 1 Introduction ## 1 Introduction It is well known that any object put within its gravitational radius forms a black hole (BH). At present time BHs can be naturally created only in the result of gravitational collapse of stars with the mass exceeding three Solar masses in the end of their evolution. From the other hand, it has long been known that primordial black hole (PBH) formation is possible in the early Universe. PBH can form when the density fluctuations become larger then unity on a scale intermediate between Jeans length and horizon size. Other possibilities are related to the dynamics of various topological defects such as the collapse of cosmic strings from the thermal second order phase transition or to the collisions of the bubble walls created at the first order phase transitions. Formally, there is no limit on the mass of PBH that forms after the collapse of highly overdense region, it is only needed to form appropriate spectrum of initial density fluctuations at the inflation . We cannot expect a priori the same ”no–mass–limit” condition in the case of PBHs formed by topological defects, because the mass of such PBHs is defined by the correlation length of the respective phase transition. In this paper we concern with the possibility to form PBHs after the self–collapse of closed domain walls created during a second order phase transition. Such PBHs should have small masses in the case of thermal second order phase transition with usual equilibrium initial conditions. It takes place because the characteristic size of walls coincides with the correlation length of phase transition, which initiate the formation of that domain walls. Usually we deal with high temperature phase transition in the early Universe, that makes the respective correlation length quite small and does not allow to store a significant amount of energy within the individual closed vacuum wall. Moreover the phase transition can have quite complicated dynamics, namely it can be accompanied by another phase transition that creates another type of topological defects, namely, vacuum strings (see discussion in Section 2). As a result we get the mesh of vacuum strings connected by walls where the probability of the existence of closed walls is strongly suppressed . We discuss here the non–equilibrium scenario of closed vacuum walls formation that opens a new possibility for the massive PBHs production in the early Universe. The starting point is rather general and evident. If a potential of a system possesses at least two different vacuum states there are two possibilities to populate that states in the early Universe. The first one is that the Universe contains both states populated with equal probability, that takes place under the usual circumstances of thermal phase transition. The other possibility corresponds to the case when two vacuum states are populated with different probability and there are islands of the less probable vacuum, surrounded by the sea of another, more preferable, vacuum. The last possibility can take place when we go beyond the equilibrium conditions, established by pure thermal dynamics. More definitely, it is necessary to redefine effectively the correlation length of the scalar field that drives a phase transition and consequently the formation of topological defects. We will show that the only necessary ingredient for it is the existence of an effectively flat direction(s), along which the scalar potential vanishes during inflation. Then the background de– Sitter fluctuations of such effectively massless scalar field could provide non– equilibrium redefinition of correlation length and give rise to the islands of one vacuum in the sea of another one. In spite of such redefinition the phase transition itself takes place deeply in the Friedman– Robertson-Walker (FRW) epoch. After the phase transition two vacua are separated by a wall, and such a wall can be very big. The motion of closed vacuum walls has been first driven analytically in . At some moment after crossing horizon they start shrinking due to surface tension. As a result, if the wall does not release the significant fraction of its energy in the form of outward scalar waves, almost the whole energy of such closed wall can be concentrated in a small volume within its gravitational radius what is the necessary condition for PBH formation. The mass spectrum of the PBHs which can be created by such a way depends on the scalar field potential which parametrizes the flat direction during inflation and triggers the phase transition at the FRW stage. Through the paper we will deal with so called pseudo Nambu–Goldstone (PNG) potential, that is quite common for the particle physics models. The plan of this paper is as following. In Section 2 the origin of vacuum walls at the PNG potential is discussed. The possibility to form the non–equilibrium conditions for the second order phase transition at the inflantional stage is considered in Section 3. In Section 4 the minimal conditions of PBHs formation from collapsing closed vacuum walls are discussed. The mass spectrum of PBHs is evaluated numericaly for the certain choises of parameters. ## 2 Domain Walls at the PNG Potential Under Thermal Approach Domain walls are planar defects in the vacuum alignment over space. They can appear if the manifold of degenerate (or nearly degenerate) vacua of the theory is disconnected. This is generally the case if there is a discrete symmetry in the potential. Such type of symmetry can be fundamental as, for example, in the double well potential or can arise dynamically during the breakdown induced by instanton effects of a continuous symmetry (see for example ). Moreover, discrete symmetry can be got as a result of explicit breaking of a continuous symmetry due to Yukawa–type couplings as in the case of classical PNG boson or due to higher–order nonrenormalizable interactions that often takes plase in string–inspired models . In all this cases the resulting term of potential has the form $$V=\mathrm{\Lambda }^4\left(1\mathrm{cos}\frac{\varphi }{f}\right)$$ (1) where $`f`$ is the scale of $`U(1)`$ symmetry breaking. The potential describing a spontaneously broken $`U(1)`$ symmetry has standard Mexican hat form with bottom radius $`f`$ and reads $$V(\phi )=\lambda (\left\|\phi \right\|^2f^2/2)^2,$$ (2) where $`\phi =\frac{f}{\sqrt{2}}\mathrm{exp}\left(i\varphi /f\right)`$. The term (1) breaks the continuous symmetry down to the symmetry $`\varphi \varphi +2\pi nf`$, where $`n`$ is integer. Thus potential (1) has a number of discrete degenerate minima. That minima also can be non–equivalent if the manifold spanned by $`\varphi `$ is non–compact. Moreover, it is not even necessary for the non–equivalent minima to be exactly degenerate, so the symmetries we consider may be only approximate ones, although non–degeneracy may change the details of wall evolution . However, as we have noted, many reasonable models can give rise to cosine potentials like (1), therefore we are going to use it in the explicit calculations in this paper. The equation of motion that corresponds to (1) admits kink–like, domain wall solution, which interpolates between two adjacent vacua. If we will consider that wall as lying in the perpendicular to $`x`$ axis then the static domain wall solution between two vacua $`\varphi _1=0`$ and $`\varphi _2=2\pi f`$ is given by $$\varphi _{wall}(x,x_0)=4f\mathrm{arctan}\left(\mathrm{exp}\frac{xx_0}{d}\right)$$ (3) where $`x_0`$ is the location point of the center of the wall, $`d`$ is its width, $$d=\frac{f}{\mathrm{\Lambda }^2}=m^1$$ (4) and $`m`$ is the mass of PNG field. The surface energy density of the wall reads $$\sigma =_{\mathrm{}}^+\mathrm{}\mathrm{\Lambda }^4\left(1\mathrm{cos}\frac{\varphi _{wall}}{f}\right)𝑑x8f\mathrm{\Lambda }^2$$ (5) The thermal dynamics of a system invoking two successive second order phase transitions, with spontaneous $`U(1)`$ symmetry breaking and with the explicit one, results in the formation of a quite sophisticated system of topological defects, network of walls bounded by strings (see for review ). We suppose that $`f>>\mathrm{\Lambda }`$. This means that first the spontaneous $`U(1)`$ symmetry breaking takes place in the early Universe at the temperature $`Tf`$. After that moment the Universe is filled with global $`U(1)`$ strings and the phase $`\theta =\varphi /f`$ of complex scalar field $`\varphi `$ obeys the standard equation of motion $$\ddot{\theta }+3H\dot{\theta }+\frac{dV}{d\theta }=0$$ (6) As long as the Hubble expansion is fast enough to provide the domination of friction term over the gradient one, the explicit symmetry is restored. It takes place until the moment when the temperature falls down below the value of $`\mathrm{\Lambda }`$. At that period the dynamics of global string network implies that there is approximately only one string per Hubble horizon , that defines the average distance between strings. The correlation length for the second phase transition with explicit $`U(1)`$ symmetry breaking is just equal to the average distance between global strings at the moment, when it takes place and triggers the formation of domain walls . Thus the typical size of domain walls created during last phase transition corresponds to the size of Hubble horizon at the temperature $`T\mathrm{\Lambda }`$. As it has been pointed out in the Introduction, the closed domain walls could collapse to the black hole by radiating away their asphericity due to the surface tension forces . If we concern the fate of closed domain walls that could be formed after the last phase transition, we have to compare the gravitational radius $`r_g=2E/m_p^2`$ of the energy $`E`$ stored in such a closed shell, with the minimal radius that can be achieved in the result of such collapse. The typical size of a closed wall is given by the correlation length of the last phase transition and reads $`L_c=0.3m_p/(\sqrt{g_{}}\mathrm{\Lambda }^2)`$ ($`g_{}`$ counts the effectively number degrees of freedom in the plasma and according to the standard model can not be much less then $`10^2`$). Combining it with (5) the gravitational radius will read $`r_g10^3\frac{f}{\mathrm{\Lambda }^2}`$. From the other hand the numerical simulation of the collapse of closed domain wall at the PNG potential shows that the minimal radius of shrinking cannot be much smaller than the wall’s width (4). This means that it is impossible to concentrate the energy of closed walls under the gravitational radius by self–collapse if we consider the domain walls as the result of usual thermal phase transition. Moreover the system of topological defects contains preferably walls bounded by strings and open walls, while the probability of closed walls formation is strongly suppressed . To escape such no–go results of equilibrium dynamics we need to enlarge essentially the distance between strings and to increase the probability of closed walls formation. The first problem can be solved by the inflational blow up of the typical distance of string’s separation. As well we show below that the non–equilibrium conditions imposed by background inflation allow us to bias the preference of the theory to choose one vacuum state over another that increase essentially the probability of very large closed walls formation. ## 3 Phase Transition With Non–Equilibrium Initial Conditions We consider the Universe that, due to the existence of an inflaton, goes through a period of inflation and then settles down to the standard FRW geometry. In addition we introduce a complex scalar field $`\phi `$, not the inflaton, with a large radial mass $`\sqrt{\lambda }f>H_i`$, that has got Mexican hat potential (2). It makes sure that $`U(1)`$ symmetry is already broken spontaneously at the beginning of inflation in our cosmological horizon. Thereby we deal only with the phase of that complex field $`\theta =\varphi /f`$, which paramertrizes potential (1). Under this condition we come to the conclusion, that the correlation length of second order phase transition with spontaneously broken $`U(1)`$ symmetry exceeds the present cosmological horizon, and all global $`U(1)`$ strings are beyond our horizon <sup>2</sup><sup>2</sup>2Also we can consider, a complex scalar field with a large initial radial positive mass that during inflation undergo a phase transition to a spontaneously broken $`U(1)`$ symmetry . In such a case the correlation length can be smaller then our cosmological horizon, but significantly larger then in the case of usual thermal dynamics. Furthermore, we assume that $$m<<H_i,$$ (7) where $`H_i`$ is the Hubble constant during inflation. This condition implies that during inflation the potential energy of field $`\varphi `$ is much smaller than the cosmological friction term what justifies neglecting the potential until the Universe goes deeply into the FRW phase. The dynamics of such effectively massless scalar field $`\theta `$ existing on de Sitter space breaks into two parts (see for review ). First, there is a classical field $`\theta _0`$, which satisfies the standard classical equation of motion (6). During inflation, and long afterward, $`H_i`$ is very large (by assumption) compared to the potential (1). It follows that we can drop the gradient term in the equation of motion (6) and resulting equation is solved by $`\theta _0=\theta _{N_{max}}`$, where $`\theta _{N_{max}}`$ is an arbitrary constant. We emphasize that there is no information in the theory that can fix the value of $`\theta _{N_{max}}`$. It is completely arbitrary. We make standard assumption that our present horizon that has been nucleated at the $`N_{max}`$ e– folds before the end of inflationary epoch is embedded in an enormous inflation horizon, created by exponential blow up of a single causal horizon. It follows that $`\theta _{N_{max}}`$ will be the same over the inter inflationary horizon. We put $`\theta _{N_{max}}<\pi `$ without loss of generality. Second we must consider the quantum fluctuations of the phase $`\theta `$ at the de Sitter background. It is well known , that there are quantum fluctuations produced on the vacuum state of $`\theta `$ due to the boundary conditions of de Sitter space. These fluctuations are sometimes referred to as contribution to the ”Hawking temperature” of de Sitter space but, in fact, there are not true thermal effects . It makes the dynamics of phase $`\theta `$ strongly non–equilibrium leading to the non–thermal distribution of scales populated with different vacuums in the postinflationary Universe. Such fluctuations can be described by ”quasiclassical” scalar field, which contains components with wavelengths ranging from the size of particle horizon during inflation $`H_i^1`$, all the way out to the inflation horizon $`H_i^1e^{N_{max}}`$. Thus, when the wavelength of a particular fluctuation becomes greater than $`H_i^1`$, the average amplitude of this fluctuation freezes out at some value due to the large friction term in the equation of motion, whereas its wavelength grows exponentially. In the other words such a frozen fluctuation is equivalent to classical field, discussed earlier, that does not vanish after averaging over macroscopic space intervals. The vacuum contains fluctuations of every wavelengths and hence inflation leads to the creation of new regions containing the classical field of different amplitudes with scale greater than $`H_i^1`$. The average amplitude of such fluctuations for massless field generated during each time interval $`H_i^1`$ is $$\delta \theta =\frac{H_i}{2\pi f}$$ (8) In the other words the phase $`\theta `$ makes quantum step $`\frac{H_i}{2\pi f}`$ during every e–fold and in every space volume with characteristic size of the order of the $`H_i^1`$. The total number of steps during time interval $`\mathrm{\Delta }t`$ is given by $`N=H_i\mathrm{\Delta }t`$. Such a motion looks like the one–dimensional Brownian motion. The $`\theta _{N_{max}}`$ plays the role of the starting point for this Brownian motion. Thus, the initial domain containing phase $`\theta _{N_{max}}`$ increases its volume in $`e^3`$ times after one e– fold by definition and hence contains $`e^3`$ separate, causally disconnected domains of size $`H_i^1`$. Each domain is characterized by everage phase value $$\theta _{N_{max}1}=\theta _{N_{max}}\pm \delta \theta .$$ (9) In half of these domains the phases evolve toward $`\pi `$ while in the other domains they move toward zero. This process is duplicated in each volume of size $`H^1`$ during next e–fold. Now at any given scale $`l=k^1`$ the size distribution of the phase value $`\theta `$ can be described by Gauss’ law . $$P(\theta _l,l)=\frac{1}{\sqrt{2\pi }\sigma _l}\mathrm{exp}\left\{\frac{(\theta _{N_{max}}\theta _l)^2}{2\sigma _l^2}\right\},$$ (10) where the dispersion could be expressed in the following manner $$\sigma _l^2=\frac{H^2}{4\pi ^2}\underset{k_{min}}{\overset{k}{}}d\mathrm{ln}k=\frac{H^2}{4\pi ^2}\mathrm{ln}\frac{l_{max}}{l}=\frac{H^2}{4\pi ^2f^2}(N_{max}N_l),$$ (11) As long as condition (7) is satisfied the vacuum is well defined by the distribution (10) of fluctuations around chosen constant $`\theta _{N_{max}}`$. Eventually, long after the end of inflational epoch, $`H(t)`$ decreases so significantly that the gradient term in the equation of motion (6) begins to dominate over the friction term, and $`\theta `$ starts to oscillate around the one of degenerated minimum. We refer the time $`t_c`$ moment at which the condition (7) is not valid anymore as the time moment when the phase transition is triggered. The key point of our consideration is based on the following observation: At $`t_c`$ the field configuration of phase $`\theta `$ within the volume of contemporary horizon is uniform and defined by $`\theta _{N_{max}}`$. There are, however, fluctuations due to the quasi–classical random field (10) inside this volume and the spectrum of this fluctuations contains all wavelentghs up to biggest cosmological scale. Thus the field $`\theta `$ at $`t_c`$ feels the tilt of potential and must decide to which of the two vacua $`\theta _{min}=0`$ or $`\theta _{min}=2\pi `$ it should roll down at the beginning of oscillation. We deal with the situation when at the beginning of inflation the Universe contains the phase $`\theta _{N_{max}}<\pi `$ and hence the final state of the main part of the present particle horizon is $`\theta _{min}=0`$. On the contrary, there will be the islands with $`\theta >\pi `$, which are the results of fluctuations (9) with positive sign. The phases inside that islands will move to the final state $`\theta _{min}=2\pi `$ after the triggering of the phase transition. Thus, the strongly non–equilibrium distribution of the phase, causing by inflation, leads to the formation of islands with vacuum $`2\pi `$ in the space with zero phase at the time moment $`t_c`$. Both states are separated by closed walls at $`\theta =\pi `$. The distribution on size for such closed domain walls resembles the size distribution of fluctuations (9) that consists on the steps with positive sign and move the phase to the region $`\theta >\pi `$ during the inflational stage. The probability to have such fluctuations can be significant that provides us with the possibility to form quite large number of very big closed domain walls. ## 4 Results and Discussion As we have seen in the preceding section all regions with phase $`\theta >\pi `$ are converted into islands with vacuum $`\theta _{\mathrm{min}}=2\pi `$ surrounding by the closed walls. The size distribution of closed walls imprints the size distribution of domains filled with phase coming from fluctuations that have crossed the point $`\pi `$ during the one dimensional brownian motion. The physical size that leaves the horizon during e– fold number N ($`NN_{max}`$) reads $$l=H_i^1e^N$$ (12) This scale becomes comparable to the FRW particle horizon at the moment $$t_h=H^1e^{2N}$$ (13) It is clear that we will start to observe the selfcollapse of a closed domain wall when its size is causally connected. Approximately at the same time the wall is acquiring spherical form due to the surface tension. Thus, if the amount of energy stored in a such vacuum configuration $$E\sigma t_h^2$$ (14) is large enough, the BH can be formed in the result of its selfcollapse. More definitely, the gravitation radius of configuration should exceed the minimal size up to which it can collapse. In our case the collapse of closed domain wall, coming from potential (1), changes on repulsion at the size comparable with wall’s width . This establish cut off for the PBH’s mass spectrum at the small masses range. To evaluate numerically mass spectrum of PBHs we have to calculate the size distribution of domains that contains phases, which are at the range $`\theta >\pi `$. Suppose that at e–fold $`N=\mathrm{ln}(lH_i)`$ before the end of inflation the volume $`V(\overline{\theta },N)`$ has been filled with phase value $`\overline{\theta }`$. Then at the e–fold $`N1`$ the volume filled with average phase $`\overline{\theta }`$ obeys following iterative expression $$V(\overline{\theta },N1)=e^3V(\overline{\theta },N)+(V_U(N)e^3V(\overline{\theta },N))P(\overline{\theta },N1)\delta \theta ,$$ (15) here the $`V_U(N)e^{3N}H_i^3`$ is the volume of the Universe at $`N`$ e–fold. We applied here distribution (10). Now one can easily calculate the size distribution of domains filled with appropriate value of phase corresponding to $`N`$, with the use of expression (15). For our numerical calculations we have chosen the following reasonable values for inflational Hubble constant $`H=10^{13}`$ GeV and for the radius of PNG potential $`f=10^{14}`$ GeV (see for example ). Also we suppose that $`N_{max}=60`$ and that the total energy of the wall is smaller than the total energy of the medium inside it. The simulation has been performed for two cases that depend on the $`\mathrm{\Lambda }`$ and $`\theta _{N_{max}}`$. The results are following. \| $`\mathrm{log}_{10}\frac{M_{PBH}}{1g.}`$ \| 12 \| 13 \| 14 \| 15 \| 16 \| 17 \| 20 \| 21 \| 22 \| \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| \| $`\mathrm{log}_{10}n_{PBH}`$ \| 20.4 \| 18.4 \| 16.3 \| 14.2 \| 12.0 \| 9.79 \| 5.18 \| 2.77 \| 0.30 \| \| $`\mathrm{\Lambda }=10^8`$GeV (see for example) , $`\theta _{N_{max}}=0.65`$ \| \| \| \| \| \| \| \| \| \| \| Table 1. The total number $`n_{PBH}`$ of PBHs against their masses $`M_{PBH}`$. \| \| \| \| \| \| \| \| \| \| The results of Table 1. are fitted in such a manner to satisfy the most stronger astrophysical constraints. It is known that only PBHs with masses larger then $`10^{15}`$ g. can survive in respect to Hawking evaporation. The observations of diffused gamma ray background establish strong limit on the density fraction $`\mathrm{\Omega }_{PBH}<10^9`$ of PBHs with masses $`10^{14}÷10^{15}`$ g.. The next limit, which has to be checked is the limit on the abundance of PBHs with masses $`10^{12}÷10^{13}`$ g.. Although such light PBHs disappeared already due to Hawking radiation it could produce a large amount of entropy at the epoch of nucleosinthesis . Actually the constraints is valued for the PBH’s mass range $`10^9÷10^{13}`$ g. and reads $`\beta <10^{15}(10^9g./M_{PBH})`$, here $`\beta `$ is the density fraction of PBHs at the moment of their formation. It means that at the moment of full evaporation of such PBHs $`\tau =M_{PBH}^3/(g_{}m_p^4)`$ their contribution into the total density of the Universe should satisfy the following restriction $`\alpha <10^4`$. To check our spectrum we express the density fraction of PBHs at the moment of their full evaporation in the terms of $`n_{PBH}`$, it gives $`\alpha 10^8n_{PBH}(1g./M_{PBH})^2`$. Thus in the case of Table 1. the entropy production limit is also satisfied. The more careful consideration of particle content of $`10^{10}÷10^{13}`$g. PBH decay products shows that such PBHs are very effective sources of antiprotons at the period of their evaporation . According to the theory on non–equilibrium nucleosynthesis (see and references therein) it can change dramatically the abundance of such elements as $`{}_{}{}^{3}He`$, $`{}_{}{}^{6}Li`$ and $`{}_{}{}^{7}Li`$. This fact establishes the most stringent constraints on the density fraction of PBHs in the mass rage $`10^{10}÷10^{13}`$g. <sup>3</sup><sup>3</sup>3Here we take into account only the standard model set of particle spices. The invoking of the some supersymmetric extensions of the standard model with a large number of unstable moduli fields which could come from evaporating PBHs also would influence on the primordial chemical content of the Universe (see for example ).. Thus, taking into account the effects of non–equilibrium nucleosynthesis of $`{}_{}{}^{3}He`$ and lithium we have got following limits on the density fraction of PBHs with mass range indicated above: $`\alpha _{{}_{}{}^{3}He}<10^{15}\mathrm{\Omega }_b(M_{PBH}/1g.)k^1`$, $`\alpha _{Li}<10^{19}\mathrm{\Omega }_b(M_{PBH}/1g.)^{5/4}k^1`$, where the $`\mathrm{\Omega }_b`$ is the fraction of baryon density and $`k`$ varies from $`1/4`$ to $`1/6`$ for PBHs of different masses . These limits are also satisfied our model in Table 1. As well we can see from the Table 1. that the number of BH with masses more than $`10^{22}`$ g. is negligible, whereas BHs with masses smaller than $`10^{12}`$ g. were not produced at all because of their small gravitational radius. Another interesting case is small $`\mathrm{\Lambda }`$ limit that comes from QCD . It gives rise to massive BH. \| $`\mathrm{log}_{10}\frac{M_{PBH}}{1g.}`$ \| 28.3 \| 29.5 \| 30.6 \| 31.8 \| 33.0 \| 34.1 \| 35.3 \| 36.4 \| \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| \| $`\mathrm{log}_{10}n_{PBH}`$ \| 17.4 \| 13.0 \| 11.1 \| 9.2 \| 7.2 \| 5.1 \| 3.0 \| 0.66 \| \| $`\mathrm{\Lambda }=1`$GeV (see for exapmle ) , $`\theta _{N_{max}}=0.8`$ \| \| \| \| \| \| \| \| \| \| Table 2. The total number $`n_{PBH}`$ of PBHs against their masses $`M_{PBH}`$. \| \| \| \| \| \| \| \| \| In this case the wall width is very large and so that only rather massive BHs could be formed. This fact can be applied for explanation of giant BHs origin in the centres of galaxies (see for example ) The maximal PBH masses in the Table 2 are at least by 3–4 orders of the magnitude smaller than the masses of BHs, assumed to be present in the centres of galaxies. However the account for possible strong concentration of small mass PBHs around close-to-maximal mass PBHs can provide the evolution and collapse of PBH systems into black holes with the mass, exceeding by several orders of the magnitude the maximal PBH mass. It provides the possibility to evolve the approach to the origin of AGNs on the base of our model. To complete our discussion let us mention following. We assumed that the coherent field oscillations around the true vacua, which start after the triggering of phase transition with the amplitude proportional to the residual between the local value of phase and the local vacuum phase, are damped due to relatively strong dissipation, so that no energy density is stored in the form of such oscillations. If the dissipation is not effective, as it is the case for invisible axion , the inhomogeneity of the phase distribution before the phase transition results in the inhomogeneity of the energy density distribution of coherent field oscillations. Then such energy density contributes into the total cosmological density and its large scale inhomogeneity induces effects of anisotropy of relic radiation. In this case one can strongly constrain the model parameters from the observational upper limits on the total cosmological density and on the possible anisotropy of thermal background radiation induced by isocurvature perturbations . Leaving a more completely study of the last two issues to future publications, we can conclude that the non–equilibrium dynamics of scalar field with effectively flat direction can give a new spin to the PBHs formation mechanisms applying collapse of domain walls. Acknowledgements. The work of SGR and MYuK was partially performed in the framework of Section ”Cosmoparticle physics” of Russian State Scientific Technological Program ”Astronomy. Fundamental Space Research”, with the support of Cosmion-ETHZ and Epcos-AMS collaborations. ASS and MYuK acknowledge supporte from Khalatnikov–Starobinsky school (grant 00–15–96699).
warning/0005/gr-qc0005017.html	ar5iv	text	# 1. Introduction ## 1. Introduction Recent progress in quantum gravity has made it possible for the first time to directly compute the Bekenstein-Hawking entropy of a black hole by counting microscopic states. Such calculations have been performed both in string theory and in “quantum geometry” , and while significant questions remain in each approach, the outlook seems promising. The agreement between these very different approaches, and the agreement of both with semiclassical calculations that know nothing of the details of quantum gravity, has suggested that the asymptotic behavior of the density of states may be determined by some simple universal feature, such as the algebra of diffeomorphisms at the horizon . As Kaul and Majumdar have recently stressed , however, different quantum theories of gravity may lead to different higher order corrections to the Bekenstein-Hawking entropy. These corrections may display differences, or—less probably but more interestingly—relations among quantizations. In Ref. , Kaul and Majumdar compute the lowest order corrections to the Bekenstein-Hawking entropy in a particular formulation of the “quantum geometry” program of Ashtekar et al. They find that the leading correction is logarithmic, with $$S\frac{A}{4G}\frac{3}{2}\mathrm{ln}\left(\frac{A}{4G}\right)+\text{const.}+\mathrm{}$$ (1.1) Although the existing computations of black hole entropy have very different physical starting points, most use techniques from two-dimensional conformal field theory at an intermediate stage. This may be no more than a useful trick—the Cardy formula makes it particularly easy to count states in a two-dimensional conformal field theory—but there are suggestions that such conformal field theories really provide a universal description of low-energy black hole thermodynamics. Whatever the origin of the conformal symmetry, however, the same trick that allows us to determine the asymptotic density of states also permits a direct computation of the leading quantum corrections to the Bekenstein-Hawking entropy. In this paper, I compute the logarithmic terms in the Cardy formula and use the results to obtain quantum corrections to black hole entropy. I examine a number of approaches, including Strominger’s asymptotic symmetry analysis of the (2+1)-dimensional BTZ black hole ; the methods of Ref. , which are based on the behavior of symmetries at the black hole horizon in any dimension; and the string theoretical counting of D-brane states for BPS black holes . In all cases, I find qualitative agreement with the “quantum geometry” result (1.1), with answers that differ by a factor of two in the coefficient of the logarithm and (sometimes) an additional term that depends on conserved charges. I suggest that these differences may be traced back to an ambiguity in the treatment of angular momentum and other conserved quantities, which may lead to the counting of different sets of states. If this is the case, these results represent a surprising new universality in the logarithmic corrections to the Bekenstein-Hawking entropy. ## 2. Logarithmic Corrections to the Cardy Formula The recent “first principles” computations of black hole entropy, whatever their physical starting point, typically rely at some critical stage on the Cardy formula for the density of states in a two-dimensional conformal field theory. I will start by reviewing the derivation of this formula, in order to obtain logarithmic corrections to the density of states. Let us begin with an arbitrary two-dimensional conformal field theory with central charge $`c`$, with the standard Virasoro algebra $`[L_m,L_n]`$ $`=`$ $`(mn)L_{m+n}+{\displaystyle \frac{c}{12}}m(m^21)\delta _{m+n,0}`$ $`[\overline{L}_m,\overline{L}_n]`$ $`=`$ $`(mn)\overline{L}_{m+n}+{\displaystyle \frac{c}{12}}m(m^21)\delta _{m+n,0}`$ (2.1) $`[L_m,\overline{L}_n]`$ $`=`$ $`0`$ for the generators $`L_n`$, $`\overline{L}_n`$ of holomorphic and antiholomorphic diffeomorphisms. The partition function on the two-torus of modulus $`\tau =\tau _1+i\tau _2`$ is defined to be $$Z(\tau ,\overline{\tau })=\mathrm{𝑇𝑟}e^{2\pi i\tau L_0}e^{2\pi i\overline{\tau }\overline{L}_0}=\rho (\mathrm{\Delta },\overline{\mathrm{\Delta }})e^{2\pi i\mathrm{\Delta }\tau }e^{2\pi i\overline{\mathrm{\Delta }}\overline{\tau }}.$$ (2.2) For a unitary theory, $`\rho `$ is the number of states with eigenvalues $`L_0=\mathrm{\Delta }`$, $`\overline{L}_0=\overline{\mathrm{\Delta }}`$, as can be seen by inserting a complete set of states into the trace. For a nonunitary theory, $`\rho `$ is the difference between the number of positive- and negative-norm states with appropriate eigenvalues. If we could somehow determine the partition function, we could extract the density of states by contour integration. Treat $`\tau `$ and $`\overline{\tau }`$ as independent complex variables (this is not necessary, but it simplifies the computation), and let $`q=e^{2\pi i\tau }`$ and $`\overline{q}=e^{2\pi i\overline{\tau }}`$. Then $$\rho (\mathrm{\Delta },\overline{\mathrm{\Delta }})=\frac{1}{(2\pi i)^2}\frac{dq}{q^{\mathrm{\Delta }+1}}\frac{d\overline{q}}{\overline{q}^{\overline{\mathrm{\Delta }}+1}}Z(q,\overline{q}),$$ (2.3) where the integrals are along contours that enclose $`q=0`$ and $`\overline{q}=0`$. Of course, it is rare that we actually know $`Z(q,\overline{q})`$. But Cardy has shown that it is still possible to relate the behavior of the partition function at high “energy” to its simpler behavior at low “energy,” thus giving us some control over the integral (2.3). Cardy’s basic result is that the quantity $$\mathrm{𝑇𝑟}e^{2\pi i(L_0\frac{c}{24})\tau }e^{2\pi i(\overline{L}_0\frac{c}{24})\overline{\tau }}=e^{\frac{\pi c}{6}\tau _2}Z(\tau ,\overline{\tau })$$ (2.4) is modular invariant, and in particular invariant under the large diffeomorphism $`\tau 1/\tau `$ that interchanges the circumferences of the torus. The argument is universal, involving only some general properties of conformal field theory. We can use this result to attempt to evaluate the integral (2.3) by steepest descent. To do so, let $`\mathrm{\Delta }_0`$ be the lowest eigenvalue of $`L_0`$ (often but not always zero), and define $$\stackrel{~}{Z}(\tau )=\rho (\mathrm{\Delta })e^{2\pi i(\mathrm{\Delta }\mathrm{\Delta }_0)\tau }=\rho (\mathrm{\Delta }_0)+\rho (\mathrm{\Delta }_1)e^{2\pi i(\mathrm{\Delta }_1\mathrm{\Delta }_0)\tau }+\mathrm{}$$ (2.5) (For simplicity, I have suppressed the $`\overline{\tau }`$ dependence.) It is then straightforward to show<sup></sup><sup></sup>*See for details. A closely related but more general result, obtained from a sophisticated number theoretical analysis that yields a much more complete description of the asymptotics, is discussed in . that $$\rho (\mathrm{\Delta })=𝑑\tau e^{2\pi i\mathrm{\Delta }\tau }e^{2\pi i\mathrm{\Delta }_0\frac{1}{\tau }}e^{\frac{2\pi ic}{24}\tau }e^{\frac{2\pi ic}{24}\frac{1}{\tau }}\stackrel{~}{Z}(1/\tau ).$$ (2.6) By construction, $`\stackrel{~}{Z}(1/\tau )`$ approaches a constant, $`\rho (\mathrm{\Delta }_0)`$, for large $`\tau _2`$, so the integral (2.6) can safely be evaluated by steepest descents provided that the imaginary part of $`\tau `$ is large at the saddle point. The integral we need has the form $$I[a,b]=𝑑\tau e^{2\pi ia\tau +\frac{2\pi ib}{\tau }}f(\tau ).$$ (2.7) The argument of the exponent is extremal at $`\tau _0=\sqrt{b/a}`$, and expanding around this extremum, we find $$I[a,b]𝑑\tau e^{4\pi i\sqrt{ab}+\frac{2\pi ib}{\tau _0^3}(\tau \tau _0)^2}f(\tau _0)=\left(\frac{b}{4a^3}\right)^{1/4}e^{4\pi i\sqrt{ab}}f(\tau _0).$$ (2.8) In particular, if $`\mathrm{\Delta }_0`$ is small ($`\mathrm{\Delta }_0c`$) and $`\mathrm{\Delta }`$ is large, the integral (2.6) yields $$\rho (\mathrm{\Delta })\left(\frac{c}{96\mathrm{\Delta }^3}\right)^{1/4}\mathrm{exp}\left\{2\pi \sqrt{\frac{c\mathrm{\Delta }}{6}}\right\}.$$ (2.9) The exponential term in (2.9) gives the standard Cardy formula, but we have now found the leading correction as well. We must next ask how reliable this approximation is. For $`f(\tau )`$ constant, the integral (2.7) can be performed explicitly, yielding a Bessel function, whose asymptotic behavior agrees with (2.8) with additional terms that are exponentially suppressed. Corrections from the nonconstancy of $`f(\tau )`$ may also be computed, and for $`f(\tau )=\stackrel{~}{Z}(1/\tau )`$, it is easy to check that these are again exponentially suppressed. For large $`\mathrm{\Delta }`$, the expression (2.9) thus gives a reliable first-order correction to the standard Cardy formula. It should be stressed that the central charge $`c`$ appearing in (2.9) is the full central charge of the conformal field theory. In general, $`c`$ will consist of a “classical” term $`c_{\text{class}}`$, which already appears in the Poisson brackets of the $`L_n`$, plus a quantum correction that can change the exponent in (2.9) from its “classical” value. In contrast to the prefactor in (2.9), this correction is likely to be highly model-dependent. Nevertheless, a bit can be said about its general features; see Appendix B. ## 3. The BTZ Black Hole As our first application of Eqn. (2.9), let us evaluate the logarithmic corrections to Strominger’s derivation of the entropy of the BTZ black hole. This (2+1)-dimensional black hole has a metric $$ds^2=N^2dt^2+N^2dr^2+r^2\left(d\varphi +N^\varphi dt\right)^2$$ (3.1) with $$N=\left(8GM+\frac{r^2}{\mathrm{}^2}+\frac{16G^2J^2}{r^2}\right)^{1/2},N^\varphi =\frac{4GJ}{r^2}(\|J\|M\mathrm{}),$$ (3.2) and solves the vacuum Einstein field equations with a cosmological constant $`\mathrm{\Lambda }=1/\mathrm{}^2`$. The spacetime is thus asymptotically anti-de Sitter, and has an outer (event) and an inner horizon at $$r_\pm {}_{}{}^{2}=4GM\mathrm{}^2\{1\pm [1\left(\frac{J}{M\mathrm{}}\right)^2]^{1/2}\},$$ (3.3) i.e., $$M=\frac{r_+{}_{}{}^{2}+r_{}^2}{8G\mathrm{}^2},J=\frac{r_+r_{}}{4G\mathrm{}}.$$ (3.4) As Brown and Henneaux first noted , the asymptotic symmetries of Einstein gravity in 2+1 dimensions with negative $`\mathrm{\Lambda }`$ are described by a pair of Virasoro algebras, with central charges $$c=\overline{c}=\frac{3\mathrm{}}{2G}.$$ (3.5) Intuitively, these are the symmetries of the adS “cylinder at infinity,” obtained by restricting diffeomorphisms in the bulk to those that preserve adS boundary conditions. The central charges (3.5) are classical, but they will presumably be reflected in any quantum theory of gravity. Thus in any quantum theory of gravity that has the correct classical limit, the fields should transform under representations of these Virasoro algebras. Given some plausible assumptions —for example, that $`\mathrm{\Delta }_0`$ is small—one should therefore be able to use the Cardy formula to compute the asymptotic density of states. Now, the generators of the Brown-Henneaux Virasoro algebras can be computed explicitly: they are simply the Hamiltonian and momentum constraints of general relativity smeared against appropriate vector fields. For the BTZ black hole, one finds that up to an ambiguous additive constant , $$\mathrm{\Delta }=\frac{(r_++r_{})^2}{16G\mathrm{}},\overline{\mathrm{\Delta }}=\frac{(r_+r_{})^2}{16G\mathrm{}}.$$ (3.6) As Strominger observed, Eqns. (3.5) and (3.6) can be used to evaluate the exponent in (2.9), yielding $$2\pi \sqrt{\frac{c\mathrm{\Delta }}{6}}+2\pi \sqrt{\frac{\overline{c}\overline{\mathrm{\Delta }}}{6}}=\frac{2\pi r_+}{4G},$$ (3.7) giving the standard Bekenstein-Hawking entropy for the (2+1)-dimensional black hole. It is now easy to read off the logarithmic corrections to the entropy. From (2.9), $$\rho (\mathrm{\Delta },\overline{\mathrm{\Delta }})\frac{8G\mathrm{}^2}{(r_+^2r_{}^2)^{3/2}}\mathrm{exp}\left\{\frac{2\pi r_+}{4G}\right\}.$$ (3.8) Thus $$S\frac{2\pi r_+}{4G}\frac{3}{2}\mathrm{ln}\left(\frac{r_+^2r_{}^2}{G^2}\right)+\text{const.}=\frac{2\pi r_+}{4G}\frac{3}{2}\mathrm{ln}\frac{2\pi r_+}{G}\frac{3}{2}\mathrm{ln}\kappa \mathrm{}+\text{const.}$$ (3.9) where $$\kappa =\frac{r_+^2r_{}^2}{\mathrm{}^2r_+}$$ (3.10) is the surface gravity. The logarithmic terms in (3.9) have the same form as those found by Kaul and Majumdar for the nonrotating (3+1)-dimensional black hole. Because of the term involving $`\kappa `$, however, the coefficients are different; in particular, if one restricts to zero angular momentum ($`r_{}=0`$), one finds logarithmic term that differs from (1.1) by a factor of two. There is an alternative derivation of the BTZ black hole entropy, first proposed in Ref. , that directly counts states of an induced $`\text{SL}(2,\text{I}\text{R})\times \text{SL}(2,\text{I}\text{R})`$ Wess-Zumino-Witten model at the black hole horizon. While the horizon radius $`r_+`$ has a natural expression in such a WZW model, it is difficult to fix $`r_{}`$, so one instead fixes a conjugate variable (essentially a component of the triad at the horizon). As explained in Appendix B of Ref. , the resulting partition function can be viewed as a functional Fourier transform of the partition function (2.2), where the central charge in (2.2) for an $`\text{SL}(2,\text{I}\text{R})`$ WZW model is $`c3`$. One finds $$Z=\underset{N}{}\rho (N)\mathrm{exp}\left\{2\pi i\tau \left(N\frac{k^2r_+^2}{\mathrm{}^2}\right)\right\}$$ (3.11) with $$\rho (N)\underset{n=0}{\overset{N}{}}\rho _0(N)\rho _0(Nn),$$ (3.12) where $`\rho _0`$ is the partition function for an $`\text{SL}(2,\text{I}\text{R})`$ WZW model and $`k=\mathrm{}/4G`$. In the large $`k`$ limit, the three oscillators of $`\text{SL}(2,\text{I}\text{R})`$ can be treated independently, and (2.9) gives $$\rho (N)\underset{n=0}{\overset{N}{}}N^{3/4}(Nn)^{3/4}\mathrm{exp}\left\{\sqrt{2}\pi \left(\sqrt{n}+\sqrt{Nn}\right)\right\}.$$ (3.13) We can evaluate this expression by approximating the sum as an integral and using the method of steepest descents, obtaining $$\rho (N)N^{3/4}e^{2\pi \sqrt{N}}.$$ (3.14) In the formalism of Ref. , the partition function (3.11) is subject to a physical state condition $`\mathrm{\Delta }=0`$, that is, $`N=k^2r_+^2/\mathrm{}^2`$. We thus obtain a density of states $$\rho (N)(r_+/G)^{3/2}\mathrm{exp}\left\{\frac{2\pi r_+}{4G}\right\}.$$ (3.15) The resulting logarithmic correction to the entropy agrees exactly with that of Kaul and Majumdar. It is interesting to note that logarithmic corrections of this sort are absent in the Euclidean path integral approach to BTZ black hole entropy. The first-order corrections were calculated in that formalism in Ref. ; they give an exponentially suppressed contribution to the density of states, with no power law prefactor that would translate into a logarithmic correction to the entropy. ## 4. Horizon Conformal Field Theory The conformal field theory derivations of the preceding section rely on special features of the (2+1)-dimensional black hole. The results are more general than they might appear at first sight, since many black holes in string theory have a near-horizon structure that looks like that of a BTZ black hole . Others do not, however, and the methods do not directly generalize directly to higher dimensions. A different conformal field theory approach to black hole entropy has recently been proposed, based on a possible universal Virasoro algebra at the horizon . This algebra is obtained by treating the horizon as a boundary and considering the behavior of the algebra of diffeomorphisms of the “$`r`$$`t`$ plane” near the horizon. The proper choice of boundary conditions is not entirely clear, but several different choices give rise to a Virasoro algebra with central charge $$c=\frac{3A}{2\pi G}\frac{\beta }{\kappa }$$ (4.1) and an $`L_0`$ eigenvalue $$\mathrm{\Delta }=\frac{A}{16\pi G}\frac{\kappa }{\beta },$$ (4.2) where $`A`$ is the horizon area (in any dimension), $`\kappa `$ is the surface gravity, and $`\beta `$ is an undetermined periodicity. An analysis of the Liouville theory near the horizon obtained from dimensional reduction of Einstein gravity gives a similar result . It is easy to see that these values of $`c`$ and $`\mathrm{\Delta }`$, inserted into the Cardy formula, give the standard Bekenstein-Hawking entropy. But we can now go further, and compute the logarithmic corrections: Eqn. (2.9) yields $$\rho (\mathrm{\Delta })\frac{c}{12}\left(\frac{A}{8\pi G}\right)^{3/2}\mathrm{exp}\left\{\frac{A}{4G}\right\}.$$ (4.3) If we can now choose $`\beta `$ to be such that $`c`$ is a universal constant, independent of $`A`$, we find an entropy $$S\frac{A}{4G}\frac{3}{2}\mathrm{ln}\left(\frac{A}{4G}\right)+\text{const.}+\mathrm{},$$ (4.4) in agreement with the result (1.1) of Kaul and Majumdar. ## 5. String Theory Much of the current interest in black hole entropy was sparked by the discovery by Strominger and Vafa that for extremal (BPS) black holes in string theory, one could compute the Bekenstein-Hawking entropy by counting D-brane states. The relevant configurations are obtained by compactifying a suitable string theory on a manifold with the topology $`M\times S^1`$$`M`$ is $`K3`$ in the case considered in Ref. , but may be different for other black holes—and considering a collection of D-branes wrapped around cycles of $`M\times S^1`$. To count states, one takes the radius of the $`S^1`$ factor to be large compared to $`M`$, and describes the low-energy excitations of the D-branes in terms of fields moving on $`S^1`$. This description involves a weak-coupling approximation, and is unrealistic for a true black hole. For BPS configurations, however, one can argue that the density of states is protected by nonrenormalization theorems when the coupling is increased. The sigma model describing the excitations on $`S^1`$ is a two-dimensional conformal field theory, albeit a conformal field theory very different from those considered in the preceding sections of this paper. Thus Cardy’s formula may be used to count states, and we may again appeal to Eqn. (2.9) to find the logarithmic corrections to the entropy. For the five-dimensional black holes investigated by Strominger and Vafa, the relevant conformal field theory has central charge $$c3Q_F^2$$ (5.1) and $`L_0`$ eigenvalue $$\mathrm{\Delta }=Q_H,$$ (5.2) where $`Q_F`$ is the Ramond-Ramond charge and $`Q_H`$ is the momentum around the $`S^1`$. These quantities translate into charges of the associated black holes, and the entropy obtained from the exponential term in the Cardy formula turns out to be $`A/4G`$, where the horizon area is $$A=8\pi \sqrt{\frac{Q_HQ_F^2}{2}}.$$ (5.3) By (2.9), the leading correction to the entropy is thus $$S\frac{A}{4G}\frac{3}{2}\mathrm{ln}\left(\frac{A}{4G}\right)+2\mathrm{ln}Q_F+\text{const.}+\mathrm{}$$ (5.4) We again obtain a logarithmic correction identical to that of Kaul and Majumdar, along with an extra term depending on the Ramond-Ramond charge. A similar analysis can be performed for the four-dimensional black holes of Horowitz et al. , which are obtained from a string theory compactified on a six-torus, with charges $`Q_2`$, $`Q_5`$, $`Q_6`$, and $`n`$ carried by two-branes, five-branes, six-branes, and strings wrapped around cycles of the torus. In the extremal limit, one obtains a conformal field theory with $$c=6Q_2Q_5Q_6,\mathrm{\Delta }=n,$$ (5.5) corresponding to a black hole with horizon area $$A=2\pi \sqrt{Q_2Q_5Q_6n}.$$ (5.6) The leading correction to the entropy is thus $$S\frac{A}{4G}\frac{3}{2}\mathrm{ln}\left(\frac{A}{4G}\right)+\mathrm{ln}(Q_2Q_5Q_6)+\text{const.}+\mathrm{}$$ (5.7) Again, we obtain an area term of the form (1.1), plus corrections that depend on conserved charges. Larsen has proposed a related conformal field theoretical picture for a large class of five-dimensional charged, rotating black holes . Such black holes are characterized by a mass $`M`$, three conserved charges $`Q_i`$, and two angular momenta $`J_{R,L}`$. The mass and charges can be parametrized as $$M=\frac{1}{2}\mu \underset{i=1}{\overset{3}{}}\mathrm{cosh}2\delta _i,Q_i=\frac{1}{2}\mathrm{sinh}2\delta _i.$$ (5.8) Larsen shows that many of the properties of such black holes can be understood in terms of a left- and right-moving conformal field theory, each with central charge $`c=6`$, with $`L_0`$ eigenvalues $$\mathrm{\Delta }_{R,L}=\frac{1}{4}\mu ^3\left(\underset{i}{}\mathrm{cosh}\delta _i\underset{i}{}\mathrm{sinh}\delta _i\right)^2J_{R,L}^2.$$ (5.9) The Bekenstein-Hawking entropy follows from Cardy’s formula, and can be written in the form $$S_{BH}=S_R+S_L,S_{R,L}=2\pi \sqrt{\frac{c\mathrm{\Delta }_{R,L}}{6}}.$$ (5.10) Since the central charge $`c`$ is just a number, it is easy to use Eqn. (2.9) to obtain the logarithmic corrections to (5.10). Using the results of , we obtain $$S\frac{A_+}{4G}\frac{3}{2}\mathrm{ln}(S_LS_R)+\text{const.}+\mathrm{}=\frac{A_+}{4G}\frac{3}{2}\mathrm{ln}\left[\left(\frac{A_+}{4G}\right)^2\left(\frac{A_{}}{4G}\right)^2\right]+\text{const.}+\mathrm{},$$ (5.11) where $`A_\pm `$ are the areas of the outer and inner horizons. This result may be compared to the nearly identical expression (3.9) for the BTZ black hole. ## 6. Speculations Although the models of black hole entropy considered here involve very different physical pictures of microscopic states, all use two-dimensional conformal field theory as a crucial tool. This makes it possible to compute the leading logarithmic corrections to the Bekenstein-Hawking entropy in a simple and systematic manner. The resulting entropy takes the general form $$S\frac{A}{4G}\frac{3}{2}\mathrm{ln}\left(\frac{A}{4G}\right)+\mathrm{ln}F(Q)+\text{const.}+\mathrm{},$$ (6.1) where $`F(Q)`$ is some function of angular momentum and other conserved charges. The existence of logarithmic corrections of the form $`\mathrm{ln}(A/4G)`$ is thus a general feature of black hole entropies obtained in this manner. The interesting question is whether the factor of $`3/2`$ in (6.1), which also appears in the results of Kaul and Majumdar, is also universal. The problem, of course, is that the charges in $`F(Q)`$ and the horizon area are not, in general, independent, so there is some ambiguity in the division of the right-hand side of (6.1) into separate terms. To explore this issue further, let us return to the Cardy formula (2.9), and note that it can be rewritten as $$\mathrm{ln}\rho (\mathrm{\Delta })=S_0\frac{3}{2}\mathrm{ln}S_0+\mathrm{ln}c+\text{const.},$$ (6.2) where $`S_0=2\pi \sqrt{c\mathrm{\Delta }/6}`$. Hence if the entropy is obtained from a single conformal field theory, and if the central charge $`c`$ is “universal” in the sense of being independent of the horizon area, the factor of $`3/2`$ will always appear. This sort of “universality” of the central charge seems to be a sensible requirement for any fundamental conformal field theoretical explanation of black hole entropy: one should surely be looking for a single conformal field theory to describe black holes with arbitrary masses and charges. The conformal theories of Refs. do not satisfy this demand, but they are presumably not yet the final word on horizon boundary conditions. The requirement of a single conformal field theory is less clear. One might reasonably argue that there should be left-moving and a right-moving sectors, as there are for the asymptotic symmetries of the BTZ black hole. Larsen, for instance, has suggested that these sectors could be associated with the inner and outer horizon . The presence of two conformal field theories changes the form of the leading correction in (6.2), which becomes $$\mathrm{ln}\rho (\mathrm{\Delta }_L)+\mathrm{ln}\rho (\mathrm{\Delta }_R)=S_L+S_R\frac{3}{2}\mathrm{ln}(S_LS_R)+\mathrm{ln}(c_Lc_R)+\text{const.}$$ (6.3) It is evident that if $`S_L+S_R`$ gives the standard Bekenstein-Hawking entropy, the logarithmic correction is no longer $`\frac{3}{2}\mathrm{ln}S_{BH}`$, but rather takes a form more like that of Eqn. (5.11). This difference appears to account for much of the variation among the logarithmic terms found in this paper. The difference may be related to the choice of how to treat angular momentum and other conserved charges when counting states, and thus implicitly to the choice of which states to count. For the (2+1)-dimensional black hole, for example, Strominger’s approach starts with a black hole with a fixed angular momentum, and gives an expression of the general form (6.3). The earlier approach of Ref. , on the other hand, does not require a specified angular momentum, and yields an expression of the form (6.2). It will be interesting to see how the results of Kaul and Majumdar , which are based on boundary conditions specific to a nonrotating black hole, change when more general rotating boundary conditions are incorporated. Acknowledgements This work was supported in part by Department of Energy grant DE-FG03-91ER40674. ## Appendix A. The String Level Density and the Cardy Formula The expression (2.9) for the density of states differs from the “level density” of string theory , which counts the oscillator states of a string theory. In this appendix I investigate the differences, while simultaneously obtaining a useful check of the methods of section 2. Let us begin by using the procedure that led to Eqn. (2.9) to evaluate the “partition function” $`p(n)`$ of Ramanujan and Hardy , the number of partitions of an integer $`n`$ into a sum of smaller integers. It is easy to see that $`p(n)`$ is also the number of oscillator states of a bosonic string in one transverse dimension. The generating function for $`p(n)`$ is $$G(\tau )=p(n)e^{2\pi in\tau }=\underset{n=1}{\overset{\mathrm{}}{}}\left(1e^{2\pi in\tau }\right)^1=e^{\frac{2\pi i\tau }{24}}\eta ^1(\tau ),$$ (A.1) where $`\eta (\tau )`$ is the Dedekind eta function. This is almost a chiral partition function $`Z(\tau )`$ for a $`c=1`$ conformal field theory, but not quite: the modular transformation properties of $`\eta (\tau )`$ differ slightly from the standard conformal field theory form. Indeed, $$\eta ^1(1/\tau )=(i\tau )^{1/2}\eta ^1(\tau ),$$ (A.2) so unlike Cardy’s partition function (2.4), the quantity $`e^{2\pi i\tau /24}G(\tau )`$ is not exactly modular invariant. It is straightforward to incorporate the extra factor in the transformation (A.2) into the integral (2.6), however: one must simply replace $`\stackrel{~}{Z}(1/\tau )`$ by $`(i\tau )^{1/2}\stackrel{~}{Z}(1/\tau )`$. This translates into a term $`f(\tau _0)=(i\tau _0)^{1/2}=(b/a)^{1/4}`$ in Eqn. (2.8), giving $$p(n)\frac{1}{\sqrt{48}n}e^{\pi \sqrt{\frac{2n}{3}}},$$ (A.3) in exact agreement with the asymptotic results of Ramanujan and Hardy. The generalization to a partition with $`N`$ “colors” is immediate. The relevant generating function is now $$G_N(\tau )=p_N(n)e^{2\pi in\tau }=\underset{n=1}{\overset{\mathrm{}}{}}\left(1e^{2\pi in\tau }\right)^N=\left(e^{\frac{2\pi i\tau }{24}}\eta ^1(\tau )\right)^N,$$ (A.4) giving a factor of $`(i\tau _0)^{N/2}`$ in Eqn. (2.8) and an asymptotic behavior $$p_N(n)\frac{1}{\sqrt{2}}\left(\frac{N}{24}\right)^{\frac{N+1}{4}}n^{\frac{N+3}{4}}\mathrm{exp}\left\{2\pi \sqrt{\frac{Nn}{6}}\right\}.$$ (A.5) This can be recognized as the “level density” for a bosonic string in $`N`$ transverse dimensions, as first computed by Huang and Weinberg in the context of the Veneziano model . This expression and its superstring generalization have been used by Solodukhin to examine logarithmic corrections in the string-black hole correspondence . Let us now try to understand the reason for the differences between the densities of states (2.9) in conformal field theory and the level density (A.5). The key observation is that the partition function for a scalar field in conformal field theory is not actually given by Eqn. (A.1), but is rather $$Z(\tau )=\tau _2^{1/4}e^{\frac{2\pi i\tau }{24}}\eta ^1(\tau ),$$ (A.6) where the extra factor of $`\tau _2^{1/4}`$ comes from zeta function regularization of a determinant, that is, from zero-modes of the boson. Under the transformation $`\tau 1/\tau `$, we have $`\tau _2\tau _2/\|\tau \|^2`$, and it is easily checked that the resulting factor of $`\|\tau \|^2`$ is just what is needed to compensate for the transformation (A.2) of $`\eta (\tau )`$, restoring the Cardy formula (2.9). The role of the zero-modes can be explored further by considering a string compactified on a circle of radius $`R`$. The Virasoro generators $`L_0`$ and $`\overline{L}_0`$ are then $`L_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{r}{2R}}+sR\right)^2+N`$ $`\overline{L}_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{r}{2R}}sR\right)^2+\overline{N},`$ (A.7) where $`r`$ and $`s`$ are integer-valued momentum and winding numbers and $`N`$ and $`\overline{N}`$ are the usual oscillator number operators. The level density (A.3) counts oscillator states alone, implicitly taking $`r=s=0`$. To incorporate the winding states, we should instead consider the sum $$\rho (\mathrm{\Delta },\overline{\mathrm{\Delta }})\underset{r,s}{}p\left(\mathrm{\Delta }\frac{1}{2}\left(\frac{r}{2R}+sR\right)^2\right)p\left(\overline{\mathrm{\Delta }}\frac{1}{2}\left(\frac{r}{2R}sR\right)^2\right).$$ (A.8) Approximating the sums by integrals, we can write (A.8) in terms of modified Bessel functions $`I_1(2\pi \sqrt{\mathrm{\Delta }/6})`$ and $`I_1(2\pi \sqrt{\overline{\mathrm{\Delta }}/6})`$, and it is straightforward to check that the asymptotic behavior of (A.8) is precisely that of a $`c=1`$ conformal field theory, as given by Eqn. (2.9). The Cardy formula (2.9) can thus be understood as a total density of states, including both the oscillator states counted by (A.5) and the winding states or zero-modes that can contribute to $`\mathrm{\Delta }`$ and $`\overline{\mathrm{\Delta }}`$. ## Appendix B. Corrections to the Central Charge As noted at the end of section 2, the central charge in Eqn. (2.9) generally takes the form $$c=c_{\text{class}}+\delta c,$$ (B.1) where the “classical” central charge $`c_{\text{class}}`$ is already present in the Poisson brackets of the Virasoro generators. The quantum correction $`\delta c`$ can be evaluated in some models, and can offer an interesting new test of the statistical mechanical picture. For certain black holes in string theories compactified on Calabi-Yau manifolds, for example, the one-loop contribution to $`\delta c`$ gives rise to an area-independent term in the Cardy formula , which can be interpreted in Wald’s Noether charge formalism as the entropy due to a curvature-squared term in the effective action. Since the coefficient of this term in the action is independently calculable in string theory, this agreement provides a delicate test of the statistical mechanical formalism. Similarly, it has been argued that the central charge (3.5) for the BTZ black hole should be shifted to $$c=\overline{c}=\frac{3\mathrm{}}{2G}+\beta +\frac{4G}{\mathrm{}}\gamma $$ (B.2) for constants $`\beta `$ and $`\gamma `$ that depend on details of the relevant conformal field theory . For an appropriate choice of $`\beta `$, the resulting correction to the Cardy formula reproduces the shift in entropy found in Ref. from a one-loop path integral computation. There seems to be no reason to expect the quantum corrections $`\delta c`$ to have any special universal properties. We should therefore ask whether the logarithmic corrections to the entropy described in this paper are really the leading corrections. Although I know of no completely rigorous argument, for large black holes it seems quite likely that they are: Recall first that the Bekenstein-Hawking formula for entropy can be obtained from quantum field theory in a classical black hole background . Such a result will receive quantum gravitational corrections, of course, but since the curvature near the horizon decreases as the horizon area $`A`$ increases, one expects these corrections to be small for large black holes. Equivalently, a term in the four-dimensional effective action with $`n`$ powers of the curvature will have a coupling constant of dimension $`[L]^{2n4}`$, which can multiply at most a factor of $`A^{2n}`$. Quantum corrections to the entropy involve $`n>1`$ powers of the curvature, and should thus be suppressed by powers of the horizon area. Such an effective action picture may fail for small black holes, but it should be approximately valid in the semiclassical regime, that is, for black holes whose horizons are large compared to the Planck scale. Now consider an expansion of $`\delta c/c_{\text{class}}`$ in powers of $`A`$. The leading correction to the entropy will come from the exponent of Eqn. (2.9), yielding $$S\frac{A}{4\mathrm{}G}\left(1+\frac{\delta c}{c_{\text{class}}}\right)^{1/2}.$$ (B.3) Positive powers of $`A`$ in $`\delta c/c_{\text{class}}`$ would dominate this expression for large $`A`$, leading to a breakdown of the Bekenstein-Hawking formula in precisely the regime in which it should be most reliable. As argued above, such terms are thus unlikely. The next terms in the expansion, those of order $`A^0`$, certainly can occur. Their main effect in Eqn. (B.3) will be to shift the coefficient $`1/4\mathrm{}G`$ in the Bekenstein-Hawking formula. But at least for corrections induced by ordinary matter fields, such shifts can be shown to merely correspond to renormalizations of Newton’s constant . If this is the case in general—as, once again, one would expect from the semiclassical derivations of the Bekenstein-Hawking formula—then the first observable effects of quantum corrections to the central charge will come at order $`A^1`$ in $`\delta c/c_{\text{class}}`$. Such corrections will, for large $`A`$, give area-independent terms in the entropy of the sort discussed in Refs. . As noted above, these are certainly important as independent tests of statistical mechanical formulations. But for large $`A`$, they will normally be dominated by the logarithmic corrections to the Cardy formula discussed in this paper.
warning/0005/hep-ph0005179.html	ar5iv	text	# 1 Introduction ## 1 Introduction The observation of the atmospheric neutrino by SuperKamiokande has shown the existence of the neutrino masses and the neutrino mixings. In particular, the data show that the mixing between $`\nu _\mu `$ and $`\nu _\tau `$ is favored and $`\mathrm{sin}^22\theta _{atm}1`$ and $`\mathrm{\Delta }_{atm}^23.5\times 10^3\mathrm{eV}^2`$. The solar neutrino problem is now considered to be due to the $`\nu _e`$ and $`\nu _\mu `$ oscillation, but the information on masses and mixing angles are ambiguous. Now four solutions are given. The another crucial information is given by CHOOZ group that $`\|V_{13}\|<0.16`$ We interpret these in the neutrino mixing matrix in the standard form $`V_{SF}=\left(\begin{array}{ccc}c_{12}c_{13}& s_{12}c_{13}& s_{13}e^{i\delta }\\ s_{12}c_{23}c_{12}s_{23}s_{13}e^{i\delta }& c_{12}c_{23}s_{12}s_{23}s_{13}e^{i\delta }& s_{23}c_{13}\\ s_{12}s_{23}c_{12}c_{23}s_{13}e^{i\delta }& c_{12}s_{23}s_{12}c_{23}s_{13}e^{i\delta }& c_{23}c_{13}\end{array}\right),`$ (1) where $`s_{ij}=\mathrm{sin}\theta _{ij}`$ and $`c_{ij}=\mathrm{cos}\theta _{ij}`$. The large atmospheric neutrino mixing requires $`\|s_{23}\|\|c_{23}\|1/\sqrt{2}`$ and the CHOOZ data gives the bound $`\|s_{13}\|<0.16`$. Now we face the following questions: (1) Why $`\|s_{23}\|\|c_{23}\|1/\sqrt{2}`$ ? (2) What is the size of $`s_{12}`$ ? (3) Why $`s_{13}`$ is so small ? (4) What is the size of the CP violation phase $`\delta `$ ? It is quite hard to construct the neutrino mass matrix which answers all these questions. In this note, we focus on the questions, (1), (2) and (4), by using the democratic-type mass matrix for the neutrino,. Throughout of this note, we consider the neutrino mass matrix in the basis where charged leptons mass matrix is diagonal. ## 2 The model construction We are interested in the CP violation phase. A famous model that predicts the CP violation phase is the Tri-maximal mixing scheme $`V_T={\displaystyle \frac{1}{\sqrt{3}}}\left(\begin{array}{ccc}1& 1& 1\\ \omega & \omega ^2& 1\\ \omega ^2& \omega & 1\end{array}\right),`$ (2) where $`\omega `$ is the element of $`Z_3`$ and is given by $`\omega =e^{i2\pi /3}`$, i.e., $`\omega ^3=1`$. This matrix predicts $`\|V_{13}\|=1/\sqrt{3}`$ in conflict with the CHOOZ data. However, this scheme has an quite interesting predictions; (i) The maximal CP phase, $`\delta =\pi /2`$ and (ii) the maximal CP violation, $`\|J_{CP}\|_{Tri}=1/6\sqrt{3}`$. (a) Deformation from the Tri-maximal mixing We considered the deformation of the Tri-maximal mixing matrix by the orthogonal matrix $`O`$, $`V=V_TO`$ and found quite significant predictions, $`\|s_{23}\|=\|c_{23}\|={\displaystyle \frac{1}{\sqrt{2}}},\delta ={\displaystyle \frac{\pi }{2}}.`$ (3) (b) The possible origin of the mixing matrix in the form of $`V=V_TO`$ We define the new basis ($`\psi `$ basis) which is related to the flavor eigenstate basis by $`(\psi _1,\psi _2,\psi _3)^T=V_T^{}(\nu _{eL},\nu _{\mu L},\nu _{\tau L})^T`$. Now, the matrix to realize the mixing matrix $`V=V_TO`$ is some real matrix in the $`\psi `$ basis, $`V_T^Tm_\nu V_T=\left(\begin{array}{ccc}m_1^0+\stackrel{~}{m}_1& \stackrel{~}{m}_3& \stackrel{~}{m}_2\\ \stackrel{~}{m}_3& m_2^0+\stackrel{~}{m}_2& \stackrel{~}{m}_1\\ \stackrel{~}{m}_2& \stackrel{~}{m}_1& m_3^0+\stackrel{~}{m}_3\end{array}\right),`$ (4) where $`m_i^0`$ and $`\stackrel{~}{m}_i`$ are real parameters. By inverting, we obtain $`m_\nu `$ in the flavor eigenstate basis as $`m_\nu `$ $`=`$ $`{\displaystyle \frac{m_1^0}{3}}\left(\begin{array}{ccc}1& \omega ^2& \omega \\ \omega ^2& \omega & 1\\ \omega & 1& \omega ^2\end{array}\right)+{\displaystyle \frac{m_2^0}{3}}\left(\begin{array}{ccc}1& \omega & \omega ^2\\ \omega & \omega ^2& 1\\ \omega ^2& 1& \omega \end{array}\right)+{\displaystyle \frac{m_3^0}{3}}\left(\begin{array}{ccc}1& 1& 1\\ 1& 1& 1\\ 1& 1& 1\end{array}\right)`$ (5) $`+\stackrel{~}{m}_1\left(\begin{array}{ccc}1& 0& 0\\ 0& \omega & 0\\ 0& 0& \omega ^2\end{array}\right)+\stackrel{~}{m}_2\left(\begin{array}{ccc}1& 0& 0\\ 0& \omega ^2& 0\\ 0& 0& \omega \end{array}\right)+\stackrel{~}{m}_3\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right),`$ which we called the democratic-type mass matrix. (c) The dimension five effective Lagrangian with the $`Z_3`$ symmetry The democratic-type mass matrix is derived from the Lagranfian by imposing the $`Z_3`$ symmetry. We define the irreducible representations of $`Z_3`$ as $`\mathrm{\Psi }_1=\frac{1}{\sqrt{3}}(\mathrm{}_e+\omega ^2\mathrm{}_\mu +\omega \mathrm{}_\tau )`$, $`\mathrm{\Psi }_2=\frac{1}{\sqrt{3}}(\mathrm{}_e+\omega \mathrm{}_\mu +\omega ^2\mathrm{}_\tau )`$ and $`\mathrm{\Psi }_3=\frac{1}{\sqrt{3}}(\mathrm{}_e+\mathrm{}_\mu +\mathrm{}_\tau )`$, where $`\mathrm{}_i`$ is the left-handed lepton doublet defined by, say, $`\mathrm{}_e^T=(\nu _{eL},e_L)`$. The fields $`\mathrm{\Psi }_i`$ transform under the permutation $`\mathrm{}_e`$, $`\mathrm{}_\mu `$ and $`\mathrm{}_\tau `$, $`\mathrm{\Psi }_1\omega \mathrm{\Psi }_1`$, $`\mathrm{\Psi }_2\omega ^2\mathrm{\Psi }_2`$ and $`\mathrm{\Psi }_3\mathrm{\Psi }_3`$. With the definition, $`\mathrm{\Psi }_i=(\psi _i,e_i)`$, the above relation is viewed as the transformation from the flavor eigenstate basis to the $`\psi `$ basis. If we introduce two up-type Higgs doublets that behave as $`H_{u1}\omega ^2H_{u1}`$ and $`H_{u2}\omega H_{u2}`$, then we can construct the $`Z_3`$ invariant dimension five effective Lagrangian as $`_\mathrm{y}`$ $`=`$ $`((m_1^0+\stackrel{~}{m}_1)\overline{(\mathrm{\Psi }_1)^C}\mathrm{\Psi }_1{\displaystyle \frac{H_{u1}H_{u1}}{u_1^2}}+(m_2^0+\stackrel{~}{m}_2)\overline{(\mathrm{\Psi }_2)^C}\mathrm{\Psi }_2{\displaystyle \frac{H_{u2}H_{u2}}{u_2^2}}`$ (6) $`+(m_3^0+\stackrel{~}{m}_3)\overline{(\mathrm{\Psi }_3)^C}\mathrm{\Psi }_3{\displaystyle \frac{H_{u1}H_{u2}}{u_1u_2}})`$ $`2\left(\stackrel{~}{m}_1\overline{(\mathrm{\Psi }_2)^C}\mathrm{\Psi }_3{\displaystyle \frac{H_{u1}H_{u1}}{u_1^2}}+\stackrel{~}{m}_2\overline{(\mathrm{\Psi }_1)^C}\mathrm{\Psi }_3{\displaystyle \frac{H_{u2}H_{u2}}{u_2^2}}+\stackrel{~}{m}_3\overline{(\mathrm{\Psi }_1)^C}\mathrm{\Psi }_2{\displaystyle \frac{H_{u1}H_{u2}}{u_1u_2}}\right),`$ where $`u_i`$ is the vacuum expectation value of the neutral component of $`H_{ui}`$. After the Higgs fields acquire the vacuum expectation values, the neutrino mass matrix in the $`\psi `$-basis defined by Eq.(4), and thus the democratic-type neutrino mass matrix in Eq.(5) is obtained. ## 3 A restricted model -one Higgs case- Here we consider the case with only one up-type Higgs by keeping $`H_{u1}`$. Then, from Eq.(6), we only have the $`\stackrel{~}{m}_1`$ and $`m_1^0`$ terms. As we can see from Eq.(4) with $`\stackrel{~}{m}_2=\stackrel{~}{m}_2=0`$ and $`m_2^0=m_3^0=0`$, we observe that the matrix is diagonalized by $`V_1`$ $`=`$ $`V_T\left(\begin{array}{ccc}1& 0& 0\\ 0& \frac{1}{\sqrt{2}}& \frac{1}{\sqrt{2}}\\ 0& \frac{1}{\sqrt{2}}& \frac{1}{\sqrt{2}}\end{array}\right)=\left(\begin{array}{ccc}1& 0& 0\\ 0& \omega & 0\\ 0& 0& \omega ^2\end{array}\right)\left(\begin{array}{ccc}\sqrt{\frac{1}{3}}& \sqrt{\frac{2}{3}}& 0\\ \frac{1}{\sqrt{3}}& \frac{1}{\sqrt{6}}& \frac{1}{\sqrt{2}}\\ \frac{1}{\sqrt{3}}& \frac{1}{\sqrt{6}}& \frac{1}{\sqrt{2}}\end{array}\right)\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& i\end{array}\right).`$ This mixing matrix is a new type and predicts $`\mathrm{sin}^22\theta _{sol}=\frac{8}{9}`$ and $`\mathrm{sin}^22\theta _{atm}=1`$ in contrast to the Bi-maximal mixing which predicts $`\mathrm{sin}^22\theta _{sol}=1`$ and $`\mathrm{sin}^22\theta _{atm}=1`$, and the democratic mixing does $`\mathrm{sin}^22\theta _{sol}=1`$ and $`\mathrm{sin}^22\theta _{atm}=8/9`$. Unfortunately, this model predicts the degenerate masses $`m_2=m_3(=\stackrel{~}{m}_1)`$. In order to remedy this deficit, we introduce the symmetry breaking terms which preserve the $`Z_2`$ symmetry $`\mathrm{\Psi }_1\mathrm{\Psi }_1`$. Among the interaction terms in Eq.(14), the $`Z_2`$ symmetry excludes $`\stackrel{~}{m}_2`$ and $`\stackrel{~}{m}_3`$ terms and thus we obtain the neutrino mass matrix including four parameters $`\stackrel{~}{m}_1`$, $`m_1^0`$, $`m_2^0`$ and $`m_3^0`$. This matrix is diagonalized by $`V=\left(\begin{array}{ccc}1& 0& 0\\ 0& \omega & 0\\ 0& 0& \omega ^2\end{array}\right)\left(\begin{array}{ccc}\frac{1}{\sqrt{3}}& \sqrt{\frac{2}{3}}c^{}& i\sqrt{\frac{2}{3}}s^{}\\ \frac{1}{\sqrt{3}}& \frac{1}{\sqrt{6}}(c^{}+i\sqrt{3}s^{})& \frac{1}{\sqrt{6}}(\sqrt{3}c^{}+is^{})\\ \frac{1}{\sqrt{3}}& \frac{1}{\sqrt{6}}(c^{}i\sqrt{3}s^{})& \frac{1}{\sqrt{6}}(\sqrt{3}c^{}is^{})\end{array}\right)\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& i\end{array}\right),`$ (7) where $`s^{}=\mathrm{sin}\theta ^{}`$, $`c^{}=\mathrm{cos}\theta ^{}`$ and $`\mathrm{tan}\theta ^{}=\mathrm{\Delta }_{}/(\stackrel{~}{m}_1+\sqrt{\stackrel{~}{m}_1^2+\mathrm{\Delta }_{}^2})`$. The neutrino masses are given by $`m_1=m_1^0+\stackrel{~}{m}_1`$, $`m_2=m_2^0+\mathrm{\Delta }_{}+\sqrt{\stackrel{~}{m}_1^2+\mathrm{\Delta }_{}^2}`$, $`m_3=m_2^0+\mathrm{\Delta }_{}\sqrt{\stackrel{~}{m}_1^2+\mathrm{\Delta }_{}^2}`$. From the mixing matrix in Eq.(7), we find $`\delta =\frac{\pi }{2}`$ and $`\mathrm{tan}\theta _{12}=\sqrt{23\mathrm{sin}^2\theta _{13}}`$, while $`s_{13}`$ is left a free parameter. Now we have $`\mathrm{sin}^22\theta _{sol}=\frac{8}{9}c^2`$ and $`\mathrm{sin}^22\theta _{atm}=1\frac{4}{9}s^4`$. If we impose the CHOOZ bound, $`\|\sqrt{2/3}s^{}\|<0.16`$, we have $`0.87<\mathrm{sin}^22\theta _{sol}<8/9`$ and $`\mathrm{sin}^22\theta _{atm}>0.999`$. If the future precision experiments on the neutrino mixing show that the atmospheric neutrino mixing is very close to the maximal value and the solar neutrino mixing is large but not maximal, our model would be a very good candidate. Here, the phase matrix $`\mathrm{diag}(1,\omega ,\omega ^2)`$ is absobed into charged lepton phases and the phase matrix $`\mathrm{diag}(1,1,i)`$ represents the CP violating Majorana phase matrix,, which is relevant to the purely lepton number violationg processes such as the neutrinoless double beta decay. As for the CP violation, the Jarlskog parameter is $`{\displaystyle \frac{\|J_{CP}\|_{ourmodel}}{\|J_{CP}\|_{max}}}=\sqrt{6}\left\|s_{13}c_{13}^2\sqrt{1{\displaystyle \frac{3}{2}}s_{13}^2}\right\|<0.37,`$ (8) where $`\|s_{13}\|=\sqrt{2/3}s^{}`$. The reduction rate from the maximum CP violation is solely dependent on the mixing angle $`s_{13}`$. If we use the CHOOZ data $`\|s_{13}\|<0.16`$, we obtain the ratio is smaller than 0.37, which is good enough to be observed in the near future laboratory experiments on the neutrino oscillations.
warning/0005/quant-ph0005122.html	ar5iv	text	# Bayesian Reconstruction of Approximately Periodic Potentials at Finite Temperature ## 1 Introduction A successful application of quantum mechanics to real world systems relies essentially on an adequate reconstruction of the underlying potential, describing the forces governing the system. The reconstruction of potentials or forces from available observational data defines an empirical learning task. It also constitutes a typical example of an inverse problem. Such problems are notoriously ill–defined in the sense of Tikhonov Tikhonov-Arsenin-1977 ; Kirsch-1996 ; Vapnik-1998 ; Honerkamp-1998 . In that case additional a priori information is required to yield a unique and stable solution. A Bayesian framework is especially well suited to include both, observational data and a priori information, in a quite flexible manner. Inverse scattering theory Newton-1989 ; Chadan-Sabatier-1989 ; Chadan-Colton-Paivarinta-Rundell-1997 and inverse spectral theory Gelfand-Levitan-1951 ; Kac-1966 ; Marchenko-1986 ; Zakhariev-Chabanov-1997 are two classical research fields which deal in particular with the reconstruction of potentials from spectral data. Both theories describe the kind of data which are necessary, in addition to a given spectrum, to determine a potential uniquely. In inverse scattering theory these additional data are for example phase shifts, obtained far away from the scatterer. For the bound state problems studied in inverse spectral theory these additional data may consist of a second spectrum obtained for boundary conditions different from those for the first spectrum. The approach of Bayesian Inverse Quantum Mechanics (BIQM) we will refer to in the following is not exclusively designed for spectral data but is able to work with quite arbitrary observational data Lemm-IQS-2000 . It can thus be easily adapted to a large variety of different reconstruction scenarios Lemm-BFT-1999 ; Lemm-TDQ-2000 ; Lemm-IHF-2000 . The basics of a Bayesian framework are summarized in Section 2. Setting up a Bayesian approach for a specific application area requires the definition of two basic probabilistic models. First, a likelihood model is needed giving, for each possible potential, the probability of the observational data. The likelihood model of quantum statistics is discussed in Section 3. Second, a prior model has to be chosen to implement available a priori information. Prior models which are useful for inverse quantum statistics are presented in Section 4. Technically the most convenient prior models are Gaussian processes, presented in Section 4.1. Section 4.2 shows how covariance and mean of a Gaussian process can be related to a priori information about approximate symmetries of the potentials to be reconstructed. Section 4.3 concentrates on approximate periodicity, Section 4.4 on potentials with discontinuities. Prior models are made more flexible by using hyperparameters (Section 4.5), or more general hyperfields, being function hyperparameters (Section 4.6). Related non–Gaussian priors are the topic of Section 4.7. Having defined liklihood and prior models Section 5 discusses the equations to be solved for reconstructing a potential. Finally, Section 6 presents numerical applications. ## 2 Bayesian approach Empirical learning is based on observational data $`D`$. In particular, we will distinguish “dependent” variables $`x`$, representing measurement results, and “independent” variables $`O`$, characterizing the kind of measurement performed. In the context of inverse quantum mechanics the latter denotes the observables which are measured. Such observables may for example be the position, the momentum, or the energy of a quantum particle. Variables $`x`$ and $`O`$ are assumed to be measurable and represent therefore visible variables. Observational data will be assumed to consist of $`n`$ pairs $`D`$ = $`\{(x_i,O_i)\|1in\}`$ = $`(x_T,O_T)`$, where $`x_T`$ and $`O_T`$ denote the vectors with components $`x_i`$ or $`O_i`$, respectively. Such data will also be called training data. In empirical learning one tries to extract a “general law” from observations. In this paper the quantum potential $`V`$ to be reconstructed will represent this “general law”. (Similarly, in the Bayesian reconstruction of quantum states the object to be reconstructed is the density operator of an unknown state Helstrom:1976 ; Holevo:1982 ; Tan:1997 ; Buzek-Drobny-Derka-Adam-Wiedemann:1998 .) Potentials, considered not to be directly observable, represent in our context the hidden or latent variables. We will now use the Bayesian framework to relate unobservable potentials to observational data. The Bayesian approach is a general probabilistic framework to deal with empirical learning problems Bayes-1763 ; Berger-1980 ; Loredo-1990 ; Bernado-Smith-1994 ; Gelman-Carlin-Stern-Rubin-1995 ; Sivia-1996 ; Carlin-Louis-1996 ; Lemm-BFT-1999 . Predicting results of future measurements on the basis of given training data is achieved by means of the predictive probability $`p(x\|O,D)`$ (or predictive density for continuous $`x`$), which is the probability of finding the value $`x`$ when measuring observable $`O`$ under the condition that the training data $`D`$ are given. To calculate the predictive probability a probabilistic model is needed which describes the measurement process. Such a model is specified by giving the probability $`p(x\|O,V)`$ of finding $`x`$ when measuring observable $`O`$ for each possible potential $`V`$. As $`p(x\|O,V)`$, considered as function of $`V`$ for fixed $`x`$ and $`O`$, is known as likelihood of $`V`$, we will call this the likelihood model. For inverse quantum problems the likelihood model is given by the axioms of quantum mechanics and will be discussed in Section 3. According to the rules of probability theory the predictive probability can now be written as an integral over the space of all possible potentials $`V`$, $$p(x\|O,D)=𝑑Vp(x\|O,V)p(V\|D).$$ (1) We note that in Eq.(1) we have assumed that the probability of $`x`$ is completely determined by giving potential and observable and does not depend on the training data, $`p(x\|O,V,D)`$ = $`p(x\|O,V)`$, and that the probability of the potential given the training data does not depend on the observables selected in the future, $`p(V\|O,D)`$ = $`p(V\|D)`$. If the set of possible potentials is a space of functions, the integral in (1) is a functional integral. As the likelihood model is assumed to be given, learning consists in the determination of $`p(V\|D)`$, known as the posterior for $`V`$. To this end, we relate the posterior for $`V`$ to the likelihood of $`V`$ under the training data by applying Bayes’ theorem, $$p(V\|D)=\frac{p(x_T\|O_T,V)p(V)}{p(x_T\|O_T)},$$ (2) assuming $`p(V\|O_T)`$ = $`p(V)`$, analogous to Eq. (1). In the numerator of Eq. (2) appears, besides the likelihood, the so called prior $`p(V)`$. This prior gives the probability of $`V`$ before training data have been collected. Hence it has to comprise all a priori information available for the potential. The need for a prior model, complementing the likelihood model, is characteristic for a Bayesian approach. The denominator in Eq. (2) plays the role of a normalization factor and can be obtained from likelihood and prior by integration over $`V`$ as $`p(x_T\|O_T)`$ = $`𝑑Vp(x_T\|O_T,V)p(V)`$. From a Bayesian perspective learning appears as updating the probability for $`V`$ caused by the arrival of new data $`D`$. If more data become available this process can be iterated, the old posterior becoming the new prior which is then updated yielding a new posterior. In practice, a major difficulty is the calculation of the integral over all possible $`V`$ to get the predictive probability (1). Even if one resorts to a discrete approximation for $`x`$ the integral (1) is typically still very high dimensional. The key point is thus to find a feasible approximation for that integral. Two approaches are common in Bayesian statistics. The first one is an evaluation of the integral by Monte Carlo methods Gelman-Carlin-Stern-Rubin-1995 ; Metropolis-Rosenbluth-Rosenbluth-Teller-Teller-1953 ; Binder-Heermann-1988 ; Neal-1997 . The second one, which we will pursue in the following, is the so called maximum a posteriori approximation (MAP), being a variant of the saddle point method Berger-1980 ; Gelman-Carlin-Stern-Rubin-1995 ; De-Bruijn-1981 ; Bleistein-Handelsman-1986 ; Girosi-Jones-Poggio-1995 ; Lemm-1996 ; Lemm-1998 . In MAP one assumes the posterior to be sufficiently peaked around the potential $`V^{}`$ which maximizes the posterior, so that approximately $$p(x\|O,D)p(x\|O,V^{}),$$ (3) with $$V^{}=\mathrm{argmax}_{V𝒱}p(V\|D)=\mathrm{argmax}_{V𝒱}p(x_T\|O_T,V)p(V).$$ (4) Maximizing the posterior with respect to $`V𝒱`$ means, according to Eq. (2) with the denominator independent of $`V`$, maximizing the product of likelihood and prior. The Bayesian framework discussed so far can analogously be applied to a variety of different contexts, including regression, density estimation and classification problems Lemm-BFT-1999 . The case of a Gaussian likelihood with fixed variance, for example, is known as regression problem, while problems with general likelihoods are known as density estimation. ## 3 Likelihood model of quantum statistics The first step in applying the Bayesian framework to inverse problems of quantum mechanics or quantum statistics is the definition of the likelihood model Lemm-IQS-2000 . This is easily obtained from the axioms of quantum mechanics. Consider a system prepared in a state described by a density operator $`\rho `$. As our aim will be to reconstruct potentials $`V`$ from observational data, we have to choose a $`\rho `$ which depends on the potential. The probability to find value $`x`$, when measuring an observable represented by the Hermitian operator $`O`$, is given by $$p(x\|O,V)=\mathrm{Tr}\left(P_O(x)\rho (V)\right),$$ (5) where $`P_O(x)`$ = $`_\zeta \|x,\zeta x,\zeta \|`$ denotes the projector on the space of (orthonormalized) eigenfunctions $`\|x,\zeta `$ of $`O`$ with eigenvalue $`x`$ and the variable $`\zeta `$ distinguishes eigenfunctions with degenerate eigenvalues. In particular, for a canonical ensemble at temperature $`1/\beta `$ (setting Boltzmann’s constant equal to 1) the density operator reads $$\rho =\frac{e^{\beta H}}{\mathrm{Tr}e^{\beta H}}.$$ (6) To be specific, we will study in the following Hamiltonians of the form $`H`$ = $`T+V`$, with kinetic energy $`T`$ = $`(1/2m)\mathrm{\Delta }`$, (with Laplacian $`\mathrm{\Delta }`$, mass $`m`$, and setting $`\mathrm{}`$ = $`1`$) and a local potential $$V(x,x^{})=v(x)\delta (xx^{}),$$ (7) defined by the function $`v(x)`$. Note that the formalism presented in the following works with nonlocal potentials as well, numerical calculations, however, would in that case be more demanding. For the likelihood models corresponding to time–dependent quantum systems and to many–body systems in Hartree–Fock approximation we refer to Lemm-TDQ-2000 ; Lemm-IHF-2000 . In the following we will study observational data consisting of $`n`$ position measurements $`x_i`$. This corresponds to choosing the position operator for the observables $`O_i`$ = $`\widehat{x}`$ with $`\widehat{x}\|x_i`$ = $`x_i\|x_i`$. Hence, for a canonical ensemble, the likelihood (5) becomes for a single position measurement $$p(x_i\|\widehat{x},v)=\underset{\alpha }{}p_\alpha \|\varphi _\alpha (x_i)\|^2=\|\varphi (x_i)\|^2$$ (8) with (non–degenerate) eigenfunctions $`\varphi _\alpha `$ of $`H`$ and energies $`E_\alpha `$, i.e., $`H\|\varphi _\alpha `$ = $`E_\alpha \|\varphi _\alpha `$. Angular brackets $`\mathrm{}`$ denote a thermal expectation under the probabilities $`p_\alpha `$ = $`\mathrm{exp}(\beta E_\alpha )/Z`$ with $`Z`$ = $`_\alpha \mathrm{exp}(\beta E_\alpha )`$ according to Eq. (6). For independent data $`D_i`$ = $`(x_i,O_i)`$, $$p(x_T\|O_T,v)=\underset{i=1}{\overset{n}{}}p(x_i\|\widehat{x},v)=\underset{i=1}{\overset{n}{}}\|\varphi (x_i)\|^2.$$ (9) A quantum mechanical measurement changes the state of the system, i.e., it changes $`\rho `$. Hence, to obtain independent data under constant $`\rho `$ requires the density operator to be restored before each measurement. For a canonical ensemble this means to wait between two consecutive observations until the system is thermalized again. Choosing a parametric family of potentials $`v(x;\xi )`$ one could now maximize the likelihood with respect to the parameters $`\xi `$, and choose as reconstructed potential $$v^{}(x)=v(x;\xi ^{})\text{with}\xi ^{}=\text{argmax}_\xi p(x_T\|O_T,v(\xi )).$$ (10) This is known as maximum likelihood approximation and works well if the number of data is large compared to the flexibility of the selected parametric family of potentials. This method does however not yield a unique optimal potential if the flexibility is too large for the available number of observations. (A possible measure of the “flexibility” of a parametric family is given by the Vapnik-Chervonenkis dimension Vapnik-1998 or variants thereof.) In such cases, the inclusion of additional restrictions on $`v`$ in form of a priori information is essential. This holds especially for nonparametric approaches, where each number $`v(x)`$ is treated as individual degree of freedom. Including a priori information generalizes the maximum likelihood approximation of Eq. (10) to the MAP of Eq. (4). ## 4 Prior models ### 4.1 Gaussian processes A finite number of observational data cannot completely determine a function $`v(x)`$. Hence, besides observational data, additional a priori information is necessary to reconstruct a potential in BIQM. In nonparametric approaches it is advantageous to formulate a priori information directly in terms of the function $`v(x)`$ itself. A convenient choice for a prior is a Gaussian process, $$p(v)=\left(det\frac{𝐊_0}{2\pi }\right)^{\frac{1}{2}}e^{\frac{1}{2}vv_0\|𝐊_0\|vv_0},$$ (11) where $$vv_0\|𝐊_0\|vv_0=$$ (12) $$𝑑x𝑑x^{}[v(x)v_0(x)]𝐊_0(x,x^{})[v(x^{})v_0(x^{})].$$ The function $`v_0`$ is the mean or regression function, representing a reference potential or template for $`v`$. The inverse covariance $`𝐊_0`$ is a real symmetric, positive (semi)definite operator which acts on potentials rather than on wave functions and defines a distance measure on the space of potentials. For technical convenience one may introduce explicitly a factor $`\lambda `$ multiplying $`𝐊_0`$ to balance the influence of the prior against the likelihood term. A Gaussian prior as in Eq. (11) is already a quite flexible tool for implementing a priori knowledge. A bias towards smooth functions $`v(x)`$, for instance, can be implemented by choosing the negative Laplacian as inverse covariance $`𝐊_0`$ = $`\mathrm{\Delta }`$. Including higher derivatives in $`𝐊_0`$ would result in even smoother potentials, in the sense that higher derivatives of $`v(x)`$ become continuous. For example, a common smoothness prior used for regression problems is the Radial Basis Function prior $`𝐊_0`$ = $`\mathrm{exp}(\sigma _{\mathrm{RBF}}^2\mathrm{\Delta }/2)`$ Girosi-Jones-Poggio-1995 . ### 4.2 Covariances and approximate symmetries Prior information on potentials $`v`$ can often be related to approximate invariance under specific transformations Lemm-BFT-1999 . Typical examples of such transformations are symmetry operations like translations or rotations. To be specific, assume that a (not necessarily local) potential $`V`$ commutes approximately, but not exactly, with some unitary operator $`S`$, $$VS^{}VS=𝐒V,$$ (13) which defines an operator $`𝐒`$ acting on $`V`$. In particular, we may choose a prior $`p(V)\mathrm{exp}\{E_S(V)\}`$ with a prior energy $$E_S=\frac{1}{2}V𝐒V\|V𝐒V=\frac{1}{2}V\|𝐊_0\|V.$$ (14) This shows that the expectation of an approximate symmetry of $`V`$ under $`S`$ can be implemented by choosing a Gaussian prior with inverse covariance operator $$𝐊_0=(𝐈𝐒)^{}(𝐈𝐒),$$ (15) where $`𝐈`$ denotes the identity operator. Symmetry operations $`S(\theta )`$, with corresponding $`𝐒(\theta )`$, may depend on a parameter (vector) $`\theta `$. Approximate invariance under $`S(\theta _i)`$ for several $`\theta _i`$ can be implemented by using the sum (or integral, for continuous variables) $`E_S`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}V𝐒(\theta _i)V\|V𝐒(\theta _i)V`$ (16) $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}V\|𝐊_0(\theta _i)\|V.`$ Alternatively, one may require approximate symmetry for only one value of $`\theta `$, not fixed a priori. For example, one may expect an approximately periodic potential with unknown periodicity length $`\theta `$ which also has to be determined from the data. Such $`\theta `$ are known as hyperparameters and will be discussed in Section 4.5. Lie groups are continuously parameterized transformations $$𝐒(\theta )=e^{_i\theta _i𝐬_i},$$ (17) where $`\theta _i`$ are the real parameters and the $`𝐬_i`$ = $`𝐬_i^T`$ (the superscript <sup>T</sup> denoting the transpose) are antisymmetric operators representing the generators of the infinitesimal transformations of the Lie–group. We can define a prior energy as an error measure with respect to an infinitesimal transformation, $`E_S`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}{\displaystyle \frac{V(1+\theta _i𝐬_i)V}{\theta _i}}\|{\displaystyle \frac{V(1+\theta _i𝐬_i)V}{\theta _i}}`$ (18) $`=`$ $`{\displaystyle \frac{1}{2}}V\|{\displaystyle \underset{i}{}}𝐬_i^T𝐬_i\|V.`$ For instance, a Laplacian smoothness prior for a local potential $`v(x)`$ can be related to an approximate symmetry under infinitesimal translations. For the group of $`d`$–dimensional translations which is generated by the gradient operator $``$ this can be verified by recalling the multidimensional Taylor formula for expanding $`v`$ around $`x`$ $$𝐒(\theta )v(x)=e^{_i\theta _i_i}v(x)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{\left(_i\theta _i_i\right)^k}{k!}v(x)=v(x+\theta ).$$ (19) Up to first order $`𝐒1+_i\theta _i_i`$. Hence, for infinitesimal translations, the error measure of Eq. (18) becomes $`E_S`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}{\displaystyle \frac{v(1+\theta _i_i)v}{\theta _i}}\|{\displaystyle \frac{v(1+\theta _i_i)v}{\theta _i}}`$ (20) $`=`$ $`{\displaystyle \frac{1}{2}}v\|\mathrm{\Delta }\|v,`$ assuming vanishing boundary terms. This is the classical Laplacian smoothness term. ### 4.3 Approximate periodicity In this paper we will in particular be interested in potentials which are approximately periodic. To measure the deviation from exact periodicity for a local potential $`v(x)`$ let us define the difference operators $`\left(_\theta ^Rv\right)(x)`$ $`=`$ $`v(x+\theta )v(x).`$ (21) $`\left(_\theta ^Lv\right)(x)`$ $`=`$ $`v(x)v(x\theta ),`$ (22) For periodic boundary conditions $`(_\theta ^L)^T`$ = $`_\theta ^R`$, where $`(_\theta ^L)^T`$ denotes the transpose of $`_\theta ^L`$. Hence, the operator $$\mathrm{\Delta }_\theta =_\theta ^L_\theta ^R=(_\theta ^R)^T_\theta ^R$$ (23) defined in analogy to the negative Laplacian, is positive (semi)definite, and a possible prior energy is an error term which measures the deviation from exact periodicity for given period $`\theta `$, $`E_S`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle 𝑑x\|v(x)v(x+\theta )\|^2}`$ (24) $`=`$ $`{\displaystyle \frac{1}{2}}_\theta ^Rv\|_\theta ^Rv`$ $`=`$ $`{\displaystyle \frac{1}{2}}v\|\mathrm{\Delta }_\theta \|v.`$ Discretizing $`v`$ the operator $`_\theta ^R`$ for periodic boundary conditions becomes, for example on a mesh with six points and $`\theta `$ = $`2`$, the matrix $$_\theta ^R=\left(\begin{array}{cccccc}1& \text{0}& 1& \text{0}& \text{0}& \text{0}\\ \text{0}& 1& \text{0}& 1& \text{0}& \text{0}\\ \text{0}& \text{0}& 1& \text{0}& 1& \text{0}\\ \text{0}& \text{0}& \text{0}& 1& \text{0}& 1\\ 1& \text{0}& \text{0}& \text{0}& 1& \text{0}\\ \text{0}& 1& \text{0}& \text{0}& \text{0}& 1\end{array}\right),$$ (25) so that $$\mathrm{\Delta }_\theta =\left(\begin{array}{cccccc}\text{2}& \text{0}& 1& \text{0}& 1& \text{0}\\ \text{0}& \text{2}& \text{0}& 1& \text{0}& 1\\ 1& \text{0}& \text{2}& \text{0}& 1& \text{0}\\ \text{0}& 1& \text{0}& \text{2}& \text{0}& 1\\ 1& \text{0}& 1& \text{0}& \text{2}& \text{0}\\ \text{0}& 1& \text{0}& 1& \text{0}& \text{2}\end{array}\right).$$ (26) As every periodic function with $`v(x)=v(x+\theta )`$ is in the null space of $`\mathrm{\Delta }_\theta `$ typically another error term has to be added to get a unique maximum of the posterior. For example, combining a prior energy (24) with a Laplacian smoothness term yields a Gaussian prior of the form (11) with inverse covariance $`𝐊_0`$ = $`\lambda (\mathrm{\Delta }+\gamma \mathrm{\Delta }_\theta )`$ and prior energy $$E_S=\frac{\lambda }{2}v\|\mathrm{\Delta }+\gamma \mathrm{\Delta }_\theta \|v,$$ (27) with weighting factors $`\lambda `$, $`\gamma `$. In case the period $`\theta `$ is not known, it can be treated as hyperparameter as will be discussed in Section 4.5. Clearly, a nonzero reference potential $`v_0`$ can be included in Eq. (27). In Eq. (24), one may also sum over several periods $$E_S=\frac{1}{2}\underset{k=1}{\overset{k_{\mathrm{max}}}{}}w(k)𝑑x\|v(x)v(x+k\theta )\|^2,$$ (28) where $`w(k)`$ is a weighting function, decreasing for larger $`k`$. Prior energies as in (28) enforce approximate periodicity over longer distances than a prior energy of the form (24). The latter, on the other hand, is more robust than (28) with respect to local deviations from periodicity, like a locally varying frequency. Instead of choosing an inverse covariance $`𝐊_0`$ with symmetric functions in its null space, approximate symmetries can be implemented by using explicitly a symmetric reference function $`v_0`$ = $`𝐒v_0`$ for the Gaussian prior (11). For approximate periodicity, this would mean to choose a periodic reference potential $`v_0(x)`$ = $`v_0(x+\theta )`$ in the prior energy $`E_S=\frac{1}{2}vv_0\|𝐊_0\|vv_0`$ where $`𝐊_0`$ could be for example the identity or a differential operator. Thus a periodic reference potential favors a specific form for the reconstructed potential, including a specific frequency and phase. This is different for the covariance implementation (24) of approximate periodicity where only the frequency is relevant and reference potentials can still be chosen arbitrarily. They may, for example be nonperiodic functions or functions with even higher symmetry like in Eq. (27) where $`v_00`$ is invariant under all translations. Flexible reference potentials will be studied in Section 4.5. ### 4.4 Potentials with discontinuities Smooth potentials $`v(x)`$ with discontinuities can either be approximated by using discontinuous templates $`v_0(x;\theta )`$ or by eliminating matrix elements of the inverse covariance which connect the two sides of the discontinuity. For example, consider the discrete version of a negative Laplacian with unit lattice spacing and periodic boundary conditions, $$𝐊_0=\mathrm{\Delta }=\left(\begin{array}{cccccc}\text{2}& 1& \text{0}& \text{0}& \text{0}& 1\\ 1& \text{2}& 1& \text{0}& \text{0}& \text{0}\\ \text{0}& 1& \text{2}& 1& \text{0}& \text{0}\\ \text{0}& \text{0}& 1& \text{2}& 1& \text{0}\\ \text{0}& \text{0}& \text{0}& 1& \text{2}& 1\\ 1& \text{0}& \text{0}& \text{0}& 1& \text{2}\end{array}\right).$$ (29) Decomposing the matrix (29) into square roots we write $`𝐊_0`$ = $`𝐖^T𝐖`$ (see also Section 4.6) where a possible square root is $$𝐖=_1^R=\left(\begin{array}{cccccc}1& 1& \text{0}& \text{0}& \text{0}& \text{0}\\ \text{0}& 1& 1& \text{0}& \text{0}& \text{0}\\ \text{0}& \text{0}& 1& 1& \text{0}& \text{0}\\ \text{0}& \text{0}& \text{0}& 1& 1& \text{0}\\ \text{0}& \text{0}& \text{0}& \text{0}& 1& 1\\ 1& \text{0}& \text{0}& \text{0}& \text{0}& 1\end{array}\right).$$ (30) Similarly, the derivative operator $`/x`$ represents a square root of the negative Laplacian for periodic boundary conditions. Two regions can now be disconnected by deleting all lines of $`𝐖`$ which have matrix elements in both regions. For instance, the first three points in the six–dimensional space of Eq. (30) can be disconnected from the last three points by setting $`𝐖(3,)`$ and $`𝐖(6,)`$ to zero, $$\stackrel{~}{𝐖}=\left(\begin{array}{c}\text{−1}\text{1}\text{0}\text{0}\text{0}\text{0}\\ \text{0}& 1& 1& \text{0}& \text{0}& \text{0}\\ \text{0}& \text{0}& \text{0}& \text{0}& \text{0}& \text{0}\\ \text{0}& \text{0}& \text{0}& 1& 1& \text{0}\\ \text{0}& \text{0}& \text{0}& \text{0}& 1& 1\\ \text{0}& \text{0}& \text{0}& \text{0}& \text{0}& \text{0}\end{array}\right).$$ (31) Squaring of $`\stackrel{~}{𝐖}`$ yields a positive semidefinite operator $$\stackrel{~}{𝐊}_0=\stackrel{~}{𝐖}^T\stackrel{~}{𝐖}=\left(\begin{array}{c}\text{1}\text{−1}\text{0}\text{0}\text{0}\text{0}\\ 1& \text{2}& 1& \text{0}& \text{0}& \text{0}\\ \text{0}& 1& \text{1}& \text{0}& \text{0}& \text{0}\\ \text{0}& \text{0}& \text{0}& \text{1}& 1& \text{0}\\ \text{0}& \text{0}& \text{0}& 1& \text{2}& 1\\ \text{0}& \text{0}& \text{0}& \text{0}& 1& \text{1}\end{array}\right)$$ (32) resulting in a smoothness prior which is ineffective between points from different regions. In contrast to using discontinuous templates, the height of the jump at the discontinuity has not to be given in advance when working with disconnected Laplacians (or other disconnected inverse covariances). On the other hand training data are then required for all separated regions to determine the free constants which correspond to the zero modes of the local Laplacians. The reconstruction of discontinuous functions with non–Gaussian priors will be discussed in Section 4.7. ### 4.5 Hyperparameters Parameters of the prior are known as hyperparameters Lemm-BFT-1999 ; Carlin-Louis-1996 ; Bishop-1995b . Like potentials $`v`$, hyperparameters $`\theta `$ are not directly observable and represent hidden variables. In the presence of hyperparameters a prior for $`v`$ can be decomposed as follows $$p(v)=𝑑\theta p(v\|\theta )p(\theta ),$$ (33) where $`p(\theta )`$ is known as hyperprior. The likelihood does not depend on $`\theta `$, the predictive probability (1), however, contains then an integral over $`\theta `$, $$p(x\|O,D)=$$ (34) $$\frac{1}{p(x_T\|O_T)}𝑑v𝑑\theta p(x\|O,v)p(x_T\|O_T,v)p(v\|\theta )p(\theta ).$$ Like the integral over $`v`$, the integral over $`\theta `$ can be calculated either by Monte Carlo methods or in MAP. We remark that, when a $`\theta `$–dependent prior is written in terms of a corresponding prior energy $`p(v\|\theta )e^{E(v\|\theta )}`$, the normalization $`𝑑ve^{E(v\|\theta )}`$ is independent of $`v`$ but does in general depend on $`\theta `$. Hyperparameters $`\theta `$ can be single numbers or vectors. They can describe continuous transformations, like translation, rotation or scaling of template functions and scaling of inverse covariance operators. For real $`\theta `$ and differentiable posterior, stationarity conditions can be found by differentiating the posterior with respect to $`\theta `$. Instead of continuous transformations of templates or inverse covariances one can consider a finite collection of alternative reference potentials $`v_i`$ or alternative inverse covariances $`𝐊_i`$. For example, a potential to be reconstructed may be expected to be similar to one reference potential out of a small number of possible alternatives $`v_i`$. The “class” variables $`i`$ are then nothing else but hyperparameters $`\theta `$ with integer values. Binary parameters allow to select from two reference functions or two inverse covariances that one which fits the data best. Indeed, writing $`v_0(\theta )`$ $`=`$ $`(1\theta )v_1+\theta v_2,`$ (35) $`𝐊_0(\theta )`$ $`=`$ $`(1\theta )𝐊_1+\theta 𝐊_2,`$ (36) a binary $`\theta \{0,1\}`$ implements hard switching between alternative templates or inverse covariances, corresponding to a conditional prior $$p(v\|\theta )e^{(1\theta )E_1(v)\theta E_2(v)}$$ (37) with $`E_1(v)`$ $`=`$ $`{\displaystyle \frac{1}{2}}vv_1\|𝐊_1\|vv_1,`$ (38) $`E_2(v)`$ $`=`$ $`{\displaystyle \frac{1}{2}}vv_2\|𝐊_2\|vv_2.`$ (39) Similarly, a real $`\theta [0,1]`$ in (35) or (36) yields soft mixing. In that case, however, the mixing of templates in (35) is not equivalent to a mixing of prior energies as in (37) because for real $`\theta `$ Eqs. (35) and (36) lead to mixed terms, like $`(1\theta )\theta vv_1\|𝐊_0\|vv_2/2`$ for $`𝐊_1`$ = $`𝐊_2`$. When $`\theta `$ takes integer values the integral $`𝑑\theta `$ becomes a sum $`_\theta `$ so that prior, posterior, and predictive probability have the form of a finite mixture with components $`\theta `$ lemm-mixture-1999 . For a moderate number of components one may be able to include all of the mixture components in the calculations. If the number of mixture components is too large one must select some of the components, for example by creating a random sample using Monte Carlo methods, or by solving for the $`\theta ^{}`$ with maximal posterior. In contrast to typical optimization problems for real variables, the corresponding integer optimization problems are usually not very smooth with respect to $`\theta `$ (with smoothness defined in terms of differences instead of derivatives), and are therefore often much harder to solve. There exists a variety of deterministic and stochastic integer optimization algorithms, which may be combined with ensemble methods like genetic algorithms Holland-1975 ; Goldberg-1989 ; Michalewicz-1992 ; Schwefel-1995 ; Mitchell-1996 , and with homotopy methods like simulated annealing Kirkpatrick-Gelatt-Vecchi-1983 ; Mezard-Parisi-Virasoro-1987 ; Aarts-Korts-1989 ; Gelfand-Mitter-1991 ; Yuille-Kosowski-1994 . Annealing methods are similar to (Markov chain) Monte Carlo methods, which aim at sampling many points from a specific distribution (i.e., for example at fixed temperature). For Monte Carlo methods it is important to have (nearly) independent samples and the correct limiting distribution for the Markov chain. For annealing methods the aim is to find the correct minimum by smoothly changing the temperature from a finite value to zero. For the latter it is thus less important to model the distribution for nonzero temperatures exactly, but it is important to use an adequate cooling scheme for lowering the temperature. ### 4.6 Hyperfields The hyperparameters $`\theta `$ considered so far have been real or integer numbers, or vectors with real or integer components $`\theta _i`$. In this section we will discuss priors parameterized by functions, called hyperfields Lemm-BFT-1999 , resulting in a still larger flexibility of the formalism. In numerical calculations where functions have to be discretized hyperfields stand for high dimensional hyperparameter vectors. Using hyperfields one has to keep in mind that a gain in flexibility at the same time tends to lower the influence of the prior. For example, consider as hyperfield a completely adaptive reference potential $`\theta (x)`$ = $`v_0(x)`$ within a Gaussian prior (11). Then, for any $`v(x)`$ the prior energy vanishes for $`v_0(x)`$ = $`v(x)`$. In the absence of additional hyperpriors $`p(\theta )`$ the corresponding MAP solution for the hyperfield $`\theta (x)`$ = $`v_0(x)`$ is thus $`\theta ^{}(x)`$ = $`v(x)`$ for which the Gaussian prior (11) becomes uniform in $`v(x)`$. Hence the price to be paid for the additional flexibility introduced by hyperfields are weaker priors and a large number of additional degrees of freedom. This can considerably complicate calculations and requires sufficiently restrictive hyperpriors for the hyperfields. Let us define local hyperfields $`\theta (x)`$ to be hyperfields depending on the position variable $`x`$. (In general hyperfields can be introduced which depend on other real variables or on several position variables.) Local hyperfields can be used, for example, to adapt templates or inverse covariances locally. To this end, we express real symmetric, positive (semi)definite inverse covariances by square roots or (real) filter operators $`𝐖`$, so that $$𝐊_0=𝐖^T𝐖.$$ (40) In components $$𝐊_0(x,x^{})=𝑑x^{\prime \prime }𝐖^T(x,x^{\prime \prime })𝐖(x^{\prime \prime },x^{}),$$ (41) and therefore $`vv_0\|𝐊_0\|vv_0`$ $`=`$ $`{\displaystyle 𝑑x𝑑x^{}𝑑x^{\prime \prime }[v(x)v_0(x)]}`$ (42) $`\times 𝐖^T(x,x^{})𝐖(x^{},x^{\prime \prime })`$ $`\times [v(x^{\prime \prime })v_0(x^{\prime \prime })]`$ $`=`$ $`{\displaystyle 𝑑x\|\omega (x)\|^2},`$ where we define the filtered difference $$\omega (x)=𝑑x^{}𝐖(x,x^{})[v(x^{})v_0(x^{})].$$ (43) For instance, a square root (30) of the discrete negative Laplacian (29) corresponds for $`v_00`$ to a filtered difference $`\omega (x)`$ = $`v(x+1)v(x)`$. The exponent of a Gaussian prior for a local potential $`v`$ can thus be written as an integral over $`x`$, $$p(v)e^{E(v)};E(v)=\frac{1}{2}𝑑x\|\omega (x)\|^2.$$ (44) In contrast to Eqs. (35) and (36) the representation (44) is well suited for introducing local hyperfields. For instance, an adaptive prior $$p(v\|\theta )=e^{E(v\|\theta )},$$ (45) with a real local hyperfield $`\theta (x)[0,1]`$ can be obtained by mixing locally two alternative filtered differences $$\omega (x;\theta )=[1\theta (x)]\omega _1(x)+\theta (x)\omega _2(x),$$ (46) where the two $`\omega _i`$ may differ in their filters and/or reference potentials. In that case the hyperfield $`\theta (x)`$ can locally select the best mixture of the filtered differences $`\omega _i`$, i.e., that one which yields in (45) the largest probability or smallest prior energy $`E(v\|\theta )`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle 𝑑x\|\omega (x;\theta )\|^2}+\mathrm{ln}Z_𝒱(\theta )`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle 𝑑x\left\|[1\theta (x)]\omega _1(x)+\theta (x)\omega _2(x)\right\|^2}+\mathrm{ln}Z_𝒱(\theta ).`$ Here the normalization factor $$Z_𝒱(\theta )=_{v𝒱}𝑑ve^{\frac{1}{2}{\scriptscriptstyle 𝑑x\|\omega (x;\theta )\|^2}},$$ (48) depends in general on $`\theta `$ if the filters of the $`\omega _i`$ differ. Clearly, allowing an unbounded $`\mathrm{}\theta (x)\mathrm{}`$ any function $`\omega (x;\theta )`$ can be written in the form of Eq. (46), provided $`\omega _1(x)\omega _2(x)`$ for all $`x`$. In contrast to soft mixing with real functions $`\theta (x)`$ a binary local hyperfield $`\theta (x)\{0,1\}`$ implements hard switching between alternative filtered differences. Since in the binary case $`\theta ^2`$ = $`\theta `$, $`(1\theta )^2`$ = $`(1\theta )`$, and $`\theta (1\theta )`$ = $`0`$, Eq. (4.6) becomes \[compare Eq. (37)\] $`E(v\|\theta )`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle }dx([1\theta (x)]\|\omega _1(x)\|^2`$ (49) $`+\theta (x)\|\omega _2(x)\|^2)+\mathrm{ln}Z_𝒱(\theta ),`$ while for real $`\theta (x)`$ Eq. (4.6) includes a mixed term in $`\omega _1\omega _2`$. It is sometimes helpful to transform an unrestricted real hyperfield $`\mathrm{}g(x)\mathrm{}`$ into a bounded real hyperfield $`\theta (x)[0,1]`$ by $$\theta (x)=\sigma (g(x)\vartheta ),$$ (50) with threshold $`\vartheta `$ and sigmoidal transformation $$\sigma (x)=\frac{1}{1+e^{2\nu x}}=\frac{1}{2}(\mathrm{tanh}(\nu x)+1).$$ (51) In the limit $`\nu \mathrm{}`$ the transformation $`\sigma (x)`$ of (51) approaches the step function $`\mathrm{\Theta }(x)`$ and (50) results in a binary $`\theta (x)`$ = $`\mathrm{\Theta }(g(x)\vartheta )\{0,1\}`$. Analogous to the global mixing or global switching in Eq. (35) and Eq. (36), the alternative filtered differences $`\omega _i(x)`$ at position $`x`$ in Eq. (46) can be constructed by local mixing or switching between template functions $`v_1(x^{})`$, $`v_2(x^{})`$ or filters $`𝐖_1(x,x^{})`$, $`𝐖_2(x,x^{})`$ using a local hyperfield $`\theta (x)`$, $`v_x(x^{};\theta )`$ $`=`$ $`[1\theta (x)]v_1(x^{})+\theta (x)v_2(x^{}),`$ (52) $`𝐖(x,x^{};\theta )`$ $`=`$ $`[1\theta (x)]𝐖_1(x,x^{})+\theta (x)𝐖_2(x,x^{}).`$ (53) It is important to note that the local templates or reference potentials $`v_x(x^{};\theta )`$ are functions of $`x^{}`$ and $`x`$. Indeed, to obtain a filtered difference $`\omega (x;\theta )`$ at position $`x`$, a reference function $`v_x`$ is needed for all $`x^{}`$ for which the corresponding $`𝐖(x,x^{})`$ is nonzero, since $$\omega (x;\theta )=𝑑x^{}𝐖(x,x^{})[v(x^{})v_x(x^{};\theta )].$$ (54) In this way the whole template function $`v_x(x^{};\theta )`$, rather than individual function values $`v_0(x,\theta )`$, is adapted individually for every local filtered difference. In particular, the local reference potentials of Eq. (52) have to be distinguished from one global, locally adapted reference potential $$v_0(x^{};\theta )=[1\theta (x^{})]v_1(x^{})+\theta (x^{})v_2(x^{}),$$ (55) which at first glance seems to be the natural generalization of Eq. (35) to local hyperfields. Only in Gaussian prior terms with the identity $`𝐈`$ as covariance, local template functions $`v_x(x^{},\theta )`$ are not required. In that case $`v_x(x^{};\theta )`$ is only needed for $`x`$ = $`x^{}`$ and we may directly write $`v_x(x^{};\theta )`$ = $`\stackrel{~}{v}_0(x^{};\theta )`$, skipping the variable $`x`$, and obtain the prior energy $$\frac{1}{2}𝑑x\|\omega (x;\theta )\|^2=\frac{1}{2}v\stackrel{~}{v}_0(\theta )\|v\stackrel{~}{v}_0(\theta ).$$ (56) We remark that one can also generalize Eq. (52), which uses the same $`v_1(x^{})`$, $`v_2(x^{})`$ for all $`x`$, by working with reference potentials $`v_{1,x}(x^{})`$, $`v_{2,x}(x^{})`$ which vary with the position $`x`$ at which the filtered difference $`\omega (x)`$ is required. This yields $$v_x(x^{};\theta )=[1\theta (x)]v_{1,x}(x^{})+\theta (x)v_{2,x}(x^{}).$$ (57) For binary $`\theta (x)`$ Eq. (53) corresponds to an inverse covariance $`𝐊_0(\theta )`$ $`=`$ $`{\displaystyle 𝑑x𝐊_x(\theta )}={\displaystyle 𝑑xW_x(\theta )W_x^T(\theta )}`$ (58) $`=`$ $`{\displaystyle 𝑑x\left([1\theta (x)]W_{1,x}W_{1,x}^T+\theta (x)W_{2,x}W_{2,x}^T\right)}`$ with $$𝐊_x(\theta )=W_x(\theta )W_x^T(\theta )$$ (59) written as dyadic product of the vector $`W_x(\theta )`$ = $`𝐖(x,;\theta )`$ and with analogously defined $`W_{i,x}`$ = $`𝐖_i(x,)`$. For $`\theta `$–dependent inverse covariances the normalization factors $`Z_𝒱(\theta )`$ become $`\theta `$–dependent. They have to be included when integrating over $`\theta `$ or solving for the optimal $`\theta `$ in MAP. In Eqs. (52) and (53) it is straightforward to introduce two binary hyperfields $`\theta `$, $`\theta ^{}`$, one for the reference potential $`v_x`$ and one for the filter $`𝐖`$. This results in a conditional prior $`p(v\|\theta ,\theta ^{})`$ $``$ $`e^{\frac{1}{2}{\scriptscriptstyle 𝑑xvv_x(\theta )\|𝐊_x(\theta ^{})\|vv_x(\theta )}}`$ (60) $`=`$ $`e^{\frac{1}{2}{\scriptscriptstyle 𝑑x\|\omega (x;\theta ,\theta ^{})\|^2}}.`$ Here we can write $`{\displaystyle 𝑑x\|\omega (x;\theta ,\theta ^{})\|^2}`$ $`=`$ $`vv_0(\theta ,\theta ^{})\|𝐊_0(\theta ^{})\|vv_0(\theta ,\theta ^{})`$ (61) $`+{\displaystyle 𝑑xv_x(\theta )\|𝐊_x(\theta ^{})\|v_x(\theta )}`$ $`v_0(\theta ,\theta ^{})\|𝐊_0(\theta ^{})\|v_0(\theta ,\theta ^{}),`$ with an effective template $`v_0(\theta ,\theta ^{})`$ given by $$v_0(\theta ,\theta ^{})=𝐊_0(\theta ^{})^1𝑑x𝐊_x(\theta ^{})v_x(\theta ),$$ (62) and effective inverse covariance $`𝐊_0(\theta ^{})`$ = $`𝑑x𝐊_x(\theta ^{})`$ as in Eq. (58). Since the last two terms in Eq. (61) are $`v`$–independent constants (only depending on $`\theta `$, $`\theta ^{}`$) we see that for fixed hyperfields this prior is minimized by $`v`$ = $`v_0(\theta ,\theta ^{})`$. For given hyperparameters $`\theta `$, $`\theta ^{}`$ we can write $`p(v\|\theta ,\theta ^{})e^{E(v\|\theta ,\theta ^{})}`$ with a prior energy of the form $`E(v\|\theta ,\theta ^{})`$ = $`\frac{1}{2}vv_0(\theta ,\theta ^{})\|𝐊_0(\theta ^{})\|vv_0(\theta ,\theta ^{})`$. As the product of Gaussians is again a Gaussian several Gaussian prior factors can easily be combined. In this way one can implement a nonlocal property like smoothness and still avoid local template functions $`v_x(x^{},\theta )`$ by combining a Gaussian prior with $`𝐊_0`$ = $`𝐈`$ as in (56) with a Gaussian prior with nondiagonal covariance and zero (or fixed) template, $$E(v\|\theta )=\frac{1}{2}v\stackrel{~}{v}_0(\theta )\|v\stackrel{~}{v}_0(\theta )+\frac{1}{2}v\|𝐊\|v.$$ (63) Combining both terms yields $`E(v\|\theta )`$ $`=`$ $`{\displaystyle \frac{1}{2}}(vv_0(\theta )\|𝐊_0\|vv_0(\theta )`$ (64) $`+\stackrel{~}{v}_0(\theta )\|𝐈𝐊_0^1\|\stackrel{~}{v}_0(\theta )),`$ with the second term being independent of $`v`$ and with effective template and effective inverse covariance $$v_0(\theta )=𝐊_0^1\stackrel{~}{v}_0(\theta ),𝐊_0=𝐈+𝐊.$$ (65) For differential operators $`𝐊_0`$ the effective $`v_0(\theta )`$ is thus a smoothed version of $`\stackrel{~}{v}_0(\theta )`$. The extreme case would be to treat $`v_0`$ and $`𝐖`$ itself as unrestricted hyperfields. As already discussed, this just eliminates the corresponding prior term. Hence, to restrict the flexibility, typically a smoothness hyperprior may be imposed to prevent highly oscillating functions $`\theta (x)`$. For real $`\theta (x)`$, for example, a smoothness prior like a Laplacian prior $`\theta \|\mathrm{\Delta }\|\theta /2`$ can be used in regions where it is defined. (The space of functions for which a smoothness prior with discontinuous templates is defined depends on the locations of the discontinuities.) An example of a non–Gaussian hyperprior is $$p(\theta )e^{\frac{\tau }{2}{\scriptscriptstyle 𝑑xC_\theta (x)}},$$ (66) where $`\tau `$ is a constant and $$C_\theta (x)=\sigma \left(\left(\frac{\theta }{x}\right)^2\vartheta _\theta \right),$$ (67) with a sigmoid $`\sigma (x)`$ as in (51). For $`\nu \mathrm{}`$ the sigmoid approaches a step function and $`C_\theta (x)`$ becomes zero at locations where the square of the first derivative is smaller than a certain threshold $`0\vartheta _\theta <\mathrm{}`$, and one otherwise. For discrete $`x`$ one can analogously count the number of jumps larger than a given threshold. One can then penalize the number $`N_d(\theta )`$ of discontinuities where $`\left(\theta /x\right)^2`$ = $`\mathrm{}`$ and use $$p(\theta )e^{\frac{\tau }{2}N_d(\theta )}.$$ (68) In the case of a binary field this corresponds to counting the number of times the field changes its value. The expression $`C_\theta `$ of Eq. (67) can be generalized to $$C_\theta (x)=\sigma \left(\|\omega _\theta (x)\|^2\vartheta _\theta \right),$$ (69) where, analogous to Eq. (43), $$\omega _\theta (x)=𝑑x^{}𝐖_\theta (x,x^{})[\theta (x^{})t_\theta (x^{})],$$ (70) with template $`t_\theta (x^{})`$ representing the expected form for the hyperfield, and a filter operator $`𝐖_\theta `$ defining a distance measure for hyperfields. Parameters of the hyperprior like $`\tau `$ in Eq. (66) or Eq. (68) can be treated as higher level hyperparameters. ### 4.7 Non–Gaussian priors and auxiliary fields As an alternative to introducing hyperfields $`\theta (x)`$ one can work with priors which are explicitly non–Gaussian with respect to $`v`$. This can be done by introducing auxiliary fields $`B(x;v)`$ whose function values are not considered as independent variables but are directly defined as functionals of $`v`$. (For the sake of simplicity we will for $`B(x;v)`$ also write $`B(x)`$ or $`B(v)`$, depending on the context.) Like hyperfields, auxiliary fields can select locally the best adapted filtered difference from a set of alternative $`\omega _i`$. For instance, consider the auxiliary field \[compare with Eqs. (50) or (69)\] $$B(x)=\sigma \left(u(x)\vartheta \right),$$ (71) where $$u(x)=\|\omega _1(x)\|^2\|\omega _2(x)\|^2,$$ (72) $`\vartheta `$ represents a threshold, $`\sigma (x)`$ a sigmoidal function as in (51), and the $`\omega _i`$ are filtered differences defined in terms of $`v`$ according to Eq. (43). Again a binary field $`B(x)`$ is obtained by letting the sigmoid approach the step function. Because the $`\omega _i`$ depend on $`v`$, it is clear from the definition (71) that the auxiliary field $`B(x)`$ is no independent hyperfield but has values being functionals of $`v`$. Notice that $`B(x)`$ is nonlocal with respect to $`v(x)`$ if $`\omega _i(x)`$ is nonlocal; a value $`B(x)`$ then depends on more than one $`v(x)`$–value. For a negative Laplacian prior in one–dimension Eq. (71) reads, $$B(x)=\sigma \left(\left\|\frac{(vv_1)}{x}\right\|^2\left\|\frac{(vv_2)}{x}\right\|^2\vartheta \right).$$ (73) While auxiliary fields $`B(x)`$ are directly determined by $`v`$, hyperfields are indirectly coupled to $`v`$ through the MAP stationarity equations. Conversely, an auxiliary field $`B(x)`$ can be treated formally as independent hyperfield if a Lagrange multiplier term $`\lambda \left[B(x)\sigma \left(u(x)\vartheta \right)\right]`$ is added to the prior energy in the limit $`\lambda \mathrm{}`$. Like hyperfields $`\theta (x)`$ auxiliary fields $`B(x)`$ can be used to adapt reference potentials $`v_0`$ or filters $`𝐖`$. However, a prior as in Eq. (11) is non–Gaussian with respect to $`v`$ if $`v_0(B)`$ and $`𝐊_0(B)`$ depend on $`B`$ and thus also on $`v`$. Furthermore, analogous to hyperpriors $`p(\theta )`$, additional prior terms $`p(B(v))\mathrm{exp}(E_B(v))`$ for $`v`$ can be included, formulated in terms of an auxiliary field $`B(x)`$. As in Eq. (49) a binary $`B(x)`$ can switch between two filtered differences $$\|\omega (x;B)\|^2=[1B(x)]\|\omega _1(x)\|^2+B(x)\|\omega _2(x)\|^2,$$ (74) within a (non–Gaussian) prior for $`v`$ $$p(v)e^{E(v)E_B(v)},$$ (75) where the normalization factor $`Z`$ = $`𝑑ve^{E(v)E_B(v)}`$ of (75) is by definition independent of $`v`$. Hence it can be skipped for MAP calculations also for non–Gaussian $`p(v)`$. In Eq. (75) $$E(v)=\frac{1}{2}𝑑x\left([1B(x)]\|\omega _1(x)\|^2+B(x)\|\omega _2(x)\|^2\right),$$ (76) according to Eq. (74), while $`E_B(v)`$ depends on $`v`$ only through $`B(v)`$. For example, the number of switchings can be restricted by taking $$E_B(v)=\frac{\tau }{2}N_d(B),$$ (77) where $`N_d(B)`$ counts the number of discontinuities of $`B(x)`$. Other choices, for real $`B(x)`$, are quadratic energies $$E_B(v)=\frac{\tau }{2}𝑑x\|\omega _B(x)\|^2$$ (78) or non–quadratic energies of the form $$E_B(v)=\frac{\tau }{2}𝑑xC_B(x)$$ (79) where, similar to (69), $$C_B(x)=\sigma \left(\|\omega _B(x)\|^2\vartheta _B\right).$$ (80) and $$\omega _B(x)=𝑑x^{}𝐖_B(x,x^{})[B(x^{})t_B(x^{})],$$ (81) is a filtered difference of $`B`$ with filter operator $`𝐖_B`$ and template $`t_B`$. Let us compare a non–Gaussian prior built of prior energies (76) and (77) for a binary auxiliary field (71) $$p(v)e^{\frac{1}{2}{\scriptscriptstyle 𝑑x\left([1B(x)]\|\omega _1(x)\|^2+B(x)\|\omega _2(x)\|^2\right)}\frac{\tau }{2}N_d(B)},$$ (82) with the similar–looking combination of Gaussian prior (49) with hyperprior (68) for a binary hyperfield, $$p(v,\theta )=p(v\|\theta )p(\theta )$$ (83) $$e^{\frac{1}{2}{\scriptscriptstyle 𝑑x\left[(1\theta (x))\|\omega _1(x)\|^2+\theta (x)\|\omega _2(x)\|^2\right]}\frac{\tau }{2}N_d(\theta )\mathrm{ln}Z_𝒱(\theta )}.$$ Eq. (83) works with conditional probabilities $`p(v\|\theta )`$, hence the corresponding normalization factors are in general $`\theta `$–dependent and have to be included for MAP calculations. Typically, MAP solutions for $`B`$, $`N_d(B)`$ and $`C_B`$ being directly defined in terms of the corresponding MAP solution for $`v`$ are different from the MAP solutions for $`\theta `$, $`N_d(\theta )`$ and $`C_\theta `$, respectively. However, if the filtered differences $`\omega _i`$ in Eq. (83) differ only in their templates, the normalization term can be skipped. Then assuming $`\vartheta `$ = $`0`$, $`p(\theta )1`$, $`p(B)1`$ the two equations are equivalent for $`\theta (x)`$ = $`\mathrm{\Theta }\left(\|\omega _1(x)\|^2\|\omega _2(x)\|^2\right)`$. In the absence of hyperpriors, it is indeed easily seen that this is a selfconsistent solution for $`\theta `$ for every given $`v`$. In general, however, when hyperpriors are included, another solution for $`\theta `$ may have a larger posterior. Hyperpriors $`p(\theta )`$ or additional auxiliary prior terms $`p(B)`$ can be useful to enforce specific global constraints for $`\theta (x)`$ or $`B(x)`$. In natural images, for example, discontinuities are expected to form closed curves. Priors or hyperpriors, organizing discontinuities along lines or closed curves, are thus important for image segmentation or image restoration Geman-Geman-1984 ; Poggio-Torre-Koch-1985 ; Marroquin-Mitter-Poggio-1987 ; Geiger-Girosi-1991 ; Zhu-Yuille-1996 . A similar method has been used in the determination of piecewise smooth relaxation time spectra from rheological data Roths-Maier-Friedrich-Marth-Honerkamp-2000 . Another useful class of non–Gaussian priors generalizing (44) has the form Winkler-1995 ; Zhu-Mumford-1997 ; Zhu-Wu-Mumford-1997 $$p(v)e^{\frac{1}{2}{\scriptscriptstyle 𝑑x\psi [\omega (x)]}},$$ (84) where $`\psi `$ is a non–quadratic function. This function $`\psi `$ can be fixed in advance for a given problem or adapted using hyperparameters. Typical choices to allow discontinuities are symmetric “cup” functions with minimum at zero and flat tails for which one large step is cheaper than many small ones (see Fig. 1). Table 1 summarizes the basic variants of prior energies discussed in the paper. ## 5 Stationarity equations To reconstruct a local potential $`v`$ in MAP we have to maximize the posterior $`p(v\|D)`$ with respect to $`v`$. If the functional derivative of the posterior with respect to $`v`$ exists, the reconstructed potential can be found by solving the stationarity equation $$\delta _v\mathrm{ln}p(v\|D)=0,$$ (85) where we have chosen the logarithm for technical convenience, and $`\delta _v`$ denotes the functional derivative with respect to $`v`$. For observational data consisting of $`n`$ independent position measurements the posterior (2) reads $$p(v\|D)p(v)\underset{i=1}{\overset{n}{}}p(x_i\|\widehat{x},v).$$ (86) To formulate the stationarity equation (85) we have to calculate the functional derivatives of likelihood and prior. For inverse quantum statistics Lemm-IQS-2000 the likelihood for position measurements (8) on a canonical ensemble (6) depends on the eigenfunctions and eigenvalues of the $`v`$–dependent Hamiltonian $`H(v)`$. We thus have to find the functional derivatives of the eigenfunctions $`\varphi _\alpha `$ and eigenvalues $`E_\alpha `$. Those can be obtained by taking the functional derivative of the eigenvalue equation $`H\|\varphi _\alpha `$ = $`E_\alpha \|\varphi _\alpha `$, where we will assume the eigenfunctions to be orthonormalized. Choosing $`\varphi _\alpha \|\delta _{v(x)}\varphi _\alpha `$ = 0 and utilizing $$\delta _{v(x)}H(x^{},x^{\prime \prime })=\delta _{v(x)}V(x^{},x^{\prime \prime })=\delta (xx^{})\delta (x^{}x^{\prime \prime }),$$ (87) we find for nondegenerate eigenfunctions $`\delta _{v(x)}E_\alpha `$ $`=`$ $`\varphi _\alpha \|\delta _{v(x)}H\|\varphi _\alpha =\|\varphi _\alpha (x)\|^2,`$ (88) $`\delta _{v(x)}\varphi _\alpha (x^{})`$ $`=`$ $`{\displaystyle \underset{\gamma \alpha }{}}{\displaystyle \frac{1}{E_\alpha E_\gamma }}\varphi _\gamma (x^{})\varphi _\gamma ^{}(x)\varphi _\alpha (x).`$ (89) It follows for the functional derivative of the likelihood $`\delta _{v(x)}p(x_i\|\widehat{x},v)`$ $`=`$ $`\left(\delta _{v(x)}\varphi ^{}(x_i)\right)\varphi (x_i)`$ (90) $`+\varphi ^{}(x_i)\delta _{v(x)}\varphi (x_i)`$ $`\beta (\|\varphi (x_i)\|^2\|\varphi (x)\|^2`$ $`\|\varphi (x_i)\|^2\|\varphi (x)\|^2).`$ Having obtained Eq. (90) for the likelihood we now have to find the functional derivative of the prior. For the Gaussian prior (11) one gets directly $$\delta _v\mathrm{ln}p(v)=𝐊_0(vv_0).$$ (91) If hyperparameters $`\theta `$ are included and treated in MAP (i.e., not integrated out by Monte Carlo techniques), the posterior has to be maximized simultaneously with respect to $`v`$ and $`\theta `$. We have already mentioned that $`\theta `$–dependent inverse covariances lead to normalization factors which are independent of $`v`$ but depend on $`\theta `$. Such factors have to be included when maximizing with respect to $`\theta `$. As a non–Gaussian example consider a prior where two filtered differences are mixed by an auxiliary field $`B(x)`$ and an additional prior factor $`p(B)`$ is included, for example to prevent fast oscillations of $`B(x)`$. With $`B(x)`$ = $`\sigma (u(x)\vartheta )`$, threshold $`\vartheta `$, sigmoidal function $`\sigma (x)`$ as in Eq. (51), and $`u(x)`$ = $`\|\omega _1(x)\|^2\|\omega _2(x)\|^2`$ this gives $$p(v)e^{\frac{1}{2}{\scriptscriptstyle 𝑑x\left\|[1B(x)]\omega _1(x)+B(x)\omega _2(x)\right\|^2}E_B}.$$ (92) Analogous to Eq. (75), the term $$E_B=𝑑xE_B(x),$$ (93) represents an auxiliary prior energy formulated in terms of the mixing function $`B(x)`$. Like $`\omega (x)`$ the function value $`E_B(x)`$ may depend on the whole function $`B`$ and not necessarily only on the function value $`B(x)`$. Using $`\omega _i(x)`$ = $`x\|𝐖_i(vv_i)`$ we find $$\delta _{v(x)}\omega _i(x^{})=𝐖_i(x^{},x),$$ (94) and thus $$\delta _{v(x)}u(x^{})=2\left(𝐖_1^T(x,x^{})\omega _1(x^{})𝐖_2^T(x,x^{})\omega _2(x^{})\right).$$ (95) Furthermore, we obtain for the functional derivative of $`E_B`$ $$\delta _{v(x)}E_B(x^{})=𝑑x^{\prime \prime }\left[\delta _{v(x)}B(x^{\prime \prime })\right]\left[\delta _{B(x^{\prime \prime })}E_B(x^{})\right],$$ (96) where with Eq. (71) $$\delta _{v(x)}B(x^{\prime \prime })=\sigma ^{}(u(x^{\prime \prime })\vartheta )\delta _{v(x)}u(x^{\prime \prime }),$$ (97) and $`\sigma ^{}(u)`$ = $`d\sigma (u)/du`$. For a prior energy as in (78) which is quadratic in $`B(x)`$ $$E_B(x)=\|\omega _B(x)\|^2,$$ (98) $`\omega _B(x)`$ defined in Eq. (81), the functional derivative with respect to $`B(x)`$ becomes $$\delta _{B(x)}E_B(x^{})=2𝐖_B^T(x,x^{})\omega _B(x^{}).$$ (99) For a non–Gaussian prior with energy (79) an additional derivative of the sigmoid appears. Now all terms can be collected and inserted into the functional derivative of the prior (92) $`\delta _v\mathrm{ln}p(v)`$ $`=`$ $`{\displaystyle }dx([[1B(x)]\omega _1(x)+B(x)\omega _2(x)]`$ (100) $`\times ([1B(x)]\delta _v\omega _1(x)+B(x)\delta _v\omega _2(x)`$ $`+\delta _vB(x)[\omega _2(x)\omega _1(x)])`$ $`+\delta _vE_B(x)).`$ The Bayesian approach to inverse quantum mechanics is not restricted to position measurements, but allows to deal with all kinds of observations for which the likelihood can be calculated. To have better information about the depth of a potential it is useful to include information on the ground state energy of a system. For instance, including a noisy measurement of the average energy $$U=E=\underset{\alpha }{}p_\alpha E_\alpha ,$$ (101) yields an additional factor in the posterior of the form $$p_Ue^{E_U},E_U=\frac{\mu }{2}(U\kappa )^2.$$ (102) In the noise free limit $`\mu \mathrm{}`$ this yields $`U\kappa `$. Calculating the functional derivative of $`U`$ with respect to a local potential $$\delta _{v(x)}U=\delta _{v(x)}E\beta E\delta _{v(x)}E+\beta E\delta _{v(x)}E,$$ (103) it is straightforward to obtain $$\delta _{v(x)}E_U=\mu \left(U\kappa \right)\|\varphi (x)\|^2\left[1\beta \left(EU\right)\right].$$ (104) Stationarity equations are typically nonlinear and have to be solved by iteration. A possible iteration scheme is $`v^{(r+1)}`$ $`=`$ $`v^{(r)}+\eta 𝐀^1[\delta _v\mathrm{ln}p(v^{(r)})`$ (105) $`+{\displaystyle \underset{i}{}}\delta _v\mathrm{ln}p(x_i\|\widehat{x},v^{(r)})\delta _vE_U^{(r)}].`$ Here $`\eta `$ is a step width which can be optimized by a line search algorithm and the positive definite operator $`𝐀`$ distinguishes different learning algorithms. ## 6 Numerical examples As numerical application of BIQM and to test several variants of implementing a priori information we will study the reconstruction of an approximately periodic, one–dimensional potential. Such a potential may represent a one–dimensional surface where a periodic structure, e.g. that of a regular crystal, is distorted by impurities, located at unknown positions and of unknown form. To test the quality of reconstruction algorithms, artificial data will be sampled from a model with known “true” potential $`v_{\mathrm{true}}`$. Selecting a specific prior model and applying the corresponding Bayesian reconstruction algorithm to the sampled data, we will be able to compare the reconstructed potential with the original one. In particular, we will take as true potential the following perturbed periodic potential $$v_{\mathrm{true}}(x)=\{\genfrac{}{}{0pt}{}{\mathrm{sin}\left(\frac{2\pi }{6}x\right);1x12,\mathrm{\hspace{0.17em}\hspace{0.33em}25}x36,}{\mathrm{sin}\left(\frac{2\pi }{12}x\right);13x24,}$$ (106) using for the numerical calculations a mesh of size 36. Considering a system prepared as canonical ensemble the potential $`v_{\mathrm{true}}`$ defines a corresponding canonical density operator $`\rho `$ as given in Eq. (6). Artificial data $`D`$ can then be sampled according to the likelihood model of quantum mechanics (5). For the following examples, $`n`$ = 200 data points representing position measurements have been sampled using the transformation method Press-Teukolsky-Vetterling-Flannery-1992 . In all calculations we used periodic boundary conditions for quantum mechanical wave functions while the potential $`v`$ has been set to zero at the boundaries. We will now discuss the results of a Bayesian reconstruction under varying prior models. As first example, consider a simple Gaussian prior (11) with negative Laplacian inverse covariance $`𝐊_0`$ = $`\lambda \mathrm{\Delta }`$, zero reference potential $`v_00`$, and an additional prior factor (102) representing a noisy measurement of the average energy. The reconstruction results are shown in Fig. 3. In particular, the figure on top compares the reconstructed likelihood $`p_{\mathrm{BIQM}}(x\|\widehat{x},v_{\mathrm{BIQM}})`$ with the true likelihood $`p_{\mathrm{true}}(x\|\widehat{x},v_{\mathrm{true}})`$ and with the empirical density, i.e., the relative frequencies of the sampled data $$p_{\mathrm{emp}}(x)=\frac{1}{n}\underset{i=1}{\overset{n}{}}\delta (xx_i).$$ (107) Similarly, the lower figure compares the reconstructed potential $`v_{\mathrm{BIQM}}`$ with the true potential $`v_{\mathrm{true}}`$. Since information on the average energy was available the depth of the potential is well approximated at least at one of its minima. This is sufficient to fulfill the noisy average energy condition. However, because only smoothness and no periodicity information is implemented by the prior the reconstructed potential is too flat. The effect is stronger near the maxima than near the minima of the potential because near the maxima only few data points are available and hence the reconstructed potential is there dominated by the zero reference potential in the smoothness prior. To include information on approximate periodicity we have replaced in the next example the zero reference potential $`v_00`$ by the strictly periodic reference potential $$v_0(x)=\mathrm{sin}\left(\frac{2\pi }{6}x\right),$$ (108) shown as dashed line in the following figures of potentials. A reconstruction with the periodic reference potential (108) but without average energy information, and starting the iteration with the reference potential as initial guess $`v^{(0)}`$ = $`v_0`$ is shown in Fig. 4. Due to missing average energy information the depth of the potential is not well approximated. It is also clearly visible in Fig. 4 that the smoothness prior does not favor solutions which are similar to the reference $`v_0`$ itself but solutions which have derivatives similar to that of $`v_0`$. Fig. 4 also displays that the reconstruction of the potential does clearly identify the impurity. As the reference potential is not adapted to the impurity region the reconstruction is there poorer than in the regular region. Furthermore, it is worth emphasizing that the reconstructed likelihood fits the empirical density well, even slightly better than the true likelihood does. This is due to the flexibility of a nonparametric approach which allows to fit the fluctuations of the empirical density caused by the finite sample size. The effect is well known in empirical learning and leads to so called “overfitting” if the influence of the prior becomes to small. Since observational data influence the reconstruction only through the likelihood, the reconstruction of potentials is in general a more difficult task than the reconstruction of likelihoods. This indicates the special importance of a priori information when reconstructing potentials. Indeed, even if the complete likelihood is given, the problem of determining the potential can still be ill–defined in regions where the likelihood is small Zhu-Rabitz-1999 . A prior model with periodic reference potential can be made more flexible by adapting amplitude, frequency, and phase of the reference potential (108). For this purpose one can introduce a hyperparameter vector $`\theta `$ = $`(\theta _1,\theta _2,\theta _3)`$ parameterizing amplitude, frequency, and phase and take as reference potential $$v_0(x;\theta )=\theta _1\mathrm{sin}\left(\frac{2\pi }{\theta _2}x+\theta _3\right).$$ (109) The corresponding maximization of the posterior with respect to $`\theta `$ is easy in that case and does not change the results of Fig. 4 where the hyperparameters are already optimally adapted. Including an additional noisy energy measurement (102) Fig. 5 shows that the depth of the potential is indeed better approximated than in Fig. 4. To avoid local maxima of the posterior the solution of Fig. 4 has been used as initial guess and the factor $`\mu `$ multiplying the average energy term has been slowly increased to its final value. Fig. 5 still only represents a local and no global maximum of the posterior, as can be by seen by starting with a different initial guess $`v^{(0)}`$. In Fig. 6 a better solution for the same parameters is presented where the initial guess has been selected using a priori information about the location of the impurity region. Alternatively to a Gaussian prior with periodic reference, approximate periodicity can be enforced by the inverse covariance of a Gaussian prior. In this case the prior favors periodicity but no special form of the potential. The prior is thus less specific than a prior with explicit periodic reference function. Corresponding BIQM results for the inverse covariance (27) are shown in Fig. 7. Indeed while the potential is well approximated in regions where many observations have been collected, it is not as well approximated in regions where no or only few data are available. These are the regions where the prior dominates the observational data. In particular, in the case presented in Fig. 7, the zero reference function $`v_00`$ of an additional Laplacian smoothness prior implements a tendency to flat potentials. If impurities are expected, a prior with one fixed periodic reference potential for the whole region is no adequate choice. Near impurities one would like to switch off the standard periodic reference potential which in these regions will be misleading. Because it is usually not known in advance where a given reference should be used and where not, those regions must be identified during learning. As first example we study a prior energy similar to Eq. (63), $`E(v)`$ $`=`$ $`{\displaystyle \frac{\lambda _1}{2}}{\displaystyle 𝑑x\|v(x)v_0(x)\|^2[1B(x)]}{\displaystyle \frac{\lambda _2}{2}}v\|\mathrm{\Delta }\|v,`$ which allows to switch off a given reference locally by means of a binary switching function defined as $`B(x)`$ = $`\mathrm{\Theta }\left(\|v(x)v_0(x)\|^2\vartheta \right)`$. (An average energy term $`E_U`$ = $`\frac{\mu }{2}(U\kappa )^2`$ could easily be included.) In the prior energy (LABEL:eq1) the reference $`v_0`$ is only used if $`\|v(x)v_0(x)\|^2`$ is smaller than the given threshold $`\vartheta `$. Starting with a smoothed version of Eq. (LABEL:eq1) with a real mixing function $`B(x)`$ = $`\sigma \left(\|v(x)v_0(x)\|^2\vartheta \right)`$, the results of Fig. 8 have been obtained by changing during iteration $`\sigma (x)`$ slowly from a sigmoid to a step function. Using a step function for $`B`$ directly from the beginning leads to nearly indistinguishable results. Compared to Fig. 7 the reconstruction in Fig. 8 is improved mainly in the unperturbed region where the algorithm can now use the correct reference potential. An additional advantage is that the final auxiliary field $`B(x)`$ directly shows the identified impurity regions. One sees in Fig. 8 that the auxiliary field $`B(x)`$ is always switched off if the solution $`v(x)`$ is similar enough to the template $`v_0(x)`$. The two $`v`$–dependent terms in Eq. (LABEL:eq1) can be combined \[compare Eqs. (63) and (64)\]. Skipping a term which only depends on $`v`$ through $`B(x)`$, one arrives at another prior which also implements local switching. More general, choosing the prior energy (76) for switching between two filtered differences with two reference potentials $`v_1`$ and $`v_2`$ leads to $`E(v)`$ $`=`$ $`{\displaystyle \frac{\lambda _1}{2}}{\displaystyle 𝑑x[1B(x)]\|\omega _1(x)\|^2}`$ (111) $`+{\displaystyle \frac{\lambda _2}{2}}{\displaystyle 𝑑xB(x)\|\omega _2(x)\|^2},`$ where the switching is controlled by the binary function $`B(x)`$ = $`\mathrm{\Theta }\left(\|\omega _1(x)\|^2\|\omega _2(x)\|^2\vartheta \right)`$ defined in terms of the filtered differences $`\omega _i(x)`$ = $`(/x)[v(x)v_i(x)]`$. A prior energy (111) with two different nonzero reference potentials $`v_1`$ and $`v_2`$ is obtained, for example, when a different nonzero reference potential is given for the unperturbed and the perturbed region. The number of changes in the switching function $`B(x)`$ = $`\mathrm{\Theta }\left(\|\omega _1(x)\|^2\|\omega _2(x)\|^2\vartheta \right)`$, can be controlled by adding a prior term $`p(B)`$ penalizing the number of times the function $`B(x)`$ changes its value. To avoid local minima for binary $`B(x)`$, simulated annealing techniques are useful. We have obtained an initial guess for $`v`$, and thus for $`B(x)`$, by writing $`v(x)`$ = $`[1c(x)]v_1(x)+c(x)v_2(x)`$ and optimizing the binary function $`c(x)`$ by simulated annealing with respect to the likelihood and the additional prior $`p(B)`$. In particular, starting from $`c(x)`$ = $`0`$, new trial functions have been generated by selecting two points $`x_1`$, $`x_2`$ randomly and exchanging the function values zero and one in between (see Fig. 2). A new trial function has been accepted or rejected using the Metropolis rule $`p(`$accept) = min$`[1,\mathrm{exp}(\beta _{\mathrm{ann}}\mathrm{\Delta }E_{\mathrm{ann}})]`$ with $`\mathrm{\Delta }E_{\mathrm{ann}}`$ denoting the difference in the error between actual function and new trial function. In the present case we have $`E_{\mathrm{ann}}(v)`$ = $`_iE(x_i\|\widehat{x},v)`$ \+ $`E_B(v)`$ where $`E(x_i\|\widehat{x},v)`$ = $`\mathrm{ln}p(x_i\|\widehat{x},v)`$ and $`p(B)\mathrm{exp}(E_B)`$. The annealing temperature $`1/\beta _{\mathrm{ann}}`$ decreases during optimization. Fig. 9 shows the reconstruction results using the following two reference potentials $`v_1(x)`$ $`=`$ $`{\displaystyle \frac{2}{3}}\mathrm{sin}\left({\displaystyle \frac{2\pi }{6}}x\right),`$ (112) $`v_2(x)`$ $`=`$ $`\mathrm{sin}^2\left({\displaystyle \frac{2\pi }{6}}x\right)\mathrm{sign}\left[\mathrm{sin}\left({\displaystyle \frac{2\pi }{6}}x\right)\right].`$ (113) Compared to Fig. 8 the reconstruction is improved in the perturbed region, where the algorithm can now rely on a useful reference potential. Finally, the switching function can be introduced as local hyperfield. As an example for a prior with hyperfield, Fig. 10 shows the reconstruction with the prior energy $$E(v,\theta )=\frac{\lambda _1}{2}vv_0(\theta )\|vv_0(\theta )\frac{\lambda _2}{2}v\|\mathrm{\Delta }\|v\mathrm{ln}p(\theta ),$$ (114) where $`v_0(x;\theta )`$ = $`v_1(x)[1\theta (x)]+v_2(x)\theta (x)`$ with the reference potentials of Eq. (112) and Eq. (113). A hyperprior $`p(\theta )`$ has been used penalizing the number of discontinuities of the hyperfield $`\theta (x)`$, analogous to $`p(B)`$ for Fig. 9. The $`E(v\|\theta )`$ part of the prior energy (114) is of the form (49) with $`\theta `$–independent covariances. Hence the $`\theta `$–independent normalization factor can be skipped. An initial guess for the local hyperfield $`\theta (x)`$ has been obtained by simulated annealing as described for Fig. 9. As in this case optimization is required only with respect to the $`\theta `$–dependent parts of the posterior, optimizing $`\theta (x)`$ for given $`v`$ is faster than optimizing $`v`$ through $`c(x)`$ which requires diagonalization of the hamiltonian $`H`$ for every new trial function. However, as $`\theta (x)`$ is independent of $`v`$, the hyperfield has to be updated during iteration which has also been done by simulated annealing. As expected, a reconstruction with the non–Gaussian prior corresponding to the prior energy (111) is very similar to a reconstruction using hyperfields as in Eq. (114). ## 7 Conclusion A nonparametric Bayesian approach has been developed and applied to the inverse problem of reconstructing potentials of quantum systems from observational data. Relying on observational data only the problem is typically ill–defined. It is therefore essential to include adequate a priori information. Since reconstructed potentials obtained by Bayesian Inverse Quantum Mechanics (BIQM) depend sensitively on the implemented a priori information, flexible prior models are required which can be adapted to the specific situation under study. In particular, the use of hyperparameters, hyperfields, and non–Gaussian priors with auxiliary fields has been discussed in detail. In this paper we have focussed on the implementation of approximate periodicity for potentials in inverse problems of quantum statistics. The presented prior models, however, can be useful for many empirical learning problems, including for example regression or general density estimation. Several variants of implementing a priori information on approximate periodicity have been tested and compared numerically.
warning/0005/math0005114.html	ar5iv	text	# On subgroups of R. Thompson’s group 𝐹 and other diagram groups ## Introduction This paper is devoted to further study of the so called diagram groups. The definition of diagram groups was first given by Meakin and Sapir in 1995. Their student Vesna Kilibarda obtained first results about diagram groups in her thesis (see also her paper ). Further results about diagram groups have been obtained in our paper . Here we survey the main results of that paper (see for details). Diagram groups reflect certain important properties of semigroup presentations. For instance we showed that three definitions of asphericity given by Pride are in fact equivalent and are equivalent to the triviality of all diagram groups over the presentation. One can say that diagram groups measure the non-asphericity of semigroup presentations. On the other hand, it turned out that the class of diagram groups is interesting even if we forget about its connection with semigroup presentations. These groups have nice algorithmic properties: the word problem in every diagram group is solvable in time $`O(n^{2+\epsilon })`$ for every $`\epsilon >0`$. This does not depend on whether the word problem for the corresponding semigroup presentation is solvable or not. If the word problem of this semigroup presentation is solvable then the conjugacy problem is solvable in the corresponding diagram group. If a group is representable by diagrams (i.e. it is a subgroup of a diagram group) then one can use geometry of planar graphs to deduce certain properties of the group. Diagrams can be viewed as “2-dimensional words” and in , we developed a calculus called “combinatorics on diagrams”, which is parallel to the well known combinatorics on words (see Lothaire ). Geometry of diagrams allows one to consider many homomorphism from diagram groups into the group of piecewise linear homeomorphisms of the real line. Thus we have a connection between groups representable by diagrams and groups representable by piecewise linear functions. This connection can be used in both directions. We showed that the class of diagram groups is wide. It contains the free groups, free abelian groups, the R. Thompson group $`F`$ and its generalizations found by Brown . This class is closed under finite direct products, arbitrary free products and some other constructions. Note that the Thompson group is the diagram group over the following simple presentation $`xx^2=x`$. In we obtained several previously unknown results about Thompson’s group, essentially using its representation as a diagram group. * The conjugacy problem in $`F`$ is solvable. * The centralizer of each element of $`F`$ is a finite direct product of groups each of which is either a copy of $`F`$ or an infinite cyclic group $`𝐙`$. Let us give a short summary of the content of this paper. Section 1 contains the list of the main concepts used in this paper. In Section 2, we introduce the concept of diagram product of groups. It is defined as the fundamental group of a certain 2-complex of groups. Theorem 4, the main result of this section, states that the class of diagram groups is closed under diagram products. It turns out that all “products” considered before (the free product, the direct product, etc.) are particular cases of the diagram product. Examples 58, 10, 12 show applications of Theorem 4. In partricular we prove that the class of diagram groups is closed under countable direct powers (Theorem 9), wreath products with $`𝐙`$ (Theorem 11), and a certain special construction $`𝒪(G,H)`$ (Theorem 13), whose role will be clear later. In Section 3, we show that nilpotent subgroups of diagram groups are abelian (Corollary 15), and abelian subgroups are free abelian (Theorem 16). We finish this section with a description of sets of pairwise commuting diagrams (Theorem 17). In Section 4, we prove that the Thompson group $`F`$ contains subgroups isomorphic to the restricted wreath product of two infinite cyclic groups and soluble subgroups of arbitrary degree (this was first proved by Brin). It turns out that for any subgroup of piecewise linear functions (including $`F`$) there exists a dichotomy: either it contains $`𝐙wr𝐙`$ or it is abelian (Theorem 21). This implies, in particular, that a non-abelian subgroup of the group of piecewise linear functions cannot be a one-relator group (Corollary 23). This result strengthens the well known fact that the group of piecewise linear functions does not contain free non-abelian subgroups. In Section 4, we also give necessary and sufficient conditions for a diagram group to contain a copy of $`𝐙wr𝐙`$ as a subgroup (Theorem 24). We study the question when a diagram group over some semigroup presentation $`𝒫`$ contains a copy of the Thompson group $`F`$. We prove that if the semigroup given by $`𝒫`$ contains an idempotent then a diagram group over this presentation contains a copy of $`F`$. (Theorem 25). The interesting question of whether the converse statement holds is open. In Section 5, we present a counterexample to the Subgroup Conjecture. This conjecture stated that every subgroup of a diagram group is a diagram group itself. It was motivated by the similarity between diagram groups and free groups. At first we thought that the conjecture is easy to disprove and the derived subgroup $`F^{}`$ of $`F`$ is a counterexample. But it turned out that $`F^{}`$ is a diagram group (Theorem 26). This solves several problems from . We asked whether a diagram group can coincide with its derived subgroup and whether every diagram group has an LOG-presentation. Corollary 27 gives a positive answer to the first question and a negative answer to the second question. But the main result of this Section is Theorem 28 which shows that the one-relator group $`x,yxy^2x=yx^2y`$ is not a diagram group but is isomorphic to a subgroup of a diagram group. In the proof, we use the construction $`𝒪(G,H)`$ from Section 2. This gives a counterexample to the Subgroup Conjecture. Nevertheless, we think that in many partricular cases this conjecture is true and we pose several open question in this regard. At the end of this section we study the following series of groups $$G_n=x_1,\mathrm{},x_n[x_1,x_2]=[x_2,x_3]=\mathrm{}=[x_{n1},x_n]=[x_n,x_1]=1.$$ For $`n4`$ these groups are diagram groups; we prove (Theorem 30), that for odd $`n5`$ this is not so. It is not known whether these groups are representable by diagrams. If so, this will give us new counterexamples to the Subgroup Conjecture. The last Section 6 is devoted to the distortrion of subgroups in diagram groups. In a recent paper Burillo prroved that for every natural $`n`$ the Thompson group $`F`$ contains subgroups isomorphic to $`F\times 𝐙^n`$ and quasi-isometrically (without distortrion) embedded into $`F`$. A similar fact is true for an embedding of $`F\times F`$. We prove (Theorem 34) that every centralizer of an element in $`F`$ is embedded into $`F`$ without distortion. Centralizers of elements of $`F`$ can be arbitrary finite direct products of copies of $`F`$ and copies of $`𝐙`$. Burillo also proved that every cyclic subgroup of $`F`$ is embedded without distortion (this fact is also an immediate corollary of Lemma 15.29 in ). We prove (Theorem 33) that not only cyclic but arbitrary finitely generated abelian subgroups of any diagram groups are undistorted. Finally we found solvable subgroups of $`F`$ which are distorted. Theorem 38 shows that for every natural $`d2`$ there exists a finitely generated solvable subgroup $`K_d`$ in $`F`$ such that its distortion function is at least $`n^d`$. Acknowledgements. The authors thank M. Brin and S. Pride for helpful discussions of the results of this paper. ## 1 Preliminaries For an alphabet $`\mathrm{\Sigma }`$ let $`\mathrm{\Sigma }^+`$ denote the free semigroup over $`\mathrm{\Sigma }`$, and let $`\mathrm{\Sigma }^{}`$ denote the free monoid. Elements of the free monoid are called words. The identity element, i.e. the empty word, is denoted by $`1`$. Let $`𝒫=\mathrm{\Sigma }`$ be a presentation of a semigroup where $`\mathrm{\Sigma }`$ is an alphabet, $``$ is a set of pairs of non-empty words over $`\mathrm{\Sigma }`$. The semigroup $`S`$ given by $`𝒫`$ is the factor-semigroup $`\mathrm{\Sigma }^+/`$ where $``$ is the smallest congruence on $`\mathrm{\Sigma }^+`$ containing $``$. Elements of $`\mathrm{\Sigma }`$ are called generators, pairs $`(u,v)`$ written also as $`u=v`$ are called defining relations. Left and right parts of defining relations are called defining words. We shall assume that all presentations are anti-symmetric that is if $`u=v`$ then $`v=u`$. In particular $``$ does not contain relations $`u=u`$. With any semigroup presentation $`𝒫`$, we associate the following graph $`\mathrm{\Gamma }(𝒫)`$. The vertices are all words in $`\mathrm{\Sigma }^+`$. Edges are the elements of $`\mathrm{\Sigma }^{}\times ^{\pm 1}\times \mathrm{\Sigma }^{}`$. We shall denote edges by $`(x,uv,y)`$, where $`x,y\mathrm{\Sigma }^{}`$, and either $`(u,v)`$ or $`(v,u)`$. If $`e=(x,uv,y)`$ then the inverse edge is defined by $`e^1=(x,vu,y)`$. The initial vertex of $`e`$ is $`\iota (e)=xuy`$ and the terminal vertex is $`\tau (e)=xvy`$. Thus vertices of this graph are words and edges are elementary transformations of words (i.e. substitutions of defining words by their pairs). The graph $`\mathrm{\Gamma }(𝒫)`$ describes all derivations over $`𝒫`$: two non-empty words $`w_1`$, $`w_2`$ are equal modulo $`𝒫`$ if and only if there exists a path in the graph connecting $`w_1`$ and $`w_2`$. This path is called a derivation of $`w_2`$ from $`w_1`$. With every derivation over $`𝒫`$, one can associate a geometric object, a semigroup diagram over $`𝒫`$. Semigroup diagrams were first introduced by E. V. Kashintsev and then rediscovered by Remmers and others(see ). We do not give an exact definition here (see also ), the definition will be clear from the following example. ###### Example 1 Let $`𝒫=a,b,cabc=ba,bca=cb,cab=ac`$. Consider the following derivation over $`𝒫`$: $$(1,accab,cb^2ca)(c,abcba,b^2ca)(cbab,bcacb,1)(cb,abcba,b).$$ The corresponding diagram $`\mathrm{\Delta }`$ over $`𝒫`$ is the following: Let us introduce the terminology associated with diagrams. Every diagram over $`𝒫`$ is a planar graph. It has vertices, edges and cells. In Example 1, the diagram $`\mathrm{\Delta }`$ has 13 vertices, 16 edges and 4 cells. The number of cells is equal to the length of the corresponding derivation. Each positive edge has a label from $`\mathrm{\Sigma }`$, positive edges are oriented from left to right. The label of an edge $`e`$ is denoted by $`\phi (e)`$. We shall consider only positive paths in $`\mathrm{\Delta }`$, that is paths consisting of positive edges. For every path $`p`$ in a diagram $`\mathrm{\Delta }`$, its label $`\phi (p)`$ is the word read on the path. Any diagram $`\mathrm{\Delta }`$ has the initial verrtex $`\iota (\mathrm{\Delta })`$ and the terminal vertex $`\tau (\mathrm{\Delta })`$, the top path $`\text{top}(\mathrm{\Delta })`$ and the bottom path $`\text{bot}(\mathrm{\Delta })`$ connecting the initial and terminal vertices. The diagram $`\mathrm{\Delta }`$ lies between its top and bottom paths. This notation is illustrated by the following example. Notice that paths $`\text{top}(\mathrm{\Delta })`$ and $`\text{bot}(\mathrm{\Delta })`$ can have common edges. Every cell $`\pi `$ of a diagram is a diagram itself, so we can define the notation $`\iota (\pi )`$, $`\tau (\pi )`$, $`\text{top}(\pi )`$, $`\text{bot}(\pi )`$ and the corresponding concepts. If words $`u`$ and $`v`$ are labels of the top and the bottom paths of a cell $`\pi `$ then either $`u=v`$ or $`v=u`$ is a defining relation (that is it belongs to $``$). In this case we call $`\pi `$ a $`(u,v)`$-cell. For every non-empty word $`w`$, there exists a trivial diagram $`\epsilon (w)`$ without cells whose top and bottom paths coincide and have label $`w`$. We do not distinguish isotopic diagrams. The notation $`\mathrm{\Delta }_1\mathrm{\Delta }_2`$ means that $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_2`$ are isotopic. If the label of $`\text{top}(\mathrm{\Delta })`$ is $`w_1`$ and the label of $`\text{bot}(\mathrm{\Delta })`$ is $`w_2`$ then $`\mathrm{\Delta }`$ is called a $`(w_1,w_2)`$-diagram. Let $`w_1`$, $`w_2`$, $`w_3`$ be any three vertices of the graph $`\mathrm{\Gamma }(𝒫)`$ and let $`p_i`$ ($`i=1,2`$) be paths in the graph $`\mathrm{\Gamma }(𝒫)`$ from $`w_i`$ to $`w_{i+1}`$. By $`\mathrm{\Delta }_i`$, we denote the diagram corresponding to the path $`p_i`$ ($`i=1,2`$). It is easy to see that the product of paths $`p_1`$ and $`p_2`$ in the graph corresponds to the diagram $`\mathrm{\Delta }`$ obtained from $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_2`$ by identifying the bottom path of $`\mathrm{\Delta }_1`$ and the top path of $`\mathrm{\Delta }_2`$. The resulting diagram $`\mathrm{\Delta }`$ will be called the composition of diagrams $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_2`$, we denote it by $`\mathrm{\Delta }_1\mathrm{\Delta }_2`$. Thus $``$ is a partial operation on the set of all diagrams over $`𝒫`$. The composition of a $`(w_1,w_2)`$-diagram and a $`(w_2,w_3)`$-diagram is a $`(w_1,w_3)`$-diagram. For every word $`w\mathrm{\Sigma }^+`$, the set of all $`(w,w)`$-diagrams over $`𝒫`$ is a semigroup with respect to the operation $``$. Diagrams of this form will be called spherical diagrams with base $`w`$. This semigroup has the identity element $`\epsilon (w)`$. We define also another associative operation on the set of all diagrams over $`𝒫`$. Namely the sum $`\mathrm{\Delta }_1+\mathrm{\Delta }_2`$ of diagrams $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_2`$ is the diagram obtained by identifying $`\tau (\mathrm{\Delta }_1)`$ and $`\iota (\mathrm{\Delta }_2)`$. These two operations are illustrated by the following figure: Suppose that a diagram $`\mathrm{\Delta }`$ contains a $`(u,v)`$-cell and a $`(v,u)`$-cell such that the top path of the first cell is the bottom path of the second cell. Then we say that these two cells form a dipole. In this case we can remove these two cells by first removing their common path, and then identifying the bottom path of the first cell with the top path of the second cell. A diagram is called reduced if it does not contain dipoles. One can get a reduced diagram from any diagram by removing dipoles. Kilibarda proved that every diagram has a unique reduced form. We call two diagrams $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_2`$ equivalent, written as $`\mathrm{\Delta }_1\mathrm{\Delta }_2`$, if their reduced forms are the same. It is easy to see that if $`\mathrm{\Delta }_1\mathrm{\Delta }_2`$, $`\mathrm{\Delta }_3\mathrm{\Delta }_4`$ then $`\mathrm{\Delta }_1\mathrm{\Delta }_3\mathrm{\Delta }_2\mathrm{\Delta }_4`$ and $`\mathrm{\Delta }_1+\mathrm{\Delta }_3\mathrm{\Delta }_2+\mathrm{\Delta }_4`$. Therefore on the set $`𝒟(𝒫,w)`$ of all equivalence classes of $`(w,w)`$-diagrams one can define a product, by setting $`[\mathrm{\Delta }_1][\mathrm{\Delta }_2]=[\mathrm{\Delta }_1\mathrm{\Delta }_2]`$, where square brackets denote equivalence classes. The product of a $`(w_1,w_2)`$-diagram $`\mathrm{\Delta }`$ and the $`(w_2,w_1)`$-diagram $`\mathrm{\Delta }^{}`$, which is a mirror image of $`\mathrm{\Delta }`$ is obviously equivalent to the trivial diagram $`\epsilon (w_1)`$. The diagram $`\mathrm{\Delta }^{}`$ will be denoted by $`\mathrm{\Delta }^1`$. For simplicity we shall call equivalent diagrams equal, use “=” instead of “$``$”, and drop square brackets and the multiplication sign. So for every $`\mathrm{\Delta }𝒟(𝒫,w)`$ $`\mathrm{\Delta }\mathrm{\Delta }^1=\epsilon (w)`$. As a result $`𝒟(𝒫,w)`$ turns out to be a group which is called the diagram group over the semigroup presentation $`𝒫`$ with base $`w`$. Since every equivalence class contains a unique reduced diagram, one can assume that $`𝒟(𝒫,w)`$ consists of reduced diagrams with the natural multiplication ($`\mathrm{\Delta }_1\mathrm{\Delta }_2`$ is the reduced form of $`\mathrm{\Delta }_1\mathrm{\Delta }_2`$). In what follows, the term diagram group means a diagram group over some presentation with some base. We shall use the standard notation for conjugation in groups: $`a^b=b^1ab`$, and for the commutator: $`[a,b]=a^1a^b=a^1b^1ab`$. If $`A`$ and $`B`$ are subgroups of a group $`G`$ then $`[A,B]`$ denotes the subgroup generated by all commutators $`[a,b]`$ where $`aA`$, $`bB`$. Now let us shortly describe some results about diagram groups obtained earlier. The diagram group corresponding to the presentation $`𝒫=xxx=x`$ with base $`x`$ is the famous R. Thompson’s group $`F`$, which has the following presentation: $$x_0,x_1,\mathrm{}x_j^{x_i}=x_{j+1}(j>i).$$ (see \[14, Example 6.4\]. This group has several interesting propertries and is studied by mathematicians working in different areas of mathematics ($`\lambda `$-calculus, functional analysis, homological algebra, homotopy theory, group theory). It was discovered by R. Thompson in 1965, and was rediscovered later by other authors. presents a survey of results about $`F`$. Since $`F`$ is one of the most important diagram groups, and since we are going to present some new results about it in this paper, let us recall some known properties of this group. These propertries can be found in , and . 1. The group $`F`$ is isomorphic to the group of all increasing continuous piecewise linear maps of the interval $`[0,1]`$ onto itself such that the singularities occur at finitely many dyadic points (points of the form $`m/2^n`$) and all slopes are powers of 2. The group operation is the composition of functions (we shall write function symbols to the right of the argument). 2. In the previous paragraph, one can replace the interval $`[0,1]`$ by $`[0,+\mathrm{}]`$, adding the assumption that the slop on $`+\mathrm{}`$ is $`1`$. The resulting group is also isomorphic to $`F`$. 3. $`F`$ does not satisfy any non-trivial identity. 4. $`F`$ does not contain any free non-abelian subgroups. Every subgroup of $`F`$ either is abelian or contains an infinite direct power of $`𝐙`$. 5. $`F`$ is finitely presented, it has a presentation with two generators and two defining relations. The word problem and the conjugacy problem are solvable in $`F`$. It has a polynomial isoperimetric function . There exists a clear connection between representation of elements of $`F`$ by diagrams and normal form of elements in $`F`$. Recall that every element in $`F`$ is uniquely representable in the following form: $$x_{i_1}^{s_1}\mathrm{}x_{i_m}^{s_m}x_{j_n}^{t_n}\mathrm{}x_{j_1}^{t_1},$$ (1) where $`i_1\mathrm{}i_mj_n\mathrm{}j_1`$; $`s_1,\mathrm{},s_m,t_1,\mathrm{}t_n0`$, and if $`x_i`$ and $`x_i^1`$ occur in (1) for some $`i0`$ then either $`x_{i+1}`$ or $`x_{i+1}^1`$ also occurs in (1). This form is called the normal form of elements in $`F`$. (Note that in , we constructed another normal form for elements of $`F`$, our normal forms are locally testable.) Let us show how given an $`(x,x)`$-diagram over $`𝒫=xxx=x`$ one can get the normal form of the element represented by this diagram. We simply describe the procedure providing no proofs. Details can be deduced from \[14, Example 6.4\]. ###### Example 2 Every diagram $`\mathrm{\Delta }`$ over $`𝒫`$ can be divided by its longest positive path from its intitial vertex to its terminal vertex into two parts, positive and negative, denoted by $`\mathrm{\Delta }^+`$ and $`\mathrm{\Delta }^{}`$, respectively. So $`\mathrm{\Delta }=\mathrm{\Delta }^+\mathrm{\Delta }^{}`$. It is easy to prove by induction on the number of cells that all cells in $`\mathrm{\Delta }^+`$ are $`(x,x^2)`$-cells, all cells in $`\mathrm{\Delta }^{}`$ are $`(x^2,x)`$-cells. This implies that the numbers of cells in $`\mathrm{\Delta }^+`$ and in $`\mathrm{\Delta }^{}`$ are the same. Denote this number by $`k`$. Let us number the cells of $`\mathrm{\Delta }^+`$ by numbers from $`1`$ to $`k`$ by taking every time the “rightmost” cell, that is, the cell which is to the right of any other cell attached to the bottom path of the diagram formed by the previous cells. The first cell is attached to the top path of $`\mathrm{\Delta }^+`$ ($`=\text{top}(\mathrm{\Delta })`$). The $`i`$th cell in this sequence of cells corresponds to an edge of the graph $`\mathrm{\Gamma }(𝒫)`$, which has the form $`(x^\mathrm{}_i,xx^2,x^{r_i})`$, where $`\mathrm{}_i`$ ($`r_i`$) is the length of the path from the initial (resp. terminal) vertex of the diagram (resp. the cell) to the initial (resp. terminal) vertex of the cell (resp. the diagram), and contained in the bottom path of the diagram formed by the first $`i1`$ cells. If $`\mathrm{}_i=0`$ then we label this cell by 1. If $`\mathrm{}_i0`$ then we label this cell by the element $`x_{r_i}`$ of $`F`$. Multiplying the labells of all cells, we get the “positive” part of the normal form. For example, the diagram on the next picture the positive part is equal to $`x_0x_2^2x_4x_5`$ (cells 1 and 3 were labelled by the identity element). In order to find the “negative” part of the normal form, consider $`(\mathrm{\Delta }^{})^1`$, number its cells as above and label them as above. In our example, we get the word $`x_1x_3^2x_4`$ (cells 1, 2, 4 are labelled by 1). Thus the “negative” part of the normal form is $`(x_1x_3^2x_4)^1`$, and it remains to multiply the positive and negative parts. In our example, the normal form is $`x_0x_2^2x_4x_5x_4^1x_3^2x_1^1`$. Diagrams presented below are generators $`x_0`$, $`x_1`$ of the group $`F`$. They generate the whole $`F`$. One can ask several natural general questions about diagram groups. Which groups are diagram groups? Which groups are representable by diagrams (are subgroups of diagram groups)? Do these two classes coinside? How to compute a diagram group over a given presentation $`𝒫`$ and with a given base? There exists a well developed technology for computing diagram groups. The starting point for computing diagram groups is Kilibarda’s theorem about fundamental groups of Squier complexes. In order to formulate this important result, let us define the structure of a 2-complex on the graph $`\mathrm{\Gamma }(𝒫)`$ for any presentation $`𝒫`$. First notice that although for every path in $`\mathrm{\Gamma }(𝒫)`$ there exists a unique diagram associated with this path, the same diagram can be associated with many paths. Consider the following typical case: Let $`𝒫=\mathrm{\Sigma }`$ where $``$ contains two defining relations $`\mathrm{}_i=r_i`$ ($`i=1,2`$), and let $`u`$, $`v`$, $`z`$ be arbitrary words in $`\mathrm{\Sigma }^{}`$. Consider the following paths in $`\mathrm{\Gamma }(𝒫)`$: $$(u,\mathrm{}_1r_1,z\mathrm{}_2v)(ur_1z,\mathrm{}_2r_2,v),$$ (2) $$(u\mathrm{}_1z,\mathrm{}_2r_2,v)(u,\mathrm{}_1r_1,zr_2v).$$ (3) It is easy to see that $`(u\mathrm{}_1z\mathrm{}_2v,ur_1zr_2v)`$-diagrams corresponding to these paths are equal. This diagram is shown on the following picture: This situation hints to a homotopy relation on the set of paths in the graph $`\mathrm{\Gamma }(𝒫)`$: paths (2) and (3) should be called homotopic. In order to define the homotopy relation we need the structure of a 2-complex on $`\mathrm{\Gamma }(𝒫)`$. For every 5-tuple $`(u,\mathrm{}_1=r_1,z,\mathrm{}_2=r_2,v)`$, where $`u,v,z\mathrm{\Sigma }^{}`$, $`(\mathrm{}_1=r_1),(\mathrm{}_2=r_2)`$ we have a $`2`$-cell whose defining path is $`p_1p_2^1`$ where $`p_1`$, $`p_2`$ are the paths (2) and (3), respectively. The resulting $`2`$-complex is called the Squier complex of the semigroup presentation $`𝒫`$. It is denoted by $`𝒦(𝒫)`$. It was implicitely defined by Squier in . The same complex was independently constructed by Kilibarda and Pride . The important role of this complex is justified by the fact that equal diagrams over $`𝒫`$ correspond to homotopic paths in $`𝒦(𝒫)`$. The following Kilibarda’s theorem plays an important role in this paper: The diagram group $`𝒟(𝒫,w)`$ is isomorphic to the fundamental group $`\pi _1(𝒦,w)`$ of the Squier complex $`𝒦=𝒦(𝒫)`$. ## 2 Diagram Product of Groups In , we considered several group-theoretical operations such that the class of diagram groups is closed under them. These operations were: finite direct products (result due to Kilibarda ), any free products, and also some special operation $``$ which we used for constructing an example of a diagram group that was finitely generated but not finitely presented. In this Section we introduce a quite general operation on groups, the diagram product. We show that the class of diagram groups is closed under this operation. All the above listed constructions are partial cases of this new operation. We also consider some concrete applications of this construction. They will be essentially used in the later Sections. Let us recall the definition of a graph. A graph (in the sense of Serre ) is an ordered tuple $`\mathrm{\Gamma }=V,E,^1,\iota ,\tau `$ where $`V`$, $`E`$ are disjoint sets, <sup>-1</sup> is an involution on $`E`$, $`\iota `$, $`\tau `$ are mappings from $`E`$ to $`V`$. The following axioms hold: * $`e^1e`$ for any $`eE`$; * $`\iota (e^1)=\tau (e)`$, $`\tau (e^1)=\iota (e)`$. Elements of the sets $`V`$ and $`E`$ are called vertices and edges of the graph respectively. If $`eE`$, then $`\iota (e)`$ is called the initial vertex of the edge $`e`$, and $`\tau (e)`$ is called the terminal vertex of the edge $`e`$. A path on the graph $`\mathrm{\Gamma }`$ is either a vertex, or a nonempty sequence of edges $`e_1`$, $`e_2`$, …, $`e_n`$ such that $`\tau (e_i)=\iota (e_{i+1})`$ for each $`i=1,\mathrm{},n1`$. Usually a path is written in the form $`p=e_1e_2\mathrm{}e_n`$. If a path $`p`$ consists of a vertex $`v`$, then it is called an empty path and we denote it by $`1_v`$. If $`p=e_1e_2\mathrm{}e_n`$ is a path, then the inverse path $`p^1`$ is the path $`e_n^1e_{n1}^1\mathrm{}e_1^1`$. An empty path coincides with its inverse. A path $`p`$ is called closed whenever $`\iota (p)=\tau (p)`$. An orientation on the graph $`\mathrm{\Gamma }`$ is a subset $`E^+`$ of the set $`E`$ of all edges such that, for any edge $`eE`$, there is exactly one of the edges $`e`$, $`e^1`$ that belongs to $`E^+`$. The edges in $`E^+`$ are called positive and the edges in $`E^{}=EE^+`$ are called negative. A path on an oriented graph is called positive whenever it involves positive edges only. (An empty path is always positive.) For any path $`p`$, there are defined its initial vertex $`\iota (p)`$ and its terminal vertex $`\tau (p)`$: if $`p=1_v`$, then $`\iota (p)=\tau (p)=v`$; if $`p=e_1\mathrm{}e_n`$, then $`\iota (p)=\iota (e_1)`$, $`\tau (p)=\tau (e_n)`$. For any two paths $`p`$, $`q`$ such that $`\tau (p)=\iota (q)`$, one can naturally define a product $`pq`$ of the paths $`p`$ and $`q`$: for $`p=e_1\mathrm{}e_n`$, $`q=f_1\mathrm{}f_m`$ we put $`pq=e_1\mathrm{}e_nf_1\mathrm{}f_m`$. If $`p`$ ($`q`$) is empty, then $`pq=q`$ ($`pq=p`$). An orineted graph is by definition a graph $`\mathrm{\Gamma }`$ with a fixed orientation $`E^+`$. It is clear that any graph admits an orientation. The concept of a graph of groups will play an important role. Let us have an oriented graph $`\mathrm{\Gamma }`$, where $`E^+`$ is the set of positive edges. We say that a graph of groups structure on the graph $`\mathrm{\Gamma }`$ is given whenever to each edge $`eE^+`$ we assign a group $`G_e`$, to each vertex $`vV`$ we assign a group $`G_v`$, and we fix embeddings $`\iota _e:G_eG_{\iota (e)}`$, $`\tau _e:G_eG_{\tau (e)}`$ for any $`eE^+`$. In the construction described below, we will have a 2-complex structure on $`\mathrm{\Gamma }`$ together with the graphs of groups structure. This means that we have a set $`F`$ which is disjoint from $`V`$ and $`E`$. This set is called the set of $`2`$-cells. We also have a mapping that assigns a closed path in $`\mathrm{\Gamma }`$ to each element in $`F`$. This path is called the defining path of the $`2`$-cell. Given a 2-complex, we define the homotopy relation on the set of paths in a standard way. Also one can define the concept of the fundamental group of $`𝒦`$ with basepoint $`w`$. We denote this group by $`\pi _1(𝒦,w)`$. We will consider $`2`$-complexes that have a graph of groups structure on their $`1`$-skeletons. We shall call such structures $`2`$-complexes of groups. The concept of a $`2`$-complex of groups already exists and it is used widely in many papers (see ). Every 2-complex of groups in our sense is a 2-complex of groups in the sense of , but not vice versa. (In general, a $`2`$-complex of groups is a structure that has not only vertex groups $`G_v`$ ($`vV`$) and edge groups $`G_e`$ ($`eE`$) but also cell groups of the form $`G_f`$ ($`fF`$) that are assigned to $`2`$-cells. In our situation all the cell groups $`G_f`$ are trivial.) So, let $`𝒢`$ be a $`2`$-complex of groups. Now we define the fundamental group of $`𝒢`$. One can define it in different ways. We will define it as a fundamental group of an ordinary $`2`$-complex $`{}_{𝒦}{}^{}(𝒢)`$ with a basepoint. Here is the description of the complex $`{}_{𝒦}{}^{}(𝒢)`$. We add new edges and new $`2`$-cells to the $`2`$-complex $`𝒦`$. For any vertex $`vV`$ and for any element $`gG_v`$ we add an edge denoted by $`g_v`$ that has $`v`$ as both initial and terminal vertex. The new $`2`$-cells are of two types. $`2`$-cells of the first type have defining paths $`g_vh_v(gh)_v^1`$ for any vertex $`vV`$ and for any elements $`g,hG_v`$. $`2`$-cells of the second type have defining paths $`e^1g_{\iota (e)}eh_{\tau (e)}^1`$, where $`g=\iota _e(x)`$, $`h=\tau _e(x)`$, $`xG_e`$, $`eE^+`$. (The $`2`$-cells of the second type correspond to all pairs of the form $`(e,x)`$, where $`eE^+`$, $`xG_e`$.) Recall that $`\iota _e`$, $`\tau _e`$ are embeddings of the group $`G_e`$ into the groups $`G_\iota (e)`$, $`G_{\tau (e)}`$, respectively. The $`2`$-complex obtained from $`𝒦`$ by adding new $`2`$-cells will be denoted by $`{}_{𝒦}{}^{}(𝒢)`$. For any vertex $`vV`$, the fundamental group $`\pi _1({}_{𝒦}{}^{}(𝒢),v)`$ will be called the fundamental group of $`2`$-complex of groups $`𝒢`$ with basepoint $`v`$. It will be denoted by $`\pi _1(𝒢,v)`$. A standard way to compute the fundamental group of a $`2`$-complex (see ) can be also applied to compute the fundamental group of a $`2`$-complex of groups. Let we have a structure of a graph of groups $`𝒢`$ on the $`1`$-skeleton of a $`2`$-complex $`𝒦`$. Consider the connected component $`𝒦_w`$ of $`𝒦`$ that contains vertex $`w`$ and let us choose some maximal subtree $`𝒯_w`$ in this component. It will also be a maximal subtree of the connected component of the new $`2`$-complex $`{}_{𝒦}{}^{}(𝒢)`$ that contains vertex $`w`$. Let us take the set of positive edges $`E_w^+`$ of this connected component together with elements of the form $`g_v`$, $`vV_w`$, $`gG_v`$, where $`V_w`$ is the set of vertices of this connected component. The union of these sets is the set of generators of the fundamental group. For defining relations, we take all relations of the form $`e=1`$, where $`e`$ belongs to the tree $`𝒯`$, and also all relations of the form $`r=1`$, where $`r`$ is the defining path of any $`2`$-cell of the $`2`$-complex $`{}_{𝒦}{}^{}(𝒢)`$. The group given by the described presentation is isomorphic to the fundamental group $`\pi _1(𝒢,w)`$ of our $`2`$-complex of groups. Let us give one more equivalent description. It is clear that the fundamental group $`\pi _1(𝒦,w)`$ with basepoint $`w`$ of the original $`2`$-complex $`𝒦`$ can be computed in the same way, choosing $`𝒯_w`$ as a maximal subtree of the connected component of $`w`$. Now each edge $`e`$ of this component uniquely defines an element in $`\pi _1(𝒦,w)`$, namely, the equivalence class of the path $`p_{\iota (e)}ep_{\tau (e)}^1`$, where $`p_v`$ denotes the geodesic path from $`w`$ to $`v`$ in the subtree $`𝒯_w`$. Then $$\pi _1(𝒢,w)\underset{v}{}G_v\pi _1(𝒦,w)/𝒩,$$ (4) where the free product of groups $`G_v`$ is taken over all vertices $`v`$ of the connected component of $`𝒦`$ that contains $`w`$, and $`𝒩`$ is the normal closure of the following relations: $$g_{\iota (e)}^e=h_{\tau (e)}\text{ for every }eE_w^+,xG_e,\text{ where }g=\iota _e(x),h=\tau _e(x).$$ (5) This description will be often used below. Let us give one of the main definitions. ###### Definition 3 Let $`X`$ be an alphabet, let $`H_x`$ ($`xX`$) be an arbitrary family of groups and let $`𝒬=X𝒮`$ be a semigroup presentation, $`wX^+`$. Consider the Squier complex $`𝒦=𝒦(𝒬)`$ and introduce the structure of graph of groups on its $`1`$-skeleton in the following way. Let $`E^+`$ be the set of all positive edges of $`𝒦`$, that is the set of triples of the form $`e=(u,st,v)`$, where $`u,vX^{}`$, $`(s=t)𝒮`$. For any word $`z=x_1x_2\mathrm{}x_n`$, where $`x_1,x_2,\mathrm{},x_nX`$, let $$H_u=H_{x_1}\times H_{x_2}\times \mathrm{}H_{x_n};$$ if $`u`$ is empty, then $`H_u=1`$. For any vertex $`u𝒦`$ let $`G_u=H_u`$; for any edge $`e=(u,st,v)E^+`$, where $`u,vX^{}`$, $`(s=t)𝒮`$, let $`G_e=H_u\times H_v`$. We have the natural embeddings $`\iota _e:G_eG_{\iota (e)}`$ as the embedding $`H_u\times H_vH_{usv}=H_u\times H_s\times H_t`$ and $`\tau _e:G_eG_{\tau (e)}`$ as an embedding $`H_u\times H_vH_{utv}=H_u\times H_t\times H_t`$. This gives us a $`2`$-complex of groups, which will be denoted by $`𝒦_H`$. The fundamental group $`\pi _1(𝒦_H,w)`$ of this $`2`$-complex of groups with basepoint $`w`$ is called the diagram product of the family $`H_X=\{H_x(xX)\}`$ of groups over the presentation $`𝒬=X𝒮`$ with base $`w`$. It will be denoted by $`𝒟(H_X;𝒮,w)`$. It is easy to see that the diagram product $`𝒟(H_X;𝒮,w)`$ coincides with the diagram group $`𝒟(𝒬,w)`$ in the case when the groups $`H_x`$ are trivial for all $`xX`$. The main result about this construction is that the diagram product of any family of diagram groups over a semigroup presentation is again a diagram group. Let us formulate this result in its general form giving the description of the presentation, for which the corresponding diagram group is the diagram product. ###### Theorem 4 Let $`𝒬=X𝒮`$ be a semigroup presentation, $`wX^+`$. To each $`xX`$ we assign a diagram group $`G_x=𝒟(𝒫_x,w_x)`$, where $`𝒫_x=\mathrm{\Sigma }_x_x`$ are semigroup presentations, $`w_x\mathrm{\Sigma }_x^+`$ $`(xX)`$. Let $`A=\{a_xxX\}`$ be some alphabet. Assume that the alphabets $`X`$, $`A`$, $`\mathrm{\Sigma }_x`$ $`(xX)`$ are disjoint. Let $$\mathrm{\Sigma }=\underset{xX}{}\mathrm{\Sigma }_x,=\underset{xX}{}_x,𝒲=\underset{xX}{}\{x=a_xw_xa_x\}.$$ Consider the presentation $$𝒫=X\mathrm{\Sigma }A𝒮𝒲.$$ We claim that the diagram group $`𝒟(𝒫,w)`$ is isomorphic to the diagram product $`𝒟(G_X;𝒮,w)`$ of the family $`G_X=\{G_x(xX)\}`$ of groups over the presentation $`𝒬=X𝒮`$ with base $`w`$. In particular, the diagram product of any family of groups over a semigroup presentation is a diagram group. Let us consider a geometric description of this construction. In the above notation, each group $`G_x`$ is isomorphic to the diagram group over the presentation $`\widehat{𝒫}_x`$ that is obtained from $`𝒫_x`$ by adding letters $`x`$, $`a_x`$ to the alphabet $`\mathrm{\Sigma }_x`$ and adding new relation $`x=a_xw_xa_x`$ to the set $`_x`$. We have $`G_x𝒟(\widehat{𝒫}_x,x)`$. Let us take any $`(w,w)`$-diagram over $`𝒬`$ and consider some of its edges. Let $`xX`$ be its label. We can cut the diagram along this edge and insert any $`(x,x)`$-diagram over $`\widehat{𝒫}_x`$ in the resulting hole. We can do this with all edges of the diagram. (If we insert trivial $`(x,x)`$-diagram, then nothing changes.) After these transformations, we obtain some $`(w,w)`$-diagram over $`𝒫`$. One can show that diagrams obtained in this way form the whole group $`𝒟(𝒫,w)`$. From this point of view, diagrams that represent elements in $`𝒟(𝒫,w)`$ are obtained from $`(w,w)`$-diagrams over $`𝒬`$ by insertions of $`(x,x)`$-diagrams, which represent elements in $`G_x`$. Proof. Let us construct the Squier complex for the presentation $`𝒫`$. By Kilibarda’s Theorem, the fundamental group of this complex with basepoint $`w`$ is isomorphic to the diagram group $`𝒟(𝒫,w)`$. Our goal is to transform the Squier complex into some new $`2`$-complex with the same fundamental group. We will need to check that it will be isomorphic to the fundamental group of a certain $`2`$-complex of groups, that is, to the diagram products of our groups. Let $`𝒦_w(𝒫)`$ be the connected component of the Squier complex over $`𝒫`$, which contains the vertex $`w`$. A vertex $`v`$ in the same component is an arbitrary word that equals $`w`$ modulo $`𝒫`$. This word can be uniquely decomposed into the product $`v=v_1\mathrm{}v_\mu `$ ($`\mu =\mu (v)`$) of subwords $`v_1`$, …, $`v_\mu `$ in such a way that each of them will be either a letter in $`X`$, or a word of the form $`a_xua_x`$, where $`u`$ is a word over $`\mathrm{\Sigma }_x`$ that equals $`w_x`$ modulo $`𝒫_x`$. The words $`v_1`$, …, $`v_\mu `$ will be called the factors of the word $`v`$. To each factor $`v_i`$ ($`1i\mu `$), we assign a letter in the alphabet $`X`$. This letter will be denoted by $`\pi (v_i)`$. If $`v_iX`$, then we put $`\pi (v_i)=v_i`$ and if $`v_i`$ has a form $`a_xua_x`$, where $`u`$ equals $`w_x`$ modulo $`𝒫_x`$, then we put $`\pi (v_i)=x`$. Let us extend the function $`\pi `$, putting by definition $`\pi (v)=\pi (v_1)\mathrm{}\pi (v_\mu )`$ for any word $`v`$ that equals $`w`$ modulo $`𝒫`$. The function $`\pi `$ will be called the projection. There is a natural two-sided action of the free monoid $`M=\left(X\mathrm{\Sigma }A\right)^{}`$ on the Squier complex. It can be defined by the following rule: for any $`m_1,m_2M`$ and for any vertex $`v`$ of the Squier complex, let $`m_1vm_2`$ be the vertex $`m_1vm_2`$ and let $`m_1em_2`$ be the edge $`(m_1u,pq,vm_2)`$, for any edge $`e=(u,pq,v)`$. The images of a given subgraph of $`𝒦(𝒫)`$ under this action will be called the shifts of this subgraph. Let $`\widehat{}_x=_x\{x=a_xw_xa_x\}`$. We introduce a presentation $`\widehat{𝒫}_x=\mathrm{\Sigma }_x,x,a_x\widehat{}_x`$ for all $`xX`$. It is clear that $`𝒟(\widehat{𝒫}_x,x)G_x`$. Let $`v`$ be a vertex of the complex $`𝒦_w(𝒫)`$ and let $`v=v_1\mathrm{}v_\mu `$ be the decomposition of the word $`v`$ into factors. It is obvious that for any $`1i\mu `$, the letter $`\pi (v_i)=x`$ is equal to the word $`v_i`$ modulo $`\widehat{𝒫}_x`$. Also it is clear that $`v`$ equals $`\pi (v)`$ modulo $`𝒫`$. Note that Squier complexes for presentations $`\widehat{𝒫}_x`$ and their shifts can be regarded as subcomplexes of the Squier complex of $`𝒫`$. For each $`xX`$ we choose a maximal subtree $`𝒯_x`$ in the connected component $`𝒦(\widehat{𝒫}_x,x)`$ of the Squier complex of the presentation $`\widehat{𝒫}_x`$ that contains the vertex $`x`$. Let us also choose a maximal subtree $`𝒯_Q`$ in the connected component $`𝒦(𝒬,w)`$ of the Squier complex of the presentation $`𝒬`$ that contains vertex $`w`$. Let $`v`$ be an arbitrary vertex of the complex $`𝒦_w(𝒫)`$ and let $`v=v_1\mathrm{}v_\mu `$ be the decomposition of $`v`$ into factors. Let $`x_i=\pi (v_i)`$ ($`1i\mu `$). Consider the following subgraphs of $`𝒦_w(𝒫)`$: $$1𝒯_{x_1}x_2\mathrm{}x_\mu ,v_1𝒯_{x_2}x_3\mathrm{}x_\mu ,\mathrm{},v_1\mathrm{}v_{\mu 1}𝒯_{x_\mu }1.$$ (6) Now consider the subgraph $`𝒯`$ of $`𝒦_w(𝒫)`$ that is a union of subgraphs (6) for all vertices $`v`$ that are equal to $`w`$ modulo $`𝒫`$, together with the subgraph $`𝒯_Q`$. Let us prove that $`𝒯`$ is a maximal subtree of $`𝒦_w(𝒫)`$. First of all we will establish that $`𝒯`$ is a connected subgraph that contains all vertices of $`𝒦_w(𝒫)`$, that is, for any vertex $`v`$ of our component, we will find a path in $`𝒯`$ from $`w`$ to $`v`$. Let $`v=v_1\mathrm{}v_\mu `$ be the decomposition of $`v`$ into a product of factors. By $`p`$ we denote the geodesic path in $`𝒯_Q`$ from $`w`$ to $`\pi (v)=x_1\mathrm{}x_\mu `$, where $`x_i=\pi (v_i)`$ for all $`i`$ from $`1`$ to $`\mu `$. For each $`i`$, let $`p_i`$ be the geodesic path from $`x_i`$ to $`v_i`$ in the graph $`𝒯_{x_i}`$. For each $`i`$ from $`1`$ to $`\mu `$ we consider the path $`\stackrel{~}{p}_i=v_1\mathrm{}v_{i1}p_ix_{i+1}\mathrm{}x_\mu `$. Obviously, it connects vertices $`v_1\mathrm{}v_{i1}x_i\mathrm{}x_\mu `$ and $`v_1\mathrm{}v_ix_{i+1}\mathrm{}x_\mu `$ in the graph $`𝒯`$. The product $`p\stackrel{~}{p}_1\mathrm{}\stackrel{~}{p}_\mu `$ is a path in $`𝒯`$ from $`w`$ to $`v_1\mathrm{}v_\mu =v`$. Now let us prove that $`𝒯`$ has no nontrivial cycles. We argue by contradiction. Suppose that a nontrivial cycle exists. If it does not consist of edges that belong to subgraphs of the form (6), then it has an edge from $`𝒯_Q`$. Since $`𝒯_Q`$ has no nontrivial cycles, our cycle must contain edges from subgraphs of the form (6). Hence our cycle has a nontrivial cyclic subpath $`\rho `$ that is a loop at some vertex in $`𝒦_w(𝒬)`$ and all its edges are from subgraphs of the form (6). Let us make a simple but important observation: the endpoints of each edge that belong to any shift of the subcomplex $`𝒦(\widehat{𝒫}_x)`$ ($`xX`$), have equal projection. Indeed, applying relations of the form $`x=a_xw_xa_x`$ does not change the projection, and applying relations from $`_x`$ occurs within a factor of the form $`a_xua_x`$, where $`u`$ is a word over $`\mathrm{\Sigma }_x`$. This also does not change the projection, (in the last case one can see the role of the auxiliary alphabet $`A`$). Thus projections of all vertices of the cyclic path $`\rho `$ coincide, that is, $`\rho `$ is a nontrivial cycle in $`𝒯`$ that consists of edges from subgraphs of the form (6). So in any case there is a nontrivial cycle $`\rho `$ with the above property. Without loss of generality, one can assume that $`\rho `$ does not contain occurrences of adjacent edges that are mutually inverse. Let $`v=v_1\mathrm{}v_\mu `$ be the decomposition of $`v`$ into factors, where $`\rho `$ is the loop at $`v`$. Each edge of the path $`\rho `$ touches exactly one of the factors, as we could see above. Let $`j`$ be the greatest number such that an edge of $`\rho `$ touches $`j`$th factor. Let $`x_i=\pi (v_i)`$ ($`1i\mu `$). It follows from the structure of subgraphs (6) that $`v_i=x_iX`$ for all $`j<i\mu `$. Since the $`j`$th factor occurs in the process of application of relations from $`\widehat{}_j`$ to it, the path $`\rho `$ or one of its cyclic shifts has a maximal subpath $`\rho ^{}`$ that consists of edges that touch the $`j`$th factor only. Let $`v^{}`$ and $`v^{\prime \prime }`$ be the initial and the terminal points of $`\rho ^{}`$ respectively. We claim that $`v^{}=v^{\prime \prime }`$. Suppose this is not true. It is clear that $`v^{}`$ and $`v^{\prime \prime }`$ differ by the $`j`$th factor only and so one can say that the $`j`$th factor is not equal to $`x_j`$ either in $`v^{}`$ or in $`v^{\prime \prime }`$. Assume that the $`j`$th factor of $`v^{\prime \prime }`$ is not equal $`x_j`$. By our assumption, $`\rho ^{}`$ has fewer edges than $`\rho `$ (otherwise $`v^{}=v^{\prime \prime }`$ automatically). So there is an edge $`e`$ such that $`\rho ^{}e`$ is a subpath of some cyclic shift of the path $`\rho `$. Since $`\rho ^{}`$ was chosen maximal, the edge $`e`$ does not touch the $`j`$th factor. It also cannot touch a factor with a number greater than $`j`$ because $`j`$ is maximal with this property. But it also cannot touch a factor with a number less than $`j`$ because it belongs to a subgraph of the form (6), and the $`j`$th factor of the initial point of $`e`$ does not belong to $`X`$, a contradiction. So $`\rho ^{}`$ is a nontrivial cycle that belongs to a shift of the tree $`𝒯_{x_j}`$. However, this is impossible since a shift of a tree is a tree itself. This contradiction shows that $`𝒯`$ has no nontrivial cycles. Applying what we have said above, we conclude that $`𝒯`$ is a maximal subtree in $`𝒦_w(𝒫)`$. Now we need to calculate the fundamental group $`G=\pi _1(𝒦_w(𝒫))`$ by using the maximal subtree $`𝒯`$. All edges of the complex $`𝒦_w(𝒫)`$ are regarded as elements of the group $`G`$, and the edges from $`𝒯`$ equal the identity in $`G`$. Paths in this complex, regarded as products of edges, are just elements of the group $`G`$. To understand how the other relations in $`G`$ look like, we need to describe the $`2`$-cells in $`𝒦_w(𝒫)`$. First of all let us mention some important property. Recall that $`M`$ is the free monoid over the alphabet of the presentation $`𝒫`$, and $`M`$ acts both from the left and from the right on the complex $`𝒦(𝒫)`$. Let $`s`$, $`t`$, $`u`$, $`v`$ be elements in $`M`$, each decomposed into the product of factors, and let $`s`$ equals $`t`$ modulo $`𝒫`$, $`usv`$ equals $`w`$ modulo $`𝒫`$. Let us take an arbitrary path $`p`$ in $`𝒦(𝒫)`$ that connects vertices $`s`$ and $`t`$. It is clear that the paths $`upv`$ and $`\pi (u)p\pi (v)`$ belong to $`𝒦_w(𝒫)`$. We claim that the following equality holds $$upv=\pi (u)p\pi (v)$$ (7) in the group $`G`$. This is what we are going to prove. To prove that, one can consider the contours of the corresponding $`2`$-cells as words in the generators of the group $`G`$ which are equal to the identity in $`G`$. However, we think that one can check equality (7) easier, using the Kilibarda Theorem. Namely, to prove the equality (7), it suffices to use the fact that $`G`$ is isomorphic to the diagram group over $`𝒫`$ with base $`\pi (usv)`$. So let us find the diagrams over $`𝒫`$ that represent the elements in $`G`$ from both sides of equality (7), and then let us check that the diagrams are equal. Let $`x=x_1\mathrm{}x_m`$ be any word in $`M`$ decomposed into the product of its factors. For any $`i`$ from $`1`$ to $`m`$ let $`q_i`$ be the geodesic path from $`\pi (x_i)`$ to $`x_i`$ in the tree $`𝒯_{\pi (x_i)}`$. Then $$p_x=(1q_1\pi (x_2\mathrm{}x_m))(x_1q_2\pi (x_3\mathrm{}x_m))\mathrm{}(x_1\mathrm{}x_{m1}q_m1)$$ is a path from $`\pi (x)`$ to $`x`$. Let $`\mathrm{\Delta }_x`$ be the diagram represented by it. Let us consider such paths and diagrams for all $`x\{s,t,u,v\}`$. Also let $`q`$ be the geodesic path from $`\pi (usv)`$ to $`\pi (utv)`$ in the tree $`𝒯_Q`$, and let $`\mathrm{\Delta }`$, $`\mathrm{\Psi }`$ be the diagrams represented by $`p`$, $`q`$, respectively. To find the spherical diagram with base $`\pi (usv)`$ represented by the path $`upv`$ (via the isomorphism of the diagram group and the group $`G`$), one needs to concatenate three diagrams: the $`\mathrm{\Delta }_1`$ that corresponds to the path in the tree $`𝒯`$ from $`\pi (usv)`$ to $`usv`$, the diagram $`\mathrm{\Delta }_2=\epsilon (u)+\mathrm{\Delta }+\epsilon (v)`$ (that corresponds to the path $`upv`$), and the diagram $`\mathrm{\Delta }_3`$ that corresponds to the path in the tree $`𝒯`$ from $`utv`$ to $`\pi (usv)`$. So consider the path $`(1p_u\pi (sv))(up_s\pi (v))(usp_v1)`$. It follows from the description of subgraphs (6) that this path is contained in $`𝒯`$. It corresponds to the diagram $$\mathrm{\Delta }_1=(\mathrm{\Delta }_u+\epsilon (\pi (sv)))(\epsilon (u)+\mathrm{\Delta }_s+\epsilon (\pi (v)))(\epsilon (us)+\mathrm{\Delta }_v).$$ Further, the path $`(utp_v^11)(up_t^1\pi (v))(1p_u^1\pi (tv))`$ is contained in $`𝒯`$. Multiplying it by the path $`q^1`$ on the right, we obtain the path in $`𝒯`$ from $`utv`$ to $`\pi (usv)`$. This path is represented by the diagram $$\mathrm{\Delta }_3=(\epsilon (ut)+\mathrm{\Delta }_v^1)(\epsilon (u)+\mathrm{\Delta }_t^1+\epsilon (\pi (v)))(\mathrm{\Delta }_u^1+\epsilon (\pi (tv)))\mathrm{\Psi }^1.$$ Let us now multiply the diagrams $`\mathrm{\Delta }_1`$, $`\mathrm{\Delta }_2`$ and $`\mathrm{\Delta }_3`$. It is easy to see that the subdiagram $`\mathrm{\Delta }_u`$ cancels with $`\mathrm{\Delta }_u^1`$ in this product, and $`\mathrm{\Delta }_v`$ cancels with $`\mathrm{\Delta }_v^1`$ (see the picture below). After cancelling $`\mathrm{\Delta }_u`$ and $`\mathrm{\Delta }_u^1`$, $`\mathrm{\Delta }_v`$ and $`\mathrm{\Delta }_v^1`$, we obtain a diagram that is a product $$(\epsilon (\pi (u))+\mathrm{\Delta }_s+\epsilon (\pi (v)))(\epsilon (\pi (u))+\mathrm{\Delta }+\epsilon (\pi (v)))(\epsilon (\pi (u))+\mathrm{\Delta }_t^1+\epsilon (\pi (v)))\mathrm{\Psi }^1.$$ (8) Repeating the arguments of the above paragraph for the path $`\pi (u)p\pi (v)`$, it is easy to see that this path is represented by the diagram (8) in the diagram group over $`𝒫`$ with the base $`\pi (usv)`$. This proves the equality (7). Let us consider an arbitrary edge $`(u,st,v)`$ of the complex $`𝒦_w(𝒫)`$. Let $`(s=t)𝒮`$. The words $`u`$, $`v`$ can be decomposed into products of factors, and the equality $`(u,st,v)=(\pi (u),st,\pi (v))`$ holds in $`G`$. The right-hand side of this equality can be regarded as an element of the group $`\pi _1(𝒦_w(𝒬))`$, where $`𝒯_Q`$ is the maximal subtree in $`𝒦_w(𝒬)`$. Now let $`(s=t)𝒮`$. In this case there exists an element $`xX`$ and words $`u_1`$, $`v_1`$, $`u_2`$, $`v_2`$ such that $`u=u_1u_2`$, $`v=v_2v_1`$, where $`u_1`$ ($`v_1`$) is the maximal prefix (suffix) of the word $`u`$ (resp. $`v`$) that can be decomposed into a product of factors. Here $`u_2sv_2`$ equals $`x`$ modulo $`𝒫`$. Then by (7) we have the equality $`e=(u_1u_2,st,v_2v_1)=\pi (u_1)(u_2,st,v_2)\pi (v_1)`$. For any $`xX`$, let us consider the fundamental group $`\pi _1(𝒦(\widehat{𝒫}_x),x)G_x`$ that can be calculated using the maximal subtree $`𝒯_x`$ in the connected component of the Squier complex over $`𝒫_x`$ that contains $`x`$. The edges of this component will be thus the elements of a group isomorphic to $`G_x`$ so the edge $`e`$ will belong to an isomorphic copy of this group that is generated by edges obtained as a result of shifts. Namely, let $`U,VX^{}`$, $`xX`$. We consider the group denoted by $`UG_xV`$. It is generated by edges of the form $`UfV`$, where $`f`$ runs over edges that generate the group $`\pi _1(𝒦(\widehat{𝒫}_x),x)G_x`$. In this sense, the edge $`e`$ belongs to the group $`\pi (u_1)G_x\pi (v_1)`$. The argument of this paragraph can be summarized as follows: the group $`G`$ is generated by the subgroup $`\pi _1(𝒦(𝒬),w)𝒟(𝒬,w)`$ and groups of the form $`uG_xv`$, where $`u,vX^{}`$, $`xX`$, and $`uxv`$ equals $`w`$ modulo $`𝒫`$. Now it remains to analyze all $`2`$-cells of $`𝒦_w`$ and to find out what will be the relations between the generators of $`G`$ described above. According to the description of $`2`$-cells in a Squier complex given in Section 1, let a $`2`$-cell be given by a $`5`$-tuple $`(u,\mathrm{}_1r_1,z,\mathrm{}_2r_2,v)`$, where $`(\mathrm{}_1,r_1)`$, $`(\mathrm{}_2,r_2)`$ belong to $`𝒮𝒲`$. Note that the word $`u\mathrm{}_1z\mathrm{}_2v`$ equals $`w`$ modulo $`𝒫`$. Let us consider several cases depending on the defining relations involved. The relation between edges that is obtained from the given $`2`$-cell, has the form $$(u,\mathrm{}_1r_1,z\mathrm{}_2v)(ur_1z,\mathrm{}_2r_2,v)=(u\mathrm{}_1z,\mathrm{}_2r_2,v)(u,\mathrm{}_1r_1,zr_2v).$$ (9) a) Let $`(\mathrm{}_1,r_1)`$, $`(\mathrm{}_2,r_2)`$ both belong to $`𝒮`$. Then each of the words $`u`$, $`v`$, $`z`$ can be decomposed into the product of factors. Using equality (7), one can replace in (9) the words $`u`$, $`v`$, $`z`$ by their projections (we use the fact that each of the words $`\mathrm{}_j`$, $`r_j`$ ($`j=1,2`$) coincides with its projection). Thus one can assume that the words $`u`$, $`v`$, $`z`$ in (9) belong to $`X^{}`$. Then (9) is a defining relation of the group $`\pi _1(𝒦(𝒬),w)𝒟(𝒬,w)`$ (calculated by using the maximal subtree $`𝒯_Q`$). b) Suppose that none of the relations $`(\mathrm{}_1,r_1)`$, $`(\mathrm{}_2,r_2)`$ belongs to $`𝒮`$. Suppose also that these relations are applied to different factors of the word $`u\mathrm{}_1z\mathrm{}_2v`$. This means that there exist letters $`x,yX`$ and decompositions of the form $`u=u_1u_2`$, $`z=z^{}z_0z^{\prime \prime }`$, $`v=v_2v_1`$, where $`u_1`$ is the maximal prefix of $`u`$ that is a product of factors, $`v_1`$ is the maximal suffix of $`v`$ that is a product of factors, and $`z_0`$ is a maximal subword of the word $`z`$ that is a product of factors. (It is not hard to see that this word can be found uniquely.) Here $`x`$ equals $`u_2\mathrm{}_1z^{}`$ and $`y`$ equals $`z^{\prime \prime }\mathrm{}_2v_2`$ (equalities are considered modulo $`\widehat{𝒫}_x`$ and $`\widehat{𝒫}_y`$, respectively). Let us substitute the decompositions of the words $`u`$, $`v`$, $`z`$ in the equality (9) using the fact that $`\pi (u_2\mathrm{}_1z^{})=\pi (u_2r_1z^{})=x`$, $`\pi (z^{\prime \prime }\mathrm{}_2v_2)=\pi (z^{\prime \prime }r_2v_2)=y`$. We obtain that the elements $`\pi (u_1)(u_2,\mathrm{}_1r_1,z^{})\pi (z_0)y\pi (v_1)`$ $`\pi (u_1)x\pi (z_0)(z^{\prime \prime },\mathrm{}_2r_2,v_2)\pi (v_1)`$ commute. The first of them belongs to the group $`\pi (u_1)G_x\pi (z_0)y\pi (v_1)`$, and the second one belongs to the group $`\pi (u_1)x\pi (z_0)G_y\pi (v_1)`$. Let $`U`$, $`V`$, $`Z`$ be arbitrary words over $`X`$ and let $`x,yX`$ be arbitrary letters such that the word $`UxZyV`$ equals $`w`$ modulo $`𝒫`$. We can conclude that any element in $`UG_xZyV`$ commutes with any element in $`UxZG_yV`$ since for any edges $`e`$ and $`f`$ that belong to the generating sets of the groups $`G_x`$ and $`G_y`$ respectively, one can find a suitable $`2`$-cell of the form described above in such a way that the defining relations obtained from it will be the relation of commutativity of $`UeZyV`$ and $`UxZfV`$. Thus we get relations of the form $`[UG_xZyV,UxZG_yV]=1`$, where $`UxZyV`$ equals $`w`$ modulo $`𝒫`$, $`U,V,ZX^{}`$, $`x,yX`$. c) Again, let none of the relations $`(\mathrm{}_1,r_1)`$, $`(\mathrm{}_2,r_2)`$ belong to $`𝒮`$ but assume now that the relations $`(\mathrm{}_1=r_1)`$ and $`(\mathrm{}_2=r_2)`$ are applied to the same factor of the word $`u\mathrm{}_1z\mathrm{}_2v`$. This means that there exists a letter $`xX`$ and decompositions $`u=u_1u_2`$, $`v=v_2v_1`$, where $`u_1`$ is the maximal prefix of the word $`u`$ that is a product of factors, $`v_1`$ is the maximal suffix of the word $`v`$ that is a product of factors. Now $`x`$ equals $`u_2\mathrm{}_1z\mathrm{}_2v_2`$ modulo $`\widehat{𝒫}_x`$. Consider a $`2`$-cell of the Squier complex over $`\widehat{𝒫}_x`$ that corresponds to the $`5`$-tuple $`(u_2,\mathrm{}_1r_1,z,\mathrm{}_2r_2,v_2)`$. All cells of this form lead to the defining relations of a group isomorphic to $`G_x`$. Acting on this $`2`$-cell by the element $`\pi (u_1)`$ on the left and by the element $`\pi (v_1)`$ on the right, we get a $`2`$-cell of the complex $`𝒦_w(𝒫)`$. The relation written on its contour is equivalent to (9) if one takes the equality (7) into account. Thus we get the defining relations of all groups of the form $`UG_xV`$, where $`UxV`$ equals $`w`$ modulo $`𝒫`$, $`U,VX^{}`$, $`xX`$. d) Suppose that one of the relations $`(\mathrm{}_1,r_1)`$, $`(\mathrm{}_2,r_2)`$ belongs to $`𝒮`$ and the other one does not. First of all, let $`(\mathrm{}_1,r_1)𝒮`$. Then we have decompositions of the form $`z=z_0z^{\prime \prime }`$, $`v=v_2v_1`$, where $`z_0`$, $`v_1`$ are products of factors that are chosen to be minimal with respect to this property, as above. Now $`z^{\prime \prime }\mathrm{}_2v_2`$ equals $`x`$ modulo $`\widehat{𝒫}_x`$ for some letter $`xX`$. Let $`f=(z^{\prime \prime },\mathrm{}_2r_2,v_2)`$. Substituting the decompositions of words $`v`$, $`z`$ in (9), taking into account that $`\pi (z^{\prime \prime }\mathrm{}_2v_2)=\pi (z^{\prime \prime }r_2v_2)=x`$ and applying (7), we obtain the following equality: $$(\pi (u),\mathrm{}_1r_1,\pi (z_0)x\pi (v_1))(\pi (u)r_1\pi (z_0)f\pi (v_1))=$$ $$(\pi (u)\mathrm{}_1\pi (z_0)f\pi (v_1))(\pi (u),\mathrm{}_1r_1,\pi (z_0)x\pi (v_1)).$$ Thus for each $`xX`$ and for any words $`U`$, $`Z`$, $`V`$ over $`X`$ such that $`(\mathrm{}_1,r_1)𝒮`$ and $`U\mathrm{}_1ZxV`$ equals $`w`$ modulo $`𝒫`$, we obtain the relations $$Ur_1ZfV=(U\mathrm{}_1ZfV)^e,$$ (10) where $`e=(U,\mathrm{}_1r_1,ZxV)`$, and $`f`$ runs over the generating set of the group $`\pi _1(\widehat{𝒫}_x,x)G_x`$. Analogously, if $`(\mathrm{}_2,r_2)𝒮`$, then we get the relations $$UfZr_2V=(UfZ\mathrm{}_2V)^e,$$ (11) where $`U`$, $`Z`$, $`V`$ are words over $`X`$, $`xX`$, $`e=(UxZ,\mathrm{}_2r_2,V)`$, and $`f`$ runs over the generating set of the group $`\pi _1(\widehat{𝒫}_x,x)G_x`$. To compute the diagram product, let us define a structure of a graph of groups on the $`1`$-skeleton of the Squier complex $`𝒦(𝒬)`$. Our diagram product is the fundamental group of the corresponding $`2`$-complex of groups. We will apply the above described procedure of computing a fundamental group of a $`2`$-complex of groups and we will then compare it with the presentation of $`G`$. Let $`U`$, $`V`$ be words over $`X`$, $`xX`$. By $`H(UxV)`$ we denote the group $`UG_xV`$ isomorphic to $`G_x`$. Then, for any word $`u=u_1\mathrm{}u_m`$, where $`u_iX`$ ($`1im`$), we denote by $`H_u`$ the free product of the groups of the form $$H(u_1\mathrm{}u_{i1}u_iu_{i+1}\mathrm{}u_m)$$ over all $`i`$ from $`1`$ to $`m`$. These groups are assigned to vertices of $`𝒦(𝒫)`$. Now let us take an edge $`e=(u,st,v)`$, where $`u,vX^{}`$, $`(s=t)𝒮`$. The group $`H_e=H_u\times H_v`$ is assigned to it. The maps $`\iota _e`$, $`\tau _e`$ naturally embed $`H_e=H_u\times H_v`$ into the groups $`H_{usv}H_u\times H_s\times H_v`$ and $`H_{utv}H_u\times H_t\times H_v`$, respectively (here $`H_u`$ maps onto $`H_u`$, nd $`H_v`$ maps onto $`H_v`$). It is clear that, instead of presenting edge groups of the form $`G_e`$, one can present an isomorphism of some subgroup of $`H_{\iota (e)}`$ to some subgroup of $`H_{\tau (e)}`$, for each egde $`e`$. In our case this isomorphism is very simple: it maps the subgroup $`H_u\times \{\mathrm{\hspace{0.17em}1}\}\times H_v`$ $`H_{usv}`$ onto $`H_u\times \{\mathrm{\hspace{0.17em}1}\}\times H_v`$ $`H_{utv}`$. In this case we will speak about the isomorphism induced by an edge $`e`$. The groups of the form $`H(UxV)`$ will be presented as groups generated by edges of the form $`UfV`$, where $`f`$ runs over the set of edges of the connected component of the Squier complex $`𝒦(\widehat{𝒫}_x)`$ that contain vertex $`x`$. Here edges satisfy the relations $`UfV=1`$ whenever $`f`$ belongs to the tree $`𝒯_x`$, and also relations $`UrV=1`$, where $`r`$ is the defining path of a $`2`$-cell of this complex. These relations of the group $`G`$ were obtained in subsection c). For the direct products of the groups of the form $`H(UxV)`$, we introduce relations of commutativity: each element in the group $`H(UxZyV)`$ commutes with each element in the group $`H(UxZyV)`$. For the group $`G`$, such relations were obtained in subsection b). Let us have an edge $`e=(u,st,v)`$ that belongs to the complex $`𝒦_w`$. Let $`x,yX`$, $`u=u_1xu_2`$, $`v=v_1yv_2`$. The isomorphism induced by the edge $`e`$, takes $`H(usv_1yv_2)`$ to $`H(utv_1yv_2)`$, and it takes $`H(u_1xu_2yv)`$ to $`H(u_1xu_2yv)`$. So, according to (5, the conjugation by the edge $`e`$ leads to the following relations $$(usv_1fv_2)^e=(utv_1fv_2),$$ (12) $$(u_1fu_2sv)^e=(u_1fu_2tv),$$ (13) where $`f`$ runs over the set of edges that generate the corresponding group in each of the cases. These relations coincide with relations (10) and (11) of the group $`G`$ from subsection d). Finally, we represent the group $`\pi _1(𝒦(𝒬),w)`$ as a group generated by edges of $`𝒦_w(𝒬)`$, claiming that the edges in $`𝒯_Q`$ are equal to the identity and adding relations that correspond to the defining paths of $`2`$-cells of this complex. Such relations of the group $`G`$ are described in subsection a). Thus the quotient group of the free product of the group $`\pi _1(𝒦(𝒬),w)`$ and groups of the form $`H_u`$ for all vertices $`u`$ of $`𝒦_w(𝒬)`$, by the normal closure of relations (12) and (13), is given by the same generators and defining relations as $`G`$. This means that the diagram product $`𝒟(G_X;𝒮,w)`$ of the family $`G_X=\{G_x(xX)\}`$ of groups over the presentation $`𝒬=X𝒮`$ with base $`w`$ is isomorphic to the group $`G=\pi _1(𝒦(𝒫,w)`$, that is, to the diagram group $`𝒟(𝒫,w)`$. The Theorem is proved. Now let us consider a few applications of Theorem 4. The first three of them deal with already known constructions. We give them to demonstrate that all group-theoretical constructions for diagram groups we dealt with earlier (see \[14, Section 8\]) are examples of diagram products. Then we show that the class of diagram groups is closed under some new operations: countable direct powers, restricted wreath products with the group $`𝐙`$, and also under some new special construction that will be used in Section 5. ###### Example 5 Let $`X=\{x_1,\mathrm{},x_n\}`$ be a finite alphabet. Consider the presentation $`𝒬=X\mathrm{}`$ with empty set of defining relations and let $`w=x_1\mathrm{}x_n`$. To each letter $`x_i`$, we assign an arbitrary group $`G_i`$ ($`1in`$). It is obvious that the connected component of the Squier complex of $`𝒬`$, which contains $`w`$, consists of exactly one vertex $`w`$. In the corresponding graph of groups, we have the group $`G_w=G_1\times \mathrm{}\times G_n`$. Obviously, it is the fundamental group of the $`2`$-complex of groups from the definition of a diagram product. Thus the diagram product $`𝒟(G_X;𝒮,w)`$ of the family $`G_X=\{G_i(1in)\}`$ of groups over the presentation $`𝒬`$ with base $`w`$ is the direct product $`G_1\times \mathrm{}\times G_n`$. ###### Example 6 Let $`I`$ be a nonempty set and let $`G_i`$ ($`iI`$) be an arbitrary family of groups. Let us consider an alphabet $`X=\{x\}\{x_i(iI)\}`$ and let $`𝒬=X𝒮`$, where $`𝒮`$ consists of relations of the form $`x=x_i`$ for all $`iI`$. Let $`G_X`$ be a family of groups that assigns the trivial group to the letter $`x`$ and the group $`G_i`$ to the letter $`x_i`$ ($`iI`$). The connected component of the Squier complex $`𝒦(𝒬)`$ containing $`x`$ is a tree in which the vertex $`x`$ is connected by edges with all vertices labelled by $`x_i`$ ($`iI`$). Let us consider the structure of a graph of groups on the $`1`$-skeleton of the connected component of this Squier complex according to the definition 3. It is easy to see that all edge groups are trivial. From this, using description (4), it is easy to see that the fundamental group of the resulting $`2`$-complex of groups is the free product of groups $`G_i`$ ($`iI`$). So the diagram product $`𝒟(G_X;𝒮,x)`$ of the family $`G_X=\{G_i(iI)\}`$ of groups over the presentation $`𝒬`$ with base $`x`$ is the free product $`G_i`$, $`iI`$. ###### Example 7 Let $`G`$, $`H`$ be any groups. Let us consider the presentation $`𝒬=X𝒮`$, where $`X=\{x,y,z\}`$, $`𝒮=\{x=xy,z=yz\}`$. Let $`G_x=G`$, $`G_y=1`$, $`G_z=H`$ and consider the diagram product $`𝒟(G_X;𝒮,xz)`$ of the family $`G_X=\{G_x,G_y,G_z\}`$ of groups over the presentation $`𝒬`$ with base $`xz`$. We obtain that it is isomorphic to the group $`GH`$, where $``$ is the operation defined in . One can check this directly by comparing the Theorem 4 and the definition of the operation $``$ in . Indeed, the group $`GH`$ can be described in the following way. Consider countable number of copies $`G_i`$ of the group $`G`$ and countable number of copies $`H_i`$ of the group $`H`$ ($`i𝐙`$). An infinite cyclic group $`z`$ acts on the group $$(G_i)\times (H_i)$$ (14) (free products are taken over all $`i𝐙`$) permuting the factors: it takes $`G_i`$ to $`G_{i+1}`$ and $`H_i`$ to $`H_{i+1}`$ for all integers $`i`$. The group $`GH`$ is the semidirect product of the group (14) and the group $`z`$. ###### Example 8 Let $`G`$ be an arbitrary group. Let us consider the presentation $`𝒬=X𝒮`$, where $`X=\{x,y\}`$, $`𝒮=\{x=xy\}`$. Let $`G_x=1`$, $`G_y=G`$. Consider the diagram product $`𝒟(G_X;𝒮,x)`$ of the family $`G_X=\{G_x,G_y\}`$ of groups over the presentation $`𝒬`$ with base $`x`$. Let us show that it is isomorphic to the countable direct power of the group $`G`$. The connected component $`𝒦_x`$ of the Squier complex over $`𝒬`$ containing $`x`$ has the following form: Here vertices are all words of the form $`xy^i`$ ($`i0`$), positive edges have the form $`e_i=(1,xxy,y^i)`$ ($`i0`$), and the maximal subtree $`𝒯`$ includes all these edges. This complex has no $`2`$-cells. Thus it is obvious that its fundamental group is trivial. Using our convention that (given a maximal subtree) all edges are regarded as elements of the fundamental groups, we have equalities $`e_i=1`$ for all $`i0`$. Let us consider the structure of the graph of groups on the $`1`$-skeleton of the complex $`𝒦_x`$ according to the definition of a diagram product. We will obtain that the group $`H_v=G_{n1}\times \mathrm{}\times G_{nn}`$, where $`G_{ni}`$ ($`1in`$) are groups isomorphic to $`G`$, corresponds to the vertex $`v=xy^n`$ ($`n0`$). Let us consider a positive edge $`e=e_n=(1,xxy,y^n)`$ ($`n0`$). The group $`G_e`$ is isomorphic to $`G^n`$, the $`n`$th direct power of $`G`$. The embedding $`\iota _e`$ maps $`G^n`$ isomorphically onto $`G_{n1}\times \mathrm{}\times G_{nn}`$, and the mapping $`\tau _e`$ maps $`G^n`$ isomorphically onto the last $`n`$ factors of the direct product $`G_{n+1,1}\times G_{n+1,2}\mathrm{}\times G_{n+1,n+1}`$. The relations from the description (5), together with the equality $`e=1`$, allow to identify corresponding elements of $`G_{n1}`$ and $`G_{n+1,2}`$, …, $`G_{nn}`$ and $`G_{n+1,n+1}`$. Thus we can introduce the following notation: $`G_0=G_{11}=G_{22}=\mathrm{}`$, $`G_1=G_{21}=G_{32}=\mathrm{}`$, …, $`G_n=G_{n+1,1}=G_{n+2,2}=\mathrm{}`$, … . Each of these groups is isomorphic to $`G`$. They generate a countable direct power of $`G`$. So our diagram product is the countable direct power of $`G`$. Theorem 4 implies the following result. ###### Theorem 9 The class of diagram groups is closed under countable direct powers. ###### Example 10 Let $`G`$ be arbitrary group. Consider the presentation $`𝒬=X𝒮`$, where $`X=\{x,y,z\}`$, $`𝒮=\{x=xy,z=yz\}`$. Let $`G_x=1`$, $`G_y=G`$, $`G_z=1`$ and consider the diagram product $`𝒟(G_X;𝒮,xz)`$ of the family $`G_X=\{G_x,G_y,G_z\}`$ of groups over the presentation $`𝒬`$ with the base $`xz`$. Let us show that it is isomorphic to the (restricted) wreath product $`Gwr𝐙`$. The connected component $`𝒦_{xz}`$ of the Squier complex over $`𝒬`$ that contains $`xz`$ has the following form: Here vertices are all words of the form $`xy^iz`$ ($`i0`$), positive edges have the form $`e_i=(1,xxy,y^iz)`$, $`f_i=(xy^i,zyz,1)`$ ($`i0`$), and the maximal subtree $`𝒯`$ consists of the edges $`e_i`$, $`i0`$. All $`2`$-cells can be described as follows. Let $`i0`$. Consider the vertex $`xy^iz`$. The edges $`e_i`$, $`f_i`$ going out of it correspond to independent transformations of words. So the given pair of edges defines two homotopic paths $`e_if_{i+1}`$ and $`f_ie_{i+1}`$ that define a $`2`$-cell. According to the convention that edges are regarded as elements of the fundamental group $`\pi _1(𝒦,xz)`$, we have $`e_i=1`$ ($`i0`$) in the group. The equalities $`e_if_{i+1}=f_ie_{i+1}`$ that hold in this group imply $`f_i=f_{i+1}`$ for all $`i0`$. Let $`f=f_0=f_1=f_2=\mathrm{}`$. According to the definition of a diagram product, let us consider the structure of the graph of groups on the $`1`$-skeleton of $`𝒦_{xz}`$. The group $`H_v=G_{n1}\times \mathrm{}\times G_{nn}`$ is assigned to the vertex $`v=xy^nz`$ ($`n0`$), where $`G_{ni}`$ ($`1in`$) is a group isomorphic to $`G`$. Let us consider relations (5) that correspond to positive edges. Let $`e=e_n=(1,xxy,y^n)`$ ($`n0`$). As in the previous example, using the equality $`e=1`$, we identify corresponding elements of the groups $`G_{ni}`$ and $`G_{n+1,i+1}`$ ($`1in`$) and introduce the notation $`G_0=G_{11}=G_{22}=\mathrm{}`$, $`G_1=G_{21}=G_{32}=\mathrm{}`$, …, $`G_n=G_{n+1,1}=G_{n+2,2}=\mathrm{}`$, … . As above, these groups generate a countable direct power of the group $`G`$. Now let $`e=f_n=(xy^n,zyz,1)`$ ($`n0`$). Consider relations of the form (5) that correspond to these edges. The group $`G_e`$ is still the $`n`$th power of $`G`$. The embedding $`\iota _e`$ maps $`G^n`$ onto $`G_{n1}\times \mathrm{}\times G_{nn}`$ isomorphically, and the embedding $`\tau _e`$ maps $`G^n`$ isomorphically onto the first $`n`$ factors of the direct product $`G_{n+1,1}\times G_{n+1,2}\mathrm{}\times G_{n+1,n+1}`$. So the relations that correspond to the edge $`f`$ show that conjugation by $`f`$ takes $`G_{ni}`$ to $`G_{n+1,i}`$ ($`1in`$). Using our notation, we obtain that the conjugation by $`f`$ takes the group $`G_k`$ to the group $`G_{k+1}`$ for all $`k0`$. Thus the diagram product we are considering is generated by groups $`G_0`$, $`G_1`$, … and the element $`f`$. From this, one can deduce that we have the restricted wreath product $`Gwr𝐙`$. Applying Theorem 4, we get one more result. ###### Theorem 11 The class of diagram groups is closed under restricted wreath products with the infinite cyclic group $`𝐙`$, that is, if $`G`$ is a diagram group, then $`Gwr𝐙`$ is also a diagram group. Note that if we take R. Thompson’s group $`F`$ represented by diagrams over $`uuu=u`$ with base $`u`$ and consider the presentation $`x,u,zxu=x,uz=z,uu=u`$, then the diagram group over it with base $`xz`$ will be isomorphic not to $`Fwr𝐙`$ but to $`F`$. This can be checked directly. To get the group $`Fwr𝐙`$, one needs to represent the group $`F`$ by diagrams according to the statement of Theorem 4. Namely, one has to take the diagram group with base $`y`$ over the presentation $`u,a,yy=aua,uu=u`$. Then the diagram group with base $`xz`$ over $`𝒫=x,y,z,a,uxy=x,yz=z,y=aua,uu=u`$ will be isomorphic to $`Fwr𝐙`$. The reader can easily list the presentations that lead to diagram groups of the form $`(\mathrm{}((𝐙wr𝐙)wr𝐙)\mathrm{})wr𝐙`$. Let us make one more remark. In we constructed an example of a diagram group that was finitely generated but not finitely presented (Theorem 10.5). We took the group $`𝐙𝐙`$ for this purpose. It has a presentation with three generators $$𝐙𝐙=a,b,t[a^{t^n},b]=1(n0).$$ Now we can also take the group $`𝐙wr𝐙`$ as an example of a finitely generated but not finitely presented group (the fact that $`𝐙wr𝐙`$ has no finite presentations can be easily proved using either HNN-extensions or representation of groups by transformations). We have the following presentation with two generators for this group: $$𝐙wr𝐙=a,b[a^{b^n},a]=1(n1).$$ In the next example we deal with a more complicated construction. At first sight one can think it is quite artificial. However, we will efficiently use it later, in Section 5. Let us consider the following group-theoretical construction. Take two groups, $`G`$ and $`H`$. We assign to them a new group denoted by $`𝒪(G,H)`$. Let us consider a countable family of copies $`G_i`$ of the group $`G`$, and a coutable family of copies $`H_i`$ of $`H`$ ($`i𝐙`$). For any $`i𝐙`$, let $`g_i`$ ($`h_i`$) be the element that corresponds to $`gG`$ ($`hH`$). By $`G^{\mathrm{}}`$ ($`H^{\mathrm{}}`$) we denote a coutable direct power of the groups $`G_i`$ ($`H_i`$) taken over all $`i𝐙`$. Let $$𝒪(G,H)=G^{\mathrm{}}H^{\mathrm{}}c/𝒩,$$ (15) where $`𝒩`$ is the normal closure of the set of relations of the two forms: $$g_i^t=g_{i+1},h_i^t=h_{i+1}\text{ for all }i𝐙,gG,hH;$$ (16) $$[g_i,h_j]=1\text{ for all }i,j𝐙,gG,hH\text{ such that }ij.$$ (17) ###### Example 12 Let $`G`$, $`H`$ be arbitrary groups. Let us consider the presentation $`𝒬=X𝒮`$, where $`X=\{x,y,\overline{y},z,p,q,r\}`$, $`𝒮=\{x=xyp,z=r\overline{y}z,pyq=q\overline{y}r\}`$. Let $`G_y=G`$, $`G_{\overline{y}}=H`$, $`G_x=G_z=G_p=G_q=G_r=1`$. Consider the diagram product $`𝒟(G_X;𝒮,w)`$ of the family $`G_X=\{G_x,G_y,G_{\overline{y}},G_z,G_p,G_q,G_r\}`$ of groups over the presentation $`𝒬`$ with base $`w=xyq\overline{y}z`$. Let us show that it is isomorphic to $`𝒪(G,H)`$. The connected component $`𝒦_w`$ of the Squier complex over $`𝒬`$ that contains $`w`$, has the following form: Here the vertices are words $`w_{ij}=x(yp)^iyq\overline{y}(r\overline{y})^jz`$ ($`i,j0`$). The positive edges have the form $`e_{ij}=(1,xxyp,(yp)^iyq\overline{y}(r\overline{y})^jz)`$, $`f_{ij}=(x(yp)^iyq\overline{y}(r\overline{y})^j,zr\overline{y}z,1)`$ and $`g_{ij}=(x(yp)^iy,pyqq\overline{y}r,\overline{y}(r\overline{y})^jz)`$ ($`i,j0`$). We choose the maximal subtree $`𝒯`$ formed by the edges $`e_{ij}`$ for all $`i,j0`$ and also by the edges $`f_{0j}`$ for $`j0`$. Thus our convention that the choice of $`𝒯`$ makes the edges to be elements of the fundamental group $`\pi _1(𝒦,w)`$, leads to the equalities $`e_{ij}=1`$ ($`i,j0`$), $`f_{0j}=1`$ ($`j0`$) in the fundamental group of the complex. Let us describe all $`2`$-cells of the complex $`𝒦_w`$. Recall that there are defining relations of three types in $`𝒬`$: $`x=xyp`$, $`pyq=q\overline{y}r`$, $`z=r\overline{y}z`$. If we have two independent applications of elementary transformations to words of the form $`w_{ij}`$, then it is easy to see that they belong to different types because each of the letters $`x`$, $`q`$, $`z`$ occurs into the word $`w_{ij}`$ only once. Therefore, we have exactly three situations. 1) The relations applied in the independent transformations are $`x=xyp`$ and $`z=r\overline{y}z`$ (see the picture below). This diagram corresponds to the two paths in the Squier complex: $`e_{ij}f_{i+1,j}`$ and $`f_{ij}e_{i,j+1}`$. This leads to relations $`e_{ij}f_{i+1,j}=f_{ij}e_{i,j+1}`$. Simplifying, we have $`f_{i+1,j}=f_{ij}`$ for all $`i,j0`$. It is obvious that $`f_{ij}`$ does not depend on $`i`$. So $`f_{0j}=1`$ gives $`f_{ij}=1`$ for all $`i,j0`$. 2) The relations are $`x=xyp`$ and $`pyq=q\overline{y}r`$ (see the picture below). In this case, we have the equality $`e_{i+1,j}g_{i+1,j}=g_{i,j}e_{i,j+1}`$, that is, $`g_{i+1,j}=g_{ij}`$ for all $`i,j0`$. This means that $`g_{ij}`$ does not depend on $`i`$. 3) The relations are $`pyq=q\overline{y}r`$ and $`z=r\overline{y}z`$ (see the picture below). Here we have the equality $`g_{i,j}f_{i,j+1}=f_{i+1,j}g_{i,j+1}`$. Hence $`g_{ij}=g_{i,j+1}`$ for all $`i,j0`$. Therefore, $`g_{ij}`$ depends on neither $`i`$ nor $`j`$. For convenience, let $`c=g_{ij}`$ for all $`i,j0`$. Let us take an arbitrary vertex $`v=w_{ij}=x(yp)^iyq\overline{y}(r\overline{y})^jz`$ for some $`i,j0`$. In the graph of groups that corresponds to the diagram product, the product of $`(i+1)`$th power of $`G`$ and the $`(j+1)`$th power of $`H`$ will correspond to the vertex $`G_{w_{ij}}`$ (the number of factors is just the number of occurrences of $`y`$ and $`\overline{y}`$ in $`v`$, respectively). Thus we can present the group $`G_{w_{ij}}`$ in the form $$K_{ij}=L_{iji}\times \mathrm{}L_{ij0}\times H_{ij0}\times \mathrm{}H_{ijj},$$ where the factors of the form $`L_{ijk}`$ are isomorphic to $`G`$, and the factors of the form $`H_{ijk}`$ are isomorphic to $`H`$. Relations (5) that correspond to a positive edge $`e=(u,st,v)`$ will be studied with respect to the type of the involved defining relation $`(s=t)𝒮`$ (there are three types of them). 1) $`s=x`$, $`t=xyp`$. We have $`e=(1,xxyp,v)`$, where $`v=(yp)^iyq\overline{y}(r\overline{y})^jz`$ for some $`i,j0`$. The group $`G_e`$ is isomorphic to $`G^{i+1}\times H^{j+1}`$. It maps isomorphically onto $`K_{ij}`$ under $`\iota _e`$. Note that $`K_{i+1,j}`$ is the direct product of $`L_{i+1,j,i+1}`$ and the isomorphic image of $`K_{ij}`$ under $`\tau _e`$. Using the fact that $`e=e_{ij}=1`$ we see that relation (5) identifies some subgroups. Let us write down these identifications as equalities. By these equalities we mean that the corresponding elements of equal groups are identified. We have: $`L_{ijk}=L_{i+1,j,k}`$ for $`0ki`$ and $`H_{ijk}=H_{i+1,j,k}`$ for $`0kj`$. 2) $`s=z`$, $`t=r\overline{y}z`$. Now $`e=(u,zr\overline{y}z,1)`$, where $`u=x(yp)^iyq\overline{y}(r\overline{y})^j`$ for some $`i,j0`$. Arguing analogously to the previous case and taking into account that $`e=f_{ij}=1`$, we conclude that relation (5) leads to the following identifications of subgroups: $`L_{ijk}=L_{i,j+1,k}`$ for $`0ki`$ and $`H_{ijk}=H_{i,j+1,k}`$ for $`0kj`$. Summarizing what we got in the first two cases of relations, we see that groups $`L_{ijk}`$, $`H_{ijk}`$ depend of $`k`$ only. In other words, one can introduce groups $`L_k`$, $`H_k`$ ($`k0`$) in such a way that the equalities $`L_{ijk}=L_k`$ for all $`ik`$, $`j0`$, and $`H_{ijk}=H_k`$ for all $`i0`$, $`jk`$ hold in our diagram product. 3) $`s=pyq`$, $`t=q\overline{y}r`$. In this case $`e=(u,pyqq\overline{y}r,v)`$, where $`u=x(yp)^iy`$, $`v=\overline{y}(r\overline{y})^j`$ for some $`i,j0`$. In the fundamental group $`\pi _1(𝒦,w)`$, the equality $`e=g_{ij}=c`$ holds, as it was shown above, the group $`G_e`$ is isomorphic to $`G^{i+1}\times H^{j+1}`$. We have $$G_{\iota (e)}=L_{i+1}\times \mathrm{}\times L_0\times H_0\times \mathrm{}\times H_j,$$ $`\iota _e`$ maps $`G^{i+1}`$ onto the direct product $`L_{i+1}\times \mathrm{}\times L_1`$ and it maps $`H^{j+1}`$ onto the direct product $`H_0\times \mathrm{}\times H_j`$. Analogously, $$G_{\tau (e)}=L_i\times \mathrm{}\times L_0\times H_0\times \mathrm{}\times H_{j+1},$$ $`\tau _e`$ maps $`G^{i+1}`$ onto the direct product $`L_i\times \mathrm{}\times L_0`$ and it maps $`H^{j+1}`$ onto the direct product $`H_1\times \mathrm{}\times H_{j+1}`$. Therefore, conjugating by the element $`c`$ takes $`L_{i+1}`$, …, $`L_1`$ to $`L_i`$, …, $`L_0`$ respectively. The subgroups $`H_0`$, …, $`H_j`$ are taken to $`H_1`$, …, $`H_{j+1}`$ respectively, under this conjugation. Briefly, we can write $`L_{i+1}^c=L_i`$, $`H_j^c=H_{j+1}`$ for any $`i,j0`$. Thus the equalities $`L_i=L_0^{c^i}`$, $`H_j=H_0^{c^j}`$ hold for any nonnegative integers $`i`$, $`j`$. Let us extend these equalities to the case of negative $`i`$, $`j`$ regarding these equalities as definitions. Note that elements from different subgroups of the form $`L_i`$ ($`i0`$) commute. So the analogous fact is true for all integers $`i`$. The same fact is true for subgroups $`H_j`$ for all $`j𝐙`$. Let $`G_i=L_i`$ ($`i𝐙`$) by definition. Obviously, $`G_i^c=G_{i+1}`$, $`H_i^c=H_{i+1}`$ for all $`i𝐙`$. This means that relations (16) hold. We also have conditions that any element in $`L_0`$, $`L_1`$, …commutes with any element in $`H_0`$, $`H_1`$, … . In particular, $`[L_0,H_{ji}]=1`$ for any $`ji`$. Taking into account that $`G_0=L_0`$ and conjugating by the element $`c^i`$, we obtain $`[G_i,H_j]=1`$ for $`ij`$, that is, relations (17) hold. It is easy to see that these relations are in fact equivalent to the condition that $`[L_i,H_j]=1`$ for any $`i,j0`$. Indeed, the inequality $`ij`$ and relations (16) allow us to conclude that $`[G_i,H_j]=1`$, where $`G_i`$ is $`L_i`$. Thus we see that the diagram product we have calculated is in fact the group given by relations (16) and (17), that is, it is isomorphic to $`𝒪(G,H)`$. Using Theorem 4, we have the following result. ###### Theorem 13 If $`G`$, $`H`$ are diagram groups, then $`𝒪(G,H)`$ is also a diagram group. The previous example shows in details how, given two diagram groups $`G`$ and $`H`$, one can construct a presentation and a base, for which $`𝒪(G,H)`$ will be a diagram group. ## 3 Nilpotent and Abelian Subgroups of Diagram Groups We know from the previous section that soluble subgroups of any degree can be subgroups of diagram groups. Contrary to that, we shall prove in this Section that any nilpotent subgroup of a diagram group is abelian. We will also establish the fact that all abelian subgroups of diagram groups are free abelian. This will generalize the result that any abelian diagram group is free abelian. Finally, we shall describe finite sets of pairwise commuting diagrams, generalizing a description of pairs of commuting diagrams from . We will use some concepts from combinatorics on diagrams from \[14, Section 15\]. For reader’s convenience, let us recall some definitions. A spherical diagram is called absolutely reduced if any positive integer power of it is reduced (does not contain dipoles). A spherical diagram is called normal if it cannot be decomposed into a sum of two non-spherical diagrams. We proved \[14, Theorem 15.14\] that for any spherical diagram $`\mathrm{\Delta }`$ there exists an absolutely reduced normal spherical diagram $`\widehat{\mathrm{\Delta }}`$ (that may have different base, in general) and some (not necessarily spherical) diagram $`\mathrm{\Psi }`$ such that $`\mathrm{\Delta }=\mathrm{\Psi }^1\widehat{\mathrm{\Delta }}\mathrm{\Psi }`$. ###### Theorem 14 Let $`H`$ be an arbitrary subgroup of a diagram group $`𝒟(𝒫,w)`$. Then the centre of $`H`$ and the commutator subgroup of $`H`$ intersect trivially that is, $`Z(H)H^{}=1`$. Proof. Let $`G=𝒟(𝒫,w)`$ be a diagram group and let $`H`$ be a subgroup of $`G`$. Suppose that $`Z(H)H^{}1`$. Consider a nontrivial element $`gZ(H)H^{}`$ and let $`\mathrm{\Delta }`$ be a diagram representing it. Applying \[14, Lemma 15.10c\], we find an absolutely reduced diagram $`\mathrm{\Delta }_0`$ that is conjugated to $`\mathrm{\Delta }`$. Let $`\mathrm{\Delta }_0=\mathrm{\Psi }^1\mathrm{\Delta }\mathrm{\Psi }`$, where $`\mathrm{\Psi }`$ is a $`(w,w_0)`$-diagram. Conjugation by $`\mathrm{\Psi }`$ is an isomorphism that takes the group $`G`$ to the group $`G_0=𝒟(𝒫,w_0)`$. Under this isomorphism, the subgroup $`H`$ is taken to a subgroup $`H_0`$, where $`g_0Z(H_0)H_0^{}`$, and the element $`g_0`$ is represented by an absolutely reduced $`(w_0,w_0)`$-diagram $`\mathrm{\Delta }_0`$. Let us decompose the diagram $`\mathrm{\Delta }_0`$ into a sum of components: $`\mathrm{\Delta }_0=A_1+\mathrm{}+A_m`$, where $`A_i`$ is a spherical $`(w_i,w_i)`$-diagram ($`1im`$). As in \[14, Theorem 15.35\], we conclude that the centralizer of $`g_0`$ is the direct sum of centralizers of the elements represented by diagrams $`A_1`$, …, $`A_m`$. More precisely, if $`\mathrm{\Gamma }`$ is a spherical $`(w_i,w_i)`$-diagram that commutes with $`\mathrm{\Delta }_0`$ in the group $`G_0`$, then $`\mathrm{\Gamma }=B_1+\mathrm{}+B_m`$, where $`B_i`$ is a $`(w_i,w_i)`$-diagram that commutes with $`A_i`$. By the assumption, any diagram representing an element in $`H_0`$, commutes with $`\mathrm{\Delta }_0`$ since $`g_0`$ belongs to the centre of $`H_0`$. Since $`\mathrm{\Delta }_0`$ represents a nontrivial element, there exists an integer $`i`$ between $`1`$ and $`m`$ such that the diagram $`A_i`$ is nontrivial and so it is a simple absolutely reduced diagram. Its centralizer is cyclic (see the proof of Theorem 15.35 in ). Let us now take two diagrams $`\mathrm{\Gamma }`$, $`\mathrm{\Xi }`$ that represent elements in $`H_0`$. By the arguments of the above paragraph, there are decompositions of the form $`\mathrm{\Gamma }=B_1+\mathrm{}+B_m`$, $`\mathrm{\Xi }=C_1+\mathrm{}+C_m`$, where $`B_i`$, $`C_i`$ are $`(w_i,w_i)`$-diagrams that commute with $`A_i`$. Cyclicity of the centralizer of $`A_i`$ implies that $`B_i`$ and $`C_i`$ commute, that is, $`[B_i,C_i]=\epsilon (w_i)`$. Therefore $`[\mathrm{\Gamma },\mathrm{\Xi }]=[B_1,C_1]+\mathrm{}+[B_m,C_m]=\mathrm{\Delta }^{}+\epsilon (w_i)+\mathrm{\Delta }^{\prime \prime }`$, where $`\mathrm{\Delta }^{}`$, $`\mathrm{\Delta }^{\prime \prime }`$ are spherical diagrams with bases $`w_1\mathrm{}w_{i1}`$, $`w_{i+1}\mathrm{}w_m`$. It is clear that the product of diagrams of the form $`\mathrm{\Delta }^{}+\epsilon (w_i)+\mathrm{\Delta }^{\prime \prime }`$ is again a diagram of the same form. Hence any element of the commutator subgroup of the group $`H_0`$ has the form $`\mathrm{\Delta }^{}+\epsilon (w_i)+\mathrm{\Delta }^{\prime \prime }`$. This contradicts the condition $`\mathrm{\Delta }_0=A_1+\mathrm{}+A_m`$, where $`A_i\epsilon (w_i)`$. The Theorem is proved. ###### Corollary 15 Any nilpotent subgroup of a diagram group is abelian. Proof. Let $`𝒟(𝒫,w)`$ be a diagram group and let $`K`$ be its nilpotent subgroup. If $`K`$ is not abelian then $`K`$ has a (non-abelian) nilpotent subgroup $`H`$ of degree $`2`$. This means that the commutator subgroup of $`H`$ is contained in its centre, that is, $`H^{}Z(H)`$. Theorem 14 claims that the centre of $`H`$ and its commutator subgroup have trivial intersection so $`H^{}=1`$, that is, $`H`$ is abelian, a contradiction. Let us now describe all abelian subgroups of diagram groups. It turns out that all of them are free abelian. Note that if there was an abelian subgroup in a diagram group that was not free abelian, then we would immediately disprove the Subgroup Conjecture, because it is known from that the quotient of any diagram group by its commutator subgroup is free abelian. ###### Theorem 16 Any abelian subgroup of a diagram group is free abelian. Proof. Let $`𝒫=\mathrm{\Sigma }`$ be a semigroup presentation, let $`G=𝒟(𝒫,w)`$ be a diagram group and let $`HG`$ be an abelian subgroup of $`G`$. If $`H=1`$ then we have nothing to prove. Let $`H1`$. Consider an element $`hH`$, $`h1`$. Using certain conjugation and replacing the subgroup by an isomorphic one, we can assume by \[14, Lemma 15.14\] that $`h`$ is represented by an absolutely reduced normal diagram $`\mathrm{\Delta }`$ that is decomposed into the sum of components: $`\mathrm{\Delta }=\mathrm{\Delta }_1+\mathrm{}+\mathrm{\Delta }_m`$, where $`\mathrm{\Delta }_i`$ is a spherical diagram with base $`u_i`$ ($`1im`$). Since $`H`$ is abelian, it is contained in the centralizer of $`h`$. Thus any element $`g`$ in $`H`$ decomposes into a sum of $`(u_i,u_i)`$-diagrams. Denote by $`\psi _i(g)`$ the $`i`$th summand of this decomposition. It is easy to see that $`\psi _i`$ is a homomorphism of the group $`H`$ into the diagram group over $`𝒫`$ with base $`u_i`$. Let $`1km`$ be a number such that $`\mathrm{\Delta }_k`$ is nontrivial. The centralizer of $`\mathrm{\Delta }_k`$ is cyclic \[14, Theorem 15.35\]. Consider the homomorphism $`\psi _k`$. Firstly, $`\psi _k(h)=\mathrm{\Delta }_k\epsilon (u_i)`$. Secondly, the image of $`\psi _k`$ is contained in the centralizer of $`\mathrm{\Delta }_k`$, that is, in a cyclic group. We thus proved that for any $`hH`$, $`h1`$, there exists a homomorphism $`\psi :H𝐙`$ such that $`\psi (h)1`$. This means that $`H`$ is residually cyclic, that is, it embeds into a Cartesian power of the infinite cyclic group. An easy argument in spirit of linear algebra (using the Choice Axiom) shows that a Cartesian power $`𝐙`$ is a free abelian group. Therefore, $`H`$ is also free abelian. In conclusion of this Section we give a simple but useful generalization of Theorem 15.34 from . We will need it later in Section 6. ###### Theorem 17 Let $`𝒫=\mathrm{\Sigma }`$ be a semigroup presentation and let $`G=𝒟(𝒫,w)`$ for some $`w\mathrm{\Sigma }^+`$. Suppose that $`A_1`$, …, $`A_m`$ are spherical diagrams with base $`w`$ that pairwise commute in $`G`$. Then there exist a word $`v=v_1\mathrm{}v_n`$, spherical $`(v_j,v_j)`$-diagrams $`\mathrm{\Delta }_j`$ $`(1jn)`$, integers $`d_{ij}`$ $`(1im`$, $`1jn)`$ and some $`(w,v)`$-diagram $`\mathrm{\Gamma }`$ such that $$\mathrm{\Gamma }^1A_i\mathrm{\Gamma }=\mathrm{\Delta }_1^{d_{i1}}+\mathrm{}+\mathrm{\Delta }_n^{d_{in}}$$ for all $`1im`$. We can additionally assume that each of the diagrams $`\mathrm{\Delta }_1`$, …, $`\mathrm{\Delta }_n`$ is either trivial or simple absolutely reduced. Proof. First of all let us show that the additional assumption about $`\mathrm{\Delta }_i`$ can be proved provided the main statement of the theorem is proved. If we have already found the decompositions of diagrams from the main statement of the Theorem, then by Lemma 15.14 from we can find diagrams $`\mathrm{\Psi }_j`$ such that diagarams $`\mathrm{\Psi }_j\mathrm{\Delta }_j\mathrm{\Psi }_j^1`$ are normal and absolutely reduced. Each of these diagarams decomposes into a sum of components that are either trivial or simple absolutely reduced. After that we apply additional conjugation by $`\mathrm{\Psi }_1+\mathrm{}+\mathrm{\Psi }_n`$ replacing each of the $`n`$ summands by a sum of components. The rest will be proved by induction on $`m`$. If $`m=1`$ or $`m=2`$ then it is proved in (Lemma 15.10c and Theorem 15.34). So we assume that $`m>2`$, and the statement is true for all values less than $`m`$. Consider the diagram $`A_m`$. Applying Lemma 15.10c, we find a $`(w,u)`$-diagram $`\mathrm{\Psi }`$ such that the diagram $`A_m^{}=\mathrm{\Psi }^1A_m\mathrm{\Psi }`$ will be absolutely reduced and normal. It can be decomposed into the sum of components: $`A_m^{}=B_1+\mathrm{}+B_k`$, where $`B_j`$ is a $`(u_j,u_j)`$-diagram for some $`u_j`$ ($`1jk`$) and $`u=u_1\mathrm{}u_k`$. Each of the diagrams $`B_j`$ ($`1jk`$) is either trivial or simple and all of them are absolutely reduced and normal. Let $`A_i^{}=\mathrm{\Psi }^1A_i\mathrm{\Psi }`$ for all $`1im1`$. It is clear that all diagrams $`A_1^{}`$, …, $`A_m^{}`$ pairwise commute. Since $`A_1^{}`$, …, $`A_{m1}^{}`$ also commute with $`A_m^{}`$, each of them can be decomposed into a sium of $`(u_j,u_j)`$-diagrams ($`1jn)`$ by \[14, Theorem 15.35\]. We have $`A_i^{}=C_{i1}+\mathrm{}+C_{in}`$, where $`1im1`$, each of the diagrams $`C_{ij}`$ is a $`(u_j,u_j)`$-diagram and $`C_{ij}`$ commutes with $`B_j`$ ($`1im1`$, $`1jn`$). Suppose that $`r`$ ($`1rk`$) is a number such that the component $`B_r`$ is nontrivial. Since $`B_r`$ is a simple absolutely reduced diagram, it has a cyclic centralizer (see the proof of Theorem 15.35). In particular, there exists a spherical diagram $`\mathrm{\Delta }_r`$ with base $`u_r`$ such that any of the diagrams $`C_{1r}`$, …, $`C_{m1,r}`$, $`B_r`$ is a power of $`\mathrm{\Delta }_r`$. Let $`C_{ir}=\mathrm{\Delta }_r^{d_{ir}}`$ ($`1im1`$), $`B_r=\mathrm{\Delta }_r^{d_{mr}}`$. We do that with each $`r`$ ($`1rk`$), for which the component $`B_r`$ is nontrivial. Now let $`1rk`$ be such that the component $`B_r`$ is trivial, that is, $`B_r=\epsilon (u_r)`$. Applying the inductive assumption to diagrams $`C_{1r}`$, …, $`C_{m1,r}`$, we find a word $`v_r=v_{1r}\mathrm{}v_{n_r,r}`$, spherical $`(v_j,v_j)`$-diagrams $`\mathrm{\Delta }_{jr}`$ $`(1jn_r)`$, integers $`d_{ijr}`$ $`(1im1`$, $`1jn_r)`$ and some $`(u_r,v_r)`$-diagram $`\mathrm{\Gamma }_r`$ such that $$\mathrm{\Gamma }_r^1C_{ir}\mathrm{\Gamma }_r=\mathrm{\Delta }_{1r}^{d_{i1r}}+\mathrm{}+\mathrm{\Delta }_{n_r,r}^{d_{i,n_r,r}}$$ (18) for all $`1im1`$. Now $`\mathrm{\Gamma }_r^1B_r\mathrm{\Gamma }_r=\epsilon (v_r)=\epsilon (v_{1r})+\mathrm{}+\epsilon (v_{n_r,r})`$, and one can put $`d_{m1r}=\mathrm{}=d_{m,n_r,r}=0`$. Then equality (18) will be true also for $`i=m`$, if we put $`C_{mr}=B_r`$. For numbers $`r`$ such that $`B_r`$ is nontrivial we put $`v_r=u_r`$, and take the trivial diagram for $`\mathrm{\Gamma }_r`$. In this case, we also need to put $`n_r=1`$, $`\mathrm{\Delta }_{1r}=\mathrm{\Delta }_r`$, $`d_{i1r}=d_{ir}`$. Then the equalities (18) are true for all $`1rk`$, $`1im`$. Putting $`\mathrm{\Gamma }=\mathrm{\Psi }(\mathrm{\Gamma }_1+\mathrm{}+\mathrm{\Gamma }_k)`$, we see that $$\mathrm{\Gamma }^1A_i\mathrm{\Gamma }=\underset{r=1}{\overset{k}{}}\underset{j=1}{\overset{n_r}{}}\mathrm{\Delta }_{jr}^{d_{ijr}},$$ for $`1im`$, that is, we have got the required decomposition of diagrams into a sum. The proof is complete. In we established that the conjugacy problem is decidable for any diagram group $`𝒟(𝒫,w)`$, where $`𝒫=\mathrm{\Sigma }`$ is a semigroup presentation with decidable word problem. In particular, this implies the decidability of the conjugacy problem in R. Thompson’s group $`F`$. We can pose a more general problem — a uniform conjugacy problem for sequences. ###### Problem 1 Let $`𝒫=\mathrm{\Sigma }`$ be a semigroup presentation with decidable word problem, $`w\mathrm{\Sigma }^+`$, $`G=𝒟(𝒫,w)`$. Does there exist an algorithm that decides, given two sequences of elements $`x_1`$, …, $`x_n`$ and $`y_1`$, …, $`y_n`$ of the group $`G`$ $`(`$elements are represented by diagrams$`)`$, whether there is an element $`zG`$ such that $`x_i^z=y_i`$ for all $`i`$ from $`1`$ to $`n`$? In particular, is this problem decidable for R. Thompson’s group $`F`$? Note that there is some analogy between diagram groups and matrix groups (we remarked about this in ). The corresponding question for matrix groups was solved positively in and independently in . ## 4 Soluble Subgroups in Diagram Groups In this section we shed some light on the structure of soluble subgroups in diagram groups. First of all let us consider an example that demonstrates that there exist soluble subgroups of any degree in R. Thompson’s group $`F`$. Let $`𝒫=\{xxx=x\}`$. All groups $`𝒟(𝒫,x^k)`$, where $`k=1,2,\mathrm{}`$ , are isomorphic to $`F`$. Consider any nontrivial $`(x,x)`$-diagram $`\mathrm{\Delta }`$. Then the diagrams $`\mathrm{\Delta }_1=\epsilon (x^2)+\mathrm{\Delta }+\epsilon (x)`$ and $`\mathrm{\Delta }_2=\epsilon (x)+\mathrm{\Delta }+\epsilon (x^2)`$ are conjugate because they are sums of components that conjugate respectively. Secondly, the diagrams commute which can be seen directly and are conjugated by the diagram $`\mathrm{\Gamma }=(x^2x)+\epsilon (x)+(xx^2)`$. Denoting by $`a`$, $`b`$ the elements of $`𝒟(𝒫,x^4)`$ that represent diagrams $`\mathrm{\Delta }_1`$ and $`\mathrm{\Gamma }`$, respectively, we can see that $`a`$ and $`b`$ generate the group $`𝐙wr𝐙`$, a (restricted) wreath product of two infinite cyclic groups. In the following theorem we present a more general form of the above construction. ###### Theorem 18 Let $`𝒫=\mathrm{\Sigma }`$ be a semigroup presentation. Suppose that there exist nonempty words $`x`$, $`y`$, $`z`$ over $`\mathrm{\Sigma }`$ such that $`xy=x`$, $`yz=z`$ modulo $`𝒫`$, and suppose that the diagram group $`𝒟(𝒫,y)`$ is nontrivial. Then the group $`G=𝒟(𝒫,xyz)`$ contains a subgroup isomorphic to $`𝐙wr𝐙`$. Namely, let $`\mathrm{\Delta }`$ be any nontrivial $`(y,y)`$-diagram, let $`\mathrm{\Gamma }_1`$ be arbitrary $`(xy,x)`$-diagram, and let $`\mathrm{\Gamma }_2`$ be arbitrary $`(z,yz)`$-diagram. Then elements $`a`$ and $`b`$, represented by diagrams $`\epsilon (x)+\mathrm{\Delta }+\epsilon (z)`$ and $`\mathrm{\Gamma }_1+\mathrm{\Gamma }_2`$, respectively, generate in $`G`$ a subgroup isomorphic to $`𝐙wr𝐙`$. Proof. First of all, let us mention that elements $`xz`$, $`xyz`$, $`xy^2z`$, … are equal modulo $`𝒫`$ so the diagram groups with these bases over $`𝒫`$ will be isomorphic to each other. To get the diagrams from the above example, one needs to put $`x=y=z`$ and then go from the group with base $`xyz=x^3`$ to the group with base $`xy^2z=x^4`$ using conjugation by the element $`(xx^2)+\epsilon (x^2)`$. To show that elements $`a`$, $`b`$ of some group $`G`$ generate $`𝐙wr𝐙`$, it suffices to show that the elements $`a_i=a^{b^i}`$ ($`i𝐙`$) form a free basis of a free abelian group. To check this, it suffices to show that for any positive integer $`n`$, the elements $`a_0`$, $`a_1`$, …, $`a_n`$ form a basis of the free abelian group they generate. We will explicitly find the elements $`a_i`$ ($`0in`$) of the corresponding diagram group $`𝒟(𝒫,xyz)`$. For convenience, we will go to the diagram group over $`𝒫`$ with base $`xy^{n+1}z`$ using conjugation by the diagram $$\mathrm{\Psi }=(\mathrm{\Gamma }_1^1+\epsilon (yz))(\mathrm{\Gamma }_1^1+\epsilon (y^2z))\mathrm{}(\mathrm{\Gamma }_1^1+\epsilon (y^nz)).$$ One can check directly that $`c_i=\mathrm{\Psi }^1a_i\mathrm{\Psi }=\epsilon (xy^{ni})+\mathrm{\Delta }+\epsilon (y^iz)`$ for all $`0in`$. It is clear that the elements $`c_0`$, $`c_1`$, …, $`c_n`$ pairwise commute. The obvious formula $$c_0^{d_0}c_1^{d_1}\mathrm{}c_n^{d_n}=\epsilon (x)+\mathrm{\Delta }^{d_0}+\mathrm{\Delta }^{d_1}+\mathrm{}+\mathrm{\Delta }^{d_n}+\epsilon (z)$$ shows that the elements $`c_0`$, $`c_1`$, …, $`c_n`$ form a basis of the free abelian subgroup in $`𝒟(𝒫,xy^{n+1}z)`$. So the elements $`a_0`$, $`a_1`$, …, $`a_n`$ also form a basis of a free abelian subgroup of $`𝒟(𝒫,xyz)`$ as desired. We will return to the group $`𝐙wr𝐙`$ later. Now we shall prove a simple fact that will imply that $`F`$ contains soluble subgroups of any degree. ###### Lemma 19 $`(`$Restricted$`)`$ wreath product $`Fwr𝐙`$ is a subgroup of $`F`$. Proof. We will use some known properties of R. Thompson’s group $`F`$ mentioned in Section 1. As we mentioned above, there are several representations of $`F`$ by piecewise linear functions. Let us consider the representation by functions on $`[0,\mathrm{})`$. For any positive integer $`k`$ we consider the functions from $`F`$ that are identical outside $`[k,k+1]`$. By $`\mathrm{\Phi }_k`$ we denote the set of all these functions. It is obvious that they form a group isomorphic to the group of all piecewise linear functions on $`[0,1]`$ (with the properties mentioned in Section 1), that is, it is isomorphic to $`F`$. It is also easy to see that elements in different subgroups $`\mathrm{\Phi }_k`$ commute with each other. Therefore, the groups $`\mathrm{\Phi }_k`$ ($`k1`$) generate a direct power of the group $`F`$. Conjugation by the element $`x_0`$, represented by the function given by $`tx_0=2t`$ ($`t[0,1]`$), $`tx_0=t+1`$ ($`t1`$), takes $`\mathrm{\Phi }_k`$ to $`\mathrm{\Phi }_{k+1}`$. It is now clear that $`t`$ and $`\mathrm{\Phi }_k`$ ($`k1`$) generate the restricted wreath product $`Fwr𝐙`$ in $`F`$. It is not hard to find generators of the subgroup $`Fwr𝐙`$ of $`F`$ in a diagram form and also in a normal form. In particular, the subgroup in $`F`$ generated by elements $`x_0`$, $`x_1x_2x_1^2`$, $`x_1^2x_2x_1^3`$ will be isomorphic to $`Fwr𝐙`$. The reader can easily draw the diagrams representing these elements. Let us define a sequence of groups by induction: $`H_1=𝐙`$, $`H_{n+1}=H_nwr𝐙`$. Thus groups $`H_n=(\mathrm{}(𝐙wr𝐙)wr\mathrm{})wr𝐙`$, where $`𝐙`$ occurs $`n`$ times, are diagram groups by Theorem 11. The group $`H_n`$ is soluble of degree $`n`$. Using Lemma 19 and elementary properties of wreath products, we have the following result that can be proved by induction on $`n`$. This statement was obtained by M. Brin (private communication), see also . ###### Corollary 20 For any $`n`$, the group $`H_n=(\mathrm{}(𝐙wr𝐙)wr\mathrm{})wr𝐙`$ is a soluble subgroup of degree $`n`$ in R. Thompson’s group $`F`$. Talking about wreath products, we would like to mention a fact about subgroups of R. Thompson’s group $`F`$. It was shown in that any subgroup of $`F`$ is either metabelian or contains an infinite direct power of the group $`𝐙`$. (In fact, one can replace the word “metabelian” by “abelian”, see .) The proof given in uses representations of $`F`$ by piecewise linear functions. Actually, the result is obtained for subgroups of some group which is bigger than $`F`$. It turns out that one can extract a stronger fact from this proof. Consider all piecewise linear continuous transformations of the unit interval $`I=[0,1]`$ onto itself. We consider only mappings that preserve orientation and have finitely many breaks of the derivative. All these functions form a group with respect to composition. Let us denote this group by $`PL_0(I)`$. It contains $`F`$ as a subgroup. We have the following alternative for subgroups of the group $`PL_0(I)`$. ###### Theorem 21 Any subgroup of $`PL_0(I)`$ is either abelian, or contains an isomorphic copy of $`𝐙wr𝐙`$. Proof. Our proof basicaly follows the proof or a weaker alternative from . For $`fPL_0(I)`$ by $`suppf`$ we denote the set of all $`tI`$, for which $`tft`$. Let $`G`$ be a non-abelian subgroup of $`PL_0(I)`$. Consider functions $`f,gG`$ such that $`fggf`$. Let $`J=suppfsuppg`$. It is obvious that $`J`$ is a union of finitely many disjoint intervals $`J_k=(a_k,b_k)`$, $`1km`$. By definition, $`[f,g]1`$ in $`PL_0(I)`$. Then on some of intervals $`J_1`$, …, $`J_m`$ our function $`[f,g]`$ is not the identity. Denote by $`\nu (f,g)`$ the number of such intervals. Without loss of generality, we can assume that the elements $`f,gG`$ which do not commute are chosen in such a way that the value $`\nu (f,g)`$ is the smallest possible. Let $`H`$ be a subgroup of $`PL_0(I)`$ generated by $`f`$ and $`g`$. By definition, the endpoints of $`J_1`$, …, $`J_m`$ are stable under $`f`$ and $`g`$ so each of these intervals is $`H`$-invariant. An easy argument shows that for any $`x,yJ_k`$ ($`1km`$), where $`x<y`$, there exists a function $`wH`$ such that $`xw>y`$. Let us take the greatest upper bound $`z`$ of the set $`\{xhhH\}`$. It is clear that $`a_i<zb_i`$. If $`zb_i`$ then either $`zfz`$ or $`zgz`$ by definition of the set $`J`$. Without loss of generality let $`zfz`$. This inequality also holds in a small neighbourhood of the point $`z`$. Therefore one of the numbers $`zf`$ or $`zf^1`$ is greater than $`z`$, a contradiction. Thus $`z=b_i`$. This implies that acting by some element of $`H`$ one can make the image of $`x`$ as close to $`b_i`$ as one wishes which is what we had to prove. Let us take an interval $`(a_i,b_i)`$ ($`1im`$) such that $`[f,g]`$ is not identical on it. It is easy to ee that the function $`[f,g]`$ is identical in some neighbourhood of each of the points $`a_i`$, $`b_i`$. Thus $`supp[f,g]`$ is nonempty and it is contained in $`[c_0,d_0]`$, where $`a_i<c_0<d_0<b`$. According to the above, there exists a function $`wH`$ such that $`d_0<c_0w<b`$. Let us denote $`[f,g]`$ by $`h_0`$. For any $`n1`$, let $`c_n=c_0w^n`$, $`d_n=d_0w^n`$, $`h_n=h_0^{w^n}`$. It is obvious that $`c_0<d_0<c_1<d_1<\mathrm{}`$, and $`supph_nJ_i[c_n,d_n]`$. Therefore, for any $`i,j0`$, the commutator $`[h_i,h_j]`$ is identical on $`J_i`$. Suppose that $`[h_i,h_j]1`$ for some $`i`$, $`j`$. Since all the intervals $`J_1`$,…, $`J_m`$ are $`H`$-invariant, it is clear that all the functions $`h_1`$, $`h_2`$, … are identical on all the intervals $`J_k`$ ($`1km`$) where the function $`h_0`$ is identical. Therefore, we can replace $`f`$, $`g`$ by $`h_i`$, $`h_j`$ obtaining $`\nu (h_i,h_j)<\nu (f,g)`$. This contradicts the choice of $`f`$, $`g`$. This proves that $`[h_i,h_j]=1`$ for any $`i,j0`$. So far we very closely followed the proof given in . The conclusion in , is that elements $`h_0`$, $`h_1`$, $`h_2`$, … form a basis of a free abelian group. To prove a stronger statement of our theorem, it remains now to add that $`h_n^w=h_{n+1}`$ for all $`n0`$, so the elements $`h_0`$, $`w`$ generate $`𝐙wr𝐙G`$. The Theorem is proved. Using Theorem 16, we have the following alternative for subgroups of R. Thompson’s group $`F`$. ###### Corollary 22 Any subgroup of R. Thompson’s group $`F`$ is either free abelian or contains the restricted wreath product $`𝐙wr𝐙`$. Note that in the group of all piecewise linear functions, not every abelian subgroup is free abelian. We can extract one more corollary from Theorem 21. ###### Corollary 23 A non-abelian group with one defining relation cannot be a subgroup of the group $`PL_0(I)`$ $`(`$in particular, it cannot be a subgroup of R. Thompson’s group $`F)`$. Proof. In , A. A. Chebotar described all subgroups of one-relator groups that do not contain free subgroups of rank $`2`$. They are: a) abelian subgroups, b) free product $`𝐙_2𝐙_2`$ of the cyclic group of order $`2`$ by itself, and c) Baumslag – Solitar groups $`B_{1k}=a,bb^1ab=a^k`$. It is clear that the group $`𝐙wr𝐙`$ does not occur in this list. Thus a non-abelian one-relator group cannot be a subgroup of $`PL_0(I)`$ by Theorem 21. The Corollary is proved. Now let us consider the following interesting question: under what condition on a semigroup presentation $`𝒫`$ a diagram group over this presentation contains $`zzwr𝐙`$ as a subgroup? The answer is given in the following theorem. ###### Theorem 24 Let $`𝒫=\mathrm{\Sigma }`$ be a semigroup presentation and let $`G=𝒟(𝒫,w)`$ be a diagram group. Then the following three conditions are equivalent. 1. The group $`G=𝒟(𝒫,w)`$ contains $`𝐙wr𝐙`$ as a subgroup. 2. The group $`G`$ contains elements $`a`$, $`b`$ such that $`[a,b]1`$, $`[a,a^b]=1`$ $`(`$in other words, there are two distinct elements in $`G`$ that are conjugate and commute$`)`$. 3. There are words $`x`$, $`y`$, $`z`$ such that the equalities $`xy=x`$, $`yz=z`$, $`xz=w`$ hold modulo $`𝒫`$, and $`𝒟(𝒫,y)1`$. Proof. The proof uses the following scheme: $`(1)(2)(3)(1)`$. The implication $`(1)(2)`$ is obvious and holds for any group $`G`$. The implication $`(3)(1)`$ was proved in Theorem 18. It remains to show that $`(2)(3)`$. Suppose that $`G=𝒟(𝒫,w)`$ has elements $`a`$, $`b`$ such that $`[a,b]1`$, $`[a,a^b]=1`$. By Theorem 17, one can pass from the base $`w`$ to some base $`v`$ that equals $`w`$ modulo $`𝒫`$ in such a way that the diagarams representing the two given commuting elements will be absolutely reduced and normal. Without loss of generality we can assume that $`a`$ is represented by a diagram $`A=A_1+\mathrm{}+A_m`$ decomposed into the sum of components, and the element $`a^b`$ is represented by a diagarm $`C`$ that has a decomposition into the sum of the same number of components: $`C=C_1+\mathrm{}+C_m`$ by \[14, Lemma 15.15\]. By the same lemma, the element $`b`$ which conjugates $`A`$ and $`C`$ is represented by a diagram $`B`$ of the form $`B_1+\mathrm{}+B_m`$, where $`C_i=B_i^1A_iB_i`$. Let $`v=v_1\mathrm{}v_m=v_1^{}\mathrm{}v_m^{}`$, where $`A_i`$, $`C_i`$ are spherical diagrams with bases $`v_i`$, $`v_i^{}`$, respectively, and let $`B_i`$ be a $`(v_i,v_i^{})`$-diagram ($`1im`$). It is obvious that $`v_i=v_i^{}`$ modulo $`𝒫`$ for all $`i`$. We would like to prove that there exists an $`i`$ from $`1`$ to $`m`$ such that the diagram $`A_i`$ (and $`C_i`$ as well) is nontrivial and the occurrences of the words $`v_i`$, $`v_i^{}`$ in the word $`v`$ do not overlap and do not contain each other. This would imply condition $`(3)`$. Indeed, without loss of generality, let $`v=pv_iqv_i^{}r`$. Then the equalities $`p=v_1\mathrm{}v_{i1}=v_1^{}\mathrm{}v_{i1}^{}=pv_iq`$, $`r=v_{i+1}^{}\mathrm{}v_m^{}=v_{i+1}\mathrm{}v_m=qv_i^{}r=qv_ir`$ hold modulo $`𝒫`$. One can put $`x=p`$, $`y=v_iq`$, $`z=v_ir`$, and then the equalities $`x=p=pv_iq=xy`$, $`z=v_ir=v_iqv_ir=yz`$, $`w=v=pv_iqv_i^{}r=pv_iqv_ir=pv_ir=xz`$ will hold modulo $`𝒫`$ (that is, in the semigroup $`S`$). The diagram group over $`𝒫`$ with base $`y=v_iq`$ will be definitely nontrivial because there exists a nontrivial spherical diagram $`A_i+\epsilon (q)`$ with this base. Let us prove the existence of an $`i`$ such that $`A_i`$ is nontrivial, and the occurrences of the word $`v_i`$, $`v_i^{}`$ in the word $`v`$ have no common letters. Let us consider an arbitrary $`1im`$ such that $`A_i`$ is nontrivial. Since $`C`$ commutes with $`A`$, the diagram $`C`$ can be presented as a sum $`C^{}+D+C^{\prime \prime }`$, where $`C^{}`$, $`D`$, $`C^{\prime \prime }`$ are spherical diagrams with bases $`v_1\mathrm{}v_{i1}`$, $`v_i`$, $`v_{i+1}\mathrm{}v_m`$, respectively, and $`D`$ commutes with $`A_i`$. Therefore, the diagram $`D`$ consists of one component. The same is true for the diagram $`C_i`$. So, if the occurrences of $`v_i`$, $`v_i^{}`$ have common letters, then diagrams $`D`$ and $`C_i`$ must coincide. This implies that $`A_i`$, $`C_i`$ are powers of the same element $`\mathrm{\Delta }_i`$, and are conjugate. It follows from results of \[14, Section 15\] that $`A_i=C_i`$, and the occurences of $`v_i`$ and $`v_i^{}`$ coincide. It is now obvious that $`A=C`$. But this contradicts the assumption that $`aa^b`$. The proof is complete. Let us make two remarks about the theorem we have just proved. Firstly, in the third condition we cannot avoid the condition that the diagram over $`𝒫`$ with base $`y`$ is nontrivial. Without this condition, the diagram group may coincide with $`𝐙`$. Note that the algorithm to verify whether the diagram group over a given finite presentation with given base is nontrivial, is unknown. Secondly, we have to note that if elements $`a`$, $`b`$ of a diagram group are such that $`[a,b]1`$, $`[a,a^b]=1`$, then the subgroup isomorphic to $`𝐙wr𝐙`$, is not necessarily contained in the subgroup generated by $`a`$, $`b`$. An example illustrating that is given below in Section 5. Finishing this Section, let us give a sufficient condition for a diagram group to contain R. Thompson’s group $`F`$ as a subgroup. ###### Theorem 25 Let $`𝒫=\mathrm{\Sigma }`$ be a semigroup presentation, and let the semigroup $`S`$ presented by $`𝒫`$ contain an idempotent. Then there is a word $`w`$ such that the diagram group $`G=𝒟(𝒫,w)`$ contains R. Thompson’s group $`F`$ as a subgroup. Moreover, one can take any word $`w`$ that represents an idempotent in $`S`$ for such a base. Proof. The proof is quite easy. It is based on the fact that all proper homomorphic images of the group $`F`$ are abelian (see ). Let us take a word $`w`$ such that $`ww=w`$ modulo $`𝒫`$. Let us consider a reduced $`(w^2,w)`$-diagram $`\mathrm{\Delta }`$ over $`𝒫`$. Now we construct a homomorphism from $`F`$ to $`𝒟(𝒫,w)`$ in the following way. We use the fact that $`F`$ is a diagram group over $`𝒬=xx^2=x`$. It is convenient to take the element $`x^5`$ as a base. To any diagram over $`𝒬`$, we assign a diagram over $`𝒫`$, replacing the label $`x`$ by $`w`$ and filling in the cells of the form $`x^2=x`$ by diagrams $`\mathrm{\Delta }`$. This rule defines a homomorphism from $`F`$ to $`𝒟(𝒫,w)`$. Taking into account what we have said above, it is enough to check that the image of this homomorphism is not abelian. To show this, it suffices to compute the image of the commutator $`[x_0,x_1]`$. In the group $`𝒟(𝒬,x^5)`$, this commutator is represented by the diagram $`\epsilon (x)+(x^2=x)+(x=x^2)+\epsilon (x)`$. It is obvious that after we replace all $`(x^2,x)`$-cells in this diagram by copies of $`\mathrm{\Delta }`$, we get a diagram without dipoles. Thus the image of this commutator is not equal to the identity, so the image of $`F`$ under the homomorphism is not abelian. The Theorem is proved. We do not know whether the condition on $`S`$ to have an idempotent is also sufficient. ###### Problem 2 Let $`𝒫=\mathrm{\Sigma }`$ be a semigroup presentation and let $`G=𝒟(𝒫,w)`$ be a diagram group. Suppose that $`G`$ contains R. Thompson’s group $`F`$ as a subgroup. Is it true that the semigroup $`S`$ presented by $`𝒫`$ contains an idempotent? ## 5 The Subgroup Conjecture In this Section, we construct a counterexample to the Subgroup Conjecture, that is, we will construct a subgroup in a diagram group that is not a diagram group itself. Note that the first candidate to disprove the Subgroup Conjecture was the group $`F^{}`$ — the commutator subgroup of R. Thompson’s group $`F`$. However, it turned out that $`F^{}`$ is a diagram group. This answered some open questions about diagram groups. Before proving the corresponding Theorem, let us make the following two remarks. The first remark is about properties of semigroup diagrams over certain special presentations. Let we have a semigroup presentation of the form $`X`$, where all relations in $``$ have the form $`u=V`$, where $`uX`$, $`VX^+`$. We assume that all left-hand sides of the relations are distinct and the right-hand sides contain more than one letter. It is known (see, for instance ) that any reduced diagram $`\mathrm{\Delta }`$ over such a presentation can be uniquely decomposed into a concatenation: $`\mathrm{\Delta }=\mathrm{\Delta }_1\mathrm{\Delta }_2`$, where $`\mathrm{\Delta }_1`$ corresponds to a derivation where only applications of relations of the form $`u=V`$ from $``$ are used, and $`\mathrm{\Delta }_2`$ corresponds to a derivation where only applications of relations of the form $`V=u`$ are used, $`(u=V)`$. This fact can be easily proved by choosing the longest positive path from $`\iota (\mathrm{\Delta })`$ to $`\tau (\mathrm{\Delta })`$. It is easy to see that all cells “above” this path will correspond to relations $`u=V`$, and all cells “below” the path will correspond to $`V=u`$. We will call $`\mathrm{\Delta }_1`$ (resp. $`\mathrm{\Delta }_2`$) the positive (resp. negative) part of diagram $`\mathrm{\Delta }`$. Note that the presentation $`xx=x^2`$ satisfies the above conditions. The same holds for the presentation below from the statement of Theorem 26. The second remark is about the structure of a commutator subgroup of a diagram group. It is described in \[14, Theorem 11.3\]. Let us recall this description. Let $`\mathrm{\Delta }`$ be a $`(w,w)`$-diagram over $`𝒫=\mathrm{\Sigma }`$. By $`M`$, we denote the monoid presented by $`𝒫`$. We consider the free abelian group $`𝒜`$ with $`M\times \times M`$ as a free basis. For each vertex $`\mu `$ of diagram $`\mathrm{\Delta }`$ we take any positive path from $`\iota (\mathrm{\Delta })`$ to $`\mu `$. Its label defines an element in $`\mathrm{}(\mu )M`$. It is easy to show that this element does not depend on the choice of a path. Analogously, we define the element $`r(\mu )M`$ as the value of the label of any positive path from $`\mu `$ to $`\tau (\mathrm{\Delta })`$. Now to each cell $`\pi `$ of the diagram $`\mathrm{\Delta }`$ we assign an element $`\delta (\mathrm{}(\iota (\pi )),u=v,r(\tau (\pi )))`$, where $`\delta =1`$, if $`u=\phi (\text{top}(\pi ))`$, $`v=\phi (\text{bot}(\pi ))`$, $`(u,v)`$ and $`\delta =1`$, if $`v=\phi (\text{top}(\pi ))`$, $`u=\phi (\text{bot}(\pi ))`$, $`(u,v)`$. By $`\rho (\mathrm{\Delta })`$ we denote the sum of elements assigned to all the cells of diagram $`\mathrm{\Delta }`$. Thus $`\rho `$ defines a homomorphism from the group $`G=𝒟(𝒫,w)`$ into $`𝒜`$. As it is shown in , the kernel of $`\rho `$ is exactly $`G^{}`$ — the commutator subgroup of $`G`$. ###### Theorem 26 The commutator subgroup of R. Thompson’s group $`F`$ is a diagram group. Namely, $`F^{}𝒟(𝒬,a_0b_0)`$, where $$𝒬=x,a_i,b_i(i0)x=xx,a_i=a_{i+1}x,b_i=xb_{i+1}(i0).$$ Proof. Here is the direct proof of this proposition. Let us construct a mapping from $`H=𝒟(𝒬,a_0b_0)`$ to $`𝒟(𝒫,x^2)`$, where $`𝒫=xx=x^2`$. To each spherical diagram over $`𝒬`$ with base $`a_0b_0`$ we assign a diagram that is obtained from the previous one replacing all its labels by $`x`$. It is clear that we get a spherical diagram over $`𝒫`$ with base $`x^2`$. Obviously, this induces a homomorphism $`\psi :HF`$ because $`𝒟(𝒫,x^2)F`$. Our aim is to prove that the homomorphism $`\psi `$ is injective and its image is $`F^{}`$. Let $`\mathrm{\Delta }`$ be a nontrivial reduced $`(a_0b_0,a_0b_0)`$-diagram over $`𝒬`$. Its image under $`\psi `$ cannot contain dipoles. Otherwise, the preimages of the cells that form a dipole in $`\mathrm{\Delta }`$, would form a dipole themselves. This implies that $`\psi `$ is injective. Let us check that $`\psi (\mathrm{\Delta })F^{}`$ for any reduced diagram $`\mathrm{\Delta }`$ in $`H`$. The monoid $`M`$ presented by $`xx=xx`$ consists of two elements $`1`$ and $`x`$. A cell $`\pi `$ of diagram $`\mathrm{\Delta }^{}=\psi (\mathrm{\Delta })`$ satisfies $`\mathrm{}(\iota (\pi ))=1`$ if and only if $`\iota (\pi )=\iota (\mathrm{\Delta }^{})`$. Analogously, $`r(\tau (\pi ))=1`$ if and only if $`\tau (\pi )=\tau (\mathrm{\Delta }^{})`$. As we know, the diagram $`\mathrm{\Delta }^{}`$ is reduced. It can be decomposed into a product $`\mathrm{\Delta }_1^{}\mathrm{\Delta }_2^{}`$ of its positive and negative part according to the first remark before the statement of the Theorem. It is easy to see that there are no cells $`\pi `$ of the diagram $`\mathrm{\Delta }^{}`$ can satisfy both of the conditions $`\iota (\mathrm{\Delta })=\iota (\pi )`$, $`\tau (\mathrm{\Delta })=\tau (\pi )`$ simultaneously (recall that the base of $`\mathrm{\Delta }^{}`$ is $`x^2`$). Consider the diagram $`\mathrm{\Delta }`$ and decompose it into a concatenation of positive and negative part: $`\mathrm{\Delta }=\mathrm{\Delta }_1\mathrm{\Delta }_2`$ (this is possible because $`𝒬`$ satisfies conditions of the first remark before the statement of the theorem). Let $`a_n`$ ($`n0`$) be the label of the first edge of the path that cuts $`\mathrm{\Delta }`$ into a positive part and a negative part. Then it is easy to extract from the form of the defining relations that all labels of edges which start at $`\iota (\mathrm{\Delta })`$, if one reads them from the top to the bottom of the diagram, are $`a_0`$, $`a_1`$, …, $`a_n`$, …, $`a_1`$, $`a_0`$. From this, it follows that the number of cells $`\pi `$ such that $`\iota (\pi )=\iota (\mathrm{\Delta })`$, is the same for $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_2`$. ¿From this fact, we immediately conclude that it is the same for $`\mathrm{\Delta }^{}`$, if we compare the number of these cells in $`\mathrm{\Delta }_1^{}`$ and $`\mathrm{\Delta }_2^{}`$. All these cells from $`\mathrm{\Delta }_1^{}`$ map to $`(1,x=x^2,x)`$ under $`\rho `$ and all cells from $`\mathrm{\Delta }_2^{}`$ map to $`(1,x=x^2,x)`$ (we emphasize the fact that we are talking about the cells whose initial vertices coincide with the initial vertex of the diagram). An analogous argument can be applied to the cells whose terminal vertices coincide with the terminal vertex the diagram. Here the list of labels of the edges that come into $`\tau (\mathrm{\Delta })`$, if one reads them from top to bottom, is $`b_0`$, $`b_1`$, …, $`b_m`$, …, $`b_1`$, $`b_0`$ for some $`m0`$. Now one can use the fact that the cells of $`\mathrm{\Delta }_1^{}`$ which we deal with in this paragraph map to $`(x,x=x^2,1)`$ and the cells of $`\mathrm{\Delta }_2^{}`$ map to $`(x,x=x^2,1)`$. It remains to note that $`\mathrm{\Delta }_1^{}`$ has the same number of cells as $`\mathrm{\Delta }_2^{}`$ since $`\mathrm{\Delta }^{}`$ is spherical and all relations have the same form. Therefore, the number of cells $`\pi `$ of $`\mathrm{\Delta }_1^{}`$ that satisfy $`\iota (\pi )\iota (\mathrm{\Delta }^{})`$ and $`\tau (\pi )\tau (\mathrm{\Delta }^{})`$ is the same as the number of corresponding cells in $`\mathrm{\Delta }_2^{}`$. However, the first ones map to $`(x,x=x^2,x)`$ and the second ones map to $`(x,x=x^2,x)`$. Hence $`\rho (\mathrm{\Delta }^{})=0`$. This proves that $`\psi (\mathrm{\Delta })F^{}`$. It remains to show that every element in $`F^{}`$ belongs to the image of $`\psi `$. In order to do that, let us take a reduced spherical diagram $`\mathrm{\Delta }^{}`$ with base $`x^2`$ over $`xx=x^2`$. It follows from $`\mathrm{\Delta }^{}F^{}`$ that $`\rho (\mathrm{\Delta }^{})=0`$. We consider separately the cells of three types: those that map into a) $`\pm (1,x=x^2,x)`$, b) $`\pm (x,x=x^2,1)`$, c) $`\pm (x,x=x^2,x)`$ under $`\rho `$, respectively. Each cell belongs to exactly one of the three types. So it is clear that the sum over all cells of each of the types equals zero. This means that in the decomposition $`\mathrm{\Delta }^{}=\mathrm{\Delta }_1^{}\mathrm{\Delta }_2^{}`$ into positive and negative part, the number of cells of each of the types in $`\mathrm{\Delta }_1^{}`$ will be the same as the number of cells of the same type in $`\mathrm{\Delta }_2^{}`$. Let us have $`n0`$ cells of the first type in each of the parts. Let us rename labels of the edges that go out of $`\iota (\mathrm{\Delta }^{})`$, replacing them by $`a_0`$, $`a_1`$, …, $`a_n`$, …, $`a_1`$, $`a_0`$, respectively, from top to bottom. Analogously, let us have $`m0`$ cells of the second type in each of the parts. We rename labels of the edges that come into $`\tau (\mathrm{\Delta }^{})`$, replacing them in the same way by $`b_0`$, $`b_1`$, …, $`b_m`$, …, $`b_1`$, $`b_0`$, respectively. The diagram we get as a result will be denoted by $`\mathrm{\Delta }`$. It remains to note that $`\mathrm{\Delta }`$ will be a spherical diagram over $`𝒬`$ with base $`a_0b_0`$. Indeed, any cell that has the same initial vertex as the one of $`\mathrm{\Delta }`$, corresponds to a relation of the form $`a_i=a_{i+1}x`$ ($`i0`$) or its inverse. If a cell $`\pi `$ of the positive part is taken, then $`\text{top}(\pi )=e`$, $`\text{bot}(\pi )=e_1e_2`$, where $`e`$, $`e_1`$, $`e_2`$ are edges of the diagram. By our construction, the label of $`e`$ equals $`a_i`$ for some $`i0`$. It follows from the way we renamed the labels that $`e_1`$ has label $`a_{i+1}`$. Note that the initial vertex of $`e_2`$ is not $`\iota (\mathrm{\Delta })`$ because $`e_1`$ cannot be a loop. Also the terminal vertex of $`e_2`$ is not $`\tau (\mathrm{\Delta })`$. Otherwise the edge $`e`$ connects the initial and the terminal vertex of $`\mathrm{\Delta }`$ but this is impossible. Therefore, the label of $`e_2`$ is $`x`$. The arguments for the negative part of the diagram are analogous. Of course, any cell that has the same terminal vertex as $`\mathrm{\Delta }`$, corresponds to a relation of the form $`b_i=xb_{i+1}`$ ($`i0`$) or its inverse. It is clear that $`\psi `$ takes $`\mathrm{\Delta }`$ into $`\mathrm{\Delta }^{}`$. This completes the proof. ###### Corollary 27 A diagram group can be simple. In particular, there exist nontrivial diagram groups that coincide with their commutator subgroups and so they do not admit an LOG-presentation. The group $`F^{}`$ is simple (see ). We proved in 26 that $`F^{}`$ is a diagram group. In \[14, Section 17\] we asked if a nontrivial diagram group may coincide with its commutator subgroup. We have given a positive answer. This is interesting if to compare this result with \[14, Theorem 12.1\]. It was proved there that if all diagram groups over a semigroup presentation coincide with their commutator subgroups, then all of them are trivial. As we see, certain diagram groups may coincide with their commutator subgroups. As a by-product, we gave an answer to Problem 17.1 of the same paper: is it true that any diagram group admits an LOG-presentation? Recall that an LOG-presentation is a group presentation such that all defining relations have form $`a=b^c`$, where $`a`$, $`b`$ and $`c`$ are generators. The groups that admit such a presentation are called LOG-groups (this concept was introduced in , where these groups were characterized in terms of labelled oriented graphs). In Russian papers, one can often meet an equivalent terminology “C-group”. Some interesting characterization of these groups was recently obtained by Yu. V. Kuzmin . We have already shown that any diagram group over a complete presentation (see ) admits an LOG-presentation (cf R. Thompson’s group $`F`$). We also proved that any diagram group is a retract of an LOG-group. Since any LOG-group has $`𝐙`$ as its homomorphic image, it cannot coincide with its commutator subgroup. Thus we proved that a diagram group may not have an LOG-presentation. To construct a counterexample to the Subgroup Conjecture, we strongly use results of Section 2. In particular, we need Theorem 13 and Example 12. ###### Theorem 28 There exist subgroups of digaram groups that are not diagram groups themselves. For instance, the following one-relator group $$x,yxy^2x=yx^2y$$ can be isomorphically represented by diagrams over a semigroup presentation but it is not a diagram group itself. Proof. Let $`L=x,yxy^2x=yx^2y`$. We shall prove that $`L`$ is not a diagram group. Consider the group $`L_{\mathrm{}}`$ given by $$L_{\mathrm{}}=z_i(i𝐙)[z_i,z_{i+1}]=1(i𝐙).$$ The mapping $`\psi `$ that takes $`z_i`$ to $`z_{i+1}`$ for all $`i𝐙`$, obviously induces an automorphism of the group $`L_{\mathrm{}}`$. Consider HNN-extension of the group $`L_{\mathrm{}}`$ with stable letter $`t`$ via automorphism $`\psi `$ (this will be also a semidirect product of $`L_{\mathrm{}}`$ and $`𝐙`$). We obtain the group $`L_{\mathrm{}},tt^1z_it=z_{i+1}(i𝐙)`$ that can be simplified to one-relator group $`t,z_0[z_0,z_0^t]=1`$. Using Tietze transformations, one can transform it into $`L`$ ($`x=z_0t`$, $`y=t^1`$). So the group $`L_{\mathrm{}}`$ is a subgroup of $`L`$. Note that we could consider the group $$L_0=z_i(i=0,1,2,\mathrm{})[z_i,z_{i+1}]=1(i=0,1,2,\mathrm{})$$ instead of $`L_{\mathrm{}}`$. The mapping $`\psi `$, where $`\psi (z_i)=z_{i+1}`$ ($`i=0,1,2,\mathrm{}`$), induces a monomorphism of the group $`L_0`$ into itself. Indeed, the mapping $`\theta `$ such that $`\theta (z_0)=1`$, $`\theta (z_i)=z_{i1}`$ ($`i=1,2,\mathrm{}`$) induces an endomorphism of the group $`L_0`$ and $`\theta (\psi (z))=z`$ for any $`zL_0`$. If we take an HNN-extension of the group $`L_0`$ with stable letter $`t`$ via monomorphism $`\psi `$, then we get the group $`L_0,tt^1z_it=z_{i+1}(i=0,1,2,\mathrm{})`$, that can be transformed into $`L`$ after simplifications. Remark that $`L_0`$ is obviously non-abelian (it maps onto a free group of rank $`2`$ by the homomorphism which maps $`z_1`$ to $`1`$, and $`z_i`$ to $`1`$ for $`i=3,4,\mathrm{}`$). Hence $`L`$ is also non-abelian, that is, $`xyyx`$. For $`a=yx`$, $`b=x`$ we have equality $`[a,a^b]=1`$ in the group $`L`$ and $`[a,b]1`$. If $`L`$ is a diagram group then it satisfies Condition $`2`$ of Theorem 24. Thus it also satisfies Condition $`1`$, that is it contains $`𝐙wr𝐙`$ as a subgroup. As we have mentioned above, in the proof of Corollary 23, the group $`𝐙wr𝐙`$ cannot be a subgroup of a one-relator group because of the result of . The contradiction we have obtained shows that $`L`$ is not a diagram group. It remains to show that $`L`$ can be isomorphically embedded into a diagram group. We apply Theorem 13. It follows from it that the group $`K=𝒪(𝐙,𝐙)`$ is a diagram group. It follows from the description given in Section 2 that $`K`$ has a presentation in terms of generators $`g_i`$, $`h_i`$ ($`i𝐙`$), $`t`$ and defining relations $`[g_i,g_j]=[h_i,h_j]=1`$, $`g_i^t=g_{i+1}`$, $`h_j^t=h_{j+1}`$, where $`i,j𝐙`$ and $`[g_i,h_j]=1`$ for $`ij`$, $`i,j𝐙`$. So it suffices to prove that the group $`L=x,y[xy,yx]=1`$ is a subgroup in the diagram group $`K=𝒪(𝐙,𝐙)`$. Consider the group $$K_0=g_i,h_i(i0)[g_i,g_j]=[h_i,h_j]=1(i,j0),[g_i,h_j]=1(ji0).$$ The map $`g_ig_{i+1}`$, $`h_jh_{j+1}`$ ($`i,j0)`$ can be extended to an endomorphism $`\psi `$ of the group $`K_0`$. It is a monomorphism because the map $`g_0,h_01`$, $`g_ig_{i1}`$, $`h_jh_{j1}`$ ($`i,j1)`$ can be also extended to an endomorphism $`\theta `$ and $`\theta (\psi (z))=z`$ for any $`zK_0`$. Therefore, one can consider an HNN-extension of the group $`K_0`$ with a stable letter $`t`$ via monomorphism $`\psi `$. We obtain the group $`K_0,t\psi (z)=z^t(zK_0)`$ that has almost the same presentation as $`K`$ with the only difference that the subscripts of the presentation of $`K`$ run over all $`𝐙`$. Adding new generators $`g_i=g_0^{t^i}`$, $`h_j=h_0^{t^j}`$ for negative $`i`$, $`j`$, we easily transform the presentation obtained above to the presentation of $`K`$. So it suffices to prove the following Lemma. ###### Lemma 29 The subgroup in $`K`$ generated by elements $`g_ih_{i+1}`$ $`(i0)`$ and $`t`$ is isomorphic to $`L`$. Proof. The group $`K`$ is an HNN-extension of the group $`K_0`$, that is, we add the letter $`t`$ and relations $`g_i^t=g_{i+1}`$, $`h_j^t=h_{j+1}`$ ($`i,j0`$) to its presentation. Define the sets $`_k`$ ($`k0`$) of defining relations over the alphabet $`\{z_0,h_0,z_1,h_1,\mathrm{}\}`$. For $`_0`$ we take the set of relations of the group $`L_0`$, that is, $`_0=\{[z_i,z_{i+1}]=1(i0)\}`$. Further, for $`k1`$ we put $$_k=\{z_i^{h_k}=z_i(0i<k)\}\{z_i^{h_k}=z_i^{z_{k1}}(ik)\}\{h_j^{h_k}=h_j(1j<k)\}.$$ It is clear that $`L_0=z_i(i0)_0`$. Let $$L_k=z_i(i0),h_j(1jk)_0_1\mathrm{}_k$$ for $`k1`$. We shall prove that for any $`k0`$, the group $`L_{k+1}`$ can be obtained from $`L_k`$ by a suitable HNN-extension. Consider the map $`\psi _k`$ given by the following rules: $`\psi _k(z_i)=z_i`$ for $`0ik`$, $`\psi _k(z_i)=z_i^{z_k}`$ for $`i>k`$, $`\psi _k(h_j)=h_j`$ for $`1jk`$. Let us extend it to a homomorphism of the corresponding free group into the group $`L_k`$. Let us check that all relations of the group $`L_k`$ will be equalities in $`L_k`$ under $`\psi _k`$. First of all we shall check that $`\psi _k([z_i,z_{i+1}])=1`$ for all $`i0`$. If $`0i<k`$, then $`\psi _k([z_i,z_{i+1}])=[z_i,z_{i+1}]=1`$ in $`L_k`$. The equality $`\psi _k(z_i)=z_i^{z_k}`$ holds also for $`i=k`$. Hence for all $`ik`$ we also have $`\psi _k([z_i,z_{i+1}])=[z_i^{z_k},z_{i+1}^{z_k}]=[z_i,z_{i+1}]^{z_k}=1`$. Now consider the other relations of $`L_k`$. They have one of the following three forms: $`z_i^{h_j}=z_i`$ for $`0i<jk`$; $`z_i^{h_j}=z_i^{z_{j1}}`$ for $`ij`$, $`1jk`$; $`[h_i,h_j]=1`$ for $`1i<jk`$. Considering three cases, we map each of the relations by $`\psi _k`$. If $`0i<jk`$ then we have $`\psi _k(z_i^{h_j})=\psi _k(z_i)^{\psi _k(h_j)}=z_i^{h_j}=z_i=\psi _k(z_i)`$. In the second case we will consider two subcases: $`ik`$ and $`i>k`$. In the first subcase, that is, $`1jik`$, we get $`\psi _k(z_i^{h_j})=\psi _k(z_i)^{\psi _k(h_j)}=z_i^{h_j}=z_i^{z_{j1}}=\psi _k(z_i)^{\psi _k(z_{j1})}=\psi _k(z_i^{z_{j1}}`$. In the second subcase, that is, $`1jk<i`$, we have: $`\psi _k(z_i^{h_j})=\psi _k(z_i)^{\psi _k(h_j)}=(z_i^{z_k})^{h_j}=(z_i^{h_j})^{z_k^{h_j}}=(z_i^{z_{j1}})^{z_k^{z_{j1}}}=(z_i^{z_k})^{z_{j1}}=\psi _k(z_i)^{\psi _k(z_{j1})}=\psi _k(z_i^{z_{j1}})`$ (equality $`z_k^{h_j}=z_k^{z_{j1}}`$ in the group $`L_k`$ we used in these calculations, is a partial case of the relation of the second form for $`i=k`$). In the third case everything is easy: $`\psi _k([h_i,h_j])=[\psi _k(h_i),\psi _k(h_j)]=[h_i,h_j]=1`$ for $`1i<jk`$. So $`\psi _k`$ induces an endomorphism of the group $`L_k`$. Let us also introduce the map $`\theta _k`$ given by the rules $`\theta _k(z_i)=z_i`$ for $`0ik`$, $`\theta _k(z_i)=z_i^{z_k^1}`$ for $`i>k`$, $`\psi _k(h_j)=h_j`$ for $`1jk`$. One can analogously check that $`\theta _k`$ induces an endomorphism of the group $`L_k`$. It is obvious that $`\theta _k(\psi _k(z))=\psi _k(\theta _k(z))`$ for any $`zL_k`$. This means that $`\psi _k`$ and $`\theta _k`$ are mutually inverse automorphisms of the group $`L_k`$. Consider an HNN-extension of the group $`L_k`$ with stable letter $`h_{k+1}`$ via automorphism $`\psi _k`$ of the group $`L_k`$. Its presentation is obtained from the one of $`L_k`$ by adding $`h_{k+1}`$ to the set of generators and adding relations of the form $`\psi _k(z)=z^{h_{k+1}}`$ to the set of defining relations, where $`z`$ runs over all generators of $`L_k`$. These new relations form exactly the set $`_{k+1}`$. Therefore, this HNN-extension is the group $`L_{k+1}`$. In addition, we also have a natural embedding of $`L_k`$ into $`L_{k+1}`$ for $`k0`$. We have a sequence of embedded subgroups $$L_0<L_1<\mathrm{}<L_k<L_{k+1}<\mathrm{},$$ that give the group $$\widehat{L}=z_i(i0),h_j(j1)_0_1\mathrm{}_k\mathrm{}$$ as a union of them. Let $`H`$ be a subgroup generated by $`z_0`$ and all $`h_j`$ ($`j1`$). Adding $`h_0`$ as a stable letter, we construct an HNN-extension of the group $`\widehat{L}`$ via identical endomorphism of $`H`$ onto itself. That is, we add a new generator $`h_0`$ and relations $`[z_0,h_0]=1`$, $`[h_j,h_0]=1`$ for all $`j1`$. Let us describe explicitly the group $`\overline{L}`$ that we get as a result. It has generators $`z_i`$, $`h_i`$ ($`i0`$) subject to the following defining relations: $$[z_i,z_{i+1}]=1(i0),$$ $$[h_i,h_j]=1(i,j0),$$ $$[z_i,h_j]=1(0i<j),$$ $$[z_i,z_{j1}h_j^1]=1(1ji),$$ $$[z_0,h_0]=1.$$ (we took the relations $`_k`$ for all $`k0`$ together with the relations added at the last step). Let us introduce new generators $`g_i=z_ih_{i+1}^1`$ ($`i0`$). The elements $`g_i`$, $`h_i`$ generate $`\overline{L}`$ so our aim is to describe relations of the group $`\overline{L}`$ in terms of these generators. Replacing elements $`z_i`$ by $`g_ih_{i+1}`$ in the defining relations of the group $`\overline{L}`$, we get: $$[g_ih_{i+1},g_{i+1}h_{i+2}]=1(i0),$$ (19) $$[h_i,h_j]=1(i,j0),$$ (20) $$[g_ih_{i+1},h_j]=1(0i<j),$$ (21) $$[g_ih_{i+1},g_{j1}]=1(1ji),$$ (22) $$[g_0h_1,h_0]=1.$$ (23) Since elements of the form $`h_i`$ ($`i0`$) pairwise commute, we can simplify (21), getting $`[g_i,h_j]=1`$ for $`0i<j`$. Then in (22) the elements $`g_{j1}`$ $`h_{i+1}`$ commute for $`1ji`$ and so (22) reduces to $`[g_i,g_{j1}]=1`$ for all $`1ji`$. This means that all elements $`g_i`$ ($`i0`$) pairwise commute. Let us simplify (19). Note that $`h_{i+2}`$ commutes with the other three elements so it can be excluded. The equality $`[g_ih_{i+1},g_{i+1}]=1`$ we get is equivalent to $`[h_{i+1},g_{i+1}]=1`$ since $`g_i`$ commutes with the other elements. Thus, simplifying (23), we obtain equalities $`[g_i,h_i]=1`$ for all $`i0`$. Let us summarize the above. The group $`\overline{L}`$ has generators $`g_i`$, $`h_i`$ ($`i0`$), where elements $`g_i`$ pairwise commute. Elements of the form $`h_i`$ also pairwise commute and $`g_i`$ commutes with $`h_j`$ whenever $`ij`$. This means that the group $`\overline{L}`$ coincides with $`K_0`$. The group $`L_0`$, naturally embedded into $`\overline{L}`$, is generated by elements $`z_i`$ ($`i0`$). So the subgroup of $`\overline{L}`$ generated by $`g_ih_{i+1}=z_i`$ ($`i0`$) is isomorphic to $`L_0`$. The group $`K`$ is an HNN-extension of the group $`K_0`$ via the monomorphism $`\psi :K_0K_0`$ with $`t`$ as a stable letter. Let us have a subgroup $`M_0`$ of $`K_0`$ such that $`\psi (M_0)M_0`$. In this case, it is easy to show that the subgroup generated by $`M_0`$ and $`t`$ will be the HNN-extension of $`M`$ via the restriction of $`\psi `$ on $`M_0`$. Indeed, let us take such an HNN-extension. It has the form $`M=M_0,tz^t=\psi (z)(zM_0)`$. The map $`tt`$, $`zz`$ for $`zM_0`$ induces a homomorphism $`\varphi `$ from $`M`$ to $`K`$. Since $`zt=t\psi (z)`$ we can represented any element of the group $`M`$ in the form $`t^\alpha zt^\beta `$, where $`zM_0`$, $`\alpha ,\beta 𝐙`$, $`\beta 0`$. Therefore, any element $`m`$ in $`M`$ is conjugated to an element of the form $`t^\gamma z`$ for some $`\gamma 𝐙`$, $`zM_0`$. If $`m1`$, then either $`\gamma 0`$ or $`z1`$. The element $`t^\gamma z`$ maps to an element in $`K`$ of the same form under $`\psi `$. It follows from the elementary properties of HNN-extensions that it is not equal to $`1`$ in $`K`$. Therefore, $`\varphi `$ is an embedding of $`M`$ into $`K`$ and its image is exactly the subgroup of $`K`$ generated by $`M_0`$ and $`t`$. Note that the subgroup of $`K_0`$ generated by elements $`z_i=g_ih_{i+1}`$ ($`i0`$) is invariant under $`\psi `$ because $`\psi (z_i)=z_{i+1}`$ for all $`i0`$. So one can regard this subgroup (isomorphic to $`L_0`$) as $`M_0`$ and apply the arguments from the above paragraph. The corresponding HNN-extension of it is the subgroup of $`K`$ generated by $`t`$ and $`g_ih_{i+1}`$ ($`i0`$). On the other hand, as we have mentioned in the beginning, this HNN-extension is exactly $`L`$. The Lemma and Theorem 28 are proved. There are not many known counterexamples to the Subgroup Conjecture. So it is natural to try to prove this conjecture under some restrictions on the subgroup. With respect to Theorem 26, we would like to ask a few questions. ###### Problem 3 Is it true that any subgroup of R. Thompson’s group $`F`$ is a diagram group? ###### Problem 4 Is it true that the commutator subgroup of any diagram group is a diagram group? It is easy to see that the commutator subgroup of the group $`F`$ satisfies the following condition. Let $`\mathrm{\Delta }`$ be a diagram representing an element in $`F^{}`$, and suppose that a conjugate diagram $`\mathrm{\Psi }^1\mathrm{\Delta }\mathrm{\Psi }`$ is a sum $`\mathrm{\Gamma }_1+\mathrm{\Gamma }_2`$ of two nontrivial spherical diagrams with bases $`v_1`$, $`v_2`$, respectively. Then the diagrams $`\mathrm{\Delta }_1=\mathrm{\Psi }(\mathrm{\Gamma }_1+\epsilon (v_2))\mathrm{\Psi }^1`$ and $`\mathrm{\Delta }_2=\mathrm{\Psi }(\epsilon (v_1)+\mathrm{\Gamma }_2)\mathrm{\Psi }^1`$ also belong to $`F^{}`$. Consider any subgroup $`H`$ of a diagram group $`G`$ that satisfies this condition. We shall say that $`H`$ is closed in $`G`$. It is easy to see that a subgroup $`H`$ is closed if whenever $`\mathrm{\Delta }H`$ and $`\mathrm{\Delta }=\mathrm{\Delta }_1\mathrm{\Delta }_2`$, where $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_2`$ commute and do not belong to the same cyclic subgroup, we have $`\mathrm{\Delta }_1,\mathrm{\Delta }_2H`$. ###### Problem 5 Let $`H`$ be a closed subgroup in a diagram group $`G`$. Is it true that $`H`$ is a diagram group? If the answer to the next problem is positive, then this would imply that all word hyperbolic diagram groups are free. ###### Problem 6 Let $`H`$ be a subgroup in a diagram group $`G`$. Suppose that for any $`hH`$, $`h1`$, the centralizer $`C_G(h)`$ of $`h`$ in $`G`$ is cyclic. Does it imply that $`H`$ is free $`(`$at least for the particular case $`H=G)`$? At the end of this Section let us consider an interesting family of groups. Let $$G_n=x_1,\mathrm{},x_n[x_1,x_2]=[x_2,x_3]=\mathrm{}=[x_{n1},x_n]=[x_n,x_1]=1.$$ It is easy to see that $`G_1=𝐙`$, $`G_2=𝐙\times 𝐙`$, $`G_3=𝐙\times 𝐙\times 𝐙`$, $`G_4=_2\times _2`$, where $`_2`$ is the free group of rank $`2`$. All these groups can be obtained from $`𝐙`$ using finite direct and free products. So these are diagram groups. However, the group $`G_5`$ is not a diagram group. ###### Theorem 30 The groups $`G_n`$ are not diagram groups for odd $`n5`$. Proof. Let $`𝒫=\mathrm{\Sigma }`$ be a semigroup presentation and let $`G=𝒟(𝒫,w)`$ be a diagram group. Consider any element $`gG`$ presented by a diagram $`\mathrm{\Delta }`$. Let us decompose $`\mathrm{\Delta }`$ into the sum of spherical components. It was proved in that the number of these components is an invariant of a diagram with respect to conjugation: see the remark after the proof of Lemma 15.15. One can see from the same Lemma that the number of nontrivial components is also an invariant. Thus one can introduce a function $`𝐜𝐨𝐦𝐩`$, denoted by $`𝐜𝐨𝐦𝐩(g)`$, the number of nontrivial components of a diagram that represents an element $`gG`$. Let us introduce a partial binary relation $``$ on $`G`$. Let $`g_1,g_2G`$ be such that $`𝐜𝐨𝐦𝐩(g_1)=𝐜𝐨𝐦𝐩(g_2)=1`$ (in particular, $`g_1`$, $`g_2`$ are nontrivial). We put $`g_1g_2`$ whenever diagrams $`\mathrm{\Delta }_1`$, $`\mathrm{\Delta }_2`$ that represent elements $`g_1`$, $`g_2`$ respectively, satisfy the following condition: there are words $`x,y,z\mathrm{\Sigma }^{}`$, $`u,v\mathrm{\Sigma }^+`$, some $`(w,xuyvz)`$-diagaram $`\mathrm{\Gamma }`$, simple absolutely reduced spherical diagrams $`\mathrm{\Psi }_1`$ and $`\mathrm{\Psi }_2`$ with bases $`u`$, $`v`$ respectively such that $`\mathrm{\Gamma }^1\mathrm{\Delta }_1\mathrm{\Gamma }=\epsilon (x)+\mathrm{\Psi }_1+\epsilon (yvz)`$, $`\mathrm{\Gamma }^1\mathrm{\Delta }_2\mathrm{\Gamma }=\epsilon (xuy)+\mathrm{\Psi }_2+\epsilon (z)`$. It follows from this definition that if $`g_1g_2`$, then elements $`g_1`$, $`g_2`$ commute and generate a subgroup isomorphic to $`𝐙\times 𝐙`$ in $`G`$. In particular, they do not belong to the same cyclic subgroup. Thus the relation $``$ is antireflexive, that is, condition $`gg`$ never holds for $`gG`$. Let us establish a few properties of $``$. ###### Lemma 31 Let $`𝒫=\mathrm{\Sigma }`$ be a semigroup presentation and let $`G=𝒟(𝒫,w)`$ be a diagram group. The relation $``$ is transitive, that is, for any $`g_1,g_2,g_3G`$ such that $`𝐜𝐨𝐦𝐩(g_1)=𝐜𝐨𝐦𝐩(g_2)=𝐜𝐨𝐦𝐩(g_3)=1`$, conditions $`g_1g_2`$ and $`g_2g_3`$ imply $`g_1g_3`$. Proof. Let element $`g_i`$ in $`G`$ be represented by a diagram $`\mathrm{\Delta }_i`$ ($`i=1,2,3`$). Since $`g_1g_2`$, there are words $`x,y,z\mathrm{\Sigma }^{}`$, $`u,v\mathrm{\Sigma }^+`$, $`(w,xuyvz)`$-diagram $`\mathrm{\Gamma }`$, simple absolutely reduced spherical diagrams $`\mathrm{\Psi }_1`$ and $`\mathrm{\Psi }_2`$ with bases $`u`$, $`v`$ respectively, such that $`\mathrm{\Gamma }^1\mathrm{\Delta }_1\mathrm{\Gamma }=\epsilon (x)+\mathrm{\Psi }_1+\epsilon (yvz)`$, $`\mathrm{\Gamma }^1\mathrm{\Delta }_2\mathrm{\Gamma }=\epsilon (xuy)+\mathrm{\Psi }_2+\epsilon (z)`$. Since $`g_2g_3`$, there are words $`x^{},y^{},z^{}\mathrm{\Sigma }^{}`$, $`u^{},v^{}\mathrm{\Sigma }^+`$, $`(w,x^{}u^{}y^{}v^{}z^{})`$-diagram $`\mathrm{\Gamma }^{}`$, simple absolutely reduced spherical diagrams $`\mathrm{\Psi }_2^{}`$ and $`\mathrm{\Psi }_3^{}`$ with bases $`u^{}`$, $`v^{}`$ respectively, such that $`(\mathrm{\Gamma }^{})^1\mathrm{\Delta }_2\mathrm{\Gamma }^{}=\epsilon (x^{})+\mathrm{\Psi }_2^{}+\epsilon (y^{}v^{}z^{})`$, $`(\mathrm{\Gamma }^{})^1\mathrm{\Delta }_3\mathrm{\Gamma }^{}=\epsilon (x^{}u^{}y^{})+\mathrm{\Psi }_3^{}+\epsilon (z^{})`$. It follows from these conditions that $`\epsilon (xuy)+\mathrm{\Psi }_2+\epsilon (z)=\mathrm{\Theta }^1(\epsilon (x^{})+\mathrm{\Psi }_2^{}+\epsilon (y^{}v^{}z^{}))\mathrm{\Theta }`$, where $`\mathrm{\Theta }=(\mathrm{\Gamma }^{})^1\mathrm{\Gamma }`$. Diagrams $`\epsilon (xuy)+\mathrm{\Psi }_2+\epsilon (z)`$ and $`\epsilon (x^{})+\mathrm{\Psi }_2^{}+\epsilon (y^{}v^{}z^{})`$ are conjugate by an element $`\mathrm{\Theta }`$. It follows from \[14, Lemma 15.15\] that the components of these diagrams are conjugate respectively. In particular, words $`x^{}`$ and $`z`$ are nonempty. Applying this Lemma, we conclude that $`\mathrm{\Theta }=\mathrm{\Theta }_1+\mathrm{\Theta }_2+\mathrm{\Theta }_3`$, where $`\mathrm{\Theta }_1`$, $`\mathrm{\Theta }_2`$, $`\mathrm{\Theta }_3`$ are $`(x^{},xuy)`$-, $`(u^{},v)`$\- and $`(y^{}v^{}z^{},z)`$-diagrams, respectively. Let us now take the diagram $`\mathrm{\Xi }=\mathrm{\Gamma }(\epsilon (xuyv)+\mathrm{\Theta }_3^1)`$. It is clear that $`\mathrm{\Xi }=\mathrm{\Gamma }^{}\mathrm{\Theta }(\epsilon (xuyv)+\mathrm{\Theta }_3^1)=\mathrm{\Gamma }^{}(\mathrm{\Theta }_1+\mathrm{\Theta }_2+\mathrm{\Theta }_3)(\epsilon (xuyv)+\mathrm{\Theta }_3^1)=\mathrm{\Gamma }^{}(\mathrm{\Theta }_1+\mathrm{\Theta }_2+\epsilon (y^{}v^{}z^{}))`$. We have $`\mathrm{\Xi }^1\mathrm{\Delta }_1\mathrm{\Xi }=(\epsilon (xuyv)+\mathrm{\Theta }_3^1)^1(\mathrm{\Gamma }^1\mathrm{\Delta }_1\mathrm{\Gamma })(\epsilon (xuyv)+\mathrm{\Theta }_3^1)=(\epsilon (xuyv)+\mathrm{\Theta }_3^1)^1(\epsilon (x)+\mathrm{\Psi }_1+\epsilon (yvz))(\epsilon (xuyv)+\mathrm{\Theta }_3^1)=\epsilon (x)+\mathrm{\Psi }_1+\epsilon (yvy^{}v^{}z^{})`$, and $`\mathrm{\Xi }^1\mathrm{\Delta }_3\mathrm{\Xi }=(\mathrm{\Theta }_1+\mathrm{\Theta }_2+\epsilon (y^{}v^{}z^{}))^1(\mathrm{\Gamma }^{})^1\mathrm{\Delta }_3\mathrm{\Gamma }^{}(\mathrm{\Theta }_1+\mathrm{\Theta }_2+\epsilon (y^{}v^{}z^{}))=(\mathrm{\Theta }_1^1+\mathrm{\Theta }_2^1+\epsilon (y^{}v^{}z^{}))(\epsilon (x^{}u^{}y^{})+\mathrm{\Psi }_3^{}+\epsilon (z^{}))(\mathrm{\Theta }_1+\mathrm{\Theta }_2+\epsilon (y^{}v^{}z^{}))=\epsilon (xuyvy^{})+\mathrm{\Psi }_3^{}+\epsilon (z^{})`$. Thus conjugating diagrams $`\mathrm{\Delta }_1`$ and $`\mathrm{\Delta }_3`$ by a $`(w,xuyvy^{}v^{}z^{})`$-diagram $`\mathrm{\Xi }`$, we represent them in the form enabling us to conclude that $`g_1g_3`$. The proof is complete. Lemma 31 implies that the relation $``$ is also antisymmetric, that is $`g_1g_2`$ excludes $`g_2g_1`$. Let us establish one more property of $``$. ###### Lemma 32 Let $`𝒫=\mathrm{\Sigma }`$ be a semigroup presentation and let $`G=𝒟(𝒫,w)`$ be a diagram group. We claim that for any commuting elements $`g_1,g_2G`$ that do not belong to the same cyclic subgroup and satisfy $`𝐜𝐨𝐦𝐩(g_1)=𝐜𝐨𝐦𝐩(g_2)=1`$, exactly one of the following conditions holds: $`g_1g_2`$ or $`g_2g_1`$. Proof. Let $`A_i`$ be a diagram that represents an element $`g_iG`$ ($`i=1,2`$). Since $`[g_1,g_2]=1`$, we can apply Theorem 17 and find a word $`v=v_1\mathrm{}v_n`$, spherical $`(v_j,v_j)`$-diagrams $`\mathrm{\Delta }_j`$ ($`1jn`$), integers $`d_{ij}`$ ($`1i2`$, $`1jn`$) and some $`(w,v)`$-diagram $`\mathrm{\Gamma }`$ such that $`\mathrm{\Gamma }^1A_i\mathrm{\Gamma }=\mathrm{\Delta }_1^{d_{i1}}+\mathrm{}+\mathrm{\Delta }_n^{d_{in}}`$, where diagrams $`\mathrm{\Delta }_j`$ ($`1jn`$) are either trivial or simple absolutely reduced. The condition $`𝐜𝐨𝐦𝐩(g_1)=1`$ means that there is exactly one number $`j`$ from $`1`$ to $`n`$ such that diagram $`\mathrm{\Delta }_j^{d_{1j}}`$ is not trivial. Analogously, condition $`𝐜𝐨𝐦𝐩(g_2)=1`$ means that there exists exactly one number $`k`$ from $`1`$ to $`n`$ such that diagram $`\mathrm{\Delta }_k^{d_{1k}}`$ is not trivial. If $`j=k`$ then diagrams $`A_1`$, $`A_2`$ belong to the same cyclic subgroup of $`G`$ but this is impossible. If $`j<k`$, then $`g_1g_2`$ by definition. If $`k<j`$ then $`g_2<g_1`$. The proof is complete. Let us continue the proof of Theorem 30. Let $`n=2k+1`$, $`k2`$. Suppose that $`G_n=𝒟(𝒫,w)`$ is a diagram group over $`𝒫=\mathrm{\Sigma }`$ with base $`w`$. First of all let us prove that $`𝐜𝐨𝐦𝐩(x_i)=1`$ for all generators $`x_i`$ of $`G_n`$. Let us establish that the centralizer $`C(x_i)`$ of the element $`x_i`$ ($`1in`$) in $`G_n`$ is the subgroup generated by elements $`x_{i1}`$, $`x_i`$ and $`x_{i+1}`$, isomorphic to the direct product $`_2\times 𝐙`$ (subscripts are taken modulo $`n`$). By symmetry, it suffices to consider the centralizer of $`x_n`$. Suppose that $`C(x_n)gpx_1,x_{n1},x_n`$. Consider a group word $`W`$ of minimal length in $`x_1`$, …, $`x_n`$ such that $`WC(x_n)`$, $`Wgpx_1,x_{n1},x_n`$. In particular, the word $`W`$ is nonempty and it has neither nonempty initial nor nonempty terminal segment that belongs to $`gpx_1,x_{n1},x_n`$. The group $`G_n`$ is an HNN-extension with base $$\mathrm{\Gamma }=x_1,\mathrm{},x_{n1}[x_1,x_2]=[x_2,x_3]=\mathrm{}=[x_{n2},x_{n1}]=1,$$ and stable letter $`x_n`$, with respect to the identical automorphism of the subgroup $`gpx_1,x_{n1}`$ of $`\mathrm{\Gamma }`$. Consider the element $`x_n^1Wx_nW^1`$ of this HNN-extension. It equals $`1`$ in the group $`G_n`$ since $`WC(x_n)`$. By Britton’s Lemma (see ), the word $`x_n^1Wx_nW^1`$ has a subword of the form $`U=x_n^\delta Vx_n^\delta `$, where $`\delta =\pm 1`$, $`Vgpx_1,x_{n1}`$ is a word that does not contain $`x_n^{\pm 1}`$. Since the word $`W`$ is chosen to have minimal length, $`U`$ is not contained in $`W^{\pm 1}`$. Otherwise the occurrence of $`U`$ can be replaced by an occurrence of the word $`V`$ that is equal to $`U`$ in $`G_n`$, decreasing the length of $`W`$. It is clear that $`V`$ is nonempty since $`W`$ cannot begin or end with $`x_n^{\pm 1}`$. Thus $`V`$ is neither initial nor terminal segment of $`W^{\pm 1}`$. So it is clear that $`U`$ does not occur in $`x_n^1Wx_nW^1`$. We got a contradiction. Thus $`C(x_n)=gpx_1,x_{n1},x_n`$. Consider a mapping of the alphabet $`\{x_1,\mathrm{},x_n\}`$ into $`G_n`$, sending each of the elements $`x_1`$, $`x_{n1}`$, $`x_n`$ to itself and senging all the other elements to $`1`$. Extending this mapping to a homomorphism of the corresponding free group into $`G_n`$, we see that all defining relations of the group $`G_n`$ are sent to $`1`$. Thus we have an induced homomorphism $`\varphi :G_nG_n`$. It is obvious that it is a retraction, that is, $`\varphi ^2=\varphi `$. On the one hand, the subgroup $`\varphi (G_n)`$ of $`G_n`$ equals $`gpx_1,x_{n1},x_n`$; on the other hand, this group is presented by relations of the group $`G_n`$ with additional conditions $`x_2=\mathrm{}=x_{n2}=1`$. Thus for any $`n>3`$ we have $$gpx_1,x_{n1},x_n=\varphi (G_n)=x_1,x_{n1},x_n[x_{n1},x_n]=[x_n,x_1]=1_2\times 𝐙,$$ as desired. So $`C(x_i)_2\times 𝐙`$ for all $`i`$ from $`1`$ to $`n`$; in particular, the centre of $`C(x_i)`$ is cyclic. If $`x_i`$ were represented by a diagram with more than one nontrivial component, then its centralizer would have at least two direct summands isomorphic to $`𝐙`$ by \[14, Theorem 15.35\]. So its centre would not be cyclic. Taking into account that $`x_i`$ is nontrivial, we conclude that $`𝐜𝐨𝐦𝐩(x_i)=1`$. (Notice that we have not used yet that $`n`$ is odd.) Apply Lemma 32 and suppose without loss of generality that $`x_1x_2`$. Suppose that $`x_2x_3`$. Then Lemma 31 would imply that $`x_1x_3`$, so elements $`x_1`$ and $`x_3`$ commute. It is clear that these elements do not commute in $`G_n`$. The contradiction we obtain allows to apply Lemma 32 again and to conclude that $`x_3x_2`$. We will obtain a contradiction again if we suppose that $`x_4x_3`$. So in fact $`x_3x_4`$. Continuing in this way, we shall conclude that $`x_{2k+1}x_{2k}`$, $`x_{2k+1}x_1`$, $`x_2x_1`$. We have a contradiction. The Theorem is proved. It is reasonable to pose a question with respect to Theorem 30. ###### Problem 7 For which $`n`$ the groups $`G_n`$ are diagram groups? For which $`n`$ they are isomorphically representable by diagrams? If there is an odd $`n5`$ such that the group $`G_n`$ is representable by diagrams, then we have one more counterexample to the Subgroup Conjecture. Otherwise we would have a generalization of Theorem 30. ## 6 Distortion of Subgroups in Diagram Groups The problems that concern distortion in groups form a branch of geometric group theory under development (see ). Let us recall some definitions. Let $`A`$ be a group with finite set of generators $`X`$. In this case, for any $`gG`$ there exists an $`n0`$ and $`x_1,\mathrm{},x_nX^{\pm 1}`$ such that $`g=x_1\mathrm{}x_n`$. The least $`n`$ with this property is called the length of the element $`g`$ with respect to the generating set $`X`$ and it is denoted by $`\|g\|_X`$. If there are two functions $`\varphi `$, $`\psi `$ from $`G`$ to the set of all nonnegative integers, then we shall write $`\varphi \psi `$, whenever there is a positive integer constant $`C`$ such that $`\varphi (g)C\psi (g)`$ for all $`gG`$. If it holds $`\varphi \psi `$ and $`\psi \varphi `$ for the two functions simultaneously, then we call these functions equivalent and denote this fact by $`\varphi \psi `$. Obviously, $``$ is in fact the equivalence relation (one does not have to mix it with another equivalence relation that is often used when the Dehn functions are discussed). So, for the two functions one has $`\varphi \psi `$ if and only if there exisets a positive integer constant $`C`$ such that $$\frac{\varphi (g)}{C}\psi (g)C\varphi (g)\text{ for all }gG.$$ If $`X`$ and $`Y`$ are finite sets of generators of the same group $`A`$, then elementary arguments show that functions $`\|\|_X`$ and $`\|\|_Y`$ are equivalent. Let we have two finitely generated groups $`A`$ and $`B`$ such that $`A`$ is a subgroup of $`B`$. Let us fix some finite system of generators $`X`$ for the group $`A`$ and some finite system of generators $`Y`$ for the group $`B`$. For any element $`gA`$ we define two numbers: $`\|g\|_X`$ and $`\|g\|_Y`$. Functions $`\|\|_X`$, $`\|\|_Y`$ can be regarded as functions on $`A`$. Later we will compare functions on two groups one embedded into another, with respect to $``$, taking the corresponding restrictions of these functions. It follows from elementary reasons that $`\|\|_Y\|\|_X`$. If the converse is true, that is, $`\|\|_X\|\|_Y`$ holds, the we say that a subgroup $`A`$ embeds into $`B`$ quasiisometrically or without distortion (this happens, if $`\|\|_X\|\|_Y`$). Note that the equivalence of two length functions $`\|\|_X`$ and $`\|\|_Y`$ does not depend on the choice of finite systems of generators $`X`$ and $`Y`$. If to consider functions up to equivalence, then one can introduce length functions $`\mathrm{}_A`$ and $`\mathrm{}_B`$ in finitely generated groups $`A`$ and $`B`$, respectively, that will depend of $`A`$ and $`B`$ only. The quasiisometricity of an embedding of $`A`$ into $`B`$ means that $`\mathrm{}_A\mathrm{}_B`$. Now consider a more general situation of an embedding of $`A`$ into $`B`$ for two finitely generated groups, $`AB`$. Let we distinguish some finite generating sets $`X`$ and $`Y`$ in groups $`A`$ and $`B`$ respectively. One can consider the function $$disto(n)=\underset{\|g\|_Yn}{\mathrm{max}}\|g\|_X,$$ that describes distortion of the subgroup $`A`$ embedded into $`B`$. It is called the distortion function of the subgroup $`A`$ in $`B`$. It is easy to find out that if we change the generating sets, then the distortion function $`disto(n)`$ is not essentially changed. The reader can easily write down the corresponding inequalities. Thus we can talk about linear, quadratic, polynomial, exponential etc distortion. The question about distortion is aslo interesting with respect to the so called membership problem. Let we have two finitely generated groups $`A`$ and $`B`$, where $`AB`$. The membership problem of elements of the group $`B`$ into the subgroup $`A`$ is the question on the existence of an algorithm that decides, given a word on the generators of $`B`$, whether the element of $`B`$ presented by this word belongs to $`A`$. The membership problem of elements of $`B`$ into a subgroup $`A`$ is decidable if and only if the distortion function $`disto(n)`$ defined above is recursive (equivalently, has a recursive upper bound). It is interesting to find the conditions under which all finitely generated subgroups of a given group will embed into it (or any its finitely generated subgroup) without distortion. Free groups and abelian groups have this property. Even for the case of nilpotent groups the situation is quite different: in any nilpotent (non-abelian) torsion-free group there are cyclic subgroups that have distortion in them. (Note that diagram groups may have distorted subgroups in general: due to the classical result of Mikhailova , the group $`_2\times _2`$ has finitely generated subgroups with undecidable membership problem, so they are distorted.) Let us mention two recent results of Burillo : he proved that every cyclic subgroup of R. Thompson’s group $`F`$ is embedded into it without distortion. Also he gave examples of quasi-isometric embeddings of groups $`F\times 𝐙^n`$ ($`n1`$) and $`F\times F`$ into $`F`$. It is thus natural to ask whether every finitely generated subgroup embeds quasi-isometrically into $`F`$. We give a negative answer. Namely, for any integer $`d2`$ we construct a finitely generated subgroup of $`F`$ with distortion at least $`n^d`$. The fact that any cyclic subgroup of every finitely generated group representable by diagrams (including the case of $`F`$) embeds quasi-isometrically into it, follows easily from \[14, Lemma 15.29\]. We shall prove a more general result. ###### Theorem 33 Let $`B`$ be a finitely generated subgroup of a diagram group $`G`$ and let $`A`$ be a finitely generated abelian subgroup of $`B`$. Then $`A`$ embeds into $`B`$ quasi-isometrically. Proof. Note that diagram groups are torsion-free \[14, Theorem 15.11\] and so $`A`$ is isomorphic to $`𝐙^m`$ for some integer $`m`$. (This also follows from Theorem 15.) Let $`G=𝒟(𝒫,w)`$, where $`𝒫=\mathrm{\Sigma }`$ is a semigroup presentation. By $`A_1`$, …, $`A_m`$ we denote spherical diagrams with base $`w`$ presenting free generators of $`A𝐙^m`$. Since these elements are pairwise commutative in $`G`$, we can apply Theorem 17 to them. Thus we have a word $`v=v_1\mathrm{}v_n`$, sperical $`(v_j,v_j)`$-diagrams $`\mathrm{\Delta }_j`$ $`(1jn)`$, integers $`d_{ij}`$ $`(1im`$, $`1jn)`$ and some $`(w,v)`$-diagram $`\mathrm{\Gamma }`$ such that $$\mathrm{\Gamma }^1A_i\mathrm{\Gamma }=\mathrm{\Delta }_1^{d_{i1}}+\mathrm{}+\mathrm{\Delta }_n^{d_{in}}$$ for all $`1im`$. Each of the diagrams $`\mathrm{\Delta }_1`$, …, $`\mathrm{\Delta }_n`$ is either trivial or simple absolutely reduced. Conjugation by diagram $`\mathrm{\Gamma }`$ is an isomorphism of groups $`G=𝒟(𝒫,w)`$ and $`𝒟(𝒫,v)`$. Since the property of a subgroup to be embeddable quasi-isometrically is an invariant under isomorphism, we can assume without loss of generality that $`B`$ is a subgroup of $`𝒟(𝒫,v)`$. It suffices to prove that $`\mathrm{}_A\mathrm{}_B`$. Consider the diagrams $$\mathrm{\Delta }_j^{}=\epsilon (v_1\mathrm{}v_{j1})+\mathrm{\Delta }_j+\epsilon (v_{j+1}\mathrm{}v_n)$$ for all $`j`$ from $`1`$ to $`n`$. By $`\mathrm{\Psi }_1`$, …, $`\mathrm{\Psi }_r`$ we denote those of diagrams $`\mathrm{\Delta }_1^{}`$, …, $`\mathrm{\Delta }_n^{}`$ that are nontrivial. The form a basis $`X`$ of a free abelian group $`C`$. Since $`A`$ embeds into $`C`$ quasi-isometrically, we have $`\mathrm{}_A\mathrm{}_C`$. For any element in the group $`𝒟(𝒫,v)`$ presented by a reduced diagram $`\mathrm{\Delta }`$, we denote by $`\mathrm{\#}(\mathrm{\Delta })`$ the number of cells in $`\mathrm{\Delta }`$. Thus $`\mathrm{\#}`$ is a function on the diagram group. Let $`Y`$ be a finite generating set of the group $`B`$ and let $`K`$ be the greatest number of cells for diagrams in $`Y`$. Then it is obvious that $`\mathrm{\#}(\mathrm{\Delta })K\|\mathrm{\Delta }\|_Y`$ for any diagram $`\mathrm{\Delta }`$ in $`B`$. So we have $`\mathrm{\#}\mathrm{}_B`$. Let $`s_1`$, …, $`s_r`$ be arbitrary integers. Consider the element $`\mathrm{\Delta }=\mathrm{\Psi }_1^{s_1}\mathrm{}\mathrm{\Psi }_r^{s_r}`$ in $`C`$. All diagrams $`\mathrm{\Delta }_j`$ ($`1jn`$) are absolutely reduced. So it follows easily from the definition of diagrams $`\mathrm{\Psi }_k`$ ($`1kr`$) that $`\mathrm{\#}(\mathrm{\Delta })=\|s_1\|\mathrm{\#}(\mathrm{\Psi }_1)+\mathrm{}+\|s_r\|\mathrm{\#}(\mathrm{\Psi }_r)`$. Since $`\|\mathrm{\Delta }\|_X=\|s_1\|+\mathrm{}+\|s_r\|`$, we can deduce inequality $`\|\mathrm{\Delta }\|_X\mathrm{\#}(\mathrm{\Delta })K^{}\|\mathrm{\Delta }\|_X`$, where $`K^{}`$ is the greatest number of cells for the diagrams in $`X`$. Therefore, $`\mathrm{}_C\mathrm{\#}`$. Summarizing what we have said above, we conclude that $`\mathrm{}_A\mathrm{}_C\mathrm{\#}\mathrm{}_B`$. Now obvious inequality $`\mathrm{}_B\mathrm{}_A`$ gives us the equivalence $`\mathrm{}_A\mathrm{}_B`$. The Theorem is proved. Now consider R. Thompson’s group $`F`$, take any its element $`gF`$ and its centralizer $`C_F(g)`$ in $`F`$. In \[14, Corollary 15.36\] we gave the description of centralizers in $`F`$: they are finite direct products of groups that isomorphic to either $`F`$ or $`𝐙`$. In particular, all of them are finitely generated. Remark that if an element $`gF`$ is presented by a diagram $`\mathrm{\Delta }`$, then to find its centralizer, one needs to find an absolutely reduced diagram $`\mathrm{\Delta }^{}`$ conjugated to $`\mathrm{\Delta }`$ (this can be done effectively by \[14, Lemma 15.14\]) and then decompose $`\mathrm{\Delta }^{}`$ into a sum of (spherical) components. To each trivail component we assign the group $`F`$ and to each nontrivial one we assign $`𝐙`$. Then we take direct product of these groups. It is easy to see that the groups we get in this way are exactly groups of the form $`F^m\times 𝐙^n`$, where $`0mn+1`$. The Theorem below generalizes Burillo’s results from , where it is shown that $`F`$ has quasi-isometrically embedded subgroups isomorphic to $`F\times F`$ (Proposition 9), and for every $`n1`$ there are quasi-isometrically embedded subgroups isomorphic to $`F\times 𝐙^n`$ (Corollary 6). (Although the group $`F\times F`$ cannot be a centralizer of an element in $`F`$, it is embeddable without distortion into $`F^2\times 𝐙`$, which is a centralizer of some element in $`F`$. This implies the first of results quoted above.) ###### Theorem 34 For any element $`g`$ in R. Thompson’s group $`F`$, the centralizer $`C_F(g)`$ of this element embeds into $`F`$ quasi-isometrically. First of all we need a lemma that can be deduced easily as a consequence of \[7, Proposition 2\]. But we give a direct proof of the fact we need. Let $`𝒫=xx^2=x`$. For any $`k1`$, by $`\mathrm{\#}_k(g)`$ we denote the number of cells in in the (reduced) spherical diagram with base $`x^k`$ that presents the element $`gF𝒟(𝒫,x^k)`$. The number $`\|g\|`$ denotes the length of $`gF`$ with respect to the set $`\{x_0,x_1\}`$ of generators. ###### Lemma 35 For any $`gF`$, the following inequalities hold: $$\frac{\|g\|}{3}\mathrm{\#}_3(g)2\|g\|.$$ For any $`k`$ the function $`\mathrm{\#}_k`$ is equivalent to the length function $`\|\|`$. Proof. Diagrams that correspond to the paths $`(x^2,xx^2,1)(1,x^2x,x^2)`$ and $`(x,xx^2,x)(1,x^2x,x^2)`$ have two cells each. They present the elements $`x_0`$, $`x_1`$ of R. Thompson’s group $`F𝒟(𝒫,x^3)`$. If $`g`$ is an element of length $`n`$, then it can be presented by a diagram with base $`x^3`$ that has at most $`2n`$ cells. The second inequality is thus proved. Let us prove the first inequality. Let an element $`g`$ is presented by a diagram $`\mathrm{\Delta }`$ with base $`x^3`$. The longest positive path from $`\iota (\mathrm{\Delta })`$ to $`\tau (\mathrm{\Delta })`$ cuts this diagram into two parts: positive and negative one. The number of cells in each of the parts is the same, let it be equal to $`m`$. Then $`\mathrm{\#}_3(g)=2m`$. It is easy to see that the longest positive path has length $`m+3`$. Represent $`g`$ as a normal form $`g=g_1g_2^1`$, where each of the elements $`g_1`$, $`g_2`$ is positive, that is, it is a product of positive exponents of generators. Since $`\|g\|\|g_1\|+\|g_2\|`$, it suffices to estimate the length of $`g_1`$. (The lenght of $`g_2`$ can be estimated analogously.) So let $`g_1=x_0^jx_{i_1}\mathrm{}x_{i_s}`$, where $`s0`$ and $`1i_1\mathrm{}i_s`$ is the normal form of $`g_1`$. For any $`i1`$, replace $`x_i`$ by $`x_0^{1i}x_1x_0^{i1}`$, which is equal to it in $`F`$. Then we have that $`g_1`$ equals in $`F`$ to the word $$x_0^{ji_1+1}x_1x_0^{i_1i_2}x_1\mathrm{}x_0^{i_{s1}i_s}x_1x_0^{i_s1},$$ that has length $$\|ji_1+1\|+(i_2i_1)+\mathrm{}+(i_si_{s1})+i_s1+s=2i_si_1+\|ji_1+1\|+s1.$$ (If $`s=0`$, then the length is just $`j`$.) according to the procedure described in 1 (Example refTGNF), we have inequality $`s+jm`$. If $`s1`$, then the element $`x_{i_s}`$ corresponds to the edge $`(x^t,xx^2,x^{i_s})`$ in the Squier complex, where $`t1`$, so $`\mathrm{\Delta }`$ has a positive path labelled by $`x^{t+2+i_s}`$. This implies $`t+2+i_sm+3`$ hence $`i_sm`$. Let us consider two cases. ) $`ji_11`$ or $`s=0`$. We have $`\|g_1\|2i_s+j+s2i_13m2`$ for $`s1`$. If $`s=0`$, then $`\|g_1\|=jm`$. ) $`s0`$, $`j<i_11`$. In this case $`\|g_1\|2i_s+sj23m2`$. Summarizing, we conclude that $`\|g_1\|3m`$ for all cases. Also $`\|g_2\|3m`$. Therefore, $`\|g\|6m=3\mathrm{\#}_3(g)`$, what we had to prove. Now it remains to note that $`\|\mathrm{\#}_k(g)\mathrm{\#}_3(g)\|2\|k3\|`$ since the diagram that presents $`g`$ in $`𝒟(𝒫,x^k)`$ can be obtained from $`\mathrm{\Delta }`$ conjugating it by a diagram of $`\|k3\|`$ cells. From what follows that functions $`\mathrm{\#}_k`$ and $`\mathrm{\#}_3`$ are equivalent (one needs to use that $`\mathrm{\#}_k(g)=0`$ if and only if $`g=1`$.) The Lemma is proved. Proof of Theorem 34. Let $`gF`$ be an arbitrary element. Let us present it by an $`(x,x)`$-diagram and reduce this diagram to absolutely reduced form by conjugation. We get some diagram $`\mathrm{\Delta }`$ with base $`x^k`$. By Lemma 35, the length function in $`F`$ is equivalent to $`\mathrm{\#}_k`$. Let $`\mathrm{\Delta }=\mathrm{\Delta }_1+\mathrm{}+\mathrm{\Delta }_m`$ is a decomposition of $`\mathrm{\Delta }`$ into the sum of components, where $`\mathrm{\Delta }_i`$ is a spherical diagram with base $`z_i`$ ($`1im`$). Any element in the centralizer of $`\mathrm{\Delta }`$ is equal to a sum of $`(z_i,z_i)`$-diagrams, and the $`i`$th summand commutes with $`\mathrm{\Delta }_i`$ ($`1im`$). For any $`i`$ from $`1`$ to $`m`$, let $`G_i=F`$, if $`\mathrm{\Delta }_i`$ is trivial and $`G_i=𝐙`$, if $`\mathrm{\Delta }_i`$ is nontrivial. Thus $`C_F(g)G_1\times \mathrm{}\times G_m`$. Choose a system of generators in each of the groups $`G_i`$: if $`G_i=F`$, then the system consists of $`x_0`$, $`x_1`$, and for $`G_i=𝐙`$ the system consists of one element. these systems of generators form a generating set $`Y`$ of the centralizer of $`\mathrm{\Delta }`$. It is clear that any element $`h`$ in the centralizer can be uniquely presented in the form $`h_1\mathrm{}h_m`$, where $`h_iG_i`$ for $`1im`$, and $`\|h\|_Y=\|h_1\|_1+\mathrm{}+\|h_m\|_m`$ (by $`\|\|_i`$ we denote the length in $`G_i`$ with respect to the generating set we have chosen). The number of cells in $`\mathrm{\Delta }`$ is equal to the sum of numbers of cells in diagrams $`\mathrm{\Delta }_i`$ ($`1im`$). So it follows from the equivalence of the length function and the number of cells that the function $`\|\|_Y`$ is equivalent to the length function in $`F𝒟(𝒫,x^k)`$. This means that the embedding of $`C_F(g)`$ into $`F`$ is quasi-isometric. ( ᫨ $`G_i=F`$, then we apply Lemma 35. In the case $`G_i=𝐙`$ the equivalence of the length function and the number of cells is obvious.) The Theorem is proved. Before going to the proof of the next result about distorted subgroups in $`F`$, let us consider the following construction that has its preimage in . Let $`H`$ be a group generated by a finite set $`X`$ and let $`R`$ be a finite subset in $`H`$. By $`N`$ we denote the normal closure of the set $`R`$ in $`H`$. Consider the subgroup $`K`$ in $`H\times H`$ generated by all elements of the form $`(x,x)`$, where $`xX`$, and also all elements of the form $`(r,1)`$, where $`rR`$. it is easy to see that for any $`g,hH`$, the element $`(g,h)`$ is in $`K`$ if and only if the cosets of $`g`$ and $`h`$ by the subgroup $`N`$ are equal. (This is proved in the same way as in ; see also .) It is possible to consider an analog of the Dehn function in this situation. For any element $`gN`$ by $`k(g)`$ we denote the least $`k`$ such that the element $`g`$ is equal in $`H`$ to a product of $`k`$ elements conjugated in $`H`$ to elements in $`R`$ or their inverses. Let $$\mathrm{\Phi }(n)=\underset{\|g\|n}{\mathrm{max}}k(g),$$ where $`\|g\|`$ is the length of $`g`$ with respect to the set $`X`$ of generators. This function can be call a (relative) Dehn function of presentation $`XR`$ with respect to $`H`$; it is clear that if $`H`$ is free, then we have standard Dehn function. Let $`Y`$ denote the above set of generators of $`K`$. Suppose that the element $`(g,1)K`$ can be presented as a product of $`m`$ elements from $`Y^{\pm 1}`$. Then we have an equality $$(g,1)=(u_0,u_0)(r_1,1)^{\epsilon _1}(u_1,u_1)\mathrm{}(r_m,1)^{\epsilon _m}(u_m,u_m),$$ that holds in $`K`$, where $`u_0,u_1,\mathrm{},u_mH`$, $`r_1,\mathrm{},r_mR`$, $`\epsilon _1,\mathrm{},\epsilon _m=\pm 1`$. Therefore, equalities $`g=u_0r_1^{\epsilon _1}u_1\mathrm{}r_m^{\epsilon _m}u_m`$, $`1=u_0u_1\mathrm{}u_m`$ hold in $`H`$. Then $$g=u_0u_1\mathrm{}u_mr_1^{\epsilon _1u_1\mathrm{}u_m}\mathrm{}r_m^{\epsilon _mu_m}=r_1^{\epsilon _1u_1\mathrm{}u_m}\mathrm{}r_m^{\epsilon _mu_m}$$ . So the inequality $`k(g)m`$ holds. In particular, representing $`(g,1)`$ as a product of the least number of generators in $`Y^{\pm 1}`$, we get the inequality $`k(g)\|(g,1)\|_K`$. For each positive integer $`n`$ we have an element $`gH`$ such that $`\|g\|_Xn`$ and $`\mathrm{\Phi }(n)=k(g)`$. The group $`H\times H`$ has the following natural set of generators: $`Z=(X\times \{1\})(\{1\}\times X)`$. It is clear that $`\|(g,1)\|_Z\|g\|_Xn`$, but we have $`\|(g,1)\|_Kk(g)=\mathrm{\Phi }(n)`$. It follows from the definition of the distortion function that $`disto(n)\mathrm{\Phi }(n)`$, where we embed $`K`$ into $`H\times H`$. We proved the following lemma. ###### Lemma 36 Let $`H`$ be a group generated by a finite set $`X`$, let $`R`$ be a finite subset of $`H`$. Consider the subgroup $`K`$ of $`H\times H`$ generated by the set $`Y`$ that consists of all elements $`(x,x)`$ $`(xX)`$ and all elements of the form $`(r,1)`$ $`(rR)`$. Then inequality $`disto(n)\mathrm{\Phi }(n)`$ holds, where $`disto(n)`$ is the distortion function for the embedding of the group $`K`$ generated by $`Y`$ into the group $`H\times H`$ generated by $`Z=(X\times \{1\})(\{1\}\times X)`$. Here $`\mathrm{\Phi }(n)`$ is the relative Dehn function of presentation $`XR`$ with respect to $`H`$. An important property of R. Thompson’s group $`F`$ is that $`F\times F`$ can be embedded into $`F`$. So in order to obtain distorted subgroups in $`F`$ we need to take such a subgroup $`H`$ with finite generating set $`X`$ and a finite subset $`R`$ of $`H`$ such that the Dehn function of $`XR`$ with respect to $`H`$ will be overlinear. Then, in the above notation, we shall get an embedding of $`K`$ into the group $`H\times H`$, which is in turn embeddable into $`F\times F`$ (and so it embeds into $`F`$). The embedding of $`K`$ into $`F`$ will not be quasi-isometric. Note that if to take $`H=F`$, $`R=\{[x_0,x_1]\}`$, then the relative Dehn function will be linear though the standard Dehn function (with respect to the free group on $`\{x_0,x_1\}`$) will be quadratic. Now we will give an example how to construct a subgroup in $`F`$ with at least quadratic distortion. The last Theorem in this Section will be a generalization of this example. ###### Example 37 Let $`H=𝐙wr𝐙`$ be a subgroup of $`F`$ constructed in Section 4. Denote its generators by $`a`$ and $`b`$. Let the elements $`a_n=a^{b^n}`$ ($`n𝐙`$) form a basis of the free abelian subgroup. For $`R`$ we take the set of a single element $`[a,b]=a_0^1a_1`$. Conjugating this element by all elements in $`H`$, we shall get all elements of the form $`c_i=a_i^1a_{i+1}`$ ($`i𝐙`$). It is obvious that all elements of the form $`c_i`$ are also a basis of the free abelian group. Let $`g_n=[a^n,b^n]=a_0^na_n^nH`$. The length of $`g_n`$ with respect to $`\{a,b\}`$ does not exceed $`4n`$. At the same time, we have equality $`g_n=c_0^nc_1^n\mathrm{}c_{n1}^n`$, which shows that $`g_n`$ can be presented as a product of $`n^2`$ elements of the form $`c_i`$ ($`i𝐙`$). Since the elements $`c_i`$ form a basis of a free abelian subgroup, it follows that $`g_n`$ cannot be presented as a product of less than $`n^2`$ elements of the form $`c_i^{\pm 1}`$. Therefore, the Dehn function $`\mathrm{\Phi }(n)`$ of presentation $`a,b[a,b]`$ with respect to $`H`$ satisfies inequality $`\mathrm{\Phi }(4n)n^2`$. Let $`K`$ be a subgroup of $`H\times H`$ generated by $`(a,a)`$, $`(b,b)`$, $`([a,b],1)`$. Lemma 36 shows that the distortion function $`disto`$ that characterizes the embedding of $`K`$ into $`H\times H`$, is at least quadratic. In particular, $`K`$ embeds into $`H\times H`$ with distortion (that is, the embedding is not quasi-isometric). It remains to embed $`H\times H`$ into $`F\times F`$ and then into $`F`$. Taking into account that $`\mathrm{}_F\mathrm{}_{H\times H}`$, we obtain that $`K`$ embeds into $`F`$ with distortion. One can give explicit expressions (in terms of normal forms) of the generators of $`K`$ as a subgroup in $`F`$. The elements $`a=x_1x_2x_1^2`$ and $`b=x_0`$ generate in $`F`$ a subgroup isomorphic to $`𝐙wr𝐙`$. The rules $`(x_0,1)x_1x_2x_1^2`$, $`(x_1,1)x_1^2x_2x_1^3`$, $`(1,x_0)x_2x_3x_2^2`$, $`(1,x_1)x_2^2x_3x_2^3`$ give an embedding of $`F\times F`$ into $`F`$. Using that, it is easy to compute the generators of $`K`$. The following elements of $`F`$ generate the subgroup isomorphic to $`K`$: $$x_1^2x_2^2x_6^2x_7^2x_8^1x_7^1x_6^2x_3^1x_2^1x_1^2,$$ $$x_1x_2x_4x_5x_4^2x_1^2,$$ $$x_1^3x_2^2x_5x_6x_5^2x_3^1x_2^1x_1^3.$$ Now let us prove the result about distorted subgroups of $`F`$ in its general form. ###### Theorem 38 For any $`d2`$, there exists a finitely generated subgroup $`K_d`$ of R. Thompson’s group $`F`$ such that the corresponding distortion function satisfies inequality $`n^ddisto(n)`$. Proof. Define the groups $`H_k`$ ($`k0`$) by induction on $`k`$ in the following way. Let $`H_0=1`$, $`H_{k+1}=H_kwra_{k+1}`$ for $`k0`$, where all groups $`a_k`$ are infinite cyclic. According to Corollary 20, all of them are embeddable into $`F`$. Let us fix an integer $`d2`$ and consider the group $`H_d\times H_d`$, which is also embeddable into $`F`$. For any integers $`k`$, $`n`$ define an element $`g_k(n)`$ as a left-normalized commutator $$g_k(n)=[a_1^n,a_2^n,\mathrm{},a_k^n],$$ defined by induction on $`k`$: $`g_1(n)=a_1^n`$, $`g_{k+1}(n)=[g_k(n),a_{k+1}^n]`$ for $`k1`$. The element $`g_k(1)`$ will be denoted by $`g_k`$. For $`R_d`$ we take the set of a single element $`g_d=[a_1,a_2,\mathrm{},a_d]`$. The elements in $`H_d`$ of the form $$a_i(t_1,\mathrm{},t_r)=a_i^{a_{i+1}^{t_1}\mathrm{}a_{i+r}^{t_r}},$$ where $`1id`$, $`0rdi`$, $`t_1,\mathrm{},t_r𝐙`$, will be called basic. If $`r=0`$ then we just have $`a_i`$. Obviously, $`a_i=a_i(0)=a_i(0,0)=\mathrm{}`$ and so on. In general, we can add zeroes on the right to the sequence $`t_1`$, …, $`t_r`$ in such a way that the total number of arguments in brackets after $`a_i`$ do not exceed $`di`$. Consider two elements $`w_i=a_i(s_{i+1},\mathrm{},s_d)`$ and $`w_j=a_j(t_{j+1},\mathrm{},t_d)`$, where $`1ijd`$, and $`s_k`$ ($`i<kd`$), $`t_k`$ ($`j<kd`$) are integers. (It is easy to see that each pair of basic elements can be presented in this form.) It is clear that if the sequence $`t_{j+1}`$, …, $`t_d`$ is not an end of the sequence $`s_{i+1}`$, …, $`s_d`$, then the elements $`w_i`$ and $`w_j`$ commute. Indeed, in this case one can choose the biggest $`k`$ such that $`s_kt_k`$. Conjugation by the inverse element to $`a_{k+1}^{s_{k+1}}\mathrm{}a_d^{s_d}=a_{k+1}^{t_{k+1}}\mathrm{}a_d^{t_d}`$ takes elements $`w_i`$, $`w_j`$ into the elements $`w_i^{}=a_i(s_{i+1}\mathrm{}s_k)H_{k1}^{a_k^{s_k}}`$, $`w_j^{}=a_j(t_{j+1}\mathrm{}t_k)H_{k1}^{a_k^{t_k}}`$, respectively. But it is clear from the elementary properties of wreath products that the subgroups $`G^{z^s}`$ and $`G^{z^t}`$ of $`Gwrz`$, where $`z`$ generates $`𝐙`$, commute elementwise for any $`st`$. Now, if the sequence $`t_{j+1}`$, …, $`t_d`$ is the end of $`s_{i+1}`$, …, $`s_d`$, that is, $`s_k=t_k`$ for $`j<kd`$, then elements $`w_i`$ and $`w_j`$ coincide in the case $`i=j`$; in the case $`i<j`$ one can write them as $`w_i=a_i(s_{i+1},\mathrm{},s_j)^v`$, $`w_j=a_j^v`$, where $`v=a_{j+1}^{s_{j+1}}\mathrm{}a_d^{s_d}`$. Then for each $`\mathrm{}𝐙`$ one has equalities $`w_i^{w_j^{\mathrm{}}}`$ $`=`$ $`\left(a_i(s_{i+1},\mathrm{},s_j)^v\right)^{a_j^\mathrm{}v}=\left(a_i(s_{i+1},\mathrm{},s_j)^{a_j^{\mathrm{}}}\right)^v=a_i^{a_{i+1}^{s_{i+1}}\mathrm{}a_j^{s_j+\mathrm{}}v}`$ (24) $`=`$ $`a_i(s_{i+1},\mathrm{},s_j+\mathrm{},s_{j+1},\mathrm{},s_d).`$ So we have a rule how to conjugate one basic element by another basic element. Let $`1kd`$. Consider the normal closure $`M_k`$ of the element $`a_1`$ in $`H_k`$. It follows from the above that $`M_k`$ is an abelian group freely generated by the set of elements $$a_1(s_2,\mathrm{},s_k)=a_1^{a_2^{s_2}\mathrm{}a_k^{s_k}},$$ where $`s_2,\mathrm{},s_k𝐙`$. It is possible to define a homomorphism $`\varphi _k:M_k𝐙`$ from $`M_k`$ into the additive group $`𝐙`$ as follows: $`\varphi _k(a_1(s_2,\mathrm{},s_k))=s_2\mathrm{}s_k`$. From this definition we have that for any $`k>1`$, $`hM_{k1}`$ and for any $`\mathrm{}𝐙`$ the following equality holds: $$\varphi _k(h^{a_k^{\mathrm{}}})=\mathrm{}\varphi _{k1}(h).$$ In particular, $`\varphi _k`$ equals zero on $`M_{k1}`$. Our aim is to establish the two facts. 1) For any $`hH_k`$ the equality $`\varphi _k(g_k^h)=1`$ holds. Note that $`g_k`$ obviously belongs to $`M_k`$ so we can apply $`\varphi _k`$ to any element conjugated to $`g_k`$. 2) If $`1kd`$, then the element $`g_k(n)`$ belongs to the normal closure of the element $`g_k`$ and $`\varphi _k(g_k(n))=n^k`$. First we shall deduce the conclusion of our Theorem from these facts. The elements $`a_1`$, …, $`a_d`$ generate the subgroup $`H_d`$. The length of $`g_d(n)`$ with respect to these generators does not exceed $`Dn`$, where $`D=32^{d1}2`$ is a constant that does not depend on $`n`$. From the above two facts it is clear that the element $`g_k(n)`$, being a product of conjugates to $`g_d`$, cannot be presented as a product of less than $`n^d`$ factors that are conjugates to $`g_d`$ or their inverses. In the notation of Lemma 36, this gives inequality $`\mathrm{\Phi }(Dn)n^d`$. Applying this Lemma, we get $`n^ddisto(n)`$. So let us prove the first of the above facts. We proceed by induction on $`k`$. If $`k=1`$, then $`g_1=a_1M_1`$ and $`g_1^h=a_1`$ for any $`hH_1=a_1`$. By definition, $`\varphi _1(a_1)=1`$. Let $`k>1`$, $`hH_k`$. Then $`g_k=[g_{k1},a_k]=g_{k1}^1g_{k1}^{a_k}M_k`$ since $`g_{k1}M_{k1}`$ by the inductive assumption. We have equalities $$\varphi _k(g_k^h)=\varphi _k\left([g_{k1},a_k]^h\right)=\varphi _k(g_{k1}^hg_{k1}^{a_kh})=\varphi _k(g_{k1}^h)+\varphi _k(g_{k1}^{a_kh}).$$ Since $`\varphi _k=0`$ on $`M_{k1}`$, the first summand equals zero. Further, the elements $`g_{k1}^{a_k}H_{k1}^{a_k}`$ and $`hH_{k1}`$ commute, what follows from the definition of a wreath product. Therefore, $`g_{k1}^{a_kh}=g_{k1}^{a_k}`$. It follows from the above properties of $`\varphi _k`$ that for any $`gM_{k1}`$ we have $`\varphi _k(g^{a_k})=\varphi _{k1}(g)`$. So the second summand equals $`\varphi _k(g_{k1}^{a_k})=\varphi _{k1}(g_{k1})=1`$ because $`g_{k1}M_{k1}`$. As a result, $`\varphi _k(g_k^h)=1`$, what we had to prove. Let us prove the second fact. By $`N_k`$ we denote the normal closure of $`g_k`$ in $`H_k`$. Let us prove by induction on $`k`$ that $`g_k(n)N_k`$. This is obvious for $`k=1`$ since $`g_1(n)=a_1^n=g_1^n`$. Let $`k>1`$, and let the fact is true for all values of the parameter less than $`k`$. Since $`g_k=[g_{k1},a_k]`$, we have equality $`g_{k1}^{a_k}=g_{k1}`$ modulo $`N_k`$. In the group $`H_k`$, any element in $`H_{k1}^{a_k}`$ commutes with any element in $`H_{k1}`$. Therefore, $`g_{k1}`$ centralizes $`H_{k1}`$ modulo $`N_k`$. Then, modulo $`N_k`$, any element in the normal closure of $`g_{k1}`$ is some power of $`g_{k1}`$. In particular, this is true for the element $`g_{k1}(n)`$ by the inductive assumption. Since $`g_{k1}`$ and $`a_k`$ commute modulo $`N_k`$, we deduce that $`g_k(n)=[g_{k1}(n),a_k^n]`$ equals $`1`$ in the quotient group $`H_k/N_k`$, that is, $`g_k(n)N_k`$. Now we prove that $`\varphi _k(g_k(n))=n^k`$ for $`1kd`$ by induction on $`k`$. For $`k=1`$ we get $`\varphi _1(g_1(n))=\varphi _1(a_1^n)=n\varphi _1(a_1)=n`$. Let $`k>1`$; suppose that $`\varphi _{k1}(g_{k1}(n))=n^{k1}`$. Then $`\varphi _k(g_k(n))=\varphi _k([g_{k1}(n),a_k^n])=\varphi _k(g_{k1}(n)^1)+\varphi _k(g_{k1}(n)^{a_k^n})=n\varphi _{k1}(g_{k1}(n))=nn^{k1}=n^k`$ (we have used the properties of $`\varphi _k`$, the fact that $`g_{k1}(n)M_{k1}`$ and the inductive assumption). The Theorem is proved. It is an interesting question what else functions may be distortion functions of finitely generated subgroups of $`F`$. In particular, it is very interesting if such a distortion function may not have a recursive upper bound. Let us give an equivalent form of this problem. ###### Problem 8 Does R. Thompson’s group $`F`$ have a finitely generated subgroup with unsolvable membership problem?
warning/0005/hep-th0005115.html	ar5iv	text	# 1 Introduction ## 1 Introduction The $`c`$-theorem establishes the irreversibility of the renormalization-group flow in two dimensions ; it shows that there exists a positive function of the coupling constants which monotonically decreases along the flows and is stationary at the fixed points. The existence of an analogous theorem in four and higher dimensions has been extensively discussed in the recent literature: several non-trivial flows have been checked in four-dimensional supersymmetric gauge theories (most notably in the “conformal window”) , and via the AdS/CFT correspondence ; moreover, two proofs of the theorem have been presented in the Refs.. Most of the examples and conjectures identify the $`c`$-function at the fixed points with the coefficient of the Euler density $`G`$ in the trace anomaly (such that $`\chi =G`$ is the Euler characteristic). This coefficient, called $`a`$, has definite sign for all unitary theories , and verify the inequality $`a_{UV}>a_{IR}`$ for all known examples of flows connecting pairs of UV and IR fixed points. According to the analysis of Ref., the trace anomaly in any even dimension $`d=2k`$ contains three parts: the topological invariant term, given by the Euler density; the Weyl invariant terms, made by polynomials of the Weyl tensor and its derivatives, and finally the variations of local counterterms which can be neglected. On conformally-flat geometries, like the $`d`$-dimensional sphere $`S^{2k}`$, the Weyl tensor vanishes and the trace anomaly is completely expressed in terms of the Euler characteristic. In this paper we calculate the trace anomaly on the sphere $`S^{2k}`$ for the free massless scalar, Dirac spinor and antisymmetric tensor fields in any dimension and obtain the corresponding values of $`a`$. We use the standard zeta-function regularization of the partition function . Next, we analyse the behaviour of $`a`$ upon varying the spin and the dimension and test its interpretation as a measure of degrees of freedom in field theory. In the large-$`d`$ limit, we find that all anomalies, once properly normalized, go to zero; this is consistent with the expected semiclassical behaviour. However, the ratios $`a(\sigma )/a(0)`$ of the spinful to the scalar values, which have the meaning of counting functions, do not approach the classical number of fields components: each component of the antisymmetric tensor field weights of order $`O(d^3)`$ with respect to the scalar field, and each fermion component weights $`O(d)`$. Therefore, the anomaly $`a`$ yields a “quantum” measure of degrees of freedom which is rather far from the classical intuition, with the antisymmetric tensor dominating the counting for any $`d4`$. At first sight, we could doubt of this interpretation as a measure of degrees of freedom; however, a closer inspection shows that this is consistent with known facts and renormalization-group flows. The same enhancement is found for the gravitational chiral anomaly and for another term in the trace anomaly, whose coefficient $`c`$ also normalizes the stress-tensor two-point function<sup></sup><sup></sup>* $`c`$ is always positive, but is not always decreasing . . Furthermore, the observed dominance of the antisymmetric tensor field is crucial for the existence of theories with $`ca`$ in any dimension ; these theories naturally arise in the AdS/CFT correspondence and have been conjectured to satisfy a version of $`c`$-theorem with all its two-dimensional features . ## 2 Trace anomaly of the scalar and Dirac fermion fields The stress tensor $`T_{\mu \nu }`$ of a $`d`$-dimensional field theory is defined by the variation of the partition function with respect to the the background metric $`g_{\mu \nu }`$: $$\delta \mathrm{log}Z[g_{\mu \nu }]=\frac{1}{2V_{d1}}𝑑x^d\sqrt{g}T_{\mu \nu }\delta g^{\mu \nu },V_{d1}=\frac{2\pi ^{d/2}}{\mathrm{\Gamma }(\frac{d}{2})},$$ (2.1) where the conventional normalization factor $`V_{d1}`$ is the volume of the unit sphere $`S^{d1}`$. We shall compute the partition function of $`d`$-dimensional free field theories on the sphere $`S^{2k}`$ ($`d=2k`$) and obtain the integrated trace anomaly by the scale variation $`g_{\mu \nu }\mathrm{exp}(2\alpha )g_{\mu \nu }`$ ($`\alpha =\mathrm{const}.`$): $$\frac{d}{d\alpha }\mathrm{log}Z\left[S^{2k}\right]=\frac{1}{V_{2k1}}_{S^{2k}}\mathrm{\Theta }=\frac{V_{2k}r^{2k}}{V_{2k1}}\mathrm{\Theta }.$$ (2.2) In this Equation, $`\mathrm{\Theta }`$ is a short-hand notation for the trace $`T_\mu ^\mu `$ and $`r`$ is the radius of the sphere. The (Euclidean) action for bosonic free fields is written in terms of the covariant Laplacian $`\mathrm{\Delta }`$, whose explicit form depends on the spin of the field: $$S[g,\varphi ]=\frac{1}{2}\sqrt{g}\varphi (x)\mathrm{\Delta }\varphi (x).$$ (2.3) The partition function, i.e. the determinant of the Laplacian, can be obtained from the analytic continuation of the zeta-function associated to the spectrum of $`\mathrm{\Delta }`$ on $`S^{2k}`$, $$\zeta _\mathrm{\Delta }(s)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{\delta _n}{\lambda _n^s};$$ (2.4) this function is defined in terms of the eigenvalues $`\lambda _n`$ and their degeneracies $`\delta _n`$. We consider theories which are invariant under conformal (Weyl) transformations $`g_{\mu \nu }(x)\mathrm{exp}2\sigma (x)g_{\mu \nu }(x)`$ at the classical level ; the corresponding Laplacian is conformally covariant, $`\mathrm{\Delta }\mathrm{exp}((\delta +2)\sigma )\mathrm{\Delta }\mathrm{exp}(\delta \sigma )`$, with $`\delta `$ the dimension of $`\varphi `$. In these theories, the scale dependence of the partition function (2.2) is purely anomalous, i.e. is induced by the regularization of the determinant and it is given by the zeta-function analytically continued at $`s=0`$ : $$\frac{d}{d\alpha }\mathrm{log}Z\left[S^{2k}\right]=\zeta _\mathrm{\Delta }(0)=\frac{1}{V_{2k1}}_{S^{2k}}\mathrm{\Theta }.$$ (2.5) In the case of fermionic fields, we consider the action (2.3) with $`\mathrm{\Delta }`$ the square of the Dirac operator and obtain the anomaly (2.5) in terms of $`\zeta _\mathrm{\Delta }(0)`$. We now evaluate the trace anomaly for the scalar field This calculation and the following one for the Dirac fermion are not new , but we present here rather compact results., whose conformal covariant Laplacian is : $$\mathrm{\Delta }_0=^2+\frac{k1}{2(2k1)},$$ (2.6) with $``$ the scalar curvature. The eigenvalues and degeneracies of this operator on $`S^{2k}`$ are given by $`\lambda _l`$ $`=`$ $`(l+k)(l+k1),l=0,1,2,\mathrm{}`$ $`\delta _l`$ $`=`$ $`(2l+2k1){\displaystyle \frac{(l+2k2)!}{l!(2k1)!}}.`$ (2.7) The zeta-function of the scalar Laplacian $`\zeta _{\mathrm{\Delta }_0}(s)`$ (2.4,2.7) can be analytically continued at $`s=0`$ by expressing it in terms of ordinary Riemann zeta-functions $`\zeta (s)`$. It is convenient to use the analytic continuation formula found in Ref.: $`{\displaystyle \underset{\sigma =n+1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{2\sigma +1}{\left[(\sigma n)(\sigma +n+1)\right]^s}}`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Gamma }(s)}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\mathrm{\Gamma }(s+m1)}{\mathrm{\Gamma }(m+1)}}(2s+m2)(2n+1)^m`$ (2.8) $`\times \left[\zeta (2s+m1){\displaystyle \underset{l=1}{\overset{2n+1}{}}}l^{2sm+1}\right].`$ We rewrite the zeta-function of the scalar Laplacian as follows (for $`k2`$): $$\zeta _{\mathrm{\Delta }_0}(s)=\frac{1}{(2k1)!}\underset{i=1}{\overset{k1}{}}\alpha _i\underset{\sigma =k1}{\overset{\mathrm{}}{}}\frac{2\sigma +1}{[\sigma (\sigma +1)]^{si}},$$ (2.9) in terms of the coefficients $`\alpha _i`$ defined by $$\frac{(\sigma +k1)!}{(\sigma k+1)!}=\underset{i=0}{\overset{k2}{}}\left[\sigma (\sigma +1)i(i+1)\right]=\underset{i=1}{\overset{k1}{}}\alpha _i[\sigma (\sigma +1)]^i,(k2),$$ (2.10) where $`\sigma =l+k1`$. We use the analytic continuation (2.8) for all the terms in (2.9), which reduce to finite sums involving the Bernoulli numbers $`B_n`$ in the limit $`s=0`$: $$\zeta _{\mathrm{\Delta }_0}(0)=\frac{1}{(2k1)!}\underset{i=1}{\overset{k1}{}}\alpha _i\underset{m=0}{\overset{i+1}{}}\frac{i!}{m!(im+1)!}B_{2i+2m}.$$ (2.11) This expression can be cast into the following compact form: $$\zeta _{\mathrm{\Delta }_0}(0)=\frac{1}{(2k1)!}_0^{B(1+B)}𝑑t\underset{i=0}{\overset{k2}{}}\left[ti(i+1)\right],(d=2k4),$$ (2.12) where it is understood that $`B^n`$ should be replaced with $`B_n`$ after evaluation of the integral. The values of $`\zeta _{\mathrm{\Delta }_0}`$ for the lowest dimensions $`d=\{4,6,\mathrm{},14\}`$ are: $$\zeta _{\mathrm{\Delta }_0}(0)=\{\frac{1}{90},\frac{1}{756},\frac{23}{113400},\frac{263}{7484400},\frac{133787}{20432412000},\frac{157009}{122594472000}\}.$$ (2.13) Finally, the trace anomaly of the scalar field is obtained by replacing (2.12) into (2.2) (the result is in agreement with Ref.). We postpone the discussion of these numbers to Section $`4`$ and move to the analogous calculation for the Dirac field. The spinor Laplacian is obtained from the square of the Dirac operator in the $`g_{\mu \nu }`$ background, $$\mathrm{\Delta }_{1/2}=/^{2}=^2+\frac{}{4},$$ (2.14) and it is conformally covariant ; its eigenvalues and degeneracies on the sphere are : $`\lambda _l`$ $`=`$ $`(l+k)^2,l=0,1,2,\mathrm{}`$ $`\delta _l`$ $`=`$ $`2^{k+1}{\displaystyle \frac{(2k+l1)!}{l!(2k1)!}}.`$ (2.15) The analytic continuation of the corresponding zeta-function is again obtained by splitting it into simpler functions ($`k1`$): $$\zeta _{\mathrm{\Delta }_{1/2}}(s)=\frac{2^{k+1}}{(2k1)!}\underset{l=0}{\overset{\mathrm{}}{}}\frac{(l+2k1)!}{l!}\frac{1}{(l+k)^{2s}}=\frac{2^{k+1}}{(2k1)!}\underset{i=0}{\overset{k1}{}}\beta _i\zeta (2s2i1),$$ (2.16) where the coefficients $`\beta _i`$ are defined by $$\frac{(\sigma +k1)!}{(\sigma k)!}=\sigma \underset{i=1}{\overset{k1}{}}\left(\sigma ^2i^2\right)=\underset{i=0}{\overset{k1}{}}\beta _i\sigma ^{2i+1},$$ (2.17) with $`\sigma =l+k`$. The limit $`s=0`$ can be taken into Eq.(2.17) and the result can be written again in compact form: $$\zeta _{\mathrm{\Delta }_{1/2}}(0)=\frac{2^k}{(2k1)!}_0^{B^2}𝑑t\underset{i=1}{\overset{k1}{}}(ti^2),\left(B^nB_n\right)$$ (2.18) (for $`k=1`$ the integrand is one). The first few values for $`d=\{4,6,\mathrm{},14\}`$ are: $$\zeta _{\mathrm{\Delta }_{1/2}}(0)=\{\frac{11}{90},\frac{191}{3780},\frac{2497}{113400},\frac{14797}{1496880},\frac{92427157}{20432412000},\frac{36740617}{17513496000}\}.$$ (2.19) ## 3 Trace anomaly of the antisymmetric tensor field Another conformal invariant free theory in $`2k`$ dimensions is given by the antisymmetric $`(k1)`$-form, called antisymmetric tensor field , which generalizes the four-dimensional vector field. Antisymmetric tensor fields commonly appear in supergravity and in superstring theory; a well known example is the low-energy world-volume theory of the $`M5`$-brane described by a $`d=6`$, $`𝒩=(2,0)`$ tensor multiplet that contains one antisymmetric tensor $`B_{\mu \nu }`$. In general, one could consider any $`p`$-form $`A`$, with $`p=1,\mathrm{},k1`$, and write its action in terms of the field strength $`F=dA`$ as follows (using the notations of Ref.): $$S=\frac{1}{2p!}FF;$$ (3.1) however, this is conformal invariant for $`p=k1`$ only. For general $`p`$, this action is invariant under the gauge transformation $`AA+d\mu `$, with $`\mu `$ a $`(p1)`$-form, thus we should introduce gauge-fixing terms and ghost fields. The ghost action is also gauge invariant, so we should consider ghosts of ghosts: the resulting tower of ghost fields contains all the forms of lower degree, till the $`0`$-forms which have no further gauge invariance. This quantization problem is well understood and the resulting quantum action takes the form: $$S=\frac{1}{2}\underset{i=0}{\overset{p}{}}\frac{1}{(pi)!}A_{pi}\mathrm{\Delta }_{pi}^{i+1}A_{pi},$$ (3.2) where the $`(pi)`$-form $`A_{pi}`$ for $`i`$ even (odd) is a bosonic (fermionic) field and $`\mathrm{\Delta }_{pi}`$ is the Hodge-de Rham operator $`\mathrm{\Delta }=d\delta +\delta d`$ expressed in terms of the exterior derivative $`d`$ and the coderivative $`\delta `$. It follows that the partition function for the antisymmetric tensor theory ($`p=k1`$) in $`d=2k`$ is given by the following product of determinants: $$Z_{AT}=\frac{1}{det^{1/2}(\mathrm{\Delta }_{k1})}\frac{det(\mathrm{\Delta }_{k2})}{det^{3/2}(\mathrm{\Delta }_{k3})}\mathrm{}\left(\frac{det^{(k1)/2}(\mathrm{\Delta }_1)}{det^{k/2}(\mathrm{\Delta }_0)}\right)^{(1)^{k1}}.$$ (3.3) On the geometry of the sphere, the eigenvalues and degeneracies for exact $`(p+1)`$-forms and coexact $`p`$-forms coincide , so we can simplify the ratios in the previous expression by restricting the determinants to the spectra of coexact forms: $$Z_{AT}=\left[\frac{1}{det(\mathrm{\Delta }_{k1}^{ce})}\frac{det(\mathrm{\Delta }_{k2}^{ce})}{det(\mathrm{\Delta }_{k3}^{ce})}\mathrm{}\left(\frac{det(\mathrm{\Delta }_1^{ce})}{det(\mathrm{\Delta }_0^{ce})}V_{2k}\right)^{(1)^{k1}}\right]^{1/2},$$ (3.4) where the superscript in $`\mathrm{\Delta }^{ce}`$ means the restriction to coexact forms. Note the volume factor in the last term that arises from the regularization of the zero mode of the $`0`$-form, which does not cancel out in the last ratio of Eq.(3.3). The eigenvalues and degeneracies of the Hodge-de Rham operator on coexact $`p`$-forms in $`S^{2k}`$ are given by : $`\lambda _{p,l}`$ $`=`$ $`(l+p)(l+2kp1),l=1,2,\mathrm{}`$ $`\delta _{p,l}`$ $`=`$ $`{\displaystyle \frac{(l+2k1)!}{p!(2kp1)!(l1)!}}{\displaystyle \frac{2l+2k1}{(l+p)(l+2kp1)}}.`$ (3.5) The calculation of the trace anomaly thus involves the following alternate sum of zeta-functions: $$\frac{d}{d\alpha }\mathrm{log}Z_{AT}\zeta _{AT}(0)=\underset{p=0}{\overset{k1}{}}(1)^{kp1}\zeta _{\mathrm{\Delta }_p^{ce}}(0)+(1)^{k1}k.$$ (3.6) Each zeta-function can be analytically continued with the help of Eq. (2.8); after some simplifications, the result reads: $$\zeta _{AT}(0)=\underset{q=1}{\overset{k}{}}\frac{(1)^q}{(kq)!(k+q1)!}\underset{i=1}{\overset{k}{}}\alpha _i(q)\underset{m=0}{\overset{i}{}}(2q1)^m\frac{(i1)!}{m!(im)!}B_{2im},$$ (3.7) where the coefficients $`\alpha _i(q)`$, with $`q=kp`$, are defined by the expression: $`{\displaystyle \frac{(\sigma +k)!}{(\sigma k)!}}`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{k1}{}}}\left[(\sigma +q)(\sigma q+1)i(i+1)q(1q)\right]`$ (3.8) $`=`$ $`{\displaystyle \underset{i=1}{\overset{k}{}}}\alpha _i(q)\left[(\sigma +q)(\sigma q+1)\right]^i.`$ The result (3.7) can again be written in the compact form (using $`B^nB_n`$): $$\zeta _{AT}(0)=\underset{q=1}{\overset{k}{}}\frac{(1)^q}{(kq)!(k+q1)!}_0^{B(B+2q1)}\frac{dt}{t}\underset{i=0}{\overset{k1}{}}\left[ti(i+1)q(1q)\right].$$ (3.9) The first few values of this quantity for $`d=\{4,6,\mathrm{},12\}`$ are: $$\zeta _{AT}(0)=\{\frac{62}{90},\frac{221}{210},\frac{8051}{5670},\frac{1339661}{748440},\frac{525793111}{243243000},\frac{3698905481}{1459458000}\}.$$ (3.10) For $`k=2`$ this matches the well-known trace anomaly of the vector field , and for $`k=3`$ it checks the recent result for the two-form gauge field in Ref.. ## 4 Anomaly versus number of degrees of freedom The trace anomaly in $`2k`$ dimensions contains the Euler density $`G`$ and a number of conformal covariant polynomials which involve the Weyl tensor $`W`$ and its derivatives, and have dimension $`2k`$ : $$\mathrm{\Theta }(x)=\lambda \left(aG(x)cW(x)\mathrm{\Delta }^{k2}W(x)+O(W^3)+\mathrm{}\right).$$ (4.1) In this Equation, we called $`a`$ the coefficient of the Euler density $`G`$ and $`c`$ that of the term quadratic in the Weyl tensor, with $`\mathrm{\Delta }`$ a conformal covariant Laplacian, and we left unspecified the other conformal covariant polynomials. The general form (4.1) can be integrated on conformally flat manifolds $``$ to yield: $$\frac{d}{d\alpha }\mathrm{log}Z\left[\right]=\frac{\lambda }{V_{2k1}}a\chi =\zeta _\mathrm{\Delta }(0)\frac{\chi }{2},$$ (4.2) with $`\chi =G`$ the Euler characteristic. In the right part of this Equation, we also matched the general expression with the previous calculation on the sphere, for which $`\chi \left(S^{2k}\right)=2`$. The normalization coefficient $`\lambda `$ remains to be chosen for the complete definition of $`a`$. In the following, we shall argue that there exist an absolute, i.e. dimension and scale independent normalization for $`a`$ (call it $`a=\widehat{a}`$), owing to the topological invariance of the Euler characteristic, and a relative, $`d`$-dependent normalization suitable for using $`a`$ as a counting function. Let us discuss the first definition. The Euler characteristic $`\chi \left(\right)`$ has a natural normalization because it takes integer values: it counts the number $`b_p`$ of non-trivial zero modes of the Hodge-De Rham operator $`\mathrm{\Delta }_p`$ on the manifold $``$, more precisely their alternating sum $`b_0b_1+b_2\mathrm{}+b_{2k}`$ . On the other side of Eq.(4.2), the scale derivative of $`\mathrm{log}Z`$ is also a number free from ambiguities, being the finite part after regularization of the “number of positive modes” of a given Laplacian on $``$. The ratio of these two numbers is an interesting dimension and scale independent quantity; thus, we identify it with $`a=\widehat{a}\zeta _\mathrm{\Delta }(0)`$, upon choosing the normalization $`\lambda =V_{2k1}/2`$. This discussion suggests a similarity between the topological part of the trace anomaly and the chiral anomaly, which counts the number of certain zero modes on $``$, as a consequence of the index theorem for the Dirac operator . We have seen that the Laplacians of the free fields in Sections $`2`$ and $`3`$ do not have zero modes on the sphere, but nevertheless their trace anomaly is proportional to $`\chi `$ which counts zero modes. Suppose instead that $`\mathrm{\Delta }`$ had $`n`$ zero modes: they would be regularized by an infra-red cut-off $`\rho `$ and would contribute to the partition function a factor $`Z\rho ^n`$, leading nicely to $`\widehat{a}=n`$. For the actual $`\mathrm{\Delta }`$ without zero modes, we can think that the ultra-violet regularization of the non-zero modes produces “effective” zero modes. In conclusion, $`\widehat{a}=\zeta _\mathrm{\Delta }(0)`$ can be interpreted as counting effective zero modes of the Laplacian of the theory. We now discuss the dependence of this quantity on the dimension and the spin. Equation (2.13) shows some values of $`\widehat{a}`$ for the scalar theory; one sees that their sign alternates with $`k`$ and that they decrease (exponentially) with the dimension. Actually, the large $`d`$ limit is equivalent to the semiclassical limit for the scalar theory, and the anomaly should go to zero. The semiclassical limit can be easily understood by defining the theory on a space-time lattice: the field variables are located on the lattice sites, and the discretized Laplacian of the field at one point $`\mathrm{\Delta }\varphi (x_o)`$ is the sum of all the $`2d`$ field variables on the nearest neighbour sites. As the value of $`d`$ increases, the action approaches its mean-field approximation. By the same argument we expect that the values of $`\widehat{a}`$ for the fermion and antisymmetric tensor fields should also vanish for large $`d`$. This is indeed the case: we should take the values in Eqs.(2.19,3.10) and divide them for the corresponding number of on-shell field components (number of independent polarizations), which are : $$n(F)=2^k,n(AT)=\frac{(2k2)!}{\left[(k1)!\right]^2}.$$ (4.3) In conclusion, the trace anomaly $`\widehat{a}`$ is a measure of effective zero modes and decreases for large $`d`$ in absolute units as expected. We now analyse the coefficient $`a`$ as a measure of degrees of freedom in field theory; this interpretation would be implied by the eventual proof of the $`c`$-theorem. In this case, we should normalize $`a(S)=1`$ for the scalar field in any dimension, by choosing $`\lambda `$ in (4.2) accordingly. The values of $`a`$ in this normalization are reported in the first three lines of Table(1); the last three lines instead contain the relative weights per field component: $$r(\sigma )=\frac{a(\sigma )}{n(\sigma )},\sigma =S,F,AT.$$ (4.4) The last column in Table (1) gives their asymptotic behaviour for large dimension, which were determined numerically from Eqs.(2.12,2.18,3.9) (more precisely, $`r(AT)0.943(2k)^3`$). We find that these ratios do not go to one, as we would naively expect in the semiclassical limit, but grow, in particular for the antisymmetric tensor: already at $`d=4`$ the vector field has a rather large weight. Therefore, this “quantum” measure of degrees of freedom is always different from the classical number of field components – a fact that is rather counter-intuitive. Apart from this physical interpretation, the behaviour of $`a`$ is not surprising because it also occurs for other terms in the trace and gravitational anomalies. This means that the spinful theories approach the semiclassical limit more slowly than the scalar theory, or that they interact more strongly with the background metric. Let us first discuss the other trace anomaly coefficient $`c`$ which was introduced in Eq. (4.1); this number also gives the normalization of the stress tensor correlator in flat space and has been already computed in the Refs. for the free theories considered here. In the same normalization $`c(S)=1`$, the results are: $$c(F)=2^{k1}(2k1),c(AT)=\frac{(2k)!}{2[(k1)!]^2},(d=2k).$$ (4.5) In Table (2) we display the first few values of $`c`$ as well as the corresponding weights per field component $`r^{}(\sigma )=c(\sigma )/n(\sigma )`$. As we anticipated, these ratios grow with the dimension, although with milder asymptotic behaviours. The same enhancement is observed in the chiral anomaly of pure gravitational origin, which occurs for dimensions $`d=2+4n`$. The comprehensive study of Ref. considers the following chiral fields: the Weyl fermion (one half of the Dirac fermion), the gravitino and the self-dual antisymmetric tensor (one half of that considered here). The complete expression of this chiral anomaly contains several independent terms, made by the different traces of the product of $`(2+2n)`$ Riemann two-forms, and their coefficients span a wide range of values . The larger coefficient is found for the genuine $`d=2+4n`$ term, which is the $`(n+1)`$-th Pontryagin class $`p_{n+1}()=\mathrm{Tr}(R^{2n+2})+\mathrm{}`$. We have computed the corresponding weights per field component for the three field types considered in Ref. and for $`d=6,10,14`$: these weights again grow with the dimension, the faster the higher spin value; moreover, the relative growth between the antisymmetric tensor and the fermion is much larger than the order $`O(d^2)`$ found for the coefficient $`a`$. The origin of this enhancement is not immediately apparent in the diagrammatic calculation of Ref., because the coupling of spin to gravity does not have a standard group-theoretic form valid for all spin values. Nevertheless, the number of external legs in the anomalous loops increase with the dimension and the algebraic complexity enhances the spin dependence. ## 5 Discussion Although puzzling, the measure of degrees of freedom given by $`a`$ could be consistent with the renormalization group flow. Here we can check a couple of standard renormalization-group patterns and discuss a $`c`$-theorem conjecture which is consistent with the observed enhancement. A well-known flow is provided by a (generalized) gauge theory which has massive gauge fields in the infrared, together with massless scalars due to spontaneous symmetry breaking of a global symmetry; our results imply a very large value for $`a_{UV}`$, which better fulfills the $`c`$-theorem inequality $`a_{UV}>a_{IR}`$. Another example is the Higgs phenomenon, in which scalar degrees of freedom become extra components for the massive antisymmetric tensor. For large $`d`$, these extra components contribute to $`a`$ far less than the original components of the massless field: anyhow, no contradiction arises in going back to the UV limit, because the massless limit of the massive gauge theory is different from the massless theory. In the Reference , it has been conjectured that the canonical form of the two-dimensional $`c`$-theorem could extend to the class of higher dimensional theories characterized by $`c=\nu a`$, with $`\nu (d)`$ a given constant. The field theories described by the AdS/CFT correspondence belong to this class ; for instance, the vector multiplet of $`𝒩=4`$ supersymmetry in four dimensions. In our notations $`a(S)=c(S)=1`$, the equation $`c=\nu a`$ of Ref. reads: $$c(2k+1)\frac{\mathrm{\Gamma }(2k)}{\mathrm{\Gamma }(k)^2}\left\|\zeta _{\mathrm{\Delta }_0}(0)\right\|a=0,(d=2k).$$ (5.1) Since the scalar, fermion and antisymmetric tensor fields contribute differently to $`a`$ and $`c`$, and with definite sign, there is no a-priori guarantee that this equation admits solutions in higher dimensions. We can easily obtain the explicit form of this equation in each dimension by using the results in the two Tables; for $`d=4,6,8,10`$, we find: $$\begin{array}{ccccccc}\hfill d=4:& & & \hfill 2N_S+& \hfill 7N_F& \hfill 26N_{AT}=& \hfill 0;\\ \hfill d=6:& & & \hfill 13N_S+& \hfill 169N_F& \hfill 2358N_{AT}=& \hfill 0;\\ \hfill d=8:& & & \hfill 67N_S+& \hfill 2543N_F& \hfill 110620N_{AT}=& \hfill 0;\\ \hfill d=10:& & & \hfill 817N_S+& \hfill 81535N_F& \hfill 9994610N_{AT}=& \hfill 0,\end{array}$$ (5.2) where $`(N_S,N_F,N_{AT})`$ are the field multiplicities. We see that the dominance of $`a(AT)`$ yields a crucial minus sign in these equations, which ensures a solution in any dimension. Therefore, we have found that the behaviour of $`a`$ makes the $`c`$-theorem conjecture possible in any dimension. The integer equations ( 5.2) has two independent solutions $`(N_S,N_F,N_{AT})`$, which can be thought of as being vectors generating a two-dimensional sub-lattice of a three-dimensional integer lattice. Some sample solutions, with minimal values for $`N_{AT}`$, are the following: $$\begin{array}{cccc}d=4:\hfill & & (6,2,1),\hfill & (13,0,1);\hfill \\ d=6:\hfill & & (161,169,13),\hfill & (265,161,13);\hfill \\ d=8:\hfill & & (835,65,2),\hfill & (2524,64,3);\hfill \\ d=10:\hfill & & (11950,248,3),\hfill & (47095,141,5).\hfill \end{array}$$ Acknowledgments We thank D. Anselmi, A. Coste, R. Guida, N. Magnoli and A. Schwimmer for interesting discussions on the $`c`$-theorem and the trace anomaly. This work is supported in part by the European Community Network grant FMRX-CT96-0012.
warning/0005/gr-qc0005050.html	ar5iv	text	# Photon surfaces ## 1. Introduction The exterior region of the maximally extended Schwarzschild space-time is described by the metric $$𝐠=\left(1\frac{2m}{r}\right)dt^2+\left(1\frac{2m}{r}\right)^1dr^2+r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2\right),r>2m.$$ (1) For any null geodesic in this exterior region the null geodesic equations give $$\frac{d^2r}{d\lambda ^2}=\left(r3m\right)\left\{\left(\frac{d\theta }{d\lambda }\right)^2+\mathrm{sin}^2\theta \left(\frac{d\varphi }{d\lambda }\right)^2\right\}$$ (2) where $`\lambda `$ is an affine parameter along the geodesic. The right side here is evidently positive for $`r>3m`$ and negative for $`r:2m<r<3m`$. It follows that any future endless null geodesic in the maximally extended Schwarzschild space-time starting at some point with $`r>3m`$ and initially directed outwards, in the sense that $`dr/d\lambda `$ is initially positive, will continue outwards and escape to infinity. Any future endless null geodesic in the maximally extended Schwarzschild space-time starting at some point with $`r:2m<r<3m`$ and initially directed inwards, in the sense that $`dr/d\lambda `$ is initially negative, will continue inwards and fall into the black hole. The hypersurface $`\{r=3m\}`$, known as the Schwarzschild photon sphere, thus distinguishes the borderline between these two types of behaviour; any null geodesic starting at some point of the photon sphere and initially tangent to the photon sphere will remain in the photon sphere. (See Darwin for a detailed analysis of the behaviour of null and timelike geodesics in Schwarzschild space-time.) The Schwarzschild photon sphere also has physical significance for massive bodies. For any timelike geodesic in the exterior region the geodesic equations give $$\frac{d^2r}{ds^2}=\frac{m}{r^2}+\left(r3m\right)\left\{\left(\frac{d\theta }{ds}\right)^2+\mathrm{sin}^2\theta \left(\frac{d\varphi }{ds}\right)^2\right\}$$ (3) where $`s`$ is arc length along the geodesic. At any point with $`r>3m`$ one may arrange for the two terms on the right of (3) to cancel and so obtain a timelike geodesic at constant $`r`$. For $`r:2m<r<3m`$ the right hand side of (3) is evidently negative. Thus any future endless timelike geodesic in the maximally extended Schwarzschild space-time starting at some point between the event horizon at $`r=2m`$ and the photon sphere at $`r=3m`$ and initially directed inwards, in the sense that $`dr/ds`$ is initially negative, will continue inwards and fall into the black hole. Any observer who traverses a Schwarzschild photon sphere must therefore engage some form of propulsion or else be drawn in to the black hole to meet an inevitable fate. A photon sphere has been defined by Virbhadra & Ellis as a timelike hypersurface of the form $`\{r=r_0\}`$ where $`r_0`$ is the closest distance of approach for which the Einstein bending angle of a light ray is unboundedly large. These authors subsequently considered the Einstein deflection angle for a general static spherically symmetric metric and obtained an equation for a photon sphere. The existence of a photon sphere in a space-time has important implications for gravitational lensing. In any space-time containing a photon sphere, gravitational lensing will give rise to relativistic images . The Schwarzschild photon sphere may be usefully be compared with the concept of a closed trapped surface. Any null geodesic originating from any point on a closed trapped surface in Schwarzschild space-time is drawn into the singularity at $`r=0`$. By contrast, any null geodesic originating from any point on the photon sphere will be drawn into the singularity if and only if it is initially directed inwards. The main objectives of the present paper are to give a geometric definition of a photon surface in a general space-time and of a photon sphere in a general static spherically symmetric space-time. An evolution equation is obtained for the cross-sectional area of a photon surface in a dynamic spherically symmetric space-time. It is shown, subject to suitable energy conditions, that in any static spherically symmetric space-time a black hole must be surrounded by a photon sphere, and a photon sphere must surround either a black hole, a naked singularity or more than a certain amount of matter. Many examples are given of photon spheres in static spherically symmetric space-times. Photon surface evolution is considered for the dynamic space-time example of Vaidya null dust collapse to a naked singularity. ## 2. Photon surfaces The hypersurface $`S:=\{r=3m\}`$ in Schwarzschild space-time has two main properties, first that any null geodesic initially tangent $`S`$ will remain tangent to $`S`$, and second that $`S`$ does not evolve with time. The following general definition of a photon surface is based on only the first of these properties. A more restrictive class of photon surfaces may be defined when the space-time admits a group of symmetries (see Definition 2.3). ###### Definition 2.1. A photon surface of $`(M,𝐠)`$ is an immersed, nowhere-spacelike hypersurface $`S`$ of $`(M,𝐠)`$ such that, for every point $`pS`$ and every null vector $`𝐤T_pS`$, there exists a null geodesic $`\gamma :(ϵ,ϵ)M`$ of $`(M,𝐠)`$ such that $`\dot{\gamma }(0)=`$k$`,\|\gamma \|S`$. Any null hypersurface is trivially a photon surface. Photon surfaces are conformally invariant structures. If $`S`$ is a photon surface of $`(M,𝐠)`$ then $`S`$ is a photon surface of $`(M,\mathrm{\Omega }^2𝐠)`$ for any smooth function $`\mathrm{\Omega }:M(0,\mathrm{})`$. Note that Definition 2.1 is entirely local. In particular, a photon surface $`S`$ need contain no endless null geodesics of $`(M,𝐠)`$. Moreover, a photon surface need only be immersed, rather than embedded in $`M`$, and so may have self-intersections. If $`(M,𝐠)`$ is of dimension $`n+1`$ ($`n2`$) then, through each point $`p`$ of a photon surface $`S`$ in $`(M,𝐠)`$, there is an $`(n2)`$-parameter family of null geodesics of $`(M,𝐠)`$ that lie entirely in $`S`$. The paper will be principally concerned with photon surfaces in space-times of $`3+1`$ dimensions. The exceptions are Examples 1 and 3 which give photon surfaces in space-times of dimension $`2+1`$ and $`4+1`$ respectively. ###### Example 1 (Minkowski 3-space). In Minkowski 3-space $`𝕄^3`$, consider the single-sheeted hyperboloid $`S`$ given by $$t^2+x^2+y^2=a^2$$ (4) for some constant $`a>0`$. This surface is doubly ruled, the rulings being given by $$\gamma _\theta ^\pm \left(t\right):=a(0,\mathrm{cos}\theta ,\mathrm{sin}\theta )+at(1,\mathrm{sin}\theta ,\pm \mathrm{cos}\theta )$$ (5) ($`\mathrm{}<t<\mathrm{},0\theta <2\pi `$), where $`\theta `$ identifies the intersection points with $`\{t=0\}`$ and $`t`$ is the parameter along the ruling lines. The tangents $`\dot{\gamma }_\theta ^\pm \left(t\right)`$ to the ruling lines are null with respect to the $`𝕄^3`$ metric. Clearly they are geodesics in $`𝕄^3`$. At each point of $`S`$ there can be just two null directions tangent to $`S`$. These must therefore be the directions of the two ruling lines through that point. Hence $`S`$ is a photon surface in the sense of Definition 2.1 (see Fig. 1). Note that for any circle of the form $$C=\{t_0,x_0+r\mathrm{cos}\theta ,y_0+r\mathrm{sin}\theta \},r>0,$$ (6) and any future-directed timelike vector field $`𝐗`$ along $`C`$ that respects the symmetry of $`C`$, in the sense of $$𝐗=(X^t,X^r\mathrm{cos}\theta ,X^r\mathrm{sin}\theta ),$$ (7) for constant $`X^t>0`$, $`X^r`$ such that $`(X^t)^2>(X^r)^2`$, there is a unique single-sheeted hyperboloid $`S`$ through $`C`$ such that $`𝐗`$ is tangent to $`S`$ along $`C`$. In the case $`a=0`$, equation (4) gives the null cone through the origin. The complement of $`p`$ in this null cone is a null photon surface of $`𝕄^3`$. ###### Example 2 (Minkowski 4-space). One may generalize Example 1 to Minkowski 4-space $`𝕄^4`$ as follows. Let $`S`$ be a timelike hypersurface in $`𝕄^4`$ of the form $$t^2+x^2+y^2+z^2=a^2$$ (8) for some constant $`a>0`$. The two-parameter family of lines $$\gamma _{\theta ,\varphi }^\pm \left(t\right)=a(0,\mathrm{cos}\theta ,\mathrm{sin}\theta \mathrm{sin}\varphi ,\mathrm{sin}\theta \mathrm{cos}\varphi )+at(1,\mathrm{sin}\theta ,\pm \mathrm{cos}\theta \mathrm{sin}\varphi ,\pm \mathrm{cos}\theta \mathrm{cos}\varphi )$$ (9) foliate $`S`$ and are null geodesics with respect to the $`𝕄^4`$ metric. For each $`pS`$, the tangents at $`p`$ to those $`\gamma _{\theta ,\varphi }^\pm \left(t\right)`$ that pass through $`p`$ can be shown to generate the null cone of $`T_pS`$. Hence $`S`$ is a photon surface in the sense of Definition 2.1. In terms of the double null coordinates $`u`$ $`:=t+r`$ (10) $`v`$ $`:=tr`$ (11) for $`r:=(x^2+y^2+z^2)^{1/2}`$, equation (8) assumes the simple form $$uv=a^2.$$ (12) In the future direction $`S`$ tends asymptotically to the null hypersurface $`\{v=0\}`$, whilst in the past direction $`S`$ tends asymptotically to the null hypersurface $`\{u=0\}`$. ###### Example 3 (De Sitter space). De Sitter space-time may be regarded as a single-sheeted hyperboloid in Minkowski $`5`$-space $`𝕄^5`$. By analogy with Examples 1 and 2, de Sitter space-time is thus realized as a photon surface in $`𝕄^5`$. ###### Example 4 (The Robertson-Walker models). Since all Robertson-Walker models are conformally flat and therefore locally conformally transformable to Minkowski space, the photon surfaces of any such model may thus be obtained, at least locally, by conformal transformations of Minkowski space. ###### Theorem 2.2. Let $`S`$ be a timelike hypersurface of $`(M,𝐠)`$. Let $`𝐧`$ be a unit normal field to $`S`$ and let $`h_{ab}`$ be the induced metric on $`S`$. Let $`\chi _{ab}`$ be the second fundamental form on $`S`$ and let $`\sigma _{ab}`$ be the trace-free part of $`\chi _{ab}`$. Then the following are equivalent: 1. $`S`$ is a photon surface; 2. $`\chi _{ab}k^ak^b=0`$ $``$ null $`𝐤T_pS`$ $`pS`$; 3. $`\sigma _{ab}=0`$; 4. every affine null geodesic of $`(S,𝐡)`$ is an affine null geodesic of $`(M,𝐠)`$. Proof. (i) $``$ (ii). Suppose $`S`$ is a photon surface. Let $`pS`$ and let $`𝐤T_pS`$ be null. There exists an affine null geodesic $`\gamma :(ϵ,ϵ)M`$ of $`(M,𝐠)`$ such that $`\dot{\gamma }\left(0\right)=𝐤,\|\gamma \|S`$. One has $$\chi _{ab}\dot{\gamma }^a\dot{\gamma }^b=n_{a;b}\dot{\gamma }^a\dot{\gamma }^b=\left(n_a\dot{\gamma }^a\right)_{;b}\dot{\gamma }^b=0$$ (13) along $`\gamma `$. At $`p`$ this gives $`\chi _{ab}k^ak^b=0`$. (ii) $``$ (iii). Let $`pS`$. By (ii) one has $`\sigma _{ab}k^ak^b=\chi _{ab}k^ak^b=0`$ $``$ null $`𝐤T_pS`$. Let $`\{𝐞_{(0)},𝐞_{(1)},𝐞_{(2)}\}`$ be an orthonormal basis for $`T_pS`$ with $`𝐞_{(0)}`$ timelike and $`𝐞_{(1)}`$, $`𝐞_{(2)}`$ spacelike. Any null $`𝐤T_pS`$, normalized such that $`g(𝐤,𝐞_{(0)})=1`$, has components $`k^a=(1,\mathrm{cos}\psi ,\mathrm{sin}\psi )`$ with respect to $`\{𝐞_{(0)},𝐞_{(1)},𝐞_{(2)}\}`$ for some $`\psi [0,2\pi )`$. A calculation gives $$\begin{array}{c}\sigma _{ab}k^ak^b=(\sigma _{00}+\frac{1}{2}\sigma _{11}+\frac{1}{2}\sigma _{22})+2\sigma _{01}\mathrm{cos}\psi +2\sigma _{02}\mathrm{sin}\psi \hfill \\ \hfill +\frac{1}{2}(\sigma _{11}\sigma _{22})\mathrm{cos}2\psi +\sigma _{12}\mathrm{sin}2\psi .\end{array}$$ (14) This must vanish for all $`\psi [0,2\pi )`$. One thus has $`\sigma _{01}=\sigma _{02}=\sigma _{12}=0`$ and $`\sigma _{00}=\sigma _{11}=\sigma _{22}`$. Since $`\sigma _{ab}`$ is trace-free one must also have $`\sigma _{00}=\sigma _{11}+\sigma _{22}`$. There follows $`\sigma _{ab}=0`$. (iii) $``$ (iv). For any curve in $`S`$ with null tangent $`𝐤`$ one has $$k^a\text{}_{\|\|b}k^b=h^a\text{}_ck^c\text{}_{;b}k^b=k^a\text{}_{;b}k^b+(\sigma _{bc}k^bk^c)n^a$$ (15) where $`\|\|`$ denotes covariant differentiation in $`S`$ with respect to $`𝐡`$. The second term on the right of $`(\text{15})`$ vanishes by hypothesis. If $`𝐤`$ is tangent to an affine null geodesic of $`(S,𝐡)`$ then the term on the left of $`(\text{15})`$ also vanishes and so $`𝐤`$ is tangent to an affine null geodesic of $`(M,𝐠)`$. (iv) $``$ (i). Let $`pS`$ and let $`𝐤T_pS`$ be null. Let $`\gamma :(ϵ,ϵ)S`$ be an affine null geodesic of $`(S,𝐡)`$ such that $`\dot{\gamma }(0)=𝐤`$. Then, by (iv), $`\gamma `$ is an affine null geodesic of $`(M,𝐠)`$ such that $`\dot{\gamma }(0)=𝐤`$, $`\|\gamma \|S`$. Condition (iii) of Theorem 2.2 is equivalent to a requirement that $`\chi _{ab}`$ is pure trace in the sense of $$\chi _{ab}=\frac{1}{3}\mathrm{\Theta }h_{ab},$$ (16) where $`\mathrm{\Theta }:=h^{cd}\chi _{cd}`$ is the expansion of the unit normal to $`S`$. For Example 1 one has $`\mathrm{\Theta }=2/a`$; for Example 2 one has $`\mathrm{\Theta }=3/a`$. (Note that, by a standard abuse of notation, $`h_{ab}`$ denotes both the induced metric on $`S`$ and the symmetric tensor field of rank $`(0,2)`$ along $`S`$ in $`M`$ which satisfies $`h_{ab}n^b=0`$ and pulls back to the induced metric on $`S`$.) It is clear from condition (iii) of Theorem 2.2 that a space-time must be specialized in some respect in order to admit any timelike photon surfaces in the sense of Definition 2.1. For this reason it is helpful to restrict attention to space-times which admit groups of symmetries. ###### Definition 2.3. Suppose $`(M,𝐠)`$ admits a group $`G`$ of isometries. A photon surface $`S`$ of $`(M,𝐠)`$ that is invariant under $`G`$, in the sense that each $`gG`$ maps $`S`$ onto itself, will be called a $`G`$-invariant photon surface. Clearly any $`G`$-invariant null hypersurface is a $`G`$-invariant photon surface. In particular, if $`G=`$ or $`G=𝕊`$, then any Killing horizon is a $`G`$-invariant photon surface. ## 3. Dynamic Spherical symmetry: General theory By definition, a general spherically symmetric space-time admits an $`SO(3)`$ isometry group for which the group orbits are spacelike $`2`$-spheres. The following result describes the evolution of the cross-sectional area of an $`SO(3)`$-invariant photon surface in a spherically symmetric space-time. ###### Theorem 3.1. Let $`(M,𝐠)`$ be a spherically symmetric space-time. Let $`S`$ be an $`SO(3)`$-invariant timelike hypersurface of $`(M,𝐠)`$ and let $`𝐗`$ be the $`SO(3)`$-invariant unit future-directed timelike tangent vector field along $`S`$ orthogonal to the $`SO(3)`$-invariant $`2`$-spheres in $`S`$. Let $`𝒯`$ be one such $`SO(3)`$-invariant $`2`$-sphere in $`S`$ and let $`𝒯_s`$ be the $`SO(3)`$-invariant $`2`$-sphere in $`S`$ at arc-length $`s`$ from $`𝒯`$ along the integral curves of $`𝐗`$. Then $`S`$ is a photon surface of $`(M,𝐠)`$ iff the area $`\text{}^{(2)}A_s`$ of $`𝒯_s`$ satisfies $$\frac{d^2}{ds^2}\text{}^{(2)}A_s=\frac{1}{4\text{}^{(2)}A_s}\left(\frac{d}{ds}\text{}^{(2)}A_s\right)^2+\text{}^{(2)}A_s\left(\frac{1}{3}\mathrm{\Theta }^2G_{ab}n^an^b\right)4\pi $$ (17) where $`n^a`$ is the unit normal to $`S`$, $`\mathrm{\Theta }`$ is the expansion of $`n^a`$ and $`G_{ab}:=R_{ab}\frac{1}{2}Rg_{ab}`$ is the Einstein tensor of $`(M,𝐠)`$. Proof. Let $`h_{ab}`$ be the induced Lorentzian $`3`$-metric on $`S`$ and, for each $`s`$, let $`\text{}^{(2)}h_{ab}`$ be the induced Riemannian $`2`$-metric on $`𝒯_s`$. The expansion of $`𝐗`$ in $`(S,𝐡)`$ is given by $`\text{}^{(2)}\mathrm{\Theta }=\text{}^{(2)}h^{ab}X_{a;b}`$ where the covariant derivative is that of $`(M,𝐠)`$. Since $`𝐗`$ is both shear-free and vorticity-free in $`(S,𝐡)`$, the Raychaudhuri equation for $`𝐗`$ in $`(S,𝐡)`$ assumes the form $$\frac{d}{ds}\text{}^{(2)}\mathrm{\Theta }=\frac{1}{2}(\text{}^{(2)}\mathrm{\Theta })^2\text{}^{(3)}R_{ab}X^aX^b$$ (18) where $`\text{}^{(3)}R_{ab}`$ is the Ricci tensor of $`(S,𝐡)`$. From first principles one has $`\text{}^{(2)}R`$ $`=\text{}^{(3)}R+2\text{}^{(3)}R_{ab}X^aX^b(\text{}^{(2)}\chi ^a\text{}_a)^2+\text{}^{(2)}\chi ^a\text{}_b\text{}^{(2)}\chi ^b\text{}_a`$ (19) $`\text{}^{(3)}R`$ $`=R2R_{ab}n^an^b+\left(\chi ^a\text{}_a\right)^2\chi ^a\text{}_b\chi ^b\text{}_a`$ (20) where $`\text{}^{(2)}\chi _{ab}`$ is the second fundamental form of each $`𝒯_s`$ in $`(S,𝐡)`$. Since $`X^a`$ is shear-free and vorticity free in $`(S,𝐡)`$ one has $`\text{}^{(2)}\chi _{ab}=\frac{1}{2}\text{}^{(2)}\mathrm{\Theta }\text{}^{(2)}h_{ab}`$. The second fundamental form of $`S`$ admits the canonical decomposition $`\chi _{ab}=\frac{1}{3}\mathrm{\Theta }h_{ab}+\sigma _{ab}`$. Equations (19) and (20) therefore give $`\text{}^{(2)}R`$ $`=\text{}^{(3)}R+2\text{}^{(3)}R_{ab}X^aX^b\frac{1}{2}\text{}^{(2)}\mathrm{\Theta }^2`$ (21) $`\text{}^{(3)}R`$ $`=R2R_{ab}n^an^b+\frac{2}{3}\mathrm{\Theta }^2\sigma ^a\text{}_b\sigma ^b\text{}_a`$ (22) which combine to yield $$2\text{}^{(3)}R_{ab}X^aX^b=\text{}^{(2)}R+2G_{ab}n^an^b\frac{2}{3}\mathrm{\Theta }^2+\frac{1}{2}(\text{}^{(2)}\mathrm{\Theta })^2+\sigma ^a\text{}_b\sigma ^b\text{}_a.$$ (23) One may now substitute for the second term on the right of (18) to obtain $$\frac{d}{ds}\text{}^{(2)}\mathrm{\Theta }=\frac{3}{4}(\text{}^{(2)}\mathrm{\Theta })^2+\frac{1}{3}\mathrm{\Theta }^2\frac{1}{2}\text{}^{(2)}RG_{ab}n^an^b\frac{1}{2}\sigma ^a\text{}_b\sigma ^b\text{}_a.$$ (24) From first principles one has $`\text{}^{(2)}\mathrm{\Theta }=\frac{d}{ds}\mathrm{ln}\text{}^{(2)}A`$, and the Gauss-Bonnet theorem gives $`\text{}^{(2)}R\text{}^{(2)}A=8\pi `$. Substituting for $`\text{}^{(2)}\mathrm{\Theta }`$ and $`\text{}^{(2)}R`$ in (24) one obtains $$\frac{d^2}{ds^2}\text{}^{(2)}A_s=\frac{1}{4\text{}^{(2)}A_s}\left(\frac{d}{ds}\text{}^{(2)}A_s\right)^2+\text{}^{(2)}A_s\left(\frac{1}{3}\mathrm{\Theta }^2G_{ab}n^an^b\frac{1}{2}\sigma ^a\text{}_b\sigma ^b\text{}_a\right)4\pi .$$ (25) This agrees with (17) iff $`\sigma ^a\text{}_b\sigma ^b\text{}_a=0`$. Construct, for the tangent bundle $`TS`$ of $`S`$, an orthonormal basis field of the form $`\{𝐗,𝐞_{(1)},𝐞_{(2)}\}`$, with $`𝐞_{(1)}`$ and $`𝐞_{(2)}`$ unit spacelike. With respect to this basis one has $$\sigma ^a\text{}_b\sigma ^b\text{}_a=(\sigma ^0\text{}_0)^2+(\sigma ^1\text{}_1)^2+(\sigma ^2\text{}_2)^2+2(\sigma ^1\text{}_2)^22(\sigma ^1\text{}_0)^22(\sigma ^2\text{}_0)^2.$$ (26) By spherical symmetry the vector field $`\sigma ^a\text{}_bX^b`$ must be proportional to $`X^a`$. Hence one has $`\sigma ^1\text{}_0=\sigma ^2\text{}_0=0`$. The vanishing of $`\sigma ^a\text{}_b\sigma ^b\text{}_a`$ is thus equivalent to the vanishing of $`\sigma _{ab}`$. One has $`\sigma _{ab}=0`$ iff $`S`$ is a photon surface. A spherically symmetric metric is locally expressible in the form $$g_{ab}=\left(\begin{array}{cccc}g_{00}& g_{01}& 0& 0\\ g_{10}& g_{11}& 0& 0\\ 0& 0& g_{\theta \theta }& 0\\ 0& 0& 0& g_{\theta \theta }\mathrm{sin}^2\theta \end{array}\right)$$ (27) with respect to coordinates $`(x^0,x^1,\theta ,\varphi )`$ adapted to the spherical symmetry, where $`g_{00}`$, $`g_{01}=g_{10}`$, $`g_{11}`$ and $`g_{\theta \theta }>0`$ depend only on $`x^0`$ and $`x^1`$. It is often convenient to introduce a radial coordinate $`r`$, depending only on $`x^0`$ and $`x^1`$, such that $`g_{\theta \theta }`$ is a function of $`r`$ only. One is free to specify $`g_{\theta \theta }`$ as a function of $`r`$ alone since to do so is, in effect, a definition of the coordinate $`r`$. This will be assumed to be done throughout this paper. The following result is useful in the locating of $`SO(3)`$-invariant photon surfaces in dynamic spherically symmetric space-times. ###### Lemma 3.2. Let $`(M,𝐠)`$ be a spherically symmetric space-time. Let $`S`$ be an $`SO(3)`$-invariant timelike hypersurface of $`(M,𝐠)`$ and let $`𝐗`$ be the $`SO(3)`$-invariant unit future-directed timelike tangent vector field along $`S`$ that is orthogonal to the $`SO(3)`$-invariant $`2`$-spheres in $`S`$. Then $`S`$ is a photon surface of $`(M,𝐠)`$ iff $$X^a\text{}_{;b}X^b=\frac{1}{2}(g^{\theta \theta }n^b_bg_{\theta \theta })n^a$$ (28) holds along $`S`$, where $`n^a`$ is the unit normal field to $`S`$ in $`(M,𝐠)`$. Proof. By spherical symmetry, and since $`𝐗`$ is unit timelike, the vector field $`_𝐗𝐗`$ must be proportional to $`𝐧`$. Hence it suffices to show that $`S`$ is a photon surface iff along $`S`$ one has $$𝐧_𝐗𝐗=\frac{1}{2}g^{\theta \theta }n^b_bg_{\theta \theta }$$ (29) or equivalently $$\chi _{ab}X^aX^b=\frac{1}{2}g^{\theta \theta }n^a_ag_{\theta \theta }.$$ (30) Construct for $`TS`$ a local orthonormal basis field of the form $`\{𝐗,𝐞_{(\theta )},𝐞_{(\varphi )}\}`$. With respect to this basis field the components of $`\chi ^a\text{}_b`$ form a diagonal matrix with $$\chi ^\theta \text{}_\theta =\chi ^\varphi \text{}_\varphi =n^\theta \text{}_{;\theta }=\frac{1}{2}g^{\theta \theta }n^a_ag_{\theta \theta }.$$ (31) Equation (30) is thus equivalent to $`\chi ^0\text{}_0=\chi ^\theta \text{}_\theta =\chi ^\varphi \text{}_\varphi `$ which is in turn equivalent to $`\sigma ^0\text{}_0=\sigma ^\theta \text{}_\theta =\sigma ^\varphi \text{}_\varphi `$. In view of the trace-free property of $`\sigma ^a\text{}_b`$, equation (30) is thus equivalent to $`\sigma ^a\text{}_b=0`$. From Theorem 2.2 one has that $`\sigma ^a\text{}_b=0`$ holds along $`S`$ iff $`S`$ is a photon surface. Let us continue to work with respect to the coordinate system $`\{x^0,x^1,\theta ,\varphi \}`$ employed in (27). Let $`x^a(s)`$ be an integral curve of the vector field $`𝐗`$ in Lemma 3.2. One has $$\frac{dx^a}{ds}=X^a$$ (32) and equation $`(\text{28})`$ becomes $$\frac{d^2x^a}{ds^2}+\mathrm{\Gamma }_{bc}^a\frac{dx^b}{ds}\frac{dx^c}{ds}=\frac{1}{2}(g^{\theta \theta }n^b_bg_{\theta \theta })n^a.$$ (33) Since $`𝐗=\frac{dx^0}{ds}(1,\frac{dx^1}{dx^0},0,0)`$ is unit timelike one has $$\left(\frac{ds}{dx^0}\right)^2=g_{ab}\frac{dx^a}{dx^0}\frac{dx^b}{dx^0}.$$ (34) The $`x^0`$ and $`x^1`$ components of equation $`(\text{33})`$ thus combine to give $$\frac{d^2x^1}{(dx^0)^2}=\frac{1}{2}g^{\theta \theta }n^a_ag_{\theta \theta }\left(n^1\frac{dx^1}{dx^0}n^0\right)g_{bc}\frac{dx^b}{dx^0}\frac{dx^c}{dx^0}+\left(\frac{dx^1}{dx^0}\mathrm{\Gamma }_{ab}^0\mathrm{\Gamma }_{ab}^1\right)\frac{dx^a}{dx^0}\frac{dx^b}{dx^0}$$ (35) where the components of $`𝐧`$ are given by $$n^0=\psi g_{1a}\frac{dx^a}{dx^0};n^1=\psi g_{0a}\frac{dx^a}{dx^0}$$ (36) for $$\psi :=\left(\mathrm{\Delta }\right)^{\frac{1}{2}}\left(g_{ab}\frac{dx^a}{dx^0}\frac{dx^b}{dx^0}\right)^{\frac{1}{2}}$$ (37) where $$\mathrm{\Delta }:=g_{00}g_{11}(g_{01})^2$$ (38) is the determinant of the time-space part of $`g_{ab}`$ in (27). Equation (35) is the coordinate equivalent of (28) and provides for the easy determination of $`SO(3)`$-invariant photon surfaces (see Example 10). ## 4. Static Spherical Symmetry: General Theory By definition, a spherically symmetric space-time is static if it admits an $`SO(3)\times `$ group of isometries such that the $``$ orbits are generated by a Killing field $`𝐊`$ which is both hypersurface orthogonal and orthogonal to the $`SO(3)`$ orbits. The present section will be concerned with $`SO(3)\times `$-invariant photon surfaces in static spherically symmetric space-times. Such surfaces may be termed photon spheres because, as will be seen, they are a natural generalization to general static spherically symmetric space-times of the Schwarzschild photon sphere concept. The term “photon sphere” will be regarded as applicable only in static spherically symmetric space-times. For clarity, the term “$`SO(3)\times `$-invariant photon surface” will usually be employed in preference to “photon sphere”. Although the space-times of Examples 1 and 2 are static and spherically symmetric, the photon surfaces in these space-times are not $`SO(3)\times `$-invariant and so are not photon spheres. Of the Robertson-Walker space-times of Example 4, only the Einstein cylinder is both spherically symmetric and static. None of the photon surfaces of the Einstein cylinder are $`SO(3)\times `$-invariant. Thus the Einstein cylinder has no photon spheres. One may characterize an $`SO(3)\times `$-invariant photon surface, or photon sphere, in a static spherically symmetric space-time by means of the following special case of Theorem 3.1. ###### Theorem 4.1. Let $`(M,𝐠)`$ be a static spherically symmetric space-time with Killing field $`𝐊`$ and let $`S`$ be an $`SO(3)\times `$-invariant timelike hypersurface of $`(M,𝐠)`$. Then $`S`$ is an $`SO(3)\times `$-invariant photon surface of $`(M,𝐠)`$ if there exists an $`SO(3)`$-invariant $`2`$-sphere $`𝒯S`$ satisfying $$A\left(\mathrm{\Theta }^23G_{ab}n^an^b\right)=12\pi $$ (39) where $`A`$ is the area of $`𝒯`$, $`n^a`$ is the unit normal to $`S`$ and $`\mathrm{\Theta }`$ is the trace of the second fundamental form of $`S`$. Conversely, if $`S`$ is an $`SO(3)\times `$-invariant photon surface of $`(M,𝐠)`$ then (39) holds for every $`SO(3)`$-invariant $`2`$-sphere $`𝒯S`$. Proof. Note that the unit future-directed timelike tangent field $`𝐗`$ along $`S`$ in Theorem 3.1 is proportional to the restriction to $`S`$ of the Killing field $`𝐊`$. Suppose first that there exists an $`SO(3)`$-invariant $`2`$-sphere $`𝒯S`$ such that (39) holds. The quantities $`A`$, $`\mathrm{\Theta }`$ and $`G_{ab}n^an^b`$ remain constant as $`𝒯`$ is mapped along the flow lines of the Killing field $`𝐊`$. So they also remain constant as they are mapped along the flow lines of $`𝐗`$. Hence (17) holds with the term on the left and the first term on the right both zero. Thus, by Theorem 3.1, $`S`$ is a photon surface of $`(M,𝐠)`$. By hypothesis $`S`$ is $`SO(3)\times `$-invariant. For the converse, suppose that $`S`$ is an $`SO(3)\times `$-invariant photon surface of $`(M,𝐠)`$. Then (17) holds for every $`SO(3)`$-invariant $`2`$-sphere $`𝒯_sS`$. Since $`𝐊`$ induces groups of local isometries, the area $`A_s`$ of $`𝒯_s`$ is independent of the parameter $`s`$. Hence the term on the left and the first term on the right of (17) both vanish and one obtains (39). ###### Corollary. If one has $`G_{ab}Y^aY^b0`$ $``$ vectors $`𝐘`$ and $`S`$ is an $`SO(3)\times `$-invariant timelike photon surface of $`(M,𝐠)`$, then for any $`SO(3)`$-invariant $`2`$-sphere $`𝒯S`$ one has $$A\mathrm{\Theta }^212\pi $$ (40) with equality holding iff $`G_{ab}n^an^b=0`$ along $`S`$. ∎ For Schwarzschild space-time where (see Example 5) the only timelike photon sphere is at $`r=3m`$, one has $`A=4\pi (3m)^2`$, $`\mathrm{\Theta }=1/(\sqrt{3}m)`$ and $`G_{ab}=0`$ which verifies (39) and (40) for this case. If the Einstein equations hold with a zero cosmological constant then, in the corollary to Theorem 4.1, the hypothesis $`G_{ab}Y^aY^b0`$ for all vectors $`𝐘`$ is equivalent to $`T_{ab}Y^aY^b0`$ for all vectors $`𝐘`$. This is a physically reasonable energy condition. In particular, for a perfect fluid with density $`\rho `$ and pressure $`p`$, it is equivalent to a condition that $`\rho `$ and $`p`$ are both non-negative. More generally, the condition holds for an energy tensor with a single timelike eigenvector (type I in the classification of Hawking & Ellis ) iff each energy tensor eigenvalue is non-negative. The characterization of timelike photon surfaces provided by Theorem 4.1 involves derivatives of the metric components up to second order. The following result (Theorem 4.2) provides an entirely different characterization of $`SO(3)\times `$-invariant photon surfaces in terms of derivatives of the metric components up to only first order. Let a general static spherically symmetric metric $`𝐠`$ be expressed in the form (27) with $`g_{\theta \theta }`$ a function of $`r`$ only. The Killing equation $`K_{(a;b)}=0`$ and the orthogonality of $`𝐊`$ to $`_\theta `$ and $`_\varphi `$ gives $`K^a_ag_{\theta \theta }=0`$ and hence $`_𝐊r=0`$ where $`r`$ is to be regarded as a scalar field on $`M`$. Since $`r`$ is independent of $`\theta `$ and $`\varphi `$ it follows that any $`SO(3)\times `$-invariant hypersurface $`S`$ of $`(M,𝐠)`$ must be of the form $`\{r=\text{const.}\}`$. If $`S`$ is also a timelike hypersurface then $`r^{;a}`$ is a spacelike vector field along $`S`$ and $`𝐊`$ is a timelike vector field in a neighbourhood of $`S`$. If the coordinates $`x^0`$ and $`x^1`$ implicit in (27) are chosen such that $`𝐊=_{x^0}`$ then the Killing equation gives that all the metric components $`g_{ab}`$ in (27) are independent of $`x^0`$. Since $`r`$ is constant along the integral curves of $`𝐊`$, the coordinate $`x^1`$ must be a function of $`r`$ only. A natural choice is $`x^1=r`$. One may redefine $`x^0`$ according to $`x^0x^0(g_{01}/g_{00})𝑑x^1`$. This diagonalizes the time-space part of $`g_{ab}`$ and leaves the components of $`g_{ab}`$ independent of $`x^0`$. Furthermore the curves $`\{x^1,\theta ,\varphi =\text{const.}\}`$ are unchanged except that they are re-parametrized. The vector field $`_{x^0}`$ then becomes a conformal Killing field. Define the tensor field $$ϵ^{ab}:=(\mathrm{\Delta })^{1/2}\left(\begin{array}{cccc}0& 1& 0& 0\\ 1& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\end{array}\right)$$ (41) on $`M`$, where the components are given with respect to the coordinate basis employed in (27) and $`\mathrm{\Delta }`$ is the determinant of the time-space part of $`𝐠`$ in (27), as in (38). ###### Theorem 4.2. Let $`(M,𝐠)`$ be a static spherically symmetric space-time with $`𝐠`$ of the form (27), with $`g_{\theta \theta }`$ a function of the coordinate $`r`$ only. Let $`S`$ be an $`SO(3)\times `$-invariant timelike hypersurface of $`(M,𝐠)`$ and suppose that $`r`$ is nowhere-zero along $`S`$. Then $`S`$ is an $`SO(3)\times `$-invariant photon surface of $`(M,𝐠)`$ iff $$2g_{\theta \theta }ϵ^{ab}ϵ^{cd}r_{;ac}r_{;b}r_{;d}+r^{;a}r_{;a}r^{;c}_cg_{\theta \theta }=0$$ (42) holds along $`S`$. Proof. Since $`(M,𝐠)`$ is both spherically symmetric and static, the surface $`S`$ is of the form $`\{r=\text{const.}\}`$. The unit spacelike normal to $`S`$ is therefore given by $$n^a=\eta r^{;a}$$ (43) for $$\eta :=(r^{;a}r_{;a})^{1/2}.$$ (44) The second fundamental form of $`S`$ is given by $$\chi _{ab}:=\eta h_a\text{}^ch_b\text{}^dr_{;cd}.$$ (45) The vector fields $`𝐗^a`$ $`:=(\mathrm{\Delta }g^{bc}r_{;b}r_{;c})^{1/2}(r_{;1},r_{;0},0,0)`$ (46) $`𝐞_{(\theta )}`$ $`:=(g_{\theta \theta })^{1/2}_\theta `$ (47) $`𝐞_{(\varphi )}`$ $`:=(g_{\theta \theta }\mathrm{sin}^2\theta )^{1/2}_\varphi `$ (48) form an orthonormal frame field along $`S`$, with $`𝐞_{(\theta )}`$ and $`𝐞_{(\varphi )}`$ unit spacelike and $`𝐗`$ unit timelike. One has $$\chi _{ab}X^a𝐞_{(\theta )}^b=\chi _{ab}X^a𝐞_{(\varphi )}^b=\chi _{ab}𝐞_{(\theta )}^a𝐞_{(\varphi )}^b=0$$ (49) $$\chi _{ab}𝐞_{(\theta )}^a𝐞_{(\theta )}^b=\chi _{ab}𝐞_{(\varphi )}^a𝐞_{(\varphi )}^b=\frac{\eta r^{;a}_ag_{\theta \theta }}{2g_{\theta \theta }}$$ (50) $$\chi _{ab}X^aX^b=\eta ^3ϵ^{ab}ϵ^{cd}r_{;ac}r_{;b}r_{;d}.$$ (51) Condition (iii) of Theorem 2.2 holds iff $`\chi _{ab}`$ is proportional to $`h_{ab}`$ and hence iff $$\chi _{ab}X^aX^b=\chi _{ab}𝐞_{(\theta )}^a𝐞_{(\theta )}^b=\chi _{ab}𝐞_{(\varphi )}^a𝐞_{(\varphi )}^b.$$ (52) This is equivalent to $$\eta ^2ϵ^{ab}ϵ^{cd}r_{;ac}r_{;b}r_{;d}=\frac{r^{;a}_ag_{\theta \theta }}{2g_{\theta \theta }}$$ (53) which is in turn equivalent to (42). Equation (42) can have solutions such that $`\{r=\text{const.}\}`$ is a spacelike hypersurface and therefore not a photon surface (see e.g. Example 7). It is therefore always necessary in the use of Theorem 4.2 to check that the hypersurface $`\{r=\text{const.}\}`$ is in fact timelike or null. Note that in a region of space-time where the Killing field $`𝐊`$ is spacelike, for example behind the event horizon of Schwarzschild space-time, the hypersurfaces $`\{r=\text{const.}\}`$ are necessarily spacelike and so cannot be photon spheres. Case 1. ($`x^1:=r`$, components of $`g_{ab}`$ independent of $`x^0`$.) As discussed previously, for a static spherically symmetric metric it is possible to choose $`x^1:=r`$, with $`g_{\theta \theta }`$ depending only upon $`r`$ and with all the components of $`g_{ab}`$ independent of $`x^0`$. In this case (42) reduces to $$g_{00}_rg_{\theta \theta }=g_{\theta \theta }_rg_{00}.$$ (54) This agrees with an equation obtained by Virbhadra & Ellis on the basis of a different definition of a photon sphere. Note that even though the components $`g_{rr},g_{0r}=g_{r0}`$ do not appear in (54), they are not assumed to vanish. A particular sub-case of interest is that of time-space coordinates, with $`t:=x^0`$ timelike in the sense of $`g^{tt}<0`$. Another sub-case of interest is that of single null (radiation) coordinates, with $`u:=x^0`$ null in the sense of $`g^{uu}=0`$. Case 2. (Double null coordinates $`u,v`$.) Let $`x^0:=u`$, $`x^1:=v`$ be double null coordinates in the sense of $`g^{uu}=g^{vv}=0`$. The radial coordinate $`r`$ is to be regarded as a function of $`u`$ and $`v`$. Then (42) assumes the form $$g_{\theta \theta }\{r_{;uu}(r_{;v})^22r_{;uv}r_{;u}r_{;v}+r_{;vv}(r_{;u})^2\}+2(r_{;u}r_{;v})^2_rg_{\theta \theta }=0.$$ (55) The metric components $`g_{uv}=g_{vu}`$ enter here through the covariant derivatives of $`r`$. Equations (42), (54) and (55) may be referred to as photon sphere equations since they give the location of timelike photon spheres in static spherically symmetric space-times. In order to facilitate further progress, a general static spherically symmetric metric will be written in such a form as to cast the Einstein tensor in a particularly simple and convenient form. One has from previous remarks that a general static spherically symmetric metric is locally expressible in the form $$𝐠=g_{tt}dt^2+g_{rr}dr^2+r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)$$ (56) where $`g_{tt}`$ and $`g_{rr}`$ are functions of $`r`$ only. Let $`m(r)`$ $`:=\frac{1}{2}r\left(1{\displaystyle \frac{1}{g_{rr}}}\right)`$ (57) $`\mu (r)`$ $`:=\mathrm{ln}(g_{tt}g_{rr}).`$ (58) Then the metric assumes the form $$𝐠=\left(1\frac{2m(r)}{r}\right)\mathrm{e}^{\mu (r)}dt^2+\left(1\frac{2m(r)}{r}\right)^1dr^2+r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)$$ (59) and the Einstein tensor is given by $$G^a\text{}_b=8\pi \left(\begin{array}{cccc}\rho (r)& 0& 0& 0\\ 0& p_1(r)& 0& 0\\ 0& 0& p_2(r)& 0\\ 0& 0& 0& p_2(r)\end{array}\right)$$ (60) for $`8\pi \rho (r)`$ $`:={\displaystyle \frac{2m^{}(r)}{r^2}}`$ (61) $`8\pi p_1(r)`$ $`:={\displaystyle \frac{1}{r^2}}\left\{(r2m(r))\mu ^{}(r)2m^{}(r)\right\}`$ (62) $`8\pi p_2(r)`$ $`:={\displaystyle \frac{1}{4r^2}}\{(2(r+m(r)3rm^{}(r))\mu ^{}(r)+r(r2m(r))(\mu ^{}(r))^2`$ $`4rm^{\prime \prime }(r)+2r(r2m(r))\mu ^{\prime \prime }(r)\}.`$ (63) where a prime denotes differentiation with respect to $`r`$. For the right sides of (61), (62) and (63) to be defined at some radius $`\widehat{r}>0`$ one evidently needs $`m(r)`$ and $`\mu (r)`$ to be twice differentiable at $`r=\widehat{r}`$. Equations (61) and (62) combine to give $$\mu ^{}(r)=\frac{8\pi r(\rho (r)+p_1(r))}{\left(1\frac{2m(r)}{r}\right)}$$ (64) whereby one may rewrite (63) in the more convenient form $$\frac{2}{r}(p_2(r)p_1(r))=p_1^{}(r)+\frac{\left(\frac{m(r)}{r^2}+4\pi rp_1(r)\right)\left(\rho (r)+p_1(r)\right)}{\left(1\frac{2m(r)}{r}\right)}.$$ (65) In the perfect fluid case $`p_1(r)=p_2(r)=:p(r)`$ equation (65) reduces to the Tolman-Oppenheimer-Volkoff equation $$p^{}(r)=\frac{\left(\frac{m(r)}{r^2}+4\pi rp(r)\right)\left(\rho (r)+p(r)\right)}{\left(1\frac{2m(r)}{r}\right)}.$$ (66) By means of equations (61) and (64), the photon sphere equation (54) becomes $$1\frac{3m(r)}{r}4\pi r^2p_1(r)=0.$$ (67) For $`r`$ such that $`2m(r)<r`$, and hence such that the hypersurface $`\{r=\text{const.}\}`$ is timelike, equation (67) gives the location of the $`SO(3)\times `$-invariant timelike photon surfaces for the metric (59). Equation (67) is the basis for the following result which shows that, subject to a suitable energy condition, any black hole in a static spherically symmetric space-time must be surrounded by an $`SO(3)\times `$-invariant photon surface. For the purpose of this and subsequent results, a function $`f:I`$ on an interval $`I`$ will be said to be piecewise $`C^r`$ if $`I`$ is the disjoint union of a locally finite collection of intervals $`I_i`$ such that $`f\|I_i`$ is $`C^r`$. Each interval $`I_i`$ may be open, closed or half-open. ###### Theorem 4.3. Suppose the metric $`𝐠`$ has the form (59) for $`r_0<r<\mathrm{}`$, for some $`r_0>0`$, with $`m(r)`$ and $`\mu (r)`$ both $`C^0`$, piecewise $`C^2`$ functions of $`r(r_0,\mathrm{})`$. Suppose the following hold: 1. $`\rho (r)`$ and $`p_1(r)`$ are bounded functions of $`r(r_0,\mathrm{})`$; 2. $`2m(r)<r`$ $`r(r_0,\mathrm{})`$; 3. $`\rho (r)0`$, $`p_1(r)0`$ $`r(r_0,\mathrm{})`$; 4. $`lim_r\mathrm{}4\pi r^2p_1(r)=lim_r\mathrm{}4\pi r^2\rho (r)=0`$; 5. for each value of $`t`$ the $`2`$-surfaces $`𝒯_{t,r}:=\{t=\text{const.}\}\{r=\text{const.}\}`$, $`r_0<r<\mathrm{}`$, are such that $`𝒯_t:=lim_{rr_0}𝒯_{t,r}`$ exists as an embedded spacelike $`2`$-sphere in $`(M,𝐠)`$ and is marginally outer trapped. Then $`(M,𝐠)`$ admits an $`SO(3)\times `$-invariant timelike photon surface of the form $`\{r=r_1\}`$ for some $`r_1(r_0,\mathrm{})`$. Proof. Fix $`t`$ and let $`𝐤`$ be the outward future-directed null normal field along each $`𝒯_{t,r}`$, $`r_0<r<\mathrm{}`$, normalized such that $`𝐠(𝐤,𝐧)=1`$, where $`𝐧=\left(1\frac{2m(r)}{r}\right)^{1/2}_r`$ is the outward radial unit tangent to $`\{t=\text{const.}\}`$. Since $`𝐤`$ is parallelly propagated along each of the geodesic integral curves of $`𝐧`$, one has that $`lim_{rr_0}𝐤`$ is a well-defined, nowhere-zero null vector field along $`𝒯_t:=lim_{rr_0}𝒯_{t,r}`$. For $`r(r_0,\mathrm{})`$ the vector field $`𝐤`$ has the form $`k^a=(k^t,a(r)k^t,0,0)`$ for $$a(r):=(g_{tt}/g_{rr})^{1/2}=\left(1\frac{2m(r)}{r}\right)\mathrm{e}^{\mu (r)/2}.$$ (68) The expansion of $`𝐤`$ is given by $$\mathrm{\Theta }_{\text{out}}=\text{}^{(2)}h^b\text{}_ak^a\text{}_{;b}=\frac{2a(r)}{r}k^t.$$ (69) The condition that $`𝒯_t`$ is marginally outer trapped therefore implies $$0=\underset{rr_0}{lim}a(r)=\underset{rr_0}{lim}\left(1\frac{2m(r)}{r}\right)\mathrm{e}^{\mu (r)/2}.$$ (70) The non-negativity of $`\rho (r)`$ and $`p_1(r)`$ gives, by means of (61) and (64), that $`m(r)`$ and $`\mu (r)`$ are non-decreasing functions of $`r(r_0,\mathrm{})`$. Thus (70) holds iff at least one of $$\underset{rr_0}{lim}\left(1\frac{2m(r)}{r}\right)=0;\underset{rr_0}{lim}\mu (r)=\mathrm{}$$ (71) holds. Suppose the first of (71) fails. Then the second must hold and one has $$\underset{rr_0}{lim}\left(1\frac{2m(r)}{r}\right)^1<\mathrm{}.$$ (72) From the boundedness of $`\rho (r)`$ and $`p_1(r)`$ on $`r(0,\mathrm{})`$ one has, by means of (64) and (72), that $`lim\; sup_{rr_0}\mu ^{}(r)`$ is finite. This is incompatible with the second of (71). Hence the first of (71) must hold. Let $`f:(r_0,\mathrm{})`$ be the left side of (67). By the non-negativity of $`p_1(r)`$ and the first of (71) one has $`lim_{rr_0}f(r)\frac{1}{2}`$. By condition (4), equation (61) and l’Hôpital’s rule one has $`lim_r\mathrm{}m(r)/r=lim_r\mathrm{}r^2p_1(r)=0`$ and hence $`lim_r\mathrm{}f(r)=1`$. Hence there exists some $`r_1(r_0,\mathrm{})`$ such that $`f(r_1)=0`$. The hypersurface $`\{r=r_1\}`$ is an $`SO(3)\times `$-invariant photon surface of $`(M,𝐠)`$. Condition (3) of Theorem 4.3 may be expressed more succinctly as $`G_{ab}Y^aY^b0`$ $``$ vectors $`𝐘`$. With regard to condition (5) of Theorem 4.3, to have required $`𝒯_t`$ to be contained in the hypersurface $`\{t=\text{const.}\}`$ would have been too strong since, for Schwarzschild space-time, no spacelike hypersurface of the form $`\{t=\text{const.}\}`$ in the exterior region contains a marginally outer trapped $`2`$-surface. Theorem 4.3 may be interpreted to the effect that, subject to the energy conditions expressed in condition (3), any static spherically symmetric black hole must be surrounded by an $`SO(3)\times `$-invariant timelike photon surface. The following result may then be regarded as a partial converse in that it shows, subject to a suitable energy condition, that if there exists an $`SO(3)\times `$-invariant timelike photon surface then there must be a naked singularity or a black hole, or more than a certain amount of matter. ###### Proposition 4.4. Suppose the metric $`𝐠`$ has the form (59) for $`0<r<\mathrm{}`$, with $`m(r)`$ and $`\mu (r)`$ both $`C^0`$, piecewise $`C^2`$ functions of $`r(0,\mathrm{})`$. If the following all hold: 1. $`\rho (r)`$ is a non-increasing, bounded function of $`r(0,\mathrm{})`$; 2. $`lim_{r0}m(r)=0`$; 3. $`4m(r)<r`$ $`r(0,\mathrm{})`$; 4. $`p_1(r)\rho (r)/3`$ $`r(0,\mathrm{})`$, then $`(M,𝐠)`$ can contain no $`SO(3)\times `$-invariant timelike photon surfaces. Proof. By conditions (1) and (2) with equation (61) one has $`m(r)(4\pi /3)r^3\rho (r)`$ $`r>0`$. By condition (4) one therefore has $`4\pi r^2p_1(r)(4\pi /3)r^2\rho (r)m(r)/r`$ $`r>0`$. The left side of (67) is thus bounded from below by $`14m(r)/r`$ $`r>0`$. This is positive by condition (3). The left side of (67) is therefore non-vanishing for all $`r>0`$. Note that this result is valid even for negative $`p_1(r)`$ and $`\rho (r)`$. Condition (2) of Proposition 4.4 prohibits any curvature singularity at $`r=0`$. Condition (3) may be interpreted as a requirement that there is no black hole and less than a certain amount of matter. The result shows that one of these two conditions must fail if there is an $`SO(3)\times `$-invariant timelike photon surface and conditions (1) and (4) both hold. When the matter is a perfect fluid it is possible to improve condition (3) of Proposition 4.4 to condition (3) of the following result. ###### Theorem 4.5. Suppose the metric $`𝐠`$ has the form (59) for $`0<r<\mathrm{}`$, with $`m(r)`$ and $`\mu (r)`$ both $`C^1`$, piecewise $`C^2`$ functions of $`r(0,\mathrm{})`$. If the following all hold: 1. the matter is a perfect fluid with pressure $`p(r)`$ and density $`\rho (r)`$; 2. $`lim_r\mathrm{}4\pi r^2p(r)=lim_r\mathrm{}4\pi r^2\rho (r)=0`$; 3. $`(24/7)m(r)<r`$ $`r(0,\mathrm{})`$; 4. $`p(r)\rho (r)/3`$ $`r(0,\mathrm{})`$, then $`(M,𝐠)`$ can contain no $`SO(3)\times `$-invariant timelike photon surfaces. Proof. Let $`f:(0,\mathrm{})`$ be the left side of equation (67). Since $`m(r)`$ and $`\mu (r)`$ are $`C^1`$, piecewise $`C^2`$ functions of $`r(0,\mathrm{})`$, one has by equation (62) that $`f(r)`$ is a $`C^0`$, piecewise $`C^1`$ function of $`r(0,\mathrm{})`$. The function $`f^{}(r)`$ is then a piecewise $`C^0`$ function of $`r(0,\mathrm{})`$ which, by means of of the Tolman-Oppenheimer-Volkoff equation (66), is given by $$rf^{}(r)=3\frac{m(r)}{r}12\pi r^2\rho (r)8\pi r^2p(r)+\frac{\left(\frac{m(r)}{r}+4\pi r^2p(r)\right)}{\left(12\frac{m(r)}{r}\right)}4\pi r^2(\rho (r)+p(r)).$$ (73) For $`r=r_1(0,\mathrm{})`$ such that $$0=f(r_1)=13\frac{m(r_1)}{r_1}4\pi r_1^2p(r_1)$$ (74) equation (73) reduces to $$r_1f^{}(r_1)=18\pi r_1^2(\rho (r_1)+p(r_1)).$$ (75) From condition (3) and equation (74) one has $`4\pi r_1^2p(r_1)>1/8`$, whence by condition (4) one has $`4\pi r_1^2\rho (r_1)>3/8`$. Thus (75) gives $`f^{}(r_1)<0`$. By condition (2), equation (61) and l’Hôpital’s rule one has $`lim_r\mathrm{}m(r)/r=lim_r\mathrm{}4\pi r^2p(r)=0`$ and hence $`lim_r\mathrm{}f(r)=1`$. Since it has been established that $`f^{}(r)`$ is negative for all $`r(0,\mathrm{})`$ such that $`f(r)=0`$, one must therefore have $`f(r)>0`$ for all $`r(0,\mathrm{})`$. Hence the space-time can contain no $`SO(3)\times `$-invariant timelike photon surfaces. Note that, as for Proposition 4.4, Theorem 4.5 is valid even for negative pressure and density. One would like to remove the insufficient matter parts of condition (3) of Proposition 4.4 and condition (3) of Theorem 4.5, in other words to weaken these to a no-black-hole condition $`2m(r)<r`$ $`r>0`$. But no result to this effect is forthcoming. On the other hand no counterexample is known. To conclude this section it will be shown that the physical significance of the photon sphere in Schwarzschild space-time, as discussed in the Introduction, carries over to the general static spherically symmetric case. Suppose the metric has the form (59) for $`r_0<r<\mathrm{}`$ and is asymptotically flat in the limit $`r\mathrm{}`$. Assume $`p_1(r)0`$, $`m(r)0`$ $`r>r_0`$. The matter need not be a perfect fluid. Denote the left side of (67) by $`f(r)`$. The condition of asymptotic flatness gives $`lim_r\mathrm{}f(r)=1`$ so, if there are any $`SO(3)\times `$-invariant timelike photon surfaces, there will be an outermost such surface $`S`$. For simplicity assume $`f^{}(r)0`$ at $`S`$. Let $`_{\text{ext}}`$ be the connected component of $`\{qM:f(q)>0\}`$ that has $`S`$ as its inner boundary and extends to $`r=\mathrm{}`$. Let $`_{\text{int}}`$ be the connected component of $`\{qM:f(q)<0\}`$ that has $`S`$ as its outer boundary. Consider first the case of a future endless affine null geodesic $`\gamma (\lambda )`$. The null geodesic equations for the metric (59) give $$\frac{d^2r}{d\lambda ^2}=rf(r)\left\{\left(\frac{d\theta }{d\lambda }\right)^2+\mathrm{sin}^2\theta \left(\frac{d\varphi }{d\lambda }\right)^2\right\}\frac{\mu ^{}(r)}{2}\left(\frac{dr}{d\lambda }\right)^2.$$ (76) At any point $`p\|\gamma \|_{\text{ext}}`$ such that $`dr/d\lambda =0`$ one has $`d^2r/d\lambda ^2>0`$. At any point $`p\|\gamma \|_{\text{int}}`$ such that $`dr/d\lambda =0`$ one has $`d^2r/d\lambda ^2<0`$. Thus if $`\gamma `$ starts outside $`S`$ (i.e. in $`_{\text{ext}}`$) and is initially directed outwards, in the sense that $`dr/d\lambda `$ is initially positive, then $`\gamma `$ will continue outwards. If $`\gamma `$ starts in $`_{\text{int}}`$ and is initially directed inwards, in the sense that $`dr/d\lambda `$ is initially negative, then $`\gamma `$ will continue inwards until it falls either into a singularity or through an $`SO(3)\times `$-invariant photon surface other than $`S`$. Consider now a unit speed timelike geodesic $`\xi (s)`$. The timelike geodesic equations give $$\frac{d^2r}{ds^2}=\frac{m(r)}{r^2}4\pi rp_1(r)+rf(r)\left\{\left(\frac{d\theta }{ds}\right)^2+\mathrm{sin}^2\theta \left(\frac{d\varphi }{ds}\right)^2\right\}\frac{\mu ^{}(r)}{2}\left(\frac{dr}{ds}\right)^2.$$ (77) For any point $`p_{\text{ext}}`$ one can arrange for the first three terms on the right of (77) to cancel and so obtain a unit speed timelike geodesic $`\xi (s)`$ through $`p_{\text{ext}}`$ at constant $`r`$. For $`p_{\text{int}}`$ the first three terms on the right of (77) are evidently negative. For $`p\|\xi \|_{\text{int}}`$ such that $`dr/ds=0`$ one has $`d^2r/ds^2<0`$. Thus if $`\xi `$ starts in $`_{\text{int}}`$ and is initially directed inwards, in the sense that $`dr/ds`$ is initially negative, then $`\xi `$ will continue inwards until it falls either into a singularity or through an $`SO(3)\times `$-invariant photon surface other than $`S`$. ## 5. Spherical symmetry: Examples The following are some examples of $`SO(3)\times `$-invariant and $`SO(3)`$-invariant photon surfaces in familiar space-times. ###### Example 5 (Schwarzschild space-time). The metric of Schwarzschild space-time in single null (radiation) coordinates has the form $$𝐠=\left(1\frac{2m}{r}\right)du^2+2dudr+r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2\right).$$ (78) In this case equation $`(\text{54})`$ reduces to $`r=3m`$. The timelike hypersurface $`\{r=3m\}`$ is thus an $`SO(3)\times `$-invariant photon surface, or photon sphere, as expected. There are no other $`SO(3)\times `$-invariant timelike photon surfaces. For a non-zero cosmological constant $`\mathrm{\Lambda }`$, the Schwarzschild metric (78) generalizes to the Schwarzschild-de Sitter metric $$𝐠=\left(1\frac{2m}{r}+\frac{\mathrm{\Lambda }r^2}{3}\right)du^2+2dudr+r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2\right).$$ (79) One finds that equation (54) reduces to $`r=3m`$, independent of the value of $`\mathrm{\Lambda }`$. This is surprising since Schwarzschild and Schwarzschild-de Sitter space-times are not conformally related. ###### Example 6 (Schwarzschild interior solution). The Schwarzschild interior solution describes a spherically symmetric distribution of perfect fluid of radius $`R`$, bounded pressure $`p`$ and constant density $`\rho _0>0`$. The solution is to have a metric of the form (59) and is to be matched at $`r=R`$ to a Schwarzschild vacuum solution in such a way that the pressure is a continuous function of $`r`$. One thus has $`p_1(r)=p_2(r)=:p(r)`$ for all $`r:0r<\mathrm{}`$, $`\rho (r)=\rho _0`$ for all $`r:0rR`$, $`\rho (r)=0`$ for all $`r:R<r<\mathrm{}`$ and $`p(r)=0`$ for all $`r:Rr<\mathrm{}`$. The pressure $`p(r)`$ for $`r:0rR`$ is to be obtained by an integration of the Tolman-Oppenheimer-Volkoff equation (66) subject to the boundary condition $`p(R)=0`$. This yields $`m(r)`$ $`={\displaystyle \frac{4\pi }{3}}\rho _0r^3,`$ $`0rR,`$ (80) $`\mathrm{e}^{\mu (r)}`$ $`={\displaystyle \frac{(3u(r))^2}{4u^2(r)}},`$ $`0rR,`$ (81) $`p(r)`$ $`={\displaystyle \frac{u(r)1}{3u(r)}}\rho _0,`$ $`0rR,`$ (82) for $`u(r)`$ $`:=\left({\displaystyle \frac{38\pi \rho _0r^2}{38\pi \rho _0R^2}}\right)^{\frac{1}{2}},`$ $`0rR.`$ (83) The spherically symmetric system described by the Schwarzschild interior solution can exist in a state of stable equilibrium iff $`m(R)/R<4/9`$ (see in Stephani ). This condition is equivalent to $`8\pi \rho _0R^2<8/3`$, which implies $`p(r)0`$ $`r0`$, and implies that the absence of a black hole is a general feature of the Schwarzschild interior solution. The left side of equation (67) now assumes the form $$1\frac{3m(r)}{r}4\pi r^2p(r)=\{\begin{array}{cc}1\frac{8\pi \rho _0r^2}{3u\left(r\right)}\hfill & :0<rR,\hfill \\ 1\frac{3m\left(R\right)}{r}\hfill & :rR.\hfill \end{array}$$ (84) For $`m(R)/R<1/3`$ one has that (84) is positive for all $`r>0`$, so there are no timelike photon spheres. For $`m(R)/R=1/3`$ there is a single timelike photon sphere which lies at the boundary $`r=R`$ of the matter. For $`1/3<m(R)/R<4/9`$ there is one timelike photon sphere outside the matter at $`r=3m(R)>R`$ and one timelike photon sphere inside the matter at $$r=\left(\frac{13\pi \rho _0R^2}{\pi \rho _0(38\pi \rho _0R^2)}\right)^{1/2}=\frac{2R}{3}\left(\frac{1\frac{9m\left(R\right)}{4R}}{\left(1\frac{2m\left(R\right)}{R}\right)\frac{m\left(R\right)}{R}}\right)^{1/2}<R.$$ (85) For fixed $`R`$ the radius of the outer photon sphere is a strictly increasing function of $`m(R)/R`$ whilst the radius of the inner photon sphere is a strictly decreasing function of $`m(R)/R`$. Thus a Schwarzschild interior solution matched to a Schwarzschild vacuum exterior solution contains no black hole, and contains one timelike photon sphere iff $`1/3=m(R)/R`$ and two timelike photon spheres iff $`m(R)/R`$ lies in the range $`1/3<m(R)/R<4/9`$. However for such values of $`m(R)/R`$ the space-time is unphysical in that the pressure at the center is $`p(0)\rho _0/\sqrt{3}>\rho _0/3`$. Therefore under the reasonable energy condition $`p(r)\rho _0/3`$ $`r:0r<\mathrm{}`$ there are no photon spheres in this example. The energy condition $`0p(r)\rho (r)/3`$ $`r[0,\mathrm{})`$ is in fact satisfied iff $`m(R)/R`$ lies in the range $`0m(R)/R5/18`$. The corresponding space-times (in view of $`5/18<7/24`$) satisfy all of conditions (1) to (4) of Theorem 4.5, except the smoothness of $`\rho (r)`$ at $`r=R`$. Thus considering a smoothed family of solutions approximating the solution above and having it as a strict limit, but each with smooth $`\rho (r)`$ at $`r=R`$, we can apply Theorem 4.5 to show no photon spheres exist in all the cases discussed in this section where this energy condition is satisfied. On the other hand Proposition 4.4 is directly applicable without smoothing, but only applies to the subset of these cases with $`0m(R)/R<1/4`$. ###### Example 7 (Reissner-Nordström space-time). The general static spherically symmetric solution to the Einstein-Maxwell equations is the Reissner-Nordström solution comprising a metric $`𝐠`$ and an electromagnetic field $`F_{ab}`$ given by $`𝐠`$ $`=\left(1{\displaystyle \frac{2m}{r}}+{\displaystyle \frac{e^2}{r^2}}\right)du^2+2dudr+r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2\right)`$ (86) $`F_{tr}`$ $`=F_{rt}={\displaystyle \frac{e}{r^2}},\text{all other components vanishing,}`$ (87) where $`m`$ is the ADM mass and $`e`$ is the electric charge. We assume $`m>0`$. There is an event horizon at $`r=r_+:=m+\sqrt{m^2e^2}`$ and a Cauchy horizon at $`r=r_{}:=m\sqrt{m^2e^2}`$. The event horizon exists for $`0(e/m)^21`$, and the Cauchy horizon exists for $`0<(e/m)^21`$. For $`(e/m)^2=1`$ they both lie at $`r=m`$. There can be no timelike photon spheres between the event horizon and the Cauchy horizon because the Killing field $`𝐊`$ is spacelike there. Outside the event horizon the Killing field $`𝐊`$ is timelike, so every hypersurface of the form $`\{r=\text{const.}\}`$ which lies outside the event horizon and satisfies the photon sphere equation (54) is necessarily a timelike photon sphere. Equation (54) assumes the form $$r^23mr+2e^2=0$$ (88) which has solutions $`r_{\text{ps}}^\pm `$ given by $$r_{\text{ps}}^\pm /m=\frac{3\pm \sqrt{98(e/m)^2}}{2}.$$ (89) The hypersurface $`S^+:=\{r=r_{\text{ps}}^+\}`$ exists for $`0(e/m)^29/8`$ and lies outside the event horizon and is therefore a timelike photon sphere. The hypersurface $`S^{}:=\{r=r_{\text{ps}}^{}\}`$ exists for $`0<(e/m)^29/8`$ but lies outside the event horizon only for $`1<(e/m)^29/8`$, and so is a timelike photon sphere only then. The hypersurfaces $`S^+`$ and $`S^{}`$ coincide for $`(e/m)^2=9/8`$. The Cauchy horizon and event horizon are always null photon spheres. For $`0(e/m)^21`$ the curvature singularity at $`r=0`$ is locally naked but hidden behind an event horizon which lies strictly inside the only timelike photon sphere. For $`1<(e/m)^2`$ the singularity at $`r=0`$ is globally naked and is surrounded by two timelike photon spheres in the case $`1<(e/m)^2<9/8`$, one timelike photon sphere in the case $`(e/m)^2=9/8`$ and by no photon spheres, either timelike or null, in the case $`(e/m)^2>9/8`$ (see Fig. 2). ###### Example 8 (Janis-Newman-Winicour space-time). The most general static spherically symmetric solution to the Einstein massless scalar field equations for a scalar field $`\mathrm{\Phi }`$ satisfying $`\mathrm{}\mathrm{\Phi }=0`$ was obtained by Janis, Newman and Winicour . The Ricci tensor has the form $`R_{ab}=8\pi \mathrm{\Phi }_{;a}\mathrm{\Phi }_{;b}`$. The solution is known to be expressible in the form $`𝐠`$ $`=\left(1{\displaystyle \frac{b}{r}}\right)^\nu dt^2+\left(1{\displaystyle \frac{b}{r}}\right)^\nu dr^2+\left(1{\displaystyle \frac{b}{r}}\right)^{1\nu }r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2\right)`$ (90) $`\mathrm{\Phi }`$ $`={\displaystyle \frac{q}{b\sqrt{4\pi }}}\mathrm{ln}\left(1{\displaystyle \frac{b}{r}}\right)`$ (91) for $`r:b<r<\mathrm{}`$ where the constants $`b`$, $`\nu `$ are related to the ADM mass $`m`$ and scalar charge $`q`$ by $$\nu =\frac{2m}{b},b=2\sqrt{m^2+q^2}.$$ (92) We assume $`b>0`$. There is a curvature singularity at $`r=b`$. In order to obtain $`m0`$ one must assume $`0\nu 1`$. For $`q=0`$ the solution reduces to the Schwarzschild solution. Since all hypersurfaces of the form $`\{r=\text{const.}\}`$ are timelike, one has from the photon sphere equation $`(\text{54})`$ that the only timelike photon sphere is at $$r=\frac{b\left(2\nu +1\right)}{2},$$ (93) which exists only for $`\nu :\frac{1}{2}<\nu 1`$, i.e. for $`0q^2<3m^2`$. For $`\frac{1}{2}<\nu 1`$ it is known that a photon coming from infinity is deflected through an unboundedly large angle, i.e. the photon passes increasingly many times around the singularity as the closest distance of approach tends to the right side of $`(\text{93})`$. ###### Example 9 (Charged dilaton space-time). A static spherically symmetric space-time with a charged dilaton field was obtained by Horne & Horowitz . It comprises a metric $`𝐠`$, a dilaton field $`\mathrm{\Phi }`$ and an electromagnetic field $`F_{ab}`$ given by $`𝐠`$ $`=\left(1{\displaystyle \frac{r_+}{r}}\right)\left(1{\displaystyle \frac{r_{}}{r}}\right)^\omega dt^2+\left(1{\displaystyle \frac{r_+}{r}}\right)^1\left(1{\displaystyle \frac{r_{}}{r}}\right)^\omega dr^2`$ $`+\left(1{\displaystyle \frac{r_{}}{r}}\right)^{1\omega }r^2d\mathrm{\Omega }^2`$ (94) $`\mathrm{e}^{2\mathrm{\Phi }}`$ $`=\left(1{\displaystyle \frac{r_{}}{r}}\right)^{(1\omega )/\beta }`$ (95) $`F_{tr}`$ $`=F_{rt}={\displaystyle \frac{e}{r^2}},\text{all other components vanishing,}`$ (96) where $`r_+`$ and $`r_{}`$ are related to the ADM mass $`m`$ and electric charge $`e`$ by $`r_++\omega r_{}`$ $`=2m`$ (97) $`r_+r_{}`$ $`=e^2(1+\beta ^2)`$ (98) and $`\beta `$ is a free parameter which controls the coupling strength between the dilaton and Maxwell fields, with $`\omega `$ defined in terms of $`\beta ^2`$ by $$\omega :=\frac{1\beta ^2}{1+\beta ^2}.$$ (99) We assume $`m>0`$. For $`\beta =0`$ the solution reduces to the Reissner-Nordström solution considered in Example 7. For $`\beta =0`$ and $`e=0`$ the solution reduces to the Schwarzschild solution. The solution also reduces to the Schwarzschild solution for $`e=r_{}=0`$ and arbitrary $`\beta `$. Here we shall consider the case $`\beta ^2=1`$. In this case one has $`(r_+,r_{})=(2m,e^2/m)`$. There is an event horizon at $`r=r_+=2m`$ and a curvature singularity at $`r=r_{}=e^2/m`$. For $`0(e/m)^2<2`$ the singularity at $`r=r_{}`$ lies inside a black hole whilst for $`(e/m)^2>2`$ it is globally naked. For $`\beta ^2=1`$ the photon sphere equation (54) reduces to $$r_{\text{ps}}^\pm /m=\frac{(6+(e/m)^2)\pm \sqrt{36+(e/m)^420(e/m)^2}}{4}.$$ (100) For $`0(e/m)^2<2`$ one has $`r_{\text{ps}}^{}r_{}<r_+<r_{\text{ps}}^+`$ so there is a single timelike photon sphere. For $`(e/m)^2=2`$ one has $`r_{\text{ps}}^{}=r_{}=r_+=r_{\text{ps}}^+`$ so there are no timelike photon spheres. For $`2<(e/m)^2<18`$ both $`r_{\text{ps}}^+`$ and $`r_{\text{ps}}^{}`$ are complex so there are no timelike photon spheres. For $`(e/m)^218`$ one has $`r_{\text{ps}}^{}r_{\text{ps}}^+<r_{}`$ so there are again no timelike photon spheres. Thus in the black hole case $`0(e/m)^2<2`$ there is a single timelike photon sphere, whilst in the naked singularity case $`(e/m)^2>2`$ there are no timelike photon spheres. (See Fig. 3.) ###### Example 10 (Vaidya null dust collapse). The Vaidya null dust collapse model is a non-static, spherically symmetric space-time with a metric which, in terms of single null (radiation) coordinates $`(u,r,\theta ,\varphi )`$, assumes the form $$𝐠=\left(1\frac{2m(u)}{r}\right)du^2+2dudr+r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)$$ (101) where $`m(u)`$ is a freely specifiable function of $`u`$. Setting $`x^0:=u`$, $`x^1:=r`$ in (35) one obtains $`{\displaystyle \frac{d^2r}{du^2}}`$ $`={\displaystyle \frac{1}{r}}\left(g_{uu}+{\displaystyle \frac{dr}{du}}\right)\left(g_{uu}+2{\displaystyle \frac{dr}{du}}\right)\frac{3}{2}{\displaystyle \frac{dr}{du}}_rg_{uu}\frac{1}{2}\left(g_{uu}_rg_{uu}+_ug_{uu}\right)`$ (102) $`={\displaystyle \frac{1}{r}}\left\{\left(1{\displaystyle \frac{3m(u)}{r}}\right)\left(1{\displaystyle \frac{2m(u)}{r}}3{\displaystyle \frac{dr}{du}}\right){\displaystyle \frac{dm(u)}{du}}+2\left({\displaystyle \frac{dr}{du}}\right)^2\right\}.`$ (103) This is the evolution equation for a spherically symmetric photon surface of the Vaidya collapse metric (101). Consider the special case $$m(u)=\{\begin{array}{cc}0\hfill & :\mathrm{}<u<0\hfill \\ \lambda u\hfill & :0uu_1\hfill \\ m_1:=\lambda u_1\hfill & :u_1<u<\mathrm{}\hfill \end{array}$$ (104) for given constants $`\lambda >0`$, $`u_1>0`$. For $`u<0`$ the space-time is locally Minkowskian, for $`u:0uu_1`$ there is inward falling null dust, and for $`u>u_1`$ the space-time is locally isometric to Schwarzschild space-time with ADM mass $`m_1>0`$. It is well-known there is a curvature singularity at $`r=0`$ and that for $`\lambda :0<\lambda \frac{1}{16}`$ the part of this singularity at $`u=0`$ is locally naked. Fix $`\lambda :0<\lambda \frac{1}{16}`$. For $`u>u_1`$ equation (103) gives, as expected, that there is a photon surface at $`r=3m_1`$. We seek to evolve this photon surface backwards in time, though the in-falling null dust, to obtain a maximally extended photon surface $`S`$. The boundary conditions are $$\begin{array}{c}r=3m_1\\ \frac{dr}{du}=0\end{array}\}\text{at }u=u_1.$$ (105) The results are shown in Fig. 4 for $`\lambda =1/16`$ and $`\lambda =1/32`$ for selected values of $`u_1`$. One sees that in all cases $`S`$ tends in the past direction to a null hypersurface of the form $`\{u=\text{const.}\}`$. The conformal diagram must therefore be of the form sketched in Fig. 5. It is evident from Fig. 5 that the naked central singularity is enclosed within the photon surface $`S`$ in the sense that any partial Cauchy surface $``$ extending to spatial infinity must intersect $`S`$ in a 2-sphere. The physical significance of this may warrant further investigation. ## 6. Concluding remarks The definition of a photon surface given in Section 2 is valid in an arbitrary space-time. However the result that a photon surface must have a second fundamental form which is pure trace indicates that a space-time must be specialized in some respect if it is to contain any photon surfaces. For spherically symmetric space-times there are always photon surfaces that respect the spherical symmetry. For space-times that are not spherically symmetric, the definitions of a photon surface and $`G`$-invariant photon surface may seem too restrictive. The problem is that, in general, one may not have orbiting null geodesics at a fixed radius. In Kerr space-time for example, although there are orbiting null geodesics in the equatorial plane, those null geodesics which move in the direction of rotation do so at a different radius than those which move in the opposite direction. But it seems implausible that a concept which is physically important in the case of exact spherical symmetry should become invalid when even a small amount of angular momentum is introduced. A non-trivial generalization of the concepts of photon surface and $`G`$-invariant photon surface, at least to axially symmetric space-times, is thus required. ## Acknowledgements This research was supported by the NFR of Sweden and NRF of South Africa.
warning/0005/hep-ph0005323.html	ar5iv	text	# Studying Indirect Violation of CP, T and CPT in a 𝐵-factory ## 1 Introduction The concept of indirect violation of discrete symmetries in a neutral meson system makes reference to the non-invariance properties of the effective hamiltonian which governs the time evolution of mesons . Those properties can be analyzed by studying the symmetries in the problem of mixing during the time evolution of meson states, while the possible effects in direct decays must be excluded. In these phenomena of time evolution, T violation is evidenced if, for a given time interval, one detects a difference between the probabilities for $`if`$ and $`fi`$, being $`i`$ and $`f`$ shortnames for neutral meson states. CPT violation would be seen as a difference between $`if`$ and $`\overline{f}\overline{i}`$. For the kaon system, this study has been performed by the CP-LEAR experiment from the preparation of definite flavour states $`K^0`$-$`\overline{K}^0`$. These tagged mesons evolve in time and their later decay to a semileptonic final state projects them again on a definite flavour state. The study of this flavour-to-flavour evolution allows the construction of observables which violate CP and T, or CP and CPT. The results are interpreted in terms of non-invariances of the effective hamiltonian, although some discussion remains on the possible CPT violation in semileptonic decays. Nevertheless, these T- and CPT-odd observables would be zero, even in presence of T and CPT fundamental violation, if there were no absorptive components in the effective hamiltonian. For the kaon system, the different lifetimes of physical states, $`K_S`$ and $`K_L`$, ensures this is not the case. On the contrary, in the case of the $`B_d`$-system, the width difference $`\mathrm{\Delta }\mathrm{\Gamma }`$ between the physical states is expected to be negligible. Then the T- and CPT-odd observables proposed for kaons, which are based on flavour tag, vanish for a $`B`$-meson system with $`\mathrm{\Delta }\mathrm{\Gamma }=0`$. But the $`B_d`$ entangled states can be used to construct alternative observables which are sensitive to T and CPT independently of the value of $`\mathrm{\Delta }\mathrm{\Gamma }`$ . In this paper we study three different types of observables that can be constructed from the entangled states of $`B_d`$ mesons in order to study indirect violation of CP, T and CPT. We separate the observables into two different cathegories 1. genuine asymmetries, characterized by the fact that they are pure symmetry observables, constructed by comparing the probabilities of two conjugated processes, so that the asymmetry vanishes if the relating symmetry is conserved; 2. non-genuine asymmetries, which do not correspond to purely conjugated pairs of processes, so that its non-vanishing value can be mimic by the presence of absorptive parts. Within the first cathegory we can distinguish two types of observables depending on whether they need the support of absorptive parts in order to give non-vanishing asymmetries. To start with we give in section 2 a brief overview of the formalism and notation used to study the problem of indirect violation of discrete symmetries in the $`B`$-system. Then, in section 3, we will review how to construct the flavour asymmetries, the analogous in the $`B`$-system to those observables measured for kaons. They turn out to be proportional to $`\mathrm{\Delta }\mathrm{\Gamma }`$, and their value is expected to be so small that they will not be useful to study the symmetry properties in the case of the $`B_d`$-system. They constitute the first type of observables we are studying, and belong to the genuine cathegory. In section 4 we construct alternative asymmetries, based on a CP tag. In that way we find CP-, T- and CPT-odd time-dependent asymmetries, whose values do not cancel even for $`\mathrm{\Delta }\mathrm{\Gamma }=0`$. This will be the second type of observables in our classification, offering the best chances to study the symmetries in $`B_d`$-mixing. The resulting asymmetries are also genuine observables. The limit of small $`\mathrm{\Delta }\mathrm{\Gamma }`$ causes the time-reversal operation and the exchange of decay products to be equivalent. We exploit this fact in section 5 to build the third kind of asymmetries that involve only the $`J/\mathrm{\Psi }K_S`$ final state, experimentally easy to identify. These observables are non-genuine, since their equivalence to the true T- and CPT-odd asymmetries only holds in the exact limit $`\mathrm{\Delta }\mathrm{\Gamma }=0`$. Therefore in section 6 we will see how the introduction of $`\mathrm{\Delta }\mathrm{\Gamma }0`$ affects the results in sections 4 and 5, and we will discuss whether it is still possible to extract information on the CP, T and CPT parameters from the non-genuine asymmetries even in presence of fake effects due to absorptive parts. We emphasize that the time dependence of the observables plays a relevant role in order to separate genuine and fake effects. But it is still possible to extract some limited information with measurements that have no time resolution, as we will discuss in section 7. Finally, in section 8 we will summarize our conclusions. ## 2 The entangled state of $`B`$-mesons In a $`B`$ factory operating at the $`\mathrm{{\rm Y}}(4S)`$ peak, correlated pairs of neutral $`B`$-mesons are produced through the reaction: $$e^+e^{}\mathrm{{\rm Y}}(4S)B\overline{B}.$$ In the CM frame, the resulting $`B`$-mesons travel in opposite directions, and each one will evolve with the effective hamiltonian of the neutral mesons system. Due to the intrinsic spin of $`\mathrm{{\rm Y}}(4S)`$, the so formed state has definite $`L=1`$. Bose statistics requires the physical $`B^0`$-$`\overline{B}^0`$ state to be symmetric under $`\mathrm{C}𝒫`$, being $`𝒫`$ the operator which permutes the spatial coordinates. Together with this requirement, $`\mathrm{C}=`$ implies that $`𝒫=`$, so that the initial state may be written as $$\|i>=\frac{1}{\sqrt{2}}\left(\|B^0(\stackrel{}{k}),\overline{B}^0(\stackrel{}{k})>\|\overline{B}^0(\stackrel{}{k}),B^0(\stackrel{}{k})>\right).$$ (1) As a consequence, one can never simultaneously have two identical mesons at both sides of the detector. This permits the performance of a flavour tag: if at $`t=0`$ one of the mesons decays through a channel which is only allowed for one flavour of the neutral $`B`$, the other meson in the pair must have the opposite flavour at $`t=0`$. The correlation (1) between both sides of the entangled state holds at any time after the production, until the moment of the first decay. The entangled $`B\overline{B}`$ state can also be expressed in terms of the CP eigenstates $`\|B_\pm \frac{1}{\sqrt{2}}\left(\|B^0\pm \mathrm{CP}\|B^0\right)`$ as $$\|i>=\frac{1}{\sqrt{2}}\left(\|B_{}(\stackrel{}{k}),B_+(\stackrel{}{k})>\|B_+(\stackrel{}{k}),B_{}(\stackrel{}{k})>\right).$$ (2) Thus, if the CP operator is well defined, it is also possible to carry out a CP tag. It is enough for that to have a CP-conserving decay into a definite CP final state, so that its detection allows us to identify the decaying meson as a $`B_+`$ or a $`B_{}`$. As in the flavour case, it is possible to perform such a tag at any time after the production of the entangled state. In Ref. we described how the determination of the CP operator is possible and unambiguous to $`𝒪(\lambda ^3)`$, which is sufficient to discuss both CP-conserving and CP-violating amplitudes in the effective hamiltonian for $`B_d`$ mesons. Here $`\lambda `$ is the flavour-mixing parameter of the CKM matrix . The determination is based on the requirement of CP conservation, to $`𝒪(\lambda ^3)`$, in the $`(sd)`$ and $`(bs)`$ sectors. To this order, however, CP-violation exists in the $`(bd)`$ sector, and it can be classified by referring it to the CP-conserving direction. A $`B_d`$ decay that is governed by the couplings of the $`(sd)`$ or $`(bs)`$ unitarity triangles, or by the $`V_{cd}V_{cb}^{}`$ side of the $`(bd)`$ triangle, will not show any CP violation to $`𝒪(\lambda ^3)`$. We may say that such a channel is free from direct CP violation. Two complex parameters, $`\epsilon _{1,2}`$, describe the CP mixing in the physical states $`\|B_1`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{1+\|\epsilon _1\|^2}}}\left[\|B_++\epsilon _1\|B_{}\right],`$ $`\|B_2`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{1+\|\epsilon _2\|^2}}}\left[\|B_{}+\epsilon _2\|B_+\right].`$ (3) They are invariant under rephasing of the meson states, and physical when the CP operator is well defined . In case of CPT conservation, $`\epsilon _1=\epsilon _2`$. There is another pair of parameters, $`\epsilon `$ and $`\delta `$, which can be alternatively used and has a simpler interpretation in terms of symmetries. These parameters are defined as $$\epsilon \frac{\epsilon _1+\epsilon _2}{2},\delta \epsilon _1\epsilon _2.$$ (4) Their explicit expression in terms of the effective hamiltonian matrix elements, when CPT violation is introduced perturbatively and we may neglect terms which are quadratic in $`\mathrm{\Delta }H_{22}H_{11}`$, is given by $$\epsilon =\frac{\sqrt{H_{12}CP_{12}^{}}\sqrt{H_{21}CP_{12}^{}}}{\sqrt{H_{12}CP_{12}^{}}+\sqrt{H_{21}CP_{12}^{}}},\delta =\frac{2\mathrm{\Delta }}{\left(\sqrt{H_{12}CP_{12}^{}}+\sqrt{H_{21}CP_{12}^{}}\right)^2}.$$ (5) The rephasing invariance <sup>3</sup><sup>3</sup>3We are using the notation $`H_{ij}`$, $`\mathrm{CP}_{ij}`$, etc. to represent the matrix elements of the corresponding operators in the flavour basis, for instance $`H_{12}B^0\|H\|\overline{B}^0`$. of $`\epsilon `$ and $`\delta `$ is apparent from Eq. (5). In the following we will only keep linear terms in $`\delta `$. When we pay attention to the restrictions imposed by discrete symmetries on the effective mass matrix, $`H=M\frac{i}{2}\mathrm{\Gamma }`$, we see that: * CP conservation imposes $`\mathrm{Im}(M_{12}\mathrm{CP}_{12}^{})=\mathrm{Im}(\mathrm{\Gamma }_{12}\mathrm{CP}_{12}^{})=0`$ and $`H_{11}=H_{22}`$; * CPT invariance requires $`H_{11}=H_{22}`$; * T invariance imposes $`\mathrm{Im}(M_{12}\mathrm{CP}_{12}^{})=\mathrm{Im}(\mathrm{\Gamma }_{12}\mathrm{CP}_{12}^{})=0`$. As a consequence, CPT invariance leads to $`\mathrm{\Delta }=0`$ and thus $`\delta =0`$, irrespective of the value of $`\epsilon `$. Similarly, T invariance leads to $`\epsilon =0`$, independently of the value of $`\delta `$. CP conservation requires both $`\epsilon =\delta =0`$. Therefore we have four real parameters which carry information on the symmetries of the effective mass matrix, according to the following list * $`\mathrm{Re}(\epsilon )0`$ signals CP and T violation, with $`\mathrm{\Delta }\mathrm{\Gamma }0`$; * $`\mathrm{Im}(\epsilon )0`$ indicates CP and T violation; * $`\mathrm{Re}(\delta )0`$ means that CP and CPT violation exist; * $`\mathrm{Im}(\delta )0`$ shows CP and CPT violation, with $`\mathrm{\Delta }\mathrm{\Gamma }0`$. To extract information on these symmetry parameters we may study the time evolution of the entangled state (1). As can be found in the literature, the special features of this system can be used to extract information on CP and CPT violation in $`B`$ mesons. We use for the final state the notation $`(X,Y)`$, where $`X`$ is the decay product observed with momentum $`\stackrel{}{k}`$ at a time $`t_0`$, and $`Y`$ the product detected at a later time $`t`$ with momentum $`\stackrel{}{k}`$. The time variables $`\mathrm{\Delta }t=tt_0`$ (with definite positive sign) and $`t^{}=t_0+t`$ are used to describe the process instead of $`t_0`$ and $`t`$. The probability to find an arbitrary final state $`(X,Y)`$ from the initial state (1) is given by $`\|(X,Y)\|^2`$ $`=`$ $`{\displaystyle \frac{1}{8}}\left({\displaystyle \frac{1+\|\epsilon _1\|^2}{1\epsilon _1\epsilon _2}}\right)^2\|{\displaystyle \frac{1+\epsilon _2}{1+\epsilon _1}}\|^2\|X\|B_1\|^2\|Y\|B_1\|^2e^{\mathrm{\Gamma }t^{}}\times `$ (6) $`\times \{(\eta _++\eta _{})\mathrm{cosh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }\mathrm{\Delta }t}{2}}\right)(\eta _+\eta _{})\mathrm{cos}(\mathrm{\Delta }m\mathrm{\Delta }t)`$ $`+2\eta _{\mathrm{re}}\mathrm{sinh}\left({\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }\mathrm{\Delta }t}{2}}\right)2\eta _{\mathrm{im}}\mathrm{sin}(\mathrm{\Delta }m\mathrm{\Delta }t)\},`$ where $`\mathrm{\Gamma }`$ is the averaged width of the physical states $`B_{1,2}`$. The $`\eta `$ coefficients are defined as $`\eta _+`$ $`=`$ $`\|\eta _X+\eta _Y\|^2,`$ $`\eta _{}`$ $`=`$ $`\|\eta _X\eta _Y\|^2,`$ $`\eta _{\mathrm{re}}`$ $`=`$ $`\mathrm{Re}[(\eta _X+\eta _Y)(\eta _X^{}\eta _Y^{})],`$ $`\eta _{\mathrm{im}}`$ $`=`$ $`\mathrm{Im}[(\eta _X+\eta _Y)(\eta _X^{}\eta _Y^{})],`$ (7) with $$\eta _X\frac{X\|B^0\frac{1\epsilon _2}{1+\epsilon _2}\mathrm{CP}_{12}^{}X\|\overline{B}^0}{X\|B^0+\frac{1\epsilon _1}{1+\epsilon _1}\mathrm{CP}_{12}^{}X\|\overline{B}^0}=\frac{1+\epsilon _1}{1+\epsilon _2}\frac{\epsilon _2X\|B_++X\|B_{}}{X\|B_++\epsilon _1X\|B_{}},$$ (8) and an analogous expression for $`\eta _Y`$. We write Eq. (8) in terms of both flavour and CP eigenstates. From the above expressions one can easily check that for $`X=Y`$ only $`\eta _+`$ remains and the probability $`\|(X,Y)\|^2`$ vanishes for $`\mathrm{\Delta }t=0`$. This is a kind of EPR correlation, imposed here by Bose statistics . The integration of the probability (6) over $`t^{}`$ between $`\mathrm{\Delta }t`$ (always positive with our definition) and $`\mathrm{}`$ gives the intensity for the chosen final state, depending only on $`\mathrm{\Delta }t`$ $$I(X,Y;\mathrm{\Delta }t)=\frac{1}{2}_{\mathrm{\Delta }t}^{\mathrm{}}𝑑t^{}\|(X,Y)\|^2.$$ (9) By comparing the intensities corresponding to different processes one builds time-dependent asymmetries that allow the extraction of the relevant parameters. ## 3 Flavour Tag: Genuine observables needing $`\mathrm{\Delta }\mathrm{\Gamma }`$ As we mentioned in the introduction, when a semileptonic final state is seen at one side of the detector, the surviving meson can be tagged as the conjugated flavour state. Its later decay through another semileptonic channel is equivalent to the projection of the evolved meson onto the corresponding definite flavour state. In our notation, we represent as $`\mathrm{}^\pm `$ the final decay product of a semiinclusive decay $`B\mathrm{}^\pm X^{}`$. Thus, the final configuration we denote $`(\mathrm{},\mathrm{})`$ is equivalent to a flavour $``$ flavour evolution, at the meson level. In Table 1 we show the equivalence for the four possible configurations with two charged leptons in the final state. From the processes of Table 1 we may construct two non-trivial asymmetries. If we compare the intensities for the first and second processes and suppose there is no violation either of the $`\mathrm{\Delta }B=\mathrm{\Delta }Q`$ rule or of CPT in the direct decay amplitudes, we obtain the asymmetry $$A(\mathrm{}^+,\mathrm{}^+)\frac{I(\mathrm{}^+,\mathrm{}^+)I(\mathrm{}^{},\mathrm{}^{})}{I(\mathrm{}^+,\mathrm{}^+)+I(\mathrm{}^{},\mathrm{}^{})}=\frac{1\left\|\frac{1\epsilon _2}{1+\epsilon _2}\right\|^2\left\|\frac{1\epsilon _1}{1+\epsilon _1}\right\|^2}{1+\left\|\frac{1\epsilon _2}{1+\epsilon _2}\right\|^2\left\|\frac{1\epsilon _1}{1+\epsilon _1}\right\|^2}\frac{4\frac{\mathrm{Re}(\epsilon )}{1+\|\epsilon \|^2}}{1+4\left(\frac{\mathrm{Re}(\epsilon )}{1+\|\epsilon \|^2}\right)^2},$$ (10) where only linear terms in $`\delta `$ have been kept, since CPT violation is treated in perturbation theory. We observe that this Kabir asymmetry is time independent. This is a genuine CP and T asymmetry, since the second process corresponds to the CP-, or T-transformed of the first one. Thus the asymmetry cannot be faked by absorptive parts in absence of true T violation. However, in the limit $`\mathrm{\Delta }\mathrm{\Gamma }=0`$ both $`\mathrm{Re}(\epsilon )`$ and $`\mathrm{Im}(\delta )`$ vanish, and this quantity will be zero, even if CP and T violation exist. So this observable, in spite of being CP- and T-odd, and constituting a true observable of CP and T violation, also needs, in order to be non-zero, the presence of $`\mathrm{\Delta }\mathrm{\Gamma }0`$. This requirement is fulfilled in the kaon system, where the asymmetry (10) has actually been measured . On the contrary, for $`B_d`$-mesons the negligible value of $`\mathrm{\Delta }\mathrm{\Gamma }`$ predicts that this asymmetry will be difficult to observe. Present limits for $`\mathrm{Re}(\epsilon )`$ are at the level of few percent . The second asymmetry to be constructed from semileptonic decays is $`A(\mathrm{}^+,\mathrm{}^{})`$ $``$ $`{\displaystyle \frac{I(\mathrm{}^+,\mathrm{}^{})I(\mathrm{}^{},\mathrm{}^+)}{I(\mathrm{}^+,\mathrm{}^{})+I(\mathrm{}^{},\mathrm{}^+)}}`$ (11) $``$ $`2{\displaystyle \frac{\mathrm{Re}\left(\frac{\delta }{1\epsilon ^2}\right)\mathrm{Sh}\frac{\mathrm{\Delta }\mathrm{\Gamma }\mathrm{\Delta }t}{2}\mathrm{Im}\left(\frac{\delta }{1\epsilon ^2}\right)\mathrm{sin}(\mathrm{\Delta }m\mathrm{\Delta }t)}{\mathrm{Ch}\frac{\mathrm{\Delta }\mathrm{\Gamma }\mathrm{\Delta }t}{2}+\mathrm{cos}(\mathrm{\Delta }m\mathrm{\Delta }t)}},`$ also to linear order in $`\delta `$. This asymmetry, contrary to that in Eq. (10), depends on time as an odd function of $`\mathrm{\Delta }t`$. Eq. (11) corresponds to a genuine CP and CPT asymmetry, as it compares the probabilities for $`B^0B^0`$ and $`\overline{B}^0\overline{B}^0`$. These transitions are self-conjugated under time reversal transformation, so T violation is not operative in this case. To get a non-zero value of the asymmetry (11), both CP and CPT violation are required, but also $`\mathrm{\Delta }\mathrm{\Gamma }0`$. The proportionality to $`\mathrm{\Delta }\mathrm{\Gamma }`$ is present in both terms of the asymmetry, explicitly in the first one and in the second term through $`\mathrm{Im}(\delta )`$. Therefore measuring a small limit for this observable does not give a straightforward bound on CPT violation, because the vanishingly small $`\mathrm{\Delta }\mathrm{\Gamma }`$ of $`B`$-mesons would hide any symmetry breaking effect. Present limits on $`\mathrm{Im}(\delta )`$ are again at the level of few percent. ## 4 CP Tag: Genuine observables which do not need $`\mathrm{\Delta }\mathrm{\Gamma }`$ As explicitly seen in Eqs. (10) and (11), the flavour asymmetries are difficult to prove discrete symmetry violation in the $`B`$-system, due to the negligible $`\mathrm{\Delta }\mathrm{\Gamma }`$ between the physical states. However, we may construct alternative asymmetries making use of the CP eigenstates, which can be identified in this system by means of a CP tag. From the entangled state (2) such a tag can be performed. Let us suppose that $`X`$, a CP eigenstate produced along the CP-conserving direction, is observed at $`t_0`$ at one side of the detector. Such a decay is free from direct CP violation, and this assures that in the moment of the decay the surviving meson of the other side had the opposite CP eigenvalue. One example of a decay channel with the properties we are looking for is given by $`X`$ $`J/\mathrm{\Psi }K_S`$, with CP$`=`$, or by $`XJ/\mathrm{\Psi }K_L`$, with CP$`=+`$, both of them governed by the “CP-allowed” side $`V_{cb}V_{cs}^{}`$ of the $`(bs)`$ triangle. Then, the detection of such a final state leads to the preparation of the remaining $`B_d`$ meson in the complementary CP eigenstate, to $`𝒪(\lambda ^3)`$. Since CP is conserved in $`K`$ and $`B_s`$ systems to $`𝒪(\lambda ^3)`$, it is correct to identify the physical kaon states with the CP eigenstates, at least to this order in the expansion. After this preparation, the CP tagged state is left to evolve for a certain time $`\mathrm{\Delta }t`$. We are interested in the transition probabilities from CP-eigenstates, $`B_+`$ and $`B_{}`$, to flavour states, $`B^0`$ and $`\overline{B}^0`$. Thus we will be looking at final configurations such as $`(J/\mathrm{\Psi }K_S,\mathrm{}^+)`$ that corresponds, at the meson level, to $`B_+B^0`$. All the transitions connected to this one by symmetry transformations, are shown in Table 2. Comparing the intensity of $`(J/\mathrm{\Psi }K_S,\mathrm{}^+)`$ with each of the processes on Table 2, we construct three genuine asymmetries of the form $$A(X,Y)=\frac{I(X,Y)I(J/\mathrm{\Psi }K_S,\mathrm{}^+)}{I(X,Y)+I(J/\mathrm{\Psi }K_S,\mathrm{}^+)}.$$ (12) Each of them is odd under the discrete transformation which connects ($`J/\mathrm{\Psi }K_S`$$`,\mathrm{}^+)`$ with $`(X,Y)`$. To linear order in $`\delta `$ and in the limit $`\mathrm{\Delta }\mathrm{\Gamma }=0`$, we get the following results: * The CP odd asymmetry, $$A_{\mathrm{CP}}A(J/\mathrm{\Psi }K_S,\mathrm{}^{})=2\frac{\mathrm{Im}(\epsilon )}{1+\|\epsilon \|^2}\mathrm{sin}(\mathrm{\Delta }m\mathrm{\Delta }t)+\frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}\frac{2\mathrm{R}\mathrm{e}(\delta )}{1+\|\epsilon \|^2}\mathrm{sin}^2\left(\frac{\mathrm{\Delta }m\mathrm{\Delta }t}{2}\right),$$ (13) has contributions from T-violating and CPT-violating terms. The first term, odd in $`\mathrm{\Delta }t`$, is governed by the T-violating $`\mathrm{Im}(\epsilon )`$, whereas the second term, which is even in $`\mathrm{\Delta }t`$, is sensitive to CPT violation through the parameter $`\mathrm{Re}(\delta )`$. The CP asymmetry corresponds to the well known “gold plate” decay and has been measured by CDF . The inclusion of CPT violation through $`\mathrm{Re}(\delta )`$ has been considered previously in Ref. . The equivalence between final configurations and mesonic transitions allows the interpretation of this observable in terms of CP-to-flavour transitions. It constitutes then a test of indirect CP, either by T violation or by CPT violation. The separation of T-odd and CPT-odd terms can be done by constructing different asymmetries, which are discussed in the following. * The T asymmetry, $$A_\mathrm{T}A(\mathrm{}^{},J/\mathrm{\Psi }K_L)=2\frac{\mathrm{Im}(\epsilon )}{1+\|\epsilon \|^2}\mathrm{sin}(\mathrm{\Delta }m\mathrm{\Delta }t)\left[1\frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}\frac{2\mathrm{R}\mathrm{e}(\delta )}{1+\|\epsilon \|^2}\mathrm{sin}^2\left(\frac{\mathrm{\Delta }m\mathrm{\Delta }t}{2}\right)\right],$$ (14) purely odd in $`\mathrm{\Delta }t`$, needs $`\epsilon 0`$, and includes CPT even and odd terms. * The CPT asymmetry, $$A_{\mathrm{CPT}}A(\mathrm{}^+,J/\mathrm{\Psi }K_L)=\frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}\frac{2\mathrm{R}\mathrm{e}(\delta )}{1+\|\epsilon \|^2}\frac{1}{12\frac{\mathrm{Im}(\epsilon )}{1+\|\epsilon \|^2}\mathrm{sin}(\mathrm{\Delta }m\mathrm{\Delta }t)}\mathrm{sin}^2\left(\frac{\mathrm{\Delta }m\mathrm{\Delta }t}{2}\right),$$ (15) needs $`\delta 0`$, and includes both even and odd time dependences, so that there is no definite symmetry under a change of sign of $`\mathrm{\Delta }t`$. There is still a fourth discrete transformation, consisting in the exchange in the order of appearance of the decay products $`X`$ and $`Y`$, i. e. $`\mathrm{\Delta }t\mathrm{\Delta }t`$. It transforms $$\begin{array}{c}B_+B^0\\ (J/\mathrm{\Psi }K_S,\mathrm{}^+)\end{array}\stackrel{\mathrm{\Delta }t}{}\begin{array}{c}\overline{B}^0B_{}\\ (\mathrm{}^+,J/\mathrm{\Psi }K_S)\end{array}.$$ (16) Looking at its effect at the meson level it turns out that $`\mathrm{\Delta }t`$ reversal cannot be associated to any of the fundamental discrete symmetries. In spite of this, it can still provide information on the symmetries of the system. In the limit $`\mathrm{\Delta }\mathrm{\Gamma }=0`$, the temporal asymmetry satisfies $$A_{\mathrm{\Delta }t}A(\mathrm{}^+,J/\mathrm{\Psi }K_S)=A(\mathrm{}^{},J/\mathrm{\Psi }K_L)A_\mathrm{T}.$$ (17) This equality is a consequence of $`\mathrm{\Delta }\mathrm{\Gamma }=0`$. In general, the equivalence of T and $`\mathrm{\Delta }t`$ inversions is only valid for hamiltonians with the property of hermiticity, up to a global (proportional to unity) absorptive part. For the kaon system, for instance, relation (17) does not hold. Measuring the presented asymmetries (13)-(17) with good time resolution, so to separate even and odd $`\mathrm{\Delta }t`$ dependences, should be enough to determine the parameters $$\frac{2\mathrm{I}\mathrm{m}(\epsilon )}{1+\|\epsilon \|^2},\frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}\frac{2\mathrm{R}\mathrm{e}(\delta )}{1+\|\epsilon \|^2},$$ (18) governing CP, T violation and CP, CPT violation, respectively, for the $`B_d`$ mixing. Contrary to what happened in section 3, the CPT and T asymmetries based on a CP tag do not vanish due to the smallness of $`\mathrm{\Delta }\mathrm{\Gamma }`$. Instead, they provide a set of observables which could separate the parameters $`\delta `$ and $`\epsilon `$. In case of CPT invariance, $`\delta =0`$, $`A_{\mathrm{CPT}}=0`$, and the four asymmetries reduce to a single independent one $`A_{\mathrm{CP}}=A_\mathrm{T}=A_{\mathrm{\Delta }t}`$. In the Standard Model this measures $`\mathrm{sin}(2\beta )`$ of the $`(bd)`$ unitarity triangle . We emphasize, however, that one has access to four experimentally different asymmetries whose results would also be different in case of CPT violation. In extended theories, efforts to find a plausible framework for CPT violation are discussed in Ref. . ## 5 Non-genuine observables only with $`J/\mathrm{\Psi }K_S`$ All the asymmetries defined in the previous section are genuine observables, since each of them compares the original process with its conjugated under a certain symmetry and is thus odd under the corresponding transformation. Nevertheless the measurement of all those quantities requires to tag both $`B_+`$ and $`B_{}`$ states. The last needs, from the experimental point of view, a good reconstruction of the decay $`B`$$`J/\mathrm{\Psi }K_L`$, not so easy to achieve as for the corresponding $`J/\mathrm{\Psi }K_S`$ channel. Therefore we may wonder how much can be learned about the symmetry parameters from the study of all possible asymmetries built from final configurations $`(X,Y)`$ with only $`J/\mathrm{\Psi }K_S`$. We show that if one considers only final states in which a $`K_S`$ is present, the same asymmetries can still be constructed, provided that the limit $`\mathrm{\Delta }\mathrm{\Gamma }=0`$ is valid. In this case there are four possible configurations of the final state, depending on the sign of the charged lepton and on the order of appearance of the decay products, $$(J/\mathrm{\Psi }K_S,\mathrm{}^+),(J/\mathrm{\Psi }K_S,\mathrm{}^{}),(\mathrm{}^+,J/\mathrm{\Psi }K_S),(\mathrm{}^{},J/\mathrm{\Psi }K_S).$$ (19) We show in Table 3 the mesonic transitions which are related to each final state. In the exact limit $`\mathrm{\Delta }\mathrm{\Gamma }=0`$, taking into account that $`\mathrm{\Delta }t`$ and T operations become equivalent, one can construct asymmetries which measure indirect violation of CP, T and CPT from the processes shown in Table 3. Thus, we obtain, to linear order in the CPT violating parameter $`\delta `$, * $$A(J/\mathrm{\Psi }K_S,\mathrm{}^{})=2\frac{\mathrm{Im}(\epsilon )}{1+\|\epsilon \|^2}\mathrm{sin}(\mathrm{\Delta }m\mathrm{\Delta }t)+2\frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}\frac{\mathrm{Re}(\delta )}{1+\|\epsilon \|^2}\mathrm{sin}^2\left(\frac{\mathrm{\Delta }m\mathrm{\Delta }t}{2}\right).$$ (20) This is the asymmetry $`A_{\mathrm{CP}}`$ defined in Eq. (13). It measures $`\mathrm{Re}(\delta )`$ or $`\mathrm{Im}(\epsilon )`$ different from zero, and constitutes indeed a genuine measurement of indirect CP violation. Since it contains $`\mathrm{\Delta }t`$ odd and $`\mathrm{\Delta }t`$ even pieces, both parameters could be separated. An equivalent observable would be given by the asymmetry between the third and fourth processes of Table 3, also related by a CP transformation. * $$A(\mathrm{}^+,J/\mathrm{\Psi }K_S)=2\frac{\mathrm{Im}(\epsilon )}{1+\|\epsilon \|^2}\mathrm{sin}(\mathrm{\Delta }m\mathrm{\Delta }t)\left[12\frac{\mathrm{Re}(\delta )}{1+\|\epsilon \|^2}\frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}\mathrm{sin}^2\left(\frac{\mathrm{\Delta }m\mathrm{\Delta }t}{2}\right)\right].$$ (21) This is the temporal asymmetry defined in the previous section. It is only different from zero if $`\mathrm{Im}(\epsilon )`$ is not null. Then it measures indirect violation of T when $`\mathrm{\Delta }\mathrm{\Gamma }=0`$. The second and fourth decays of the table are connected by a $`\mathrm{\Delta }t`$ exchange, too, and will measure the same parameters. * $$A(\mathrm{}^{},J/\mathrm{\Psi }K_S)=2\frac{\mathrm{Re}(\delta )}{1+\|\epsilon \|^2}\frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}\mathrm{sin}^2\left(\frac{\mathrm{\Delta }m\mathrm{\Delta }t}{2}\right)\frac{1}{12\frac{\mathrm{Im}(\epsilon )}{1+\|\epsilon \|^2}\mathrm{sin}(\mathrm{\Delta }m\mathrm{\Delta }t)}.$$ (22) This quantity is different from zero only if $`\mathrm{Re}(\delta )`$ is not null. It equals the CPT asymmetry $`A_{\mathrm{CPT}}`$ in Eq. 15 and can be used to measure CPT violation if $`\mathrm{\Delta }\mathrm{\Gamma }=0`$. These results point out a way of testing the discrete symmetries in the evolution of $`B`$-system, from data of entangled states decays. The observed events where a $`K_S`$ is present in the final decay products can be classified according to Eq.(19). This is enough to construct the different asymmetries sensitive to CP, T and CPT violating parameters $`\mathrm{Im}(\epsilon )`$ and $`\mathrm{Re}(\delta )`$, as long as $`\mathrm{\Delta }\mathrm{\Gamma }=0`$. The method exploits the good $`\mathrm{\Delta }t`$-time resolution accessible in asymmetric $`B`$-factories. The separation of $`\epsilon `$ and $`\delta `$ is associated with resolving odd and even functions of $`\mathrm{\Delta }t`$, respectively. ## 6 Linear $`\mathrm{\Delta }\mathrm{\Gamma }`$ corrections The asymmetries in Eqs. (21) and (22), contrary to those built in section 4, are not genuine. They do not correspond to true T- and CPT-odd observables, for the processes we are comparing are not related by a symmetry transformation. This implies that the presence of $`\mathrm{\Delta }\mathrm{\Gamma }0`$ may induce non-vanishing values of the asymmetries discussed in section 5, even if there is no true T- or CPT-violation. Although in the limit $`\mathrm{\Delta }\mathrm{\Gamma }=0`$ the results of Eqs. (21) and (22) agree, respectively, with those in Eqs. (14) and (15), experimentally one is measuring different quantities. When $`\mathrm{\Delta }\mathrm{\Gamma }0`$ is considered, absorptive parts affect all the above expressions, giving corrections to both genuine and non-genuine asymmetries. These absorptive effects appear differently in observables which yielded the same result before, and originate fake contributions to non-genuine asymmetries. To analyze them, we must take into account two types of corrections: 1. time-dependent terms in the amplitudes which are governed by the evolution $`\frac{\mathrm{\Delta }\mathrm{\Gamma }\mathrm{\Delta }t}{2}`$, 2. non-vanishing values of $`\mathrm{Re}(\epsilon )`$ and $`\mathrm{Im}(\delta )`$. To treat the problem in a systematic way, it is convenient to study how the parameters $`\mathrm{Re}(\epsilon )`$ and $`\mathrm{Im}(\delta )`$ depend on the non-vanishing value of $`\mathrm{\Delta }\mathrm{\Gamma }`$. In terms of the matrix elements of the effective hamiltonian, the relevant parameters were given by Eq. (5). From these expressions we see $$\mathrm{Re}(\epsilon )\|M_{12}\|^2\mathrm{Im}\left(\frac{\mathrm{\Gamma }_{12}}{M_{12}}\right),\mathrm{Im}(\delta )\mathrm{\Delta }\|M_{12}\|\left\|\frac{\mathrm{\Gamma }_{12}}{M_{12}}\right\|.$$ (23) Since $`\left\|\frac{\mathrm{\Gamma }_{12}}{M_{12}}\right\|`$ is proportional to $`\mathrm{\Delta }\mathrm{\Gamma }`$, and we are only interested in the first corrections to our asymmetries, we may neglect $`\mathrm{Im}(\delta )`$, which is already first order in both $`\mathrm{\Delta }\mathrm{\Gamma }`$ and $`\mathrm{\Delta }`$, and parametrize $$\epsilon x\mathrm{\Delta }\mathrm{\Gamma }+i\mathrm{Im}(\epsilon ),\delta \mathrm{Re}(\delta ).$$ (24) Using this parametrization it is now easy to throw away terms of order higher than one in $`\mathrm{\Delta }\mathrm{\Gamma }`$ or $`\mathrm{Re}(\delta )`$. The parameter $`\mathrm{Im}(\epsilon )`$ will have also $`\mathrm{\Delta }\mathrm{\Gamma }`$ corrections. However, the new terms will accompany the zero order $`\mathrm{Im}(\epsilon )`$ and have the same time dependence, so that we will be unable to separate them. Thus an explicit expansion of $`\mathrm{Im}(\epsilon )`$ is not needed for our study. ### 6.1 Corrections to flavour observables We discussed in section 3 two asymmetries that could be built from flavour tag, Eq. (10) and Eq. (11). Both of them vanished in case of $`\mathrm{\Delta }\mathrm{\Gamma }=0`$. If we include linear $`\mathrm{\Delta }\mathrm{\Gamma }`$ corrections, as described above, the first observable, $`A(\mathrm{}^+,\mathrm{}^{})`$, acquires a non zero value $$A(\mathrm{}^+,\mathrm{}^+)\frac{4x\mathrm{\Delta }\mathrm{\Gamma }}{1+\|\epsilon \|^2},$$ (25) whereas the second one, $`A(\mathrm{}^+,\mathrm{}^{})`$, is linear in both $`\mathrm{\Delta }\mathrm{\Gamma }`$ and $`\mathrm{\Delta }`$, and vanishes also at this order. ### 6.2 Corrections to genuine CP-to-flavour observables For the genuine asymmetries of section 4, based on a CP tag, we get * $`A_{\mathrm{CP}}`$ $`=`$ $`{\displaystyle \frac{I(J/\mathrm{\Psi }K_S,\mathrm{}^+)I(J/\mathrm{\Psi }K_S,\mathrm{}^{})}{I(J/\mathrm{\Psi }K_S,\mathrm{}^+)+I(J/\mathrm{\Psi }K_S,\mathrm{}^{})}}`$ (26) $``$ $`{\displaystyle \frac{2\mathrm{R}\mathrm{e}(\delta )}{1+\|\epsilon \|^2}}{\displaystyle \frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}}\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }m\mathrm{\Delta }t}{2}}\right){\displaystyle \frac{2\mathrm{I}\mathrm{m}(\epsilon )}{1+\|\epsilon \|^2}}\mathrm{sin}(\mathrm{\Delta }m\mathrm{\Delta }t)`$ $`+{\displaystyle \frac{4x\mathrm{\Delta }\mathrm{\Gamma }}{1+\|\epsilon \|^2}}\left[1+2\left({\displaystyle \frac{2\mathrm{I}\mathrm{m}(\epsilon )}{1+\|\epsilon \|^2}}\right)^2\right]\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }m\mathrm{\Delta }t}{2}}\right)`$ $`{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }\mathrm{\Delta }t}{2}}{\displaystyle \frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}}{\displaystyle \frac{2\mathrm{I}\mathrm{m}(\epsilon )}{1+\|\epsilon \|^2}}\mathrm{sin}(\mathrm{\Delta }m\mathrm{\Delta }t){\displaystyle \frac{8\mathrm{\Delta }\mathrm{\Gamma }x}{1+\|\epsilon \|^2}}\left({\displaystyle \frac{2\mathrm{I}\mathrm{m}(\epsilon )}{1+\|\epsilon \|^2}}\right)^2\mathrm{sin}^4({\displaystyle \frac{\mathrm{\Delta }m\mathrm{\Delta }t}{2}}).`$ With respect to Eq.(13), the $`\mathrm{\Delta }\mathrm{\Gamma }`$ corrections induce both $`\mathrm{\Delta }t`$ even and odd functions. The first two terms in Eq.(26) reproduce the result for $`\mathrm{\Delta }\mathrm{\Gamma }=0`$, whereas the third term has the same time dependence of the first one and would make it difficult to extract $`\mathrm{Re}(\delta )`$ from this CP asymmetry. The last two terms have new $`\mathrm{\Delta }t`$ dependences. * $`A_\mathrm{T}`$ $`=`$ $`{\displaystyle \frac{I(J/\mathrm{\Psi }K_S,\mathrm{}^+)I(\mathrm{}^{},J/\mathrm{\Psi }K_L)}{I(J/\mathrm{\Psi }K_S,\mathrm{}^+)+I(\mathrm{}^{},J/\mathrm{\Psi }K_L)}}`$ (27) $``$ $`{\displaystyle \frac{2\mathrm{I}\mathrm{m}(\epsilon )}{1+\|\epsilon \|^2}}\mathrm{sin}(\mathrm{\Delta }m\mathrm{\Delta }t)\left[1{\displaystyle \frac{2\mathrm{R}\mathrm{e}(\delta )}{1+\|\epsilon \|^2}}{\displaystyle \frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}}\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }m\mathrm{\Delta }t}{2}}\right)\right]`$ $`+{\displaystyle \frac{4x\mathrm{\Delta }\mathrm{\Gamma }}{1+\|\epsilon \|^2}}\left[12\left({\displaystyle \frac{2\mathrm{I}\mathrm{m}(\epsilon )}{1+\|\epsilon \|^2}}\right)^2\right]\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }m\mathrm{\Delta }t}{2}}\right)`$ $`+{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }\mathrm{\Delta }t}{2}}{\displaystyle \frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}}{\displaystyle \frac{2\mathrm{I}\mathrm{m}(\epsilon )}{1+\|\epsilon \|^2}}\mathrm{sin}(\mathrm{\Delta }m\mathrm{\Delta }t)+{\displaystyle \frac{8x\mathrm{\Delta }\mathrm{\Gamma }}{1+\|\epsilon \|^2}}\left({\displaystyle \frac{2\mathrm{I}\mathrm{m}(\epsilon )}{1+\|\epsilon \|^2}}\right)^2\mathrm{sin}^4({\displaystyle \frac{\mathrm{\Delta }m\mathrm{\Delta }t}{2}}).`$ Contrary to what happens for $`A_{\mathrm{CP}}`$, all the new terms in $`A_\mathrm{T}`$ from linear $`\mathrm{\Delta }\mathrm{\Gamma }`$ corrections have different $`\mathrm{\Delta }t`$ dependences as compared to those in Eq.(14). * $`A_{\mathrm{CPT}}`$ $`=`$ $`{\displaystyle \frac{I(J/\mathrm{\Psi }K_S,\mathrm{}^+)I(\mathrm{}^+,J/\mathrm{\Psi }K_L)}{I(J/\mathrm{\Psi }K_S,\mathrm{}^+)+I(\mathrm{}^+,J/\mathrm{\Psi }K_L)}}`$ (28) $``$ $`{\displaystyle \frac{2\mathrm{R}\mathrm{e}(\delta )}{1+\|\epsilon \|^2}}{\displaystyle \frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}}\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }m\mathrm{\Delta }t}{2}}\right){\displaystyle \frac{1}{1\frac{2\mathrm{I}\mathrm{m}(\epsilon )}{1+\|\epsilon \|^2}\mathrm{sin}(\mathrm{\Delta }m\mathrm{\Delta }t)}}.`$ To the order considered in our perturbation expansion, $`A_{\mathrm{CPT}}`$ has no linear $`\mathrm{\Delta }\mathrm{\Gamma }`$ corrections. The genuine character of the asymmetry would put them in higher order terms. We conclude that, in presence of $`\mathrm{\Delta }\mathrm{\Gamma }`$ corrections, the extraction of $`\mathrm{Re}(\delta )`$ should be done from $`A_{\mathrm{CPT}}`$. * $`A_{\mathrm{\Delta }t}`$ $`=`$ $`{\displaystyle \frac{I(J/\mathrm{\Psi }K_S,\mathrm{}^+)I(\mathrm{}^+,J/\mathrm{\Psi }K_S)}{I(J/\mathrm{\Psi }K_S,\mathrm{}^+)+I(\mathrm{}^+,J/\mathrm{\Psi }K_S)}}`$ (29) $``$ $`2{\displaystyle \frac{\mathrm{Im}(\epsilon )}{1+\|\epsilon \|^2}}\mathrm{sin}(\mathrm{\Delta }m\mathrm{\Delta }t)\left[12{\displaystyle \frac{\mathrm{Re}(\delta )}{1+\|\epsilon \|^2}}{\displaystyle \frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}}\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }m\mathrm{\Delta }t}{2}}\right)\right]`$ $`+{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }\mathrm{\Delta }t}{2}}{\displaystyle \frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}}`$ $`{\displaystyle \frac{2x\mathrm{\Delta }\mathrm{\Gamma }}{1+\|\epsilon \|^2}}{\displaystyle \frac{2\mathrm{I}\mathrm{m}(\epsilon )}{1+\|\epsilon \|^2}}\mathrm{sin}(\mathrm{\Delta }m\mathrm{\Delta }t)\left[12\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }m\mathrm{\Delta }t}{2}}\right)\right].`$ From this equation we observe that the “equivalence” between time reversal and $`\mathrm{\Delta }t`$ exchange does not hold any longer. On the contrary, the difference between both asymmetries is now linear in $`\mathrm{\Delta }\mathrm{\Gamma }`$. Except for $`A_{CPT}`$ in Eq. (28), all the asymmetries have linear corrections in $`\mathrm{\Delta }\mathrm{\Gamma }`$. These terms make the analysis more complicated, but do not avoid the separation of parameters, provided one has enough statistics to identify the different time dependences in the asymmetries. It is worth paying closer attention to the way linear $`\mathrm{\Delta }\mathrm{\Gamma }`$ corrections enter the different quantities. $`A_{\mathrm{CP}}`$, $`A_\mathrm{T}`$ and $`A_{\mathrm{CPT}}`$ are genuine symmetry observables, i. e. they are constructed from the comparison between processes related by a true symmetry transformation. They are in general affected by corrections in $`\mathrm{\Delta }\mathrm{\Gamma }`$, but these do not generate fake effects such that even in absence of symmetry violation a non-vanishing asymmetry remains due to $`\mathrm{\Delta }\mathrm{\Gamma }0`$. We can check that in expressions (26)-(28) all $`\mathrm{\Delta }\mathrm{\Gamma }`$ factors are multiplied either by $`\mathrm{Im}(\epsilon )`$ or by $`x`$, both vanishing in case of exact symmetries. On the contrary, $`A_{\mathrm{\Delta }t}`$, which as we have discussed does not correspond to any discrete symmetry, contains terms in $`\mathrm{\Delta }\mathrm{\Gamma }`$ when $`\epsilon =\delta =0`$. ### 6.3 Corrections to non-genuine CP-to-flavour observables The fact that $`\mathrm{\Delta }t`$ and T inversions are no longer equivalent affects the non-genuine asymmetries we have built only with the $`J/\mathrm{\Psi }K_S`$ final state. If we study how the linear terms modify the expressions in Eqs. (20)-(22), we find * the asymmetry in Eq. (20) is still a genuine CP asymmetry, and takes the same value as Eq. (26); * the asymmetry of Eq. (21) is the temporal asymmetry, so it will be equal to that of Eq. (29), and cannot be identified with $`A_\mathrm{T}`$, as was the case for $`\mathrm{\Delta }\mathrm{\Gamma }=0`$; * in the same way, Eq. (22) is not equal to the true CPT asymmetry. Instead we get $`A(\mathrm{}^{}`$ , $`J/\mathrm{\Psi }K_S){\displaystyle \frac{1}{1\frac{2\mathrm{I}\mathrm{m}(\epsilon )}{1+\|\epsilon \|^2}\mathrm{sin}(\mathrm{\Delta }m\mathrm{\Delta }t)}}[{\displaystyle \frac{2\mathrm{R}\mathrm{e}(\delta )}{1+\|\epsilon \|^2}}{\displaystyle \frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}}+{\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }\mathrm{\Delta }t}{2}}{\displaystyle \frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}}`$ (30) $`+{\displaystyle \frac{4x\mathrm{\Delta }\mathrm{\Gamma }}{1+\|\epsilon \|^2}}\mathrm{sin}^2\left({\displaystyle \frac{\mathrm{\Delta }m\mathrm{\Delta }t}{2}}\right){\displaystyle \frac{2\mathrm{I}\mathrm{m}(\epsilon )}{1+\|\epsilon \|^2}}{\displaystyle \frac{2x\mathrm{\Delta }\mathrm{\Gamma }}{1+\|\epsilon \|^2}}\mathrm{sin}(\mathrm{\Delta }m\mathrm{\Delta }t)],`$ which contains linear terms in $`\mathrm{\Delta }\mathrm{\Gamma }`$ even when $`\epsilon =\delta =0`$. ## 7 What can be seen in a symmetric factory? In previous sections we saw the crucial role played by $`\mathrm{\Delta }t`$ resolution in order to separate the symmetry parameters of the $`B`$-system. To achieve the measurement of all the proposed observables one needs to distinguish between configurations where final decay products $`X`$ and $`Y`$ occur with opposite time ordering, and also to determine the time dependence of the resulting quantities with good enough precision. In symmetric $`e^+e^{}`$ facilities at the $`\mathrm{{\rm Y}}(4S)`$ peak, $`B`$ mesons are created at rest, so that they decay essentially at the point of production. Therefore the measurement has no time resolution, but we may still extract some information on indirect symmetry violation from the available observations. ### 7.1 Flavour Tag If $`\mathrm{\Delta }t`$ is not measured, there are three different probabilities, $`\mathrm{\Gamma }_{++}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}d(\mathrm{\Delta }t)I(\mathrm{}^+,\mathrm{}^+)={\displaystyle \frac{C}{4\mathrm{\Gamma }}}\left[1+{\displaystyle \frac{4x\mathrm{\Delta }\mathrm{\Gamma }}{1+\|\epsilon \|^2}}+{\displaystyle \frac{2\mathrm{R}\mathrm{e}(\delta )}{1+\|\epsilon \|^2}}\right]{\displaystyle \frac{\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}{1+\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}},`$ $`\mathrm{\Gamma }_{}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}d(\mathrm{\Delta }t)I(\mathrm{}^{},\mathrm{}^{})={\displaystyle \frac{C}{4\mathrm{\Gamma }}}\left[1{\displaystyle \frac{4x\mathrm{\Delta }\mathrm{\Gamma }}{1+\|\epsilon \|^2}}+{\displaystyle \frac{2\mathrm{R}\mathrm{e}(\delta )}{1+\|\epsilon \|^2}}\right]{\displaystyle \frac{\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}{1+\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}},`$ $`\mathrm{\Gamma }_+`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}d(\mathrm{\Delta }t)[I(\mathrm{}^+,\mathrm{}^{})+I(\mathrm{}^{},\mathrm{}^+)]={\displaystyle \frac{C}{2\mathrm{\Gamma }}}\left[1+{\displaystyle \frac{2\mathrm{R}\mathrm{e}(\delta )}{1+\|\epsilon \|^2}}\right]{\displaystyle \frac{2+\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}{1+\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}},`$ (31) where $`C`$ is a common normalization factor for all the probabilities. The only symmetry observable which can be constructed from these quantities is the asymmetry between $`\mathrm{\Gamma }_{++}`$ and $`\mathrm{\Gamma }_{}`$. This yields the same result of Eq. (10), since all time dependence factored out from $`I(\mathrm{}^+,\mathrm{}^+)`$ and $`I(\mathrm{}^{},\mathrm{}^{})`$. ### 7.2 CP Tag When we sum over all possible values of $`\mathrm{\Delta }t`$, probabilities for channels which do not correspond to the same meson transition are added up. For instance, $`(J/\mathrm{\Psi }K_S,\mathrm{}^+)`$ and $`(\mathrm{}^+,J/\mathrm{\Psi }K_S)`$, which correspond to $`B_+B^0`$ and $`\overline{B}^0B_{}`$, contribute to the same probability. In this way we may construct four quantities $`\mathrm{\Gamma }_{S+}`$ $``$ $`{\displaystyle _0^{\mathrm{}}}[I(J/\mathrm{\Psi }K_S,\mathrm{}^+)+I(\mathrm{}^+,J/\mathrm{\Psi }K_S)]d(\mathrm{\Delta }t)`$ $`={\displaystyle \frac{C}{\mathrm{\Gamma }}}\left\{1+{\displaystyle \frac{2x\mathrm{\Delta }\mathrm{\Gamma }}{1+\|\epsilon \|^2}}{\displaystyle \frac{\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}{1+\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}}+{\displaystyle \frac{\mathrm{Re}(\delta )}{1+\|\epsilon \|^2}}\left[2+{\displaystyle \frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}}{\displaystyle \frac{\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}{1+\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}}\right]\right\};`$ $`\mathrm{\Gamma }_S`$ $``$ $`{\displaystyle _0^{\mathrm{}}}[I(J/\mathrm{\Psi }K_S,\mathrm{}^{})+I(\mathrm{}^{},J/\mathrm{\Psi }K_S)]d(\mathrm{\Delta }t)`$ $`={\displaystyle \frac{C}{\mathrm{\Gamma }}}\left\{1{\displaystyle \frac{2x\mathrm{\Delta }\mathrm{\Gamma }}{1+\|\epsilon \|^2}}{\displaystyle \frac{\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}{1+\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}}+{\displaystyle \frac{\mathrm{Re}(\delta )}{1+\|\epsilon \|^2}}\left[2{\displaystyle \frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}}{\displaystyle \frac{\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}{1+\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}}\right]\right\};`$ $`\mathrm{\Gamma }_{L+}`$ $``$ $`{\displaystyle _0^{\mathrm{}}}[I(J/\mathrm{\Psi }K_L,\mathrm{}^+)+I(\mathrm{}^+,J/\mathrm{\Psi }K_L)]d(\mathrm{\Delta }t)=`$ $`={\displaystyle \frac{C}{\mathrm{\Gamma }}}\left\{1+{\displaystyle \frac{2x\mathrm{\Delta }\mathrm{\Gamma }}{1+\|\epsilon \|^2}}{\displaystyle \frac{\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}{1+\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}}+{\displaystyle \frac{\mathrm{Re}(\delta )}{1+\|\epsilon \|^2}}\left[2{\displaystyle \frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}}{\displaystyle \frac{\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}{1+\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}}\right]\right\};`$ $`\mathrm{\Gamma }_L`$ $``$ $`{\displaystyle _0^{\mathrm{}}}[I(J/\mathrm{\Psi }K_L,\mathrm{}^{})+I(\mathrm{}^{},J/\mathrm{\Psi }K_L)]d(\mathrm{\Delta }t)=`$ (32) $`={\displaystyle \frac{C}{\mathrm{\Gamma }}}\left\{1{\displaystyle \frac{2x\mathrm{\Delta }\mathrm{\Gamma }}{1+\|\epsilon \|^2}}{\displaystyle \frac{\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}{1+\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}}+{\displaystyle \frac{\mathrm{Re}(\delta )}{1+\|\epsilon \|^2}}\left[2+{\displaystyle \frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}}{\displaystyle \frac{\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}{1+\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}}\right]\right\}.`$ Each one of these probabilities corresponds to a pair of meson transitions, according to Table 4. Nevertheless we may construct genuine asymmetries by comparing the probabilities for conjugated pairs of transitions. In this way we find $`𝒜_{\mathrm{CP}}`$ $``$ $`{\displaystyle \frac{\mathrm{\Gamma }_{S+}\mathrm{\Gamma }_S}{\mathrm{\Gamma }_{S+}+\mathrm{\Gamma }_S}}\left[{\displaystyle \frac{2x\mathrm{\Delta }\mathrm{\Gamma }}{1+\|\epsilon \|^2}}+{\displaystyle \frac{\mathrm{Re}(\delta )}{1+\|\epsilon \|^2}}{\displaystyle \frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}}\right]{\displaystyle \frac{\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}{1+\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}},`$ (33) $`𝒜_\mathrm{T}`$ $``$ $`{\displaystyle \frac{\mathrm{\Gamma }_{S+}\mathrm{\Gamma }_L}{\mathrm{\Gamma }_{S+}+\mathrm{\Gamma }_L}}{\displaystyle \frac{2x\mathrm{\Delta }\mathrm{\Gamma }}{1+\|\epsilon \|^2}}{\displaystyle \frac{\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}{1+\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}},`$ (34) $`𝒜_{\mathrm{CPT}}`$ $``$ $`{\displaystyle \frac{\mathrm{\Gamma }_{S+}\mathrm{\Gamma }_{L+}}{\mathrm{\Gamma }_{S+}+\mathrm{\Gamma }_{L+}}}{\displaystyle \frac{\mathrm{Re}(\delta )}{1+\|\epsilon \|^2}}{\displaystyle \frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}}{\displaystyle \frac{\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}{1+\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}}.`$ (35) Since no time dependence is measured, the distinction between genuine and non-genuine observables, which was based on the equivalence between $`\mathrm{\Delta }t`$ and T reversals for $`\mathrm{\Delta }\mathrm{\Gamma }=0`$, does not apply in this case. All the asymmetries here are genuine, since none of them will give a non-zero value due to $`\mathrm{\Delta }\mathrm{\Gamma }`$ corrections, unless there is also a violation of the corresponding symmetry. Looking at the expressions, we realize that $`𝒜_{\mathrm{CPT}}`$ is proportional to $`\mathrm{Re}(\delta )`$, different form zero only if CPT is not conserved. $`𝒜_\mathrm{T}`$ is proportional to the T-violating $`\mathrm{Re}(\epsilon )`$, so that it needs also the presence of absorptive parts to be detected. Finally the CP asymmetry $`𝒜_{\mathrm{CP}}`$, already considered in the literature , contains both a T-violating and a CPT-violating term, with no chance to separate $`\mathrm{Re}(\epsilon )`$ and $`\mathrm{Re}(\delta )`$. The most interesting result is therefore the access to a non-vanishing value of $`\mathrm{Re}(\delta )`$, even for $`\mathrm{\Delta }\mathrm{\Gamma }=0`$, from $`A_{\mathrm{CPT}}`$. There is a constant factor $`\frac{\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}{1+\left(\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}\right)^2}`$ common to all the asymmetries. For the $`B_d`$-system $`\frac{\mathrm{\Delta }m}{\mathrm{\Gamma }}=0.723\pm 0.032`$ , so that this factor will not dilute appreciably the value of the observables. ## 8 Conclusions In this paper we have studied in detail the possibilities to explore indirect violation of discrete symmetries CP, T, CPT in a neutral meson system. We have shown how, even in absence of relative absorptive parts, i. e., if $`\mathrm{\Delta }\mathrm{\Gamma }=0`$, this study is possible, but one needs observables beyond flavour-to-flavour transitions due to time evolution of the meson system. In addition to flavour tags, then, we have considered also CP tags, which are uniquely defined for the $`B_d`$-system at order $`𝒪(\lambda ^3)`$, enough to include CP and T violation at leading order. Possible CPT violation is included perturbatively. This consistent scheme of the treatment of symmetry violation can be tested at the $`B`$ factories, where the production of entangled states of $`B_d`$ mesons allows the preparation of the meson in a CP eigenstate. The asymmetries analyzed in this work exploit their time dependences in order to separate out two different ingredients. On one hand CP and T violation, described by $`\epsilon `$, and on the other CP and CPT violation, given by $`\delta `$. The complex parameters $`\epsilon `$ and $`\delta `$ are defined in a rephasing invariant way and, due to the well defined CP operator, they are unique and physical quantities. The observables we have described can be classified in three types: 1. Genuine asymmetries for T or CPT violation, based on flavour-to-flavour transitions at the meson level, which need the support of absorptive parts, $`\mathrm{\Delta }\mathrm{\Gamma }0`$. The T asymmetry measures $`\frac{\mathrm{Re}(\epsilon )}{1+\|\epsilon \|^2}`$ and does not depend on time. The CPT asymmetry is odd in $`\mathrm{\Delta }t`$, but involves $`\mathrm{Re}\left(\frac{\delta }{1\epsilon ^2}\right)`$ and $`\mathrm{Im}\left(\frac{\delta }{1\epsilon ^2}\right)`$ with different time dependences. Since both asymmetries are zero in absence of $`\mathrm{\Delta }\mathrm{\Gamma }`$, these are not promising observables for $`B_d`$ physics. 2. Genuine observables which do not need $`\mathrm{\Delta }\mathrm{\Gamma }`$. We have constructed three different asymmetries, signals of CP, CPT and T violation, respectively, which are based on the combination of flavour and CP tags. In the limit $`\mathrm{\Delta }\mathrm{\Gamma }=0`$ they involve $`\frac{\mathrm{Im}(\epsilon )}{1+\|\epsilon \|^2}`$ and $`\frac{1\|\epsilon \|^2}{1+\|\epsilon \|^2}\frac{\mathrm{Re}(\delta )}{1+\|\epsilon \|^2}`$, which can be separated out from the different time dependences. There is yet a fourth quantity, the temporal asymmetry, experimentally different from the others, which becomes theoretically equivalent to time reversal asymmetry only in the limit $`\mathrm{\Delta }\mathrm{\Gamma }=0`$. In practice, this kind of observables needs the identification of a semileptonic decay in one side and $`J/\mathrm{\Psi }K_S`$, $`J/\mathrm{\Psi }K_L`$ decays in the other side. 3. Making use of the equivalence between $`\mathrm{\Delta }t`$ and T reversal operations for $`\mathrm{\Delta }\mathrm{\Gamma }=0`$, we have also considered a group of non genuine observables. They are equivalent to the genuine observables described in the previous paragraph in the limit $`\mathrm{\Delta }\mathrm{\Gamma }=0`$, with the advantage that they involve only the hadronic decay $`J/\mathrm{\Psi }K_S`$. As a consequence, these asymmetries are accessible in the first generation of $`B`$ factory experimental results. These observables are non genuine in the sense that absorptive parts can mimic a non-vanishing value for them, even if violation of the fundamental symmetry is not present. The different character of genuine and non genuine observables becomes apparent when corrections due to a non zero $`\mathrm{\Delta }\mathrm{\Gamma }`$ are included. While in genuine observables the effect of $`\mathrm{\Delta }\mathrm{\Gamma }`$ is limited to correct the existing terms of $`\epsilon `$ and $`\delta `$, non genuine asymmetries acquire new terms, proportional to $`\mathrm{\Delta }\mathrm{\Gamma }`$, even for $`\epsilon =\delta =0`$. In a symmetric $`B`$ factory, all the effects which are odd in $`\mathrm{\Delta }t`$ are automatically cancelled by the integration over time implicit in all the experimentally measured quantities. Nevertheless we have seen that some relevant information survives, associated with even terms in $`\mathrm{\Delta }t`$. ## Acknowledgements M. C. B. is indebted to the Spanish Ministry of Education and Culture for her fellowship. This research was supported by CICYT, Spain, under Grant AEN-99/0692.
warning/0005/hep-ph0005167.html	ar5iv	text	# GENIUS project, neutrino oscillations and Cosmology: neutrinos reveal their nature?Presented by M. Zrałek at the Cracow Epiphany Conference on Neutrinos in Physics and Astrophysics, January 6-9, 2000, Cracow, Poland. Work supported in part by the Polish Committee for Scientific Research under Grant No. 2P03B05418. J.G. would like to thank the Alexander von Humboldt-Stiftung for fellowship. ## 1 Introduction There are two main problems in neutrino physics. First is the problem of neutrino mass. In the light of present observations this question seems to be solved, neutrinos are massive particles. The second problem is the one of the neutrino’s nature. Massive neutrinos can be either Dirac or Majorana particles. As their visible interactions are left-handed and known sources generate ultrarelativistic states, it is very difficult to distinguish experimentally between the two. Dirac spin 1/2 fermions were introduced to describe interactions which are invariant under spatial reflections. Majorana fermions were invented later without special purpose. In those times, parity was conserved and it was generally believed that Majorana particle interactions must be asymmetric under spatial reflection. Today we know that particle characters are not responsible for parity symmetry breaking . Only one property - the charge \- discriminates Dirac from Majorana massive fermions. Dirac particles carry charge. Massive Majorana particles must be chargeless and cannot carry static electric or magnetic moments. Neutrinos are “special” fermions, they have no electric charge and only one “charge” - the lepton number can characterize them. From all present terrestrial experiments it follows that family lepton numbers $`L_e`$, $`L_\mu `$ and $`L_\tau `$ are separately conserved and as a result, their sum, the total lepton number $`L=`$ $`L_e`$ $`+L_\mu `$ $`+`$ $`L_\tau `$ has the same property. For massive Dirac neutrinos, flavor lepton numbers can be broken and only L must be conserved. For Majorana neutrinos both, family and total lepton numbers are broken. It is even impossible to define these numbers in the way known from Dirac particles. We have to stress that Majorana neutrinos are more fundamental objects and naturally arise in most extensions of the Standard Model. Only in models where the lepton number ($`L`$ and $`BL`$) is conserved, neutrinos are Dirac particles. However, there are many arguments to abandon lepton number conservation. It is not a fundamental quantity, unlike electric charge and does not govern the dynamics. Also, lepton number violation is naturally induced by the presence of right-handed neutrinos $`\left(\nu _R\right)`$ which are usually necessary to form the Dirac mass term $`\left(\overline{\nu }_L\nu _R\right)`$. In spite of these obvious theoretical arguments, supporting the Majorana nature of neutrinos, finding direct experimental indications which would determine the neutrino character is very important. It is common belief that the first place where the nature of massive neutrinos will be revealed is the neutrinoless double $`\beta `$ decay of nuclei, $`\left(\beta \beta \right)_{0\nu }`$. Many experimental searches for $`\left(\beta \beta \right)_{0\nu }`$ decay of different nuclei have been done and are presently underway . Unfortunately, up to now this decay has not been found and the experimental data can only help to estimate the lower bound on the life times of $`\left(\beta \beta \right)_{0\nu }`$ decay modes. The most stringent limit was found in the germanium Heidelberg-Moscow experiments. Their latest result on the half-life time is $$T_{1/2}^{o\nu }\left(Ge\right)>5.7\times 10^{25}\text{ year (at 90\% CL)}$$ (1) from which the following upper bound on the effective Majorana mass was found $$\|m_\nu \|=\left\|\underset{i}{}U_{ei}^2m_i\right\|<0.2eV.$$ (2) The above number has been used to restrict many aspects of the neutrino mass spectrum or the solar neutrino mechanism . We propose an opposite way of thinking. Using the present data from oscillation experiments, and tritium $`\beta `$ decay we can find the modules $`\left\|U_{ei}\right\|`$ of mixing matrix elements and the possible values of neutrino masses $`m_i`$. Then we check whether the bound Eq. (2) is satisfied or not. If it is, the problem is unsolved. If, however the bound Eq. (2) is not satisfied, then neutrinos are Dirac particles. In the latter case the effective mass is calculated as $$\underset{l=1}{\overset{n}{}}U_{el}^2m_l\underset{l=1}{\overset{n}{}}\frac{1}{2}\left[\left(iU_{el}\right)^2+\left(U_{el}\right)^2\right]m_l=0$$ (3) and is strictly equal zero. Ambitious plans are to shift up the limit Eq. (1) and to move the upper limit of $`<m_\nu >`$ down to 0.02 eV (using a tank of 1 ton of Germanium, after one year) or in a further time scale even to 0.006 eV (1t, 10 years). In a previous work we have considered three neutrino mixing schemes. Here we present analytical results for both three and four neutrino mixing scenarios. Information about the same subject with numerical estimations is given in . In the next Chapter we summarize the efforts undertaken in order to find lepton number violating processes. Explanation is given of why the family and total lepton numbers are so strongly conserved. In Chapter 3 we collect all the relevant information about mixing matrix elements and masses extracted from experimental data. Four presently accepted neutrino mass schemes which cover the case of three and four neutrino mixing are discussed. All necessary information from oscillation experiments, tritium $`\beta `$ decay and cosmology is given. Chapter 4 is the main of the paper. All data are connected together with the bound on effective neutrino mass from $`(\beta \beta )_{0\nu }`$, and restrictions on various neutrino mass schemes are presented. Conclusions are to be found in Chapter 5. ## 2 Lepton numbers and neutrino character of light SM neutrino states In order to explain the lack of lepton flavor violating processes, the concept of the flavor lepton number $`L_\alpha `$ followed by the idea of the total lepton number $`L`$ have been introduced. The upper bounds on branching ratios of $`L_\alpha `$ violating processes are very small, for instance: $$\begin{array}{cc}BR(\mu ^{}e^{}\gamma )<4.910^{11},\hfill & BR(\mu ^{}e^+e^{}e^{})<1.010^{12},\hfill \\ & \\ BR(\pi ^0\mu ^{}e^+)<1.7210^8,\hfill & BR(K_L^0\mu ^{}e^+)<3.310^{11},\hfill \\ & \\ BR(\tau ^{}\mu ^{}\gamma )<4.210^6.\hfill & \end{array}$$ (4) In the frame of the SM with massless neutrinos the above processes are strictly forbidden. If neutrinos are massive, then in analogy to the quark sector neutrinos should mix and lepton numbers are not conserved. However, these effects must be very small, below sensitivity of processes given in Eq. (4). That means that the concept of leptons numbers $`L,L_\alpha `$ is still usefull, at least in all neutrino nonoscillation phenomena. For Dirac neutrinos represented by a bispinor $`\mathrm{\Psi }_D`$ it is possible to change the phase of the field $$\mathrm{\Psi }_De^{i\alpha }\mathrm{\Psi }_D.$$ (5) The charge connected with such a global gauge transformation is just the flavor charge operator. This operator can, but not necessarily must, commute with the interaction Hamiltonian, $`[L_\alpha ,H]=0`$ for a massless, $`[L_\alpha ,H]0`$ for a massive neutrino. Majorana neutrinos on the other hand are described by self-conjugate fields $$\mathrm{\Psi }_M=\mathrm{\Psi }_M^cC\overline{\mathrm{\Psi }}_M^T,$$ (6) and it is not possible to define the same kind of gauge transformation as in Eq. (5). There is then no special reason why $`L_\alpha `$ and $`L`$ should be conserved for Majorana neutrinos. All processes in Eq. (4) break $`L_\alpha `$ but not $`L`$, so they can be realized by both kind of neutrinos at the one loop level. At this level only very heavy, nonstandard, neutrinos matter . We do not go to details and focus only on direct effects connected with light, SM neutrinos. Let us mention only that in see-saw models heavy neutrino effects are also negligible, both at tree and loop levels . To make life easier and to understand how processes with Majorana (Dirac) neutrinos mimic family $`L_\alpha `$ and total lepton $`L`$ numbers conservation let us consider a tree level process of electron (positron) production using electron and muon neutrinos scattering on nuclear target $$\nu _{e(\mu )}Ne^\pm X.$$ (7) Let us define the connection between flavor $`\nu _\alpha `$ and massive $`\nu _i`$ states in the following way $$\|\nu _\alpha (\lambda =\frac{1}{2})=\underset{i}{}U_{\alpha i}\|\nu _i(\lambda =\frac{1}{2})$$ (8) for negative helicity states and $$\|\nu _\alpha (\lambda =+\frac{1}{2})=\underset{i}{}U_{\alpha i}^{}\|\nu _i(\lambda =+\frac{1}{2})$$ (9) for $`\lambda =+\frac{1}{2}`$. In the same way Weyl particle and antiparticle states are connected. Note that Eqs. (8-9) for massive particles seems to be in contradiction with the special theory of relativity. The real problem is that states on the right hand side of Eqs. (8-9) can not be defined, in general . However, the left-handed interaction cannot change the neutrino helicity and it is practically impossible to find a real frame moving faster than the neutrino itself (neutrinos are ultrarelativistic) and the relations Eqs. (8-9) can be safely used . This is what is usually considered to be true when neutrino oscillation phenomena are discussed. To be more general, let us assume that there is also a right-handed neutrino interaction <sup>1</sup><sup>1</sup>1For Dirac neutrinos there are actually two independent neutrino mixing matrices in left- and right- handed charged currents . Even for Majorana neutrinos these could be in principle different (as is e.g. the case of see-saw type models where the light neutrino mixing matrix in the right-handed current is dumped by the heavy neutrino mass scale). These simplifications do not spoil the general idea given here. $`L_{CC}`$ $`=`$ $`{\displaystyle \frac{g}{\sqrt{2}}}\left[A_L\left(\overline{N}_i\gamma ^\mu P_L(U^T)_{i\alpha }l_\alpha \right)W_{L\mu }^++A_R\left(\overline{N}_i\gamma ^\mu P_R(U^{})_{i\alpha }l_\alpha \right)\right]W_{R\mu }^++h.c.`$ Then in the ultrarelativistic regime $`(m_i<<E)`$ using the unitarity of the U matrix, the following amplitudes to the $`𝒪\left(\left(\frac{m_i}{E}\right)^2\right)`$ order are obtained: $`A\left(\nu _e(1/2)e^{}\right)`$ $`=`$ $`A(e^{})\left[A_L^{}+A_R^{}{\displaystyle \underset{i}{}}{\displaystyle \frac{m_i}{2E}}\left(U_{ei}\right)^2\right],`$ (11) $`A\left(\nu _e(+1/2)e^{}\right)`$ $`=`$ $`A(e^{})\left[A_L^{}{\displaystyle \underset{i}{}}{\displaystyle \frac{m_i}{2E}}\left(U_{ei}^{}\right)^2+A_R^{}\right],`$ (12) $`A\left(\nu _\mu (1/2)e^{}\right)`$ $`=`$ $`A(e^{})\left[A_L^{}{\displaystyle \underset{i}{}}{\displaystyle \frac{m_i^2}{8E^2}}U_{\mu i}U_{ei}^{}+A_R^{}{\displaystyle \underset{i}{}}{\displaystyle \frac{m_i}{2E}}U_{\mu i}U_{ei}\right],`$ $`A\left(\nu _\mu (+1/2)e^{}\right)`$ $`=`$ $`A(e^{})\left[A_L^{}{\displaystyle \underset{i}{}}{\displaystyle \frac{m_i}{2E}}U_{\mu i}^{}U_{ei}^{}A_R^{}{\displaystyle \underset{i}{}}{\displaystyle \frac{m_i^2}{8E^2}}U_{\mu i}^{}U_{ei}\right]`$ (14) and $`A\left(\nu _e(1/2)e^+\right)`$ $`=`$ $`A(e^+)\left[A_L{\displaystyle \underset{i}{}}{\displaystyle \frac{m_i}{2E}}\left(U_{ei}\right)^2+A_R\right],`$ (15) $`A\left(\nu _e(+1/2)e^+\right)`$ $`=`$ $`A(e^+)\left[A_L+A_R{\displaystyle \underset{i}{}}{\displaystyle \frac{m_i}{2E}}\left(U_{ei}^{}\right)^2\right],`$ (16) $`A\left(\nu _\mu (1/2)e^+\right)`$ $`=`$ $`A(e^+)\left[A_L{\displaystyle \underset{i}{}}{\displaystyle \frac{m_i}{2E}}U_{\mu i}U_{ei}+A_R{\displaystyle \underset{i}{}}{\displaystyle \frac{m_i^2}{8E^2}}U_{\mu i}U_{ei}^{}\right],`$ $`A\left(\nu _\mu (+1/2)e^+\right)`$ $`=`$ $`A(e^+)\left[A_L{\displaystyle \underset{i}{}}{\displaystyle \frac{m_i^2}{8E^2}}U_{\mu i}^{}U_{ei}+A_R^{}{\displaystyle \underset{i}{}}{\displaystyle \frac{m_i}{2E}}U_{\mu i}^{}U_{ei}^{}\right]`$ (18) where $`A(e^+)`$ and $`A(e^{})`$ are appropriate amplitudes for massless neutrinos. In the approximation $`\frac{m_i}{E}<<1`$ and $`\|A_L\|>>\|A_R\|`$ only two cross sections for electron production by a $`\nu _e(\lambda =1/2)`$ beam Eq. (11) and positron production by a $`\nu _e(\lambda =+1/2)`$ beam Eq. (16) are large enough to be seen $$\sigma \left(\nu _e(1/2)e^{}\right)\|A(e^{})\|^2,$$ (19) $$\sigma \left(\nu _e(+1/2)e^+\right)\|A(e^+)\|^2.$$ (20) All other helicity cross sections are suppressed by factors $$\left(\frac{m_i}{E}\right)^2,\frac{m_i}{2E}\|A_R\|\text{or}\|A_R\|^2,$$ (21) and for instance, for $`m_i1`$ eV and $`E1`$ MeV we have $`\left(\frac{m_i}{E}\right)^210^{12}`$. Such factors cause that the cross sections for flavor lepton number $`L_\alpha `$ violating processes Eqs. (LABEL:3a-14), Eqs. (LABEL:3b-18) are invisibly small. The total lepton $`L`$ non-conserving processes Eqs. (12,15), share the same property. Neglecting the factors from Eq. (21), our amplitudes are identical to those of massless Weyl neutrinos whose family and total lepton numbers are strictly conserved. Turning our results around we can see that processes where neutrino masses (and right-handed currents) are not important give no chance to distinguish Dirac from Majorana neutrinos. Could CP phases help? In the case of Dirac \[Majorana\] neutrinos the mixing matrix U has $`(n1)(n2)/2\left[n(n1)/2\right]`$ phases. Let us look into processes where the neutrino mass is important. Though the transition probability of neutrino oscillations depends on CP phases, the physical phases by which the neutrino mixing matrices differ do not enter into transition probabilities and the results are the same for Dirac and Majorana neutrinos . The neutrino mass distortion measured in tritium $`\beta `$ decay is a function of absolute values of mixing matrix elements (see next chapter) so it is not sensible to CP phases, either. There are also processes which do not conserve the total lepton number in which only Majorana neutrinos could participate. Since many years the most promising investigation along this line is connected with the neutrinoless double beta decay. Surprisingly, we will see that even if this process is not observed, it can solve the problem of the nature of neutrinos, when augmented with Cosmology (assuming neutrinos as Hot Dark Matter) and neutrino oscillations results. ## 3 Neutrino masses and mixing matrix $`U_{ei}`$ elements There are two completely different situations which depend on the present status of the LSND experiment. Three light neutrinos are necessary to explain solar and atmospheric anomalies. With the LSND result an extra light neutrino must be introduced. ### 3.1 Three neutrinos scenario For neutrino mixing 3 flavor states $`(\nu _e,\nu _\mu ,\nu _\tau )`$ are related to 3 eigenmass states $`(\nu _1,\nu _2,\nu _3)`$ through $$\left(\begin{array}{c}\nu _e\\ \nu _\mu \\ \nu _\tau \end{array}\right)=\left(\begin{array}{ccc}U_{e1}& U_{e2}& U_{e3}\\ U_{\mu 1}& U_{\mu 2}& U_{\mu 3}\\ U_{\tau 1}& U_{\tau 2}& U_{\tau 3}\end{array}\right)\left(\begin{array}{c}\nu _1\\ \nu _2\\ \nu _3\end{array}\right).$$ (22) Our concern is about the first row of the mixing matrix. We use the standard parameterization $$U=\left(\begin{array}{ccc}c_{12}c_{13}& s_{12}c_{13}& s_{13}e^{i\delta }\\ s_{12}c_{23}c_{12}s_{23}s_{13}e^{i\delta }& c_{12}c_{23}s_{12}s_{23}s_{13}e^{i\delta }& s_{23}c_{13}\\ s_{12}s_{23}c_{12}c_{23}s_{13}e^{i\delta }& c_{12}s_{23}s_{12}c_{23}s_{13}e^{i\delta }& c_{23}c_{13}\end{array}\right)$$ (23) To explain the solar neutrino anomaly the mass splitting between two neutrinos must be extremely tiny, $`\delta m_{sun}^210^5÷10^{11}eV^2`$. Only slightly larger mass splitting between neutrino masses is needed in the case of atmospheric oscillations, $`\delta m_{atm}^210^2÷10^3eV^2`$. These relations leave us with two possible neutrino mass scenarios: $`\delta m_{12}^2=\delta m_{sun}^2`$ , $`\delta m_{23}^2\delta m_{13}^2=\delta m_{atm}^2`$ (Scheme $`A_3`$, Fig.1) and $`\delta m_{23}^2=\delta m_{sun}^2`$, $`\delta m_{12}^2\delta m_{13}^2=\delta m_{atm}^2`$ (Scheme $`B_3`$, Fig.1) Reactor experiments are of the so-called short baseline and are able to measure the neutrino mass splitting of the order $`\delta m^2<10^3`$. Then we can neglect terms with $`\delta m_{sun}^210^3`$ and the disappearance probability for $`\overline{\nu }_e`$ reactor neutrino oscillations is (the following discussion is given for the $`A_3`$ scheme) $$P_{\overline{\nu }_e\overline{\nu }_e}=P_{\nu _e\nu _e}=1\mathrm{sin}^22\mathrm{\Theta }_{13}\mathrm{sin}^2\mathrm{\Delta }_{reactor},$$ (24) where $$\mathrm{\Delta }_{reactor}=\mathrm{\Delta }_{23}(L_{reactor},E_{reactor}),\mathrm{\Delta }_{ij}(L,E)=\frac{1.27\times \delta m_{ij}^2[eV^2]L[km]}{E[GeV]}.$$ In reactor experiments the disappearance of $`\overline{\nu }_e`$ is not seen, which means that $`\mathrm{sin}^22\mathrm{\Theta }_{13}`$ must be small. CHOOZ gives $$\mathrm{sin}^22\mathrm{\Theta }_{13}<0.18,$$ (25) and two solutions for $`\mathrm{\Theta }_{13}`$ can be found: $$\mathrm{sin}^2\mathrm{\Theta }_{13}<0.05\text{or}\mathrm{sin}^2\mathrm{\Theta }_{13}>0.95.$$ (26) The observed $`\nu _\mu `$ neutrino deficit from the atmosphere is favorable describe by a $`\nu _\mu \nu _\tau `$ transition where matter effects are not important $`\left(\mathrm{\Delta }_{atm}=\mathrm{\Delta }_{23}(L_{atm},E_{atm})\right)`$: $$P_{\nu _\mu \nu _\tau }=\mathrm{sin}^22\mathrm{\Theta }_{23}\mathrm{cos}^4\mathrm{\Theta }_{13}\mathrm{sin}^2\mathrm{\Delta }_{atm}.$$ (27) We know that the atmospheric neutrino mixing is very large $$0.72\mathrm{sin}^22\mathrm{\Theta }_{23}\mathrm{cos}^4\mathrm{\Theta }_{13}1\text{and}\delta m_{atm}^24\times 10^3eV^2$$ (28) and only a small value of $`\mathrm{sin}^2\mathrm{\Theta }_{13}`$ in Eq. (26) is compatible with the bound in Eq. (28). The recent fit to the new (830-920 days) atmospheric data of Superkamiokande gives the minimum of $`\chi ^2`$ for $$\mathrm{sin}^2\mathrm{\Theta }_{13}=0.03.$$ (29) Similar values ($`\mathrm{sin}^2\mathrm{\Theta }_{13}<0.03÷0.04`$ ) are given by the reactor data. In all solar neutrino experiments the deficit of electron antineutrinos is measured and four different solutions are possible . The first, “just so” solution, is based on the hypothesis of neutrino oscillations in vacuum (VO), $`\delta m_{sun}^210^{10}eV^2`$ in this case. The other three solutions are based on the Wolfenstein-Mikheyev-Smirnov mechanism of coherent neutrino scattering in matter (so called small mixing angle (SMA MSW), large mixing angle (LMA MSW) and low $`\delta m^2`$ (LOW MSW) solutions). For VO the $`\nu _e`$ disappearance probability is given by $`\left(\mathrm{\Delta }_{sun}=\mathrm{\Delta }_{12}(L_{sun},E_{sun})\right)`$ $`P_{\nu _e\nu _e}^{sun}`$ $`=`$ $`1{\displaystyle \frac{1}{2}}\mathrm{sin}^22\mathrm{\Theta }_{13}\mathrm{sin}^22\mathrm{\Theta }_{12}\mathrm{cos}^4\mathrm{\Theta }_{13}\mathrm{sin}^2\mathrm{\Delta }_{sun}.`$ (30) This expression can be rewritten in the form $$P_{\nu _e\nu _e}^{sun}=\mathrm{cos}^4\mathrm{\Theta }_{13}\left(1\mathrm{sin}^22\mathrm{\Theta }_{12}\mathrm{sin}^2\mathrm{\Delta }_{sun}\right)+\mathrm{sin}^4\mathrm{\Theta }_{13}.$$ (31) Taking into account that $`\mathrm{sin}^4\mathrm{\Theta }_{13}0`$ (Eq. (29)) we get $$P_{\nu _e\nu _e}^{sun}1\mathrm{sin}^22\mathrm{\Theta }_{sun}\mathrm{sin}^2\mathrm{\Delta }_{sun}.$$ (32) where $`\mathrm{\Theta }_{sun}\mathrm{\Theta }_{12}`$. Similarly we get for the case of the MSW solution: $$P_{\nu _e\nu _e}^{sun(MSW)}1\mathrm{sin}^22\stackrel{~}{\mathrm{\Theta }}_{sun}\mathrm{sin}^2\stackrel{~}{\mathrm{\Delta }}_{sun}.$$ (33) where now $$\mathrm{sin}^22\stackrel{~}{\mathrm{\Theta }}_{sun}=\frac{\mathrm{sin}^22\mathrm{\Theta }_{sun}}{\left[\left(\frac{A}{\delta m_{sun}^2}\mathrm{cos}2\mathrm{\Theta }_{sun}\right)^2+\mathrm{sin}^22\mathrm{\Theta }_{sun}\right]^{1/2}}.$$ (34) $`\stackrel{~}{\mathrm{\Delta }}_{sun}`$ includes the effective neutrino mass parameter with $`\delta m_{sun}^2`$ replaced by $$\stackrel{~}{\delta }m_{sun}^2=\delta m_{sun}^2\left[\left(\frac{A}{\delta m_{sun}^2}\mathrm{cos}2\mathrm{\Theta }_{sun}\right)^2+\mathrm{sin}^22\mathrm{\Theta }_{sun}\right]^{1/2},$$ (35) where $`A=2\sqrt{2}G_FEN_e`$ ($`N_e`$ \- electron number density). From Eq. (35) we can see that in order to fulfill the resonance condition we need for $`\delta m_{sun}^2>0`$ $$\mathrm{cos}2\mathrm{\Theta }_{sun}>0,$$ (36) and this means that $$\mathrm{cos}\mathrm{\Theta }_{12}>\mathrm{sin}\mathrm{\Theta }_{12}.$$ (37) Many fits have been done to the solar neutrino data . The results of the fit which takes into account the full set of measurements (rates, energy spectrum, day-night asymmetry in the case of the MSW solution and seasonal variation for VO solution) are presented in Table I. For VO only the best fit value $`\mathrm{sin}^22\mathrm{\Theta }_{sun}`$ is given in . For a scheme $`B_3`$ a change $`U_{e3}U_{e1}`$ must be done. ### 3.2 Four neutrinos scenario The electron (anti)neutrino appearance in the LSND experiment can be explained by $`\nu _\mu \nu _e`$ oscillation with additional large $`\delta m^2`$ scale $$\delta m_{LSND}^20.2÷2\text{eV}^2.$$ (38) In principle there are six possible four-neutrino mass schemes with three different scales of $`\delta m^2`$. They are widely discussed in literature and it is known that only two schemes (Fig.2) are accepted by reactor, LSND, solar and atmospheric neutrino data. As the parameterization of the $`4\times 4`$ neutrino mixing matrix is very complicated in the case when all entries of the mass matrix are nonzero we will use only the symbolic denotations and take $$\left(\begin{array}{c}\nu _e\\ \nu _\mu \\ \nu _\tau \\ \nu _s\end{array}\right)=\left(\begin{array}{cccc}U_{e1}& U_{e2}& U_{e3}& U_{e4}\\ U_{\mu 1}& U_{\mu 2}& U_{\mu 3}& U_{\mu 4}\\ U_{\tau 1}& U_{\tau 2}& U_{\tau 3}& U_{\tau 4}\\ U_{s1}& U_{s2}& U_{s3}& U_{s4}\end{array}\right)\left(\begin{array}{c}\nu _1\\ \nu _2\\ \nu _3\\ \nu _4\end{array}\right).$$ (39) For the short baseline experiment (and for scheme $`A_4`$) the probability of disappearance of $`\nu _e`$ neutrinos is given by $`\left(\mathrm{\Delta }_{SBL}=\mathrm{\Delta }_{32}(L_{SBL},E_{SBL})\right)`$ $$P_{\nu _e\nu _e}=1c_e(1c_e)\mathrm{sin}^2\mathrm{\Delta }_{SBL}$$ (40) where $$c_e=\|U_{e3}\|^2+\|U_{e4}\|^2.$$ Again we can implement the CHOOZ result and get $$4c_e(1c_e)<0.18,$$ (41) and there are two solutions for $`c_e`$, namely $$c_e<0.05\text{or}c_e>0.95.$$ (42) On the other hand the deficit of solar neutrinos in the VO scenario and in four-neutrino language reads $`\left(\mathrm{\Delta }_{sun}=\mathrm{\Delta }_{12}(L_{sun},E_{sun})\right)`$ $$P_{\nu _e\nu _e}^{(4)}=\left(1\|U_{e3}\|^2\|U_{e4}\|^2\right)^2\left[1\frac{4\|U_{e1}\|^2\|U_{e2}\|^2}{\left(\|U_{e1}\|^2+\|U_{e2}\|^2\right)^2}\mathrm{sin}^2\mathrm{\Delta }_{sun}\right]+\|U_{e3}\|^4+\|U_{e4}\|^4.$$ (43) By comparison with the two flavour oscillations formula $$P_{\nu _e\nu _e}=1\mathrm{sin}^22\mathrm{\Theta }_{sun}\mathrm{sin}^2\mathrm{\Delta }_{sun}$$ (44) we can see that the factor $`\left(1\|U_{e3}\|^2\|U_{e4}\|^2\right)^2`$ must be close to one, so only one solution of Eq. (42) is possible, namely $`c_e<0.05`$. For small values of $`c_e`$ the probability $`P_{\nu _e\nu _e}^{(4)}`$ is well described in the frame of the 3 neutrino scenario by $`P_{\nu _e\nu _e}^{(3)}`$ (Eq. (31)) with the following substitutions $`\|U_{e3}\|^2+\|U_{e4}\|^2`$ $``$ $`\mathrm{sin}^2\mathrm{\Theta }_{13}<0.05,`$ $`{\displaystyle \frac{\|U_{e1}\|^2}{\|U_{e1}\|^2\|U_{e2}\|^2}}`$ $`=`$ $`\mathrm{cos}^2\mathrm{\Theta }_{12},{\displaystyle \frac{\|U_{e2}\|^2}{\|U_{e1}\|^2\|U_{e2}\|^2}}=\mathrm{sin}^2\mathrm{\Theta }_{12}.`$ (45) That means that fitted parameters are the same as in 3-neutrino case $$\mathrm{sin}^2\mathrm{\Theta }_{13}=\|U_{e3}\|^2+\|U_{e4}\|^2c_e<0.05.$$ (46) The MSW and VO solutions are described by $`\mathrm{sin}^22\mathrm{\Theta }_{12}=\mathrm{sin}^22\mathrm{\Theta }_{sun}`$ with the same values as in 3 neutrino scenario given in Table I. For a scheme $`B_4`$ a change $`U_{e3(4)}U_{e1(2)}`$ must be done. ### 3.3 Tritium beta decay Other constraints on neutrino masses and mixings come from the observation of the end of the Curie plot for the tritium $`\beta `$ decay. Two collaborations from Mainz and Troitsk give similar results for the upper limit ($`95\%`$ of c.l.) on the effective electron neutrino mass $$m_{\nu _e}_\beta =\left[\underset{i=1}{\overset{n}{}}\|U_{ei}\|^2m_i^2\right]^{1/2},$$ (47) $`m_{\nu _e}_\beta `$ $`<`$ $`2.8eV\text{Mainz Collaboration }\text{[37]},`$ $`m_{\nu _e}_\beta `$ $`<`$ $`2.5eV\text{Troitsk Collaboration }\text{[38]}.`$ Both collaborations have ambitious plans to probe the mass region below 1 eV during the next five years . ### 3.4 Cosmological bounds There are also astrophysical and cosmological bounds on neutrino masses and mixings. All this information depends on many other assumptions (as e.g. nonzero cosmological constant $`\mathrm{\Lambda }`$) and is not as strict as laboratory data. We will take into account only one data which comes from the so called dark matter problem. If neutrinos compose all invisible matter in the Universe then $$m_\nu 30eV.$$ (48) If only 20 % of all dark matter is formed by neutrinos (the so called Hot Dark Matter) then $$m_\nu 6eV.$$ (49) The best fit to many cosmological quantities is obtained if around 70 % of dark matter is given by nonzero cosmological constant, 24 % by Cold Dark Matter and 6 % by Hot Dark Matter. In such a case $$m_\nu 2eV.$$ (50) ## 4 Neutrinoless double beta decay and constraints on neutrino nature Neutrinoless double beta decay is sensitive to the first element of the neutrino mass matrix $$m_{\alpha \beta }=\underset{i=1}{\overset{n}{}}U_{\alpha i}U_{\beta i}m_i$$ (51) and luckily, is very well constrained. This is not the case of other entries which are also measured in various laboratory experiments, for instance $`\|m_{e\mu }\|`$ $`\text{in }Ti+\mu ^{}Ca+e^+`$ $`\|m_{\mu \mu }\|`$ $`\text{in }K^+\pi ^{}\mu ^+\mu ^+`$ (52) $`\|m_{e\tau }\|,\|m_{\mu \tau }\|,\|m_{\tau \tau }\|`$ $`\text{in HERA from}e^{}p\nu _el^{}l^{}X.`$ All these quantities have quite large bounds, in the MeV-GeV range e.g. $`\|m_{e\mu }\|`$ $`<17\text{MeV},`$ $`\|m_{\mu \mu }\|`$ $`<500\text{GeV},`$ $`\|m_{e\tau }\|`$ $`<8.4\text{TeV}.`$ (53) The mixing matrix for Majorana neutrino has 3(6) phases for 3(4) neutrinos so we have $`\|m_\nu \|`$ $`=`$ $`\left\|\|U_{e1}\|^2m_1+e^{2i\varphi _2}\|U_{e2}\|^2m_2+e^{2i\varphi _3}\|U_{e3}\|^2m_3\right\|,`$ (54) $`\|m_\nu \|`$ $`=`$ $`\|\|U_{e1}\|^2m_1+e^{2i\varphi _2}\|U_{e2}\|^2m_2+e^{2i\varphi _3}\|U_{e3}\|^2m_3\|+e^{2i\varphi _4}\|U_{e4}\|^2m_4\|,`$ for $`n=3(4)`$ neutrinos, respectively. We should stress that all our results are obtained in the approximation in which the lightest of neutrinos $`(m_\nu )_{min}`$ is heavier than the difference of squares of neutrino masses responsible for solar neutrino oscillations $`\left((m_\nu )_{min}>>\delta m_{sun}^2\right)`$. ### 4.1 A schemes Let us first discuss the schemes $`A_3`$ and $`A_4`$. We have * for $`A_3`$ $`m_1`$ $`=`$ $`(m_\nu )_{min},`$ (56) $`m_2`$ $`=`$ $`\sqrt{(m_\nu )_{min}^2+\delta m_{sun}^2}m_1,`$ $`m_3`$ $`=`$ $`\sqrt{(m_\nu )_{min}^2+\delta m_{atm}^2+\delta m_{sun}^2}\sqrt{(m_\nu )_{min}^2+\delta m_{atm}^2},`$ * for $`A_4`$ $`m_1,m_2`$ $`\text{as for}A_3,`$ $`m_3`$ $`=`$ $`\sqrt{(m_\nu )_{min}^2+\delta m_{LSND}^2+\delta m_{sun}^2}\sqrt{(m_\nu )_{min}^2+\delta m_{LSND}^2},`$ $`m_4`$ $`=`$ $`\sqrt{m_3^2+\delta m_{atm}^2}m_3.`$ (57) Using the relation $$min\|z_1+z_2+z_3+z_4\|=\{\begin{array}{c}\|z_3+z_4\|_{min}\|z_1+z_2\|_{max}>0,\hfill \\ 0,\hfill \\ \|z_1+z_2\|_{min}\|z_3+z_4\|_{max}>0,\hfill \end{array}$$ (58) we get for both schemes $$\|m_\nu \|_{min}=\{\begin{array}{cc}s_{min}\left(\|U_{e1}\|^2m_1+\|U_{e2}\|^2m_2\right)\hfill & (m_\nu )_{min}(0,x_1^A),\hfill \\ 0\hfill & (m_\nu )_{min}(x_1^A,x_2^A),\hfill \\ \|\|U_{e1}\|^2m_1\|U_{e2}\|^2m_2\|\|s_{max}\hfill & (m_\nu )_{min}>x_2^A,\hfill \end{array}$$ (59) where $`s_{\mathrm{min}}`$ $`=`$ $`s_{\mathrm{max}}=c_em_3,(A_3)\text{scheme}`$ (60) $`s_{\mathrm{min}}`$ $`=`$ $`\left\|\left\|U_{e3}\right\|^2m_3\left\|U_{e4}\right\|^2m_4\right\|,(A_4)\text{scheme}`$ (61) $`s_{\mathrm{max}}`$ $`=`$ $`\left\|U_{e3}\right\|^2m_3+\left\|U_{e4}\right\|^2m_4,(A_4)\text{scheme}.`$ (62) $`x_1^A`$ and $`x_2^A`$ are the values of $`(m_\nu )_{min}=m_1>0`$ for which $$s_{\mathrm{min}}\left(\left\|U_{e1}\right\|^2m_1+\left\|U_{e2}\right\|^2m_2\right)=0\text{ and }\left\|\left\|U_{e1}\right\|^2m_1\left\|U_{e2}\right\|^2m_2\right\|s_{\mathrm{max}}=0$$ (63) respectively. In both schemes there is (in agreement with Eq. (37) we take $`\left\|U_{e1}\right\|^2>\left\|U_{e2}\right\|^2`$) $$\left\|U_{e1}\right\|^2m_1+\left\|U_{e2}\right\|^2m_2=(m_\nu )_{min}\left(1c_e\right)$$ (64) and $$\left\|\left\|U_{e1}\right\|^2m_1\left\|U_{e2}\right\|^2m_2\right\|=(m_\nu )_{min}\left(1c_e\right)\sqrt{1\mathrm{sin}^22\theta _{sun}}.$$ (65) In the $`A_4`$ scheme we do not know $`\left\|U_{e3}\right\|^2`$ and $`\left\|U_{e4}\right\|^2`$ separately and only approximate values for $`s_{\mathrm{max}}`$ can be found $$s_{\mathrm{max}}=\left\|U_{e3}\right\|^2m_3+\left\|U_{e4}\right\|^2\sqrt{m_3+\delta m_{atm}^2}c_e\sqrt{(m_\nu )_{min}^2+\delta m_{LSND}^2}.$$ (66) The $`s_{\mathrm{min}}`$ is unknown so the region of $`(m_\nu )_{min}(0,x_1^A)`$ cannot be checked precisely. We can find however that in both schemes $`(\delta m^2=\delta m_{atm}^2`$ or $`\delta m_{LSND}^2)`$ $`\left\|m_\nu \right\|_{\mathrm{min}}`$ (67) $`s_{\mathrm{max}}\left\|U_{e1}\right\|^2m_1\left\|U_{e2}\right\|^2m_2=c_e\sqrt{(m_\nu )_{min}^2+\delta m^2}(m_\nu )_{min}\left(1c_e\right)`$ $`<\{\begin{array}{c}c_e\sqrt{\delta m_{atm}^2}0.002eV\text{ for }A_3,\hfill \\ c_e\sqrt{\delta m_{LSND}^2}0.03eV\text{ for }A_4.\hfill \end{array}`$ (70) The region $`(m_\nu )_{min}>x_2^A`$ is more interesting. In both schemes this region occurs if $$(m_\nu )_{min}\left(1c_e\right)\sqrt{1\mathrm{sin}^22\theta _{sun}}c_e\sqrt{(m_\nu )_{min}^2+\delta m^2}0,$$ (71) from which the condition for $`\mathrm{sin}^22\theta _{sun}`$ follows $$\mathrm{sin}^22\theta _{sun}\frac{12c_e}{\left(1c_e\right)^2}.$$ (72) For such values of mixing angle $`\theta _{sun}`$ we can find $`x_2^A`$ $$x_2^A=\frac{\delta m^2c_e}{\sqrt{\left(1c_e\right)^2\left(1\mathrm{sin}^22\theta _{sun}\right)c_e}}.$$ (73) ### 4.2 B schemes For $`B_3`$ and $`B_4`$ schemes the neutrino masses are connected with the lightest neutrino mass as follows: * for $`B_3`$ $$\begin{array}{c}m_1=(m_\nu )_{min},\hfill \\ m_2=\sqrt{(m_\nu )_{min}^2+\delta m_{atm}^2},\hfill \\ m_3=\sqrt{(m_\nu )_{min}^2+\delta m_{atm}^2+\delta m_{sun}^2}m_3\hfill \end{array}$$ (74) * for $`B_4`$ $$\begin{array}{c}m_1,m_2\text{ as in }B_3,\hfill \\ m_3=\sqrt{(m_\nu )_{min}^2+\delta m_{atm}^2+\delta m_{LSND}^2},\hfill \\ m_4=\sqrt{m_3^2+\delta m_{sun}}m_3.\hfill \end{array}$$ (75) Using the relation Eq. (58) we obtain $$\left\|m_\nu \right\|_{\mathrm{min}}=\{\begin{array}{c}w_{\mathrm{min}}\left(\left\|U_{e3}\right\|^2m_3+\left\|U_{e4}\right\|^2m_4\right)>0,(m_\nu )_{min}(0,x_1^B)\hfill \\ 0,(m_\nu )_{min}(x_1^B,x_2^B)\hfill \\ \left\|\right\|U_{e3}\|^2m_3\left\|U_{e4}\|^2m_4\right\|w_{\mathrm{max}},(m_\nu )_{min}>x_2^B.\hfill \end{array}$$ (76) where ($`c_e=\left\|U_{e1}\right\|^2`$ or $`\|U_{e1}\|^2+\|U_{e2}\|^2)`$ $`w_{\mathrm{min}}`$ $`=`$ $`w_{\mathrm{max}}=c_em_1,(B_3)\text{scheme}`$ (77) $`w_{\mathrm{min}}`$ $`=`$ $`\left\|\left\|U_{e1}\right\|^2m_1\left\|U_{e2}\right\|^2m_2\right\|,(B_4)\text{scheme}`$ (78) $`w_{\mathrm{max}}`$ $`=`$ $`\|U_{e1}\|^2m_1+\|U_{e2}\|^2m_2,,(B_4)\text{scheme}.`$ (79) $`x_1^B`$ and $`x_1^B`$ are solutions of the equations $$w_{\mathrm{min}}\left(\left\|U_{e3}\right\|^2m_3+\left\|U_{e4}\right\|^2m_4\right)=0,\text{ and }\left\|\left\|U_{e3}\right\|^2m_3\left\|U_{e4}\right\|^2m_4\right\|w_{\mathrm{max}}=0,$$ (80) respectively. Now there is (once more we assume $`\left\|U_{e3}\right\|>\left\|U_{e4}\right\|`$) $$\left\|U_{e3}\right\|^2m_3+\left\|U_{e4}\right\|^2m_4=\left(1c_e\right)\sqrt{(m_\nu )_{min}^2+\delta m^2},\text{ and}$$ (81) $$\left\|\left\|U_{e3}\right\|^2m_3\left\|U_{e4}\right\|^2m_4\right\|=\left(1c_e\right)\sqrt{(m_\nu )_{min}^2+\delta m^2}\sqrt{1\mathrm{sin}^22\theta _{sun}}$$ (82) where $`\delta m^2=\delta m_{atm}^2`$ (for $`B_3)`$ and $`\delta m^2=\delta m_{atm}^2+\delta m_{LSND}^2`$ for $`B_4`$. As $`0c_e0.05`$ $$w_{\mathrm{min}}\left(\left\|U_{e3}\right\|^2m_3+\left\|U_{e4}\right\|^2m_4\right)<c_e(m_\nu )_{min}\left(1c_e\right)\sqrt{(m_\nu )_{min}^2+\delta m_{atm}^2}<0,$$ (83) and the first two regions in Eq. (76) are not present. In the $`B_4`$ scheme, as in the $`A_4`$, we do not know $`\left\|U_{e1}\right\|^2`$ and $`\left\|U_{e2}\right\|^2`$ separately. For $`w_{\mathrm{max}}`$ only the bound can be found $$c_e(m_\nu )_{min}<w_{\mathrm{max}}<c_e\sqrt{(m_\nu )_{min}^2+\delta m_{atm}^2}.$$ (84) In this case the $`\left\|m_\nu \right\|`$ satisfies $$\left\|m_\nu \right\|_{\mathrm{min}}\left(1c_e\right)\sqrt{(m_\nu )_{min}^2+\delta m^2}\sqrt{1\mathrm{sin}^22\theta _{sun}}c_em_{min}f\left[(m_\nu )_{min}\right]$$ (85) where $`m_{min}=(m_\nu )_{min}`$ for $`B_3`$ and $`m_{min}=\sqrt{(m_\nu )_{min}^2+\delta m_{atm}^2}`$ for $`B_4.`$ If condition Eq. (72) is satisfied $`f\left[(m_\nu )_{min}\right]`$ is an increasing function of $`(m_\nu )_{min}`$, if not, $`f\left[(m_\nu )_{min}\right]`$ decreases from $$f\left[0\right]=\left(1c_e\right)\sqrt{\delta m^2}\sqrt{1\mathrm{sin}^22\theta _{sun}}c_em_{min}\left[(m_\nu )_{min}=0\right]\text{ for }(m_\nu )_{min}=0$$ (86) to $$f\left[(m_\nu )_{min}\right]=0\text{ for }(m_\nu )_{min}=\left[\frac{\delta m^2\left(1\mathrm{sin}^22\theta _{sun}\right)\frac{c_e^2}{\left(1c_e\right)^2}m_{min}\left[(m_\nu )_{min}=0\right]}{\mathrm{sin}^22\theta _{sun}\frac{\left(12c_e\right)}{\left(1c_e\right)^2}}\right]^{1/2}.$$ (87) All the above analytical considerations lead us to the following conclusions (for plots see ). 1. Present bound $`\left\|m_\nu \right\|<0.2`$ $`eV.`$ (see Fig. 3) * If SMA MSW solution is the proper mechanism then: + $`B_4`$ scheme is excluded for Majorana neutrinos, + In schemes $`A_3,A_4`$ and $`B_3`$ Majorana neutrinos are accepted if $`\left(m_\nu \right)_{\mathrm{min}}<0.22`$ $`eV.`$ Above this mass all three schemes are open only for Dirac neutrinos. * For LMA and LOW MSW solutions: + the $`A_3,A_4`$ and $`B_3`$ schemes accept Majorana neutrinos only for $`\left(m_\nu \right)_{\mathrm{min}}<1.5`$ $`eV`$, an analogous limit in the $`B_4`$ scheme is $`\left(m_\nu \right)_{\mathrm{min}}<1.1`$ $`eV.`$ 2. If GENIUS I gives only a bound $`\left\|m_\nu \right\|<0.02`$ $`eV:`$ (see Fig. 3) * For SMA MSW solution: + scheme $`B_3`$ is excluded for Majorana neutrinos, + in schemes $`A_3`$ and $`A_4`$ Majorana neutrinos are accepted only for small masses $`\left(m_\nu \right)_{\mathrm{min}}<0.04`$ $`eV.`$ * For LMA and LOW solutions: + the $`B_4`$ scheme is excluded for Majorana neutrinos, + Majorana neutrinos can exist for $`\left(m_\nu \right)_{\mathrm{min}}<0.16`$ $`eV`$ $`\left(A_3\right),`$ $`\left(m_\nu \right)_{\mathrm{min}}<0.14`$ $`eV`$ $`\left(B_3\right)`$ and $`\left(m_\nu \right)_{\mathrm{min}}<0.22eV\left(A_4\right).`$ 3. If finally GENIUS II does not find the $`\left(\beta \beta \right)_{0\nu }`$ decay (see Fig. 4): * For SMA MSW solution: + Majorana neutrinos with $`\left(m_\nu \right)_{\mathrm{min}}0.02eV`$ can exist only in the $`A_3`$ and $`A_4`$ schemes * For LMA MSW and LOW solutions + the $`B_3`$ scheme is excluded for Majorana neutrinos, + Majorana neutrinos can exist only in $`A_3`$ and $`A_4`$ schemes with $`\left(m_\nu \right)_{\mathrm{min}}<0.05`$ eV and $`\left(m_\nu \right)_{\mathrm{min}}<0.12`$ eV, respectively. There are additional restrictions with assumption that neutrinos contribute to the dark matter content. Three neutrinos with almost degenerate masses $`m_\nu 0.7`$ $`eV`$ $`\left(2eV\right)`$ must exist if $`m_\nu 2`$ $`eV`$ $`\left(6eV\right).`$ This means that already the present $`\left(\beta \beta \right)_{0\nu }`$ bound closes all schemes for three Majorana neutrinos if the SMA solution is the proper one. The GENIUS I bound will close schemes for three Majorana neutrinos. For schemes $`A_4`$ and $`B_4`$ with a sterile neutrino $`\left(m_\nu \right)_{\mathrm{min}}`$ must be very small if $`m_\nu 2eV`$ and $`\left(m_\nu \right)_{\mathrm{min}}1.0`$ if $`m_\nu 6`$ $`eV`$. Then the only scheme with one sterile neutrino is accepted if the sum of all Majorana neutrinos is approximately 2 eV. If $`m_\nu 6`$ $`eV`$ and GENIUS I will give negative results only 3 or 4 Dirac neutrinos can constitute the HDM. In such case there is a problem how to explain the number of neutrino degrees of freedom from the abundance of the <sup>4</sup>He and D/N. The present highest bound is $`N_\nu <5.3`$ . ## 5 Conclusions. We have entered an exciting era in neutrino physics. Mixing in the lepton sector seems to be established. An obvious consequence of this fact is the nonconservation of lepton family number $`L_\alpha .`$ Breaking of the $`L_\alpha `$ is very weak. It is seen only in neutrino oscillation and in no other terrestrial laboratory experiments. The problem of the conservation or violation of the total lepton number L, which is connected with the Dirac or Majorana neutrino nature, is not solved up to now. Approximate conservation of $`L_\alpha `$ and $`L`$ follows from i) smallness of neutrino masses, ii)ultrarelativistic character of produced neutrinos, iii) unitarity (exact or approximate) of the mixing matrix, iv) left-handed nature of the neutrino interaction. For Majorana neutrinos this approximate $`L_\alpha `$ and $`L`$ conservation can be proved even though lepton numbers are not defined for neutral particles. The Majorana neutrino mass matrix elements $`m_{\alpha \beta }`$ ($`\alpha ,\beta =e,\mu ,\tau `$) are bounded by various experiments. Such bounds are usually in the MeV-GeV range. Only one element $`m_{ee}`$ is limited with good enough precision to play a role in the reconstruction of the mixing in the lepton sector. The $`m_{ee}`$ element is measured in double $`\beta `$ decay of various nuclei. Up to now this decay has not been observed. The contrary would establish the Majorana nature of neutrinos. However the combination of various informations about masses and mixing matrix elements from i) oscillation experiments, (ii) tritium $`\beta `$ decay and (iii) cosmology together with $`\left(\beta \beta \right)_{0\nu }`$ is able to discriminate between the accepted neutrino mass spectra allowed for Majorana or only for Dirac neutrinos. The data are not precise enough to make conclusive statements. The bound on $`m_\nu `$depends strongly on the determination of nuclear matrix elements. Our estimation was made with 95% CL. At $`3\sigma `$ which corresponds to 99% CL, a value of one for $`\mathrm{sin}^22\theta _{sun}`$ is accepted and we cannot make any discrimination between the two natures. Our estimation is interesting also for those who strongly believe that neutrinos are Majorana particles. We found the corner of the mass schemes where such neutral particles are still allowed. With the present experimental precision the room for the Majorana neutrino is bounded but still large. If next $`\left(\beta \beta \right)_{0\nu }`$ experiments give negative results the Majorana neutrino corner will become smaller and smaller. More precise informations about i) existence of sterile neutrino, ii) which solution of the solar neutrino anomaly is accepted and iii) knowledge of oscillation parameters with smaller error are urgently needed. We hope that future (already working and planned) experiments will provide us with these informations and together with the neutrino mass scheme, the neutrino character will be established.
warning/0005/hep-ph0005274.html	ar5iv	text	# Reconciliation of the Measurement of Parity-Nonconservation in Cs with the Standard Model ## Abstract Contributions from the Breit interaction in atomic-structure calculations account for 1.3$`\sigma `$ of the previously reported $`2.5\sigma `$ deviation from the Standard Model in the <sup>133</sup>Cs weak charge \[S.C. Bennett and C.E. Wieman, Phys. Rev. Lett. 82, 2484 (1999)\]. The updated corrections for the neutron distribution reduce the discrepancy further to 1.0$`\sigma `$. The updated value of the weak charge is $`Q_\mathrm{W}(^{133}\mathrm{Cs})=72.65(28)_{\mathrm{expt}}(34)_{\mathrm{theor}}`$. The present analysis is a higher-order extension of previous calculation \[A. Derevianko, E-print physics/0001046\]. Atomic parity-nonconserving (PNC) experiments combined with accurate atomic structure calculations provide powerful constraints on “new physics” beyond the Standard Model of elementary particles . Compared to high-energy experiments or low-energy scattering experiments, atomic single-isotope PNC measurements are uniquely sensitive to new isovector heavy physics . Presently, the PNC effect in atoms has been most precisely measured by Wieman and co-workers using $`{}_{}{}^{133}\mathrm{Cs}`$ . In 1999, Bennett and Wieman updated the value of the Cs weak charge by measuring a supporting quantity, the vector transition polarizability $`\beta `$, and by re-evaluating the precision of atomic structure calculations from the early 1990s. The determined weak charge differed from the prediction of the Standard Model by 2.5 standard deviations $`\sigma `$. The value of the $`{}_{}{}^{133}\mathrm{Cs}`$ weak charge from Ref. (together with other precision electroweak observables) has been employed in numerous articles. In particular, recent theoretical investigations interpret this 2.5$`\sigma `$ deviation as possible evidence for extra neutral vector $`Z`$-bosons. The main focus of the two previous ab initio relativistic calculations for the atomic structure of <sup>133</sup>Cs was the correlation contribution from the residual Coulomb interaction (i.e., beyond Dirac-Hartree-Fock level). The purpose of this work is to evaluate rigorously contributions from the Breit interaction to PNC in <sup>133</sup>Cs. The previous calculations either omitted such contributions , or evaluated them only partially . The present analysis is a higher-order extension of my recent calculation . It is found that the Breit contribution corrects the weak charge by 0.9%, reducing the 2.5$`\sigma `$ deviation from the Standard Model to 1.2$`\sigma `$. Including a correction for the neutron density distribution in the <sup>133</sup>Cs nucleus further reduces the deviation to $`1.0\sigma `$. Thus the result reported here brings the most accurate atomic PNC measurement to date into substantial agreement with the Standard Model. The Breit interaction arises due to an exchange of transverse photons between electrons. Its low-frequency limit, employed here, is given by $`B_{ij}={\displaystyle \frac{1}{2r_{ij}}}\left(\alpha _i\alpha _j+(\alpha _i\widehat{r}_{ij})(\alpha _j\widehat{r}_{ij})\right)`$It is convenient to separate the second-quantized Breit interaction into zero-, one-, and two-body parts normally ordered with respect to the core: $`B=B^{(0)}+B^{(1)}+B^{(2)}`$. The parity-nonconserving amplitude for the $`6S_{1/2}7S_{1/2}`$ transition in $`{}_{}{}^{133}\mathrm{Cs}`$ can be represented as a sum over intermediate states $`mP_{1/2}`$ $`E_{\mathrm{P}NC}={\displaystyle \underset{m}{}}{\displaystyle \frac{7S\|D\|mP_{1/2}mP_{1/2}\|H_W\|6S}{E_{6S}E_{mP_{1/2}}}}`$ (1) $`+`$ $`{\displaystyle \underset{m}{}}{\displaystyle \frac{7S\|H_W\|mP_{1/2}mP_{1/2}\|D\|6S}{E_{7S}E_{mP_{1/2}}}}.`$ (2) Here $`D`$ and $`H_\mathrm{W}`$ are electric-dipole and weak interaction matrix elements, and $`E_i`$ are atomic energy levels. It is convenient to break the total Breit correction $`\delta E_{\mathrm{P}NC}`$ into three distinct parts due to corrections in the weak interaction and dipole matrix elements, and energy denominators, respectively $$\delta E_{\mathrm{P}NC}=E_{\mathrm{P}NC}(\delta H_\mathrm{W})+E_{\mathrm{P}NC}(\delta D)+E_{\mathrm{P}NC}(\delta E).$$ (3) The overwhelming contribution from parity-violating interactions arises from the Hamiltonian $$H_\mathrm{W}=\frac{G_F}{\sqrt{8}}Q_\mathrm{W}\rho _{\mathrm{n}uc}(r)\gamma _5,$$ (4) where $`G_F`$ is the Fermi constant, $`\gamma _5`$ is the Dirac matrix, and $`\rho _{\mathrm{n}uc}(r)`$ is the neutron density distribution. To be consistent with the previous calculations the $`\rho _{\mathrm{n}uc}(r)`$ is taken to be a proton Fermi distribution employed in Ref. . The slight difference between the neutron and proton distributions is addressed in the conclusion. The PNC amplitude is expressed in units of $`10^{11}i\|e\|a_0(Q_\mathrm{W}/N)`$, where $`N=78`$ is the number of neutrons in the nucleus of <sup>133</sup>Cs. In these units the results of past calculations for $`{}_{}{}^{133}\mathrm{Cs}`$ are $`E_{\mathrm{P}NC}=0.905`$, Ref. , and $`E_{\mathrm{P}NC}=0.908`$, Ref. . The former value includes a partial Breit contribution $`+0.002`$, and the latter includes none. The reference Coulomb-correlated amplitude $$E_{\mathrm{P}NC}^C=0.9075$$ (5) is determined as an average, with the partial Breit contribution removed from the value of Ref. . Hartree-Fock analysis— Before proceeding to the correlated calculations discussed in the second part of this work, it is worth examining the Breit contribution to the PNC amplitude at the lowest-order level. The conventional Dirac-Hartree-Fock (DHF) equation reads $$\left(h_\mathrm{D}+V_{\mathrm{HF}}\right)\varphi _i=\epsilon _i\varphi _i,$$ (6) where $`h_D`$ is the Dirac Hamiltonian including the interaction of an electron in state $`i`$ with a finite-size nucleus. $`V_{\mathrm{HF}}`$ is a mean-field Hartree-Fock potential; this potential contains direct and exchange Coulomb interactions of electron $`i`$ with core electrons. A set of DHF equations is solved self-consistently for core orbitals; valence wavefunctions and energies are determined subsequently by “freezing” the core orbitals. The Breit-Dirac-Hartree-Fock (BDHF) approximation constitutes the introduction of the one-body part of the Breit interaction $`B^{(1)}`$ into the above DHF equation $$\left(h_\mathrm{D}+\stackrel{~}{V}_{\mathrm{HF}}+B^{(1)}\right)\stackrel{~}{\varphi }_i=\stackrel{~}{\epsilon _i}\stackrel{~}{\varphi _i}.$$ (7) Compared to the DHF equations, energies, wave-functions, and the Hartree-Fock potential are modified, as designated by tildes. This self-consistent BDHF approximation was used by Lindroth et al. and a related iterative analysis was considered by Johnson et al. . Both papers point out the importance of the “relaxation” effect, which leads to modification of the Hartree-Fock potential through adjustment of core orbitals. In the present work, the relaxation effect is taken into account automatically by direct integration of Eq. (7). Most of the Breit contribution to the PNC amplitude can be determined by limiting the summation over intermediate states in Eq. (2) to the two lowest valence $`P_{1/2}`$ states: $`6P_{1/2}`$ and $`7P_{1/2}`$. In the DHF approximation one then finds $`E_{\mathrm{P}NC}=0.6888`$ (90% of the total value). The lowest-order corrections to matrix elements and energy denominators calculated as differences between BDHF and DHF values are listed in Table I. The resultant BDHF corrections to $`E_{\mathrm{P}NC}`$ are: $`E_{\mathrm{P}NC}(\delta H_\mathrm{W})`$ $`=`$ $`0.0022(0.32\%),`$ (8) $`E_{\mathrm{P}NC}(\delta D)`$ $`=`$ $`0.0020(0.29\%),`$ (9) $`E_{\mathrm{P}NC}(\delta E)`$ $`=`$ $`0.0019(0.28\%).`$ (10) The sum of these three terms leads to $`\delta E_{\mathrm{P}NC}=0.0023`$ in agreement with the 0.002 correction found by Blundell et al. . Inclusion of intermediate states beyond $`6P_{1/2}`$ and $`7P_{1/2}`$ leads to a small additional modification to $`\delta E_{\mathrm{P}NC}`$ of -0.00004. Note that if experimental energies (which effectively include the Breit interaction) are used in the energy denominators of Eq. (2), then the $`E_{\mathrm{P}NC}(\delta E)`$ term must be excluded and the total correction becomes twice as large: $`\delta E_{\mathrm{P}NC}=0.0042`$. With further examination of the modifications of individual uncorrelated matrix elements presented in Table I, one notices the following. (i) Weak interaction matrix elements are each reduced in absolute value by 0.3%, which is directly reflected in a 0.3% correction to the PNC amplitude. (ii) Modification of dipole amplitudes is strongly nonuniform. There are substantial corrections only to the $`6S_{1/2}7P_{1/2}`$ (0.5%) and $`7S_{1/2}6P_{1/2}`$ (0.1%) matrix elements. The large 0.5% Breit correction to $`6S_{1/2}\|D\|7P_{1/2}`$ provides partial resolution to a long-standing discrepancy of spectroscopic experiment and ab initio calculations . The relatively large Breit correction is caused both by an accidentally small matrix element and by admixture into $`6S_{1/2}\|D\|7P_{1/2}`$ from a 30 times larger $`7S_{1/2}7P_{1/2}`$ matrix element. (iii) The largest modification in the energy denominators is 0.1% for $`E_{7S}E_{6P}`$; however, this leads to a 0.3% correction $`E_{\mathrm{P}NC}(\delta E)`$. As recently emphasized by Dzuba et al. , such large sensitivity of the resulting PNC amplitude to small variations in individual atomic properties entering Eq. (2) arises due to a cancellation of relatively large terms in the sum over states. Correlated calculations— It is well known that correlations caused by residual Coulomb interactions not included in the Hartree-Fock equations can lead to substantial modifications of the lowest-order values. For example, the weak matrix element $`6S_{1/2}\|H_\mathrm{W}\|6P_{1/2}`$ is increased by a factor of 1.8 by correlations due to residual Coulomb interactions. It will be shown that the correlations are also important for a proper description of the Breit corrections. The major correlation effects in atoms appear because of shielding of externally applied (e.g., electric) fields by core electrons and an additional attraction of the valence electron by an induced dipole moment of the core . The former effect is described by contributions beginning at second order and the latter in third order of many-body perturbation theory (MBPT). Since these two effects lead to the dominant contributions in Coulomb-correlated calculations, the third-order analysis reported here seems sufficient . MBPT calculations were performed with the two-body Breit interaction $`B^{(2)}`$ treated on equal footing with the residual Coulomb interaction. Sample many-body diagrams are presented in Fig. 1. To treat the one-body contribution $`B^{(1)}`$, an extension of the B-spline basis set technique was developed, based on the Breit-Dirac-Hartree-Fock (BDHF) equation (7). Such a formulation made it possible to handle $`B^{(1)}`$ and the associated relaxation effect exactly. Contributions of negative-energy states, discussed for example in Ref. , were also included and found to be relatively small . Two series of third-order calculations were performed, first with the Breit and Coulomb interactions fully included using the BDHF basis set, and second in the DHF basis set without the Breit interaction and negative-energy states. The obtained differences are the Breit corrections reported in Table I. Breit corrections to <sup>133</sup>Cs hyperfine-structure magnetic-dipole constants $`A`$ are discussed first, since these were considered in the literature previously. The correction to hyperfine constants is very sensitive to correlations: e.g., Ref. found a numerically insignificant modification for $`A_{6S}`$, while Ref. determined the modification to be large (-4.64 MHz), and the approach reported here yields +4.89 MHz. In the calculation of Ref. the correction was determined as a difference of the BDHF and DHF values, however such approach misses two-body Breit corrections of comparable size. In Ref. a second order perturbation analysis was used for the Breit interaction, but the important relaxation effect discussed earlier was omitted. The present calculation incorporates all mentioned diagrams and is also extended to third order. Using this same calculational scheme, the corrections to hyperfine constants for other states of <sup>133</sup>Cs are +1.16 MHz for $`7S_{1/2}`$, -0.51 MHz for $`6P_{1/2}`$, and -0.146 MHz for $`7P_{1/2}`$. These corrections improve agreement with experiments for the ab initio all-order Coulomb-correlated calculations to 0.1% for all states except $`6P_{1/2}`$ where the discrepancy becomes 0.5%. Examination of the third-order corrections listed in Table I reveals the significant effect of correlations on the Breit contribution. For example, corrections to weak interaction matrix elements become three times larger than those in the lowest order. Compared to hyperfine-structure constants there is no cancellation of various contributions to the weak interaction matrix elements. Using third-order matrix elements and second-order energies the following ab initio corrections are determined: $`E_{\mathrm{P}NC}(\delta H_\mathrm{W})=0.0043`$, $`E_{\mathrm{P}NC}(\delta D)=0.0035`$, and $`E_{\mathrm{P}NC}(\delta E)=0.0028`$. Thus the lowest-order corrections given in Eq. (9) are increased. To improve the consistency of the calculation, one can combine all-order Coulomb-correlated matrix elements and experimental energy denominators tabulated in Ref. with the present third-order Breit corrections. The results are: $`E_{\mathrm{P}NC}(\delta H_\mathrm{W})`$ $`=`$ $`0.0047(0.5\%),`$ (11) $`E_{\mathrm{P}NC}(\delta D)`$ $`=`$ $`0.0037(0.4\%).`$ (12) The Breit correction in energy-denominators $`E_{\mathrm{P}NC}(\delta E)`$ was set to zero because the experimental energies were extensively used in Ref. . For example, the experimental energies were employed in eight out of ten test cases in the scatter analysis of Ref. based on Eq. (2). The total 0.9% Breit correction, $`\delta E_{\mathrm{P}NC}=0.0084`$, is two times larger than the corresponding lowest-order modification, which is rather common in conventional Coulomb-correlated calculations. An even larger 2% Breit correction was found in related calculations of the electric-dipole-moment enhancement factor in thallium . Discussion — Combining the calculated 0.9% Breit correction with the reference Coulomb-correlated value, Eq.(5), one obtains the parity-nonconserving amplitude $`E_{\mathrm{P}NC}^{C+B}(^{133}\mathrm{Cs})=0.8991(36)\times 10^{11}i(Q_\mathrm{W}/N).`$A 0.4% theoretical uncertainty is assigned to the above result following the analysis of Ref. . Since the Breit interaction contributes at the 0.9% level to the total PNC amplitude, even a conservative 10% uncertainty in $`\delta E_{\mathrm{P}NC}`$ barely affects the accuracy of $`E_{\mathrm{PNC}}`$. When $`E_{\mathrm{P}NC}^{C+B}`$ is combined with the experimental values of the transition polarizability $`\beta `$ and $`E_{\mathrm{PNC}}/\beta `$ , one obtains for the weak charge: $`Q_\mathrm{W}(^{133}\mathrm{Cs})=72.65(28)_{\mathrm{expt}}(34)_{\mathrm{theor}}.`$This value differs from the prediction of the Standard Model $`Q_\mathrm{W}^{\mathrm{SM}}=73.20(13)`$ by 1.2$`\sigma `$, versus 2.5$`\sigma `$ of Ref. , where $`\sigma `$ is calculated by taking experimental and theoretical uncertainties in quadrature. This 1.2$`\sigma `$ deviation is slightly reduced further by taking into account corrections for the neutron nuclear distribution in $`{}_{}{}^{133}\mathrm{Cs}`$, estimated but not included in the final $`E_{\mathrm{PNC}}`$ of Ref. . Recently Pollock and Welliver determined the relevant modification to be $`\mathrm{\Delta }Q_\mathrm{W}^{\mathrm{SM}}=+0.11`$, which reduces the deviation from the Standard Model to 1.0 $`\sigma `$. The present calculation also provides a large Breit correction to the $`6S_{1/2}7P_{1/2}`$ electric-dipole matrix element. Using the ab initio all-order Coulomb-correlated value , $`6S_{1/2}D7P_{1/2}=0.279`$, and adding the 0.7% Breit correction of 0.0019, one finds $`6S_{1/2}D7P_{1/2}=0.281`$ in much better agreement with the 0.284(2) experimental value . The calculated Breit corrections bring most of the ab initio Coulomb-correlated hyperfine-structure constants for <sup>133</sup>Cs into 0.1% agreement with experimental values. To summarize, third-order many-body calculations of the contribution of the Breit interaction to the <sup>133</sup>Cs parity-nonconserving amplitude $`E_{\mathrm{P}NC}`$ and relevant atomic properties are reported. The difference between the present and the earlier calculations is due to additional inclusion of two-body Breit interaction, correlations, and the consistent use of experimental energies. The present analysis is a higher-order extension of my recent calculation . Since the major correlation effects are included, the present third-order analysis seems sufficient. The calculations reveal a 0.9% correction to $`E_{\mathrm{P}NC}`$ leading to a reduction to 1.2$`\sigma `$ of the recently reported 2.5$`\sigma `$ deviation of the <sup>133</sup>Cs weak charge from the Standard Model value. If corrections for the neutron distribution in <sup>133</sup>Cs nucleus are included, then the agreement between the atomic PNC in <sup>133</sup>Cs and the Standard Model stands at 1.0$`\sigma `$. Thus the result reported here brings the most accurate atomic PNC measurement to date into substantial agreement with the Standard Model. This work was supported by the U.S. Department of Energy, Division of Chemical Sciences, Office of Energy Research. Part of the work has been performed at Notre Dame University during a visit supported by NSF grant No. PHY-99-70666. Calculations were partially based on codes developed by Notre Dame group led by W.R. Johnson. The author is thankful to W.R. Johnson and M.S. Safronova for useful discussions and H.R. Sadeghpour for suggestions on the manuscript. Help with the manuscript and the stimulating interest of R.L. Walsworth is greatly appreciated.
warning/0005/hep-ph0005230.html	ar5iv	text	# Energy loss and 𝑥₂ scaling breakdown in 𝐽/𝜓 nuclear production ††footnote †Talk given at XXXVth Rencontres de Moriond: QCD and High Energy Hadronic Interactions, Les Arcs, France, 18-25 March 2000. ## 1 Introduction Charmonium production is one of the cleanest hard probes studied in high-energy hadron-hadron, hadron-nucleus, and nucleus-nucleus scatterings. If we look at the reaction in the center of mass frame, the parton model tells us that the production results from the fusion of a beam parton with energy-momentum fraction $`x_1`$ and a target parton with energy-momentum fraction $`x_2`$ which are converted into a $`c\overline{c}`$ pair that eventually gives the observed charmonium. In the absence of nuclear effects the integrated production rate in $`A_1`$$`A_2`$ collision would be given by $`A_1\times A_2`$ times that in $`p`$$`p`$ reaction. The deviation from this is often interpreted within a conventionnal absorption model. Recently the E866 collaboration provided the first evidence for a distinction between $`J/\psi `$ and $`\psi ^{}`$ in the differential production ratio $`^\mathrm{?}`$ $$R=\frac{\mathrm{d}\sigma (p+\mathrm{W}\psi +X)/\mathrm{d}x_F}{A\mathrm{d}\sigma (p+\mathrm{Be}\psi +X)/\mathrm{d}x_F}.$$ (1) The distinction occurs in the kinematical region where the charmonia have relatively small energies $`E_{c\overline{c}}`$ in the target rest frame. This is a definite departure from the observation of identical ratios for larger $`c\overline{c}`$ energies. It is tempting to describe this result by saying that for small $`c\overline{c}`$ energies the nucleus “acts as a detector” for the formation of the charmonium states. A possible interpretation of this is that for large $`c\overline{c}`$ energies a common precursor to $`J/\psi `$ and $`\psi ^{}`$ passes through the target and final-state formation takes place outside the nuclear volume. At low $`c\overline{c}`$ energies formation occurs within the target, leading to a weaker absorption for $`J/\psi `$ than $`\psi ^{}`$ due to the smaller size of the former. At the time E866 data became available we decided to take these ideas at face value and to test them in a simple but quantitative scenario of nuclear absorption $`^\mathrm{?}`$. A direct consequence of such an approach is an $`x_2`$ scaling. At high energy the charmonium energy amounts to $`E_{c\overline{c}}=x_1E`$ in the target rest frame where $`E`$ is the beam energy in this frame. In terms of the charmonium invariant mass $`M_{c\overline{c}}`$ one has $`E_{c\overline{c}}=M_{c\overline{c}}^2/(2m_px_2)`$. Consequently, if the production ratio is in actual facts driven by $`E_{c\overline{c}}`$ one expects $`R`$ to be independent of the beam energy at a given $`x_2`$ ($`x_2`$ scaling). Unlike what we expected, a comparison of $`J/\psi `$ differential production between 800 GeV (E866) $`^\mathrm{?}`$ and earlier less precise 200 GeV data (NA3) $`^\mathrm{?}`$ as a function of $`x_2`$ seems to indicate that $`c\overline{c}`$ production does not follow an $`x_2`$ scaling. However, it could happen that NA3 data reveal that at intermediate energies there is a mechanism involved in $`c\overline{c}`$ nuclear production beside the above mentionned nuclear suppression. In fact, it has been argued for a long time that an incoming parton might lose some energy while scattering through nuclear matter. Its momentum fraction is then shifted from $`x_1`$+$`\mathrm{\Delta }x_1`$ to $`x_1`$ at the point of fusion resulting in a modification of the effective beam parton density by a factor $`F(x_1+\mathrm{\Delta }x_1)/F(x_1)`$. Because the parton distributions $`F(x)`$ drop dramatically at intermediate and large $`x`$, the effective density of beam partons and consequently the production ratio might decrease significantly. The aim of these proceedings is to investigate whether energy loss could be a relevant mechanism to explain a breakdown of $`x_2`$ scaling in $`J/\psi `$ production. ## 2 Methods Since a corresponding quark energy loss may be expected in Drell-Yan (DY) pair production which is moreover not overwhelmed by final-state interactions, we first used the DY process to extract limits on parton energy loss. Energy loss is then combined with $`c\overline{c}`$ absorption to examine the consequences in the $`J/\psi `$ channel. To be more specific we chose to test two alternative energy loss scenarios. The first is that put forward by Gavin and Milana (GM) $`^\mathrm{?}`$. It predicts a beam parton shift in momentum fraction $$\mathrm{\Delta }x_1=\kappa _1x_1A^{1/3},$$ (2) where $`\kappa _1`$ is a free parameter.<sup></sup><sup></sup>$`\kappa _1`$ might depend on $`Q^2`$ and we needed its value in the region around $`m_\psi ^2`$. It turns out that the DY data are not precise enough and do not probe a sufficient range in $`Q^2`$ to make any definite statement about the latter dependence. We thus simply disregarded it. We stress that the shift depends on kinematics through $`x_1`$. The second model was proposed by Brodsky and Hoyer $`^\mathrm{?}`$ in order to fix a potential defect in GM. On general grounds they showed that $$\mathrm{\Delta }x_1\frac{\kappa _2}{s}A^{1/3},$$ (3) with $`\kappa _20.5`$ GeV<sup>2</sup>, an upper bound with which $`\mathrm{\Delta }x_1`$ as given by Eq. (2) is incompatible at large energies. For accessible energies it also sets an upper bound on $`\kappa _1`$. In relation with these considerations an alternative scenario, denoted BH in the following, is then obtained by assuming a shift in $`x_1`$ that saturates the bound given by (3), i.e., $$\mathrm{\Delta }x_1=\frac{\kappa _2}{s}A^{1/3},$$ where it should be noted here that the dependence on kinematics occurs via the center of mass energy. The free parameter in each model was fixed with DY data. The data set encompassed, on the one hand, the 150 and 280 GeV pion beam data on $`A=1`$ and 195 targets $`^\mathrm{?}`$ and, on the other hand, the 800 GeV proton beam data on $`A=9`$ and 184 $`^\mathrm{?}`$. At 800 GeV the data set spans a range in $`x_2`$ where shadowing is known to be sizeable. We thus used data points corrected for shadowing as given in Ref. 7. At SPS energies the data points lie in the (weak) antishadowing region and the corresponding correction is small. Calculations were done using MRST parton distributions $`^\mathrm{?}`$ and shadowing corrections were given by the EKS parameterization $`^\mathrm{?}`$. The models were then used to assess the possible importance of energy loss in $`J/\psi `$ production at 200 and 800 GeV. We assumed that $`J/\psi `$ production results from gluon fusion and took into account the color factor correction to energy loss : $`\mathrm{\Delta }x_g=\frac{9}{4}\mathrm{\Delta }x_q`$ At large $`x_1`$, $`q\overline{q}`$ annihilation into $`c\overline{c}`$ becomes predominant $`^\mathrm{?}`$. We checked that taking this channel into account did not change the overall results.. Subsequent energy loss was combined with nuclear absorption extracted from 800 GeV measurements and the result was compared to 200 GeV data $`^\mathrm{?}`$. In the $`J/\psi `$ sector shadowing corrections were not considered (see discussion at the end of Section 3). ## 3 Results and Discussion Drell-Yan production* In view to quantify parton energy loss, we first considered Drell-Yan dimuon production. At 800 GeV the wide kinematic acceptance ($`0.2x_10.95`$) and high statistics of E866/NuSea data make it possible to put stringent constraints on parton energy loss $`^\mathrm{?}`$. The fitted parameter values were $$\kappa _1=3\times 10^4(\chi ^2/\mathrm{ndf}=0.69),$$ for GM and $$\kappa _2=0.08\mathrm{GeV}^2(\chi ^2/\mathrm{ndf}=0.72),$$ for BH. Their 1$`\sigma `$ upper limits, i.e., $`\chi ^2/`$ndf$`=1`$, were $$\kappa _1+\mathrm{\Delta }\kappa _1=11\times 10^4,$$ and $$\kappa _2+\mathrm{\Delta }\kappa _2=0.77\mathrm{GeV}^2.$$ This suggested that the Drell-Yan data — corrected for shadowing — are consistent with zero energy loss within both models, but the 1$`\sigma `$ upper-limits imply that energy loss effects at a level of a few hundreds of MeV$`/`$fm cannot be excluded. Similarly, limits were estimated from a fit to the $`\pi ^{}`$ beam NA3 data at 150 and 280 GeV $`^\mathrm{?}`$ (Table 1). The lowest $`\chi ^2/`$ndf were reached with essentially no energy loss for both models — with or without inclusion of shadowing. Table 1 also indicates that the consideration of shadowing effects tended to improve the fits and correlatively to raise the upper limits of parton energy loss by a factor of four. Fitting the nuclear dependence of Drell-Yan production in $`p`$$`A`$ and $`\pi `$$`A`$ reactions leads to the conclusion that quark energy loss has to be small if any. Contrary to what we expected, the comparison between 200 and 800 GeV data did not allow us to discriminate between the different energy behaviours of GM and BH scenarios. A closer look at the data showed that this is mostly due to the large error bars in NA3 data. $`J/\psi `$ production In the parton model, production of $`c\overline{c}`$ pairs proceeds via parton fusion and might therefore be affected by initial-state effects (e.g., energy loss and shadowing), as well as final-state interactions. From the preliminary Drell-Yan study, we assumed $`\kappa _1`$ to be equal to 0.001 and $`\kappa _2`$ to be equal to 0.5 GeV<sup>2</sup> to evaluate the possible energy loss contribution to $`J/\psi `$ production. Figure 1 shows the $`J/\psi `$ production ratios $`R(\mathrm{Pt}/p)`$ at 200 GeV and $`R(\mathrm{W}/\mathrm{Be})`$ at 800 GeV. The most striking result is the large suppression seen at SPS according to BH mechanism whereas this suppression remains negligible at Fermilab energies. The $`J/\psi `$ production at 200 GeV is already affected ($`R`$ 0.9) at low $`x_1`$, and decreases when $`x_1`$ increases. The situation is different in the GM approach where the loss scales with $`x_1`$, independently of the beam energy. This can be seen in Figure 1 The difference between the two dotted curves in Figure 1 arises from the normalization with $`p`$ and Be targets at 200 GeV and 800 GeV respectively. which also demonstrates that the effect is moderate for $`\kappa _1=0.001`$ with a decrease from 1 at small $`x_1`$ to $`0.8`$ at $`x_1=0.7`$. The decrease seen in both approaches when going from small to large $`x_1`$ is the behaviour that might be reflected by NA3 $`J/\psi `$ data. Once energy loss in the $`J/\psi `$ sector was examined, $`J/\psi `$ suppression was calculated taking in addition nuclear absorption into account. The results are compared with experimental data $`^\mathrm{?}`$ on Figure 2 as a function of $`x_Fx_1x_2`$. First, Figure 2 reminds us that the model described in Ref. 2 for charmonium suppression (solid curve) is unable to describe large $`x_F`$ NA3 data. Second, the energy loss mechanism proved to reduce the disagreement with large $`x_F`$ data points. In particular, GM loss did not affect low $`x_F`$ production ratios and somehow decreased $`J/\psi `$ production at $`x_F0.5`$. By contrast, the BH scenario could almost reproduce the large suppression seen at large $`x_F`$ but tended to give too strong a suppression at small $`x_F`$. Though it is true that a first level comparison of both models would favor the GM approach, we should however stress that the results depend strongly on the values chosen for the parameters $`\kappa _1`$ and $`\kappa _2`$. This prompted us not to draw too quantitative conclusions. In summary, the results obtained may indicate that NA3 $`J/\psi `$ data result from a possible interplay between energy loss and nuclear suppression. In particular, the general trend of data is better reproduced when parton energy loss is taken into account. Consequently, the lack of $`x_2`$ scaling exhibited by NA3 and E866/Nusea data could be — at least partly — explained with such a mechanism. One limitation nevertheless arises from the remaining discrepancy between data and theoretical calculations. Further to that, we should mention that other mechanisms have been neglected in this study. In particular, gluon antishadowing calculated with EKS98 leads to an enhancement of $`J/\psi `$ production in platinum near $`x_2`$ 0.1 (i.e. $`x_F`$ 0.2 for NA3). It is hardly necessary to repeat that precise measurements of both Drell-Yan and charm differential production over a large kinematic window is mandatory to disentangle all the mechanisms involved in $`J/\psi `$ production. ## References
warning/0005/astro-ph0005545.html	ar5iv	text	# UV excess galaxies: Wolf-Rayet galaxies ## 1 INTRODUCTION Sullivan et al. (2000, hereafter S2000) published an ultraviolet (UV) flux-limited catalog of galaxies, imaged with the FOCA balloon-borne UV camera of Milliard et al. (1992). The catalog is an extension of work by Treyer et al. (1998), and contains 273 galaxies with redshifts $`z0.4`$. Vacuum UV imaging enables the direct detection of radiation from hot, massive stars and, by association, actively star forming galaxies. The luminosity function (LF) of the UV-selected galaxies is somewhat surprising. S2000 find a faint end slope in the Schechter (1976) parametrization of the LF, $`\alpha _{uv}=1.55\pm 0.11`$, steeper than typical optical LF faint end slopes (Folkes et al., 1999; Marzke et al., 1998; Geller et al., 1997, etc.). As a result, S2000 find no strong decrease in UV luminosity density between $`z0.4`$ and $`z=0`$ and thus little decline in the density of star formation over this range of cosmic epochs. In contrast, Madau et al. (1996) conclude that there is a factor of $`10`$ decline in the comoving density of star formation between redshifts $`z=1`$ and $`z=0`$. This conclusion is based on a combination of rest frame near-UV luminosities from the $`I`$-selected Canada France Redshift Survey (Lilly et al., 1995) and local H$`\alpha `$ measures from Gallego et al.’s (1995) objective prism survey. Curiously, a substantial fraction of the S2000 UV-selected galaxies have (UV$`B)`$ colors bluer than the bluest color that can be produced by starburst spectral synthesis models (Fioc & Rocca-Volmerange, 1997; Leitherer et al., 1999). S2000 could therefore not explain the color of these galaxies. Here we propose that the extremely blue galaxies are Wolf-Rayet (WR) galaxies. WR galaxies show direct signatures of WR stars in their spectra, usually a broad He II $`\lambda 4686`$ feature originating in the stellar winds of WR stars. This subtle feature is not seen in the fiber-spectra of S2000 (M. Sullivan 2000, private communication), but requires good S/N and is spatially dependent on the inclusion of starburst regions (Schaerer et al., 1999). The starburst regions dominate the broadband UV flux and, though short-lived ($`<10^7`$ yrs), WR stars can dominate the UV light of young starbursts. In this paper we describe optical photometry for 67% of the S2000 Selected Area 57 (SA57) UV-selected galaxy sample. Our goal is to compare the optical and UV colors, and to suggest that WR emission lines are the source of the extremely blue colors found by S2000. In section 2 we discuss our data, in section 3 we discuss the (V R) colors, and in section 4 we discuss the impact of Wolf-Rayet UV emission lines on the interpretation of the UV sample. We conclude in section 5. ## 2 DATA We obtained broadband $`V`$ and $`R`$ images in 1998 December and 1999 February with the MOSAIC camera on the Kitt Peak National Observatory 0.9m telescope. MOSAIC is an 8 CCD array with pixel scale 0.424 arcsec pixel<sup>-1</sup> and field of view $`59^{}`$ x $`59^{}`$ on the 0.9m. We obtained three offset 200 sec $`R`$ and 333 sec $`V`$ images at each position along a strip in RA. Our data serendipitously overlap with 148 of the 222 S2000 UV-selected galaxies in the SA57 region. We processed the images in the standard way with the IRAF MSCRED package. Total magnitudes are obtained using SExtractor (Bertin &Arnouts, 1996). Seven percent of the magnitudes are obtained with a modified version of IRAF-based GALPHOT (Freudling, 1993). GALPHOT allows photometry of galaxies poorly fit by SExtractor, including low surface brightness galaxies or galaxies with close companions. A comparison between the two software packages yields a 1-$`\sigma `$ scatter in the total galaxy magnitudes of $`\pm 0.05`$ mag. We use Landolt fields (Landolt, 1992) and the M67 cluster (Montgomery, Marschall, & Janes, 1993) for flux calibration. Depending on the nightly seeing conditions, the error in the standard star calibration is $`\pm 0.020.03`$ mag. A direct measure of our intrinsic error is the RMS scatter of measurements of identical galaxies taken on different nights: $`\pm 0.04`$ mag. The photometry includes galactic extinction corrections from Schlegel, Finkbeiner, & Davis (1998), who combine $`COBE`$/DIRBE and $`IRAS`$/ISSA data to map dust emission in the Milky Way. Typical extinction values are $`A_R=0.025`$ and $`A_V=0.03`$ mag. The photometry also includes “k corrections” for the changing rest frame bandpass (Frei & Gunn, 1994). Because we know the (V R) color and the redshift of each galaxy, we can determine accurate k corrections. Typical values for the $`V`$ and $`R`$ bandpasses are $`0.1<k_R<0.1`$ and $`0.15<k_V<0.15`$ mag. ### 2.1 Existing UV and $`B`$ measurements The S2000 UV measurements were made from a balloon-borne, 40 cm Cassegrain telescope stabilized to 2 arcsec RMS. The filter response approximates a Gaussian centered at 2015 Å, with a FWHM of 188 Å. The limiting magnitude in the UV catalog is m$`{}_{UV}{}^{}=18.5`$. S2000 state that the UV zero-point is accurate to $`0.2`$ mag; the uncertainty in the relative photometry may reach $`0.5`$ mag at the limiting magnitude. With the UV images in hand, S2000 searched for optical counterparts in digitized POSS plates. S2000 obtained accurate positions and $`B`$ photometry, with uncertainty $`\pm 0.20`$ mag, from the POSS plates. We follow S2000 and only use galaxies with a single optical counter-part when discussing UV colors. S2000 used the resulting (UV$`B)`$ colors to calculate k corrections. Finally, S2000 apply a reddening correction based on H$`\alpha `$ and H$`\beta `$ line strengths and the empirical Calzetti (1997) extinction law for starburst galaxies. They applied a mean reddening correction $`A_V=0.97`$ to galaxies lacking H$`\alpha `$ and/or H$`\beta `$ emission. The (UV$`B)_\mathrm{o}`$ colors we quote from S2000 include k corrections and reddening corrections. Our (V R) colors include k corrections but no reddening corrections, because S2000 do not provide enough information to reproduce their reddening corrections. Even if this reddening data were provided, 72% of the galaxies with (V R) colors lack both H$`\alpha `$ and H$`\beta `$ emission. For these galaxies, the mean reddening correction would add a uniform $`0.2`$ mag offset to the (V R) colors. ## 3 OPTICAL COLORS Figure 1 shows that the (V R) colors of the UV-selected galaxies are significantly bluer than those from an $`R`$-selected galaxy sample, the Century Survey (Geller et al., 1997). The Century Survey is a complete redshift survey of 102 square degrees to a limiting magnitude $`R=16.13`$, containing 1762 galaxies. The Century Survey galaxies have median $`(\mathrm{V R})_{CS}=0.58\pm 0.04`$; the UV-selected galaxies have median $`(\mathrm{V R})=0.38\pm 0.06`$. The distribution in Figure 1 appears double-peaked, but application of the non-parametric estimator of Pisani (1993) shows the second peak is not significant. The Pisani (1993) kernel probability estimator allows estimation of the probability density distribution underlying a data sample. In applying this we assume an average 1-$`\sigma `$ (V R) error $`\pm 0.066`$ mag. The S2000 (UV$`B)_\mathrm{o}`$ colors provide a measure of the star formation activity, but there is no correlation between the (V R) and (UV$`B)_\mathrm{o}`$ colors. Unlike the extreme blue (UV$`B)_\mathrm{o}`$ colors found by S2000, however, the optical (V R) colors are consistent with starburst galaxy colors. For example, the template starburst-galaxy spectral energy distributions SB1 and SB2 of Kinney et al. (1996) have synthetic $`(VR)=0.26`$. All but two of the S2000 UV-selected galaxies are redder than $`(VR)=0.23`$, and we note that these blue outliers appear atypically faint and point-like in our images. Thus $`95\%`$ of the UV galaxy (V R) colors are consistent with starburst galaxy colors. If the extreme blue galaxies were AGN, their average color should be $`(\mathrm{V R})=0.60`$ (Kotilainen et al., 1993), which is clearly not the case. ## 4 DISCUSSION An unresolved issue from the S2000 paper is the existence of galaxies with (UV$`B)_\mathrm{o}`$ colors $`2`$, colors too blue to be explained by standard spectral synthesis models. In our subsample of the S2000 UV-selected galaxies, 48 of 148 galaxies have (UV$`B)_\mathrm{o}<2`$, yet the optical (V R) colors show no significant departure from the rest of the sample. A KS test gives an 80% probability of the extreme (UV$`B)_\mathrm{o}<2`$ and the other UV-selected galaxies being drawn from the same distribution of (V R) colors. Next we examine the redshift distribution of the extremely blue galaxies. The lower panel of Figure 2 shows the redshift distribution of the extreme blue (UV$`B)_\mathrm{o}<2`$ galaxies, and the upper panel the remaining (UV$`B)_\mathrm{o}>2`$ galaxies. Figure 2 shows that the extreme blue galaxies are preferentially found at high redshift. A KS-test yields a $`10^{11}`$ probability of the two populations being drawn from the same distribution. The peaks in the redshift histograms (Figure 2) are well matched to the $`z0.3`$ limit of the Medium Deep Survey (Willmer et al., 1996), which includes the SA57 region. The Medium-Deep Survey is $`70\%`$ complete to $`B_J20.5`$ mag, and gives a sense of the existing large scale structure in the direction of UV survey. Thus the extremely blue galaxies appear to trace the same large-scale structure as “normal” galaxies. ### 4.1 Wolf-Rayet galaxies. The very blue (UV$`B)_\mathrm{o}<2`$ galaxies are an extreme population of galaxies found mostly at larger redshift. An explanation probably requires emission lines redshifting into the UV bandpass. Because S2000 and our $`(\mathrm{V R})`$ color distribution rule out AGN, it seems plausible the emission lines come from a stellar emission line source, i.e. Wolf-Rayet (WR) stars, planetary nebulae, or symbiotic stars. Because there are known galaxies with WR features, we suggest that the extreme blue (UV$`B)_\mathrm{o}<2`$ galaxies are WR galaxies. The strongest UV features come from ionized C, O, He, Si, and Fe in the stellar winds of WR stars (e.g. Hillier & Miller 1999). Our optical spectra of 10 extreme blue galaxies all show narrow emission lines indicative of active star formation. WR galaxies are characterized by the direct signatures of WR stars in their spectra, but our S/N is too poor to detect a possible $`\lambda 4686`$ WR feature convincingly. To calculate a typical (V R) color for WR galaxies, we searched the NASA/IPAC Extragalactic Database for photometery of known WR galaxies. The average optical color of 13 WR galaxies is $`\overline{(\mathrm{V R})}=0.27\pm 0.34`$ mag. This average color is in agreement with the blue color we find for the (UV$`B)_\mathrm{o}<2`$ sample, $`\overline{(\mathrm{V R})}=0.35\pm 0.11`$ mag. Figure 3 illustrates how WR emission lines are redshifted into the S2000 UV bandpass. We plot a sample WR ultraviolet spectrum, WC star HD 165763 (Hillier & Miller, 1999), at redshift $`z=0`$ and at redshift $`z=0.285`$ with the S2000 UV and $`B`$ bandpasses overlaid. Integration of the UV bandpass over the redshifted spectrum shows that the WR emission lines can boost the UV flux by a factor of $`3`$ compared to the flux at $`z=0`$. We also integrate the $`B`$ bandpass over the Hillier & Miller (1999) WC stellar spectrum, and find the (UV$`B)`$ color is shifted up to 1.3 mag bluewards over the redshift range z=0 to z=0.5 (see Figure 4), blue enough to explain $`2/3`$ of the extreme blue (UV$`B)_\mathrm{o}`$ galaxies in the SA57 UV-selected sample. The WC stellar spectrum represents a maximal effect; known WR galaxies have typical WC/WN ratios of 0.2 - 0.4 (Schaerer et al., 1999). We perform the (UV$`B)`$ calculation on the WN stellar spectrum, WR152 (Hamann & Koesterke, 1998), and find the (UV$`B)`$ color shifts up to 1.1 mag bluewards over the same redshift range (see Figure 4). We acknowledge that strong UV emission lines are not common features of known low redshift WR galaxies (Leitherer 2000, private communication). We calculate the shifts in (UV$`B)`$ for a dozen of the McQuade et al. (1995) “UV to optical spectral distributions of northern star-forming galaxies,” and find an average $`\mathrm{\Delta }(\mathrm{UV}B)0.3`$ mag (mostly due to the continua, not the emission lines). Yet good examples of the very blue objects may not be found locally. The UV selection should pick out extreme objects; known WR galaxies are optically selected. We note the very blue UV galaxies have median z=0.2. By comparison, the Schaerer et al. (1999b) catalog of known WR galaxies drops off sharply at z=0.02. We can also take advantage of $`HST`$ FOS observations of the Small Magellanic Cloud H II region N88A (Kurt et al., 1999) as a UV spectral templates for WR galaxies. When the S2000 bandpass is centered on the strong CIII and SiIII lines in this object, the UV flux doubles compared to the surrounding continuum level. S2000 compute their luminosity function (LF) with the full UV sample; here we construct a model LF for the WR galaxies alone. The impact of the emission lines on the observed UV flux changes with redshift. We explore the selection function for finding WR galaxies in a magnitude limited redshift survey by computing the change in absolute magnitude resulting from the UV emission lines. WR galaxies are blue and show active star formation. We thus use the Second Southern Sky Redshift Survey (Marzke et al., 1998, SSRS2) LF for the bluest galaxies, the Magellanic spiral and irregulars, as the basis of our model. The SSRS2 is a $`B`$-selected survey complete to $`m_b15.5`$ that contains 5404 galaxies, all with morphological classifications. The SSRS2 Sm-Irr LF has a steep faint-end slope, $`\alpha =1.81`$, similar to the S2000 LF, $`\alpha =1.55`$. We compute the selection function for the WR galaxies according to equation (1), $$N(<M_{lim})_zD_C(z)^2𝑑D_C(z)_{M=\mathrm{}}^{M_{lim}(z)}\varphi (M,z)𝑑M$$ (1) where $`D_C(z)`$ is the co-moving distance. We assume a flat cosmology and $`H_o`$=100 km s<sup>-1</sup> Mpc<sup>-1</sup>, $`q_o`$=0.5. $`\varphi (M,z)`$ is the luminosity function in the Schechter (1976) parameterization, $`\varphi (M)=0.4ln10\varphi ^{}`$ $`[10^{0.4(M^{}M)}]^{(1+\alpha )}`$ $`exp(10^{0.4(M^{}M)})`$, where $`\varphi ^{}`$ is the normalization (galaxies per unit volume), $`\alpha `$ is the faint end slope, and $`M^{}`$ is the characteristic absolute magnitude. To convert the SSRS2 Sm-Irr $`M_B^{}=19.78`$ to the UV bandpass, we add the median (UV$`B)_\mathrm{o}=0.91`$ color for the UV galaxies. To simulate WR galaxies, we add the relative change in (UV$`B)`$ over redshift (Figure 4): $$M_{\mathrm{UV}}^{}=M_B^{}+(\mathrm{UV}B)_\mathrm{o}+\mathrm{\Delta }(\mathrm{UV}B)[z]$$ (2) To find the number of galaxies in redshift bin $`(z,z+\delta z)`$, we integrate $`\varphi (M,z)`$ to the limiting absolute magnitude $`M_{lim}(z)`$ over the volume of the redshift interval $`z`$ to $`z+\delta z`$. The selection functions represent the number of galaxies expected per unit redshift if galaxies are uniformly distributed. The curves in Figure 2 are the UV selection functions normalized to the observed number of galaxies. The selection function for the S2000 LF ($`M^{}=20.59,\alpha =1.51`$) is very similar to the SSRS2 Sm-Irr LF ($`M^{}=19.78+(`$UV$`B)_\mathrm{o},`$ $`\alpha =1.81`$), and matches the (UV$`B)_\mathrm{o}>2`$ galaxy redshift distribution (Figure 2, upper panel) quite well. In contrast, the UV selection function we calculate for WR galaxies predicts more WR galaxies at large redshifts. The redshift distribution of (UV$`B)_\mathrm{o}<2`$ galaxies (Figure 2, lower panel) matches the WR selection function curves much better than the S2000 curve. The shifts we make in $`M^{}`$ reflect only the change in the flux through the UV bandpass as a function of redshift. Thus the UV emission lines of WR galaxies can explain the extreme blue (UV$`B)_\mathrm{o}`$ colors and the redshift distribution of the (UV$`B)_\mathrm{o}<2`$ galaxies. We calculate the number density of Century Survey (which includes the SA57 region) galaxies to the same effective limiting magnitude and solid angle as the UV sample, and find that the extreme blue objects are $``$1% of a “normal” ($`R`$-selected) galaxy population. It is difficult to compare the known WR galaxy fraction. We estimate a lower limit by assuming the WR catalog of Schaerer et al. (1999b) is complete to $`B`$=14 mag over the northern sky above 15 deg galactic latitude. The resulting WR galaxy fraction lower limit is 0.75%. Thus the extreme blue galaxy number density is consistent with the abundance of known WR galaxies. ## 5 CONCLUSIONS We show that the extreme blue (UV$`B)_\mathrm{o}<2`$ galaxies observed by S2000 may be WR galaxies, starburst galaxies with strong UV emission from WR stars. The model explains the observed (UV$`B)_\mathrm{o}<2`$ colors and the redshift distribution for the bluest objects. The average (V R) color of 13 known WR galaxies taken from the literature, $`\overline{(\mathrm{V R})}=0.27\pm 0.34`$, is consistent with the range of (V R) colors we observe for the extreme blue (UV$`B)_\mathrm{o}<2`$ galaxies. We measure a median optical color $`(\mathrm{V R})=0.38\pm 0.06`$ for 67% of the S2000 SA57, UV-selected sample. The UV emission lines can double or triple the flux through the S2000 UV bandpass at large redshift. As a result, the (UV$`B)`$ color of a WR galaxy can become up to 1.3 mag bluer at larger redshift. This changing (UV$`B)`$ as a function of redshift predicts a selection function weighted toward large redshifts. Not only do we find an excess of (UV$`B)_\mathrm{o}<2`$ galaxies at large redshift, compared to the (UV$`B)_\mathrm{o}>2`$ sample, but our model matches the observed redshift distribution well. We plan to test this model by obtaining high quality spectra of the suspected WR galaxies and looking for the telltale broad WR emission features. We thank R. Ellis, M. Sullivan, A. Barth, C. Leitherer, and the referee, H. Ferguson, for their thoughtful comments and careful reading of the paper. We thank A. Mahdavi for the use of his “kernel” code to do the Pisani (1993) statistical calculation. This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with NASA. This research was supported by the Harvard-Smithsonian Center for Astrophysics.
warning/0005/cond-mat0005161.html	ar5iv	text	# Geometrical Approach for the Mean-Field Dynamics of a Particle in a Short Range Correlated Random Potential ## I Introduction Understanding of the out-of-equilibrium dynamics of glassy systems, including spin glasses, structural glasses, supraconductor vortex glasses, etc, is a challenging problem. The need for exact, but non trivial results led to the introduction of ideal spin-glasses, like the celebrated Sherrington-Kirkpatrick model . The spherical $`p`$-spin glass is such a model, where a closed set of equations exists for the time correlation functions of the dynamical variables, or soft spins . In the thermodynamical limit, each spin becomes coupled to an infinity of other spins. Each variable experiences gaussian fluctuations in the effective environment created by all the other spins. The dynamics then simplifies drastically, and reduces to a set of correlation functions which have to be determined self-consistently. The result is the “dynamical mean-field solution” of the model. The mean-field dynamics of spin-glasses has revealed extremely rich, the most striking feature being the existence of a non-trivial aging relaxation regime at low temperature . For instance, the solution of the mean-field equations in this out-of-equilibrium regime demonstrate the existence of a generalised fluctuation-dissipation theorem (i.e. connecting correlation and response functions) whose validity seems now to extend to many realistic, non mean-field, models . Then, mean-field solutions are valuable for explaining the experimental aging of disordered magnetic systems . Finally, a sustained interest has followed the discovery of a deep formal analogy between the mode coupling description of structural glasses (supercooled liquids) and the mean field treatment of the spherical $`p`$-spin glass . A crucial shortcoming of the mean-field description, however, is its inability to take into account properly thermally activated motion over energy barriers, leading to a sharp dynamical transition – divergence of an internal relaxational time scale – whereas the corresponding “finite dimensional” behaviour is only a strong but progressive slowing down of the dynamics. Despite of this last point, mean-field dynamics remains a major issue in the study of out-of-equilibrium statistical physics of disordered systems, and any approach providing a physical insight on its aging mechanism is of interest. A major step in that direction was made by J.Kurchan and L.Laloux who investigated the zero temperature relaxation of systems including ferromagnets and spin glasses. The zero temperature limit makes it possible to consider the energy landscape, rather than a ill-defined “free-energy” landscape, without reducing the dynamics to anything trivial. In this work, we extend further their approach, and apply it to another system of interest: the mean-field dynamics of a particle in a short-range correlated random potential. The out-of-equilibrium, aging dynamics of this model has been first studied in , and thoroughly investigated in . Its glassy behaviour belongs to the same universality class than the spherical $`p`$-spin model. What makes this model interesting is its natural extension, when a finite and constant external force is applied to the particle. Then, it becomes a paradigm of “driven glassy system”, in which a non-linear response to the force as well as a significant violation of the fluctuation-dissipation theorem are expected, as shown by Horner . In this paper, we present the dynamical mean-field equations in the presence of a constant force, allowing for an arbitrary mean displacement along this one. These equations are then numerically solved, in the zero temperature limit, for an exponentially decreasing correlator, in the absence, and in the presence of a weak external force. The corresponding numerical results are presented, restricting ourselves to the linear response regime. Then, we start our geometrical analysis of the zero temperature relaxation by a simple random matrix calculation that we believe to describe satisfactorily the hessian of the exponentially correlated gaussian potential. Next, we perform an “instantaneous normal mode” analysis of the relaxational motion. The key observable turns out to be the (intensive) energy difference between the energy $`(t)`$ at a given time $`t`$, and its asymptotic value $`lim_t\mathrm{}(t)`$. We subsequently analyse the waiting time dependence of two characteristic time scales $`t_f,t_b`$, that we relate to $`(t)(\mathrm{})`$. This work is preliminary to the study of the stationary driven situation in the presence of a finite force, which will be the subject of a forthcoming publication, and where the velocity-force characteristics, and the cross-over between linear and non-linear response will be exposed . ## II The out-of-equilibrium dynamics in the mean-field approximation We introduce in this section the mean-field dynamics equations and discuss the low temperature aging solution in the absence of force. Let $`𝐱(t)`$ be the position of a particle, obeying a usual Langevin equation: $$\dot{𝐱}(t)=\mathbf{}V(𝐱(t))++𝜻(t),$$ (1) where are introduced the random potential $`V(𝐱)`$, the external force $``$, a white Langevin noise $`𝜻(t)`$ corresponding to a temperature $`T`$, and a friction coefficient equal to 1. Quantities $`𝐱,\mathbf{},,𝜻`$ are $`N`$-dimensional vectors. Three sub-cases of the dynamics defined by (1) are of interest: 1/ the “isolated” dynamics, without force: $`=0`$; 2/ the driven relaxational dynamics, which is the zero temperature limit of (1): $`0`$ and $`T0`$; 3/ the relaxational “isolated” dynamics: $`=0,T0`$. The potential $`V(𝐱)`$ is a quenched disorder, chosen from a gaussian distribution. All the averages with respect to it will be denoted with an over-line $`\overline{}`$, while the average over the thermal noise (if any) $`\zeta `$ will be denoted by the brackets $``$. We suppose that the motion starts at $`t=0`$ and $`𝐱(t=0)=0`$. After averaging over the quenched disorder, this choice becomes equivalent to start with a random, “infinite temperature” distribution of initial positions. We expect that the process (1) is self-averaging with respect to $`V(𝐱)`$ in the infinite dimensional limit. One introduces the correlator $`f(y)`$ of the gaussian disorder, explicitly dependent on the dimension $`N`$ of the configuration space $`\{𝐱\}`$. $$\overline{V(𝐱)V(𝐱^{})}=Nf\left(\frac{𝐱𝐱^{}^2}{N}\right);\overline{V(𝐱)}=0.$$ (2) This form ensures a meaningful $`N\mathrm{}`$ limit, in which each coordinate $`𝐱_i(t)`$, or gradient component $`_iV(𝐱)`$, remains of order 1, while the norms $`𝐱(t)`$, $`\mathbf{}V`$ scale like $`N^{1/2}`$. As a consequence, the external force must scale ($`𝐞_1`$ being a unit vector) like: $$=N^{1/2}F𝐞_1.$$ (3) One expects a displacement $`\overline{𝐱(t)}=N^{1/2}u(t)𝐞_1`$, and possibly a mean velocity $`\overline{\text{d}𝐱(t)/\text{d}t}=N^{1/2}v𝐞_1`$. From now onwards, we arrange that $`𝐞_1`$ coincides with the first coordinate axis $`i=1`$. In the present paper, we restrict ourselves to the exponentially correlated potential: $$f(y)=\mathrm{exp}(y).$$ (4) This is a special case of short range correlated random potential, characterised by $`lim_y\mathrm{}f(y)<\mathrm{}`$. The average difference $`\overline{[V(𝐱)V(𝐱^{})]^2}`$ is bounded when $`𝐱𝐱^{}`$ grows, and this ensures the existence of a normal diffusion regime at temperatures high enough. Another common choice is the power-law correlator: $`f(y)=2/(\gamma 1)(1+y)^{(1\gamma )/2};\gamma >1`$ . Choice (4) is also a particular case of $`f(y)=U_p^2\mathrm{exp}(y/\xi ^2)`$, with a pinning energy $`U_p`$ and correlation length $`\xi `$ set to 1, thanks to a simple rescaling, without loss of generality. The Langevin dynamics is handled with the help of a Martin-Siggia-Rose (MSR)-like functional integral, convenient for averaging over the gaussian disorder . All the technical details corresponding to the saddle-point equations as $`N\mathrm{}`$, are given in and the action is given in appendix A. The crucial point is that the limit $`N\mathrm{}`$ is taken first, before any other limit $`T0`$ or $`t\mathrm{}`$. As a result, we obtain a general effective quadratic action $`S[x_j(t),i\stackrel{~}{x}_j(t)]`$, involving the original field $`x_j(t)`$, and the MSR auxiliary field $`i\stackrel{~}{x}_j(t)`$. Three among the four following correlation functions appear explicitly in the action $`S[x_j(t),i\stackrel{~}{x}_j(t)]`$: $`u(t)`$ $`=`$ $`N^{1/2}\overline{x_1(t)};`$ (5) $`r(t,t^{})`$ $`=`$ $`N^1{\displaystyle \underset{j=1}{\overset{N}{}}}\overline{x_j(t)i\stackrel{~}{x}_j(t^{})};`$ (6) $`b(t,t^{})`$ $`=`$ $`N^1{\displaystyle \underset{j=2}{\overset{N}{}}}\overline{(x_j(t)x_j(t^{}))^2};`$ (7) $`d(t,t^{})`$ $`=`$ $`N^1{\displaystyle \underset{j=1}{\overset{N}{}}}\overline{(x_j(t)x_j(t^{}))^2};`$ (8) $`=`$ $`b(t,t^{})+[u(t)u(t^{})]^2.`$ (9) These are the displacement $`u(t)`$, the response function $`r(t,t^{})`$, and the correlation functions $`b(t,t^{})`$ and $`d(t,t^{})`$. The Dyson equations for $`r,b,d,u`$ form a closed system of coupled integro-differential equations. For $`tt^{}`$ one has to solve: $`_tr(t,t^{})`$ $`=`$ $`\delta (tt^{})`$ (11) $`{\displaystyle _0^t}\text{d}s4f^{\prime \prime }(d(t,s))r(t,s)[r(t,t^{})r(s,t^{})];`$ $`_tb(t,t^{})`$ $`=`$ $`2T{\displaystyle _0^t}\text{d}s4f^{}(d(t,s))[r(t,s)r(t^{},s)]`$ (13) $`{\displaystyle _0^t}\text{d}s4f^{\prime \prime }(d(t,s))r(t,s)[b(t,s)+b(t,t^{})b(s,t^{})];`$ $`_tu(t)`$ $`=`$ $`F{\displaystyle _0^t}\text{d}s4f^{\prime \prime }(d(t,s))r(t,s)[u(t)u(s)].`$ (14) Equations (11-LABEL:eq:Dyson-2temps:u) are original ones, and allow for a non uniform displacement $`u(t)`$. The aging, isolated, situation corresponds to the limit $`u(t)=F=0`$, and $`d(t,t^{})b(t,t^{})`$ in the above system. The stationary limit, investigated by Horner amounts to write $`r(t,t^{})=R(tt^{})`$, $`b(t,t^{})=B(\|tt^{}\|)`$, $`u(t)=vt`$, and to reject the lower bound of the time integrals $`\text{d}t\text{d}s`$ to $`\mathrm{}`$. Three relevant observables: the energy $`(t)`$, the curvature $`(t)`$ and the pinning force $`F_p(t)`$ can be expressed with the help of these correlation functions. $`(t)`$ $`=`$ $`N^1\overline{V(𝐱(t))},`$ (16) $`=`$ $`{\displaystyle _0^t}\text{d}s2f^{}(d(t,s))r(t,s);`$ (17) $`(t)`$ $`=`$ $`N^1{\displaystyle \underset{j=1}{\overset{N}{}}}\overline{_{jj}^2V(𝐱(t))},`$ (18) $`=`$ $`{\displaystyle _0^t}\text{d}s4f^{\prime \prime }(d(t,s))r(t,s);`$ (19) $`F_p(t)`$ $`=`$ $`N^{1/2}\overline{_1V(𝐱(t))},`$ (20) $`=`$ $`{\displaystyle _0^t}\text{d}s4f^{\prime \prime }(d(t,s))r(t,s)[u(t)u(s)].`$ (21) The pinning force is such that $`\overline{\text{d}𝐱(t)/\text{d}t}=_p(t)+`$ with $`_p(t)=N^{1/2}F_p(t)𝐞_1`$. The pinning force $`F_p`$ and the driving force $`F`$ have opposite signs. A proper study of the mean-field equilibrium phase diagram requires an extra quadratic confinement potential $`\mu 𝐱^{\mathrm{\hspace{0.17em}2}}/2`$. This ensures the existence of a true thermal equilibrium in the high temperature phase, while correlation functions reach their asymptotic values exponentially fast. Then, a transition line $`T_d(\mu )`$, called dynamical temperature, separates the high temperature ergodic phase, from a low temperature, aging and non ergodic phase . At high temperature, the system reaches a true stationary state, and the dynamics becomes time-translationally invariant (TTI), i.e the 2-times correlation functions depend only on the difference $`tt^{}`$, while the 1-time expectation values are constant. In this stationary situation, it is convenient to introduce the TTI correlation functions $`B(tt^{})=lim_{t,t^{}\mathrm{}}b(t,t^{})\|_{tt^{}\mathrm{finite}};R(tt^{})=lim_{t,t^{}\mathrm{}}r(t,t^{})\|_{tt^{}\mathrm{finite}}`$. The fluctuation-dissipation theorem (FDT) holds and reads: $$\text{d}B(t)/\text{d}t=2TR(t).$$ (22) As a consequence, the equal-time correlation functions coincide with their thermodynamical (canonical ensemble) counterparts. Taking the limit $`\mu 0`$ does not lead to any singular result . Provided the contribution from the harmonic potential has been subtracted off, the energy $`(t)`$ behaves smoothly as $`\mu `$ tends to 0. When $`\mu `$ exactly equals 0, the system cannot be at equilibrium, and instead, one expects a long time behaviour corresponding to a normal diffusion situation, with a finite diffusivity $`D=lim_t\mathrm{}𝐱^{\mathrm{\hspace{0.17em}2}}(t)/(2Nt)`$, a finite mobility $`\eta ^1=lim_t\mathrm{}u(t)/(Ft)`$, and the Einstein relation $`D=T\eta ^1`$. Kinzelbach and Horner described the dynamics in the stationary, high temperature phase . They found that these correlation functions behave in the same way than those of the well known mode-coupling theories for supercooled liquids, as expected on general grounds . The non-linearities of the self-consistent equations cause a dramatic slowing down of the dynamics as $`T_d`$ is approached from above, leading to a sharp transition at $`T=T_d`$. For instance, the function $`B(t)`$, after a fast increase at short times $`t1`$, has a long plateau near a characteristic value $`B(tt_f)q`$, before eventually reaching its asymptotic, long time regime $`B(t)=\widehat{B}(t/t_b)`$. Both $`t_f`$ and $`t_b`$ diverge like power laws of the difference $`\|TT_d\|`$ . The low temperature region however corresponds to an out-of-equilibrium situation. In the absence of external force, this is meant by the loss of both time-translational invariance (TTI) and fluctuation-dissipation theorem (FDT). The 2-time correlation functions cannot be reduced any more to functions of the time differences $`tt^{}`$, and there is a domain in the $`(t,t^{})`$ plane, where the system ages . The addition of a weak, constant external force leads to a somewhat different picture. As proposed by Horner, the system is expected to reach a stationary state (TTI), but the FDT remains definitively lost . It turns out, however (cf next section) , that when the force is switched on at a time $`t=0`$, there is a finite time interval during which the dynamics can be successfully described as a perturbation around the aging isolated ($`F=0`$) regime, with a linear response approach. The extent of this linear response regime is inversely related to the magnitude of the force. The aging dynamics of the isolated particle has been exhaustively treated in . The fluctuation-dissipation theorem is violated and must be replaced by: $$X(t,t^{})_t^{}b(t,t^{})=r(t,t^{}).$$ (23) In the time sector $`tt^{}`$ finite; $`t^{}\mathrm{}`$ of the $`(t,t^{})`$ plane, the behaviour is very similar to the one observed just above $`T_d`$, and the value of $`X(t,t^{})`$ is very close to its equilibrium value $`1/(2T)`$. When the time separation $`tt^{}`$ ceases to be small relatively to a characteristic time $`t_f(t^{})`$ which has to be determined, $`X(t,t^{})`$ departs from its equilibrium value, decreasing its magnitude $`\|X\|`$. The analytical study of the equations (9) has only been possible in the asymptotic limit $`t,t^{}\mathrm{}`$, by dropping out sub-leading terms presumably of order $`1/t,1/t^{}`$. In this limit, the authors have shown (this is the crucial point) that it was possible to parametrise the dynamics with the help of the correlation function $`b(t,t^{})`$ of the system itself. This implies that $`X(t,t^{})`$ becomes a one variable function $`X[b(t,t^{})]`$, and it turns out that all short range correlated models can be solved thanks to the ansatz : $`b(t,t^{})<q`$ $``$ $`X[b]=1/2T;`$ (24) $`b(t,t^{})>q`$ $``$ $`X[b]=\chi .`$ (25) This extension of the FDT is called “quasi-fluctuation-dissipation theorem” (QFDT). In this paper, we rather use the function $`\overline{T}(t,t^{})=1/(2X(t,t^{}))`$. In the aging regime, this effective temperature $`\overline{T}=1/(2\chi )`$ is higher than the thermostat temperature $`T`$, and remains finite in the limit $`T0`$. The physical meaning of these “two or many temperatures systems” is discussed in . In the case we are interested in, i.e. in the absence of confinement ($`\mu 0`$), for a correlator $`f(y)=\mathrm{exp}(y)`$, $`\chi `$ and $`q`$ are for any temperature $`T<T_d`$, solutions of the system : $$\{\begin{array}{ccccc}T& =& q\sqrt{f^{\prime \prime }(q)}& =& qe^{q/2};\\ \chi & =& \frac{\sqrt{f^{\prime \prime }(q)}}{2f^{}(q)}& =& e^{q/2}/2;\end{array}$$ (26) and in the low temperature limit: $$\{\begin{array}{ccc}q& & T+T^2/2\mathrm{};\\ \chi & & \left(\frac{1}{2}+\frac{T}{4}+\frac{3T^2}{16}\mathrm{}\right)\end{array}$$ (27) In the same way, given a triplet $`t_1<t_2<t_3`$, in the time domain where $`(t_1,t_2,t_3)\mathrm{}`$ and $`t_1/t_2`$, $`t_2/t_3`$ finite, $`b(t_3,t_2)`$ is uniquely determined by the knowledge of $`b(t_2,t_1)`$ and $`b(t_3,t_1)`$. Again, the explicit dependence can be carried out exactly when the correlator is exponential. The result is : $$b(t_3,t_2)q=b(t_3,t_1)q[b(t_2,t_1)q].$$ (28) A well known shortcoming of this approach, is that any reference to the original times $`t,t^{}`$ is definitively lost. The asymptotic solution cannot distinguish between $`b(t,t^{})`$ and $`b(h(t),h(t^{}))`$ where $`th(t)`$ can be any suitable reparametrization of the time variable. As a by-product, the previous analysis predicts only the more general form of the solution, in the aging regime $`t/t^{}1`$: $$b(t,t^{})=\stackrel{~}{B}\left[\mathrm{ln}\left(\frac{h(t)}{h(t^{})}\right)\right]+q.$$ (29) For exponentially correlated potentials, the master function is known , and without loss of generality: $$b(t,t^{})=\mathrm{ln}(h(t))\mathrm{ln}(h(t^{}))+q.$$ (30) In is made the conjecture $`h(t)=t^\delta `$, compatible with the results found below. In what follows, we will refer to this solution as the time-reparametrization invariant (TRI) solution. At the beginning of the aging regime, for $`t`$ and $`t^{}`$ such that $`(tt^{})/t^{}`$ is finite but small compared to $`1`$, the scaling form (29), reads: $$b(t,t^{})=\stackrel{~}{B}\left((tt^{})t_b^1(t^{})+\mathrm{}\right)+q.$$ (31) Here, $`t_b(t^{})`$ is the characteristic time of the aging regime, defined by $`t_b(t^{})=h(t^{})/h^{}(t^{})`$. This is the typical time needed by the particle for diffusing over a distance $`b(t,t^{})q1`$. Non exponential correlators have a non analytic scaling function $`\stackrel{~}{B}(m)`$ around $`m=0`$ and the r.h.s of (31) is singular in $`tt^{}`$ . The time-reparametrization invariant solution describes a situation where the time scales for the FDT regime ($`t_0=1`$) and for the aging regime ($`t_b(t^{})`$) are well separated (i.e. $`t_b(t^{})1`$), which implicitly assumes $`t,t^{}\mathrm{}`$. In order to go further, one needs to take into account the times derivatives $`_t,_t^{}`$ neglected in the asymptotic regime of the TRI solution. It is enough, in principle, to fix up the reparametrization function $`h(t)`$. Moreover, the TRI solution does not say how the parameter $`\overline{T}`$ goes from its FDT value ($`b(t,t^{})<q`$) to its QFDT value ($`b(t,t^{})>q`$). One defines for this purpose the new time scale $`t_f(t)`$, such that, for instance, $`\overline{T}(t,tt_f(t))`$ takes a given value between $`T`$ and $`1/(2\chi )`$. We shall see below that $`t_f(t)`$ is much smaller than $`t_b(t)`$. ## III Results from the numerical integration of mean-field equations The mean-field equations, with $`F`$ and $`u(t)`$ equal to zero were first numerically integrated by Franz and Mézard . The quadrature scheme is of order one in the time step $`h`$, but reveals itself surprisingly robust as $`h`$ is increased up to value as large as $`0.3`$. The authors of report being able to reach $`t1000`$ at the best. Our investigations have shown that the quality of our solutions gets worst if $`h`$ in increased above $`0.2`$, and we present results up to $`t400`$. For reasons detailed in the next section, we have only considered the exponential correlator case (4). We set $`T`$ to 0 in (9) and took the initial value $`C(0,0)=0`$. The information coming from the numerics may be pigeonholed in three categories. 1/ Results related to the TRI solution. First of all, we must check that the quasi-fluctuation dissipation relation (24) is true by plotting the integrated response versus the correlation function, on Figure (1). The observed value of $`\chi `$ is close to 0.46, while the predicted value is $`1/2`$. The TRI solution predicts also $`q0`$ and $`lim_t\mathrm{}b(t,0)=\mathrm{}`$, in the absence of confinement. The measured asymptotic energy $`(\mathrm{})`$ and mean curvature $`(\mathrm{})`$ are found to be in excellent agreement with the predicted values $`2`$ and $`+4`$ respectively. 2/ Beyond the TRI solution, without external force. This includes for instance the algebraic decay of the energy $`(t)=2+c_1t^\kappa `$. The exponent $`\kappa `$ is determined by plotting $`\mathrm{log}(2+(t))`$ versus $`\mathrm{log}(t)`$, and also by computing directly the logarithmic derivative, as shown on Figure (2). The exponent $`\kappa `$ lies between $`0.66`$ and $`0.67`$ and our best estimate is $`c_1=1.08`$. Also concerned are the characteristic times of the aging regime, and the precise nature of the cross-over from equilibrium to quasi-equilibrium fluctuation dissipation theorem. We are interested here in finding the characteristic time $`t_f(t)`$ as a function of $`t`$, defined by: $$_{tt_f}^t\text{d}s2f^{}(b(t,s))r(t,s)=1,$$ (32) or alternatively: $$_{tt_f^{}}^t\text{d}sr(t,s)=\frac{1}{2f^{}(0)}=\frac{1}{2}.$$ (33) Equation (32) comes from the fact that the equilibrium $`X=1/2T`$ and aging $`X=\chi `$ time sectors contribute for $`1`$ each to the energy. The value $`t_f`$ which solves (32) separates the equilibrium regime ($`b(t,s)<b(t,tt_f)`$) from the aging one ($`b(t,s)>b(t,tt_f)`$). The equivalence between (32) and (33) is a straightforward consequence of (24) One generalises (33) in: $$_{tt_a}^t\text{d}sr(t,s)=a,$$ (34) For $`a<1/2`$, $`t_a`$ must tend to a constant as $`t\mathrm{}`$, while for $`a>1/2`$, the asymptotic scaling (29) predicts: $`a1/2`$ $`=`$ $`\chi \stackrel{~}{B}\left[\mathrm{ln}\left({\displaystyle \frac{h(t)}{h(tt_a)}}\right)\right];`$ (35) $``$ $`\chi \stackrel{~}{B}(t_a/t_b),`$ (36) where terms $`(t_a/t_b)^2`$ have been neglected in the last expression, and $`t_b=h(t)/h^{}(t)`$. If $`a`$ is small enough, $`t_a`$ is simply proportional to $`t_b`$. Moreover, if $`h(t)`$ is indeed $`t^\delta `$, then $`a1/2=\chi \stackrel{~}{B}_1((1t_a/t_b)^\delta )`$, and $`t_a/t_b`$ is strictly constant. Our Figure (4) shows a plot of $`t_f`$, $`t_{a=0.55}`$, and $`t_{a=0.45}`$. The characteristic time scale $`t_f`$ tends asymptotically towards a power law $`c_2t^\alpha `$, with $`c_20.51`$ and $`\alpha 0.64`$ (according to our best estimate). The correlation function is found to grow logarithmically with $`t`$, and $`f(b(t,t^{}))=\mathrm{exp}(b(t,t^{}))`$ behaves as a power law of $`t`$. Figure (5) presents $`\mathrm{exp}(b(t,0))`$ and $`\mathrm{exp}(b(t,t^{}))`$ for a fixed $`t^{}`$. An algebraic decay $`t^\delta `$ of $`\mathrm{exp}(b(t,0))`$ is likely, while $`\mathrm{exp}(b(t,t^{}))`$ has not yet reached its asymptotic regime, but could tend to the same $`t^\delta `$ behaviour. From the asymptotic form (29, 30) we note that $`\mathrm{exp}(b(t,t^{}))h(t^{})/h(t);t,t^{}\mathrm{};t/t^{}`$ finite, and our Figure is consistent with $`h(t)=t^\delta `$, $`\delta 1.10`$. Also shown is $`\mathrm{exp}(b(t,tt_f))=t^\gamma `$, $`\gamma 0.42`$. Expanding $`\mathrm{exp}(b(t,t_f))`$ as $`(t/t_f)^\delta `$, and using $`t_ft^\alpha `$, one finds a relation $`\gamma =\delta (1\alpha )`$ between exponents. The agreement between $`\delta (1\alpha )=0.39`$ and $`\gamma 0.42`$ is acceptable. 3/ The linear displacement regime, in the presence of a driving force. A small force $`F`$ is applied and the displacement $`u(t)`$ monitored. The linear response implies that $`u(t)`$ must be proportional to $`F`$, and it is indeed the case for time intervals not too large. Figure (6) presents $`u(t)/F`$ for decreasing values of $`F`$. The curve $`F=0.05`$ is virtually indistinguishable from the integrated response $`(t)=_0^tr(t,s)\text{d}s`$, and this shows that $`lim_{F0}u_F(t)/F=(t)`$, i.e. the expected linear response behaviour. The other curves depart from the integrated response after a time $`t_F`$ decreasing with $`F`$. When starting from the isolated and aging situation $`F=0`$, the linear response only holds during a finite time interval $`0<tt_F`$. What happens later is the onset of a stationary state, with a well defined velocity $`v`$ and a non-linear dependence in the force as advocated by Horner . A study of this regime is to be published soon . In the linear response regime, Figure (6) is compatible with: $$u(t)=F(c_3+c_4\mathrm{ln}(t)),$$ (37) $`c_30.71`$ and $`c_4=0.54`$. ## IV The geometrical approach In this section, we transpose to the particle in a random potential some of the ideas which have revealed fruitful when applied to the spherical $`p`$-spin model, namely the geometrical analysis of Kurchan and Laloux . We expose first the main concepts of the method, and then propose a method for computing the limit value of the dynamical energy from generic properties of the potential $`V(𝐱)`$, working only for an exponential correlator $`f(y)=\mathrm{exp}(y)`$. The $`p`$-spin model starts aging below a dynamical temperature $`T_d`$ , and encounters a thermodynamical glassy transition at $`T_s<T_d`$ . Detailed investigations have brought an appealing picture of the complex free-energy landscape of the spherical $`p`$-spin model, accounting for many features of its thermodynamics and its dynamics . The phase space of the $`p`$-spin model can be investigated with the help of a “Thouless-Anderson-Palmer” free-energy $`\mathrm{\Phi }(m_i)`$ of the magnetization $`m_i,i=1\mathrm{}N`$. At low enough temperatures, the function $`\mathrm{\Phi }`$ develops many extrema $`m_i^{(\alpha )}`$, the TAP solutions $`\alpha `$. Those extrema which are minima, i.e. the second derivative matrix $`^2\mathrm{\Phi }/m_im_j`$ is definite positive, are metastable states, as they are separated from each others by extensive free-energy barriers. A particular realization of the system, prepared in a given metastable state $`\alpha `$ remains for ever in this state in the thermodynamic limit. The stability of a metastable state is related to the lowest eigenvalue $`\lambda _{min}`$ of $`\mathrm{Spec}(^2\mathrm{\Phi }/m_im_j)`$, spectrum of the hessian matrix. $`\lambda _{min}`$ turns out to be a monotonically decreasing function of the free energy $`\mathrm{\Phi }(m_i^{(\alpha )})`$ of the state itself. This defines the free-energy $`\mathrm{\Phi }_d`$ of the marginal states as $`\lambda _{min}=0`$ for $`\mathrm{\Phi }=\mathrm{\Phi }_d`$. Magnetizations such that $`\mathrm{\Phi }(m_i)\mathrm{\Phi }_d`$ represent regions of negative curvature which does not contribute to the thermodynamics but play a role in the dynamics . The glassy dynamics of the $`p`$-spin model is observed when the stable metastable states $`\mathrm{\Phi }^{(\alpha )}<\mathrm{\Phi }_d`$ are populated, or equivalently, when the canonical Boltzmann measure is split into its metastable components $`\alpha `$. At a temperature lower than $`T_d`$, thermal equilibration requires the system to explore all the relevant metastable states $`\mathrm{\Phi }<\mathrm{\Phi }_d`$. Such an equilibration is impossible as going from one metastable state to the next one requires to go over an infinite barrier. What happens instead to a system quenched from high temperature, to $`T<T_d`$, is the onset of aging. The systems wander more and more slowly around the magnetization region $`\mathrm{\Phi }(m_i)\mathrm{\Phi }_d`$, i.e. around the marginal states. The zero temperature relaxational dynamics is simpler because the free energy reduces to the hamiltonian $`H`$ of the spins $`s_i`$. At variance with the finite temperature case, the regions with negative curvature of $`H`$ are now well defined. Taking the limit $`T0`$ in the mean-field equations does not lead to any singular behaviour. This somewhat counter-intuitive property is the consequence of sending $`N\mathrm{}`$ first, by keeping finite the times $`t`$ and $`t^{}`$. The dynamics is a pure gradient descent, but remains non trivial. In order to perform a geometrical analysis of this relaxational dynamics, it is necessary to keep the dimension $`N`$ large but finite. Then, the relaxational process occurs till the particle falls into a true minimum of the Hamiltonian $`H(s_i)`$, and, at $`T=0`$, remains stuck there indefinitely. According to the description advocated in , a system starting from a random configuration $`\{s_i(0)\}`$ will explore regions with smaller and smaller gradient $`\mathbf{}H(s_i)`$, and a decreasing number of negative eigenvalues in the spectrum of $`_{ij}`$, hessian of $`H(s_i)`$. The typical time $`t_N^I`$ for reaching regions where $`_{ij}`$ has only $`I`$ negative eigenvalues, diverges as $`N`$ is sent to $`\mathrm{}`$ by keeping $`I`$ finite . As a consequence, in the $`N\mathrm{}`$ limit, the system is unable to reach within a finite time $`t`$ a true minima, or even a saddle between two minima, and the difference $`_d`$ remains positive. To what extent does the above picture describe the particle in a random potential ? Much less is known about the properties of the metastable states, and there is no obvious equivalent of the free energy functional $`\mathrm{\Phi }(𝐱)`$ of the mean particle’s position $`𝐱^{(\alpha )}`$ in the state $`\alpha `$. Nonetheless, we expect that the basic mechanism of the dynamical transition remain the same as for the spherical $`p`$-spin model, i.e. a slow relaxation toward a region of marginal states, $`\lambda _{min}0`$. When considering the zero temperature limit, the dynamics reduces to a gradient descent $`\dot{𝐱}(t)=\mathbf{}V(𝐱(t))`$. The metastable states now correspond to local minima of the potential $`V(𝐱)`$, and their stability will depend on the spectrum of the hessian $`_{ij}=_{ij}^2V(𝐱)`$, where $`_i`$ means $`/x_i`$ and $`_{ij}^2=^2/x_ix_j`$. The purpose of the geometrical approach, at zero temperature, is to relate the values provided by the more formal field-theoretical approach, to basic properties of the potential $`V(𝐱)`$. For instance, one must be able to compute the asymptotic values of the energy $`(t)`$, curvature $`(t)`$, and the “plateau value” $`lim_{T0}q/T`$ of the correlation function. This is the first step, already outlined in . One of the original contributions of this work concerns a second step, where we justify, with some geometrical arguments, many of the fine properties of the aging behaviour, beyond the time reparametrization invariant solution. A random matrix computation of the spectrum of $``$. A challenging problem in the study of the supercooled liquids dynamics, concerns the computation of the canonically averaged spectrum of the instantaneous normal modes $`\mathrm{Spec}()`$ . Here, $``$ is the dynamical matrix, hessian of the potential energy $`V`$ of the interacting particles, function of the coordinates $`𝐫_i`$. $`\mathrm{Spec}()`$ is any representative characteristic function of the eigenvalues spectrum, e.g. the density of states. $$\mathrm{Spec}()=Z^1_𝒟\text{d}𝐱\left(e^{\beta V}\mathrm{Spec}()\right).$$ (38) $`Z`$ is the canonical partition function, and $`𝒟`$ is a bound domain eventually becoming infinite. In the mean field situation, as the energy $`V`$ is extensive, the canonical average is dominated by a saddle point value $`V(\beta )`$ of the potential $`V(𝐱)`$. The analogous of (38) becomes: $`\mathrm{Spec}()`$ $`=`$ $`W_\beta ^1{\displaystyle _𝒟}\text{d}𝐱\delta (V(𝐱)V(\beta ))\mathrm{Spec}()(𝐱),`$ (39) $`W_\beta `$ $`=`$ $`{\displaystyle _𝒟}\text{d}𝐱\delta (V(𝐱)V(\beta )).`$ (40) In the $`p`$-spin case, the spectrum of $``$ is self-averaging with $`𝐱`$, i.e. the spectrum of $``$ is a shifted semi-circle, by an amount controlled by $`V(\beta )`$ only. The dynamical energy is found to be the highest value of $`V(\beta )`$ such that $`\mathrm{Spec}()0`$ (marginality). The averaged spectrum $`\mathrm{Spec}()`$ defined by equation (40) seems to be the natural quantity to consider when looking at the zero-temperature relaxational dynamics of our mean field model. We have found that as far as exponentially correlated potentials are concerned , $`\mathrm{Spec}()`$ is, at the leading order, a non-fluctuating quantity determined by $`V_0=V(\beta )`$. More precisely, $`\mathrm{Spec}()`$ has got a semi-circular distribution of radius $`\mathrm{\Lambda }`$, centred around $`𝒟`$. $`\mathrm{\Lambda }`$ $`=`$ $`4\sqrt{f^{\prime \prime }(0)};`$ (41) $`𝒟`$ $`=`$ $`{\displaystyle \frac{2f^{}(0)}{f(0)}}{\displaystyle \frac{V_0}{N}}.`$ (42) Let us outline our demonstration. We consider first the ($`𝐫`$ independent) “annealed average”. $$\overline{\delta (V(𝐫)V_0)\mathrm{Spec}()(𝐫)}.$$ (43) In order to compute (43), it is enough to enumerate the correlations of $`V(𝐫),_{ij}V(𝐫)`$, where $`𝐫`$ is an arbitrary point. All the $`_{ij}V(𝐫)`$ are independent at the leading order $`N^1`$, whereas the $`N+1`$ remaining variables $`V(𝐫),_{ii}V(𝐫)`$ are found to be correlated. One has: $$\begin{array}{cccc}N\overline{[_{ij}V(𝐫)]^2}& =& 4f^{\prime \prime }(0)& +𝒪(N^1);\\ N\overline{[_{ii}V(𝐫)]^2}& =& 12f^{\prime \prime }(0)& +𝒪(N^1);\\ N\overline{_{ii}V(𝐫)_{jj}V(𝐫)}& =& 4f^{\prime \prime }(0)& +𝒪(N^1);\\ \overline{V(𝐫)_{ii}V(𝐫)}& =& 2f^{}(0)& +𝒪(N^1);\\ N^1\overline{[V(𝐫)]^2}& =& f(0)& +𝒪(N^1).\end{array}$$ (44) $``$ is split into a scalar part $`𝒟\delta _{ij}`$ and a fluctuating part $`^{}`$. The elements of $`^{}`$ are independent and gaussian, and its eigenspectrum has, at the leading order, a semi-circular shape of radius $`4\sqrt{f^{\prime \prime }(0)}`$ centred around 0. If $`N\mathrm{}`$ and $`V(𝐫)/N`$ finite, then $`𝒟`$ is constant, up to fluctuations of order $`N^{1/2}`$ (cf appendix B). $$𝒟=\frac{2f^{}(0)}{f(0)}\frac{V(𝐫)}{N}+𝒪(N^{1/2}).$$ (45) The resulting spectrum is the one announced in equations (41, 42). In order to bridge the gap between (40) and (43), we consider now the two-points annealed average: $$\overline{\delta (V(𝐫)V_0)\delta (V(𝐫_1)V_1)\mathrm{Spec}()(𝐫)}.$$ (46) The analysis involves now correlations between $`V(𝐫),_{ij}V(𝐫),V(𝐫_1),_{ij}V(𝐫_1)`$. One finds that for a generic correlator $`f(y)`$, $`\mathrm{Spec}()(𝐫)`$ depends on both $`V(𝐫)`$ and $`V(𝐫_1)`$. However, if $`f(y)`$ obeys $`ff^{\prime \prime }(f^{})^2=0`$, with $`f(y)=\mathrm{exp}(y)`$ as a particular case, the dependence in $`V(𝐫_1)`$ disappears, and the result (42) holds. Computing $$\begin{array}{c}\overline{\delta (V(𝐫)V_0)\delta (V(𝐫_1)V_1)\mathrm{}}\\ \overline{\times \delta (V(𝐫_n)V_n)\mathrm{Spec}()(𝐫)},\end{array}$$ (47) becomes very difficult as $`n3`$, and we were not able to find a close expression for $`\mathrm{Spec}()(𝐫)(V_0,V_1,\mathrm{},V_n)`$. However, if $`ff^{\prime \prime }(f^{})^2=0`$, again $`\mathrm{Spec}()(𝐫)`$ depends only on $`V_0`$, and (42) is valid. This shows that the spectrum of $``$ is a local quantity, independent of the environment of the particle. Because $`\mathrm{Spec}((r))`$ is a function of $`V(𝐫)`$ only, we conclude that the average (40) is described by (41, 42) and that the self-averaging property of $`\mathrm{Spec}()(𝐫)`$ and its linear dependence in $`V(𝐫)`$, which was true for the $`p`$-spin model, is still true for exponential correlators. The appendix B gives further details on the computation of (43) and (46). Now, we suppose that the trajectory $`𝐱(t)`$ explores representative regions of the potential (i.e. non-exceptional points), for which the above mentioned results hold. The lowest eigenvalue $`\mathrm{SS}`$ of the hessian, defined by (42) becomes a time-dependent function: $$\mathrm{SS}(t)=\mathrm{\Lambda }\frac{2f^{}(0)}{f(0)}(t),$$ (48) leading to the energy dependent (through $`\mathrm{SS}`$) density of eigenvalues of $`_{ij}`$. The number of eigenvalues between $`\lambda \mathrm{SS}(t)`$ and $`\lambda \mathrm{SS}(t)+\text{d}\lambda `$ is $`\rho (\lambda )\text{d}\lambda `$ (time independent). $$\rho (\lambda )=2(\pi \mathrm{\Lambda }^2)^1\sqrt{\lambda (2\mathrm{\Lambda }\lambda )}.$$ (49) The marginality condition, by definition, is $`\mathrm{SS}0`$. Equation (48) yields the “geometrical energy”, necessary for $``$ to be marginal: $$_{geom}=2\frac{\sqrt{f^{\prime \prime }(0)}f(0)}{f^{}(0)},$$ (50) and the curvature $`_{geom}`$: $`_{geom}`$ $`=`$ $`{\displaystyle \text{d}\lambda \lambda \rho (\lambda )},`$ (51) $`=`$ $`4\sqrt{f^{\prime \prime }(0)}.`$ (52) After a time $`t`$ long enough, the particle evolves in a marginal region ($`\mathrm{SS}(t)0`$) of the potential $`V(𝐱)`$, with a small gradient $`\mathbf{}V(𝐱)`$. At low temperature, the potential may be developed up to the second order by means of local coordinates $`y_i`$: $`V(𝐲)=V(0)+𝐲\mathbf{}V(0)+_{i=1,N}\lambda _iy_i^2/2`$. The plateau value “$`q`$” of the correlation function $`b(t,t^{})`$ is thus given by assuming that each direction of curvature $`\lambda _i`$ is thermalized with $`y_i^2T/\lambda _i`$, and $`q=2N^1_{i=1,N}y_i^2`$: $`q_{geom}`$ $`=`$ $`{\displaystyle _0^{8\sqrt{f^{\prime \prime }(0)}}}\text{d}\lambda {\displaystyle \frac{2T}{\lambda }}\rho (\lambda )\|_{S=0},`$ (53) $`=`$ $`{\displaystyle \frac{T}{\sqrt{f^{\prime \prime }(0)}}}.`$ (54) Let us compare now with the results from the dynamical mean-field theory, in the zero temperature limit . $`\underset{t\mathrm{}}{lim}(t)`$ $`=`$ $`{\displaystyle \frac{f^{}(0)}{\sqrt{f^{\prime \prime }(0)}}}+{\displaystyle \frac{f(0)\sqrt{f^{\prime \prime }(0)}}{f^{}(0)}},`$ (55) $`\underset{t\mathrm{}}{lim}(t)`$ $`=`$ $`4\sqrt{f^{\prime \prime }(0)},`$ (56) $`q`$ $`=`$ $`{\displaystyle \frac{T}{\sqrt{f^{\prime \prime }(0)}}}.`$ (57) Agreement holds for the curvature and $`q`$, whereas the geometrical and dynamical energy differ, unless $`f(0)f^{\prime \prime }(0)=f^{}(0)^2`$. We cannot conclude about the relevance of the geometrical approach for a generic correlator, e.g. power law, as (42) probably does not hold. However, the exponentially correlated toy-model turns out to be a very favourable model, for which the geometrical approach gives reasonable results. The following of this paper aims at demonstrating that many features of the zero temperature dynamics of this model (exponents, aging, driving with a force) can be explained with the help of geometrical arguments. ## V The distribution of the gradient’s coordinates In this section, we define an orthonormal frame “attached” to the particle. The procedure used is reminiscent from the definition of the instantaneous normal modes in the study of supercooled liquids dynamics . Then, we investigate the statistical properties of the components of $`\mathbf{}V(𝐱(t))`$ in this special frame. We find that these components are distributed according to a self-similar form, determined by the value of the exponent $`\kappa `$ of the energy decay. We develop up to the second order the potential around the actual position of the particle $`𝐱(t)`$. $`𝐐(𝐱)`$ $`=`$ $`V(𝐱(t))+{\displaystyle \underset{i}{}}(𝐱_i𝐱_i(t))_iV(𝐱(t))`$ (59) $`+1/2{\displaystyle \underset{ij}{}}_{ij}V(𝐱(t))(𝐱_i𝐱_i(t))(𝐱_j𝐱_j(t)),`$ We define an orthonormal frame of eigendirections $`\{𝐞_{\lambda _i}(t)\}`$ in which the hessian $`_{ij}(t)=_i_jV(𝐱(t))`$ is diagonal. $`\lambda _i`$ belongs to the –time independent– interval $`[0,2\mathrm{\Lambda }]`$, so that the corresponding eigenvalue of $`𝐐`$ is just $`\lambda _i\mathrm{SS}(t)`$. We follow “adiabatically” the eigenvectors $`\{𝐞_{\lambda _i}(t)\}`$ as the particle moves. A mild assumption is that the $`\{𝐞_{\lambda _i}\}`$ evolve smoothly, provided the levels $`\lambda _i`$ are allowed to freely cross each other. This choice implies that any ordering of the $`\lambda _i`$ lasts only for a short period of time. The $`\{𝐞_{\lambda _i}\}`$ define a comoving frame, in which the gradient $`\mathbf{}V`$, or equivalently the velocity, can be projected. $`\mathbf{}V(𝐱(t))`$ $`=`$ $`{\displaystyle \underset{i}{}}𝜸_i(t)𝐞_{\lambda _i}(t),`$ (60) $`=`$ $`\dot{𝐱}(t).`$ (61) There are reasons to consider that the components $`𝜸_i(t)`$ are randomly and evenly distributed, even in the deterministic zero temperature limit. First, this randomness reflects the average over the “white” initial conditions. Then, as the correlator $`\overline{_{ij}^2V(𝐱)_{ij}^2V(𝐱^{})}`$ is exponentially short range correlated, one can suppose that the comoving frame is rotating on itself in a chaotic manner, as it does in the spherical ($`p3`$)-spin model . So, during the particle’s motion, each component $`𝜸_i`$ spreads continuously over the $`N1`$ others directions. The sign of $`𝜸_i(t)`$ itself is irrelevant, because of the arbitrary definition of the frame, invariant under the reflections $`𝐞_{\lambda _i}𝐞_{\lambda _i}`$. We claim that $`𝜸_i^2(t)`$ has to be preferred to $`𝜸_i(t)`$. On physical grounds, we propose to consider only the smoothed quantity $`𝜸_i^2(t)`$, obtained by averaging locally over the few $`\sqrt{N}`$ indices $`j`$ such that $`\lambda _iN^{1/2}<\lambda _j<\lambda _i+N^{1/2}`$. This is possible because the mean interspacing between the $`\lambda _i`$ is $`𝒪(N^1)`$. As $`N`$ goes to $`\mathrm{}`$, one expects $`𝜸_i^2(t)`$ to become a smooth function of $`\lambda _i`$, varying only on the scale $`\delta \lambda 1`$ (although, rigorously, the scale of variation is $`\delta \lambda N^{1/2}`$), making the dependence in the index $`i`$ irrelevant. The function: $$g(\lambda _i,t)=𝜸_i^2(t),$$ (62) is the distribution of the gradient’s coordinates (or equivalently of the instantaneous velocity coordinates) and is a central object in the present study. In this continuous limit, the two first time derivatives of $``$ can be expressed with the help of the density $`\rho (\lambda )`$ and the distribution $`g(\lambda ,t)`$ as: $`\dot{}(t)`$ $`=`$ $`{\displaystyle \underset{i}{}}_iV(𝐱(t))_iV(𝐱(t)),`$ (63) $`=`$ $`{\displaystyle \text{d}\lambda \rho (\lambda )g(\lambda ,t)}.`$ (64) $`\ddot{}(t)`$ $`=`$ $`{\displaystyle \underset{ij}{}}_jV(𝐱(t))_{ij}V(𝐱(t))_iV(𝐱(t)),`$ (65) $`=`$ $`{\displaystyle \text{d}\lambda \rho (\lambda )(\lambda \mathrm{SS}(t))g(\lambda ,t)}.`$ (66) In was already noticed that, due to the algebraic decay of the energy $`(t)=2+1.08t^\kappa `$, the ratio $`\ddot{}(t)/\dot{}(t)`$ was $`1/t`$. From section III, we know that $`\mathrm{SS}(t)t^{0.67}`$, which implies $`\ddot{}(t)\mathrm{SS}(t)\dot{}(t)`$, and thus: $`{\displaystyle \text{d}\lambda \rho (\lambda )\lambda g(\lambda ,t)}`$ $`=`$ $`\ddot{}(t)+\mathrm{SS}(t){\displaystyle \text{d}\lambda \rho (\lambda )g(\lambda ,t)},`$ (67) $``$ $`\mathrm{SS}(t){\displaystyle \text{d}\lambda \rho (\lambda )g(\lambda ,t)}.`$ (68) The first moment of $`g(\lambda ,t)\rho (\lambda )`$, is proportional to $`\mathrm{SS}(t)`$, suggesting a self-similar scaling form for $`g(\lambda ,t)`$, valid for $`t\mathrm{}`$ and $`T=0`$ (Figure 7): $$g(\lambda ,t)=\mathrm{\Gamma }(t)\widehat{G}\left(\frac{\lambda }{\mathrm{SS}(t)}\right),$$ (69) The knowledge of the other moments of $`g`$ would be useful to confirm equation (69), but unfortunately, they are very difficult to compute, and are no more given by the next derivatives of $``$. As $`t`$ increases, only the smaller $`\lambda _i`$ keep any relevance, and the density $`\rho `$ is well approximated by its $`\lambda 0`$ equivalent $`\pi ^1(2/\mathrm{\Lambda })^{3/2}\sqrt{\lambda }`$. In this limit, the loss of energy rate becomes, from (64) and (69): $$\dot{\mathrm{SS}}(t)=\frac{2f^{}(0)}{f(0)}\dot{}(t)\mathrm{\Gamma }(t)\mathrm{SS}(t)^{3/2}.$$ (70) The knowledge of the exponent $`\kappa `$ of $`\mathrm{SS}(t)t^\kappa `$ (section III) fixes the prefactor $`\mathrm{\Gamma }`$ up to a constant, to: $$\mathrm{\Gamma }=\mathrm{SS}^{(2\kappa )/2\kappa }.$$ (71) Our section II suggests $`\kappa `$ is very close to $`2/3`$, which would imply $`\mathrm{\Gamma }\mathrm{SS}`$. The next momentum of $`g(\lambda ,t)\rho (\lambda )`$ provides information on the time correlations of the unit vector $`𝐰(t)`$ of the particle’s trajectory. One the one hand, $`\left(_t\mathbf{}V(𝐱(t))\right)^2`$ $`=`$ $`{\displaystyle \underset{i}{}}\left({\displaystyle \underset{j}{}}_jV(x(t))_{ij}V(x(t))\right),`$ (73) $`\times \left({\displaystyle \underset{k}{}}_kV(x(t))_{ik}V(x(t))\right),`$ $`=`$ $`{\displaystyle \underset{i,j,k}{}}_jV_{ji}V_{ik}V_kV,`$ (74) $`=`$ $`{\displaystyle \text{d}\lambda \rho (\lambda )(\lambda \mathrm{SS}(t))^2g(\lambda ,t)}.`$ (75) With the scaling form for $`g(\lambda ,t)`$, the right hand side is of order $`\mathrm{\Gamma }(t)\mathrm{SS}(t)^{7/2}`$. On the other hand, we perform a decomposition $`\mathbf{}V(𝐱(t))=M(t)𝐰(t)`$. The norm $`M(t)`$ equals $`(\dot{}(t))^{1/2}`$, and $`𝐰(t)`$ is the unit vector, tangent to the trajectory. The following equality holds: $$\left(_t\mathbf{}V\right)^2=\left(_tM\right)^2+M^2_t𝐰^2.$$ (76) This sum is clearly dominated by $`M^2_t𝐰^2`$, with $`M^2\mathrm{\Gamma }\mathrm{SS}^{3/2}`$. The unitary vector rotates, regardless to the actual value of $`\mathrm{\Gamma }(t)`$, at a rate $`_t𝐰\mathrm{SS}(t)`$. One expects the “director” $`𝐰(t)`$ to have changed its orientation after a typical time $`\mathrm{SS}(t)^1`$, which looks like a “persistence time” for the trajectory of the particle. Consequently, the motion of $`𝐱(t)`$ crosses over from a “ballistic” regime $`𝐱(t+\delta t)𝐱(t)M^2(\delta t)^2;\delta t\mathrm{SS}^1`$ to a diffusive regime $`𝐱(t+\delta t)𝐱(t)D\delta t;\delta t\mathrm{SS}^1`$. The existence of a diffusive regime is here inferred by the expansion (29, 31), valid only for exponentially correlated disorder, and by no means generic. ## VI A short time, quasi-static approximation We investigate here the breakdown of the fluctuation dissipation theorem, in the zero temperature limit. The fluctuation-dissipation violation is measured by the function $`\overline{T}(t,t^{})`$ (cf 24). We propose here a model for the short time evolution of $`\overline{T}(t,t^{})`$, and show that its predictions are in good agreement with the findings of section II. We approximate locally the potential $`V(𝐱)`$ around $`𝐱(t)`$ by a quadratic function (59), which may can be considered as constant provided we restrict ourselves to a time separation $`tt^{}`$ small enough. One can always find a coordinate system $`\{𝐲_i\}`$ such that this quadratic potential reads: $$𝐐(𝐲)=𝐐(0)+1/2\underset{i}{}(\lambda _i\mathrm{SS})y_i^2,$$ (77) where the coordinates $`𝐲`$ must not be confused with the original coordinates $`𝐱`$ of the relaxational motion. This section aims at demonstrating that when a particle diffuses, or relaxes in such a parabolic potential, then a characteristic time $`t_f`$ scaling like $`\mathrm{SS}^1`$ arises, which turns out to be the time scale along which the function $`\overline{T}(t,t^{})`$ departs from its equilibrium value 0, i.e the fluctuation-dissipation violation characteristic time. We consider a particle moving on the potential (77), starting at $`t_0`$, and define the time difference $`\tau =tt_0`$. The intermediate steps of the calculation make use of $`\tau ,t_0`$, while the final results are expressed in term of $`t,t^{}`$ in relation with the original out-of-equilibrium relaxation. Let us consider the same local average $`𝐲_i^2(\tau )`$ as in equation (62). The $`𝐲_i`$ are related to the gradient’s coordinates by $`𝜸_i(\tau )=(\lambda _iS)y_i(\tau )`$: $$y_i^2(\tau )=\frac{g(\lambda ,t_0+\tau )}{(\lambda \mathrm{SS})^2}.$$ (78) The initial conditions $`y_i^2(\tau =0)`$ are given by $`g(\lambda ,t_0)`$. One computes the fluctuation dissipation violation $`\overline{T}(t,t^{})`$, when the quadratic potential (77) does not evolve with time ($`\mathrm{SS}`$ fixed once for all), and with initial conditions arising from a realistic distribution $`g(\lambda ,t_0)=\mathrm{\Gamma }\widehat{G}(\lambda /\mathrm{SS})`$. $$y_i(\tau )=y_i(0)e^{(\lambda _i\mathrm{SS})(\tau )}.$$ (79) The distribution $`g(\lambda ,t_0+\tau )`$ evolves like $`𝐲_i^2(\tau )(\lambda \mathrm{SS})^2`$. One has, far all $`t>t_0`$: $`g(\lambda ,t_0+\tau )`$ $`=`$ $`g(\lambda ,t_0)e^{2(\lambda \mathrm{SS})(\tau )};`$ (80) $`_tg(\lambda ,t)`$ $`=`$ $`2g(\lambda ,t)(\lambda \mathrm{SS}).`$ (81) The usual response $`r(t,t^{})=N^1_i\delta y_i(t)/\delta \zeta _i(t^{})`$, and correlation $`b(t,t^{})=N^1_i(y_i(t)y_i(t^{}))^2`$ functions reexpress in terms of $`g(\lambda ,t^{})`$: $`r(t,t^{})`$ $`=`$ $`{\displaystyle \text{d}\lambda \rho (\lambda )e^{(\lambda \mathrm{SS})(tt^{})}};`$ (82) $`b(t,t^{})`$ $`=`$ $`{\displaystyle \text{d}\lambda \rho (\lambda )g(\lambda ,t^{})\left(\frac{1e^{(\lambda \mathrm{SS})(tt^{})}}{\lambda \mathrm{SS}}\right)^2};`$ (83) $`_t^{}b(t,t^{})`$ $`=`$ $`2{\displaystyle \text{d}\lambda \rho (\lambda )g(\lambda ,t^{})\left(\frac{1e^{(\lambda \mathrm{SS})(tt^{})}}{\lambda \mathrm{SS}}\right)};`$ (84) $`=`$ $`2\overline{T}(t,t^{}).`$ (85) By inserting $`g(\lambda ,t^{}=t_0)=\mathrm{\Gamma }\widehat{G}(\lambda /\mathrm{SS})`$ in (85), one deduces the short time, $`\tau \mathrm{SS}^1`$, value of $`\overline{T}(t,t^{})`$, $$\overline{T}(t,t^{})=(tt^{})\mathrm{\Gamma }\mathrm{SS}^{3/2},$$ (86) and the intermediate time $`\tau \mathrm{SS}^1`$ one, $`r(t,t^{})`$ $`=`$ $`\mathrm{SS}^{3/2}\mathrm{\Phi }_0(\mathrm{SS}(tt^{})),`$ (87) $`b(t,t^{})`$ $`=`$ $`2\mathrm{\Gamma }\mathrm{SS}^{1/2}\mathrm{\Phi }_1(\mathrm{SS}(tt^{})),`$ (88) $`\overline{T}(t,t^{})`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }}{\mathrm{SS}}}{\displaystyle \frac{\mathrm{\Phi }_0}{\mathrm{\Phi }_1}}(\mathrm{SS}(tt^{})).`$ (89) $`\mathrm{\Phi }_0`$,$`\mathrm{\Phi }_1`$ are scaling function presented in appendix B. Equation (89) shows that $`\mathrm{SS}^1`$ plays the role of a characteristic time for the onset of the effective temperature $`\overline{T}`$. Assuming now the very likely value $`\kappa =2/3`$ and $`\mathrm{\Gamma }=\mathrm{SS}`$, one finds that $`\overline{T}`$ becomes an order one quantity after a time $`t_f\mathrm{SS}^1`$. One concludes that the characteristic time scale $`t_f`$ should scale like $`t_f=t^\alpha t^\kappa `$, and $`\alpha =\kappa =2/3`$. Our numerics (Figure 4) lead to an estimated value $`\alpha 0.64`$, while $`\kappa 0.67`$. While we haven’t proved that $`\kappa `$ is actually $`2/3`$, we find the agreement satisfactory, and believe that the above picture describes correctly the first stage of the breaking of the “fluctuation-dissipation relation at zero temperature”. How long can the quadratic approximation (59) accurately describe the original relaxational process ? As $`\mathrm{SS}(t)`$ decreases algebraically, the necessary time $`\delta t`$ to have $`\|\mathrm{SS}(t+\delta t)\mathrm{SS}(t)\|\mathrm{SS}(t)`$ is $`t`$ itself. More seriously, we have seen in the previous section, that the unit vector of the trajectory $`\dot{𝐱}(t)`$ changes with the time scale $`\mathrm{SS}^1t^{2/3}`$. As this change is somewhat related to the frame’s chaotic motion, we deduce than $`\mathrm{SS}^1`$ must be an upper limit of validity of the quasi-static approximation. Finally, the relaxation on the saddle becomes ill-defined when $`\mathrm{SS}(tt^{})1`$, due to the exponential divergence of the functions $`r(t,t^{})`$ and $`b(t,t^{})`$, given by the equations (82, 83). We arrive to the conclusion that this quasi-static picture is not valid beyond times much greater than $`\mathrm{SS}^1`$, but provides a strong presumption in favour of $`t_f(t)\mathrm{SS}^1(t)t^{2/3}`$, in good agreement with our numerical findings (section III and Figure 4). Let us close this section by computing the typical distance covered during a time interval $`tt^{}<\mathrm{SS}^1`$, with a gradient coordinates distribution $`g(\lambda ,t^{})=\mathrm{\Gamma }\widehat{G}(\lambda /\mathrm{SS})`$. $`b(t,t^{})`$ $`=`$ $`\mathrm{\Gamma }\mathrm{SS}^{3/2}(tt^{})^2`$ (90) $`=`$ $`\mathrm{SS}^{5/2}(tt^{})^2\text{if}\kappa =2/3.`$ (91) For a time interval $`(tt^{})\mathrm{SS}^1`$, $`b(t,t^{})\mathrm{\Gamma }\mathrm{SS}^{1/2}1`$, becoming $`b(t,t^{})\mathrm{SS}^{1/2}`$ for $`\kappa =2/3`$. As $`\mathrm{SS}^{1/2}`$ tends to zero, the characteristic time $`t_b`$ of the evolution of $`b(t,t^{})`$ in the aging regime, is necessarily much greater than $`t_f\mathrm{SS}^1`$. ## VII A dynamics restricted to the downhill directions This section shows how the above approach describes the long time aging regime. The equation (70) has a simple physical interpretation. The only non-vanishingly small components $`\gamma _i`$ of $`\mathbf{}V`$ are those corresponding to $`\lambda _i\mathrm{SS}`$. Only a number $`N_0^{\mathrm{SS}}\text{d}\lambda \rho (\lambda )N\mathrm{SS}^{3/2}`$ directions $`i`$ are contributing to $`(\mathbf{}V)^2=_i𝜸_i^2`$. Each one of these $`𝜸_i`$ has a magnitude of order $`𝜸_i^2\mathrm{\Gamma }`$. As a result, $`N\dot{}`$ scales like $`_i𝜸_i^2=N\mathrm{SS}^{3/2}\mathrm{\Gamma }`$. $$\dot{}((t)_d)^{1+1/\kappa }.$$ (92) The relaxation dynamics looks like if it was controlled by the difference $`(t)_d`$. The linear response regime. A constant force $``$ is now applied, uncorrelated to the potential $`V`$. Each one of its (comoving) coordinate $`f_i`$ is random, time-dependent as the frame rotates during the particle’s motion, and has a magnitude $`f_iF`$. We suppose $`F`$ weak enough to be considered as a perturbation around the relaxational dynamics described in section VI. At any time, there are “open”, or downhill directions, with $`\lambda _i\mathrm{SS}`$ and “close”, or uphill directions with $`\lambda _i\mathrm{SS}`$. The close directions behave as confining harmonic potentials which prevent the (weak) force $`f_i`$ to drive the particle along this direction. The open directions are the one along with the external force drives efficiently the particle away. As the particle moves, “open” and “close” directions exchange their role, but the proportion of open directions remains proportional to $`\mathrm{SS}^{3/2}`$. The force $``$ induces a displacement $`\dot{𝐱}`$ whose components are $`\dot{𝐱}_if_i`$ along an open direction and $`\dot{𝐱}_i0`$ along a close direction . The average velocity $`\dot{𝐱}/`$ is given by ($`\theta `$ Heaviside function): $`{\displaystyle \frac{\dot{𝐱}(t)}{}}`$ $``$ $`{\displaystyle \underset{i}{}}f_i^2\theta (\mathrm{SS}\lambda _i)/;`$ (93) $`=`$ $`N\mathrm{SS}^{3/2}F^2/(\sqrt{N}F);`$ (94) $`=`$ $`\sqrt{N}F\mathrm{SS}^{3/2};`$ (95) that we identify to $`\sqrt{N}\dot{u}(t)`$. As a result, one finds a velocity proportional to the number of downhill directions: $$\dot{u}F\mathrm{SS}^{3/2}.$$ (96) Inserting the likely value $`\kappa =2/3`$, one finally gets a displacement $`u(t)u(t^{})F(\mathrm{ln}(t)\mathrm{ln}(t^{}))`$, well confirmed by the numerics (section III and Figure 6). This pure relaxational motion is driven by the components $`𝜸_i`$ of $`\mathbf{}V`$ along the open directions, while the external force acts with $`f_i`$ along the same open directions. One expects the linear response to hold if $`f_i^2𝜸_i^2`$ but to break down when $`f_i^2𝜸_i^2`$. This leads to a predicted cross-over time scaling like $`\mathrm{\Gamma }(t_F)=F^2`$, or $`t_FF^{4/(\kappa 1)}`$, to be investigated in a forthcoming publication . The diffusive regime. The asymptotic behaviour predicted for $`b(t,t^{})`$ is, from equations (29-31): $$b(t,t^{})=\frac{tt^{}}{t_b}+𝒪\left(\frac{tt^{}}{t_b}\right)^2.$$ (97) One recognises a simple diffusive behaviour, with effective diffusivity $`t_b^1`$. From section IV, we know that the short-time motion $`(tt^{})\mathrm{SS}^1t_f`$ is ballistic, and that the particles covers a distance of $`\mathrm{\Gamma }\mathrm{SS}^{1/2}`$ On the other hand, our results from section III show that the direction $`𝐰`$ of the trajectory $`𝐱(t)`$ uncorrelates itself after this same time $`\mathrm{SS}^1`$. Using a well-known result on correlated random walks, and assuming a free diffusive behaviour at intermediate times $`tt^{}t_b`$, as inferred by equation (97), one finds: $$b(t,t^{})\left(\frac{tt^{}}{\mathrm{SS}^1}\right)\times \left(\frac{\mathrm{\Gamma }}{\sqrt{\mathrm{SS}}}\right).$$ (98) This corresponds to ballistic steps of length $`(\mathrm{\Gamma }/\sqrt{\mathrm{SS}})`$ (91), and a cross-over time from ballistic to diffusive regime equal to $`\mathrm{SS}^1`$. As $`t_f\mathrm{SS}^1`$, $$b(t,t^{})(tt^{})\mathrm{\Gamma }\mathrm{SS}^{1/2},$$ (99) leading to the identification: $`t_b`$ $``$ $`\mathrm{\Gamma }^1\mathrm{SS}^{1/2};`$ (100) $``$ $`\mathrm{SS}^{1/\kappa }`$ (101) If $`\kappa `$ is taken to be $`2/3`$, one gets $`t_bt^{}t`$. Let us discuss different reasons to be confident in the scaling $`t_f\mathrm{SS}^1`$, $`t_b\mathrm{SS}^{3/2}`$ and $`\mathrm{SS}t^{2/3}`$. First, a matching argument similar to predicts $`t_bt_f^{3/2}`$. Then, the result $`t_bt^{}`$ is in agreement with the conjecture $`h(t)t^\delta `$. This entails a logarithmic growth of $`b(t,t^{})`$, $`t^{}`$ fixed, and we have asymptotically (i.e. $`t,t^{}1`$, and $`t/t^{}1`$) a free brownian motion in logarithmic time. $$b(t,t^{})=\delta (\mathrm{ln}t\mathrm{ln}t^{}).$$ (102) This makes $`\mathrm{exp}(b(t,t^{}))`$ as well as $`r(t,t^{})`$ decaying as a power law. While we have no demonstration of that, we think that a power-law decay of the memory function $`f(b(t,t^{}))r(t,t^{})`$ is necessary for the “fine tuned” aging solution of the system (9). Asking for a power-law decay $`f(b(t,t^{}))`$ in turn fixes $`\kappa `$ to $`2/3`$. Finally, if $`\kappa =2/3`$, $`t_b=\mathrm{SS}^{3/2}`$, and the characteristic times for the linear response regime $`\dot{u}(t)F/t_b`$, and for the diffusion regime $`b(t,t^{})(tt^{})/t_b`$ are the same, which is consistent with the persistence of an “Einstein relation” at the beginning of the aging regime. ## VIII Conclusion We have proposed a geometrical description of the mean-field relaxational dynamics of a particle, for a subclass of short-range correlated disorders. We have restricted ourselves to the isolated case, and to the driven case in the linear response regime. A numerical integration of the mean-field equations gives evidence of a power-law decay of the dynamical energy with an exponent $`\kappa `$ numerically close to $`2/3`$. We also found evidence of a logarithmic growth $`b(t,t^{})\mathrm{ln}t`$ consistent with the conjecture $`h(t)t^\delta `$ for the reparametrization function $`h`$. The exponential correlator makes it possible to compute the density of eigenvalues of the hessian $``$ associated to the random potential, and we were able to predict the correct value (i.e. $`2`$) of the dynamical energy $`_d`$. Introducing a comoving frame, reminiscent from the INM frame of a supercooled liquid, we derive an expression for the distribution $`g(\lambda ,t)`$ of the components of $`\mathbf{}V\left(𝐱(t)\right)`$. This expression is $`g=\mathrm{\Gamma }\widehat{G}(\lambda /\mathrm{SS})`$, where $`\mathrm{SS}(t)`$ is the (time dependent) lowest eigenvalue of $`\mathrm{Spec}()`$. For reasons exposed in section VII, namely the consistence with $`h(t)t^\delta `$, the requirement that $`f(b)`$ is likely to decrease as a power law, and acknowledging the numerical estimate of $`\kappa `$, we believe that $`\kappa `$ is indeed equal to $`2/3`$. This leads to the following predictions: (1) $`\mathrm{\Gamma }S`$, and a typical gradient coordinate is, along a downhill direction, $`\|𝜸_i\|\sqrt{S}t^{1/3}`$. (2) From a short time, harmonic expansion of the particle’s motion (section VI) the characteristic time $`t_f`$ leading to the appearance of an effective temperature goes like $`t_fS^1t^{2/3}`$. (3) The characteristic time at the beginning of the aging regime is $`t_bS^{3/2}t`$. Both linear response $`\dot{u}F/t_b`$ and diffusion $`b(t,t^{})(tt^{})/t_b`$ are controlled by it. We conclude that the aging mechanism of this model comes from a simultaneous decrease of the number of downhill directions (going like $`NS^{3/2}Nt^1`$) and of the typical gradient component $`\|𝜸_i\|t^{1/3}`$. We predict also that the effect of a constant force brings about a dramatic change in the dynamics after a time $`t_FF^3`$, reaching a out-of-equilibrium but stationary regime . ### Acknowledgements I especially thank L.Cugliandolo and J.Kurchan for having lent me their numerical code, and S.Scheidl, J.P Bouchaud, J.Kurchan, M.Mézard and A.Cavagna for discussions on this field. I thank D.Feinberg for suggestions and criticisms about the manuscript. I warmly thank the hospitality of the Department of Physics, IISc, Bangalore, where a part of the writing has been done. ## A The MSR action The action leading to equations (9) is: $`S[x,i\stackrel{~}{x}]`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}\text{d}t\left\{T{\displaystyle \underset{j=1}{\overset{N}{}}}(i\stackrel{~}{x}_j)^2(t)+i\stackrel{~}{x}_1(t)(\dot{x}_1(t)N^{1/2}F)+{\displaystyle \underset{j=2}{\overset{N}{}}}i\stackrel{~}{x}_j(t)\dot{x}_j(t)\right\}`$ (A2) $`+{\displaystyle _0^{\mathrm{}}}\text{d}t\text{d}s\left\{f^{}(d(t,s)){\displaystyle \underset{j=1}{\overset{N}{}}}i\stackrel{~}{x}_j(t)i\stackrel{~}{x}_j(s)+4f^{\prime \prime }(d(t,s))r(t,s){\displaystyle \underset{j=1}{\overset{N}{}}}i\stackrel{~}{x}_j(t)(x_j(t)x_j(s))\right\},`$ and the expectation value of an observable $`𝒪[x(t),i\stackrel{~}{x}(t)]`$, averaged over the disorder, is given by: $$\overline{𝒪}=𝒟x[t]𝒟\stackrel{~}{x}[t]𝒪\mathrm{exp}(S).$$ (A3) ## B The spectrum of the Hessian $``$ We consider $`\overline{\delta (V(𝐫)V_0)\mathrm{Spec}()(𝐫)}`$ for any arbitrary potential. $`V_0=N`$ is fixed, and the $`_{ij}=_{ij}^2V(𝐫)`$ are $`N(N+1)/2`$ gaussian random variables. The correlations among the $`_{ij}`$ are listed in equation (44). We define the self-averaging quantity $`𝒟=N^1_i_{ii}V(𝐫)`$ so that: $`N\overline{(_{ii}V(𝐫)𝒟)^2}`$ $`=`$ $`\mathrm{\hspace{0.33em}8}f^{\prime \prime }(0)8f^{\prime \prime }(0)/N,`$ (B1) $`N\overline{(_{ii}V(𝐫)𝒟)V(𝐫)}`$ $`=`$ $`\mathrm{\hspace{0.33em}0},`$ (B2) $`N\overline{(_{ii}V(𝐫)𝒟)(_{jj}V(𝐫)𝒟)}`$ $`=`$ $`8f^{\prime \prime }(0)/N,`$ (B3) $`=`$ $`\mathrm{\hspace{0.33em}0}+𝒪(N^1).`$ (B4) The hessian is now $`_{ij}=𝒟\delta _{ij}+_{ij}^{}`$. $`^{}`$ is a matrix of independent gaussian centred random numbers. The diagonal elements are slightly correlated (of order $`1/N^2`$) and have a different variance than the off-diagonal elements, but this does not prevent the Wigner result to apply and the spectrum of $`^{}`$ is a centred semi-circle of width $`\mathrm{\Lambda }=4\sqrt{f^{\prime \prime }(0)}`$. The determination of $`𝒟`$ follows from the fact that $``$ and $`𝒟`$ are gaussian distributed, with correlations: $`N\overline{^2}`$ $`=`$ $`f(0);`$ (B5) $`N\overline{𝒟^2}`$ $`=`$ $`4\left({\displaystyle \frac{N+2}{N}}\right)f^{\prime \prime }(0);`$ (B6) $`N\overline{𝒟}`$ $`=`$ $`2f^{}(0).`$ (B7) The joint probability distribution of $``$ and $`𝒟`$ is: $`𝒫(𝒟,)`$ $`=`$ $`{\displaystyle \frac{N}{2\pi \sqrt{cf(0)}}}\mathrm{exp}\left({\displaystyle \frac{N}{2}}\left[{\displaystyle \frac{^2}{f(0)}}+{\displaystyle \frac{(𝒟a)^2}{c}}\right]\right);`$ (B8) $`c`$ $`=`$ $`{\displaystyle \frac{4}{f(0)}}(f(0)f^{\prime \prime }(0)f^{}(0)^2)+{\displaystyle \frac{8}{N}}f^{\prime \prime }(0);`$ (B9) $`a`$ $`=`$ $`{\displaystyle \frac{2f^{}(0)}{f(0)}}.`$ (B10) Fluctuations of $`𝒟`$ are of order $`N^{1/2}`$ around its saddle point value $`2f^{}(0)/f(0)\times `$. It follows that $`\mathrm{Spec}()`$ is a semi-circle of radius $`\mathrm{\Lambda }`$ shifted by an amount $`𝒟=2f^{}(0)/f(0)\times `$. Let us consider now $`\overline{\delta (V(𝐫)V_0)\delta (V(𝐫_1)V_1)\mathrm{Spec}()(𝐫)}`$. This simple average measures the non-locality of $`\mathrm{Spec}()`$, i.e. its dependence on the values taken by the random potential $`V(𝐫^{})`$ around $`𝐫`$. The rotational invariance of the above average is broken, and $`𝐫𝐫^{}`$ plays a special role. We relabel hereafter the direction 1 to coincide with $`𝐫𝐫^{}`$, and define $`b=𝐫𝐫^{}^2/N`$. Correlations are now, in addition to (44), $`\overline{_{ii}V(𝐫)V(𝐫^{})}`$ $`=`$ $`2f^{}(b)\text{if}i2;`$ (B11) $`\overline{_{11}V(𝐫)V(𝐫^{})}`$ $`=`$ $`2f^{}(b)+4bf^{\prime \prime }(b);`$ (B12) $`\overline{_{ij}V(𝐫)V(𝐫^{})}`$ $`=`$ $`0.`$ (B13) $`𝒟`$ is again defined as $`N^1_i_{ii}V(𝐫)`$ and $`_{(i,j)2}^{}`$ is equivalent to the above situation: independent, centred, gaussian random components, and the spectrum is a centred semi-circle. Adding one row and one column of random independents elements to $`_{(i,j)2}^{}`$ must not change the density profile of eigenvalues. This is because this eigenvalue distribution is a fixed point under the change $`NN+1`$, as argued in the cavity approach of the problem. A possible trouble come from the single component $`_{11}V𝒟`$ which does not average to 0, but this does not alter the final result more than by a single isolated eigenvalue. It is possible to show, with the help of a formal field theoretical approach, that the correlations (44),(B13) indeed lead to the ordinary $`N\mathrm{}`$ saddle point for $`\mathrm{Spec}(^{})`$, i.e. a semi-circle law of radius $`\mathrm{\Lambda }`$. The computation of $`𝒟`$ follows closely the lines of the previous paragraph. We found that if $`=V(𝐫)/N`$, $`^{}=V(𝐫^{})/N`$ and $`b=𝐫𝐫^{}^2/N`$, then: $`𝒟`$ $`=`$ $`\left({\displaystyle \frac{f^{}(0)+f^{}(b)}{f(0)+f(b)}}\right)(+^{})`$ (B15) $`+\left({\displaystyle \frac{f^{}(0)f^{}(b)}{f(0)f(b)}}\right)(^{}).`$ For a generic correlator, there is an explicit dependence on $`^{}`$ (“non locality”) while for an exponential correlator $`f=\mathrm{exp}(y)`$, the above formula reduces to $`𝒟=2f^{}(0)/f(0)\times `$. This suggests that the determination of $`\mathrm{Spec}()`$ from (40) is a complex problem and the simple behaviour (42) fails for a generic $`f`$. The exponential correlator, however, has a strong property. The average: $`N\overline{\left({\displaystyle \frac{2f^{}(0)}{f(0)}}\mathrm{SS}\right)^2}`$ $`=`$ $`{\displaystyle \frac{4}{f(0)}}(ff^{\prime \prime }f^2)+{\displaystyle \frac{8}{N}}f^{\prime \prime }(0);`$ (B16) $`=`$ $`{\displaystyle \frac{8}{N}}f^{\prime \prime }(0).`$ (B17) is 0 at the order $`N^1`$. This means that, while $`b`$ is strictly positive, there is no possible fluctuations of $`𝒟`$ around $`2f^{}(0)/f^{\prime \prime }(0)\times `$. Repeating the argument for (47), $`n`$ finite, shows that $`\mathrm{Spec}()`$ depends only on $`V(𝐫)`$ and not on its local environment. Thus, we argue that the average (40) is given by (41, 42), as announced. ## C The quasi-static picture From equation (82), we derive the expression for $`r(t^{}+\tau ,t^{})`$. $$r(t^{}+\tau ,t^{})=2\frac{e^{\mathrm{SS}\tau }\text{I}_1(\mathrm{\Lambda }\tau )e^{\mathrm{\Lambda }\tau }}{\mathrm{\Lambda }\tau },$$ (C1) where $`\text{I}_1`$ is the first kind modified Bessel function. The short time expansion of (85) is: $`_t^{}b(t^{}+\tau ,t^{})`$ $`=`$ $`2\tau {\displaystyle \text{d}\lambda \rho (\lambda )g(\lambda ,t^{})};`$ (C2) $`=`$ $`2\tau \left(\dot{}(t^{})\right)^2;`$ (C3) $`=`$ $`2\tau \mathrm{\Gamma }S^{3/2}.`$ (C4) In the intermediate time separation regime, $`\tau `$ is of order $`\mathrm{SS}^1`$. The integral is dominated by $`\lambda \mathrm{SS}`$ and cut off by $`g(\lambda ,t^{})`$ for $`\lambda \mathrm{SS}1`$. $`\rho (\lambda )`$ can be replaced by its $`\lambda 0`$ equivalent. $`r(t^{}+\tau ,t^{})`$ $``$ $`\sqrt{2/\pi }e^{\mathrm{SS}\tau }(\mathrm{\Lambda }\tau )^{3/2};`$ (C5) $`=`$ $`\mathrm{SS}^{3/2}(\sqrt{2/\pi }\mathrm{\Lambda }^{3/2})(\mathrm{SS}\tau )^{3/2}e^{\mathrm{SS}\tau };`$ (C6) $`=`$ $`\mathrm{SS}^{3/2}\mathrm{\Phi }_0(\mathrm{SS}\tau ).`$ (C7) $`_t^{}b(t^{}+\tau ,t^{})`$ $`=`$ $`2\mathrm{\Gamma }\mathrm{SS}^{1/2}{\displaystyle _0^{2\mathrm{\Lambda }/\mathrm{SS}\mathrm{}}}\text{d}u(2/\mathrm{\Lambda })^{3/2}\pi ^1\sqrt{u}`$ (C9) $`\times \widehat{G}(u)\left({\displaystyle \frac{1e^{\mathrm{SS}\tau (u1)}}{u1}}\right);`$ $`=`$ $`2\mathrm{\Gamma }S^{1/2}\mathrm{\Phi }_1(\mathrm{SS}\tau ).`$ (C10) The effective temperature behaves as: $$\overline{T}(t^{}+\tau ,t^{})=\left(\frac{\mathrm{\Gamma }}{\mathrm{SS}}\right)\frac{\mathrm{\Phi }_1(\mathrm{SS}\tau )}{\mathrm{\Phi }_0(\mathrm{SS}\tau )},$$ (C11) reducing to $$\overline{T}(t,t^{})=\frac{\mathrm{\Phi }_1}{\mathrm{\Phi }_0}(\mathrm{SS}(tt^{})),$$ (C12) if $`\kappa =2/3`$ and $`\mathrm{\Gamma }=\mathrm{SS}`$. Captions FIGURE 1. Parametric plot of the integrated response $`b(t,t^{})`$ vs $`(t)=_0^t\text{d}sr(t,s)`$ at zero temperature, for time steps $`h=0.025,0.5,1.0,2.0`$. The horizontal part corresponds to the short time regime, with $`T0`$. Then, the aging regime is the straight line with a slope $`X^10.21`$, to be compared with the theoretical value $`1/X_{QFD}=2`$. The inset shows the derivative $`X^1(b)=\text{d}b/\text{d}(t)`$, stepping from 0 to 2.1. FIGURE 2. Dynamical energy $`(t)+2`$ vs time, in log-log coordinate, for $`h=0.025`$ and $`h=0.2`$. The power-law decay is unambiguous, and a fit to $`\kappa =0.67`$ has been done between the two vertical arrows. FIGURE 3. Logarithmic derivative $`\kappa `$ of the energy $`(t)+2`$, for $`h=0.025,0.5,1.0,2.0`$. The curve is noisy as $`(t)2`$. The straight line stands for $`\kappa =2/3`$ which we believe to be its exact value. Curves for $`h=0.025,0.05`$ seems to tend to 2/3 from above, while $`h=0.2`$ seems to tend to 2/3 from below. $`\kappa 2/3`$ is well realised for $`h=1`$. FIGURE 4. The characteristic times $`t_f(t)t^\alpha `$ and $`t_a(t),a=0.55`$ and $`a=0.45`$ determined from equations (32), (33) and (34). Results are shown for the time steps $`h=0.1`$ and $`h=0.2`$, and the finiteness of $`h`$ is visible at small $`t`$. A numerical estimate of $`\alpha `$ is $`0.64`$ between the first and last vertical arrows. The exponent of $`t_{a=0.55}`$ is close to $`0.93`$ while we expect 1, and $`t_{a=0.45}`$ should saturate to a constant. FIGURE 5. The functions $`\mathrm{exp}[b(t,0)]`$, $`\mathrm{exp}[b(t,tt_f(t))]`$ vs $`t`$, for $`h=0.2`$ and $`h=0.1`$, in logarithmic coordinates. The functions $`\mathrm{exp}[b(t,20)]`$ vs $`t20`$ and $`\mathrm{exp}[b(t,40)]`$ vs $`t40`$. Here, 20 and 40 are waiting times. The behaviour of $`b(t,0)`$ and $`b(t,tt_f(t))`$ is doubtless logarithmic. The slopes of $`\mathrm{exp}[b(t,0)]`$, $`\mathrm{exp}[b(t,tt_f(t))]`$, on this figure are respectively $`1.10`$ and $`0.42`$. According to the predictions of , $`\mathrm{exp}(b(t,t^{}))`$ tends to $`h(t)h(t^{})`$ for $`t,t^{}1`$; $`t/t^{}`$ finite. The curves $`\mathrm{exp}[b(t,20)]`$ and $`\mathrm{exp}[b(t,40)]`$ tend to imitate $`\mathrm{exp}[b(t,0)]`$, with a delay. FIGURE 6. A test of the linear response of the displacement $`u(t)`$. We plot $`u(t)/F`$, as a function of $`\mathrm{ln}(t)`$, for $`F=0.05,0.1,0.2`$ and $`0.4`$; $`h=0.2`$. As $`F0`$, the curves are indistinguishable from the integrated response $`(t)`$. A departure from the straight line signals the breakdown of the linear response, as the particle acquires a finite velocity, dependent (non-linearly) on the force. The suggested behaviour of $`u(t)`$ is thus: $`u(t)=F(c_3+c_4\mathrm{ln}(t))`$. FIGURE 7. The density $`\rho (\lambda )`$ of eigenvalues $`\lambda \mathrm{SS}`$ (on top). A sketch of the self-similar distribution $`g=\mathrm{SS}\widehat{G}(\lambda /\mathrm{SS})`$, assuming $`\mathrm{\Gamma }=\mathrm{SS}`$ (bottom). The tail of $`\widehat{G}`$ goes to 0 as $`\lambda /\mathrm{SS}\mathrm{}`$ , in the asymptotic limit $`\mathrm{SS}0`$.
warning/0005/cond-mat0005081.html	ar5iv	text	# Macroscopic resonant tunneling of magnetic flux ## I Introduction Bose condensates of Cooper pairs in superconductors have a remakable ability to populate a single quantum state with a macroscopically large number of particles. This property of the Cooper pair condensates leads to macroscopic quantum effects in the dynamics of Josephson junctions, i.e. makes it possible for a Josephson junction to behave as a pure quantum state of a simple quantum mechanical system while still containing macroscopically large number of Cooper pairs. Due to the strongly nonlinear character of the Josephson dynamics even in the classical regime, macroscopic quantum effects can be quite non-trivial. They have been known for the past twenty years – see reviews in , and are continuing to attract considerable interest . This interest is to a large extent stimulated by the challenge to our understanding of the foundations of the quantum theory presented by direct manifestations of quantum mechanics at the macroscopic level . Another recent motivation for studying quantum effect in Josephson dynamics is provided by possible applications to quantum computation . The most advanced macroscopic quantum effect observed experimentally up to now is resonant tunneling between quantized energy levels in the adjacent wells of the Josephson potential . The aim of this work is to develop a theory of this phenomenon in SQUID systems, where the potential contains two wells each with a different value of the average magnetic flux. We consider the regime of weak energy dissipation important for studying coherent effects in resonant tunneling. This regime has not been discussed appropriately in the existing treatments of resonant tunneling in double-well potentials or multi-well potentials corresponding to the current-biased Josephson junctions . The most essential new feature of our approach is an account of the two types of relaxation mechanisms, intrawell and interwell, that exist in the system. The two relaxation mechanisms are very different in their dependence on the parameters of the SQUID potential, and lead to different shapes of the resonant tunneling peaks. Differences in relaxation mechanisms also make macroscopic resonant tunneling of flux different from the otherwise very similar “mesoscopic” resonant tunneling between charge states of small Josephson junctions and electron states in quantum dots . Another new element of this work is the discussion of the “photon-assisted” macroscopic resonant flux tunneling under rf irradiation. We show that in contrast to tunneling under stationary-bias conditions, the peaks of the photon-assisted tunneling depend qualitatively on the strength of decoherence in the flux dynamics. In the case of coherent flux dynamics, the resonant peaks of the photon-assisted tunneling are split in two. The splitting reflects the coherent hybridization of the macroscopic flux states in the two wells of the SQUID potential and is suppressed with increasing rate of decoherence. Very recently, such a splitting of the resonant flux-tunneling peaks has been observed experimentally , demonstrating the quantum coherence of the macroscopically distinct flux states. The paper is organized as follows. In Sec. 2 we derive the evolution equations for the density matrix describing the resonant flux tunneling under stationary-bias conditions, and introduce the two relaxation mechanisms for tunneling dynamics. Using these equations, we calculate the rate of flux tunneling in Sec. 3. In Sec. 4, we extend the results of Secs. 2 and 3 to the case of photon-assisted tunneling. ## II Equations for the density matrix To derive equations for the density matrix in the regime of resonant tunneling of magnetic flux $`\mathrm{\Phi }`$ we consider the standard model of the phase dynamics in SQUIDS. The combination of the magnetic energy of the SQUID loop biased with an external flux and the Josephson coupling energy of the SQUID junctions (for details, see e.g., ) produces the double-well potential $`U(\mathrm{\Phi })`$ for $`\mathrm{\Phi }`$ evolution (shown schematically in Fig. 1 below). The main part of the Hamiltonian governing the flux dynamics consists of the potential $`U(\mathrm{\Phi })`$ and the charging energy of the junction capacitance $`C`$: $$H_0=\frac{Q^2}{2C}+U(\mathrm{\Phi }).$$ (1) The charge $`Q`$ on the junction capacitance and the flux $`\mathrm{\Phi }`$ satisfy standard commutation relations $`[\mathrm{\Phi },Q]=i\mathrm{}`$. The two wells of the potential $`U(\mathrm{\Phi })`$ have discrete energy states $`\epsilon _{jn}`$ with characteristic energy separation on the order of $`\omega _j`$, where $`j=1,2`$ is the well index, $`\omega _j`$ are the oscillation frequencies around the potential minima, and $`n=0,1,\mathrm{}`$ numbers the states within each well. The two frequencies $`\omega _j`$ have the same order of magnitude $`\omega _1\omega _2\omega _p`$. External magnetic flux controls the energy difference between the states in opposite wells. Away from the resonance conditions, when all the energies $`\epsilon _{jn}`$ are separated by large energy gaps of order $`\omega _p`$, the states $`\|jn`$ are localized within the $`j`$th well, and the amplitude of the wavefunctions $`\psi _{jn}(\mathrm{\Phi })`$ in the opposite well is very small. However, when the energies of the two states $`\|1\|1n_1`$ and $`\|2\|2n_2`$ are close, $`\|\epsilon \|\omega _p`$, where $`\epsilon \epsilon _{1n_1}\epsilon _{2n_2}`$, these states become strongly coupled, and the wavefunctions spread over the both wells. As shown in the Appendix, strong coupling of the states $`\|1,2`$ at resonance can be described by the tunneling amplitude $`\mathrm{\Delta }`$, and the Hamiltonian (1) reduces to the regular two-state form in the basis of these states: $$H_0=\frac{1}{2}[\epsilon (\|11\|\|22\|)\mathrm{\Delta }(\|12\|+\|21\|)],$$ (2) $`\mathrm{\Delta }=(\omega _1\omega _2)^{1/2}D/\pi ,`$ where $`D`$ is the quantum mechanical transparency of the barrier separating the wells. Without perturbations, the two-state dynamics described by (2) is decoupled from the other states of the Hamiltonian (1). The most important perturbation creating such coupling is the energy dissipation that induces transitions between the states $`\|1,2`$ and other states $`\|jn`$. In the relevant temperature range below superconducting energy gap of the junction electrodes, the quasiparticle tunneling is suppressed and the main source of energy dissipation is the electromagnetic environment of the system. Under the assumption that the electromagnetic modes of the environment are in equilibrium at temperature $`T`$, and are well described by linear electrodynamics, the interaction between the flux $`\mathrm{\Phi }`$ and the heat bath of these modes can be written as $$V=I_f\mathrm{\Phi }.$$ (3) Here $`I_f`$ is the fluctuating current created in the SQUID loop by the environment, with the correlation function given by the fluctuation-dissipation theorem: $$I_f(t)I_f(t+\tau )=\frac{d\omega }{\pi }\frac{\omega G(\omega )e^{i\omega \tau }}{1e^{\omega /T}},$$ (4) where the brackets $`\mathrm{}`$ denote averaging over the equilibrium density matrix of the environment, and $`G(\omega )`$ is the dissipative part of the environment conductance. Equation (4) is sufficient to characterize completely the effects of the weak energy dissipation considered in this work. For arbitrary dissipation strength, one can use the Caldeira-Leggett model to express explicitly the environment Hamiltonian and the current operator $`I_f`$ in terms of a set of harmonic oscillators. The interaction (3) induces both “vertical” transitions within each well and direct interwell transitions. In terms of the two-state dynamics with the Hamiltonian (2), the latter correspond to the modulation of the tunneling amplitude $`\mathrm{\Delta }`$ by the environment. The matrix elements of this type of interwell transitions are, however, smaller by a factor of $`\mathrm{\Delta }/\omega _p`$ than those of the intrawell transitions. The small matrix elements can be neglected under the conditions of resonance, when the flux tunneling between the two wells is dominated by the stronger resonant processes. In this approximation, we can omit the terms in the flux operator in the interaction (3) that are non-diagonal in the well index $`j`$: $$\mathrm{\Phi }=\underset{j,n,n^{}}{}\mathrm{\Phi }_{n,n^{}}^{(j)}\|jnjn^{}\|.$$ (5) Here $`\mathrm{\Phi }_{n,n^{}}^{(j)}`$ are the matrix elements of $`\mathrm{\Phi }`$ in the $`\|jn`$ basis. The perturbation (3) with the flux operator (5) has two effects on the dynamics of the states $`\|1,2`$. The first is fluctuations of the energy difference $`\epsilon `$, which induce transitions between these states and lead to the loss of mutual coherence between them. The part of the $`\mathrm{\Phi }`$ operator (5) responsible for these fluctuations can be written as $`\delta \mathrm{\Phi }(\|11\|\|22\|)`$, where $`\delta \mathrm{\Phi }`$ is half of the difference between the average flux values in the states $`\|1`$ and $`\|2`$. The remaining terms in (5) induce intrawell transitions from the states $`\|1`$ and $`\|2`$ to the other states in the corresponding wells. For weak dissipation both effects can be described quantitatively by the standard density matrix technique – see, e.g., . The description starts from the equation for the evolution of the density matrix $`\rho `$, obtained treating the coupling $`V`$ (3) in second-order perturbation theory: $$\dot{\rho }(t)=i[H_0,\rho ]^t𝑑\tau [V(t),[V(\tau ),\rho (\tau )]].$$ (6) If the environment has a large cut-off frequency $`\omega _c\epsilon ,\mathrm{\Delta }`$, the density matrix evolves slowly on the time scale of variations of $`V(t)`$, and we can make the Markov approximation $`\rho (\tau )\rho (t)`$ in the last term of eq. (6). The condition of weak dissipation also allows us to keep only the dissipative terms in eq. (6) that do not oscillate in time with frequencies of the main Hamiltonian $`H_0`$, since only these terms lead to effects that accumulate with time. Using these approximations to evaluate the dissipative part of the eq. (6), we obtain the final equation for the density matrix in the basis of resonant states $`\|1,2`$ relevant for the transfer of flux between the wells: $$\dot{\rho }=i[H_0,\rho ]+\mathrm{\Gamma }[\rho ]+\gamma [\rho ].$$ (7) The term $`\mathrm{\Gamma }`$ in this equation describes the effect of the intrawell transitions from the states $`\|1`$ and $`\|2`$: $$\mathrm{\Gamma }[\rho ]=\left(\begin{array}{cc}\mathrm{\Gamma }_1\rho _{11},& (\mathrm{\Gamma }_1+\mathrm{\Gamma }_2)\rho _{12}/2\\ (\mathrm{\Gamma }_1+\mathrm{\Gamma }_2)\rho _{21}/2,& \mathrm{\Gamma }_2\rho _{22}\end{array}\right).$$ (8) At temperatures smaller than energy separation in the wells, the total decay rate $`\mathrm{\Gamma }_j`$ in eq. (8) of the state $`\|j`$ into the states with lower energy in the same well is: $`\mathrm{\Gamma }_j={\displaystyle \underset{n<n_j}{}}{\displaystyle \frac{2\|\mathrm{\Phi }_{n,n_j}^{(j)}\|^2}{\mathrm{}^2}}(\epsilon _{n_j}\epsilon _n)G(\epsilon _{n_j}\epsilon _n).`$ The second dissipation term $`\gamma [\rho ]`$ in eq. (7) describes transitions and decoherence within the $`\|1`$, $`\|2`$ subspace, and has a simple form in the basis of energy eigenstates of the two-state Hamiltonian (2): $$\gamma [\rho ]=U^{}\left(\begin{array}{cc}\gamma _ur_{11}\gamma _dr_{22},& (\gamma +\frac{\gamma _d+\gamma _u}{2})r_{12}\\ (\gamma +\frac{\gamma _d+\gamma _u}{2})r_{21},& \gamma _dr_{22}\gamma _ur_{11}\end{array}\right)U.$$ (9) Here $`r`$ is the density matrix in the eigenstate basis: $`R=U\rho U^{}`$, and $`U`$ is the rotation matrix from this basis to the flux basis $`\|1`$, $`\|2`$: $`U=[(1\epsilon /\mathrm{\Omega })^{1/2}\sigma _z+(1+\epsilon /\mathrm{\Omega })^{1/2}\sigma _x]/\sqrt{2},`$ where the $`\sigma `$’s denote Pauli matrices, and $`\mathrm{\Omega }(\epsilon ^2+\mathrm{\Delta }^2)^{1/2}`$. The transition rates $`\gamma _{d,u}`$ and the decoherence rate $`\gamma `$ are: $`\gamma _d={\displaystyle \frac{g\mathrm{\Delta }^2}{\mathrm{\Omega }}}{\displaystyle \frac{1}{1e^{\mathrm{\Omega }/T}}}\gamma _u=\gamma _de^{\mathrm{\Omega }/T},\gamma =2gT{\displaystyle \frac{\epsilon ^2}{\mathrm{\Omega }^2}},`$ with the dimensionless parameter $`g=2G(\delta \mathrm{\Phi })^2/\mathrm{}`$ characterizing the strength of the interwell relaxation, and we assumed that $`G(\omega )`$ is constant in the small-frequency range $`\omega \mathrm{\Omega }`$. Equation (7), with the dissipation terms (8) and (9), is used below to describe resonant tunneling of flux between the two wells in various regimes. Before doing this, we discuss the relative magnitude of the two dissipation terms in this equation. Since the width of the wells is of the same order of magnitude as the barrier between them, the magnitude of the flux matrix elements of the intrawell relaxation $`\mathrm{\Gamma }[\rho ]`$, and of the interwell relaxation $`\gamma [\rho ]`$ (determined, respectively, by the “width” of the wavefunctions inside the wells and the distance between the wells) should be close. The main difference between the two relaxation mechanisms is that the intrawell transitions dissipate energy $`\omega _j`$, whereas the interwell relaxation $`\gamma [\rho ]`$ involves only much smaller energies on the order of $`\epsilon ,\mathrm{\Delta },T`$. This means that the intrawell relaxation typically dominates the flux-tunneling dynamics. In particular, even under the assumed condition of weak relaxation (which for $`\mathrm{\Gamma }[\rho ]`$ means that the rates $`\mathrm{\Gamma }_j`$ are small compared to the oscillation frequencies $`\omega _j`$) the rates $`\mathrm{\Gamma }_j`$ can still be much larger than the frequencies $`\epsilon ,\mathrm{\Delta }`$ of the two-state dynamics. The interwell relaxation $`\gamma [\rho ]`$ will generally only be stronger than the intrawell relaxation if the environment has relatively low cut-off frequency $`\omega _c\omega _p`$. ## III Stationary bias In this Section, we calculate the rate of the resonant tunneling between the wells in the situation when the external flux through the SQUID loop does not contain an ac component, and the energy dissipation drives the initial flux state in the left well towards equilibrium. At low temperatures $`T\omega _p`$ the flux stays in the ground state of the left well and tunnels into the right well out of this state (Fig. 1). In this case, the relaxation rate $`\mathrm{\Gamma }_1`$ obviously vanishes. We begin by considering the situation with only the intrawell relaxation present. Equation (7) can then be written in matrix elements as: $`\dot{\rho }_{11}`$ $`=`$ $`\mathrm{\Delta }\text{Im}\rho _{12},\dot{\rho }_{22}=\mathrm{\Delta }\text{Im}\rho _{12}2\mathrm{\Gamma }\rho _{22},`$ (10) $`\dot{\rho }_{12}`$ $`=`$ $`(i\epsilon +\mathrm{\Gamma })\rho _{12}+(i\mathrm{\Delta }/2)(\rho _{22}\rho _{11}),`$ (11) with $`\mathrm{\Gamma }_2/2\mathrm{\Gamma }`$ in this Section. After transformation to the real and imaginary parts of the off-diagonal matrix element $`\rho _{12}`$ and the sum/difference $`\rho _{11}\pm \rho _{22}`$ of the diagonal elements of the density matrix, eqs. (11) can be solved directly with the initial condition that the flux is in the left well at time $`t=0`$, $`\rho _{11}(0)=1`$: $`\rho _{11}(t)={\displaystyle \frac{1}{2}}{\displaystyle \frac{e^{\mathrm{\Gamma }t}}{\omega ^2+\lambda ^2}}[(\mathrm{\Gamma }^2+\omega ^2)((1+{\displaystyle \frac{\lambda ^2}{\mathrm{\Gamma }^2}})\mathrm{cosh}\lambda t+`$ (12) $`2{\displaystyle \frac{\lambda }{\mathrm{\Gamma }}}\mathrm{sinh}\lambda t)+(\lambda ^2\mathrm{\Gamma }^2)((1{\displaystyle \frac{\omega ^2}{\mathrm{\Gamma }^2}})\mathrm{cos}\omega t2{\displaystyle \frac{\omega }{\mathrm{\Gamma }}}\mathrm{sin}\omega t)],`$ (13) $`\rho _{22}(t)={\displaystyle \frac{1}{2}}{\displaystyle \frac{e^{\mathrm{\Gamma }t}}{\omega ^2+\lambda ^2}}[(\mathrm{\Gamma }^2+\omega ^2)(1{\displaystyle \frac{\lambda ^2}{\mathrm{\Gamma }^2}})\mathrm{cosh}\lambda t+`$ (14) $`(\lambda ^2\mathrm{\Gamma }^2)(1+{\displaystyle \frac{\omega ^2}{\mathrm{\Gamma }^2}})\mathrm{cos}\omega t].`$ (15) In eq. (15), the eigenfrequencies $`\omega `$ and $`\lambda `$ of the system of equations (11) are: $$\omega ,\lambda =\left[(\frac{(\mathrm{\Omega }^2\mathrm{\Gamma }^2)^2}{4}+\mathrm{\Gamma }^2\epsilon ^2)^{1/2}\pm \frac{\mathrm{\Omega }^2\mathrm{\Gamma }^2}{2}\right]^{1/2}.$$ (16) Equations (15) contain all the information about dynamics of the flux tunneling. When the relaxation rate $`\mathrm{\Gamma }_2`$ is much smaller than the oscillation frequency $`\mathrm{\Omega }`$, the tunneling process consists of the weakly damped coherent oscillations of the flux between the wells followed by relaxation in the right well. With increasing relaxation rate the oscillation part of this process becomes increasingly more damped and turns into incoherent jumps of the flux from the left into the right well represented by the non-oscillatory exponential decay of $`\rho _{11}`$. The matrix element $`\rho _{11}(t)`$ has the meaning of the probability for the flux to remain in the left well at time $`t`$. In the case of monotonous decay of $`\rho _{11}`$, its derivative $`f(t)=[\rho _{11}(t)]^{}`$ gives the probability density of the time for flux tunneling between the wells. This probability density can be used to find all statistical characteristics of the flux tunneling process. For instance, the average time $`\tau `$ it takes for the flux to tunnel can be calculated as $`\tau =_0^{\mathrm{}}𝑑ttf(t)=_0^{\mathrm{}}𝑑t\rho _{11}(t)`$. In the coherent regime, when $`\rho _{11}(t)`$ oscillates in time, $`f(t)`$ can be negative and cannot be interpreted as probability density, and the question of how to define the tunneling time $`\tau `$ becomes non-trivial. To define $`\tau `$, one needs to establish the event that terminates the tunneling process. The definition adopted below assumes that the tunneling process ends when the flux makes the transition from the state $`\|2`$ into one of the lower energy states in the right well. This definition is motivated by the fact that such a transition eliminates the possibility for the flux to return into the left well. With such definition, the time the flux spends in the state $`\|2`$ is included in the tunneling time $`\tau `$ which then should be calculated as $$\tau =_0^{\mathrm{}}𝑑t(\rho _{11}(t)+\rho _{22}(t)).$$ (17) From equations (15) and (17) we obtain the tunneling rate: $$\tau ^1=\frac{\mathrm{\Delta }^2\mathrm{\Gamma }_2}{2\mathrm{\Delta }^2+\mathrm{\Gamma }_2^2+4\epsilon ^2}.$$ (18) Equation (18) describes the Lorentzian peak of the resonant flux tunneling. It shows that the resonant flux tunneling under stationary-bias conditions does not allow one to distinguish qualitatively between the regimes of coherent and incoherent flux tunneling since the shape of the resonance peak (18) remains the same regardless of the magnitude of the relaxation/decoherence rate $`\mathrm{\Gamma }`$. The average tunneling rate $`\tau ^1`$ can be calculated without explicit solution of the time-dependent equations for the density matrix. Instead of attempting to describe an individual tunneling event with the time-dependent solution, we can consider a large number of these events, assuming that after each transition from the left to the right well the system is immediately returned back to its initial state and the process is repeated. This immediate return means that the system is effectively decaying from the resonant state in the right well directly into the initial state in the left well, and can be modeled by adding the term $`\mathrm{\Gamma }_2\rho _{22}`$ into the equation for $`\rho _{11}`$. With such a modification, eq. (7) has a non-trivial stationary solution $`\rho ^{(0)}`$, and the tunneling rate $`\tau ^1`$ defined by eq. (17) can be found from this solution as $$\tau ^1=\mathrm{\Gamma }_2\rho _{22}^{(0)}.$$ (19) Since this method does not require solution of the time-dependent equations for the density matrix, it considerably simplifies the calculation of the average tunneling rate. To illustrate this procedure, we consider first the same tunneling under the stationary bias conditions described by eqs. (11) with the term $`\mathrm{\Gamma }_2\rho _{22}`$ included into the equation for $`\rho _{11}`$. Solving the stationary equation for the off-diagonal element $`\rho _{12}`$ of the density matrix and plugging the solution into the equations for the diagonal elements, we get the simple rate equations: $$\dot{\rho }_{11}=\mathrm{\Gamma }^{}(\rho _{22}\rho _{11})+\mathrm{\Gamma }_2\rho _{22},\dot{\rho }_{22}=\dot{\rho }_{11},$$ (20) where the transfer rate $`\mathrm{\Gamma }^{}=\mathrm{\Delta }^2\mathrm{\Gamma }_2/(\mathrm{\Gamma }_2^2+4\epsilon ^2)`$ between the two wells can be viewed as the “Golden-rule” rate of transition with the matrix element $`\mathrm{\Delta }/2`$ into the state $`\|2`$ broadened by the relaxation $`\mathrm{\Gamma }_2`$. From the stationary solution of eq. (20) we find that $`\rho _{22}^{(0)}=\mathrm{\Gamma }^{}/(\mathrm{\Gamma }_2+2\mathrm{\Gamma }^{})`$, and see that eq. (19) indeed reproduces the tunneling rate (18). The tails of the resonant peak (18) at $`\epsilon \mathrm{\Delta },\mathrm{\Gamma }_2`$ allow for another simple interpretation. At large $`\epsilon `$, the wavefunction of the eigenstate of the two-state Hamiltonian localized in the left well has the probability amplitude $`\mathrm{\Delta }/2\epsilon `$ in the right well. The tunneling rate $`\overline{\mathrm{\Gamma }}`$ in this regime can be found then as the probability to be in the right well times the relaxation rate $`\mathrm{\Gamma }_2`$: $$\overline{\mathrm{\Gamma }}=\mathrm{\Gamma }_2\mathrm{\Delta }^2/4\epsilon ^2.$$ (21) This simple reasoning indeed reproduces the tails of the peak (18) and allows us to obtain an estimate of the tunneling rate between the resonances. As shown in the Appendix, the wavefunction amplitude in the right well between the resonances is $`\pi \mathrm{\Delta }/[2\omega _2\mathrm{sin}(\pi \epsilon /\omega _2)]`$. From this we can write: $$\tau ^1=\mathrm{\Gamma }_2\left(\frac{\pi \mathrm{\Delta }}{2\omega _2\mathrm{sin}(\pi \epsilon /\omega _2)}\right)^2.$$ (22) It should be noted that eq. (22) is only an estimate, since $`\mathrm{\Gamma }_2`$ and $`\mathrm{\Delta }`$ depend on $`\epsilon `$ for $`\epsilon \omega _2`$, and can be different from their values at resonance. However, at $`\epsilon \omega _2`$, they are constant, and eqs. (18) and (22) coincide for $`\epsilon \mathrm{\Delta },\mathrm{\Gamma }_2`$. In this range of $`\epsilon `$ the tunneling rate $`\tau ^1`$ changes as $`\epsilon ^2`$. If the interwell relaxation is non-negligible, we need to keep both relaxation terms in eq. (7). The assumption that the relaxation rates are small in comparison with $`\mathrm{\Delta }`$, allows us to use eq. (9) for the interwell relaxation and makes it convenient to consider the flux dynamics in the eigenstates basis. The stationary values of the off-diagonal elements of the density matrix $`r`$ in this basis are vanishing for weak relaxation. To find the diagonal elements of $`r`$, we transform eq. (8) for the intrawell relaxation (with $`\mathrm{\Gamma }_1=0`$ and added term $`\mathrm{\Gamma }_2\rho _{22}`$ in the evolution of $`\rho _{11}`$) into this basis. Neglecting rapidly oscillating terms, we see that the diagonal part of the weak intrawell relaxation is: $`\dot{r}_{11}=\mathrm{\Gamma }_2[{\displaystyle \frac{\epsilon }{2\mathrm{\Omega }}}+{\displaystyle \frac{1}{4}}(1+{\displaystyle \frac{\epsilon ^2}{\mathrm{\Omega }^2}})(r_{22}r_{11})],\dot{r}_{22}=\dot{r}_{11}.`$ Combining this expression with eq. (9) for the interwell relaxation, we find the stationary values of the diagonal elements of $`r`$. After transformation back to the flux basis we finally obtain the stationary element $`\rho _{22}^{(0)}`$ of the density matrix $`\rho `$ in the flux basis and the flux tunneling rate (19): $$\tau ^1=\frac{\mathrm{\Gamma }_2}{2}\frac{\mathrm{\Omega }\mathrm{coth}(\mathrm{\Omega }/2T)+\epsilon +\mu }{\mathrm{\Omega }\mathrm{coth}(\mathrm{\Omega }/2T)+\mu (1+2\epsilon ^2/\mathrm{\Delta }^2)},$$ (23) $`\mu \mathrm{\Gamma }_2/2g.`$ The parameter $`\mu `$ in eq. (23) can be interpreted as the energy at which the characteristic interwell relaxation rate (which increases with the energy difference between the two resonant states) becomes equal to the rate $`\mathrm{\Gamma }_2`$ of the relaxation in the right well. The tunneling rate (23) is plotted in Fig. 2 for zero temperature and several values of $`\mu `$. The interwell relaxation becomes stronger with decreasing $`\mu `$, making the the resonant-tunneling peak progressively more asymmetric. For negative bias $`\epsilon `$, when the state $`\|2`$ in the right well is higher in energy than the state $`\|1`$ in the left well, the interwell relaxation is suppressed for low temperatures, $`T\mathrm{\Delta }`$, and the tail of the resonant peak (23) coincides with that of the Lorentzian peak (18). On the other hand, for positive $`\epsilon `$, transitions from the state $`\|1`$ into $`\|2`$ are allowed, and if $`\mu \mathrm{\Delta }`$, the interwell relaxation dominates for all positive $`\epsilon `$. The tunneling rate decreases in this case only as $`1/\epsilon `$ with increasing $`\epsilon `$. At $`\epsilon \mathrm{\Delta }`$, and vanishing temperature $`T`$, the tunneling rate (23) determined by the interplay between the interwell and intrawell relaxation can be understood in terms of the competition between the two tunneling paths (see inset in Fig. 2). One is the direct decay within the right well out of the eigenstate of the two-state Hamiltonian localized predominantly in the left well, but with a small probability amplitude in the right well. The rate of this decay is $`\overline{\mathrm{\Gamma }}`$ (21). Another is the transition between the two eigenstates induced by the interwell relaxation that transfers the probability between the two wells and is followed by the intrawell decay out of the lower-energy eigenstate with the rate $`\mathrm{\Gamma }_2`$. The rate of the interwell transition between the eigenstates is $$\overline{\gamma }=\frac{g\mathrm{\Delta }^2}{\epsilon }.$$ (24) At $`\mu <\mathrm{\Delta }`$, the second path dominates, and for sufficiently large $`\epsilon `$, $`\epsilon \mathrm{\Delta }^2/\mu `$, the bottleneck of the relaxation process is the interwell transition between the eigenstates and the tunneling rate (23) becomes independent of the intrawell rate $`\mathrm{\Gamma }_2`$: $`\tau ^1=\overline{\gamma }`$. In general, the competition between the two types of transitions gives an expression for the tunneling rate, $`\tau ^1=\mathrm{\Gamma }_2(\overline{\gamma }+\overline{\mathrm{\Gamma }})/(\overline{\gamma }+\mathrm{\Gamma }_2)`$, that agrees with eq. (23) at $`\epsilon \mathrm{\Delta }`$. For large $`\mu `$, $`\mu \mathrm{\Delta }`$, the interwell relaxation is weak close to resonance, and for sufficiently small bias, $`\epsilon \mu `$, eq. (23) coincides with eq. (18) with $`\mathrm{\Gamma }_2\mathrm{\Delta }`$. However, at larger bias, $`\epsilon \mu `$, the interwell relaxation increases and the tunneling rate is again given by $`\overline{\gamma }`$ (24). Only when $`\mu `$ becomes comparable to the level separation $`\omega _p`$ in the wells, the interwell relaxation is completely negligible and the resonant tunneling peak has the Lorentzian shape for all relevant energies. Away from resonance, when $`\epsilon \omega _p`$, the interwell tunneling with the rate (24) corresponds to the transitions between the states localized in opposite wells that are the closest in energy, while “intrawell” tunneling (22) corresponds to the transition between the states that are at least next-nearest neighbors in energy. Although there is no qualitative difference between the two types of the transitions away from the resonance, they lead to very different shapes for the tunneling peaks close to resonance. ## IV Photon-assisted tunneling When the SQUID is irradiated with an external rf signal, the macroscopic resonant flux tunneling can go more effectively through one of the excited states in the left well of the SQUID potential rather than out of the ground state, since the amplitude of tunneling $`\mathrm{\Delta }/2`$ out of the excited state is much larger than the tunneling amplitude for the ground state. In this Section, we consider the situation when an rf signal of frequency $`\omega `$ resonantly couples the ground state $`\|0`$ in the left well of the potential to an excited state $`\|1`$ with energy $`E`$ in this well (Fig. 3). The energy $`E`$ is on the order of $`\omega _p`$, and the condition of the resonant excitation is that the detuning $`\nu =E\omega `$ is small, $`\nu \omega _p`$. If the amplitude $`a`$ of the rf excitation is also relatively small, $`a\omega _p`$, the off-resonant coupling to other states is not important, and the coupling between the states $`\|0`$ and $`\|1`$ can be described in the rotating-wave approximation. In this approximation, the terms in the coupling that oscillate rapidly (with frequencies on the order of $`\omega _p`$) are neglected, and the coupling Hamiltonian is written as: $$H_{rf}=\frac{a}{2}(\|01\|e^{i\nu t}+\|10\|e^{+i\nu t}).$$ (25) If the excited state $`\|1`$ in the left well is coupled resonantly with amplitude $`\mathrm{\Delta }/2`$ to a state $`\|2`$ in the right well that is shifted in energy by $`\epsilon `$ with respect to $`\|1`$, the total Hamiltonian for the flux dynamics in the basis of the three states $`\|0`$, $`\|1`$, and $`\|2`$ is: $$H_0=\left(\begin{array}{ccc}0& a/2& 0\\ a/2& \nu & \mathrm{\Delta }/2\\ 0& \mathrm{\Delta }/2& \nu \epsilon \end{array}\right).$$ (26) As in the previous Section, the average flux tunneling rate $`\tau ^1`$ can be calculated according to eq. (19) from the stationary density matrix $`\rho `$ of the system in the basis of states $`\|0`$, $`\|1`$, and $`\|2`$. Time evolution of $`\rho `$ is described by the same eq. (7) but with the Hamiltonian $`H_0`$ given now by eq. (26). We begin discussion of the dynamics of flux tunneling in the three-state system (26) with the case when the interwell relaxation $`\gamma [\rho ]`$ can be neglected, and the only mechanism of the energy relaxation in the system is the intrawell relaxation $`\mathrm{\Gamma }[\rho ]`$. This relaxation is characterized by the two rates $`\mathrm{\Gamma }_{1,2}`$ of the transitions from the states $`\|1,2`$ into the lower energy states of the left and right potential well, respectively (Fig. 3). An obvious generalization of eq. (8) for $`\mathrm{\Gamma }[\rho ]`$ to the three-state basis gives the off-diagonal part of eq. (7) for the time evolution of $`\rho `$ of the following form: $`\dot{\rho }_{01}`$ $`=`$ $`(i\nu \mathrm{\Gamma }_1/2)\rho _{01}+ia(\rho _{00}\rho _{11})/2i\mathrm{\Delta }\rho _{02}/2,`$ (27) $`\dot{\rho }_{12}`$ $`=`$ $`(i\epsilon +(\mathrm{\Gamma }_1+\mathrm{\Gamma }_2)/2)\rho _{12}+i\mathrm{\Delta }(\rho _{22}\rho _{11})/2ia\rho _{02}/2,`$ (28) $`\dot{\rho }_{02}`$ $`=`$ $`(i(\nu \epsilon )\mathrm{\Gamma }_2/2)\rho _{02}ia\rho _{12}/2i\mathrm{\Delta }\rho _{01}/2.`$ (29) Equations (29) allow us to express the off-diagonal elements of $`\rho `$ in terms of the diagonal ones in the stationary regime. Inserting the stationary values of the off-diagonal elements into the equations for the diagonal elements: $`\dot{\rho }_{00}`$ $`=`$ $`a\text{Im}\rho _{01}+\mathrm{\Gamma }_1\rho _{11}+\mathrm{\Gamma }_2\rho _{22},`$ (30) $`\dot{\rho }_{11}`$ $`=`$ $`a\text{Im}\rho _{01}+\mathrm{\Delta }\text{Im}\rho _{12}\mathrm{\Gamma }_1\rho _{11},`$ (31) $`\dot{\rho }_{22}`$ $`=`$ $`\mathrm{\Delta }\text{Im}\rho _{12}\mathrm{\Gamma }_2\rho _{22},`$ (32) we calculate the stationary probability $`\rho _{22}^{(0)}`$ and find the flux tunneling rate (19). Equations (31) were written under the assumption that the relaxation in the left well brings the system out of the state $`\|1`$ directly into the ground state $`\|0`$. Although this is strictly true only in the case when $`\|1`$ is the first excited state in the well, eqs. (31) can be also used to calculate the average flux tunneling rate in other situations. Indeed, when the rf signal drives the system into the state $`\|1`$ that is not the first excited state, the intermediate states in the left well that exist between the states $`\|0`$ and $`\|1`$ are populated by the process of relaxation out of $`\|1`$. If the flux tunneling out of these states is neglected (similarly to tunneling out of the state $`\|0`$), their effect on the average tunneling rate can be accounted for by inclusion of the occupation probabilities of these states in the normalization condition. Since these probabilities in the stationary regime are proportional to the stationary occupation probability $`\rho _{11}`$, this can be done through an additional factor $`\lambda `$ in the normalization condition for the state $`\|0`$, $`\|1`$, and $`\|2`$: $$\rho _{00}+\lambda \rho _{11}+\rho _{22}=1.$$ (33) To give an example, we can calculate the factor $`\lambda `$ assuming that the left well is parabolic in the relevant energy range. Then, the standard result for the linear relaxation of the harmonic oscillator (see, e.g., ) is that the oscillator makes the transitions only between the nearest-neighbor states and that the transition rate from the state $`\|m`$ into $`\|m1`$ is proportional to $`m`$. This means that in the stationary state $`\rho _{m1,m1}=m\rho _{m,m}/(m1)`$, and if the rf radiation drives the system into the $`n`$th excited state of the left well, then $`\lambda =n_{m=1}^n(1/m)`$. To avoid extra parameters, however, we assume from now on that the state $`\|1`$ is the first excited state in the left well, so that $`\lambda =1`$. Sufficiently compact analytical expressions for the tunneling rate can be obtained from eqs. (29) and (31) only in certain limits. For example, for small rf amplitude and weak relaxation, $`a\mathrm{\Gamma }_{1,2}\mathrm{\Delta }`$, we get: $$\tau ^1=\frac{\mathrm{\Gamma }_2a^2\mathrm{\Delta }^2}{(2\nu \epsilon \mathrm{\Omega })^2(2\nu \epsilon +\mathrm{\Omega })^2+4(\mathrm{\Gamma }_1(\nu \epsilon )+\mathrm{\Gamma }_2\nu )^2}.$$ (34) Equation (34) describes two peaks (discussed in more details below) in the dependence of the tunneling rate on the detuning $`\nu `$. The peaks are broadened by the relaxation, and their positions correspond to the two eigenstates of the Hamiltonian (2): $`\nu =(\epsilon \pm \mathrm{\Omega })/2`$. Another expression for the tunneling rate can be obtained for large relaxation rates $`\mathrm{\Gamma }_{1,2}a,\mathrm{\Delta }`$: $$\tau ^1=\frac{\mathrm{\Gamma }_2a^2\mathrm{\Delta }^2}{(4\nu ^2+\mathrm{\Gamma }_1^2)(4(\nu \epsilon )^2+\mathrm{\Gamma }_2^2)}.$$ (35) Depending on the relation between the energy bias $`\epsilon `$ and the relaxation rates, the tunneling rate (35) as a function of detuning $`\nu `$ contains either one (for small $`\epsilon `$) or two separate peaks (for large $`\epsilon `$). The peak positions in this case coincide with the position of the energy levels $`\|1,2`$ localized in the two wells. For arbitrary parameters it is convenient to use the stationary solution of eqs. (29) and (31) to plot the tunneling rate $`\tau ^1`$ numerically. To simplify the discussion, we assume that the relaxation rates in the two wells are the same, $`\mathrm{\Gamma }_1=\mathrm{\Gamma }_2\mathrm{\Gamma }`$. Figure 4 shows the dependence of $`\tau ^1`$ on detuning $`\nu `$ obtained in this way for several rf amplitudes $`a`$ in the regime of small relaxation rate $`\mathrm{\Gamma }`$. The main qualitative feature of Fig. 4 (that can also be seen in eq. (34)) is that for small $`a`$ the resonant peak in the tunneling rate $`\tau `$ is split into two peaks due to coherent oscillations of flux between the two wells. The appearance of such splitting can be easily understood, since a weak rf signal excites the system not into the state $`\|1`$ localized in the left well, but into the two hybridized states formed out of the states $`\|1`$ and $`\|2`$ in the two wells. The two different energies of the two hybridized states lead to the two peaks in the tunneling rate. This means that the splitting of the resonant-tunneling peak is the direct manifestation of the quantum coherent oscillations of flux between the two wells of the SQUID potential. Figure 4 illustrates also how the splitting of the resonant tunneling peak is suppressed with the increase of the rf amplitude. Suppression of the peak splitting can be understood in terms of the time $`1/\mathrm{\Delta }`$ required to establish the stationary states hybridized between the two wells. Large rf amplitude $`a`$ on the order of the tunnel amplitude $`\mathrm{\Delta }`$ causes rapid Rabi oscillation between the states $`\|0`$ and $`\|1`$ and does not allow for sufficient time to establish the two hybridized states. The system is therefore effectively excited into the state $`\|1`$ localized in the left well, and single resonant tunneling peak is formed around the energy of this state. Broadening and suppression on this single peak seen in Fig. 4 as $`a`$ increases beyond $`\mathrm{\Delta }`$, is a version of the generic “quantum Zeno” effect, when tunneling out of a metastable state is suppressed by rapid perturbation of this state. Figure 5 shows the evolution of the coherently-split tunneling peaks at small rf amplitude $`a`$ with the bias energy $`\epsilon `$ and with the relaxation rate $`\mathrm{\Gamma }`$. We see that with increasing energy bias (Fig. 5a), the peaks follow the position of the energy levels and the splitting between them increases. Simultaneously, the peak height decreases reflecting the overall suppression of the tunneling rate as one moves away from the resonance. Figure 5b shows how the double-peak structure in the small-bias regime representing the coherent mixing of the flux states in the two wells is suppressed by increasing relaxation rate $`\mathrm{\Gamma }`$. The structure is visible up to the relaxation rates $`\mathrm{\Gamma }0.5\mathrm{\Delta }`$. The peaks in Fig. 5 are shown only for $`\epsilon 0`$. The peak structure for negative $`\epsilon `$ can be understood from the “symmetry” relation $`\tau ^1(\nu ,\epsilon )=\tau ^1(\nu ,\epsilon )`$ that can be deduced from eqs. (29) and (31). Equations (29) show that in the stationary regime, changing the sign of $`\nu ,\epsilon `$ is equivalent to changing the sign of $`a,\mathrm{\Delta }`$ and replacing the off-diagonal elements of $`\rho `$ with their complex conjugate values. This transformation obviously does not change the transition rates in eq. (31), and therefore does not change the flux tunneling rate $`\tau ^1`$. Another interesting manifestation of the coherent flux tunneling between the wells can be seen in the dependence of the tunneling rate on the bias energy $`\epsilon `$ at fixed detuning $`\nu `$ (Fig. 6). For weak relaxation, the hybridized states are well-developed, and when the rf excitation energy lies between these two states, the tunneling rate is strongly suppressed. At $`\nu =0`$, when the system is excited precisely into the state $`\|1`$ localized in the left well, this condition is satisfied and tunneling rate is strongly suppressed for any energy bias $`\epsilon `$. As can be seen from eq. (34) and Fig. 6, in this case the tunneling rate as a function of $`\epsilon `$ is described by a Lorentzian centered around $`\epsilon =0`$. In contrast to resonant tunneling peaks in the $`\nu `$-dependence of the tunneling rate, which have small width proportional to the relaxation rate $`\mathrm{\Gamma }`$, the width of this Lorentzian is large, $`\mathrm{\Delta }^2/\mathrm{\Gamma }`$, and is inversely proportional to $`\mathrm{\Gamma }`$. When the detuning $`\nu `$ deviates from zero, there is an energy bias $`\epsilon `$ at which the excitation energy coincides with the energy of one of the hybridized states. The tunneling rate has a peak under such resonance conditions. For not-too-small $`\nu `$’s, this resonant peak again has a small width proportional to $`\mathrm{\Gamma }`$. In Fig. 6, one can see how the transition between the broad and narrow tunneling peaks takes place for $`\nu 0`$. As before, the results for negative detuning can be deduced from the relation $`\tau ^1(\nu ,\epsilon )=\tau ^1(\nu ,\epsilon )`$. Finally, we discuss the effect of weak interwell relaxation on the photon-assisted tunneling. We start by generalizing eq. (9) for this relaxation to the three-state situation relevant for the photon-assisted tunneling. Under the natural assumption that the average flux in the states $`\|0`$ and $`\|1`$ in the left well of the SQUID potential is the same, the part of the dissipative coupling (3) that corresponds to the interwell relaxation is: $$V=I_f\delta \mathrm{\Phi }\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right)I_f\delta \mathrm{\Phi }U.$$ (36) While eq. (36) is written in the flux basis $`\|0,\|1,\|2`$, weak relaxation is conveniently described in the basis of the eigenstates $`\|n`$ of the Hamiltonian (26). In this basis, the contribution of the interwell relaxation (36) to the evolution of the density matrix $`\rho `$ is given by the standard expression similar to eq. (9): $`\dot{\rho }_{nn}`$ $`=`$ $`{\displaystyle \underset{m}{}}(\gamma _{mn}\rho _{mm}\gamma _{nm}\rho _{nn}),`$ (37) $`\dot{\rho }_{nm}`$ $`=`$ $`[\gamma _{mn}^{}+{\displaystyle \frac{1}{2}}{\displaystyle \underset{k}{}}(\gamma _{nk}+\gamma _{mk})]\rho _{nm},nm.`$ (38) Transition and dephasing rates in these equations are: $`\gamma _{nm}={\displaystyle \frac{g\|U_{nm}\|^2(\epsilon _n\epsilon _m)}{1e^{(\epsilon _n\epsilon _m)/T}}},\gamma _{nm}^{}={\displaystyle \frac{gT}{2}}(U_{nn}U_{mm})^2,`$ where $`U_{nm}`$ are the matrix elements of the operator $`U`$ (36) in the eigenstates basis, and $`\epsilon _n`$ is the energy of the eigenstate $`\|n`$. Interwell relaxation can be included into the evolution equations for the density matrix on the basis of eq. (37) numerically. We diagonalize the Hamiltonian (26), calculate the interwell relaxation terms (37) in the eigenstates basis, and transfer them into the flux basis, where the intrawell relaxation has the simple form (29), (31). Calculating finally the stationary value of the density matrix $`\rho `$ we find the flux tunneling rate (19). Figures 7 and 8 show results of such a calculation obtained at vanishing temperature $`T`$. In Figure 7, the tunneling rate is plotted as a function of the detuning $`\nu `$ for $`\epsilon =0`$ and several values of the relative strength of the intrawell relaxation $`\mathrm{\Gamma }`$. The eigenstates of the two-state Hamiltonian at $`\epsilon =0`$ are symmetric between the two wells, and in absence of the interwell relaxation produce two symmetric resonant tunneling peaks (see Figs. 4 and 5a). As can be seen from Fig. 7, the interwell relaxation makes the tunneling peaks asymmetric. The positive-$`\nu `$ side of the double-peak structure corresponds to excitation of the system into the lower-energy eigenstate and is unaffected by the interwell relaxation at zero temperature, since there is no energy in this regime to create additional tunneling path. In contrast, the negative-$`\nu `$ side of the double-peak structure corresponds to excitation of the system into the eigenstate with larger energy, and the interwell relaxation increases the rate of tunneling out of this state. Because of this, the negative-$`\nu `$ peak in Fig. 7 is larger than the tunneling peak at positive $`\nu `$, and the tunneling rate at $`\nu <0`$ decreases much more slowly away from the peak than at positive $`\nu `$. Interwell relaxation introduces asymmetry also in the dependence of the flux tunneling rate on the bias energy $`\epsilon `$. Examples of such dependence are shown in Fig. 8 for two values of the detuning, $`\nu =\pm 2\mathrm{\Delta }`$. In both cases, the tunneling rate has a resonant peak at $`\epsilon \nu `$ similar to the peaks shown in Fig. 6 for vanishing interwell relaxation, when the peaks for the two values of detuning are symmetric. Comparison of Figs. 8a and 8b shows that the interwell relaxation makes the peaks asymmetric. In particular, the peak at $`\nu <0`$ is smaller than the peak at $`\nu >0`$. Although this asymmetry appears to be opposite to that in Fig. 7, where the negative-$`\nu `$ peak is larger, it has the same origin as in Fig. 7. The peaks at $`\nu <0`$ and $`\nu >0`$ correspond to resonant excitation of the system into, respectively, the upper and lower energy eigenstates. When the resonance occurs for $`\|\epsilon \|>\mathrm{\Delta }`$, as in Fig. 8, the eigenstates are already to a large extent localized in one or the other well. At $`\epsilon \nu >0`$, the lower eigenstate is centered in the right well and the interwell relaxation increases the tunneling rate, while at $`\epsilon \nu <0`$ the lower eigenstate is centered in the left well and the interwell relaxation brings the system back to this well suppressing the tunneling rate. As a result, the resonant tunneling peak in Fig. 8a ($`\nu >0`$) is larger than in Fig. 8b ($`\nu <0`$). The height of the negative-$`\nu `$ peak is more sensitive to the relative strength of the two relaxation mechanisms and decreases with decreasing rate $`\mathrm{\Gamma }`$ of the intrawell relaxation. The tails of the photon-assisted resonant peaks can also be described analytically. When the bias energy and detuning are not close to any resonance, $`\|\nu \|,\|\epsilon \|,\|\nu \epsilon \|a,\mathrm{\Delta }`$ both the interwell tunneling and rf excitation can be treated as perturbations. The dynamics of flux tunneling in this regime can be described as a coexistence of the two tunneling paths similar to the off-resonant tunneling discussed in Sec. 2 (see inset to Fig. 2). If $`\nu >\epsilon `$, the effective energy $`\nu \epsilon `$ of the state $`\|2`$ in the right well (i.e., the energy of this state brought down by a quantum of rf radiation, as in the Hamiltonian (26)) is above the energy of the initial state $`\|0`$ in the left well, and the only energy-allowed tunneling path is direct relaxation in the right well out of the perturbed state $`\|0`$. Similarly to eq. (21), perturbation theory in $`a,\mathrm{\Delta }`$ gives for the rate $`\overline{\mathrm{\Gamma }}`$ of this tunneling: $$\overline{\mathrm{\Gamma }}=\frac{\mathrm{\Gamma }_2a^2\mathrm{\Delta }^2}{16\nu ^2(\nu \epsilon )^2}.$$ (39) If $`\nu <\epsilon `$, the effective energy of the state $`\|2`$ is lower than the energy of the state $`\|0`$, and in addition to the tunneling (39) there is a competing tunneling process. It consists of a transition between the perturbed states $`\|0`$ and $`\|2`$ with the rate $`\overline{\gamma }`$, $$\overline{\gamma }=\frac{ga^2\mathrm{\Delta }^2}{4\nu ^2(\epsilon \nu )}.$$ (40) that is driven by interwell relaxation, followed by direct relaxation in the right well with the rate $`\mathrm{\Gamma }_2`$. As in Sec. 2, the coexistence of the two tunneling paths gives the following total tunneling rate: $$\tau ^1=\mathrm{\Gamma }_2\frac{\overline{\gamma }+\overline{\mathrm{\Gamma }}}{\overline{\gamma }+\mathrm{\Gamma }_2}.$$ (41) One can check that eqs. (39) and (41) agree, respectively, with the negative-$`ϵ`$ tail of the resonant tunneling peak in Fig. 8b and the positive-$`ϵ`$ tail of the peak in Fig. 8a. In summary, we have studied the effects of two types of relaxation mechanisms on the macroscopic resonant tunneling of flux in SQUIDs under stationary-bias conditions and with external rf irradiation. Coherent splitting of the resonant-tunneling peaks by rf radiation provides a convenient way of studying quantum coherence of flux states. This work was supported by ARO grant DAAD199910341. Appendix In the Appendix we show explicitly how the Hamiltonian of the two-well system (1) can be reduced at resonance to the two-state form (2), and derive an expression for the tunneling amplitude $`\mathrm{\Delta }`$. Assuming that the transparency $`D`$ of the barrier separating the two wells is small, $`D1`$, and that the resonance occurs between the states with large $`n`$, we can use the WKB approximation for the wavefunctions $`\psi _{jn}(\mathrm{\Phi })`$ of the Hamiltonian (1). In this approximation, the wavefunction between the right and the left turning points $`r`$ and $`l`$ is: $$\psi (\mathrm{\Phi })=\frac{A}{\sqrt{p}}\mathrm{cos}(w(\mathrm{\Phi })\delta ),$$ (42) with the WKB phase $`w(\mathrm{\Phi })=(1/\mathrm{})_l^\mathrm{\Phi }𝑑\mathrm{\Phi }^{}p(\mathrm{\Phi }^{})\pi /4`$ and momentum $`p=[2C(EU(\mathrm{\Phi }))]^{1/2}`$. In eq. (42), $`\delta `$ is a constant phase shift, and $`A`$ is a normalization constant. The phase $`w`$ is defined in such a way that $`\delta =0`$ for an isolated well, when the wavefunction decays exponentially in the classically inaccessible region $`\mathrm{\Phi }>r,\mathrm{\Phi }<l`$. When the two wells are coupled, the energy $`E`$ of the state common to them deviates from the eigenenergies $`\epsilon _n`$ of the isolated wells that are determined by the Bohr-Sommerfeld condition $`(1/\mathrm{})_l^rp𝑑\mathrm{\Phi }=2\pi (n+1/2)`$. For weak tunneling, this deviation is small in comparison to the state energy, and creates a small but non-vanishing phase shift $`\delta `$: $$\delta =\frac{1}{\mathrm{}}_l^r(p(E)p(\epsilon _n))=\pi \frac{E\epsilon _n}{\omega _p},$$ (43) where $`\omega _p`$ is the frequency of the classical oscillations in the well and which in the WKB approximation determines the spacing of the energy levels. At non-vanishing $`\delta `$, the wavefunction (42) has a part that grows exponentially in the classically inaccessible region, as one can see rewriting it as $$\psi =\frac{Ae^{i\delta }}{\sqrt{p}}\mathrm{cos}w(\mathrm{\Phi })\frac{iA\mathrm{sin}\delta }{\sqrt{p}}e^{iw(\mathrm{\Phi })}.$$ (44) According to general rules of the WKB approximation – see, e.g., , the two terms in this expression produce exponentially decaying and growing components of the wavefunctions with the amplitudes $`Ae^{i\delta }/2`$ and $`A\mathrm{sin}\delta `$, respectively. To find the energy $`E`$ we need to match the amplitudes of the wavefunctions of the right and left wells in the barrier region. Under the conditions of resonance, equating the amplitudes and keeping only the terms of the first order in $`D`$ we get: $$A_1\frac{E\epsilon _1}{\omega _1}=A_2\frac{D}{2\pi },A_2\frac{E\epsilon _2}{\omega _2}=A_1\frac{D}{2\pi },$$ (45) where $`\epsilon _j`$, $`j=1,2`$ are the energies of the resonant states, and $`D=\mathrm{exp}\{(1/\mathrm{})_{r_1}^{l_2}\|p\|𝑑\mathrm{\Phi }\}`$. The probability to be in the right/left well is directly related to the wavefunction amplitudes $`A_j`$: $$_{l_j}^{r_j}𝑑\mathrm{\Phi }\|\psi _j\|^2=\frac{B\|A_j\|^2}{\omega _j},$$ (46) where $`B`$ is the $`j`$-independent part of the normalization constant. Introducing the amplitudes $`\alpha `$ of this probability, $`\alpha _j=A_j/\omega _j^{1/2}`$, we bring eq. (45) into the form that coincides with the Schrödinger equation of the two-state system: $$E\alpha _1=\epsilon _1\alpha _1\frac{\mathrm{\Delta }}{2}\alpha _2,E\alpha _2=\epsilon _2\alpha _2\frac{\mathrm{\Delta }}{2}\alpha _1,$$ (47) with the tunneling amplitude $`\mathrm{\Delta }={\displaystyle \frac{(\omega _1\omega _2)^{1/2}}{\pi }}D.`$ Since the resonant states $`\|1,2`$ are orthogonal to all other states of the full Hamiltonian (1), this proves that their dynamics can be described by the Hamiltonian (2). Away from resonance, when $`\epsilon \epsilon _1\epsilon _2\omega _p`$, one of the amplitudes $`\alpha `$ (for the states localized in the left well, $`\alpha _2`$) is small, $`\alpha _2D`$. Keeping, as before, only the terms of the first order in $`D`$, we see from eq. (47) that in this case $`E_1=\epsilon _1`$, i.e. $`\delta _1=0`$. Matching the growing wavefunction in the right well with the decaying wavefunction in the left well we find that away from the resonance $$\alpha _2=\frac{\pi \mathrm{\Delta }}{2\omega _2\mathrm{sin}(\pi \epsilon /\omega _2)}.$$ (48) Equation (48) extrapolates smoothly between the successive resonances. Making use of the harmonic approximation for the potential $`U(\mathrm{\Phi })`$, one can show that the results of this Appendix (eqs. (47) and (48)) can be extended to the low-lying states with small $`n`$ which cannot be described with the WKB approximation. Such an extension leads only to minor modifications in the definition of the barrier transparency $`D`$ in these equations.
warning/0005/cond-mat0005465.html	ar5iv	text	# Giant thermoemf in multiterminal superconductor/normal metal mesoscopic structures. ## Abstract We considered a mesoscopic superconductor/normal metal (S/N) structure in which the N reservoirs are maintained at different temperatures. It is shown that in the absence of current between the N reservoirs a voltage difference $`V_T`$ arises between the superconducting and normal conductors. The voltage $`V_T`$ oscillates with increasing phase difference $`\phi `$ between the superconductors, and its magnitude does not depend on the small parameter $`(T/ϵ_F).`$ It is well known that if the terminals of a normal conductor are maintained at different temperatures, then in the absence of a current a thermoelectric voltage ($`V_{emf})`$ appears between the terminals. The magnitude of $`V_{emf}`$ is equal to $`c_1(T/ϵ_F)\delta T/e,`$ where $`c_1`$ is a factor of the order 1, $`ϵ_F`$ is the Fermi energy and $`\delta T`$ is the temperature difference (see for example ). In this paper we analyse the thermoelectric effect in the mesoscopic structure shown in Fig.1. We show that when a temperature difference exists between the normal (N) reservoirs, a voltage between normal and superconducting circuits $`V_T`$ appears. Unlike in the normal case, the magnitude of this voltage does not depend on a small parameter $`(T/ϵ_F)`$, and also this voltage oscillates as the phase difference $`\phi `$ between the superconductors varies. We assume that the superconductors are connected via a superconducting loop and the phase difference between them $`\phi `$ is controlled by an applied magnetic field. There is no current between the N reservoirs and the temperatures of the reservoirs are different: $`T(\pm L)=T_o\pm \delta T`$. We will calculate the electric potential in the N film and, in particular, the potential $`V_T`$ in the N reservoirs. Since we set the potential in the superconductors equal to zero, the potential $`V_T`$ is the voltage difference between the N reservoirs and superconductors which arises in the presence of the temperature difference $`\delta T`$. In order to find the potential $`V_T`$, we need to determine the distribution functions $`f_\pm `$ and the condensate wave functions $`\widehat{F}^{R(A)}`$ induced in the N film. The distribution functions $`f_\pm `$ are related to the ordinary distribution functions for electrons $`n_{}`$ and holes $`p_{}`$: $`f_+=f_+=1(n_{}+`$ $`p_{})`$ and $`f_{}=f_{}=(n_{}`$ $`p_{})`$ (we assume that there is no spin-dependent interaction in the system). The function $`f_+`$ determines the condensate current and the function $`f_{}`$ determines the quasiparticle current and electrical potential (see for example Ref. ,where $`f_+`$ and $`f_{}`$ are denoted by $`f`$ and $`f_{1\text{ }}`$respectively, and Ref. where the application of the Green’s function technique to the study of transport in S/N mesoscopic structures is discussed). These functions satisfy the kinetic equation (see Ref. ) $$L_x[M_\pm _xf_\pm (x)+J_Sf_{}(x)\pm J_{an}_xf_{}(x)]=r[\stackrel{}{A_\pm }\delta (xL_1)+\stackrel{\mathrm{\_}}{A_\pm }\delta (x+L_1)].$$ (1) where all the coefficients are expressed in terms of the retarded (advanced) Green’s functions: $`\widehat{G}^{R(A)}=G^{R(A)}\widehat{\sigma }_z+\widehat{F}^{R(A)};`$ $`M_\pm =(1G^RG^A(\widehat{F}^R\widehat{F}^A)_1)/2;`$ $`J_{an}=(\widehat{F}^R\widehat{F}^A)_z/2,J_s=(1/2)(\widehat{F}^R_x\widehat{F}^R\widehat{F}^A_x\widehat{F}^A)_z,`$ $`A_\pm =(\nu \nu _S+g_1)(f_\pm f_{S\pm })(g_{z\pm }f_S+g_zf_{});`$ $`g_{1\pm }=(1/4)[(\widehat{F}^R\pm \widehat{F}^A)(\widehat{F}_S^R\pm \widehat{F}_S^R)]_1;`$ $`g_{z\pm }=(1/4)[(\widehat{F}^R\widehat{F}^A)(\widehat{F}_S^R\pm \widehat{F}_S^R)]_z.`$ The parameter $`r=R/R_b`$ is the ratio of the resistance of the N wire $`R`$ and S/N interface resistance $`R_b`$; the functions $`\stackrel{\mathrm{\_}}{A_{}}`$ and $`\underset{+}{\overset{\mathrm{\_}}{A}}`$ coincide with $`\stackrel{}{A_{}},\underset{+}{\overset{}{A}}`$ if we make a substitution $`\phi \phi `$. We introduced above the following notations $`(\widehat{F}^R\widehat{F}^A)_1=Tr(\widehat{F}^R\widehat{F}^A),`$ $`(\widehat{F}^R\widehat{F}^A)_z=Tr(\widehat{\sigma }_z\widehat{F}^R\widehat{F}^A)`$ etc.; $`\nu ,`$ $`\nu _S`$ are the density-of states in the N film at $`x=L_1`$ and in the superconductors. The functions $`f_{S\pm }`$ are the distribution functions in the superconductors which are assumed to have equilibrium forms. This means that $`f_{S+}f_{eq}=\mathrm{tanh}(ϵ\beta _o)`$ and $`f_S=0,`$ because we set the potential of the superconductors equal to zero (no branch imbalance in the superconductors). At the reservoirs the distribution functions $`f_\pm `$ obey the boundary conditions: $`f_\pm (\pm L)=F_{V\pm }(\pm L),`$ where $`F_{V\pm }(L)=(1/2)[\mathrm{tanh}(\beta _o(ϵ+eV_T))\pm \mathrm{tanh}(\beta _o(ϵeV_T))],\beta _o=(2T_o)^1.`$ In order to clarify the physical meaning of different terms in Eq.(1), consider the equation for $`f_{}`$. The first term in this equation is the partial quasiparticle current (the quasiparticle current at a given energy) in the N film . The second term is the condensate current and the third term is also the condensate current which appears under nonequilibrium conditions. The factors $`A_{}`$ and $`\stackrel{\mathrm{\_}}{A_{}}`$ are the partial currents through the S/N interfaces. The first term in $`A_{}`$ is the partial quasiparticle current above (the term $`\nu \nu _Sf_{}`$) and below (the term $`g_1f_{}`$) the gap $`\mathrm{\Delta },`$ and the last term is the condensate current. The factors $`A_+`$ and $`\stackrel{\mathrm{\_}}{A_+}`$ are equal to zero below the gap (complete Andreev reflection). In this paper we consider the case when the S/N interface resistance is greater than or approximately equal to the resistance of the N wire ($`r<1`$). However $`r`$ should not be too small because in Eq.(1) we neglected the inelastic collision integral. This implies that the condition $$(ϵ\tau _ϵ)^1<<r1$$ (2) must be fulfilled; here $`\tau _ϵ`$ is the energy relaxation time, $`ϵ\mathrm{min}\{T,ϵ_L\},`$ and $`ϵ_L=D/L^2`$ is the Thouless energy. We consider the most interesting case of low temperatures ($`T<<\mathrm{\Delta }`$) when for characteristic energies $`ϵ<<\mathrm{\Delta }`$ the functions $`A_+`$ and $`\stackrel{\mathrm{\_}}{A_+}`$ equal zero as in this case $`\nu _S0`$ and $`\widehat{F}^R\widehat{F}^A.`$ The first integration of Eq.(1) yields $$M_+_xf_+(x)+J_Sf_{}(x)+J_{an}_xf_{}(x)=J_+$$ (3) where $`J_+`$ is a constant of integration. In the limit of the small parameter $`r`$, all the corrections related to the proximity effect are small (the functions $`\widehat{F}^{R(A)}\stackrel{}{\text{are proportional to }r\text{). Therefore}}`$ in the main approximation in the parameter $`r,`$ we have for the distribution function $`f_+(x)`$ $$f_+(x)\delta \beta ϵ(x/L)\mathrm{cosh}^2(\beta _oϵ)+f_{eq}$$ (4) In obtaining Eq.(4) we have taken into account that in the N reservoirs the function $`f_+(\pm L)`$ has an equilibrium form with different temperatures:$`f_+(\pm L)=f_{eq}\pm \delta \beta ϵ\mathrm{cosh}^2(\beta _oϵ),`$ where $`\delta \beta /\beta \delta T/T,\beta _o=(2T_o)^1`$. Eq.(4) implies that the temperature gradient leads to a flow of nonequilibrium electrons and holes in the N wire. In order to find the distribution function $`f_{}(x)`$, we integrate Eq.(1) $$M_{}_x\stackrel{}{f_{}(x)}+J_Sf_+(x)J_{an}_xf_+(x)=J_1\mathrm{\Theta }(L_1x)+J_{2+,2}\mathrm{\Theta }(xL_1)$$ (5) The constants $`J_1`$ and $`J_{2+,2}`$ are the partial currents in regions $`(L_1,+L_1)`$ , $`(L_1,L)`$ and $`(L,L_1)`$ respectively. Outside the interval $`(L_1,+L_1)`$ the ”supercurrent” $`J_S`$ is zero (this follows directly from the expressions for $`J_S`$ and for $`\widehat{F}^R`$ ). As there is no current between the N reservoirs, the integrals over energy $`ϵ`$ from $`J_{2+,2}`$ should be equal to zero. In the absence of the temperature gradient we obtain from Eq.(5) $`f_{}=0`$ and $`J_Sf_{eq}=J_{1eq}=r(g_z+g_{z+})f_{eq}.`$ According Eq.(1) the constants $`J_1`$ and $`J_{2+,2}`$ are related to each other ( Kirchoff’s law) $$J_{2+,2}\delta J_1=\pm r[g_{1+}f_{}g_{z+}\delta f_+]_{(\pm L_1)}$$ (6) where $`\delta J_1=J_1J_S`$ and $`\delta f_+=f_+f_{eq}.`$ Integrating Eq.(5) outside the interval $`(L_1,+L_1)`$ and taking into account the boundary condition for $`f_{}(x)`$ , we obtain $`J_{2+,2}=\pm [f_{}(\pm L)f_{}(\pm L_1)]/L_2`$. Using this expression and Eq.(6), one can show that in the main approximation in $`r`$ the function $`f_{}(x)`$ is almost constant along the N wire and equal to $`f_{}(x)f_{}(0)f_{}(\pm L)=eV_T\beta \mathrm{cosh}^2(\beta _oϵ).`$Integrating Eq.(6) over energies we readily find $`eV_T`$ $$eV_T=\frac{L_1}{L}\frac{\delta T}{T_o}_0^{\mathrm{}}𝑑ϵ\text{ }g_{z+}ϵ\mathrm{cosh}^2(\beta _oϵ)/_0^{\mathrm{}}𝑑ϵ\text{ }g_{1+}\mathrm{cosh}^2(\beta _oϵ)$$ (7) The potential $`V_T`$ equals approximately the voltage difference between the N reservoirs and superconducting loop. It is worth noting that $`V_T,`$ determined by Eq.(7), does not depend on the small parameter $`r`$ because both functions $`g_{z+}`$ and $`g_{1+}`$ are proportinal to $`r`$ (however one should have in mind that according to the condition (2) this parameter must not be too small). The integrand in Eq.(7) can be calculated if the function $`\widehat{F}^R`$ is known from an approximate or numerical solution of the Usadel equation. In the limit considered, of small $`r`$, the retarded (advanced) Green’s functions are readily found from the linearized Usadel equation. In this case we find $$g_{z+}=rRe(F_S^R\frac{\mathrm{sinh}^2\theta _2}{\theta \mathrm{sinh}(2\theta )})\mathrm{sin}\phi ;\text{ }g_{1+}=rIm\{F_S^R[\mathrm{sinh}(\theta _2+2\theta _1)+\mathrm{cos}\phi \mathrm{sinh}\theta _2]\frac{\mathrm{sinh}\theta _2}{\theta \mathrm{sinh}(2\theta )}\}$$ (8) where $`\theta =\theta _1+i\theta _2,\theta _{1(2)}=k_ϵL_{1(2)},k_ϵ=\sqrt{2iϵ/ϵ_L}`$,$`L_2=LL_1`$ ,$`F_S^R=\mathrm{\Delta }/\sqrt{(ϵ+i\mathrm{\Gamma })^2\mathrm{\Delta }^2}`$ is the retarded Green’s function in the superconductor. One can see that the ratio in Eq.(7) indeed does not depend on $`r.`$ Numerical analysis of the Usadel equation shows that at a characteristic energy $`ϵϵ_L`$ the difference between the linearized and exact numerical solutions of the Usadel equation is less than 10% even for $`r1.`$ It follows from Eqs.(7-8) that the voltage $`V_T`$ caused by the temperature gradient is zero when the phase difference between the superconductors is zero and oscillates with increasing $`\phi `$. One can easily estimates the order of magnitude of $`V_T`$. We find $$eV_T=\delta T(L_1/L)\mathrm{sin}\phi \{\begin{array}{c}(T/ϵ_L)C_1(\phi ),T<<ϵ_L\\ (ϵ_L/T)^{3/2}C_2(\phi ),T>>ϵ_L\end{array}$$ (9) Here $`C_{1,2}(\phi )`$ are periodic functions of the phase difference $`\phi `$ of order 1 and are not zero when $`\phi =0.`$ If we define the thermoemf $`V_T`$ as the voltage between the N reservoirs and the superconducting circuit, we can state that this thermoemf is much larger than the thermoemf between the N reservoirs in the absence of superconductors because $`V_T`$ does not contain the small parameter $`(T`$/$`ϵ_F)`$ as is the case in a normal system (see for example Ref.). In addition, $`V_T`$ oscillates with an applied magnetic field $`H`$ ($`\phi `$ is proportional to $`H`$) allowing one to detect small temperature gradients. We stress once again that the thermoemf analysed in this work arises not between the normal reservoirs, as takes place in case of the ordinary thermoelectric effect, but between the superconducting and normal circuits. In Fig.2 we plot the dependence $`V_N`$ on $`\phi `$ for various $`\beta ^1=(2T)`$. In Fig.3 the temperature dependence $`V_T`$ at $`\phi =\pi /2`$ is presented. We see that this dependence is non-monotonic with a maximum at $`Tϵ_L`$ (reentrant behaviour). We note that the influence of the proximity effect on the ordinary thermoelectric effect was studied theoretically in . We ignore this effect regarding the thermoelectric current $`\alpha T`$ as negligable, compared to the effect under consideration. In Ref. the thermoelectric voltage was measured in complicated S/N structures which differ from the simple structure considered by us. The reentrant behaviour of $`V_T`$ and its oscillations on $`\phi `$ were observed in this work. It is possible that the observed effects are related to those considered in this paper. It is worthwhile noting that the influence of the ordinary thermoelectric current $`\alpha T`$ on the Josephson effect was studied long ago . The physical explanation of the effect is the following. The temperature gradient creates a deviation of the distrubtion function $`\delta f_+=(\delta n+\delta p)`$ from the equilibrium form. On the other hand the superconductors do not affect this function because complete Andreev reflection conserves the total number of excess electrons and holes. The function $`\delta f_+`$ changes the condensate current flowing across the S/N interface. If $`\delta f_+`$ is created by an electrical current flowing between the N reserviors, it has the same sign at $`x=\pm L`$ and leads to a change in the Josephson current. In the case considered here the function $`\delta f_+`$ has different signs at these points and leads to a variation of the condesate current of the same sign at different S/N interfaces. Therefore, the potential $`V_T`$ arises in the N wire producing a subgap current $`rg_{1+}f_{}(\pm L_1)`$ which compensates the current $`\delta J_S`$. In summary, we have calculated the voltage $`V_T`$ between the superconducting and normal curcuits in a S/N mesoscopic structure where the normal reservoirs are maintained at different temperatures. This voltage arises due to a branch imbalance in the N film and oscillates with varying phase difference. Its magnitude does not contain the small parameter $`(T/ϵ_F)`$ which is present in normal systems and is of the order $`\delta T(L_1/L)/e`$. We are grateful to the Royal Society and to the EPSRC for their financial support. We are grateful to Yu.Gal’perin for his useful comments and to V.V.Pavlovskii for his assistance. FIG. Giant thermoemf in multiterminal superconductor/normal metal mesoscopic structures. Schematic view of the 4-terminal S/N/S structure under consideration. The electric potential of the superconductors is zero. The N reservoirs are disconnected from the external circuit. FIG. Giant thermoemf in multiterminal superconductor/normal metal mesoscopic structures. The dependence of the normalised thermoeletrical voltage $`\stackrel{~}{V_T}=eV_T/(\delta TL_1/L)`$ on the phase difference $`\phi `$ for various $`\beta =T/ϵ_L`$ (the parameters are $`\mathrm{\Delta }/ϵ_L=10,L_1/L=0.5,r=0.3`$). FIG. Giant thermoemf in multiterminal superconductor/normal metal mesoscopic structures. The temperature dependence of the normalised voltage $`\stackrel{~}{V_T}`$ at $`\phi =\pi /2`$, for the same parameters as in Fig. 2.
warning/0005/hep-ph0005206.html	ar5iv	text	# Contents ## 1. Introduction In this talk I will describe a paper published last year with Shelly Glashow; the work was done at the Korean Institute for Advanced Study, Seoul, in June 1999 and the topic is what is the simplest way to extend the minimal standard model, where neutrinos are by definition massless, to accommodate the compelling evidence especially from the SuperKamiokande experiment for non-zero neutrino mass. ## 2. The Situation of the Data The minimal standard model involves three chiral neutrino states, but it does not admit renormalizable interactions that can generate neutrino masses. Nevertheless, experimental evidence suggests that both solar and atmospheric neutrinos display flavor oscillations, and hence that neutrinos do have mass. Two very different neutrino squared-mass differences are required to fit the data: $$10^{11}\mathrm{eV}^2\mathrm{\Delta }_s10^5\mathrm{eV}^2\mathrm{and}\mathrm{\Delta }_a10^3\mathrm{eV}^2$$ (1) where the neutrino masses $`m_i`$ are ordered such that: $$\mathrm{\Delta }_s\|m_2^2m_1^2\|\mathrm{and}\mathrm{\Delta }_a\|m_3^2m_2^2\|\|m_3^2m_1^2\|$$ and the subscripts $`s`$ and $`a`$ pertain to solar and atmospheric oscillations respectively. The large uncertainty in $`\mathrm{\Delta }_s`$ reflects the several potential explanations of the observed solar neutrino flux: in terms of vacuum oscillations or large-angle or small-angle MSW solutions, but in every case the two independent squared-mass differences must be widely spaced with $$r\mathrm{\Delta }_s/\mathrm{\Delta }_a<10^2$$ In a three-family scenario, four neutrino mixing parameters suffice to describe neutrino oscillations, akin to the four Kobayashi-Maskawa parameters in the quark sector.<sup>1</sup><sup>1</sup>1Two additional convention-independent phases are measurable in principle, but they ordinarily do not affect neutrino oscillations Solar neutrinos may exhibit an energy-independent time-averaged suppression due to $`\mathrm{\Delta }_a`$, as well as energy-dependent oscillations depending on $`\mathrm{\Delta }_s/E`$. Atmospheric neutrinos may exhibit oscillations due to $`\mathrm{\Delta }_a`$, but they are almost entirely unaffected by $`\mathrm{\Delta }_s`$. It is convenient to define neutrino mixing angles as follows: $$\left(\begin{array}{c}\nu _e\\ \nu _\mu \\ \nu _\tau \end{array}\right)=X\left(\begin{array}{c}\nu _1\\ \nu _2\\ \nu _3\end{array}\right)$$ where $$X=\left(\begin{array}{ccc}c_2c_3& c_2s_3& s_2e^{i\delta }\\ +c_1s_3+s_1s_2c_3e^{i\delta }& c_1c_3s_1s_2s_3e^{i\delta }& s_1c_2\\ +s_1s_3c_1s_2c_3e^{i\delta }& s_1c_3c_1s_2s_3e^{i\delta }& +c_1c_2\end{array}\right)$$ with $`s_i`$ and $`c_i`$ standing for sines and cosines of $`\theta _i`$. For neutrino masses satisfying Eq.(1) the vacuum survival probability of solar neutrinos is $$P(\nu _e\nu _e)\|_s1\frac{\mathrm{sin}^22\theta _2}{2}\mathrm{cos}^4\theta _2\mathrm{sin}^22\theta _3sin^2(\mathrm{\Delta }_sR_s/4E)$$ (2) whereas the transition probabilities of atmospheric neutrinos are: $`P(\nu _\mu \nu _\tau )\|_a`$ $``$ $`\mathrm{sin}^22\theta _1\mathrm{cos}^4\theta _2\mathrm{sin}^2(\mathrm{\Delta }_aR_a/4E)`$ $`P(\nu _e\nu _\mu )\|_a`$ $``$ $`\mathrm{sin}^22\theta _2\mathrm{sin}^2\theta _1\mathrm{sin}^2(\mathrm{\Delta }_aR_a/4E)`$ $`P(\nu _e\nu _\tau )\|_a`$ $``$ $`\mathrm{sin}^22\theta _2\mathrm{cos}^2\theta _1\mathrm{sin}^2(\mathrm{\Delta }_aR_a/4E)`$ None of these probabilities depend on $`\delta `$, the measure of CP violation. ## 3. Extending the Standard Model Let us turn to the origin of neutrino masses. Among the many renormalizable and gauge-invariant extensions of the standard model that can do the trick are: * The introduction of a complex triplet of mesons ($`T^{++},T^+,T^0`$) coupled bilinearly to pairs of lepton doublets. They must also couple bilinearly to the Higgs doublet(s) so as to avoid spontaneous $`BL`$ violation and the appearance of a massless and experimentally excluded majoron. This mechanism can generate an arbitrary complex symmetric Majorana mass matrix for neutrinos. * The introduction of singlet counterparts to the neutrinos with very large Majorana masses. The interplay between these mass terms and those generated by the Higgs boson—the so-called see-saw mechanism—yields an arbitrary but naturally small Majorana neutrino mass matrix. * The introduction of a charged singlet meson $`f^+`$ coupled antisymmetrically to pairs of lepton doublets, and a doubly-charged singlet meson $`g^{++}`$ coupled bilinearly both to pairs of lepton singlets and to pairs of $`f`$-mesons. An arbitrary Majorana neutrino mass matrix is generated in two loops. * The introduction of a charged singlet meson $`f^+`$ coupled antisymmetrically to pairs of lepton doublets and (also antisymmetrically) to a pair of Higgs doublets. This simple mechanism was first proposed by Tony Zee and results (at one loop) in a Majorana mass matrix in the flavor basis ($`e,\mu ,\tau `$) of a special form: $$=\left(\begin{array}{ccc}0& m_{e\mu }& m_{e\tau }\\ m_{e\mu }& 0& m_{\mu \tau }\\ m_{e\tau }& m_{\mu \tau }& 0\end{array}\right)$$ (4) In extending the standard model one crucial requirement is anomaly cancellation (e.g. ) and this is satisfied by all the above four scenarios. We focus here on the last scenario. In particular, we adopt the Zee ansatz for $``$ without committing ourselves to the Zee mechanism for its origin. Related discussions of Eq.(4) appear elsewhere The present work is essentially a continuation of . Because the diagonal entries of $``$ are zero, the amplitude for no-neutrino double beta decay vanishes at lowest order and this process cannot proceed at an observable rate. Furthermore, the parameters $`m_{e\mu },m_{e\tau }`$ and $`m_{\mu \tau }`$ may be taken real and non-negative without loss of generality, whence $``$ becomes real as well as traceless and symmetric. With this convention, the analog to the Kobayashi-Maskawa matrix becomes orthogonal: $``$ is explicitly CP invariant and $`\delta =0`$. But as we have noted, it is well known that the mere existence of a squared-mass hierarchy virtually precludes any detectable manifestation of CP violation in the neutrino sector. The sum of the neutrino masses (the eigenvalues of $``$) vanishes: $$m_1+m_2+m_3=0$$ (5) An important result emerges when the squared-mass hierarchy Eqs.(1) is taken into account along with Eq.(5). In the limit $`r0`$, two of the squared masses must be equal. ## 4. The Logical Possibilities There are two possibilities. In case A, we have $`m_1+m_2=0`$ and $`m_3=0`$. This case arises iff at least one of the three parameters in $``$ vanishes. In case B, we have $`m_1=m_2`$ and $`m_3=2m_2>0`$. This case arises iff the three parameters in $``$ are equal to one another. Of course, $`r`$ is small but it does not vanish: in neither case can the above relations among neutrino masses be strictly satisfied. But they must be nearly satisfied. Consequently we may deduce certain approximate but useful restrictions on the permissable values of the neutrino mixing angles $`\theta _i`$. Prior to examining these restrictions, we note that Eqs.(5) and (1) exclude the possibility that the three neutrinos are nearly degenerate in mass. If the Zee ansatz is even approximately realized in nature, no neutrino mass can exceed a small fraction of an elecron volt in magnitude and neutrinos are unlikely to contribute significantly to the dark matter of the universe. We first consider case A. The relation $`m_1+m_2=0`$ may be obtained in three ways depending on which parameter in $``$ is set to zero. If $`m_{e\mu }=0`$, the quantum number $`L_\tau L_eL_\mu `$ is conserved. It follows that $`\mathrm{cos}\theta _1=0`$ and $`\theta _3=\pi /4`$. We see from Eq.(LABEL:atmos) that atmospheric $`\nu _\mu `$’s oscillate exclusively into $`\nu _e`$’s and vice versa. This subcase (or any nearby assignment of mixing angles) is strongly disfavored by SuperKamiokande data. Alternatively we may set $`m_{e\tau }=0`$ to obtain conservation of $`L_\mu L_\tau L_e`$. For this subcase, we obtain $`\mathrm{sin}\theta _1=0`$ and $`\theta _3=\pi /4`$. We see from Eq.(LABEL:atmos) that atmospheric $`\nu _\mu `$’s do not oscillate at all. This subcase (or any nearby assignment of mixing angles) is also strongly disfavored by SuperKamiokande data. The last and best () version of Case A has $`m_{\mu \tau }=0`$ and leads to conservation of $`L_eL_\mu L_\tau `$. For this subcase, we obtain $`\mathrm{sin}\theta _2=0`$ and $`\theta _3=\pi /4`$. We see from Eq.(2) that solar neutrino oscillations are maximal: $$P(\nu _e\nu _e)\|_s=1\mathrm{sin}^2(\mathrm{\Delta }_sR_s/4E)$$ (6) Moreover, we see from Eq.(LABEL:atmos) that atmospheric $`\nu _\mu `$’s oscillate exclusively into $`\nu _\tau `$’s with the unconstrained mixing angle $`\theta _1`$: $`P(\nu _\mu \nu _\tau )\|_a`$ $`=`$ $`\mathrm{sin}^22\theta _1\mathrm{sin}^2(\mathrm{\Delta }_aR_a/4E)`$ $`P(\nu _\mu \nu _e)\|_a=0`$ $`P(\nu _e\nu _\tau )\|_a=0`$ This implementation of the Zee ansatz is compatible with experiment: It predicts maximal solar oscillations (without an energy-independent term) and it is consistent with the just-so vacuum oscillation hypothesis. However, it is evidently not compatible with the small-angle MSW explanation of solar neutrino data. Neither is it compatible with the large-angle MSW solution, because it predicts virtually maximal solar neutrino oscillations. If $`m_{\mu \tau }`$ is permitted to depart slightly from zero so as to generate a small finite value of $`r\mathrm{\Delta }_s/\mathrm{\Delta }_a`$, the coefficient of the oscillatory term in Eq.(6) will depart from unity by a term of order $`r^210^4`$. The resulting solar-neutrino oscillations remain nearly maximal: they are energy independent and experimentally disfavored unless $`\mathrm{\Delta }_s`$ lies within the just-so domain. We furthermore note that the $`m_{\mu \tau }0`$ version of the Zee ansatz admits atmospheric neutrino oscillations of the type $`\nu _\mu \nu _\tau `$ with any value of the mixing angle $`\theta _1`$. At the same time, it precludes all oscillations involving atmospheric $`\nu _e`$. These results are quite in accord with SuperKamiokande data. Others have noted that the relation $`m_1^2=m_2^2`$ is preserved by radiative corrections for case A. This is not necessarily true for case B, where the relations satisfied by the neutrino masses, $`m_3=2m_2=2m_1`$, are not a consequence of a symmetry principle. In any event, we argue that case B cannot fit the data. Ref. shows that case B leads to the relation: $$\mathrm{tan}^2\theta _2=1/2$$ (8) Along with Eqs.(2), this implies: $$P(\nu _e\nu _e)\|_a=1(8/9)\mathrm{sin}^2(\mathrm{\Delta }_aR_a/4E)$$ (9) That is, we must have large (almost maximal) oscillations of atmospheric $`\nu _e`$. This result is strongly disfavored by SuperKamiokande data, so that case B can be rejected without further ado. Our conclusion is simple. We find that one (and only one!) of the realizations of the Zee ansatz incorporating a squared-mass hierarchy is compatible with both solar and atmospheric neutrino data. It corresponds to the assignments $`m_{e\mu }m\mathrm{cos}\theta _1`$, $`m_{e\tau }m\mathrm{sin}\theta _1`$ and $`m_{\mu \tau }m`$, with $`\theta _20`$ and $`\theta _3\pi /4`$. Near this domain, atmospheric electron neutrinos oscillate negligibly, while atmospheric muon neutrinos oscillate into tau neutrinos with the arbitrary mixing angle $`\theta _1`$. Solar neutrino oscillations are very nearly maximal. They can be described by vacuum oscillations, but not by MSW oscillations. It is straightforward to implement the Zee model so as to conserve $`L_eL_\mu L_\tau `$ exactly so as to obtain $`m_{\mu \tau }=0`$ to all orders. Of course, this must not be done! Some unspecified new physics (beyond the introduction of Zee’s $`f^+`$ meson) is required to lift the degeneracy between $`m_1^2`$ and $`m_2^2`$ so as to yield an extreme hierarchy of neutrino squared-mass differences, with $`\mathrm{\Delta }_s10^8\mathrm{\Delta }_a`$. ## 5. Acknowledgement This work was supported in part by the U.S. Department of Energy under grant number DE-FG02-97ER-41036.
warning/0005/math-ph0005018.html	ar5iv	text	# Reflection symmetries of almost periodic functions ## 1. Introduction In this article we are interested in global reflection symmetries in almost periodic functions. Given a bounded function $`f:`$ we denote the set of translates of $`f`$ by $`U_0`$. The function $`f`$ is said to be almost periodic if $`U_0`$ is precompact in $`\mathrm{}^{\mathrm{}}()`$. The closure of this set is called the hull of $`f`$ and will be denoted by $`U`$. Let us suppose that $`f`$ is almost periodic. As we shall see in Section 2, $`U`$ admits an abelian group structure. We use $`\mu `$ to denote the Haar measure on this compact abelian group and normalize it by $`\mu (U)=1`$. We are be interested in the set of those elements in $`U`$ which are almost symmetric about the origin. To state this more formally, we introduce the reflection map $`R:\mathrm{}^{\mathrm{}}()\mathrm{}^{\mathrm{}}()`$, $`[Rg](n)=g(n)`$ and the standard metric on $`\mathrm{}^{\mathrm{}}()`$, $`d(g,g^{})=gg^{}_{\mathrm{}^{\mathrm{}}()}`$. Then we define $$U_\mathrm{r}(\epsilon )=\{gU:d(Rg,g)\epsilon \}$$ for each $`\epsilon 0`$. There are two further notions which we wish to introduce: The translation operator maps $`U`$ onto itself by $`[Tg](n)=g(n+1)`$. A function is said to be limit periodic if it is the uniform limit of a sequence of periodic functions. With this definition we may state ###### Theorem 1. Suppose $`f`$ is almost periodic but not limit periodic. Then there exists a constant $`C`$ such that for all $`\epsilon >0`$, $`\mu (U_\mathrm{r}(\epsilon ))C\epsilon `$. Our motivation for Theorem 1 is to understand the applicability of a criterion due to Jitomirskaya and Simon . This is a criterion for certain Schrödinger operators to have no eigenvalues. To each $`gU`$ one may associate a discrete Schrödinger operator $$[H_gu](n)=u(n+1)+u(n1)+g(n)u(n)\text{acting in }\mathrm{}^2()\text{.}$$ It is well known that the spectral properties of these operators are $`\mu `$-almost surely independent of $`g`$ (in fact, the spectrum and the absolutely continuous spectrum are completely independent of $`g`$, see Last-Simon ). The most prominent example is the almost Mathieu operator family given by $$f(n)=\lambda \mathrm{cos}(2\pi \alpha n),$$ where $`\lambda >0`$ and $`\alpha (0,1)`$ is irrational (see for reviews). Note that $`f`$ is even and hence $`U_\mathrm{r}(0)`$ is non-empty. It was conjectured by Aubry and Andre that for $`\lambda >2`$, the operators should exhibit pure point spectrum. However, based on Gordon , Avron and Simon proved that for Liouville $`\alpha `$, every $`H_g`$ has empty point spectrum. The conjecture was then modified to assert pure point spectrum for Diophantine $`\alpha `$. However, Jitomirskaya and Simon proved that if $`f`$ is even, $`\sigma _{\mathrm{pp}}(H_g)=\mathrm{}`$ for a dense $`G_\delta `$ of elements $`gU`$. It is the criterion they developed to prove this that we wish to discuss. Given $`f`$ they prescribe a number $`B>0`$ and introduce sets $$U^{(n)}=\{gU:d(RT^{2n}g,g)\mathrm{exp}(Bn)\}$$ and $$U_{\mathrm{rs}}=\underset{n\mathrm{}}{lim\; sup}U^{(n)}.$$ Jitomirskaya and Simon then prove that (1) $$U_{\mathrm{rs}}\{gU:\sigma _{\mathrm{pp}}(h_g)=\mathrm{}\}$$ and, moreover, that if $`f`$ is even, then $`U_{\mathrm{rs}}`$ is a dense $`G_\delta `$ subset of $`U`$. In view of (1) it is therefore natural to study $`\mu (U_{\mathrm{rs}})`$. Turning back to the almost Mathieu family, the strongest localization result one can hope for is pure point spectrum (in fact with exponentially decaying eigenfunctions) for every $`\lambda >2`$, Lebesgue-almost every $`\alpha `$ and $`\mu `$-almost every $`gU`$. Indeed, this modified form of the original Aubry-Andre conjecture has been recently established by Jitomirskaya (we remark that using duality, one can infer consequences for the case $`\lambda <2`$ ). By (1) this in turn implies $`\mu (U_{\mathrm{rs}})=0`$ in the almost Mathieu case, for $`\lambda >2`$ and almost every $`\alpha `$. Theorem 2 below, which follows quickly from Theorem 1, establishes this directly for all non-limit periodic hulls. ###### Theorem 2. Suppose $`f`$ is almost periodic but not limit periodic. Then, $`\mu (U_{\mathrm{rs}})=0`$. This proves that, although the Jitomirskaya-Simon criterion may be used to prove absence of eigenvalues for generic $`H_g`$, the set of applicability will always have zero Haar measure. To demonstrate that in general one cannot expect this result for limit periodic $`f`$, we consider the example (2) $$f(n)=\underset{k=1}{\overset{\mathrm{}}{}}a_k\mathrm{cos}\left(\frac{2\pi }{k}n\right)$$ with rapidly decaying $`a_k`$ and show that $`U_{\mathrm{rs}}`$ may consist of the entire hull. ###### Theorem 3. Let $`f`$ be given by (2) and suppose that (3) $$\underset{m\mathrm{}}{lim\; inf}\mathrm{\hspace{0.17em}\hspace{0.17em}4}\mathrm{exp}(m!B)\underset{k=m+1}{\overset{\mathrm{}}{}}\|a_k\|<1.$$ Then we have $`U_{\mathrm{rs}}=U`$. The organization is as follows. In Section 2 we recall some basic notations and observations for almost periodic hulls $`U`$. Section 3 provides the proofs of Theorems 1 and 2. Finally, Section 4 contains a proof of Theorem 3. ## 2. Preliminaries Much of the material presented in this section is an elementary consequence of the theories of almost periodic functions and of compact abelian groups—both core topics in harmonic analysis. The appendix of provides an extremely concise development of the relation of these two topics in the case of almost periodic functions on $``$. Background material is available in Katznelson and in a more abstract setting in Hewitt-Ross . Given a bounded function $`f:`$, we denote the translates of $`f`$ by $`f_m(n)=T^mf(n)=f(n+m)`$ and let $`U_0=\{f_m:m\}`$. Recall from the introduction that the function $`f`$ is called almost periodic if $`U_0`$ is precompact in $`\mathrm{}^{\mathrm{}}()`$ and that we denote the closure of $`U_0`$ by $`U`$. The set $`U_0`$ admits a group structure, $`[f_mf_m^{}](n)=f_{m+m^{}}(n)`$, which can then be extended to $`U`$. In particular, $`[f_mg](n)=g(n+m)`$ for all $`gU`$. Henceforth we shall assume that $`f`$ is almost periodic so that $`U`$ is a compact abelian group and take $`\mu `$ to be the normalized Haar measure on $`U`$; $`\mu (U)=1`$. Let $`\widehat{U}`$ denote the dual group, that is, the group of continuous homomorphisms $`𝝃:US^1=\{z:\|z\|=1\}`$. As $`U_0`$ is a dense subset of $`U`$, each character $`𝝃\widehat{U}`$ is uniquely determined by its values on $`U_0`$, and hence by $`\xi =𝝃(f_1)S^1`$. Notice that $`𝝃(f_m)=\xi ^m`$. We shall use freely the identification of $`\widehat{U}`$ as a subset of $`S^1`$ in what follows. The set $`\widehat{U}`$ forms an orthonormal basis of $`L^2(U;d\mu )`$ and so one obtains Fourier expansions for functions on $`U`$. For example, let $`\delta _0:U`$ be defined by $`\delta _0(g)=g(0)`$. Then $$\delta _0(g)=\widehat{\delta _0}(𝝃)𝝃(g),$$ where the sum should be interpreted in the $`L^2(U;d\mu )`$ sense and $$\widehat{\delta _0}(𝝃)=_U\delta _0(g)\overline{𝝃(g)}𝑑\mu (g).$$ It is also possible to reconstruct each element of $`U`$ from its character values. To be precise, (4) $$g(n)=\delta _0(f_ng)=\widehat{\delta _0}(𝝃)\xi ^n𝝃(g),$$ where the sum may be made to converge uniformly by employing the summability method of \[10, §5.18\]. One particular case of this is $`f(n)=\widehat{\delta _0}(𝝃)\xi ^n`$ which shows that $`\widehat{\delta _0}(𝝃)`$ are the Fourier coefficients of $`f`$, more commonly written as $`\widehat{f}(\{𝝃\})`$ (cf. ). The set of $`𝝃\widehat{U}`$ with $`\widehat{f}(\{𝝃\})0`$ is called the spectrum of $`f`$ and is denoted by $`\sigma (f)`$. Since (4) shows that each $`gU`$ is determined by merely $`𝝃(g)`$, $`𝝃\sigma (f)`$, it should come as no surprise that $`\widehat{U}`$ is generated by $`\sigma (f)`$; $`\widehat{U}=\sigma (f)`$. Recall that the function $`f`$ is limit periodic if there exists a sequence of periodic functions $`g^{(n)}:`$ which converge to $`f`$ in $`\mathrm{}^{\mathrm{}}()`$. The spectral condition is given by the following proposition (cf. Theorem A.1.2 of ). ###### Proposition 2.1. $`f`$ is limit periodic $``$ $`\sigma (f)\mathrm{exp}\{2\pi i\}`$ $``$ $`\widehat{U}\mathrm{exp}\{2\pi i\}`$. Proof. If $`f`$ is limit periodic, then for each $`\eta \mathrm{exp}\{2\pi i\}`$, we have (5) $$\underset{N\mathrm{}}{lim}\frac{1}{2N+1}\underset{n=N}{\overset{N}{}}\eta ^nf(n)=0.$$ This is easily shown for periodic $`f`$; it then follows for limit periodic $`f`$ by uniform approximation. It is also easy to show that $$\underset{N\mathrm{}}{lim}\frac{1}{2N+1}\underset{n=N}{\overset{N}{}}\eta ^n\xi ^n=\{\begin{array}{cc}1,\hfill & \eta =\xi \hfill \\ 0,\hfill & \eta \xi \hfill \end{array}.$$ Applying this and (5) to (4), we see that if $`f`$ is limit periodic and $`\eta \sigma (f)`$, then $`\eta \mathrm{exp}\{2\pi i\}`$. Conversely, if $`\sigma (f)\mathrm{exp}\{2\pi i\}`$, then applying summability to (4) shows that $`f`$ is a uniform limit of periodic functions. The other equivalence is a direct consequence of $`\widehat{U}=\sigma (f)`$. $`\mathrm{}`$ ## 3. The measure of $`U_\mathrm{r}(\epsilon )`$ and the applicability of the Jitomirskaya-Simon criterion In this section we prove Theorems 1 and 2. Proof of Theorem 1. Given $`gU_\mathrm{r}(\epsilon )`$, define $`\mathrm{\Delta }_g:U`$ by $$\mathrm{\Delta }_g(h)=\delta _0(gh)\delta _0(gh^1).$$ As $`gU_\mathrm{r}(\epsilon )`$, we have $`\|\mathrm{\Delta }_g(h)\|\epsilon `$ for every $`hU_0`$ (recall that $`[gf_m](n)=g(n+m)`$). Since $`U_0`$ is dense in $`U`$ and $`\mathrm{\Delta }_g`$ is continuous, this gives $`\|\mathrm{\Delta }_g(h)\|\epsilon `$ for every $`hU`$ and hence $`\mathrm{\Delta }_g_{L^2(U;d\mu )}^2\epsilon ^2`$. Consider the Fourier expansion of $`\mathrm{\Delta }_g`$. One finds $$\widehat{\mathrm{\Delta }_g}(𝝃)=\widehat{\delta _0}(𝝃)\left[𝝃(g)\overline{𝝃(g)}\right]$$ and thus by Parseval, (6) $$\|\widehat{\delta _0}(𝝃)\|^2\|𝝃(g)\overline{𝝃(g)}\|^2=\mathrm{\Delta }_g_{L^2(U;d\mu )}^2\epsilon ^2.$$ Each $`𝝃\widehat{U}`$ is a continuous homomorphism from $`U`$ to $`S^1`$ and so each $`𝝃(U)`$ is a compact subgroup of $`S^1`$. Of course, the characters $`𝝃`$ for which $`𝝃(U)`$ is discrete are precisely those with $`\xi \mathrm{exp}\{2\pi i\}`$. Since $`f`$ is assumed not to be limit periodic, by Proposition 2.1 there exists $`\eta \sigma (f)`$ which is not an element of $`\mathrm{exp}\{2\pi i\}`$. Consequently, $`𝜼(U)=S^1`$ and (7) $$\mu \left(\{gU:\|𝜼(g)\overline{𝜼(g)}\|\alpha \}\right)=\left\|\{\theta [0,1):4\mathrm{sin}^2(2\pi \theta )\alpha ^2\}\right\|\frac{\alpha }{2},$$ where $`\|\|`$ denotes Lebesgue measure. Also, since $`\eta \sigma (f)`$, $`\widehat{\delta _0}(𝜼)=\widehat{f}(\{𝜼\})0`$. Thus we may combine (6) with (7) to obtain $$\mu (U_\mathrm{r}(\epsilon ))\frac{\epsilon }{2\left\|\widehat{f}(\{𝜼\})\right\|}.$$ This proves the theorem. $`\mathrm{}`$ Proof of Theorem 2. It is easy to check that the maps $`R,T`$ defined in Section 1 obey $`RT=T^1R`$. This and the fact that $`\mu `$ is $`T`$-invariant yield $$\mu \left(\{gU:d(RT^{2n}g,g)\mathrm{exp}(Bn)\}\right)=\mu \left(U_\mathrm{r}(e^{Bn})\right)Ce^{Bn}.$$ The assertion now follows from Borel-Cantelli. $`\mathrm{}`$ ## 4. A counterexample In this section we consider limit periodic hulls $`U`$ generated by functions of the form (2) and prove that under assumption (3), every element of the hull has reflection symmetries as required in the definition of $`U_{\mathrm{rs}}`$. Proof of Theorem 3. Let $`f`$ be given by (2) and define $`f^{(m)}(n)=_{k=1}^ma_k\mathrm{cos}(\frac{2\pi }{k}n)`$. Obviously, $`f^{(m)}`$ is even and periodic with period bounded by $`(m!)`$. Thus $`f^{(m)}`$ has infinitely many centers of global reflection symmetry and two consecutive centers are separated by a distance of at most $`(m!)`$. Consider now an arbitrary element $`g`$ in the hull $`U`$ of $`f`$. By assumption there exists a sequence $`m_j\mathrm{}`$ such that $`4\mathrm{exp}(m_j!B)\epsilon _j<1`$ for every $`j`$, where $`\epsilon _j=_{k=m_j+1}^{\mathrm{}}\|a_k\|`$. Now choose some translate $`T^{l_j}f`$ of $`f`$ which is $`\epsilon _j`$-close to $`g`$. Then, $`T^{l_j}f^{(m_j)}`$ is $`2\epsilon _j`$-close to $`g`$ and has at least one center of global symmetry in $`\{1,\mathrm{},m_j!\}`$. Thus there exists an $`r_j\{1,\mathrm{},m_j!\}`$ such that $`d(RT^{2r_j}g,g)\mathrm{exp}(m_j!B)\mathrm{exp}(r_j!B)`$. $`\mathrm{}`$ Acknowledgments. We would like to thank T. Wolff for useful discussions. D. D. was supported by the German Academic Exchange Service through Hochschulsonderprogramm III (Postdoktoranden). R. K. was supported by a Sloan Doctoral Dissertation fellowship.
warning/0005/hep-th0005076.html	ar5iv	text	# 1 Introduction and Summary ## 1 Introduction and Summary Many beautiful exact results on supersymmetric $`4D`$ gauge theories have been obtained recently, following Seiberg and Witten’s breakthough on $`N=2`$ supersymmetric theories -. One exciting related development is Seiberg’s $`N=1`$ non-Abelian duality, present in many cases, between a pair of theories with different number of color, flavor, and matter contents, which describe the same low-energy physics -. Another is the discovery of universal classes of conformally invariant theories (CFT) for different values of color, flavor, and in some cases, for appropriately tuned values of the parameters of the theory -. Still another concerns the microscopic mechanism of confinement (e.g., monopole condensation) and dynamical symmetry breaking, and study of other phases such as the oblique confinement, upon addition of a perturbation breaking the $`N=2`$ supersymmetry to $`N=1`$ and/or to $`N=0`$ ,-. In fact, an interesting phenomenon has been observed in the case of $`SU(2)`$ gauge theories with various flavors and with adjoint mass perturbation : confinement is caused by condensation of magnetic monopoles carrying nontrivial flavor quantum numbers. More explicitly, for $`n_f=2`$, monopoles in the $`(\underset{¯}{2},\underset{¯}{1})+(\underset{¯}{1},\underset{¯}{2})`$ (spinor) representation of the flavor $`[SU(2)\times SU(2)]/Z_2=SO(4)`$ group is found to condense upon $`N=1`$ perturbation $`\mu \text{Tr}\mathrm{\Phi }^2`$: the flavor symmetry is necessarily broken to $`U(2)`$. For $`n_f=3`$, monopoles in the $`\underset{¯}{4}`$ (spinor) representation of the flavor $`SO(6)`$ group condense with $`\mu 0`$ and the flavor symmetry is broken to $`U(3)`$. In such systems spontaneous global symmetry breaking is caused by the same dyamical mechanism responsible for confinement. In one of the vacua with $`n_f=3`$, though, the condensing entity carries magnetic number twice the basic unit but is flavor neutral. In this case \- which can be interpreted as oblique confinement à la ’t Hooft \- confinement is not accompanied by flavor symmetry breaking. For $`n_f=1,4`$, there is no dynamical flavor symmetry breaking. These results naturally lead to a conjecture that the condensation of magnetic monopoles with non-trivial flavor transformation property might in a general class of systems explain the confinement à la ‘t Hooft, Nambu, Mandelstam and the flavor symmetry breaking, simultaneously <sup>1</sup><sup>1</sup>1Such a possibility has been critically discussed recently in QCD . . However, a simple thought reveals a problem with this picture. For instance, the monopoles in $`USp(2n_c)`$ theories transform under the spinor representation of $`SO(2n_f)`$ flavor symmetry, and their effective low-energy Lagrangian coupled to the magnetic $`U(1)`$ gauge group would have a large accidental $`SU(2^{n_f1})`$ flavor symmetry: their condensation would lead to far too many Nambu–Goldstone multiplets. The case of $`SU(2)`$ gauge theories was special because the flavor symmetries of the monopole action precisely coincide with the symmetry of the microscopic theories, somewhat accidentally, due to the small number of flavors. It is not at all obvious how such a paradox is avoided in higher-rank theories. Argyres, Plesser and Seiberg studied higher-rank $`SU(n_c)`$ theories with $`n_f2n_c1`$ (asymptotically free) in detail. They showed how the non-renormalization theorem of the hyperKähler metric on the Higgs branch could be used to show the persistence of unbroken non-abelian gauge group at the “roots” of the Higgs branches (non-baryonic and baryonic branches) where they intersect the Coulomb branch. Some isolated points on the non-baryonic roots with $`SU(r)`$ ($`r[n_f/2]`$) gauge group as well as the baryonic root (single point) with $`SU(\stackrel{~}{n}_c)=SU(n_fn_c)`$ gauge group were found to survive the $`\mu 0`$ perturbation. Their main focus, however, was an attempt to “derive” Seiberg’s duality between $`SU(n_c)`$ and $`SU(\stackrel{~}{n}_c)`$ gauge theories relying on the so-called baryonic root. Their “derivation,” however, was incomplete as it did not produce all components of the “meson” superfield. The effective low-energy theory was perturbed by a relevant operator (the mass term for the mesons) and did not flow to the Seiberg’s $`N=1`$ magnetic theory correctly. <sup>2</sup><sup>2</sup>2 We thank P. Argyres and A. Buchel for discussions on this point. On the other hand, the issue of flavor symmetry breaking was not studied at any depth in . Their analysis also left a puzzle why there were “extra” theories at the non-baryonic roots which seemingly had nothing to do with Seiberg’s dual theories. Another paper by Argyres, Plesser and Shapere addressed similar questions, and left puzzles, in $`SO(n_c)`$ and $`USp(2n_c)`$ theories . We investigate here the microscopic mechanism of dynamical symmetry breaking, taking as theoretical laboratory the same class of theories studied by the above mentioned authors , namely, theories constructed from exactly solvable $`N=2`$ $`SU(n_c)`$ and $`USp(2n_c)`$ gauge theories with all possible numbers of flavor compatible with asymptotic freedom, by perturbing them with a small adjoint mass (reducing supersymmetry to $`N=1`$) and quark masses. The Lagrangian of the models has the structure $$=\frac{1}{8\pi }\text{Im}\tau _{cl}\left[d^4\theta \mathrm{\Phi }^{}e^V\mathrm{\Phi }+d^2\theta \frac{1}{2}WW\right]+^{(quarks)}+\mathrm{\Delta }+\mathrm{\Delta }^{},$$ (1.1) where $$\mathrm{\Delta }=d^2\theta \mu \text{Tr}\mathrm{\Phi }^2$$ (1.2) reduces the supersymmetry to $`N=1`$; $$^{(quarks)}=\underset{i}{}[d^4\theta \{Q_i^{}e^VQ_i+\stackrel{~}{Q}_{i}^{}{}_{}{}^{}e^{\stackrel{~}{V}}\stackrel{~}{Q}_i\}+d^2\theta \{\sqrt{2}\stackrel{~}{Q}_i\mathrm{\Phi }Q^i+m_i\stackrel{~}{Q}_iQ^i\}$$ (1.3) describes the $`n_f`$ flavors of hypermultiplets (“quarks”), and $$\tau _{cl}\frac{\theta _0}{\pi }+\frac{8\pi i}{g_0^2}$$ (1.4) is the bare $`\theta `$ parameter and coupling constant. The $`N=1`$ chiral and gauge superfields $`\mathrm{\Phi }=\varphi +\sqrt{2}\theta \psi +\mathrm{}`$, and $`W_\alpha =i\lambda +\frac{i}{2}(\sigma ^\mu \overline{\sigma }^\nu )_\alpha ^\beta F_{\mu \nu }\theta _\beta +\mathrm{}`$ are both in the adjoint representation of the gauge group, while the hypermultiplets are taken in the fundamental representation of the gauge group. We shall consider, besides the adjoint mass, small generic nonvanishing bare masses for the hypermultiplets (“quarks”). The advantage of doing so is that the only vacua retained are those in which the gauge coupling constant grows in the infrared. Another advantage is that all flat directions are eliminated in this way and one is left with a finite number of isolated vacua; keeping track of this number allows us to do highly nontrivial checks of our analyses at various stages. The most salient features of the result of our analysis will be as follows. The theories studied in different regimes, semiclassical, large and small adjoint and/or bare quark masse, give a mutually consistent picture as regards the number of the vacua and the dynamical properties in each of them. Various dynamical possibilities, found to be realized in the $`SU(n_c)`$ and $`USp(2n_c)`$ theories with a small finite adjoint mass and in the vanishing bare quark mass limit, are summarized in Tables 1 and 2. With small but generic bare quark masses, the order parameter of confining vacua is indeed the condensation of magnetic monopoles for every $`U(1)`$ factor à la ’t Hooft, in all cases considered. The massless limit, however, is non-trivial and exhibits a much richer range of interesting dynamical possibilities. In $`SU(n_c)`$ theories with $`n_f`$ flavors, the following diversity of dynamical scenarios are realized, according to the number of flavors and to the particular vacua considered. In the first group of vacua with finite meson or dual quark vacuum expectation values (VEVS), labelled by an integer $`r`$, $`r[n_f/2]`$, the system is in confinement phase. The nature of the actual carrier of the flavor quantum numbers depends on $`r`$. In vacua with $`r=0`$, magnetic monopoles are singlets of the global $`U(n_f)`$ group, hence no global symmetry breaking accompanies confinement. This is analogous to the oblique confinement ’t Hooft suggested for QCD around $`\theta =\pi `$. In vacua with $`r=1`$, the light particles are magnetic monopoles in the fundamental representation of $`U(n_f)`$ flavor group, and are charged under one of the color $`U(1)`$ factors. Their condensation leads to the confinement and flavor symmetry breaking, simultaneouly. In vacua labelled by $`r`$, $`2r<n_f/2`$ but $`rn_fn_c`$, the grouping of the associated singularities on the Coulomb branch, with multiplicity, $`{}_{n_f}{}^{}C_{r}^{}`$,<sup>3</sup><sup>3</sup>3In this paper, we use the traditional notation $`{}_{n}{}^{}C_{r}^{}`$ for the binomial coefficient most of time: another frequently used symbol is $`\left(\genfrac{}{}{0pt}{}{n}{r}\right)`$. at first sight suggests the condensation of monopoles in the rank-$`r`$ anti-symmetric tensor representation of the global $`SU(n_f)`$ group. Actually, this does not occur. The true low-energy degrees of freedom of these theories are (magnetic) quarks plus a number of singlet monopoles of an effective $`SU(r)\times U(1)^{n_cr}`$ gauge theory. Monopoles in higher representations of $`SU(n_f)`$ flavor group probably exist semi-classically as seen in a Jackiw–Rebbi type analysis . Such monopoles can be interpreted as “baryons” made of the magnetic quarks, which, interactions being infrared-free, however break up into magnetic quarks before they become massless at singularities on the Coulomb branch. It is the condensation of these magnetic quarks that induces confinement and flavor symmetry breaking, $`U(n_f)U(r)\times U(n_fr)`$, in these vacua. The system thus realizes the exact global symmetry of the theory in a Nambu-Goldstone mode, without having unusually many Nambu-Goldstone bosons. It is a novel mechanism for confinement and dynamical symmetry breaking. In the special cases with $`r=n_f/2`$, still another dynamical scenario takes place. In these cases, the interactions among the monopoles become so strong that the low-energy theory describing them is a nontrivial superconformal theory, with conformal invariance explicitly broken by the adjoint or quark masses. Although the symmetry breaking pattern is known ($`U(n_f)U(n_f/2)\times U(n_f/2)`$), the low-energy degrees of freedom involve relatively nonlocal fields and their interactions cannot be described in terms of a local action. Finally, in the group of vacua labelled by $`r=n_fn_c`$, the low-energy degrees of freedom are again magnetic quarks and a number of singlet monopoles of an effective infrared-free $`SU(n_fn_c)\times U(1)^{2n_cn_f}`$ gauge theory. There are two physically distinct sub-groups of vacua: one in which the magnetic quarks condense (i.e. confinement phase) with the unbroken symmetry $`U(n_fn_c)\times U(n_c)`$ is analogous to the ones found for generic $`r`$’s; the other is vacua in which magnetic-quark do not condense and remain as physically observable particles at long distances (free magnetic phase). The global $`U(n_f)`$ symmetry remains unbroken. This last phase is related to the one discovered by Seiberg in $`N=1`$ massless SQCD for the range of $`n_f`$, $`n_c+1<n_f<3n_c/2`$. Nevertherless, it should be emphasized that the $`m_i0`$ limit here is a smooth one and the symmetry properties of the vacua are independent of the way the limit is taken, while in SQCD without the adjoint chiral superfield the vacuum properties depend critically on the order in which the $`m_i`$’s approach zero, showing typically the phenomenon of the run-away vacua. All in all, we find the number $$𝒩_1=(2n_cn_f)\mathrm{\hspace{0.17em}2}^{n_f1}$$ (1.5) of $`N=1`$ vacua with finite flavor-carrying VEVS and $$𝒩_2=\underset{r=0}{\overset{n_fn_c1}{}}(n_fn_cr){}_{n_f}{}^{}C_{r}^{},$$ (1.6) of them with vanishing VEVS. The latter is present only for theories with the large number of flavors ($`n_fn_c+1`$). Their sum, $`𝒩=_{r=0}^{\text{ min}\{n_f,n_c1\}}{}_{n_f}{}^{}C_{r}^{}(n_cr)`$, correctly generalizes <sup>4</sup><sup>4</sup>4 It may be said that the fact that the generalization is given by these formulas and analogous ones Eq.(1.7), Eq.(1.8), and not simply, e.g., by $`n_f+n_c`$, reveals the richness of dynamical scenarios of these theories. the well-known number of the vacua in the $`SU(2)`$ gauge theory, $`𝒩_{SU(2)}=n_f+2.`$ In $`USp(2n_c)`$ theories, again, we find two groups of vacua, whose properties are shown in Table 2. The most interesting difference as compared to the $`SU(n_c)`$ theory is that here the entire first group of vacua correspond to SCFT. As the superconformal theory is a nontrivial one, one does not have a local effective Lagrangian description for those theories. Nonetheless, the symmetry breaking pattern can be deduced, from the analysis done at large $`\mu `$: $`SO(2n_f)`$ symmetry is always spontaneously to $`U(n_f)`$. It is most instructive to consider the equal but nonvanishing quark mass case, first. (See Table 3.) The flavor symmetry group of the underlying theory is now broken explicitly to $`U(n_f)`$. The first group of vacua split into various branches labelled by $`r`$, $`r=0,1,2,\mathrm{},[\frac{n_f1}{2}]`$, each of which is described by a local effective gauge theory of Argyres-Plesser-Seiberg , with gauge group $`SU(r)\times U(1)^{n_cr1}`$ and $`n_f`$ (dual) quarks in the fundamental representation of $`SU(r)`$. Indeed, the gauge invariant composite VEVS characterizing these theories differ by some powers of $`m`$, and the validity of each effective theory is limited by small fluctuations of order of $`m`$ around each vacuum. In the limit $`m0`$ these points in the quantum moduli space (QMS) collapse into one single point. Obviously, a smooth $`m_i0`$ limit is not possible. The location of this singularity can be obtained exactly in terms of Chebyshev polynomials. At the singularity there are mutually non-local dyons and hence the theory is at a non-trivial infrared fixed point. In the example of $`USp(4)`$ theory with $`n_f=4`$, we have explicitly verified this by determining the singularities and branch points at finite equal mass $`m`$ and then by studying the limit $`m0`$. These cases, together with the special $`r=n_f/2`$ nonbaryonic root for the $`SU(n_c)`$ theory, constitute another new mechanism for dynamical symmetry breaking: although the global symmetry breaking pattern deduced indirectly looks familiar enough, the low-energy degrees of freedom are relatively nonlocal dual quarks and dyons. It would be interesting to get a better understanding of this phenomenon. For large numbers of flavor, there are also vacua, just as in large $`n_f`$ $`SU(n_c)`$ theories, with no confinement and no dynamical flavor symmetry breaking. The low-energy particles are solitonlike magnetic quarks which weakly interact with dual (in general) non-abelian gauge fields: the system is in the free magnetic phase. The vacuum counting gives, for $`USp(2n_c)`$ theories, $$𝒩_1=(2n_c+2n_f)2^{n_f1}$$ (1.7) vacua with finite VEVS (first group of vacua), and $$𝒩_2=\underset{r=0}{\overset{n_fn_c2}{}}(n_fn_c1r){}_{n_f}{}^{}C_{r}^{}$$ (1.8) of them with vanishing VEVS (second group of vacua); the latter are present only for $`n_fn_c+2`$. The paper will be organized as follows. After discussing briefly the standard expectation on chiral symmety breaking in Sec. 2, we start in Sec. 3 a preparatory analysis, finding all isolated semi-classical vacua by minimizing the scalar potential and determining the vevs of the adjoint scalar and of the squark fields in each of them. This allows us to count the number of all possible vacua of the theory, after taking appropriate account of Witten’s index corresponding to the unbroken gauge group in each case. The first major step of our analysis is the analytic, first-principle determination of the global symmetry breaking pattern in each vacuum, done by studying the theories at large $`\mu `$ ($`\mathrm{\Lambda }`$), and $`m_i0`$ (Sec. 4). Such a determination is possible because in this case the effective superpotential can be read off from the bare Lagrangian by integrating out the heavy, adjoint fields and by adding to it the known exact instanton–induced superpotentials of the corresponding $`N=1`$ theories. By minimizing the scalar potential, we reproduce in all cases the correct number of the vacua and explicitly determine the pattern of global symmetry breaking in each them. By $`N=1`$ supersymmetry and holomorphic dependence of physics on $`\mu `$ the same symmetry breaking pattern is valid at any finite $`\mu `$. In Section 5 another related limit, $`\mu \mathrm{}`$, $`m_i`$ and $`\mathrm{\Lambda }_1\mu ^{\frac{n_c}{3n_cn_f}}\mathrm{\Lambda }^{\frac{2n_cn_f}{3n_cn_f}}`$ fixed, is studied and consistency with the known results in the standard $`N=1`$ theories without adjoint fields is checked. The properties of the quantum vacua at small $`\mu `$ and small $`m_i`$ are studied in detail, in Sections 6 \- 9. In Section 6 we show how all of the $`N=1`$ vacua arise from various classes of conformally invariant theories (CFT) upon perturbation in bare quark masses on Seiberg-Witten curves, reproducing the correct number of $`N=1`$ vacua found in the earlier analyses. In Section 7 we check and illustrate these results in the cases of rank-two gauge groups, $`SU(3)`$ and $`USp(4)`$ theories, by directly finding the associated singularities numerically. In Section 8 the effective-Lagrangian description of these $`N=1`$ theories is analysed, where we verify that the number of $`N=1`$ vacua and the symmetry breaking pattern in each of them indeed agree with what we found earlier in Sections 3-7. This leads us to a better, more microscopic understanding of the phenomena of confinement and dynamical global symmetry breaking, as summarized above. The actual perturbation theory in bare quark masses on Seiberg-Witten curves, whose outcome is quite central to the whole analysis but whose analysis per se is independent of the rest of the paper, is developed in the last section (Sec. 9). Several technical discussions are relegated to Appendices. Appendix A gives a proof of $`SO(2N)USp(2N)=U(N)`$; in Appendix B we discuss the Jackiw-Rebbi construction of flavor multiplet structure of semiclassical monopoles for $`SU(n_c)`$ and $`USp(2n_c)`$ gauge theories; we list explicit expressions for $`a_{Di}`$, $`a_i`$, $`\frac{a_{Di}}{u_j},`$ and $`\frac{a_i}{u_j}`$ in Appendix C; the proof of absence of the “nonbaryonic branch root” with $`r=\stackrel{~}{n}_c=n_fn_c`$ is given in Appendix D; the study of the monodromy around the seventeen singularities of $`SU(3)`$, $`n_f=4`$ theory is discussed in Appendix E . A shorter version of this work has appeared already . The case of $`SO(n_c)`$ theories, where some new subtleties are present, will be discussed in a separate article. ## 2 Expected Pattern of Chiral Symmetry Breaking In non-supersymmetric theories with fermions in the fundamental representation of the gauge group, global symmetry can be broken spontaneously if a fermion bilinear condensate $$\psi \psi \mathrm{\Lambda }^3$$ (2.1) forms. Apart from the requirement of gauge invariance there are no general rules which fermion pairs condense, although in QCD at large $`n_c`$ limit one can argue that $`SU_L(n_f)\times SU_R(n_f)\times U_V(1)`$ symmetry is broken to the diagonal $`SU_V(n_f)\times U_V(1)`$, which is the observed pattern of chiral symmetry breaking in Nature. Nonetheless certain general considerations can be made. For nonsupersymmetric $`SU(2)`$ theories, $`2n_f`$ fermions transform as the fundamental representation of the global $`SU(2n_f)`$ group: bilinear condensate ($`a,b=1,2`$ are color indices; $`i,j=1,2,\mathrm{}2n_f`$ are flavor indices), $`ϵ^{ab}\psi _a^i\psi _b^j,`$ is necessarily antisymmetric in $`(i,j)`$. If these condensates can be put by an $`SU(2n_f)`$ transformation into the “standard” form $$ϵ^{ab}\psi _a^i\psi _b^j=\text{const.}J^{ij},$$ (2.2) where $`J=i\sigma _2𝐈_{n_f\times n_f}`$, then such condensates would leave $`USp(2n_f)`$ invariant. Note that this is consistent with what was found in the $`N=2`$ $`SU(2)`$ QCD, broken to $`N=1`$ by a small adjoint mass term, in case of $`n_f=2`$ and in one of the vacua of $`n_f=3`$, where the symmetry of the vacuum was found to be $`U(n_f)`$ . The reason is that in the model of the global symmetry of the theory is reduced to $`SO(2n_f)`$ due to the Yukawa interaction, and the intersection of $`USp(2n_f)`$ and $`SO(2n_f)`$ is precisely $`U(n_f)`$ (see Appendix A). For another vacuum of the $`n_f=3`$ theory (where oblique confinement of ’t Hooft takes place), the monopole that condenses is a flavor singlet and chiral symmetry remains unbroken. In the standard QCD with $`SU(n_c)`$ gauge symmetry ($`n_c3`$), $`2n_f`$ fermions transform as $`(\underset{¯}{n}_f,\underset{¯}{1})+(\underset{¯}{1},\underset{¯}{n}_f^{})`$ of $`SU_L(n_f)\times SU_R(n_f)`$; condensates of the form $$\overline{\psi }_{R}^{}{}_{}{}^{i}\psi _{Lj}=v\delta _j^i,$$ (2.3) is believed to form, at least for small $`n_f`$, leaving the unbroken diagonal $`SU(n_f)`$ symmetry. Unfortunately, in the corresponding $`N=2`$ theories (with a small $`N=1`$ perturbation) the axial symmetry is explicitly broken at the tree level by the characteristic Yukawa interactions so that the global symmetry contains only the diagonal $`SU(n_f)`$, already at the tree level. Thus $`SU(n_c)`$ theories will be considered below mainly as a testing ground of our approach, in correctly identifying the quantum vacua which survive $`N=1`$ perturbation, matching the numbers of classical and quantum vacua, and in verifying in each such vacua the ’t Hooft–Mandelstam mechanism for confinement. In fact this study reveals new, unexpected ways confinement and dynamical symmetry breaking are realized in non-Abelian gauge theories. The cases of $`USp(2n_c)`$ theories are more promising. As noted above for $`SU(2)`$ gauge theories, in a nonsupersymmetric theories with $`2n_f`$ fermions in the fundamental representation, the global symmetry is $`SU(2n_f)`$. Bifermion condensates of the “standard” form $$\psi _i^a\psi _j^bJ_{ab}=vJ_{ij}$$ (2.4) would break it to $`USp(2n_f)`$. On the other hand, the corresponding $`N=2`$ models (1.1) have a smaller flavor group, $`SO(2n_f)SU(2n_f)`$ due to the Yukawa interactions. Nevertheless, there is a nontrivial overlap between its flavor group $`SO(2n_f)`$ and the $`USp(2n_f)`$ (expected invariance group for nonsupersymmetric model), which is $`U(n_f)`$. One then expects that the global chiral symmetry $`SO(2n_f)`$ is broken spontaneously to $`U(n_f)SO(2n_f)`$. We shall see below that these expectations are indeed met by quantum vacua of the $`USp(2n_c)`$ theories (in the large flavor cases, these take place in the first group of vacua, while we find also a secon group of vacua in which the chiral $`SO(2n_f)`$ symmetry remains unbroken). The proof that $`SO(2N)USp(2N)=U(N)`$ is given in Appendix A. ## 3 Semi-Classical Vacua In this section, we find all semi-classical vacua in $`N=2`$ $`SU(n_c)`$ and $`USp(2n_c)`$ theories with quark hypermultiplets in their fundamental representations, perturbed by the quark masses as well as that of the adjoint fields in the $`N=2`$ vector multiplet. The analysis is quantum mechanically valid at large $`\mu `$ and $`m_i`$. $`N=1`$ supersymmetry and holomorphy in $`\mu `$ and $`m_i`$ forbid any phase transition as one moves to smaller $`\|m_i\|`$ and $`\|\mu \|`$ hence the same number of $`N=1`$ vacua must be present in the different regimes to be considered in the subsequent sections. In particular, this analysis allows to determine the symmetry breaking pattern in the equal (and nonvanishing) mass case, in which the classical global symmetry is $`U(n_f)`$ in both $`SU(n_c)`$ and $`USp(2n_c)`$ theories. ### 3.1 Semi-Classical Vacua in $`SU(n_c)`$ In the limit $`m_i0`$ and $`\mu 0`$, the global symmetry of the model is $`U(n_f)\times Z_{2n_cn_f}\times SU_R(2).`$ The superpotential of our theory, with $`N=2`$ supersymmetry softly broken to $`N=1`$ by the adjoint mass, is given by: $$W=\mu \text{Tr}\mathrm{\Phi }^2+\sqrt{2}\stackrel{~}{Q}_i^a\mathrm{\Phi }_a^bQ_b^i+m_i\stackrel{~}{Q}_i^aQ_a^i,$$ (3.1) where $$\mathrm{\Phi }\lambda ^A\mathrm{\Phi }^A,(A=1,2,\mathrm{}N^21),$$ (3.2) $$\text{Tr}(\lambda ^A\lambda ^B)=\frac{1}{2}\delta ^{AB}.$$ (3.3) $`i=1,2,\mathrm{}n_f`$ is the flavor index; $`a,b=1,2,\mathrm{}n_c`$ are the color indices. Note that with our normalization of the $`SU(n_c)`$ generators the color Fierz relation reads $$\underset{A=1}{\overset{n_c^21}{}}(\lambda ^A)_c^d(\lambda ^A)_a^b=\frac{1}{2}\left[\delta _c^b\delta _a^d\frac{1}{n_c}\delta _c^d\delta _a^b\right].$$ (3.4) The vacuum equations are $$[\mathrm{\Phi },\mathrm{\Phi }^{}]=0;$$ (3.5) $$\nu \delta _a^b=Q_a^i(Q^{})_i^b(\stackrel{~}{Q}^{})_a^i\stackrel{~}{Q}_i^b;$$ (3.6) $$Q_a^i\stackrel{~}{Q}_i^b\frac{1}{n_c}\delta _a^b(Q_c^i\stackrel{~}{Q}_i^c)+\sqrt{2}\mu \mathrm{\Phi }_a^b=0;$$ (3.7) $$Q_a^im_i+\sqrt{2}\mathrm{\Phi }_a^bQ_b^i=0(\text{no sum over}i);$$ (3.8) $$m_i\stackrel{~}{Q}_i^a+\sqrt{2}\stackrel{~}{Q}_i^b\mathrm{\Phi }_a^b=0(\text{no sum over}i).$$ (3.9) where the quark masses have been taken diagonal by flavor rotations. Use first $`SU(n_c)`$ rotation to bring $`\mathrm{\Phi }`$ into diagonal form, $$\mathrm{\Phi }=\text{diag}(\varphi _1,\varphi _2,\mathrm{}\varphi _{n_c}),\varphi _a=0.$$ (3.10) Eq. (3.8) and Eq. (3.9) say that $`Q_a^i`$ and $`\stackrel{~}{Q}_i^b`$ are either nontrivial eigenvectors of the matrix $`\mathrm{\Phi }`$ with possible eigenvalues $`m_i`$, or null vectors. With $`\mathrm{\Phi }`$ put in the diagonal form, and with generic (so unequal) masses, the eigenvectors have simple forms, $$Q^i=(0,\mathrm{},d_i,0,\mathrm{}),$$ (3.11) each with only one nonvanishing component (similarly for $`\stackrel{~}{Q}_i^a`$). There can be at most $`n_f`$ nontrivial eigenvalues, chosen from $`m_1,m_2,\mathrm{},m_{n_f}`$; at the same time the form (3.10) allows for at most $`n_c1`$ nonzero independent elements of $`\mathrm{\Phi }`$. The solutions can thus be classified by the number of nontrivial eigenvectors, $`r=1,2,\mathrm{},\text{min}\{n_f,n_c1\}`$. There are $`\left(\begin{array}{c}n_f\\ r\end{array}\right)`$ solutions for a given $`r`$, according to which $`m_i`$’s appear as eigenvalues. The solution with eigenvalues $`m_1,m_2,\mathrm{},m_r`$ is <sup>5</sup><sup>5</sup>5 These results are slight generalization of the ones in to generic nonvanishing quark and adjoint masses. Note that the flat directions are completely eliminated.: $$Q_a^i=\left(\begin{array}{c}d_1\\ 0\\ \mathrm{}\\ \mathrm{}\\ 0\end{array}\right),\left(\begin{array}{c}0\\ d_2\\ 0\\ \mathrm{}\\ 0\end{array}\right),\left(\begin{array}{c}0\\ \mathrm{}\\ d_r\\ 0\\ \mathrm{}\end{array}\right);Q_a^i=0,i=r+1,\mathrm{},n_f.$$ (3.12) $$\stackrel{~}{Q}_i^a=\left(\begin{array}{c}\stackrel{~}{d}_1\\ 0\\ \mathrm{}\\ \mathrm{}\\ 0\end{array}\right),\left(\begin{array}{c}0\\ \stackrel{~}{d}_2\\ 0\\ \mathrm{}\\ 0\end{array}\right),\left(\begin{array}{c}0\\ \mathrm{}\\ \stackrel{~}{d}_r\\ 0\\ \mathrm{}\end{array}\right);\stackrel{~}{Q}_i^a=0,i=r+1,\mathrm{},n_f,$$ (3.13) where $$r=0,1,\mathrm{},\text{min}\{n_f,n_c1\},$$ (3.14) $$\mathrm{\Phi }=\frac{1}{\sqrt{2}}\text{diag}(m_1,m_2,\mathrm{},m_r,c,\mathrm{},c);c=\frac{1}{n_cr}\underset{k=1}{\overset{r}{}}m_k.$$ (3.15) $`d_i`$’s can be chosen real (by residual $`SU(n_c)`$); $`\stackrel{~}{d}_i`$’s are in general complex. They are given by Eqs. (3.21) and (3.18) below. Eq. (3.5) is obviously satisfied. The nondiagonal ($`ab)`$ of Eq. (3.6) is also obvious. The first $`r`$ diagonal ($`a=b)`$ equations are: $$\nu =d_i^2\|\stackrel{~}{d}_i\|^2;(i=1,2,\mathrm{}r);$$ (3.16) the others give $$\nu =0,$$ (3.17) hence $$d_i^2=\|\stackrel{~}{d}_i\|^2.$$ (3.18) Eq. (3.8) and Eq. (3.9) are satisfied by construction; Eq. (3.7) gives for $`a=b=1,2,\mathrm{}r`$: $$d_i\stackrel{~}{d}_i\frac{1}{n_c}\underset{k}{}d_k\stackrel{~}{d}_k=\mu m_i;$$ (3.19) from which one finds that $$\underset{k}{}d_k\stackrel{~}{d}_k=\frac{n_c}{n_cr}\mu \underset{k=1}{\overset{r}{}}m_k$$ (3.20) and $$d_i\stackrel{~}{d}_i=\mu m_i+\frac{1}{n_cr}\mu \underset{k=1}{\overset{r}{}}m_k(d_i>0).$$ (3.21) On the other hand, Eq. (3.7) for $`a=b=r+1,r+2,\mathrm{}n_c`$ gives $$\underset{k}{}d_k\stackrel{~}{d}_k=n_c\mu c.$$ (3.22) This is compatible with Eq. (3.20) because of (3.15). A solution with a given $`r`$ leaves a local $`SU(n_cr)`$ invariance. Thus each of them counts as a set of $`n_cr`$ solutions (Witten’s index). In all, therefore, there are $$𝒩=\underset{r=0}{\overset{\text{min}\{n_f,n_c1\}}{}}(n_cr)\left(\begin{array}{c}n_f\\ r\end{array}\right)$$ (3.23) classical solutions. (For $`r=0`$, $`Q_a^i=\stackrel{~}{Q}_i^b=0,`$ $`\mathrm{\Phi }=0`$ is obviously a solution with full $`SU(n_c)`$ invariance.) For $`n_c=2`$ the formula (3.23) reproduces the known result ($`𝒩=2+n_f`$) as can be easily verified. Note that when $`n_f`$ is equal to or less than $`n_c`$ the sum over $`r`$ is done readily, and Eq. (3.23) is equivalent to $$𝒩_1=(2n_cn_f)\mathrm{\hspace{0.17em}2}^{n_f1},(n_fn_c).$$ (3.24) ### 3.2 Semi-Classical Vacua in $`USp(2n_c)`$ The superpotential reads in this case $$W=\mu \text{Tr}\mathrm{\Phi }^2+\frac{1}{\sqrt{2}}Q_a^i\mathrm{\Phi }_b^aQ_c^iJ^{bc}+\frac{m_{ij}}{2}Q_a^iQ_b^jJ^{ab},$$ (3.25) where $`J=i\sigma _2\mathrm{𝟏}_{n_c}`$ and $$m=i\sigma _2\text{diag}(m_1,m_2,\mathrm{},m_{n_f}).$$ (3.26) In the $`m_i0`$ and $`\mu 0`$ limit, the global symmetry is $`SO(2n_f)\times Z_{2n_c+2n_f}\times SU_R(2)`$. The vacuum equations are: $$[\mathrm{\Phi },\mathrm{\Phi }^{}]=0;$$ (3.27) $$\underset{i}{}(Q_a^iQ_b^iQ_{n_c+b}^iQ_{n_c+a}^i)=0;\underset{i}{}Q_a^iQ_{n_c+b}^i=0;$$ (3.28) $$\sqrt{2}\mathrm{\Phi }_b^aQ_c^iJ^{bc}+m_{ij}Q_b^jJ^{ab}=0.$$ (3.29) $$2\mu \mathrm{\Phi }_a^b+\frac{1}{\sqrt{2}}Q_a^iJ^{bc}Q_c^i=0.$$ (3.30) To find the solutions of these equations diagonalize first $`\mathrm{\Phi }`$ by a unitary transformation: $$\mathrm{\Phi }=\text{diag}(\varphi _1,\varphi _2,\mathrm{},\varphi _{n_c},\varphi _1,\varphi _2,\mathrm{},\varphi _{n_c}).$$ (3.31) Define $$\stackrel{~}{Q}_a^iQ_{n_c+a}^i.$$ (3.32) The vacuum equations can be rewritten as: $$\underset{i}{}(Q_a^iQ_b^i\stackrel{~}{Q}_b^i\stackrel{~}{Q}_a^i)=0;\underset{i}{}Q_a^i\stackrel{~}{Q}_b^i=0;(i=1,\mathrm{}2n_f);$$ (3.33) and $$\sqrt{2}\varphi _a\stackrel{~}{Q}_a^im_i\stackrel{~}{Q}_a^{n_f+i}=0;$$ (3.34) $$\sqrt{2}\varphi _a\stackrel{~}{Q}_a^{n_f+i}+m_i\stackrel{~}{Q}_a^i=0;$$ (3.35) $$\sqrt{2}\varphi _aQ_a^i+m_iQ_a^{n_f+i}=0;$$ (3.36) $$\sqrt{2}\varphi _aQ_a^{n_f+i}m_iQ_a^i=0;$$ (3.37) $$Q_a^i\stackrel{~}{Q}_b^i=0,(ab);$$ (3.38) and $$2\sqrt{2}\mu \varphi _a+Q_a^i\stackrel{~}{Q}_a^i+Q_a^{n_f+i}\stackrel{~}{Q}_a^{n_f+i}=0.$$ (3.39) In Eq. (3.34) – Eq. (3.39) the index $`i`$ runs only over $`i=1,2,\mathrm{},n_f`$. The solutions can again be classified by the number of the nonzero $`\varphi `$’s, $`r=1,2,\mathrm{},\mathrm{min}\{n_f,n_c\}`$. The one with $`r`$ masses $`m_1,\mathrm{},m_r`$ is $$\varphi =\frac{1}{\sqrt{2}}\text{diag}(im_1,im_2,\mathrm{},im_r,0,\mathrm{},im_1,im_2,\mathrm{}im_r,0,\mathrm{});$$ (3.40) $$Q_a^i=\left(\begin{array}{cccccccccccc}d_1& & & & & & id_1& & & & & \\ & d_2& & & & & & id_2& & & & \\ & & \mathrm{}& & & & & & \mathrm{}& & & \\ & & & d_r& & & & & & id_r& & \\ & & & & 0& & & & & & 0& \\ & & & & & \mathrm{}& & & & & & \mathrm{}\\ id_1& & & & & & d_1& & & & & \\ & id_2& & & & & & d_2& & & & \\ & & \mathrm{}& & & & & & \mathrm{}& & & \\ & & & id_r& & & & & & d_r& & \\ & & & & 0& & & & & & 0& \\ & & & & & \mathrm{}& & & & & & \mathrm{}\end{array}\right)$$ (3.41) where $$d_i=\sqrt{m_i\mu }.$$ (3.42) Each solution with $`r`$ leaves unbroken $`USp(2(n_cr))`$ hence counts as $`n_cr+2`$ solutions. The total number of the classical vacua is then $$𝒩=\underset{r=0}{\overset{min\{n_c,n_f\}}{}}(n_cr+1)\left(\begin{array}{c}n_f\\ r\end{array}\right).$$ (3.43) Note that for smaller values of $`n_f`$, , the sum over $`r`$ is easily done and an equivalent formula is $$𝒩=(2n_c+2n_f)\mathrm{\hspace{0.17em}2}^{n_f1},(n_fn_c).$$ (3.44) It is amusing (and reassuring) that different formulas, Eq.(3.23) , Eq.(3.24), Eq.(3.43) and Eq.(3.44) found here reproduce correctly the formula $$𝒩_{SU(2)}=n_f+2,$$ (3.45) for the $`SU(2)`$ gauge theory (which is the simplest case, both of $`SU(n_c)`$ and $`USp(2n_c)`$), for $`n_f=04`$. ## 4 Determination of Symmetry Breaking Patterns at Large $`\mu `$ In this section, we determine the number of $`N=1`$ vacua and the pattern of flavor symmetry breaking in each of them, in the limit $`m_i0`$. This can be done most easily by studying these $`SU(n_c)`$ and $`USp(2n_c)`$ theories at large adjoint mass $`\mu `$. The advantage of considering the theories at large $`\mu `$ is that the adjoint field can be integrated out from the theory: the resulting low-energy effective theory is an exactly known $`N=1`$ supersymmetric gauge theory, perturbed by certain superpotential terms suppressed by $`1/\mu `$. The dynamics of such a theory is known, either in terms of dynamical superpotential with confined meson or baryon degrees of freedom , or in a dual description using magnetic degrees of freedom . By minimizing the potential of such low-energy effective actions we find the symmetry breaking pattern in each $`N=1`$ vacua. Supersymmetry and holomorphy imply that there are no phase transition at finite $`\|\mu \|`$: the qualitative features found here, such as the unbroken symmetry and the number of the Nambu-Goldstone bosons, are valid also at small nonvanishing $`\mu `$. This method not only allows us to cross-check the counting of the number of vacua in the previous section, but also enables us to determine the pattern of dynamical symmetry breaking from the first principles, as the limit of massless quarks can be studied exactly. To be precise, we shall be interested here in the limit $`m_i0`$ with $`\mu `$ and $`\mathrm{\Lambda }`$ ($`\mu \mathrm{\Lambda }`$) fixed. As long as $`\mu `$ is kept fixed, the limit $`m_i0`$ is smooth and it is possible to determine how the exact flavor symmetry is realized in each vacuum. This is to be contrasted to another limit, $`\mu \mathrm{}`$ with $`m_i`$ fixed, which will be discussed in the next section. This latter limit, which does not commute with the former ($`m_i0`$ first), is the relevant one for studying the decoupling the adjoint field and verifying the consistency with the known results in the standard $`N=1`$ theories. The analysis of this section will be divided in two parts, for small and for large values of $`n_f`$: this is necessary due to the emergence of the dual gauge group in the corresponding $`N=1`$ theories for relatively large values of the flavor ($`n_f>n_c+1`$ in $`SU(n_c)`$, $`n_f>n_c+2`$ in $`USp(2n_c)`$). For the ease of reading, a short summary is given at the end of this section: the reader who is more interested in the physics results than the technical aspects of the analysis, might well jump to it. ### 4.1 $`SU(n_c)`$ theories: Small Numbers ($`n_fn_c`$) of Flavor Let us first consider the cases, $`n_f<n_c`$. At large $`\mu `$ one can integrate over $`\mathrm{\Phi }`$: $$\mathrm{\Phi }^A=\frac{\sqrt{2}}{\mu }\stackrel{~}{Q}_i^a(\lambda ^A)_a^bQ_b^i;$$ (4.1) resubstituting it into the superpotential one gets $$W=\frac{1}{2\mu }\left[\text{Tr}M^2\frac{1}{n_c}(\text{Tr}M)^2\right]+\text{Tr}(Mm)+(n_cn_f)\frac{\mathrm{\Lambda }_1^{(3n_cn_f)/(n_cn_f)}}{(detM)^{1/(n_cn_f)}},$$ (4.2) where $`M_i^j\stackrel{~}{Q}_i^aQ_a^j`$ and where the known instanton–induced effective superpotential of $`N=1`$ SQCD has been added. $$\mathrm{\Lambda }_1\mu ^{\frac{n_c}{3n_cn_f}}\mathrm{\Lambda }^{\frac{2n_cn_f}{3n_cn_f}}$$ (4.3) is the scale of the $`N=1`$ SQCD. By making an ansatz for $`M`$: $$M=\text{diag}(\lambda _1,\lambda _2,\mathrm{},\lambda _{n_f}),$$ (4.4) the equations for $`\lambda _i`$ are: $$\frac{1}{\mu }\left(\lambda _i\frac{1}{n_c}\underset{j}{}\lambda _j\right)+m_i\frac{\mathrm{\Lambda }_1^{(3n_cn_f)/(n_cn_f)}}{(_j\lambda _j)^{1/(n_cn_f)}}\lambda _i^1=0.$$ (4.5) We now study the solutions of these equations in the limit $`m_i0`$, with $`\mathrm{\Lambda }`$ and $`\mu `$ ($`\mu \mathrm{\Lambda }`$) fixed. By making an ansatz, $$M=\text{diag}(\lambda _1,\lambda _2,\mathrm{},\lambda _{n_f}),$$ (4.6) and upon multiplication with $`\lambda _i`$ one finds (for each $`i`$) $$\frac{1}{\mu }(\lambda _i^2\lambda _iY)+m_i\lambda _i+X=0.$$ (4.7) where $$X\mathrm{\Lambda }_1^{(3n_cn_f)/(n_cn_f)}(\underset{j}{}\lambda _j)^{\frac{1}{n_fn_c}};Y\frac{1}{n_c}\underset{j}{\overset{n_f}{}}\lambda _j.$$ (4.8) One can take the limit $`m_i0`$ directly in Eq.(4.7), which becomes simply $$\lambda _i^2Y\lambda _i\mu X=0.$$ (4.9) It can be solved by first assuming that $`X`$ and $`Y`$ are given: $$\lambda _i=\frac{1}{2}(Y\pm \sqrt{Y^2+4\mu X}).$$ (4.10) In general, $`r`$ of $`\lambda _i`$’s can take the upper sign and the rest ($`n_fr`$) of $`\lambda _i`$’s the lower sign, $`r=0,1,2,\mathrm{},n_f`$. These solutions $`\lambda _1=\mathrm{}=\lambda _r={\displaystyle \frac{1}{2}}(Y+\sqrt{Y^2+4\mu X});`$ $`\lambda _{r+1}=\mathrm{}=\lambda _{n_f}={\displaystyle \frac{1}{2}}(Y\sqrt{Y^2+4\mu X}),`$ (4.11) must then be re-inserted into the definitions of $`X`$ and $`Y`$ to determine the latter quantities. One finds two relations $$(2rn_f)\sqrt{Y^2+4\mu X}=(2n_cn_f)Y;$$ (4.12) and $$X=\frac{\mathrm{\Lambda }_1^{(3n_cn_f)/(n_cn_f)}}{2^{n_f/(n_fn_c)}}\left(Y+\sqrt{Y^2+4\mu X}\right)^{\frac{r}{n_fn_c}}\left(Y\sqrt{Y^2+4\mu X}\right)^{\frac{n_fr}{n_fn_c}}.$$ (4.13) These two equations can be easily solved for $`X`$ and $`Y`$ and give $$Y=C\frac{\mathrm{\Lambda }_1^{(3n_cn_f)/(2n_cn_f)}}{\mu ^{(n_fn_c)/(2n_cn_f)}}e^{2\pi ki/(2n_cn_f)}=C\mathrm{\Lambda }\mu e^{2\pi ki/(2n_cn_f)},k=1,2,\mathrm{}2n_cn_f;$$ (4.14) $$X=\frac{C^{}}{\mu }Y^2.$$ (4.15) where $`C,C^{}`$ are constants depending on $`n_f`$, $`n_c`$ and $`r`$. Note that $`X`$ is uniquely determined in terms of $`Y`$. At given $`r`$, then, there are $`2n_cn_f`$ solutions for $`(X,Y)`$ hence for $`\{\lambda _i\}`$. By summing over $`r`$, taking into account $`{}_{n_f}{}^{}C_{r}^{}`$ ways of distributing $`r`$ solutions with positive sign among $`n_f`$ flavor, one appears to end up with $`(2n_cn_f)_{r=0}^{n_f}{}_{n_f}{}^{}C_{r}^{}=(2n_cn_f)2^{n_f}`$ vacua. Actually, one counts exactly twice each vacuum this way. The solution for $`\{\lambda _i\}`$’s depends on the value of $`r`$ in a non trivial manner, in general. When $`r`$ is replaced by $`n_fr`$, however, $`C`$ remains unchanged: the net change of $`(X,Y)`$ and hence of $`\{\lambda _i\}`$’s, is that the two types of roots Eq.(4.10) are precisely interchanged, as can be seen from Eq.(4.12), Eq.(4.13), giving the same set of $`\{\lambda _i\}`$’s. We find therefore $$𝒩_1=(2n_cn_f)2^{n_f1}$$ (4.16) solutions of this type (i.e., finite in the $`m_i0`$ limit). For $`n_f<n_c`$ these exhaust all possible solutions (see Eq.(3.24), Eq.(5.11)). They are classified by the value of an integer $`r`$: in a vacuum characterized by $`r`$, $`U(n_f)`$ symmetry of the theory is broken spontaneously to $`U(r)\times U(n_fr)`$ by the condensates, Eq.(4.11). The analysis for the case of $`n_f=n_c`$ is similar but can be made by using the superpotential valid for $`n_c=n_f`$ $$W=\frac{1}{2\mu }\left[\text{Tr}M^2\frac{1}{n_c}(\text{Tr}M)^2\right]+\text{Tr}(Mm)+\text{Tr}\kappa \{(detM)B\stackrel{~}{B}\mathrm{\Lambda }^{2n_c}\},$$ (4.17) where $`B`$ and $`\stackrel{~}{B}`$ are baryonlike composite, $`B=ϵ_{i_1i_2\mathrm{}i_{n_c}}ϵ^{a_1a_2\mathrm{}a_{n_c}}Q_{a_1}^{i_1}Q_{a_2}^{i_2}\mathrm{}Q_{a_{n_c}}^{i_{n_c}}`$ and analogously for $`\stackrel{~}{B}`$ in terms of $`\stackrel{~}{Q}`$’s, and $`\kappa `$ is a Lagrange multiplier. ### 4.2 $`SU(n_c)`$: $`n_f=n_c+1`$ In the case with $`n_f=n_c+1`$ the effective superpotential is $`W`$ $`=`$ $`{\displaystyle \frac{1}{2\mu }}\left[\text{Tr}M^2{\displaystyle \frac{1}{n_c}}(\text{Tr}M)^2\right]+\text{Tr}(Mm)`$ (4.18) $`+`$ $`{\displaystyle \frac{1}{\mathrm{\Lambda }_1^{2n_f3}}}\{(detM)B_i(M)_j^i\stackrel{~}{B}^j\},`$ where $$B_i=ϵ_{ii_1i_2\mathrm{}i_{n_c}}ϵ^{a_1a_2\mathrm{}a_{n_c}}Q_{a_1}^{i_1}Q_{a_2}^{i_2}\mathrm{}Q_{a_{n_c}}^{i_{n_c}}.$$ (4.19) Set $$M=\text{diag}(\lambda _1,\mathrm{},\lambda _{n_f}),$$ (4.20) $`W`$ $`=`$ $`{\displaystyle \frac{1}{2\mu }}\left[{\displaystyle \underset{i}{}}\lambda _i^2{\displaystyle \frac{1}{n_c}}\left({\displaystyle \underset{i}{}}\lambda _i\right)^2\right]+{\displaystyle \underset{i}{}}m_i\lambda _i`$ (4.21) $`+`$ $`{\displaystyle \frac{1}{\mathrm{\Lambda }_1^{2n_f3}}}\{({\displaystyle \underset{j}{}}\lambda _j)B_i\stackrel{~}{B}^i\lambda _i\}.`$ Derivation with respect to $`B_j`$ and $`\stackrel{~}{B}^i`$ yields $$\lambda _i\stackrel{~}{B}^i=0;\lambda _iB_i=0,$$ (4.22) while derivation with respect to $`\lambda _i`$ leads to $$\frac{1}{\mu }\left[\lambda _i\frac{1}{n_c}\underset{j}{}\lambda _j\right]+m_i+\frac{1}{\mathrm{\Lambda }_1^{2n_f3}}\{\underset{ji}{}\lambda _jB_i\stackrel{~}{B}^i\}=0.$$ (4.23) Multiplying with $`\lambda _i`$ one gets (still no sum over $`i`$) $$\frac{1}{\mu }\left[\lambda _i^2\frac{1}{n_c}\lambda _i\underset{j}{}\lambda _j\right]+m_i\lambda _i+\frac{1}{\mathrm{\Lambda }_1^{2n_f3}}\underset{j}{}\lambda _j=0.$$ (4.24) Set now $`Y=\frac{1}{n_c}(_j\lambda _j)`$, and $`X=_j\lambda _j,`$ and set $`\mathrm{\Lambda }_11`$ from now on. $`\lambda _i`$ satisfies $$\lambda _i^2(Y+m_i\mu )\lambda _i\mu X=0.$$ (4.25) The solution is $$\lambda _i=\frac{1}{2}(Y+m_i\mu \pm \sqrt{(Y+m_i\mu )^24\mu X}).$$ (4.26) Note that when $`m_i=0`$, $`\lambda =B=0`$, all $`i`$’s, is a solution. We now classify the solutions according to the number ($`r`$) of $`\lambda `$’ s which remain finite in the limit $`m_i=0`$. i) First consider the solution $`\lambda _i0`$, $`i`$. In this case $`B_i=\stackrel{~}{B}^i=0`$ , $`i`$. $`\lambda _i`$ are given by $$\lambda _i=\frac{1}{2}(Y\pm \sqrt{Y^24\mu X}).$$ (4.27) The rest of the argument is the same as the one given Sec. 4.1 so one has $`(2n_cn_f)2^{n_f1}`$ vacua of this kind. ii) Consider now $`\lambda _1=0`$, $`\lambda _j0,j1.`$ One finds that $$B_j=\stackrel{~}{B}^j=0,j1;$$ (4.28) $$\frac{1}{n_c}\underset{j1}{}\lambda _j\frac{\mu }{\mathrm{\Lambda }^{2n_f3}}B_1\stackrel{~}{B}^1=0;$$ (4.29) $$\lambda _j\frac{1}{n_c}\underset{k1}{}\lambda _km_j\mu =0.$$ (4.30) The last equation has the form, $$C\left(\begin{array}{c}\lambda _2\\ \lambda _3\\ \mathrm{}\\ \lambda _{n_f}\end{array}\right)\mu \left(\begin{array}{c}m_2\\ m_3\\ \mathrm{}\\ m_{n_f}\end{array}\right)=0,$$ (4.31) with $$detC=0.$$ (4.32) For generic $`m_i`$’s therefore this equation has no solutions. Here it is important that we consider the massless limit of massive theory, not directly massless theory. iii) Consider now $`\lambda _i=0`$, $`i=1,2,\mathrm{},r,`$ and $`\lambda _j0,jr+1.`$ One gets $$B_j=\stackrel{~}{B}^j=0,j=r+1,\mathrm{}n_f$$ (4.33) and $$\frac{1}{n_c}\underset{j>r}{}\lambda _j\frac{\mu }{\mathrm{\Lambda }^{2n_f3}}B_i\stackrel{~}{B}^i=0,i=1,2,\mathrm{},r;$$ (4.34) $$\lambda _j\frac{1}{n_c}\underset{k>r}{}\lambda _km_j\mu =0,j=r+1,\mathrm{}n_f.$$ (4.35) The last equations now give finite answers for $`\lambda _j`$: but they are of $`O(m_i)`$ and approach $`0`$ as $`m_i0`$. One gets then also, $$B_i,\stackrel{~}{B}^i0,m_i0.$$ (4.36) Therefore all these cases degenerate into one single solution, $$\lambda =B=\stackrel{~}{B}=0.$$ (4.37) To conclude, we find in the massless limit $`(2n_cn_f)2^{n_f1}`$ solutions with finite $`\lambda `$’s and one solution with vanishing vevs. The latter is consistent with the general formula Eq.(4.47) for the second group of vacua (with vanishing VEVS) found for $`n_f>n_f+1`$ below, since $`𝒩_2=_{r=0}^{\stackrel{~}{n}_c1}{}_{n_f}{}^{}C_{r}^{}(\stackrel{~}{n}_cr)=1,`$ for $`n_f=n_c+1`$. ### 4.3 $`SU(n_c)`$: large numbers ($`n_f>n_c+1`$) of flavor In the cases $`n_f>n_c+1`$, the effective low energy degrees of freedom are dual quarks and mesons . The effective superpotential is given by $$𝒲=\stackrel{~}{q}Mq+\text{Tr}(mM)\frac{1}{2\mu }\left[\text{Tr}M^2\frac{1}{n_c}(\text{Tr}M)^2\right],$$ (4.38) where $`q`$’s are $`n_f`$ sets of dual quarks in the fundamental representation of the dual gauge group $`SU(\stackrel{~}{n}_c)`$, with $`\stackrel{~}{n}_c=n_fn_c`$. The vacuum equations following from Eq. (4.38) are: $$M_{ij}q_j^\alpha =0;\stackrel{~}{q}_\alpha ^iM_{ij}=0;$$ (4.39) $$\stackrel{~}{q}_\alpha ^iq_j^\alpha +\delta _{ij}m_i\frac{1}{\mu }\left(M_{ij}\frac{1}{n_c}(\text{Tr}M)\delta _{ij}\right)=0.$$ (4.40) The first set of equations tell us that the meson matrix $`M`$ and the dual squarks are orthogonal in the flavor space. By using the dual color and flavor rotations (and the use of Eq. (4.40)) the dual squarks can be taken to be nonvanishing in the first $`r`$ flavors and of the form: $$q_i^\alpha =\left(\begin{array}{c}d_1\\ 0\\ \mathrm{}\\ \mathrm{}\\ 0\end{array}\right),\left(\begin{array}{c}0\\ d_2\\ 0\\ \mathrm{}\\ 0\end{array}\right),\left(\begin{array}{c}0\\ \mathrm{}\\ d_r\\ 0\\ \mathrm{}\end{array}\right);q_i^\alpha =0,i=r+1,\mathrm{},n_f.$$ (4.41) $$\stackrel{~}{q}_\alpha ^i=\left(\begin{array}{c}\stackrel{~}{d}_1\\ 0\\ \mathrm{}\\ \mathrm{}\\ 0\end{array}\right),\left(\begin{array}{c}0\\ \stackrel{~}{d}_2\\ 0\\ \mathrm{}\\ 0\end{array}\right),\left(\begin{array}{c}0\\ \mathrm{}\\ \stackrel{~}{d}_r\\ 0\\ \mathrm{}\end{array}\right);\stackrel{~}{q}_\alpha ^i=0,i=r+1,\mathrm{},n_f,$$ (4.42) where<sup>6</sup><sup>6</sup>6 Note that the value $`r=\stackrel{~}{n}_c`$ should be excluded. In this case, $`n_fr=n_c`$ and the nonvanishing meson submatrix is $`n_c\times n_c`$. Eq. (4.40) for $`i,j=r+1,\mathrm{},n_f`$ have no solutions since the matrix $`M_{ij}\frac{1}{n_c}(\text{Tr}M)\delta _{ij}`$ is of rank $`n_c1`$ while $`\delta _{ij}m_i`$ has rank $`n_c`$. $$r=0,1,2,\mathrm{},\stackrel{~}{n}_c1,$$ (4.43) and $$d_i\stackrel{~}{d}_i+m_i+\frac{1}{\mu n_c}\text{Tr}M=0.$$ (4.44) Eq. (4.40) implies also that the meson matrix is diagonal, $$M=\text{diag}(0,0,\mathrm{},0,\lambda _{r+1},\lambda _{r+2},\mathrm{},\lambda _{n_f}),$$ (4.45) $$\lambda _i\frac{1}{n_c}\lambda =m_i\mu .$$ (4.46) Clearly the last equations determine uniquely all $`\lambda `$’s: on the other hand there are $`{}_{n_f}{}^{}C_{r}^{}`$ choices of masses which enter the equations for nonzero squark VEVS. Furthermore, because the vacua with $`r`$ nonzero entries in the squark vevs leave an $`SU(\stackrel{~}{n}_cr)`$ dual gauge group unbroken ($`M`$ being singlet), each such vacuum must be counted as $`\stackrel{~}{n}_cr`$ vacua (Witten’s index). In all, then, there are $$𝒩_2=\underset{r=0}{\overset{\stackrel{~}{n}_c1}{}}{}_{n_f}{}^{}C_{r}^{}(\stackrel{~}{n}_cr)$$ (4.47) vacua, with $`\text{Rank}M<n_f`$, in which the global $`U(n_f)`$ symmetry remains unbroken, in the $`m_i0`$ limit. We seem to face a difficulty, however. The number of vacua found here is less than the known total number of vacua $`𝒩`$ (Eq.(3.23)). Where are other vacua? This apparent puzzle can be solved once the nontrivial $`SU(\stackrel{~}{n}_c)`$ dynamics are taken into account<sup>7</sup><sup>7</sup>7 In fact, a related puzzle is how Seiberg’s dual Lagrangian \- the first two terms of Eq. (4.38) - can give rise to the right number of vacua for the massive $`N=1`$ SQCD with $`n_f>n_c+1`$. By following the same method as below but with $`\mu =\mathrm{}`$, we do find the correct number ($`n_c`$) of vacua. . If the VEVS of the mesons have Rank$`M=n_f,`$ dual quarks are all massive. The theory becomes pure Yang-Mills type in the infrared, and the strong interaction effects of dual gauge dynaimcs must be properly taken into account. By integrating out the dual quarks out, we find the effective superpotential, $$W_{eff}=\frac{1}{2\mu }\left[\text{Tr}M^2\frac{1}{n_c}(\text{Tr}M)^2\right]+\text{Tr}(Mm)+\mathrm{\Lambda }_1^{(3n_cn_f)/(n_cn_f)}(detM)^{1/(n_fn_c)}.$$ (4.48) The analysis of the vacua of this effective action is similar to that of the $`n_f<n_c`$ cases, Eqs.(4.5)–(4.15): the superpotential (4.48) is actually identical to (more precisely, continuation of) (4.2)!<sup>8</sup><sup>8</sup>8 In the case of SQCD ($`\mu =\mathrm{}`$) this observation is in agreement with the well known fact that, in spite of distinct physical features at $`n_fn_c`$ and $`n_fn_c+1`$, the results for the squark and gaugino condensates, $`m_i\stackrel{~}{Q}^iQ^i=\lambda \lambda =(detm)^{1/n_c}\mathrm{\Lambda }_1^{3n_f/n_c}`$ hold true for all values of the flavor and color as long as $`n_f<3n_c`$ . One finds therefore $`𝒩_1=(2n_cn_f)2^{n_f1}`$ solutions with finite VEVS in exactly the same way as in Eqs.(4.5)–(4.15). In these vacua, classified by an integer $`r`$, $`U(n_f)`$ symmetry of the theory is broken spontaneously to various $`U(r)\times U(n_fr).`$ It is now possible to make a highly nontrivial consistency test, by the vacuum counting. By using $`\stackrel{~}{n}_c=n_fn_c`$ and changing the summation index from $`r`$ to $`n_fr`$, it is easy to show that $`𝒩_2`$ $`=`$ $`{\displaystyle \underset{r=0}{\overset{n_f}{}}}{}_{n_f}{}^{}C_{r}^{}(rn_c){\displaystyle \underset{r=0}{\overset{n_c}{}}}{}_{n_f}{}^{}C_{r}^{}(rn_c)`$ (4.49) $`=`$ $`(2n_cn_f)2^{n_f1}+{\displaystyle \underset{r=0}{\overset{n_c1}{}}}{}_{n_f}{}^{}C_{r}^{}(n_cr).`$ The total number of the $`N=1`$ vacua is therefore found to be $$𝒩_1+𝒩_2=\underset{r=0}{\overset{n_c1}{}}{}_{n_f}{}^{}C_{r}^{}(n_cr),$$ (4.50) which is the correct answer (see Eq. (3.23)) for $`n_f>n_c+1`$! ### 4.4 $`USp(2n_c)`$ Theories: Small Numbers ($`n_fn_c+1`$) of Flavor We start with the cases $`n_fn_c+1`$. At large $`\mu `$ the equation of motion for the $`\mathrm{\Phi }`$ superfield is: $$\mathrm{\Phi }^A=\frac{1}{\sqrt{2}\mu }(Q_\alpha ^iS_{\alpha \beta }^AQ_\beta ^i),$$ (4.51) where $`\mathrm{\Phi }=\mathrm{\Phi }^AS^A`$ and $`S^A`$ are the $`USp(2n_c)`$ generators<sup>9</sup><sup>9</sup>9 With respect to the $`\widehat{S}^A`$ generators defined in App. A, $`S^A=\widehat{S}^AJ`$ with $`S_{\alpha \beta }^A=S_{\beta \alpha }^A`$ and $`\text{Tr}S^AJS^BJ=\frac{1}{2}\delta ^{AB}`$.. Resubstituting it into the superpotential (3.25), and accounting for the instanton–induced contribution for $`n_fn_c`$, one gets: $$W=\frac{1}{8\mu }\text{Tr}MJMJ\frac{1}{2}\text{Tr}mM+(n_c+1n_f)\frac{\mathrm{\Lambda }_1^{3+2n_f/(n_c+1n_f)}}{(\mathrm{Pf}M)^{1/(n_c+1n_f)}},$$ (4.52) where $`M^{ij}=Q_a^iJ^{ab}Q_b^j`$’s are the meson–like composite superfields, $`m=i\sigma _2\text{diag}(m_1,\mathrm{},m_{n_f})`$ is the mass matrix, and $`\mathrm{\Lambda }_1^{2(3n_c+3n_f)}=\mu ^{2n_c+2}\mathrm{\Lambda }^{2(2n_c+2n_f)}`$. With the ansatz $`M=i\sigma _2\mathrm{diag}(\lambda _1,\lambda _2,\mathrm{},\lambda _{n_f})`$, then the superpotential is ($`\mathrm{\Lambda }_1=1`$ below): $$W=\frac{1}{4\mu }\underset{i=1}{\overset{n_f}{}}\lambda _i^2\underset{i=1}{\overset{n_f}{}}m_i\lambda _i+(n_c+1n_f)\frac{1}{(_{i=1}^{n_f}\lambda _i)^{1/(n_c+1n_f)}},$$ (4.53) and the vacuum equations become: $$\frac{1}{2\mu }\lambda _im_i\frac{1}{\lambda _i}\frac{1}{(_j\lambda _j)^{1/(n_c+1n_f)}}=0.$$ (4.54) Since the last term is common to all $`i`$, we find $$\lambda _i=\mu (m_i\pm \sqrt{m_i^2+(2X/\mu )}),X\frac{1}{(_j\lambda _j)^{1/(n_c+1n_f)}}$$ (4.55) We choose $`r`$ negative signs and $`n_fr`$ positive signs in the roots of $`\lambda _i`$. In strictly massless limit $`m=0`$, the definition of $`X`$ gives<sup>10</sup><sup>10</sup>10 Eq. (4.55) and analogous relations in $`SU(n_c)`$ case clearly show the non commutativity of the two limits, $`\mu \mathrm{}`$ first with $`\mathrm{\Lambda }_1,m_i`$ fixed to be studied in the next section, and $`m_i0`$ first with $`\mu \mathrm{\Lambda }`$ fixed, being examined here. $$X\frac{1}{(_j\lambda _j)^{1/(n_c+1n_f)}}=\frac{1}{((1)^r(2X\mu )^{n_f/2})^{1/(n_c+1n_f)}}$$ (4.56) and hence $$X_0=\mathrm{\Lambda }_1^{\frac{2(3n_c+3n_f)}{2n_c+2n_f}}(2\mu )^{\frac{n_f}{2n_c+2n_f}}e^{2\pi ik/(2n_c+2n_f)},k=1,2,\mathrm{},2n_c+2n_f,$$ (4.57) where the subscript 0 indicates that this is for the massless quark limit and we have reinstated the dependence on the scale $`\mathrm{\Lambda }_1`$. Note that $$X_0\mu \mathrm{\Lambda }^2$$ (4.58) in terms of the $`N=2`$ scale factor $`\mathrm{\Lambda }`$ and $`\mu `$. There are $`(2n_c+2n_f)`$ roots for $`X_0`$. Up to $`O(m)`$, we find $$X=X_0\left[1+\frac{2(_{j=1}^rm_j_{j=r+1}^{n_f}m_j)}{2n_c+2n_f}\left(\frac{\mu }{2X_0}\right)^{1/2}\right],$$ (4.59) and $$\lambda _i=\pm \mu \left[\left(\frac{2X_0}{\mu }\right)^{1/2}+\frac{_{j=1}^rm_j_{j=r+1}^{n_f}m_j}{2n_c+2n_f}\pm m_i\right]$$ (4.60) There appears to be $`2^{n_f}`$ choices for the signs among $`\lambda _i`$’s; actually, the signs of $`\lambda _i`$ must be such that Eq. (4.56), and not its $`2(n_c+1n_f)`$ th power, is satisfied. This restricts the choices of the signs by half: for a particular phase of $`X_0`$ with $`k`$ even or odd, the number of minus signs among $`\lambda _i`$ must be even or odd, respectively. In total, there are $$(2n_c+2n_f)\mathrm{\hspace{0.17em}2}^{n_f1}$$ (4.61) vacua, which is consistent with the previous method as well as the semi-classical method in the previous section (Eq. (3.44)). Since the massless limit gives $`SO(2n_f)`$ flavor symmetry, the choice of different signs in $`\lambda _i`$ for each flavor strongly suggests the vacua to form a spinor representation of $`SO(2n_f)`$. The constraint that the number of minus signs to be even or odd implies that it is a spinor representation of a definite chirality. In fact, all of these vacua transform among each other under $`SO(2n_f)`$ group because it can flip the signs of two eigenvalues at the same time, consistent with an irreducible representation of $`SO(2n_f)`$. The equal mass case $`m_i=m`$ has $`U(n_f)`$ flavor symmetry and the vacua above form $`{}_{n_f}{}^{}C_{r}^{}`$ multiplets (with even or odd $`r`$), consistent with the decomposition of the $`SO(2n_f)`$ spinor to even or odd-rank $`r`$ anti-symmetric tensors under $`U(n_f)`$. In each group of vacua, the global symmetry is broken to $`U(r)\times U(n_fr)`$. One of the most important results of this section is that the meson condensates in the massless limit always break $`SO(2n_f)`$ flavor symmetry as $$SO(2n_f)U(n_f),$$ (4.62) matching nicely with the expectation discussed in Section 2. When $`n_f=n_c+1`$, the large $`\mu `$ theory develops a quantum modified constraint $$W=\frac{1}{8\mu }\text{Tr}M^2\frac{1}{2}\text{Tr}mM+X(\mathrm{Pf}M\mathrm{\Lambda }_1^{2n_f}).$$ (4.63) Following the similar analysis as above, we again find the total number of vacua to be $`(2n_c+2n_f)\mathrm{\hspace{0.17em}2}^{n_f1}`$, consistent with the semi-classical method. When $`n_f=n_c+2`$, the large $`\mu `$ theory develops a superpotential $$W=\frac{1}{8\mu }\text{Tr}M^2\frac{1}{2}\text{Tr}mM+\frac{\mathrm{Pf}M}{\mathrm{\Lambda }_1^{2n_f3}}.$$ (4.64) Following the similar analysis as above, we again find the total number of vacua to be $`n_c\mathrm{\hspace{0.17em}2}^{n_f1}+1`$, consistent with the semi-classical method. The last vacuum corresponds to the case without symmetry breaking $`\lambda _i=2m_i\mu 0`$ in the massless quark limit. ### 4.5 $`USp(2n_c)`$ Theories: Large Numbers ($`n_f>n_c+2`$) of Flavor Next consider the cases $`n_f>n_c+2`$. The large $`\mu `$ theory has a description in terms of the dual magnetic gauge group $`USp(2\stackrel{~}{n}_c)=USp(2(n_fn_c2))`$ and magnetic quarks $`q`$, $$W=\frac{1}{8\mu }\text{Tr}M^2\frac{1}{2}\text{Tr}mM+\frac{1}{\mu _m}qMq,$$ (4.65) where the scale $`\mu _m`$ is the matching scale between the electric and magnetic gauge couplings. As in the $`SU(n_c)`$ cases above, it is only consistent to use this effective action to get information on the vacuum properties as long as dual quarks turn out to be light. Otherwise, the nontrivial dual gauge dynamics must be taken into acoount. By minimizing the potential of the magnetic theory with dual quarks directly, we find vacua characterized by VEVS $$q_i=\sqrt{\mu m_i},\lambda _i=0,$$ (4.66) for $`r`$ flavors and $$q_i=0,\lambda _i=m_i\mu ,$$ (4.67) for the remaining $`n_fr`$ flavors. $`r`$ is restricted to be $`r\stackrel{~}{n}_c=n_fn_c2`$. For each value of $`r`$, there are $`{}_{n_f}{}^{}C_{r}^{}`$ choices on which flavor to have non-vanishing $`q_i`$, and there is unbroken $`USp(2(n_fn_c2r))`$ gauge group. It develops the gaugino condensate, giving $`n_fn_c1r`$ vacua each. The number of vacua of this type is $$𝒩_2=\underset{r=0}{\overset{n_fn_c2}{}}{}_{n_f}{}^{}C_{r}^{}(n_fn_c1r).$$ (4.68) Changing the variable $`r`$ to $`n_fr`$, it is rewritten as $$𝒩_2=\underset{r=n_c+2}{\overset{n_f}{}}{}_{n_f}{}^{}C_{r}^{}(rn_c1)=\underset{r=0}{\overset{n_f}{}}{}_{n_f}{}^{}C_{r}^{}(rn_c1)+\underset{r=0}{\overset{n_c+1}{}}{}_{n_f}{}^{}C_{r}^{}(n_c+1r).$$ (4.69) The first sum can be computed and gives $$𝒩_2=(2n_c+2n_f)\mathrm{\hspace{0.17em}2}^{n_f1}+\underset{r=0}{\overset{n_c}{}}{}_{n_f}{}^{}C_{r}^{}(n_c+1r).$$ (4.70) Note that the sum in the second term can be stopped at $`r=n_c`$ because the argument vanishes for $`r=n_c+1`$. Therefore the total $`𝒩=𝒩_1+𝒩_2`$ agrees with the counting of classical vacua. On the other hand, this gives a nice interpretation of the number that the “extra” contribution $`𝒩_2`$ signals the emergence of the dual gauge group in the massless quark limit. In this group of vacua (present only for larger values of $`n_f`$), the chiral symmetry is not spontaneously broken in the limit, $`m_i0`$. In order to get the vacua with $`\text{Rank}M=n_f`$, one must integrate out the dual quarks first and consider the resulting effective action: $$W_{\mathrm{𝑒𝑓𝑓}}=\frac{1}{8\mu }\text{Tr}M^2\frac{1}{2}\text{Tr}mM+(n_c+1n_f)\left[\text{Pf}\frac{M}{\mu _m}\stackrel{~}{\mathrm{\Lambda }}^{3(n_fn_c1)n_f}\right]^{1/(n_fn_c1)}.$$ (4.71) By making the ansatz, $`M=i\sigma _2\mathrm{diag}(\lambda _1,\lambda _2,\mathrm{},\lambda _{n_f})`$, the vacuum equation becomes $$\left(\underset{i}{\overset{n_f}{}}\frac{\lambda _i}{\mu _m}\right)^{1/(n_fn_c1)}\stackrel{~}{\mathrm{\Lambda }}^{(3(n_fn_c1)n_f)/(n_fn_c1)}=m_i\lambda _i\frac{\lambda _i^2}{\mu }.$$ (4.72) Call the right hand side, which is flavor indepedent, $`X`$. We have two solutions for $`\lambda _i`$ for given $`X`$, $$\lambda _i=\frac{\mu }{2}\left[m_i\pm \sqrt{m_i^2+4X/\mu }\right].$$ (4.73) In the massless limit, $`\lambda _i=\sqrt{X\mu }`$, which in turns gives $$X=(X\mu )^{n_f/2(n_fn_c1)}\stackrel{~}{\mathrm{\Lambda }}^{(3(n_fn_c1)n_f)/(n_fn_c1)}\mu _m^{n_f/(n_fn_c1)}.$$ (4.74) The solution is given by $$X^{2n_c+2n_f}=\stackrel{~}{\mathrm{\Lambda }}^{2(3(n_fn_c1)n_f)}\mu _m^{2n_f}\mu ^{n_f}.$$ (4.75) This obviously gives $`2n_c+2n_f`$ solutions, for which there are $`2^{n_f1}`$ possibilities on the sign choices for each $`\lambda _i`$. Therefore we find $$𝒩_1=(2n_c+2n_f)\mathrm{\hspace{0.17em}2}^{n_f1}$$ (4.76) vacua. In this set of vacua $`X`$ depends on $`\stackrel{~}{\mathrm{\Lambda }}`$ and has a finite $`m_i0`$ limit (i.e., $`\lambda _i`$ stay non-vanishing, justifying the assumption of the maximal rank meson matrix). Since all $`\lambda _i`$’s are equal in the magnitude in this limit, the chiral $`SO(2n_f)`$ symmetry is spontaneously broken as $$SO(2n_f)U(n_f)$$ (4.77) in all vacua belonging to this group. Note that this number of vacua would precisely corresponds to that of monopole condensation in the $`SO(2n_f)`$ spinor representation at the Chebyshev points of the curve (see below). The total number of vacua at large $`\mu `$ found, $`𝒩_1+𝒩_2=_{r=0}^{n_c}{}_{n_f}{}^{}C_{r}^{}(n_c+1r)`$, agrees with that of the semiclassical theories. ### 4.6 Summary of Section 4 The number of $`N=1`$ vacua and the pattern of symmetry breaking in each of them has thus been determined in $`SU(n_c)`$ and $`USp(2n_c)`$ theories at large $`\mu `$, from the first principles. For small numbers of flavor ($`n_fn_c+1`$ for $`SU(n_c)`$; $`n_fn_c+2`$ for $`USp(2n_c)`$) the low–energy degrees of freedom are meson–like (sometimes also baryon–like) composites: their condensation lead to a definite pattern of symmetry breaking in each vacuum. $`SU(n_c)`$ theories have an exact global $`U(n_f)`$ symmetry in the equal mass (or massless) limit, which is spontaneously broken to $`U(r)\times U(n_fr)`$ in $`(n_cr){}_{n_f}{}^{}C_{r}^{}`$ vacua, $`r=0,1,\mathrm{},[n_f/2]`$. The number of the vacua with the particular pattern of symmetry breaking will match exactly with those found from the analysis of low-energy monopole/dual quark effective action, to be analyzed in the next sections. In $`USp(2n_c)`$ theories, for small numbers of flavor, the chiral symmetry ($`SO(2n_f)`$) in the massless limit is always spontaneouly broken down to an unbroken $`U(n_f)`$. This result nicely agrees with what is expected generally from bifermion condensates of the standard form in non supersymmetric theories. This fact that various vacua have exactly the same symmetry breaking pattern, has an important consequence in the physics at small $`\mu `$, to be studied below. The difference in the symmetry breaking pattern in $`SU(n_c)`$ and $`USp(2n_c)`$ theories reflects the structures of the low-energy effective actions of the respective theories, which in turn is a direct consequence of the different structure of the two types of gauge groups, see Eq.(4.2) and Eq.(4.52). For larger numbers of flavor ($`n_f>n_c+1`$ for $`SU(n_c)`$; $`n_f>n_c+2`$ for $`USp(2n_c)`$) the low-energy degrees of freedom are dual quarks and gluons, as well as some mesons. In these cases, besides the vacua with the properties mentioned above, other vacua exist in which all VEVS vanish and in which the global symmetry ($`SU(n_f)`$ for $`SU(n_c)`$; $`SO(2n_f)`$ in $`USp(2n_c)`$ theories) remains unbroken in the massless limit ($`m_i0,`$ $`i`$). ## 5 Decoupling of the Adjoint Fields In this section we discuss briefly another limit, in which $`\mu `$ is sent to $`\mathrm{}`$, keeping $`m_i`$ and $`\mathrm{\Lambda }_1\mu ^{\frac{n_c}{3n_cn_f}}\mathrm{\Lambda }^{\frac{2n_cn_f}{3n_cn_f}}`$ fixed. We first re-analyse the number of the $`N=1`$ vacua in the regime, $`\mu \mathrm{\Lambda }_1,m_i`$, and reproduce the correct multiplicity of vacua. Since the two limits ($`\mu \mathrm{}`$ first and $`m_i0`$ first) do not commute, this provides for an independent check of the vacuum counting. Subsequently, taking the decoupling limit, $`\mu \mathrm{}`$, we identify the standard supersymmetric vacua of the theories without the adjoint field $`\mathrm{\Phi }`$ . ### 5.1 $`SU(n_c)`$ First consider the cases with small number of flavors, $`n_fn_c`$, and go back to the equations for $`\lambda _i`$ $$\frac{1}{\mu }\left(\lambda _i\frac{1}{n_c}\underset{j}{}\lambda _j\right)+m_i\frac{1}{n_cn_f}\frac{\mathrm{\Lambda }_1^{(3n_cn_f)/(n_cn_f)}}{(_j\lambda _j)^{1/(n_cn_f)}}\lambda _i^1=0.$$ (5.1) following from $$W=\frac{1}{2\mu }\left[\text{Tr}M^2\frac{1}{n_c}(\text{Tr}M)^2\right]+\text{Tr}(Mm)+\frac{\mathrm{\Lambda }_1^{(3n_cn_f)/(n_cn_f)}}{(detM)^{1/(n_cn_f)}},$$ (5.2) where $`M_i^j\stackrel{~}{Q}_i^aQ_a^j`$; $`M_i^j=diag(\lambda _1,\mathrm{},\lambda _{n_f})`$. The scale of the $`N=1`$ SQCD $$\mathrm{\Lambda }_1\mu ^{\frac{n_c}{3n_cn_f}}\mathrm{\Lambda }^{\frac{2n_cn_f}{3n_cn_f}}$$ (5.3) must be kept fixed in the $`\mu \mathrm{}`$ limit, to recover the standard $`N=1`$ SQCD. In the large $`\mu `$ limit, some of the $`\lambda _i`$’s of Eq.(5.1) are of the order of $`\mu `$, while others are much smaller. The solutions can thus be classified according to the number $`r`$ of the $`\lambda _i`$’s which are of the order of $`\mu `$. The large $`\lambda _i`$’s (say $`\lambda _1,\lambda _2,\mathrm{},\lambda _r`$) satisfy (setting $`\mathrm{\Lambda }_1=1`$ from now on) $$\lambda _i\frac{1}{n_c}\underset{k=1}{\overset{r}{}}\lambda _k\mu m_i0;$$ (5.4) which nicely corresponds to Eq. (3.19). The justification for dropping the last term of Eq. (5.1) will be shortly given. The smaller eigenvalues $`\lambda _p`$ can be found as follows: substituting the approximate solutions for the large $`\lambda _i`$’s (see Eq. (3.21)) $$\lambda _i\mu m_i+\frac{1}{n_cr}\mu \underset{k=1}{\overset{r}{}}m_k=(m_i+c)\mu $$ (5.5) into Eq. (5.1) with $`p=r+1,r+2,\mathrm{},n_f`$, one finds $$(c+m_p)\lambda _p=\frac{1}{n_cn_f}\frac{1}{(_{i=1}^r\lambda _i_{q=r+1}^{n_f}\lambda _q)^{1/n_cn_f}}.$$ (5.6) This can be further put in the form, $$(\lambda _p)^{n_cn_f}\underset{q}{}\lambda _q=\frac{A_p}{_{i=1}^r\lambda _i}=\frac{B_p}{\mu ^r},$$ (5.7) where $`A_p`$ and $`B_p`$ are some finite constants depending on the masses. By taking the product over all $`p=r+1,r+2,\mathrm{},n_f`$, one gets $$(\underset{q}{}\lambda _q)^{n_cr}=\frac{\text{const.}}{\mu ^{r(n_fr)}},$$ (5.8) therefore $$\underset{q}{}\lambda _q=\frac{\text{const.}}{\mu ^{r(n_fr)/(n_cr)}}\mathrm{exp}2\pi ik/(n_cr),(k=1,2,\mathrm{},n_cr).$$ (5.9) Once $`_q\lambda _q`$ is determined, each of the small eigenvalues can be found uniquely from Eq. (5.6), so that each choice of $`r`$ large $`\lambda _i`$’s yields $`n_cr`$ solutions. It is easy to see from Eq. (5.9) that the last term of Eq. (5.1) behaves as $$\mu ^{n_c/(n_cr)}$$ (5.10) and thus is indeed negligible as compared to the terms kept in Eq. (5.4), as long as $`r<n_c`$. The total number of the vacua at large $`\mu `$ is thus $$𝒩=\underset{r=0}{\overset{n_f}{}}(n_cr)\left(\begin{array}{c}n_f\\ r\end{array}\right)$$ (5.11) which coincides with the number of the classical vacua, Eq. (3.23). With $`n_f>n_c+1`$, a naïve use of Eq.(4.38) in the limit $`\mu \mathrm{}`$ would lead to no supersymmetric vacua. The correct vacua can be found by taking into account the nontrivial dual gauge dynamics and consequently considering the effective action Eq.(4.48): the analysis of the decoupling limit is then similar to the $`n_fn_c`$ cases discussed above. A subtle new point however is that now the number of “large” eigenvalues $`r`$, can a priori exceed $`n_c`$. Actually, however, for these values of $`r`$ the last term of Eq. (5.1) becomes dominant and invalidates the solution (see Eq.(5.10)): the sum over $`r`$ must be truncated at $`r=n_c1`$.<sup>11</sup><sup>11</sup>11From Eq.(5.4) it can be seen that the case $`r=n_c`$ is also excluded. One thus finds for the number of vacua $$𝒩=\underset{r=0}{\overset{n_c1}{}}(n_cr)\left(\begin{array}{c}n_f\\ r\end{array}\right)\text{for}n_f>n_c+1,$$ (5.12) which is indeed the correct vacuum multiplicity in this case (Eq. (3.23)). In the $`\mu \mathrm{}`$ limit, all solutions except for those with $`r=0`$ have some VEVS running away to infinity. They do not belong to the space of vacua of the $`N=1`$ supersymmetric QCD. Only the $`r=0`$ solutions are characterized by finite VEVS, $$m_iQ_i\stackrel{~}{Q}_i=\text{indep. of }i=\mathrm{\Lambda }_1^{\frac{3n_cn_f}{n_c}}(\underset{j=1}{\overset{n_f}{}}m_j^{1/n_c})e^{2\pi ik/n_c},k=1,2,\mathrm{},n_c:$$ (5.13) they are indeed the well-known $`n_c`$ vacua of $`N=1`$ SQCD . ### 5.2 $`USp(2n_c)`$ For small numbers of flavors ($`n_fn_c`$) the equations for $`\lambda _i`$’s are (Eq. (4.54)): $$\frac{1}{2\mu }\lambda _im_i\frac{1}{\lambda _i}\frac{1}{(_j\lambda _j)^{1/(n_c+1n_f)}}=0.$$ (5.14) The large $`\lambda _i`$’s (say $`\lambda _1,\lambda _2,\mathrm{},\lambda _r`$) satisfy: $$\lambda _i2m_i\mu i=1,2,\mathrm{},r,$$ (5.15) where the last term of Eq. (5.14) is negligible, as will be shown shortly. The $`n_fr`$ smaller $`\lambda _p`$’s are found by substituting Eq. (5.15) into Eq. (5.14): $$m_p\lambda _p=\frac{1}{(_i\lambda _i_q\lambda _q)^{1/(n_c+1n_f)}},$$ (5.16) where $`i=1,2,\mathrm{}r`$ runs over the large $`\lambda `$’s and $`p,q=r+1,r+2,\mathrm{}n_f`$ refer to the smaller ones. This can be further rewritten as: $$(\underset{p}{}\lambda _p)^{n_c+1r}=\frac{\text{const.}}{\mu ^{r(n_fr)}},$$ (5.17) and thus $$\underset{p}{}\lambda _p=\frac{\text{const.}}{\mu ^{r(n_fr)/(n_c+1r)}}\mathrm{exp}2\pi ik/(n_c+1r)\text{with }k=1,2,\mathrm{},n_c+1r.$$ (5.18) One can see that each choice of $`r`$ large $`\lambda _i`$’s yields $`n_c+1r`$ solutions. It is easy to show that the last term of Eq. (5.14) is indeed negligible, behaving as $$\mu ^{(n_c+1)/(n_c+1r)}$$ (5.19) as long as $`rn_c`$ (see below). Summing over $`r`$ one finds for the total number of the vacua $$𝒩=\underset{r=0}{\overset{n_f}{}}(n_c+1r)\left(\begin{array}{c}n_f\\ r\end{array}\right)$$ (5.20) which agrees (for $`n_fn_c`$) with the number of the classical vacua, Eq. (3.43); the above solutions therefore exhaust all possible vacua of the theory. The analysis for the cases with larger number of flavor ($`n_f>n_c+1`$) is quite similar to the one made above for smaller values of $`n_f`$. The only difference is that for $`r>n_c`$ it is no longer correct to neglect the last term of Eq. (5.14), as can be seen from Eq. (5.19), hence the sum over $`r`$ must stop at $`r=n_c`$. The case $`r=n_c+1`$ might look subtle, but it is clear from Eq. (5.17) that no solution exists in this case either. One thus ends up with the number of vacua, $$𝒩=\underset{r=0}{\overset{n_c}{}}(n_cr+1)\left(\begin{array}{c}n_f\\ r\end{array}\right),$$ (5.21) which is the correct result (see Eq. (3.43)). Again, in the strict $`\mu =\mathrm{}`$ limit, only the vacua with finite VEVS must be retained. They are the solutions correspondintg to $`r=0`$ above: we find precisely $`n_c1`$ solutions of the $`N=1`$ $`USp(2n_c1)`$ theory without the adjoint matter fields. ## 6 Microscopic Picture of Dynamical Symmetry Breaking In this and in three subsequent sections we seek for a microscopic understanding of the mechanism of dynamical flavor symmetry breaking as well as of confinement itself, by studying these theories at small $`\mu `$ and small $`m_i`$. It is necessary to analyze the $`N=2`$ vacua on the Coulomb branch which survive the $`\mu 0`$ perturbation. The auxiliary genus $`g=n_c1`$ (or $`n_c`$) curves for $`SU(n_c)`$ ($`USp(2n_c)`$) theories are given by $$y^2=\underset{k=1}{\overset{n_c}{}}(x\varphi _k)^2+4\mathrm{\Lambda }^{2n_cn_f}\underset{j=1}{\overset{n_f}{}}(x+m_j),SU(n_c),n_f2n_c2,$$ (6.1) and $$y^2=\underset{k=1}{\overset{n_c}{}}(x\varphi _k)^2+4\mathrm{\Lambda }\underset{j=1}{\overset{n_f}{}}\left(x+m_j+\frac{\mathrm{\Lambda }}{n_c}\right),SU(n_c),n_f=2n_c1,$$ (6.2) with $`\varphi _k`$ subject to the constraint $`_{k=1}^{n_c}\varphi _k=0`$, and $$xy^2=\left[x\underset{a=1}{\overset{n_c}{}}(x\varphi _a^2)^2+2\mathrm{\Lambda }^{2n_c+2n_f}m_1\mathrm{}m_{n_f}\right]^24\mathrm{\Lambda }^{2(2n_c+2n_f)}\underset{i=1}{\overset{n_f}{}}(x+m_i^2),USp(2n_c).$$ (6.3) The connection between these genus $`g`$ hypertori and physics is made - through the identification of various period integrals of the holomorphic differentials on the curves with $`(da_{Di}/du_j,da_i/du_j)`$, where the gauge invariant parameters $`u_j`$’s are defined by the standard relation, $$\underset{a=1}{\overset{n_c}{}}(x\varphi _a)=\underset{k=0}{\overset{n_c}{}}u_kx^{n_ck},u_0=1,u_1=0,SU(n_c);$$ (6.4) $$\underset{a=1}{\overset{n_c}{}}(x\varphi _a^2)=\underset{k=0}{\overset{n_c}{}}u_kx^{n_ck},u_0=1,USp(2n_c),$$ (6.5) and $`u_2\text{Tr}\mathrm{\Phi }^2,`$ etc. The VEVS of $`a_{Di},a_i`$, which are directly related to the physical masses of the BPS particles through the exact Seiberg-Witten mass formula $$M^{n_{mi},n_{ei},S_k}=\sqrt{2}\left\|\underset{i=1}{\overset{g}{}}(n_{mi}a_{Di}+n_{ei}a_i)+\underset{k}{}S_km_k\right\|,$$ (6.6) are constructed as integrals over the non-trivial cycles of the meromorphic differentials on the curves. See Appendix C. We require that the curve is maximally singular, i.e. $`g=n_c1`$ (or $`n_c`$ for $`USp(2n_c)`$) pairs of branch points to coincide: this determines the possible values of $`\{\varphi _a\}`$’s. These points correspond to the $`N=1`$ vacua, for the particular $`N=1`$ perturbation, Eq.(1.2). Note that as we work with generic and nonvanishing quark masses (and then consider $`m_i0`$ limit), this is an unambiguous procedure to identify all the $`N=1`$ vacua of our interest. <sup>12</sup><sup>12</sup>12There are other kinds of singularities of $`N=2`$ QMS at which, for instance, three of the branch points meet. These correspond to $`N=1`$ vacua, selected out by different types of perturbations such as $`\text{Tr}\mathrm{\Phi }^3`$, which are not considered here. In fact, near one of the singularities where dyons with quantum numbers $$(n_{m1},n_{m2},\mathrm{},n_{mg};n_{e1},n_{e2},\mathrm{},n_{eg})=(1,0,\mathrm{},0;0,\mathrm{}),\mathrm{},(0,0,\mathrm{},1;0,\mathrm{})$$ (6.7) become massless, the effective superpotential reads $$𝒲=\underset{i=1}{\overset{g}{}}\left\{\sqrt{2}a_{Di}\stackrel{~}{M}_iM_i+\underset{k=1}{\overset{n_f}{}}S_k^im_k\stackrel{~}{M}_iM_i\right\}+\mu u_2(a_D,a)$$ (6.8) where $`S_k^i`$ is the $`k`$th quark number charge of the $`i`$ th dyon . Treating $`u_i`$ as independent variables (or equivalently, $`a_{Di}`$’s), the equations of the minimum are $`{\displaystyle \frac{\mu }{\sqrt{2}}}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{g}{}}}{\displaystyle \frac{a_{Di}}{u_2}}\stackrel{~}{M}_iM_i;0={\displaystyle \underset{i=1}{\overset{g}{}}}{\displaystyle \frac{a_{Di}}{u_j}}\stackrel{~}{M}_iM_i,j=3,4,\mathrm{},g+1;`$ (6.9) $`\left(\sqrt{2}a_{D1}+{\displaystyle \underset{k=1}{\overset{n_f}{}}}S_k^1m_k\right)\stackrel{~}{M}_1`$ $`=`$ $`\left(\sqrt{2}a_{D1}+{\displaystyle \underset{k=1}{\overset{n_f}{}}}S_k^1m_k\right)M_1=0;`$ $`\mathrm{}\mathrm{}`$ $`=`$ $`\mathrm{}\mathrm{}`$ $`\left(\sqrt{2}a_{Dg}+{\displaystyle \underset{k=1}{\overset{n_f}{}}}S_k^gm_k\right)\stackrel{~}{M}_g`$ $`=`$ $`\left(\sqrt{2}a_{Dg}+{\displaystyle \underset{k=1}{\overset{n_f}{}}}S_k^gm_k\right)M_g=0.`$ (6.10) The $`D`$-term constraint gives $`\|\stackrel{~}{M}_i\|=\|M_i\|`$. For generic hyperquark masses, Eqs.(6.9) require that $$\stackrel{~}{M}_i,M_i\sqrt{\mu \mathrm{\Lambda }}0,i,$$ (6.11) since $`\frac{a_{Di}}{u_j}`$ and $`\frac{a_{Di}}{u_j}`$ obey no special relations. This means that $$\sqrt{2}a_{Di}+\underset{k=1}{\overset{n_f}{}}S_k^im_k=0$$ (6.12) for all $`i`$: i.e., all $`g`$ monopoles are massless simultaneously. Condensation of each type of monopole, $`\stackrel{~}{M}_i0`$; $`M_i0`$, corresponding to the maximal Abelian subgroup of $`SU(n_c)`$ or of $`USp(2n_c)`$, amounts to confinement à la ’t Hooft-Mandelstam-Nambu. Actually, physical picture in the $`m_i0`$ limit of these theories is more subtle and is far more interesting, as is discussed especially in Section 8. It is in general difficult to determine explicitly the configurations $`\{\varphi _a\}`$ which satisfy the $`N=1`$ criterion mentioned above, although in some special cases (with $`m_i=0`$) they can be found explicitly. We approach the problem by first setting $`m_i=0,`$ $`i`$, and by perturbing the solutions for $`\{\varphi _a\}`$ by considering the effects of $`m_i`$ to first nontrivial orders. It turns out that the $`N=1`$ vacua of the $`SU(n_c)`$ and $`USp(2n_c)`$ gauge theories can all be generated from the various classes of superconformal theories with $`m_i=\mu =0`$, by perturbing them with masses $`m_i`$ (as well as with $`\mu `$). Some qualifying remarks are in order. In the case of $`SU(n_c)`$ theories we find that the first group of $`N=1`$ vacua surviving the the adjoint mass perturbation ($`\mu \text{Tr}\mathrm{\Phi }^2`$) and leading to finite VEVS, arise from the class 1 ($`r<n_f/2`$), class 3 ($`r=n_f/2`$ with $`n_cn_f/2`$ odd) and class 4 ($`r=n_f/2`$ with $`n_cn_f/2`$ even) CFT, according to the classification of Eguchi et. al . The second group of vacua, present only for $`n_f>n_c`$, on the other hand, arise from the so-called baryonic branch root (see Eq.(6.18)): these latter CFT can also be regarded as special case of class 1 theories. In these vacua the flavor symmetry is unbroken in the $`m_i0`$ limit. It is important that the number of the vacua and qualitative results about symmetry breaking in each of them, can be obtained this way even when the unperturbed solution for $`\{\varphi _a\}`$ is not known explicitly. ### 6.1 Superconformal Points and $`N=1`$ Vacua for $`SU(n_c)`$ It will be seen that the first group of vacua (with multiplicity $`𝒩_1`$) are associated to the points, $$\text{diag}\varphi =(\underset{r}{\underset{}{0,0,\mathrm{},0}},\varphi _{r+1}^{(0)},\mathrm{},\varphi _{n_c}^{(0)}),\underset{a=r}{\overset{n_c}{}}\varphi _a^{(0)}=0,$$ (6.13) in Eq.(6.1), with $`\varphi _a^{(0)}`$’s chosen such that the nonzero $`2(n_cr1)`$ branch points are paired <sup>13</sup><sup>13</sup>13Actually there is no vacuum of this type Eq.(6.14) with $`r=\stackrel{~}{n}_c`$. This will be shown in Appendix D. . The curves are of the form, $$y^2x^{2r}(x\beta _{01})^2\mathrm{}(x\beta _{0,n_cr1})^2(x\gamma _0)(x\kappa _0),r=0,1,2,\mathrm{},[n_f/2].$$ (6.14) For $`r<n_f/2`$ they correspond to the so-called class 1 superconformal theories . The special cases with $`r=n_f/2`$, with $`n_cn_f/2`$ odd, and $`r=n_f/2`$, $`n_cn_f/2`$ even, may be interpreted as belonging to class 3 and class 4 , respectively. In fact, in these special cases the explicit configuration of $`\varphi _a`$’s can be found by using the method of . The curve for bare quark masses with $`r=n_f/2`$ vanishing $`\varphi _a`$’s is: $$y^2=x^{n_f}\left[\underset{k=1}{\overset{n_cn_f/2}{}}(x\varphi _k)^24\mathrm{\Lambda }^{2n_cn_f}\right].$$ (6.15) We identify the first term in the square bracket with $`(2\mathrm{\Lambda }^{n_cn_f/2}T_{n_cn_f/2}(x/2\mathrm{\Lambda }))^2`$, where $`T_N(x)`$ is the Chebyshev polynomial of order $`N`$, $`T_N(x)=\mathrm{cos}(N\mathrm{cos}^1x),`$ implying that $$\varphi _k=2\mathrm{\Lambda }\mathrm{cos}\pi (2k1)/2(n_cn_f/2),k=1,\mathrm{},n_cn_f/2.$$ (6.16) The two terms in the square bracket combine as $`y^2`$ $`=`$ $`x^{n_f}\left[\left(2\mathrm{\Lambda }^{n_cn_f/2}T_{n_cn_f/2}\left({\displaystyle \frac{x}{2\mathrm{\Lambda }}}\right)\right)^24\mathrm{\Lambda }^{2n_cn_f}\right]`$ (6.17) $`=`$ $`4x^{n_f}\mathrm{\Lambda }^{2n_cn_f}\mathrm{sin}^2\left[\left(n_c{\displaystyle \frac{n_f}{2}}\right)\mathrm{arccos}{\displaystyle \frac{x}{2\mathrm{\Lambda }}}\right].`$ There are $`n_f`$-th order zero at $`x=0`$, and double zeros at $`x=2\mathrm{\Lambda }\mathrm{cos}\pi k/(n_cn_f/2)`$ for $`k=1,\mathrm{},n_cn_f/21`$, and single zeros at $`x=\pm 2\mathrm{\Lambda }`$. Note that for $`r=n_f/2`$, $`n_cn_f/2`$ even, the zero at $`x=0`$ is actually of order $`n_f+1`$. Since these adjoint VEVS leading to Eq.(6.14) or Eq(6.17) break the discrete symmetry spontaneously, they appear in $`2n_cn_f`$ copies.<sup>14</sup><sup>14</sup>14There is an exception to this. In the case of $`r=n_f/2`$ with $`n_f`$ even, the explicit configuration of $`\varphi _a`$’s (Section 9.1) shows that the vacuum respects $`Z_2`$ subgroup of the $`Z_{2n_cn_f}`$ symmetry, showing that it appears in $`n_cn_f/2`$ copies rather than $`2n_cn_f.`$ This fact is crucial in the vacuum counting below Eq.(6.21). When (generic) quark masses are turned on, these vacua split into $`{}_{n_f}{}^{}C_{r}^{}`$-plet of single vacua. The vacua Eq.(6.14), Eq(6.17), correspond to what is called “nonbaryonic” branch roots in . The second group of vacua will be related to the (trivial) superconformal theory $$y^2=x^{2\stackrel{~}{n}_c}(x^{n_c\stackrel{~}{n}_c}\mathrm{\Lambda }^2)^2,\stackrel{~}{n}_c=n_fn_c,$$ (6.18) corresponding to the adjoint configuration $$\text{diag}\varphi =(\underset{\stackrel{~}{n}_c}{\underset{}{0,0,\mathrm{},0}},\mathrm{\Lambda }\omega ^2,\mathrm{},\mathrm{\Lambda }\omega ^{2(n_c\stackrel{~}{n}_c}))$$ (6.19) with $`\omega =e^{\pi i/(n_c\stackrel{~}{n}_c)}`$. This is the “baryonic” root of . To justify these statements we must solve the following problem of purely mathematical nature. Suppose that a configuration $`\{\varphi _a\}=\{\varphi _a^0\}`$ has been found such that the curve of the $`m_i=0`$ theory reduces to one of the forms, Eq.(6.14) or Eq.(6.18). Now we add small generic bare hyperquark masses $`m_i`$; we want to determine the shifts of $`\{\varphi _a\}`$, $`\{\varphi _a\}=\{\varphi _a^0+\delta \varphi _a\}`$, with constraints $`_a\delta \varphi _a=0`$ and $`\delta \varphi _a0`$ as $`m_i0`$, such that the massive curve Eq.(6.1) or Eq.(6.2) is maximally singular (with $`g=n_c1`$ double branch points). How many such solutions are there? It turns out that this problem must be treated separately for several different cases: the result of this mass perturbation theory, which will be developed in the last section of this paper (Section 9), can be summarized as follows. For small number of the flavors ($`n_fn_c`$) the total number of the $`N=1`$ vacua generated by the mass perturbation from various conformal invariant points Eq.(6.14), Eq.(6.13), is: $$𝒩_1=(2n_cn_f)\underset{r=0}{\overset{(n_f1)/2}{}}{}_{n_f}{}^{}C_{r}^{}=(2n_cn_f)\mathrm{\hspace{0.17em}2}^{n_f1},(n_f=\text{odd})$$ (6.20) $$𝒩_1=(2n_cn_f)\underset{r=0}{\overset{n_f/21}{}}{}_{n_f}{}^{}C_{r}^{}+\frac{2n_cn_f}{2}{}_{n_f}{}^{}C_{n_f/2}^{}=(2n_cn_f)\mathrm{\hspace{0.17em}2}^{n_f1},(n_f=\text{even}),$$ (6.21) which exhausts $`𝒩`$, Eq.(3.24). In Eq.(6.21) we have taken into account the fact that for even $`n_f`$, the vacua with $`r=n_f/2`$ do not transform under $`Z_{2n_cn_f}`$ but only under $`Z_{n_cn_f/2}.`$ For larger $`n_f`$ ($`n_fn_c+1`$), there are also $$𝒩_2=\underset{r=0}{\overset{\stackrel{~}{n}_c1}{}}(\stackrel{~}{n}_cr){}_{n_f}{}^{}C_{r}^{}.$$ (6.22) vacua coming from the “baryonic root”, Eq.(6.18), Eq.(6.19). The total $`𝒩_1+𝒩_2`$ correctly matches the known total number of the vacua. The arithmetics is the same as in Eq.(4.49), Eq.(4.50), and will not be repeated. Actually, there is an interesting subtlety in this vacuum counting. In the first group of vacua Eq.(6.20), Eq.(6.21), the term $`r=\stackrel{~}{n}_c=n_fn_c`$ must be dropped. The vacuum Eq.(6.14) turns out to be nonexistent for $`r=\stackrel{~}{n}_c=n_fn_c`$, see Appendix D. At the same time, however, the mass perturbation around the baryonic branch root Eq.(6.18), Eq.(6.19), gives (see Sections 8, 9) $$𝒩_2+(n_c\stackrel{~}{n}_c){}_{n_f}{}^{}C_{\stackrel{~}{n}_c}^{}$$ (6.23) vacua, the second term of which compensates precisely the missing term in the sum in $`𝒩_1`$. ### 6.2 Superconformal Points and $`N=1`$ Vacua of $`USp(2n_c)`$ Theory The general curve of $`USp(2n_c)`$ theories is given by $$xy^2=\left[x\underset{a=1}{\overset{n_c}{}}(x\varphi _a^2)+2\mathrm{\Lambda }^{2n_c+2n_f}m_1\mathrm{}m_{n_f}\right]^24\mathrm{\Lambda }^{2(2n_c+2n_f)}\underset{i=1}{\overset{n_f}{}}(x+m_i^2).$$ (6.24) When $`2n_c+2n_f=0`$, the theory is superconformal and $`2\mathrm{\Lambda }^{2n_c+2n_f}`$ in the above expression is replaced by $$g(\tau )=\frac{\vartheta _2^4}{\vartheta _3^4+\vartheta _4^4}.$$ (6.25) It turns out that the first group of $`N=1`$ vacua can be generated from the CFT described by Chebyshev solutions for $`m_i=0`$ theory. When $`n_f`$ is odd, choose $`(n_f1)/2`$ eigenvalues $`\varphi _1=\mathrm{}=\varphi _{(n_f1)/2}`$ vanishing. Then the curve becomes $$y^2=x^{n_f1}\left[x\underset{k=1}{\overset{n_c(n_f1)/2}{}}(x\varphi _k^2)^24\mathrm{\Lambda }^{2(2n_c+2n_f)}\right].$$ (6.26) We identify the first term in the square bracket with $`(2\mathrm{\Lambda }^{2n_c+2n_f}T_{2n_c+2n_f}(\sqrt{x}/2\mathrm{\Lambda }))^2`$, where $`\varphi _k^2=4\mathrm{\Lambda }^2\mathrm{cos}^2\pi (2k1)/2(2n_c+2n_f)`$ for $`k=1,\mathrm{},n_c(n_f1)/2`$. Then the two terms in the square bracket combine as $`y^2`$ $`=`$ $`x^{n_f1}\mathrm{\Lambda }^{2(2n_c+2n_f)}\left[4T_{2n_c+2n_f}^2\left({\displaystyle \frac{\sqrt{x}}{2\mathrm{\Lambda }}}\right)4\right]`$ (6.27) $`=`$ $`4x^{n_f1}\mathrm{\Lambda }^{2(2n_c+2n_f)}\mathrm{sin}^2\left[(2n_c+2n_f)\mathrm{arccos}{\displaystyle \frac{\sqrt{x}}{2\mathrm{\Lambda }}}\right].`$ There are $`n_f1`$ zeros at $`x=0`$, and double zeros at $`x=4\mathrm{\Lambda }^2\mathrm{cos}\pi k/(2n_c+2n_f)`$ for $`k=1,\mathrm{},n_c(n_f1)/2`$, and a single zero at $`x=4\mathrm{\Lambda }^2`$. When $`n_f`$ is even, choose $`n_f/21`$ eigenvalues $`\varphi _1=\mathrm{}=\varphi _{n_f/21}`$ vanishing. Then the curve becomes $$y^2=x^{n_f1}\left[\underset{k=1}{\overset{n_c+1n_f/2}{}}(x\varphi _k^2)^24\mathrm{\Lambda }^{2(2n_c+2n_f)}\right].$$ (6.28) We identify the first term in the square bracket with $`(2\mathrm{\Lambda }^{2n_c+2n_f}T_{2n_c+2n_f}(\sqrt{x}/2\mathrm{\Lambda }))^2`$, where $`\varphi _a^2=4\mathrm{\Lambda }^2\mathrm{cos}^2\pi (2k1)/2(2n_c+2n_f)`$ for $`k=1,\mathrm{},n_c+1n_f/2`$. Then the two terms in the square bracket combine as $`y^2`$ $`=`$ $`x^{n_f1}\mathrm{\Lambda }^{2(2n_c+2n_f)}\left[4T_{2n_c+2n_f}^2\left({\displaystyle \frac{\sqrt{x}}{2\mathrm{\Lambda }}}\right)4\right]`$ (6.29) $`=4x^{n_f1}\mathrm{\Lambda }^{2(2n_c+2n_f)}\mathrm{sin}^2\left[(2n_c+2n_f)\mathrm{arccos}{\displaystyle \frac{\sqrt{x}}{2\mathrm{\Lambda }}}\right].`$ Since the $`\mathrm{sin}^2`$ factor gives a single zero at $`x=0`$, there are $`n_f`$ zeros at $`x=0`$, and double zeros at $`x=4\mathrm{\Lambda }^2\mathrm{cos}\pi k/(2n_c+2n_f)`$ for $`k=1,\mathrm{},n_cn_f/2`$, and a single zero at $`x=4\mathrm{\Lambda }^2`$. In the absence of quark masses, the theory is invariant under $`Z_{2n_c+2n_f}`$ symmetry: $`xe^{2\pi i/(2n_c+2n_f)}x`$, $`\varphi _a^2e^{2\pi i/(2n_c+2n_f)}\varphi _a^2`$. Therefore the Chebyshev solutions discussed here appear $`2n_c+2n_f`$ times. $`USp(2n_c)`$ theories also have special Higgs branch roots similar to the baryonic roots of the $`SU(n_c)`$ theories. This baryonic-like root is obtained in the $`m_i=0`$ limit, by setting $`\varphi _1,\mathrm{},\varphi _{n_c\stackrel{~}{n}_c}0`$, $`\varphi _{n_c\stackrel{~}{n}_c+1}=\mathrm{}=\varphi _{n_c}=0`$. Here and below, $`\stackrel{~}{n}_c=n_fn_c2`$. The curve Eq. (6.24) becomes $$y^2=x^{2\stackrel{~}{n}_c+1}\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}(x\mathrm{\Phi }_k)^24\mathrm{\Lambda }^{4n_c+42n_f}x^{n_f1}.$$ (6.30) We take $$(\mathrm{\Phi }_1^2,\mathrm{},\mathrm{\Phi }_k^2,\mathrm{},\mathrm{\Phi }_{n_c\stackrel{~}{n}_c}^2)=\mathrm{\Lambda }^2(\omega ,\mathrm{},\omega ^{2k1},\mathrm{},\omega ^{2(n_c\stackrel{~}{n}_c)1}),$$ (6.31) where $`\omega =e^{\pi i/(n_c\stackrel{~}{n}_c)}`$. Note that our $`\omega `$ is the square root of $`\omega `$ in because of later convenience. Then the product $`_{k=1}^{n_c\stackrel{~}{n}_c}(x\mathrm{\Phi }_k)`$ can be rewritten as $`x^{n_c\stackrel{~}{n}_c}+\mathrm{\Lambda }^{2(n_c\stackrel{~}{n}_c)}`$, and the curve becomes $`y^2`$ $`=`$ $`x^{2\stackrel{~}{n}_c+1}\left[(x^{n_c\stackrel{~}{n}_c}+\mathrm{\Lambda }^{2(n_c\stackrel{~}{n}_c)})^24\mathrm{\Lambda }^{4n_c+42n_f}x^{n_c\stackrel{~}{n}_c}\right]`$ (6.32) $`=`$ $`x^{2\stackrel{~}{n}_c+1}(x^{n_c\stackrel{~}{n}_c}\mathrm{\Lambda }^{2(n_c\stackrel{~}{n}_c)})^2.`$ The double zeros of the factor in the parenthesis are at $$x=\mathrm{\Lambda }^2\omega ^2,\mathrm{\Lambda }^2\omega ^4,\mathrm{},\mathrm{\Lambda }^2\omega ^{2k},\mathrm{},\mathrm{\Lambda }^2\omega ^{2(n_c\stackrel{~}{n}_c)}=\mathrm{\Lambda }^2.$$ (6.33) When the quark masses are turned on, these points split. We require again that the shift of $`\varphi `$’s be such that the full curve remains maximally sigular (with maximal possible number of double roots). The result of mass perturbation analysis, given in Section 9, can be summarized as follows. There are two groups of $`N=1`$ vacua predicted by the Seiberg-Witten curve in $`USp(2n_c)`$ theories. The Chebyshev point Eq.(6.27), Eq.(6.28), spawns $`𝒩_1=(2n_c+2n_f)2^{n_f1}`$ vacua upon mass perturbation, while the special point Eq.(6.30), Eq.(6.31), splits into $`𝒩_2=_{r=0}^{\stackrel{~}{n}_c}(\stackrel{~}{n}_cr+1){}_{n_f}{}^{}C_{r}^{}`$ vacua. Their sum coincides with the total number of $`N=1`$ vacua found from the semiclassical as well as from large $`\mu `$ analyses. ## 7 Numerical Study of $`N=1`$ Vacua in $`SU(3)`$ and $`USp(4)`$ Theories As a way of checking these results and of illustrating some of their features, we have performed a study of the rank $`2`$ theories, by numerically determing the points in QMS where the curves $`y^2=G(x)`$ of Eq.(6.1)-Eq.(6.3) become maximally singular. This has been done by solving the equation $$(G(x),\frac{dG(x)}{dx})=0,$$ (7.1) where $``$ stands for the resultant, using Mathematica. ### 7.1 $`SU(3)`$ theory with $`n_f=15`$ i) In the case with $`n_f=1`$, one expects from Eq. (3.23) $`𝒩=5`$. From the known curve, $$y^2=(x^3uxv)^2\mathrm{\Lambda }_{1}^{}{}_{}{}^{5}(x+m),$$ (7.2) one finds indeed five vacua, related by an approximate $`Z_5`$ symmetry. See Fig. 1 for $`\mathrm{\Lambda }_1=2`$ and $`m=1/64`$. ii) In the case with $`n_f=2`$ one expects eight vacua, related by an approximate $`Z_4`$ symmetry which transform $`u`$ and $`v`$ as $`uu`$ and $`v\mathrm{exp}(\frac{3}{2}i\pi )v`$. One finds indeed eight singularities, grouped into two approximate doublets, and four singlets for generic small masses. (see Fig. 2). iii) For $`n_f=3`$, one gets $`𝒩=12`$, with $`Z_3`$ symmetry. The numerical analysis with the curve: $$y^2=(x^3uxv)^2\mathrm{\Lambda }_{3}^{}{}_{}{}^{3}(x+m_1)(x+m_2)(x+m_3).$$ (7.3) shows that there are indeed twelve vacua satisfying the criterion of the two mutually local dyons becoming massless, and they are found in roughly three groups of triplets and three singlets of singularities. In the equal mass limit each of the three triplets coalesce to a point in the $`(u,v)`$ space, showing that the massless monopoles there are in the representation $`\underset{¯}{3}`$ of $`SU(3)`$, while those at the other vacua are singlets (see Fig. 3), in complete agreement with the analysis of previous sections. iv) For $`n_f=4`$, Eq. (3.23) gives $`𝒩=17`$ vacua. It is reassuring that one indeed finds from Eq.(7.1) seventeen vacua for generic and unequal masses <sup>15</sup><sup>15</sup>15As a further check, we verified the number of the vacua by using another parametrization of the curve given by Minahan et al. .. At small masses these vacua are grouped into an approximate sextet, two quartets and three singlets, suggesting the assignements of rank $`2`$, $`1`$ and $`0`$ antisymmetric representation of $`SU(4)`$ global flavor group. The number of the vacua (six) in the limit of equal masses is consistent with this assignment (see Fig. 4) In a quartet vacuum, the condensation of the monopoles breaks the $`SU(4)`$ symmetry to $`U(3)`$, while in a singlet vacua the flavor symmetry remains unbroken. Something very interesting happens in the sextet vacua. Namely, in the equal mass or massless limit, we find that four branch points in the $`x`$ plane coalesce, suggesting conformal invariant vacuum. In fact, this result was to be expected, since these “sextet” vacua are examples of the particular class (“class 3” ) of nontrivial conformal invariant theories studied in , where $`2r`$ branch points in the $`x`$ plane coalesce, where $`r=n_f/2.`$ This is known to occurs for $`SU(3)`$ theories with $`n_f=4`$ at special values of $`u`$ and $`v`$ when the quark masses are all equal ($`m`$): their precise position is $`U=3M^2,`$ and $`V=2M^3`$, where $`u=(7+24M+36M^2)/12`$ and $`v=(11+45M+54M^2+54M^3)/27`$ and $`M=m1/3,`$ and the curve becomes $$y^2=(x+m)^4(x+12m)(x12m).$$ (7.4) The values of $`u,v`$ found by us numerically from the criterion of $`N=1`$ vacua precisely match these values, showing that this particular conformal vacuum survives the $`N=1`$ perturbation. In order to determine which particles are actually present, it is necessary to study the monodromy transformation properties of $`(a_{D1},a_{D2},a_1,a_2)`$ around this particular singularity. In Appendix E we present such an analysis (the analysis is actually done for all seventeen vacua of $`SU(3)`$, $`n_f=4`$ model). Our result shows that the massless particles (in the $`\mu =0`$ limit) present at this singularity have quantum numbers $$(n_{m1},n_{m2},n_{e1},n_{e2})=(0,1,0,1),(1,2,2,0),(1,1,0,1).$$ (7.5) As these particles are relatively nonlocal, such a vacuum is conformal invariant . This is an example of a very general phenomenon, already discussed in the previous section. v) Finally, for $`n_f=5`$ we verified (for large masses) the presence of twenty–three quantum vacua, in accordance with Eq. (3.23). From the curve: $$y^2=\underset{a=1}{\overset{n_c}{}}(x\varphi _a)^24\mathrm{\Lambda }_5\underset{i=1}{\overset{5}{}}\left(x+m_i+\frac{\mathrm{\Lambda }_5}{n_c}\right)=(x^3uxv)^24\mathrm{\Lambda }_5\underset{i=1}{\overset{5}{}}\left(x+m_i+\frac{\mathrm{\Lambda }_5}{n_c}\right),$$ (7.6) in the small mass limit, the grouping of the singularities found is compatible with the assignment into a decuplet, two quintets and three singlets of (approximate) global $`SO(5)`$ symmetry. The number of vacua (six) in the equal mass case, is in agreement with this structure. In particular, the decuplet vacua corresponds to the curve, $$y^2=\left(x+m+\frac{\mathrm{\Lambda }_5}{n_c}\right)^4\left[\left(x\frac{2\mathrm{\Lambda }_5}{n_c}2m\right)^24\mathrm{\Lambda }_5\left(x+m+\frac{\mathrm{\Lambda }_5}{n_c}\right)\right],$$ (7.7) and corresponds to the class 1 (trivial) conformal field theory . Furthermore, since $`n_f>n_c+1`$ in this case, the theory belongs to the “large $`n_f`$” class of Sec. 4.3 we expect $`𝒩_2=7`$ (see Eq. (4.47)) of vacua to show particular properties. We find that indeed seven of the vacua (in the equal mass limit, a quintet and two singlets) approache the form $$y^2=x^4(x\mathrm{\Lambda })^2$$ (7.8) in the $`m_i0`$ limit, after a shift in $`x`$. This is exactly the structure of the singularities at the root of the baryonic branch. ### 7.2 $`USp(2n_c)`$ with $`2n_f`$ Flavors i) For $`USp(4)`$ with $`n_f=1`$, we do find five vacua, as shown in Fig. 5, consistently with the approximate $`Z_5`$ symmetry. ii) For $`n_f=2`$ with the same gauge group, eight quantum vacua are found to group into four doublets, showing that monopoles appear in the fundamental representation $`(2,1)`$ and $`(1,2)`$, of the flavor group $`SO(2n_f)SU(2)\times SU(2)`$. Their condensation leaves a $`U(2)`$ subgroup invariant, in accordance with the naïve expectation. In the equal mass limit one of the $`SU(2)`$ symmetries is exact. It is seen that in this case the $`SU(2)`$ symmety is broken spontaneously in four of the vacua, while in four others it is not, in perfect agreement with what was found in Sec. 4 (see Fig. 6). iii) For the $`USp(4)`$ theory with $`n_f=3`$, the formula (3.44) and the discrete symmetry suggests that monopoles form a quartet ($`12`$ vacua grouped into three quartets of nearby singularities). This is indeed the case as shown in Fig. 7. The condensation of $`\underset{¯}{4}`$ of $`SO(6)SU(4)`$ breaks the chiral symmetry to $`U(3)`$, as expected. iv) The theory with four flavors (and the same gauge group) has seventeen vacua. In the massless limit, they group in two spinors (octets) and one singlet of the gloabal symmetry $`SO(8)`$. In the case of degenerate but nonvanishing masses, the spinors $`\mathrm{𝟖}`$ of $`SO(8)`$ can split in two possible ways: $`\mathrm{𝟏}+\mathrm{𝟔}+\mathrm{𝟏}`$ of $`U(4)`$ in one case, and with $`\mathrm{𝟒}+\mathrm{𝟒}^{}`$ of $`U(4)`$ in the other. This is indeed the situation shown in Fiq. 8. This shows the correctness of our assignment and that in the massless limit the condensation of the monopole in the spinor representation of $`SO(8)`$ breaks the chiral symmetry to $`U(4)`$. v) Finally, we verified that the theory $`USp(4)`$ with five flavors has indeed twenty-three quantum vacua. ## 8 Effective Lagrangian Description of $`N=1`$ Vacua at Small $`\mu `$ A deeper insight into physics at small $`m_i`$ and $`\mu `$ can be obtained by re-examining the works of Argyres, Plesser and Seiberg and Argyres, Plesser and Shapere , who showed how the non-renormalization theorem of the hyperKähler metric on the Higgs branch could be used to show the persistence of unbroken non-abelian gauge group at the “roots” of the Higgs branches where they intersect the Coulomb branch. In fact, they found two kinds of such submanifolds, called “non-baryonic branch” (or mixed-branch) roots, and “baryonic branch” roots (these terminologies refer specifically to the $`SU(n_c)`$ theory, but the situation is similar also in $`USp(2n_c)`$ theory). The latter is present only for larger values of the flavor ($`n_f>n_c`$) while the former exists always. Below, we show how the low-energy effective action description of match our findings of Sections 3 \- 7, after correcting a few errors and clarifying some issues left unclear there. In doing so, a very clear and interesting picture of the infrared dyamics of our theories emerges, which was summarized in the Introduction. Let us discuss the $`SU(n_c)`$ theories first. ### 8.1 $`SU(n_c)`$ The non–baryonic roots are further classified into sub–branches characterized by an unbroken $`SU(r)\times U(1)^{n_cr}`$ gauge symmetry for $`r[n_f/2]`$, with $`n_f`$ flavor of massless hypermultiplets in the fundamental representation of $`SU(r)`$, as well as $`n_cr1`$ singlet “monopole” hypermultiplets having charges only in the $`U(1)^{n_cr}`$ gauge sector. Their quantum numbers are shown in Table 4 taken from . Upon turning on the $`\mu \mathrm{\Phi }^2`$ perturbation, the effective superpotential of the theory is, according to Argyres, Plesser and Seiberg , $$W_{nonbar}=\sqrt{2}\text{Tr}(q\varphi \stackrel{~}{q})+\sqrt{2}\psi _0\text{Tr}(q\stackrel{~}{q})+\sqrt{2}\underset{k=1}{\overset{n_cr1}{}}\psi _ke_k\stackrel{~}{e}_k+\mu \left(\mathrm{\Lambda }\underset{i=0}{\overset{n_cr1}{}}x_i\psi _i+\frac{1}{2}\text{Tr}\varphi ^2\right),$$ (8.1) where the last term arises from $`\mu \mathrm{\Phi }^2`$ perturbation, $`\varphi `$ referring to the $`SU(r)`$ part of the adjoint field and $`\psi _i`$ being the $`N=2`$ partner of the dual $`U(1)_i`$ gauge field; $`x_i`$ are some constants. By minimizing the potential, one finds that supersymmetric $`N=1`$ vacua exist for any $`r`$. (Actually, one assumed here that the maximum number of massless monopole–like particles $`e_k`$ exist; such vacua are called in “special points” of the non baryonic roots. Only these vacua survive the $`N=1`$ perturbation. Some flat directions remain, as expected.) We now perturb these theories at the nonbaryonic branch roots further with small hypermultiplet (quark) masses, $`m_i`$. Add the mass terms $$\mathrm{\Delta }W_{nonbar}=m_iq_i\stackrel{~}{q}_i+\underset{k=1}{\overset{n_cr1}{}}S_k^jm_je_k\stackrel{~}{e}_k,$$ (8.2) where $`S_k^j`$ represents the $`jth`$ quark number of the “monopole” $`e_k`$ . The part involving $`e_k`$, $`\stackrel{~}{e}_k`$ and $`\psi _k`$, $`k=1,2,\mathrm{}n_cr1`$ is trivial and gives $$\psi _k\{m_i\};e_k=\stackrel{~}{e}_k\sqrt{\mu \mathrm{\Lambda }}.$$ (8.3) The vacuum equations for other components are $$0=[\varphi ,\varphi ^{}];$$ (8.4) $$\nu \delta _a^b=q_a^i(q^{})_i^b(\stackrel{~}{q}^{})_a^i\stackrel{~}{q}_i^b;$$ (8.5) $$0=q_a^i(q^{})_i^a(\stackrel{~}{q}^{})_a^i\stackrel{~}{q}_i^a;$$ (8.6) $$q_a^i\stackrel{~}{q}_i^b\frac{1}{r}\delta _a^b(q_c^i\stackrel{~}{q}_i^c)+\sqrt{2}\mu \varphi _a^b=0;$$ (8.7) $$0=\sqrt{2}\varphi _a^bq_b^i+q_a^i(m_i+\sqrt{2}\psi _0);(\text{no sum over}i,a)$$ (8.8) $$0=\sqrt{2}\stackrel{~}{q}_i^b\varphi _b^a+(m_i+\sqrt{2}\psi _0)\stackrel{~}{q}_i^b(\text{no sum over}i,a).$$ (8.9) $$\sqrt{2}\text{Tr}(q\stackrel{~}{q})+\mu \mathrm{\Lambda }=0.$$ (8.10) First diagonalize the Higgs scalar by color rotations. $$\text{diag}\varphi =(\varphi _1,\varphi _2,\mathrm{}\varphi _r),\varphi _a=0.$$ (8.11) The equations Eq.(8.4)-Eq.(8.9) are formally identical to Eq.(3.5)-Eq.(3.9), with the replacements, $$Qq;,m_im_i+\sqrt{2}\psi _0;n_cr,$$ (8.12) therefore we immediately find solutions classified by an integer $`\mathrm{}`$ ($`\mathrm{}=0,1,\mathrm{},r1`$) where $$\varphi =\frac{1}{\sqrt{2}}\text{diag}(m_1\sqrt{2}\psi _0,m_2\sqrt{2}\psi _0,\mathrm{}m_{\mathrm{}}\sqrt{2}\psi _0,c,c,\mathrm{}c),\varphi _a=0;$$ (8.13) $$q_a^i=\left(\begin{array}{c}d_1\\ 0\\ \mathrm{}\\ \mathrm{}\\ 0\end{array}\right),\left(\begin{array}{c}0\\ d_2\\ 0\\ \mathrm{}\\ 0\end{array}\right),\left(\begin{array}{c}0\\ \mathrm{}\\ d_{\mathrm{}}\\ 0\\ \mathrm{}\end{array}\right);q_a^i=0,i=\mathrm{}+1,\mathrm{},n_f.$$ (8.14) $$\stackrel{~}{q}_i^a=\left(\begin{array}{c}\stackrel{~}{d}_1\\ 0\\ \mathrm{}\\ \mathrm{}\\ 0\end{array}\right),\left(\begin{array}{c}0\\ \stackrel{~}{d}_2\\ 0\\ \mathrm{}\\ 0\end{array}\right),\left(\begin{array}{c}0\\ \mathrm{}\\ \stackrel{~}{d}_{\mathrm{}}\\ 0\\ \mathrm{}\end{array}\right);\stackrel{~}{q}_i^a=0,i=\mathrm{}+1,\mathrm{},n_f,$$ (8.15) where $$c=\frac{1}{r\mathrm{}}\underset{k=1}{\overset{\mathrm{}}{}}(m_k+\sqrt{2}\psi _0)$$ (8.16) and $`d_i`$, $`\stackrel{~}{d}_i`$’s are of order of $`O(\sqrt{\mu (\psi _0+m_i)})`$. Note however that Eq.(8.10) is new: it fixes $`\psi _0`$ such that $$\frac{\sqrt{2}r\mu }{r\mathrm{}}\underset{j=1}{\overset{\mathrm{}}{}}(\psi _0+m_j)=\mu \mathrm{\Lambda }.$$ (8.17) These are the first group of $`N=1`$ solutions found in . The fact that in the limit $`m_i0`$, $$\psi _0\mathrm{\Lambda },$$ (8.18) however, shows that these solutions, involving fluctuations much larger than both $`m_i`$ and $`\mu `$, lie beyond the validity of the low-energy effective Lagrangian. They are therefore an artifact of the approximation, and must be discarded. The correct $`N=1`$ vacua are found instead by choosing $`\mathrm{}=r`$ in Eq.(8.13) and by selecting VEVS, $$\psi _0=\frac{1}{\sqrt{2}r}\underset{i}{}m_i,$$ (8.19) $$d_i\stackrel{~}{d}_i=\mu \left(m_i\frac{1}{r}\underset{j}{\overset{r}{}}m_j\right)\frac{\mu \mathrm{\Lambda }}{\sqrt{2}r};$$ (8.20) $$e_k\stackrel{~}{e}_k=\mu \mathrm{\Lambda },$$ (8.21) with no $`c`$’s in Eq.(8.13) <sup>16</sup><sup>16</sup>16 Note that an analogous solution was not possible for Eq.(3.5)-Eq.(3.9), since the quark masses are generic and the adjoint field must be traceless. . These satisfy clearly all of the vacuum equations. These vacua are smoothly related to the unperturbed ones and therefore are reliable. In the massless limit this gives $$d_i\stackrel{~}{d}_i\sqrt{\mu \mathrm{\Lambda }}.$$ (8.22) We find thus $`{}_{n_f}{}^{}C_{r}^{}(2n_cn_f)`$ vacua with the correct symmetry breaking pattern, $$U(n_f)U(r)\times U(n_fr),$$ (8.23) which is exactly what is expected from the analysis made at large $`\mu `$. The multiplicity $`{}_{n_f}{}^{}C_{r}^{}`$ arises from the choice of $`r`$ (out of $`n_f`$) quark masses used to construct the solution. For $`r=n_f/2`$, the theory at the singularity becomes a non-trivial superconformal theory. There is no description of this singularity in terms of weakly coupled local field theory. The monodromy around the singularity shows that the theory is indeed superconformal (we checked this explicitly for $`n_c=3`$ and $`n_f=4`$). Careful perturbation of the curve by the quark masses made in Section 9 shows that there are $`(n_cn_f/2){}_{n_f}{}^{}C_{n_f/2}^{}`$ vacua. The total number of $`N=1`$ vacua is then given by $`𝒩_1`$ of Eq.(6.20), Eq.(6.21), after summing over $`r`$. Actually, one must subtract the term $`r=\stackrel{~}{n}_c`$, since the nonbaryonic root with $`r=\stackrel{~}{n}_c`$ does not exist (see Appendix D). The baryonic root is an $`SU(\stackrel{~}{n}_c)\times U(1)^{n_c\stackrel{~}{n}_c}`$ theory ($`\stackrel{~}{n}_cn_fn_c`$), with $`n_f`$ massless quarks $`q`$ and $`n_c\stackrel{~}{n}_c=2n_cn_f`$ massless singlets. Their charges are summarized in Tab. 5. The effective action in this case, upon the $`N=1`$ perturbation, is $$W_{bar}=\sqrt{2}\text{Tr}(q\varphi \stackrel{~}{q})+\frac{\sqrt{2}}{\stackrel{~}{n}_c}\text{Tr}(q\stackrel{~}{q})\left(\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}\psi _k\right)\sqrt{2}\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}\psi _ke_k\stackrel{~}{e}_k+\mu \left(\mathrm{\Lambda }\underset{i=1}{\overset{n_c\stackrel{~}{n}_c}{}}x_i\psi _i+\frac{1}{2}\text{Tr}\varphi ^2\right).$$ (8.24) Again, we add the quark mass terms $$\text{Tr}(mq)\underset{k,i}{}S_k^im_ie_k\stackrel{~}{e}_k$$ (8.25) and minimize the potential. The equations are: $$q_a^iq_i^b\stackrel{~}{q}_a^i\stackrel{~}{q}_i^b=\nu \delta _a^b;$$ (8.26) $$q_a^i\stackrel{~}{q}_i^b\frac{1}{\stackrel{~}{n}_c}\delta _a^b(q_c^i\stackrel{~}{q}_i^c)+\mu \varphi _a^b=0;$$ (8.27) $$q_a^im_i+\sqrt{2}\varphi _a^bq_b^i+\frac{\sqrt{2}}{\stackrel{~}{n}_c}q_a^i\underset{k}{}\psi _k=0;$$ (8.28) $$\stackrel{~}{q}_i^am_i+\sqrt{2}\varphi _b^a\stackrel{~}{q}_i^b+\frac{\sqrt{2}}{\stackrel{~}{n}_c}\stackrel{~}{q}_i^a\underset{k}{}\psi _k=0;$$ (8.29) $$\frac{\sqrt{2}}{\stackrel{~}{n}_c}\text{Tr}(\stackrel{~}{q}q)\sqrt{2}e_k\stackrel{~}{e}_k+\mu \mathrm{\Lambda }x_k=0;$$ (8.30) $$(\psi _kS_k^im_i)e_k=0;(\psi _kS_k^im_i)\stackrel{~}{e}_k=0,$$ (8.31) and $$\varphi =\frac{1}{\sqrt{2}}\text{diag}(m_1,\mathrm{},m_r,c,\mathrm{},c);c=\frac{1}{\stackrel{~}{n}_cr}\underset{k=1}{\overset{r}{}}m_k.$$ (8.32) We find two types of vacua. The first type has $`e_k=\stackrel{~}{e}_k=(\mu \mathrm{\Lambda }x_k/\sqrt{2})^{1/2}`$ for all $`k=1,\mathrm{}n_c\stackrel{~}{n}_c`$. Minimizing the potential in this case, we find $$𝒩_2=\underset{r=0}{\overset{\stackrel{~}{n}_c1}{}}(\stackrel{~}{n}_cr){}_{n_f}{}^{}C_{r}^{}$$ (8.33) $`N=1`$ vacua, characterized by the VEVS Eq.(8.32) and $$d_i,\stackrel{~}{d}_i\sqrt{\mu m}\stackrel{m_i0}{}0,e_k,\stackrel{~}{e}_k\sqrt{\mu \mathrm{\Lambda }}.$$ (8.34) The unbroken $`SU(\stackrel{~}{n}_cr)`$ gauge group gives $`\stackrel{~}{n}_cr`$ vacua each. These vacua describe the vacua with unbroken $`U(n_f)`$ symmetry, which are known to exist from the large $`\mu `$ analysis. The total number of vacua of this group found here agrees with Eq.(6.22) at the large $`\mu `$ regime. The second type of vacua in Eqs. (8.24,8.25) has one of the $`e_k=\stackrel{~}{e}_k=0`$ (hence $`n_c\stackrel{~}{n}_c=2n_cn_f`$ choices) while $`W/\psi _k=0`$ requires quarks to condense with $`q=\stackrel{~}{q}\sqrt{\mu \mathrm{\Lambda }}`$. Dropping $`e_k=\stackrel{~}{e}_k=0`$ from the Lagrangian, it becomes the same as that of the non-baryonic root Eqs. (8.1,8.2) and gives $`(2n_cn_f){}_{n_f}{}^{}C_{\stackrel{~}{n}_c}^{}`$ vacua. This precisely compensates the exclusion of $`r=\stackrel{~}{n}_c`$ in the sum for the non-baryonic roots and the correct total number of vacua $`𝒩_1+𝒩_2`$ is obtained. We thus find that both the number and the symmetry properties of the $`N=1`$ theories at small adjoint mass $`\mu `$ match exactly those found at large $`\mu .`$ For vacua with $`r=1`$, the “quarks” in the effective Lagrangian (8.1) are nothing but the $`U(1)^{n_c1}`$ monopoles in the fundamental representation of $`U(n_f)`$; this is checked by studying the monodromy around the singularity which showed that the “quarks” are indeed magnetically charged. Therefore the standard picture of confinement and flavor symmetry breaking by condensation of flavor-non-singlet monopoles is valid for these vacua. For general vacua $`r>1`$ associated with various nonbaryonic roots, the effective Lagrangian (8.1) describes correctly the physics of $`N=1`$ vacua at small $`\mu `$, in terms of magnetic quarks of a non Abelian $`SU(r)\times U(1)^{n_cr}`$ theory. In contrast to the $`r=1`$ case, these quarks cannot be identified with the semiclassical monopoles of the maximally Abelian $`U(1)^{n_c1}`$ group. Note that the condensation of such monopoles in the rank-$`r`$ anti-symmetric tensor representation, which might be suggested from the number of the singularities which group into a nearby cluster in the small $`m_i`$ limit and at the same time from the semiclassical analysis (see Appendix B), would have yielded the correct pattern of symmetry breaking; at the same time, however, it would have led to an uncomfortably large number of Nambu-Goldstone bosons associated to the accidental $`SU({}_{n_f}{}^{}C_{r}^{})`$ symmetry. The system avoids this paradox elegantly, by having magnetic quarks as low-energy degrees of freedom and having these condensed. These facts, and the comparison of their quantum numbers, lead us to conclude that, as we approach the non-baryonic roots from semi-classical (large VEV) region on the Coulomb branch, the semi-classical monopoles in the rank-$`r`$ anti-symmetric tensor representation are smoothly matched to “baryons” of the $`SU(r)`$ theory, $$ϵ^{a_1\mathrm{}a_r}q_{a_1}^{i_1}q_{a_2}^{i_2}\mathrm{}q_{a_r}^{i_r},$$ (8.35) and break up (the system being infrared-free) into weakly coupled magnetic quarks, before becoming massless. The case $`r=n_f/2`$ is exceptional and highly non-trivial. Although the analysis leading to Eq.(8.19), Eq.(8.22), is formally valid in this case also, physics is really different. The low energy degrees of freedom involve relatively nonlocal states, arising from the nontrivial, class 3 CFT . In this case, the theory is in the same universality class as the finite $`SU(n_f/2)`$ theories. For an explicit check with monodromy, see Appendix E, for the simplest example of this type, $`r=2`$ vacua of the $`SU(3)`$ theory with $`n_f=4`$. As for the second group of vacua with vanishing VEVS found at the “baryonic root”, they are in the free-magnetic phase, with observable (magnetic) quarks, weakly interacting with $`SU(n_fn_c)`$ gauge fields. ### 8.2 $`USp(2n_c)`$ In the $`USp(2n_c)`$ gauge theories, the first type of vacua can be identified more easily by first considering the equal but novanishing quark masses. The adjoint VEVS in the curve Eq.(6.24) can be chosen so as to factor out the behavior $$y^2=(x+m^2)^{2r}[\mathrm{}],r=1,2,\mathrm{}$$ (8.36) which describes an $`SU(r)\times U(1)^{n_cr+1}`$ gauge theory with $`n_f`$ quarks. See Section 9.2. These (trivial) superconformal theories belong in fact to the same universality classes as those found in the $`SU(n_c)`$ gauge theory, as pointed out by Eguchi and others . They are therefore described by exactly the same Lagrangian Eq.(8.1). At each vacuum with $`r`$, the symmetry (of equal mass theory, $`U(n_f)`$) is broken spontaneously as $$U(n_f)U(r)\times U(n_fr).$$ (8.37) When a small mass splitting is added among $`m_i`$’s, each of the $`r`$ vacuum splits into $`{}_{n_f}{}^{}C_{r}^{}`$ vacua, leading to the total of $$(2n_c+2n_f)\underset{r=0}{\overset{(n_f1)/2}{}}{}_{n_f}{}^{}C_{r}^{}=(2n_c+2n_f)\mathrm{\hspace{0.17em}2}^{n_f1},(n_f=\text{odd})$$ (8.38) $$(2n_c+2n_f)\underset{r=0}{\overset{n_f/21}{}}{}_{n_f}{}^{}C_{r}^{}+\frac{2n_c+2n_f}{2}{}_{n_f}{}^{}C_{n_f/2}^{}=(2n_c+2n_f)\mathrm{\hspace{0.17em}2}^{n_f1},(n_f=\text{even}),$$ (8.39) vacua of this type, consistently with Eq. (3.44). These $`N=1`$ vacua seem to have been overlooked in altogether<sup>17</sup><sup>17</sup>17With a hindsight, we see that this was inevitable: as we show below, there is no local effective Lagrangian description of low-energy physics in the $`m_i=0`$ limit and Argyres, Plesser and Shapere worked precisely in such a regime. . In the massless limit the underlying theories possess a larger, flavor $`SO(2n_f)`$ symmetry. We know also from the large $`\mu `$ analysis that in the first group of vacua (with finite vevs), this symmetry is broken spontaneously to $`U(n_f)`$ symmetry always. How can such a result be consistent with Eq.(8.37) of equal (but nonvanishing) mass theory? What happens is that in the massless limit various $`N=1`$ vacua with different symmetry properties Eq.(8.37) (plus eventually other singularities) coalesce. The location of this singularity can be obtained exactly in terms of Chebyshev polynomials, and is given by Eqs. (6.26)-(6.29). At the singularity there are mutually non-local dyons and hence the theory is at a non-trivial infrared fixed point. In the example of $`USp(4)`$ theory with $`n_f=4`$, we have explicitly verified this by determining the singularities and branch points at finite equal mass $`m`$ and by studying the limit $`m0`$. There is no description in terms of a weakly coupled local field theory, just as in the case $`r=n_f/2`$ for $`SU(n_c)`$ theories. Since the global flavor symmetry is $`SO(2n_f)`$, these superconformal theories belong to a different universality class as compared to those at finite mass or to those of $`SU(n_c)`$ theories. When $`n_f`$ is even, the theory is in the same universality class as the finite $`USp(n_f2)`$ theories. When $`n_f`$ is odd, the theory is in a different universality class of “strong-coupled conformal theories.” We find this behavior resonable because the semi-classical monopoles are in the spinor representation of the $`SO(2n_f)`$ flavor group and, in contrast to the situation in $`SU(n_c)`$ theories, cannot “break up” into quarks in the vector representation. They are therefore likely to persist at the singularity and in general make the theory superconformal. We indeed checked that there are mutually non-local degrees of freedom using monodromy of the curve for $`USp(4)`$ theory with $`n_f=5`$.<sup>18</sup><sup>18</sup>18The case $`n_f=3`$, however, is special. In this case, the analysis of the curve in Section 9.2 tells us that the theory at the Chebyshev vacuum is the same as the $`SU(2)`$ theory with $`n_f=3`$. Seiberg and Witten studied this case and found that the singularity can be described in terms of massless monopoles in the spinor representation of the flavor $`SO(6)`$ group interacting locally with the magnetic photon. Note that there is no problem with unwanted Nambu–Goldstone multiplets in this case (see Section 8 in ). We believe that the situation with higher $`n_f`$ is also non-local; this is because a local effective low-energy Lagrangian of massless monopoles in the spinor representation of the flavor group would have an accidental $`U(2^{n_f1})`$ symmetry and would lead to an unacceptably large number of Nambu–Goldstone multiplets. Once the quark masses are turned on, however, the flavor group reduces to (at least) $`U(n_f)`$ and it becomes possible for monopoles to break up into quarks; this explains the behavior in the equal mass case. As for the second group of vacua, the situation is more analogous to the case of $`SU(n_c)`$ theories. The superpotential reads in this case (by adding mass terms to Eq.(5.10) of ): $`W`$ $`=`$ $`\mu \left(\text{Tr}\varphi ^2+\mathrm{\Lambda }{\displaystyle \underset{a=1}{\overset{2n_c+2n_f}{}}}x_a\psi _a\right)+{\displaystyle \frac{1}{\sqrt{2}}}q_a^i\varphi _b^aq_c^iJ^{bc}+{\displaystyle \frac{m_{ij}}{2}}q_a^iq_b^jJ^{ab}`$ (8.40) $`+`$ $`{\displaystyle \underset{a=1}{\overset{2n_c+2n_f}{}}}\left(\psi _ae_a\stackrel{~}{e}_a+S_a^im_ie_a\stackrel{~}{e}_a\right),`$ where $`J=i\sigma _2\mathrm{𝟏}_{n_c}`$ and $$m=i\sigma _2\text{diag}(m_1,m_2,\mathrm{},m_{n_f}).$$ (8.41) By minimizing the potential, we find $$𝒩_2=\underset{r=0}{\overset{\stackrel{~}{n}_c}{}}(\stackrel{~}{n}_cr+1){}_{n_f}{}^{}C_{r}^{}$$ (8.42) vacua, which precisely matches the number of the vacua of the second group, with squark VEVS behaving as $$q_i,\stackrel{~}{q}_i\sqrt{\mu m_i}\stackrel{m_i0}{}0.$$ (8.43) These are the desired $`SO(2n_f)`$ symmetric vacua. ## 9 Quark-Mass Perturbation on the Curve In this section we develop a perturbation theory in the quark masses for the singularities of the Seiberg-Witten curves of $`SU(n_c)`$ and $`USp(2n_c)`$ theories, around certain conformal points. The results of this section allows us, on the one hand, to establish the connection between the classes of CFT singularities of $`N=2`$ space of vacua and the $`N=1`$ vacua surviving the perturbation $`\mu \mathrm{\Phi }^2`$ (as discussed in Section 6), and on the other, to identify the $`N=1`$ vacua at small $`\mu `$ (whose physics was discussed in the previous section) with those found at large $`\mu `$. ### 9.1 Perturbation around CFT Points: $`SU(n_c)`$ i) Generic $`r`$ ($`r<\frac{n_f}{2})`$ and formulation of the problem. Suppose the conformal point (Eq.(6.13)), $`\text{diag}\varphi =(\underset{r}{\underset{}{0,0,\mathrm{},0}},\varphi _{r+1}^{(0)},\mathrm{},\varphi _{n_c}^{(0)}),_{a=r}^{n_c}\varphi _a^{(0)}=0,`$ where the full curve with $`m_i=0`$ $$y^2=\underset{k=1}{\overset{n_c}{}}(x\varphi _k)^2+4\mathrm{\Lambda }^{2n_cn_f}\underset{j=1}{\overset{n_f}{}}(x+m_j)\|_{m_i=0}$$ (9.1) takes the form $$y^2=x^{2r}(x\beta _{01})^2\mathrm{}(x\beta _{0,n_cr1})^2(x\gamma _0)(x\kappa _0),r=0,1,2,\mathrm{},[n_f/2]$$ (9.2) is given. For nonzero and small bare quark masses $`m_i`$ the multiple zero at the origin will split, and other zeros will be shifted. We require that the perturbed singularity, $$\text{diag}\varphi =(\lambda _1,\lambda _2,\mathrm{},\lambda _r,\varphi _{r+1}^{(0)}+\delta \varphi _{r+1},\mathrm{},\varphi _{n_c}^{(0)}+\delta \varphi _{n_c}),$$ (9.3) $$\underset{a=1}{\overset{r}{}}\lambda _a+\underset{a=r+1}{\overset{n_c}{}}\delta \varphi _a=0,$$ (9.4) be such that the curve $$y^2=F(x)=\underset{a=1}{\overset{r}{}}(x\lambda _a)^2\underset{a=r+1}{\overset{n_c}{}}(x\varphi _{0a}\delta \varphi _a)^2\underset{i=1}{\overset{n_f}{}}(x+m_i),$$ (9.5) is maximally singular, i.e., $$y^2=\underset{a=1}{\overset{r}{}}(x\alpha _a)^2\{\underset{b=1}{\overset{n_cr1}{}}(x\beta _b)^2\}(x\gamma )(x\kappa );$$ (9.6) $$beta_i=\beta _{0i}+\delta \beta _i;\gamma =\gamma _0+\delta \gamma ,\kappa =\kappa _0+\delta \kappa .$$ (9.7) The problem is to determine how many such sets $`\{\lambda _i,\alpha _i,\delta \beta _i,\delta \varphi _i,\delta \gamma ,\delta \kappa \}`$ exist. The condition that the curve is maximally singular (maximal number of double branch points) can be expressed as ($`F^{^{}}dF(x)/dx`$): $`F(\alpha _a)`$ $`=`$ $`F^{^{}}(\alpha _a)=0,a=1,2,\mathrm{}r;F(\gamma )=F(\kappa )=0`$ $`F(\beta _i)`$ $`=`$ $`F^{^{}}(\beta _i)=0,i=1,2,\mathrm{}n_cr1,`$ (9.8) There are $`2n_c`$ relations for $`2n_c`$ unknowns. Let us consider the case $`r=1`$ first. In this case there is only one set of $`\{\alpha ,\lambda \}`$. Consider first the $`2n_c2`$ relations $$F(\beta _i)=F^{^{}}(\beta _i)=0,i=1,2,\mathrm{}n_c2,F(\gamma )=F(\kappa )=0,$$ (9.9) and expand each equation in the small quantities $`\eta _A\{\lambda ,\alpha ,\delta \beta _i,\delta \varphi _i,\delta \gamma ,\delta \kappa \}`$. By assumption these quantities are zero if $`m_i=0`$. To zeroth order, (9.9) is satisfied by assumption. To first order we find a linear system of $`2n_c2`$ equations, $`{\displaystyle \underset{A}{}}{\displaystyle \frac{dF(\beta _i)}{d\eta _A}}\eta _A`$ $`=`$ $`{\displaystyle \underset{A}{}}{\displaystyle \frac{dF^{^{}}(\beta _i)}{d\eta _A}}\eta _A=0,i=1,2,\mathrm{}n_c2,`$ $`{\displaystyle \underset{A}{}}{\displaystyle \frac{dF(\gamma )}{d\eta _A}}\eta _A`$ $`=`$ $`{\displaystyle \underset{A}{}}{\displaystyle \frac{dF(\kappa )}{d\eta _A}}\eta _A=0,`$ (9.10) which determine $`\delta \beta _i,\delta \varphi _i,\delta \gamma ,\delta \kappa `$ uniquely, in terms of $`m_i`$’s and of $`\lambda ,\alpha `$. The two equations $$F(\alpha )=F^{^{}}(\alpha )=0$$ (9.11) (which have no zeroth order counterpart) cannot be linearized. In fact, they give $$C(\alpha \lambda )^2\underset{i=1}{\overset{n_f}{}}(\alpha +m_i)=0,C=O(1);$$ (9.12) $$2C(\alpha \lambda )+C^{^{}}(\alpha \lambda )^2\underset{i=1}{\overset{n_f}{}}\underset{ji}{\overset{n_f}{}}(\alpha +m_i)=0,C^{^{}}=O(1).$$ (9.13) Eq.(9.12) and Eq.(9.13) must be solved for $`\alpha ,\lambda `$. We see immediately that $`\alpha `$ must be very special. Set $$m_i=O(ϵ)1,$$ (9.14) and suppose $`\alpha O(ϵ^{^{}}).`$ In either cases, $`ϵϵ^{^{}}`$ or $`ϵϵ^{^{}}`$, we get from Eq.(9.12) and Eq.(9.13) $$\|(\alpha \lambda )^2\|ϵ^{n_f};\|\alpha \lambda \|ϵ^{n_f1};$$ (9.15) or an analogous relation with $`ϵϵ^{^{}}`$. These cannot be satisfied (in other words, there are no solutions of this type). The only way out (to get solutions) is to assume that $`ϵϵ^{^{}}`$ and take $$\alpha =m_i+\mathrm{\Delta },\mathrm{\Delta }O(ϵ^2),i=1,2,\mathrm{},n_f.$$ (9.16) There are obviously $`n_f`$ such possibilities. We find now from Eq.(9.12) and Eq.(9.13), $$\|(\alpha \lambda )^2\|\mathrm{\Delta }ϵ^{n_f1};\|\alpha \lambda \|ϵ^{n_f1}.$$ (9.17) Note that in Eq.(9.13) the terms containing $`\mathrm{\Delta }`$ are indeed smaller than the term kept. Now these can be solved and give $$\mathrm{\Delta }ϵ^{n_f1};\|\alpha \lambda \|ϵ^{n_f1},$$ (9.18) in other words, $$\alpha =m_i+O(ϵ^{n_f1});\lambda =m_i+O(ϵ^{n_f1}).$$ (9.19) We found obviously $`n_f`$ solutions with $`r=1`$, according to which $`m_i`$ is used to make the solution. A straightforward generalization to generic $`r`$, $`r<\frac{n_f}{2}`$ leads to the result that $`\lambda _a\alpha _a,`$ $`a=1,2,\mathrm{},r`$ must be chosen to be equal to $`r`$ out of $`n_f`$ masses, $`m_i`$. There are thus $`{}_{n_f}{}^{}C_{r}^{}`$ solutions of this type, according to which masses are chosen to construct the solution. In order to see the order of magnitude of $`\lambda _a\alpha _a,`$ let us write the low-energy curve (at $`x\mathrm{\Lambda }`$) as $$y^2F(x)=\underset{a=1}{\overset{r}{}}(x\lambda _a)^2+\frac{4}{\mathrm{\Lambda }^{n_f2r}}\underset{i=1}{\overset{n_f}{}}(x+m_i),$$ (9.20) and require that this curve behaves as $$F(x)=\underset{a=1}{\overset{r}{}}(x\alpha _a)^2,$$ (9.21) if we neglect zeros at $`xO(\mathrm{\Lambda })`$. Namely, we require $`F(\alpha _1)`$ $`=`$ $`{\displaystyle \underset{a=1}{\overset{r}{}}}(\alpha _1\lambda _a)^2+{\displaystyle \frac{4}{\mathrm{\Lambda }^{n_f2r}}}{\displaystyle \underset{i=1}{\overset{n_f}{}}}(\alpha _1+m_i)=0,`$ (9.22) $`F^{}(\alpha _1)`$ $`=`$ $`2{\displaystyle \underset{b=1}{\overset{r}{}}}(\alpha _1\lambda _b){\displaystyle \underset{ab}{}}(\alpha _1\lambda _a)^2+{\displaystyle \frac{4}{\mathrm{\Lambda }^{n_f2r}}}{\displaystyle \underset{j=1}{\overset{n_f}{}}}{\displaystyle \underset{ij}{}}(\alpha _1+m_i)=0,`$ and similarly for $`\alpha _2,\mathrm{},\alpha _r`$. To satisfy the first equation (9.22) down to $`O(m^{n_f+1})`$, all we need is to set $`\alpha _1=\lambda _1`$, etc. The second equation (LABEL:eq:F') is then approximated by keeping only the first power in $`(\alpha _1\lambda _1)`$ (in other words, only keeping $`b=1`$ in the sum), $$F^{}(\alpha _1)=2(\alpha _1\lambda _1)\underset{a1}{}(\alpha _1\lambda _a)^2+\frac{4}{\mathrm{\Lambda }^{n_f2r}}\underset{j=1}{\overset{n_f}{}}\underset{ij}{}(\alpha _1+m_i)=0.$$ (9.24) Barring an accidental cancellations between $`\alpha _1`$ and other $`\lambda _a`$ ($`a1`$), we find $$\alpha _1\lambda _1O(m^{n_f1}/\mathrm{\Lambda }^{n_f2n_c}/m^{2(r1)})=O(m^{n_f2r+1}/\mathrm{\Lambda }^{n_f2r}).$$ (9.25) Using this fact, the first term in (9.22) is $`O(m^{n_f2r+1}/\mathrm{\Lambda }^{n_f2r})^2\times O(m^{2(r1)}),`$ which is much smaller than the second term of $`O(m^{n_f})`$. Therefore, the second term must vanish by itself, which requires $`\alpha _1=m_1`$ etc. Repeating the same analysis for every $`\alpha _a`$, we need to choose $`r`$ masses out of $`n_f`$ and assign $`\alpha _a=m_a`$ etc. Then we should retain only the term $`j=1`$ in the second term in Eq. (9.24), and we find $$\alpha _1\lambda _1=\frac{4}{\mathrm{\Lambda }^{n_f2r}}\underset{i1}{\overset{n_f}{}}(m_1+m_i)\frac{1}{2_{a1}^r(m_1m_a)^2}.$$ (9.26) Obviously the case of equal masses is singular and beyond the validity of this analysis. ii) $`r=\frac{n_f}{2}`$ with odd $`n_cn_f/2`$ For special cases with $`r=n_f/2`$ some of the considerations above are not valid (e.g., Eq.(9.25)), and the analysis must be done ad hoc. Fortunately, in these cases it is possible to find the unperturbed configuration $`\{\varphi _0\}`$ explicitly as in Eq.(6.16), Eq.(6.17). When $`n_f`$ is even, choose $`n_f/2`$ eigenvalues $`\varphi _1=\mathrm{}=\varphi _{n_f/2}`$ vanishing. Then the curve becomes $$y^2=x^{n_f}\left[\underset{k=1}{\overset{n_cn_f/2}{}}(x\varphi _k)^24\mathrm{\Lambda }^{2n_cn_f}\right].$$ (9.27) By identifying the first term in the square bracket with $`(2\mathrm{\Lambda }^{n_cn_f/2}T_{n_cn_f/2}(x/2\mathrm{\Lambda }))^2`$, where $$\varphi _k=2\mathrm{\Lambda }\mathrm{cos}\pi (2k1)/2(n_cn_f/2),k=1,\mathrm{},n_cn_f/2,$$ (9.28) one easily finds that $`y^2`$ $`=`$ $`x^{n_f}\left[\left(2\mathrm{\Lambda }^{n_cn_f/2}T_{n_cn_f/2}\left({\displaystyle \frac{x}{2\mathrm{\Lambda }}}\right)\right)^24\mathrm{\Lambda }^{2n_cn_f}\right]`$ (9.29) $`=`$ $`4x^{n_f}\mathrm{\Lambda }^{2n_cn_f}\mathrm{sin}^2\left[\left(n_c{\displaystyle \frac{n_f}{2}}\right)\mathrm{arccos}{\displaystyle \frac{x}{2\mathrm{\Lambda }}}\right].`$ There are $`n_f`$ zeros at $`x=0`$, and double zeros at $`x=2\mathrm{\Lambda }\mathrm{cos}\pi k/(n_cn_f/2)`$ for $`k=1,\mathrm{},n_cn_f/21`$, and single zeros at $`x=\pm 2\mathrm{\Lambda }`$. In the absence of quark masses, the theory is invariant under $`Z_{2n_cn_f}`$ symmetry: $`xe^{2\pi i/(2n_cn_f)}x`$, $`\varphi _ae^{2\pi i/(2n_cn_f)}\varphi _a`$. However, the gauge invariant symmetric polynomials of $`\varphi _k`$ vanish for odd powers because of the equal number of positive and negative ones with the same absolute values. Therefore the Chebyshev solutions discussed here appear only $`n_cn_f/2`$ times. This point is crucial for the vacuum counting (see Eq.(6.21)) to work out correctly. We first study the specific case of $`2n_c=n_f+2`$. More general cases will be seen to reduce to this case. The Chebyshev solution is obtained in the massless limit by setting all of $`\varphi _a=0`$: $$y^2=x^{2n_c}4\mathrm{\Lambda }^2x^{n_f}=x^{n_f}(x+2\mathrm{\Lambda })(x2\mathrm{\Lambda }).$$ (9.30) The zero at $`x=0`$ is of degree $`n_f`$, and there are other isolated zeros at $`x=\pm 2\mathrm{\Lambda }`$. Under the perturbation by generic quark masses, we go back to the original curve. The only way that the curve can be arranged to have $`n_c1`$ double zeros as $$y^2=\underset{a=1}{\overset{n_c1}{}}(x\alpha _a)^2(x+2\mathrm{\Lambda }\beta )(x2\mathrm{\Lambda }\gamma )$$ (9.31) is by assuming $$\alpha _am,\varphi _1\mathrm{}\varphi _{n_c2}m,\varphi _{n_c1}=\varphi _{n_c}(m\mathrm{\Lambda })^{1/2}.$$ (9.32) (These behaviors have been suggested by the numerical solution of several explicit examples.) Studying the region $`xm`$, and neglecting $`\beta ,\gamma \mathrm{\Lambda }`$, $`x\varphi _{n_c}`$, we need to solve the equation $$\varphi _{n_c}^4\underset{a=1}{\overset{n_c2}{}}(x\varphi _a)^24\mathrm{\Lambda }^2\underset{i=1}{\overset{2n_c2}{}}(x+m_i)=4\mathrm{\Lambda }^2\underset{a=1}{\overset{n_c1}{}}(x\alpha _a)^2.$$ (9.33) Note that the traceless condition for $`\varphi _a`$ is not a stringent constraint because $`\varphi _{n_c1}=\varphi _{n_c}=O(m\mathrm{\Lambda })^{1/2}`$ can shift by small amount to absorb the trace of $`\varphi _a`$. By moving terms around, we find $`{\displaystyle \underset{i=1}{\overset{2n_c2}{}}}(x+m_i)`$ $`=`$ $`\left[{\displaystyle \underset{a=1}{\overset{n_c1}{}}}(x\alpha _a)+i{\displaystyle \frac{\varphi _{n_c}^2}{2\mathrm{\Lambda }}}{\displaystyle \underset{a=1}{\overset{n_c2}{}}}(x\varphi _a)\right]\left[{\displaystyle \underset{a=1}{\overset{n_c1}{}}}(x\alpha _a)i{\displaystyle \frac{\varphi _{n_c}^2}{2\mathrm{\Lambda }}}{\displaystyle \underset{a=1}{\overset{n_c2}{}}}(x\varphi _a)\right].`$ For this identity to hold for any $`x`$, $`2n_c2`$ zeros of l.h.s. $`(m_i)`$ must coincide with $`n_c2`$ zeros in each of square brackets in the r.h.s. of the equation. Therefore, we divide up $`2n_c2`$ masses into to sets $`\{p_i,i=1,\mathrm{},n_c1\}`$ and $`\{q_i,i=1,\mathrm{},n_c1\}`$. This gives us $`{}_{n_f}{}^{}C_{n_f/2}^{}`$ choices. By identifying $`p_i`$ to the first square bracket and $`q_i`$ to the second square bracket, we find $`{\displaystyle \underset{i=1}{\overset{n_c1}{}}}(x+p_i)`$ $`=`$ $`{\displaystyle \underset{a=1}{\overset{n_c1}{}}}(x\alpha _a)+i{\displaystyle \frac{\varphi _{n_c}^2}{2\mathrm{\Lambda }}}{\displaystyle \underset{a=1}{\overset{n_c2}{}}}(x\varphi _a),`$ $`{\displaystyle \underset{i=1}{\overset{n_c1}{}}}(x+q_i)`$ $`=`$ $`{\displaystyle \underset{a=1}{\overset{n_c1}{}}}(x\alpha _a)i{\displaystyle \frac{\varphi _{n_c}^2}{2\mathrm{\Lambda }}}{\displaystyle \underset{a=1}{\overset{n_c2}{}}}(x\varphi _a).`$ (9.35) Expanding both sides in terms of symmetric polynomials, $`{\displaystyle \underset{k=0}{\overset{n_c1}{}}}s_k(p)x^{n_c1k}`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{n_c1}{}}}(1)^ks_k(\alpha )x^{n_ck1}+i{\displaystyle \frac{\varphi _{n_c}^2}{2\mathrm{\Lambda }}}{\displaystyle \underset{k=0}{\overset{n_c2}{}}}(1)^ks_k(\varphi )x^{n_ck2},`$ $`{\displaystyle \underset{k=0}{\overset{n_c1}{}}}s_k(q)x^{n_c1k}`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{n_c1}{}}}(1)^ks_k(\alpha )x^{n_ck1}i{\displaystyle \frac{\varphi _{n_c}^2}{2\mathrm{\Lambda }}}{\displaystyle \underset{k=0}{\overset{n_c2}{}}}(1)^ks_k(\varphi )x^{n_ck2}.`$ Finally, we find the solutions $`(1)^ks_k(p)`$ $`=`$ $`s_k(\alpha )i{\displaystyle \frac{\varphi _{n_c}^2}{2\mathrm{\Lambda }}}s_{k1}(\varphi ),`$ $`(1)^ks_k(q)`$ $`=`$ $`s_k(\alpha )+i{\displaystyle \frac{\varphi _{n_c}^2}{2\mathrm{\Lambda }}}s_{k1}(\varphi ).`$ (9.37) Here and below, we use the notation that $`s_1=0`$. By setting $`k=1`$ and subtracting both sides, we find $`\varphi _{n_c}^2=i\mathrm{\Lambda }(s_1(p)s_1(q))=i\mathrm{\Lambda }_{k=1}^{n_f/2}(p_iq_i)`$. All $`s_k(\alpha )`$ and $`s_{k1}(\varphi )`$ are then given in terms of $`s_k(p)`$ and $`s_k(q)`$, and hence we can write order $`n_c1`$ ($`n_c2`$) polynomial equation for $`\alpha `$ ($`\varphi `$) which can always be solved to find all $`\alpha `$’s ($`\varphi `$’s). For smaller even $`n_f`$ with $`n_cn_f/2`$ odd, the problem reduces the one with $`2n_cn_f=2`$. To see this, we write the curve around the Chebyshev point as $$y^2=\underset{a=1}{\overset{n_f/2}{}}(x\varphi _a)^2\underset{k=1}{\overset{n_cn_f/2}{}}(x\varphi _k)^24\mathrm{\Lambda }^{2n_cn_f}\underset{i=1}{}(x+m_i),$$ (9.38) where $`\varphi _k=2\mathrm{\Lambda }\mathrm{cos}\pi (k1/2)/(n_cn_f/2)`$. When $`k=(n_cn_f/2+1)/2`$, $`\varphi _k=0`$. Therefore, neglecting the fluctuations of $`\varphi _k`$ ($`k(n_cn_f/2+1)/2`$), the curve reduces to $$y^2=\underset{a=1}{\overset{n_f/2+1}{}}(x\varphi _a)^2\underset{k=1,k(n_cn_f/2+1)/2}{\overset{n_cn_f/2}{}}\varphi _k^24\mathrm{\Lambda }^{2n_cn_f}\underset{i=1}{}(x+m_i),$$ (9.39) where $`\varphi _a`$ with $`a=n_f/2+1`$ is $`\varphi _k`$ with $`k=(n_cn_f/2+1)/2`$. The product of non-vanishing $`\varphi _k`$’s can be obtained as follows. Recalling $`T_N(x)=2^{N1}_{k=1}^N(xw_k)`$ with $`w_k=\mathrm{cos}\pi (k1/2)/N`$, we obtain for odd $`N`$, $`T_N^{}(0)=2^{N1}_{k(n_cn_f/2+1)/2}(w_k)`$, while $`T_N^{}(x)=N\mathrm{sin}(N\mathrm{arccos}x)/\sqrt{1x^2}`$ and hence $`T_N^{}(0)=N\mathrm{sin}N\pi /2=N(1)^{(N1)/2}`$. Therefore, $`_{k=1,k(n_cn_f/2+1)/2}^{n_cn_f/2}\varphi _k^2=(T_N^{}(0))^2\mathrm{\Lambda }^{2N2}=N^2\mathrm{\Lambda }^{2N2}`$. The low-energy curve is then $$\frac{y^2}{(n_cn_f/2)^2\mathrm{\Lambda }^{2n_cn_f2}}=\underset{a=1}{\overset{n_f/2+1}{}}(x\varphi _a)^24\frac{\mathrm{\Lambda }^2}{(n_cn_f/2)^2}\underset{i=1}{}(x+m_i).$$ (9.40) This is nothing but the curve of $`SU(n_f/2+1)`$ theory with the dynamical scale $`\mathrm{\Lambda }/(n_cn_f/2)`$. iii) $`r=\frac{n_f}{2}`$ with even $`n_cn_f/2`$ The Chebyshev solution Eq.(9.28), Eq.(9.29) is valid in this case also. Again, we first study the special case with $`2n_c=n_f+4`$. The Chebyshev solution is obtained in the massless limit by setting all but two of $`\varphi _a=0`$: $$y^2=x^{n_f}\left[(x\varphi _{n_{c1}})^2(x\varphi _{n_c})^24\mathrm{\Lambda }^4\right]=x^{n_f}(x^2\varphi _{n_c}^22\mathrm{\Lambda }^2)(x^2\varphi _{n_c}^2+2\mathrm{\Lambda }^2).$$ (9.41) By choosing $`\varphi _{n_c1}=\varphi _{n_c}=\sqrt{2}\mathrm{\Lambda }`$, the curve becomes $$y^2=x^{2n_c}4\mathrm{\Lambda }^2x^{n_f+2}=x^{n_f+2}(x2\mathrm{\Lambda })(x+2\mathrm{\Lambda }).$$ (9.42) The zero at $`x=0`$ is of degree $`n_f+2`$, and there are other isolated zeros at $`x=\pm 2\mathrm{\Lambda }`$. There is another solution with $`\varphi _{n_c1}=\varphi _{n_c}=i\sqrt{2}\mathrm{\Lambda }`$ as required by the discrete $`Z_4`$ symmetry under which the Chebyshev solution transforms as a doublet. Under the perturbation by generic quark masses, we go back to the original curve. The only way that the curve can be arranged to have $`n_c1`$ double zeros as $$y^2=\underset{a=1}{\overset{n_c1}{}}(x\alpha _a)^2(x+2\mathrm{\Lambda }\beta )(x2\mathrm{\Lambda }\gamma )$$ (9.43) is by assuming $$\varphi _am,\alpha _1\mathrm{}\alpha _{n_c2}m,\alpha _{n_c2}=\alpha _{n_c1}(m\mathrm{\Lambda })^{1/2}.$$ (9.44) (The choice $`\alpha _{n_c1}=\alpha _{n_c}`$ has been suggested from several explicit examples studied numerically.) Studying the region $`xm`$, and neglecting $`\beta ,\gamma \mathrm{\Lambda }`$, $`x\alpha _{n_c1}`$, we must solve the equation $$4\mathrm{\Lambda }^4\underset{a=1}{\overset{n_c2}{}}(x\varphi _a)^24\mathrm{\Lambda }^4\underset{i=1}{\overset{2n_c4}{}}(x+m_i)=4\mathrm{\Lambda }^2\alpha _{n_c1}^4\underset{a=1}{\overset{n_c3}{}}(x\alpha _a)^2.$$ (9.45) Note that the traceless condition for $`\varphi _a`$ is not a stringent constraint because $`\varphi _{n_c1}=\varphi _{n_c}=O(\mathrm{\Lambda })`$ can shift by small amount to absorb the trace of $`\varphi _a`$. By moving terms around, we find $`{\displaystyle \underset{i=1}{\overset{2n_c4}{}}}(x+m_i)`$ $`=`$ $`\left[{\displaystyle \underset{a=1}{\overset{n_c2}{}}}(x\varphi _a)+i{\displaystyle \frac{\alpha _{n_c1}^2}{\mathrm{\Lambda }}}{\displaystyle \underset{a=1}{\overset{n_c3}{}}}(x\alpha _a)\right]\left[{\displaystyle \underset{a=1}{\overset{n_c2}{}}}(x\varphi _a)i{\displaystyle \frac{\alpha _{n_c1}^2}{\mathrm{\Lambda }}}{\displaystyle \underset{a=1}{\overset{n_c3}{}}}(x\alpha _a)\right].`$ For this identity to hold for any $`x`$, $`2n_c4`$ zeros of l.h.s. $`(m_i)`$ must coincide with $`n_c1`$ zeros in each of square brackets in the r.h.s. of the equation. Therefore, we divide up $`2n_c4=n_f`$ masses into to sets $`\{p_i,i=1,\mathrm{},n_c2\}`$ and $`\{q_i,i=1,\mathrm{},n_c2\}`$. This gives us $`{}_{n_f}{}^{}C_{n_f/2}^{}`$ choices. By identifying $`p_i`$ to the first square bracket and $`q_i`$ to the second square bracket, we find $`{\displaystyle \underset{i=1}{\overset{n_c2}{}}}(x+p_i)`$ $`=`$ $`{\displaystyle \underset{a=1}{\overset{n_c2}{}}}(x\varphi _a)+i{\displaystyle \frac{\alpha _{n_c1}^2}{\mathrm{\Lambda }}}{\displaystyle \underset{a=1}{\overset{n_c3}{}}}(x\alpha _a),`$ $`{\displaystyle \underset{i=1}{\overset{n_c1}{}}}(x+q_i)`$ $`=`$ $`{\displaystyle \underset{a=1}{\overset{n_c2}{}}}(x\varphi _a)i{\displaystyle \frac{\alpha _{n_c1}^2}{\mathrm{\Lambda }}}{\displaystyle \underset{a=1}{\overset{n_c3}{}}}(x\alpha _a).`$ (9.47) Expanding both sides in terms of symmetric polynomials, $`{\displaystyle \underset{k=0}{\overset{n_c2}{}}}s_k(p)x^{n_c2k}`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{n_c2}{}}}(1)^ks_k(\varphi )x^{n_ck2}+i{\displaystyle \frac{\alpha _{n_c1}^2}{\mathrm{\Lambda }}}{\displaystyle \underset{k=0}{\overset{n_c3}{}}}(1)^ks_k(\alpha )x^{n_ck2},`$ $`{\displaystyle \underset{k=0}{\overset{n_c1}{}}}s_k(q)x^{n_c1k}`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{n_c2}{}}}(1)^ks_k(\varphi )x^{n_ck2}i{\displaystyle \frac{\alpha _{n_c1}^2}{\mathrm{\Lambda }}}{\displaystyle \underset{k=0}{\overset{n_c3}{}}}(1)^ks_k(\alpha )x^{n_ck2}.`$ Finally, we find the solutions $$(1)^ks_k(p)=s_k(\varphi )i\frac{\alpha _{n_c1}^2}{\mathrm{\Lambda }}s_{k1}(\alpha ),(1)^ks_k(q)=s_k(\varphi )+i\frac{\alpha _{n_c1}^2}{\mathrm{\Lambda }}s_{k1}(\alpha ).$$ (9.49) Here and below, we use the notation that $`s_1=0`$ identically. By setting $`k=1`$ and subtracting both sides, we find $`\alpha _{n_c1}^2=i\mathrm{\Lambda }(s_1(p)s_1(q))/2=i\mathrm{\Lambda }_{k=1}^{n_f/2}(p_iq_i)/2`$. All $`s_k(\varphi )`$ and $`s_{k1}(\alpha )`$ are then given in terms of $`s_k(p)`$ and $`s_k(q)`$, and hence we can write order $`n_c3`$ ($`n_c2`$) polynomial equation for $`\alpha `$ ($`\varphi `$) which can always be solved to find all $`\alpha `$’s ($`\varphi `$’s). For smaller even $`n_f`$ with $`n_cn_f/2`$ even, the curve reduces to the one with $`2n_cn_f=4`$. To see this, we write the curve around the Chebyshev point as $$y^2=\underset{a=1}{\overset{n_f/2}{}}(x\varphi _a)^2\underset{k=1}{\overset{n_cn_f/2}{}}(x\varphi _k)^24\mathrm{\Lambda }^{2n_cn_f}\underset{i=1}{}(x+m_i),$$ (9.50) where $`\varphi _k=2\mathrm{\Lambda }\mathrm{cos}\pi (k1/2)/(n_cn_f/2)`$. Therefore, neglecting the fluctuations of $`\varphi _k`$ ($`k(n_cn_f/2+1)/2`$), the curve reduces to $$y^2=\underset{a=1}{\overset{n_f/2}{}}(x\varphi _a)^2\underset{k=1}{\overset{n_cn_f/2}{}}\varphi _k^24\mathrm{\Lambda }^{2n_cn_f}\underset{i=1}{}(x+m_i).$$ (9.51) The product of $`\varphi _k`$’s can be obtained as follows. Recalling $`T_N(x)=2^{N1}_{k=1}^N(xw_k)`$ with $`w_k=\mathrm{cos}\pi (k1/2)/N`$, we obtain for even $`N`$, $`T_N(0)=2^{N1}_k(w_k)`$, while $`T_N(x)=\mathrm{cos}(N\mathrm{arccos}x)`$ and hence $`T_N(0)=\mathrm{cos}N\pi /2=(1)^{N/2}`$. Therefore, $`_{k=1}^{n_cn_f/2}\varphi _k^2=(2T_N(0))^2\mathrm{\Lambda }^{2N}=4\mathrm{\Lambda }^{2N}`$. Therefore the low-energy curve is $$\frac{y^2}{4\mathrm{\Lambda }^{2n_cn_f}}=\underset{a=1}{\overset{n_f/2}{}}(x\varphi _a)^2\underset{i=1}{\overset{n_f/2}{}}(x+m_i).$$ (9.52) Note that this is precisely the curve of the $`SU(\frac{n_f}{2})`$ theory. It has the same form as the left hand side of Eq.(9.45). iv) $`r=\stackrel{~}{n}_c=n_fn_c:`$ Root of the baryonic branch The case with $`r=\stackrel{~}{n}_c=n_fn_c`$ also requires a separate consideration since the unperturbed curve has a special form, Eq.(6.18), Eq.(6.19). The curve is $$y^2=x^{2\stackrel{~}{n}_c}\left[\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}(x\mathrm{\Phi }_k)^2+4\mathrm{\Lambda }^{2n_cn_f}x^{n_c\stackrel{~}{n}_c}\right]=x^{2\stackrel{~}{n}_c}(x^{n_c\stackrel{~}{n}_c}\mathrm{\Lambda }^2)^2.$$ (9.53) with adjoint VEVS taken as $`\varphi _1,\mathrm{},\varphi _{n_c\stackrel{~}{n}_c}0`$, $`\varphi _{n_c\stackrel{~}{n}_c+1}=\mathrm{}=\varphi _{n_c}=0`$, $$(\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_k,\mathrm{},\mathrm{\Phi }_{n_c\stackrel{~}{n}_c})=\mathrm{\Lambda }(\omega ^2,\mathrm{},\omega ^{2k},\mathrm{},\omega ^{2(n_c\stackrel{~}{n}_c)}),$$ (9.54) where $`\omega =e^{\pi i/(n_c\stackrel{~}{n}_c)}`$. The double zeros of the factor in the parenthesis are at $$x=\mathrm{\Lambda }\omega ,\mathrm{\Lambda }\omega ^3,\mathrm{},\mathrm{\Lambda }\omega ^{2k1},\mathrm{},\mathrm{\Lambda }\omega ^{2(n_c\stackrel{~}{n}_c)1}.$$ (9.55) There are two ways for maintaining the curve maximally singular, when generic bare quark masses are added. One is to keep all “large” $`n_c\stackrel{~}{n}_c`$ zeros doubled, while allowing $`2\stackrel{~}{n}_c`$ zeros at $`x=0`$ to decompose into $`\stackrel{~}{n}_c1`$ double zeros and two single zeros. The other is to take all “small” zeros doubled, while keeping only $`n_c\stackrel{~}{n}_c1`$ “large” double zeros. iv-a) Keeping all of “large” zeros doubled Upon mass perturbation the adjoint scalar VEVS take the form, $`(\varphi _a,\mathrm{\Phi }_k=\mathrm{\Lambda }\omega ^{2k}+\gamma _k),`$ with $`\gamma _k,\varphi _aO(m)`$. The constraint is therefore $$\underset{a=1}{\overset{\stackrel{~}{n}_c}{}}\varphi _a+\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}\gamma _k=0.$$ (9.56) The perturbed curve is $$y^2=\underset{a=1}{\overset{\stackrel{~}{n}_c}{}}(x\varphi _a)^2\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}(x\mathrm{\Lambda }\omega ^{2k}\gamma _k)^2+4\mathrm{\Lambda }^{2n_cn_f}\underset{i=1}{\overset{n_f}{}}(x+m_i).$$ (9.57) The zeros (9.55) are also shifted to $$x=\mathrm{\Lambda }\omega ^{2\mathrm{}1}+\delta _{\mathrm{}},\mathrm{}=1,\mathrm{},n_c\stackrel{~}{n}_c.$$ (9.58) We substitute these zeros for each $`\mathrm{}`$ into the curve and require that the r.h.s. of the curve (9.57) vanishes at $`O(m)`$. The first factor in the curve (9.57) is expanded as $$\underset{a=1}{\overset{\stackrel{~}{n}_c}{}}(\mathrm{\Lambda }\omega ^{2\mathrm{}1}+\delta _{\mathrm{}}\varphi _a)^2=(\mathrm{\Lambda }\omega ^{2\mathrm{}1})^{\stackrel{~}{n}_c}\left[1+2\underset{a=1}{\overset{\stackrel{~}{n}_c}{}}\frac{\delta _{\mathrm{}}\varphi _a}{\mathrm{\Lambda }\omega ^{2\mathrm{}1}}\right].$$ (9.59) The second factor in the curve (9.57) is expanded as $$\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}(\mathrm{\Lambda }\omega ^{2\mathrm{}1}+\delta _{\mathrm{}}\mathrm{\Lambda }\omega ^{2k}\gamma _k)^2=\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}(\mathrm{\Lambda }\omega ^{2\mathrm{}1}\mathrm{\Lambda }\omega ^{2k})^2\left[1+2\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}\frac{\delta _{\mathrm{}}\gamma _k}{\mathrm{\Lambda }\omega ^{2\mathrm{}1}\mathrm{\Lambda }\omega ^{2k}}\right].$$ (9.60) This factor needs to be simplified. We show that $$\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}(\mathrm{\Lambda }\omega ^{2\mathrm{}1}\mathrm{\Lambda }\omega ^{2k})=2\mathrm{\Lambda }^{n_c\stackrel{~}{n}_c}.$$ (9.61) This can be proven by studying the polynomial $$\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}(t\mathrm{\Lambda }\omega ^{2k})=t^{n_c\stackrel{~}{n}_c}\mathrm{\Lambda }^{n_c\stackrel{~}{n}_c}.$$ (9.62) By substituting $`t=\mathrm{\Lambda }\omega ^{2\mathrm{}1}`$, and noting that $`(\mathrm{\Lambda }\omega ^{2\mathrm{}1})^{n_c\stackrel{~}{n}_c}=\mathrm{\Lambda }^{n_c\stackrel{~}{n}_c}e^{\pi i(2\mathrm{}1)}=\mathrm{\Lambda }^{n_c\stackrel{~}{n}_c}`$. The second factor (9.60) of the curve (9.57) is now simplified to $$4\mathrm{\Lambda }^{2(n_c\stackrel{~}{n}_c)}\left[1+2\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}\frac{\delta _{\mathrm{}}\gamma _k}{\mathrm{\Lambda }\omega ^{2\mathrm{}1}\mathrm{\Lambda }\omega ^{2k}}\right].$$ (9.63) The last term in the curve (9.57) is expanded as $$4\mathrm{\Lambda }^{2n_cn_f}\underset{i=1}{\overset{n_f}{}}(\mathrm{\Lambda }\omega ^{2\mathrm{}1}+\delta _{\mathrm{}}+m_i)=4\mathrm{\Lambda }^{2n_c}(\omega ^{2\mathrm{}1})^{n_f}\left[1+\underset{i=1}{\overset{n_f}{}}\frac{\delta _{\mathrm{}}+m_i}{\mathrm{\Lambda }\omega ^{2\mathrm{}1}}\right].$$ (9.64) Note further that $$(\omega ^{2\mathrm{}1})^{n_f}=(\omega ^{2\mathrm{}1})^{2\stackrel{~}{n}_c+(n_c\stackrel{~}{n}_c)}=(\omega ^{2\mathrm{}1})^{2\stackrel{~}{n}_c}e^{\pi i(2\mathrm{}1)}=(\omega ^{2\mathrm{}1})^{2\stackrel{~}{n}_c}.$$ (9.65) The last term therefore is $$4\mathrm{\Lambda }^{2n_c}(\omega ^{2\mathrm{}1})^{2\stackrel{~}{n}_c}\left[1+\underset{i=1}{\overset{n_f}{}}\frac{\delta _{\mathrm{}}+m_i}{\mathrm{\Lambda }\omega ^{2\mathrm{}1}}\right].$$ (9.66) By putting together the expansion Eqs.(9.59,9.63,9.66) up to $`O(m)`$ into the curve (9.57) and requiring it to vanish, we find $$0=2\underset{a=1}{\overset{\stackrel{~}{n}_c}{}}\frac{\delta _{\mathrm{}}\varphi _a}{\mathrm{\Lambda }\omega ^{2\mathrm{}1}}+2\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}\frac{\delta _{\mathrm{}}\gamma _k}{\mathrm{\Lambda }\omega ^{2\mathrm{}1}\mathrm{\Lambda }\omega ^{2k}}\underset{i=1}{\overset{n_f}{}}\frac{\delta _{\mathrm{}}+m_i}{\mathrm{\Lambda }\omega ^{2\mathrm{}1}}.$$ (9.67) Multiply the equation by $`\mathrm{\Lambda }\omega ^{2\mathrm{}1}`$ to simplify it to $$0=2\underset{a=1}{\overset{\stackrel{~}{n}_c}{}}(\delta _{\mathrm{}}\varphi _a)+2\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}\frac{\delta _{\mathrm{}}\gamma _k}{1\omega ^{2k2\mathrm{}+1}}\underset{i=1}{\overset{n_f}{}}(\delta _{\mathrm{}}+m_i).$$ (9.68) We try to further simplify this equation. The sum over $`a`$ in the first term gives $$2\underset{a=1}{\overset{\stackrel{~}{n}_c}{}}(\delta _{\mathrm{}}\varphi _a)=2\stackrel{~}{n}_c\delta _{\mathrm{}}2\underset{a=1}{\overset{\stackrel{~}{n}_c}{}}\varphi _a=2\stackrel{~}{n}_c\delta _{\mathrm{}}+2\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}\gamma _k,$$ (9.69) where the constraint Eq. (9.56) was used in the last equality. To simplify the second term, we would like to prove that $$\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}\frac{1}{1\omega ^{2k2\mathrm{}+1}}=\frac{1}{2}(n_c\stackrel{~}{n}_c).$$ (9.70) Note that the sum over $`k`$ exhausts all possible odd powers in $`\omega `$ in the denominator. Therefore we can shift $`k`$ to $`k+\mathrm{}`$ and find $$\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}\frac{1}{1\omega ^{2k2\mathrm{}+1}}=\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}\frac{1}{1\omega ^{2k+1}}.$$ (9.71) Now we distinguish two cases, when $`n_c\stackrel{~}{n}_c=2m`$ (even), and $`n_c\stackrel{~}{n}_c=2m+1`$ (odd). When $`n_c\stackrel{~}{n}_c=2m`$, the sum is separated in two pieces, $$\underset{k=1}{\overset{2m}{}}\frac{1}{1\omega ^{2k+1}}=\underset{k=1}{\overset{m}{}}\frac{1}{1\omega ^{2k+1}}+\underset{k=m+1}{\overset{2m}{}}\frac{1}{1\omega ^{2k+1}}.$$ (9.72) Change the variable $`k`$ in the second sum to $`2mk+1`$, and using the definition $`\omega =e^{\pi i/2m}`$, $$=\underset{k=1}{\overset{m}{}}\frac{1}{1\omega ^{2k+1}}+\underset{k=1}{\overset{m}{}}\frac{1}{1\omega ^{(2k+1)}}.$$ (9.73) Adding terms for each $`m`$, $$=\underset{k=1}{\overset{m}{}}\frac{1\omega ^{2k+1}+1\omega ^{(2k+1)}}{1\omega ^{2k+1}\omega ^{(2k+1)}+1}=m=\frac{1}{2}(n_c\stackrel{~}{n}_c).$$ (9.74) This completes the proof of Eq. (9.70) for even $`n_c\stackrel{~}{n}_c`$. Similarly, for $`n_c\stackrel{~}{n}_c=2m+1`$, $$\underset{k=1}{\overset{2m}{}}\frac{1}{1\omega ^{2k+1}}=\underset{k=1}{\overset{m}{}}\frac{1}{1\omega ^{2k+1}}+\frac{1}{1\omega ^{2m+1}}+\underset{k=m+2}{\overset{2m+1}{}}\frac{1}{1\omega ^{2k+1}}.$$ (9.75) The middle term in the r.h.s. is $`1/(1(1))=1/2`$. Change the variable $`k`$ in the last sum to $`2m+2k`$, $$=\underset{k=1}{\overset{m}{}}\frac{1}{1\omega ^{2k+1}}+\frac{1}{2}+\underset{k=1}{\overset{m}{}}\frac{1}{1\omega ^{(2k+1)}}.$$ (9.76) Again by adding terms for each $`m`$, $$=\frac{1}{2}+\underset{k=1}{\overset{m}{}}\frac{1\omega ^{2k+1}+1\omega ^{(2k+1)}}{1\omega ^{2k+1}\omega ^{(2k+1)}+1}=\frac{1}{2}+m=\frac{1}{2}(n_c\stackrel{~}{n}_c).$$ (9.77) This completes the proof of Eq. (9.70) for odd $`n_c\stackrel{~}{n}_c`$. Therefore the identity (9.70) is proven. The second term in Eq. (9.68) is then given by $$2\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}\frac{\delta _{\mathrm{}}\gamma _k}{1\omega ^{2k2\mathrm{}+1}}=(n_c\stackrel{~}{n}_c)\delta _{\mathrm{}}2\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}\frac{\gamma _k}{1\omega ^{2k2\mathrm{}+1}}.$$ (9.78) Finally the last term in Eq. (9.68) is simply $$\underset{i=1}{\overset{n_f}{}}(\delta _{\mathrm{}}+m_i)=n_f\delta _{\mathrm{}}\underset{i=1}{\overset{n_f}{}}m_i.$$ (9.79) Putting together Eqs. (9.69,9.78,9.79) in (9.68), we find $`0`$ $`=`$ $`2\stackrel{~}{n}_c\delta _{\mathrm{}}+2{\displaystyle \underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}}\gamma _k+(n_c\stackrel{~}{n}_c)\delta _{\mathrm{}}2{\displaystyle \underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}}{\displaystyle \frac{\gamma _k}{1\omega ^{2k2\mathrm{}+1}}}n_f\delta _{\mathrm{}}{\displaystyle \underset{i=1}{\overset{n_f}{}}}m_i`$ (9.80) $`=`$ $`2{\displaystyle \underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}}\gamma _k2{\displaystyle \underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}}{\displaystyle \frac{\gamma _k}{1\omega ^{2k2\mathrm{}+1}}}{\displaystyle \underset{i=1}{\overset{n_f}{}}}m_i,`$ and $`\delta _{\mathrm{}}`$ disappeared from the equation. Finally, we add over $`\mathrm{}`$. Only the second term depends on $`\mathrm{}`$, and the sum over $`\mathrm{}`$ is simplified again by using the identity Eq. (9.70). We find $$0=2(n_c\stackrel{~}{n}_c)\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}\gamma _k2\frac{1}{2}(n_c\stackrel{~}{n}_c)\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}\gamma _k(n_c\stackrel{~}{n}_c)\underset{i=1}{\overset{n_f}{}}m_i,$$ (9.81) and therefore $$\underset{a=1}{\overset{\stackrel{~}{n}_c}{}}\varphi _a=\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}\gamma _k=\underset{i=1}{\overset{n_f}{}}m_i.$$ (9.82) The fluctuation around $`x0`$ is described by a low-energy curve by neglecting $`x,\gamma _k\mathrm{\Lambda }`$ in Eq. (9.57), $$y^2=\underset{a=1}{\overset{\stackrel{~}{n}_c}{}}(x\varphi _a)^2\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}(\mathrm{\Lambda }\omega ^{2k})^2+4\mathrm{\Lambda }^{2n_cn_f}\underset{i=1}{\overset{n_f}{}}(x+m_i),$$ (9.83) or $$\frac{y^2}{\mathrm{\Lambda }^{2(n_c\stackrel{~}{n}_c)}}=\underset{a=1}{\overset{\stackrel{~}{n}_c}{}}(x\varphi _a)^2+\frac{4}{\mathrm{\Lambda }^{n_c\stackrel{~}{n}_c}}\underset{i=1}{\overset{n_f}{}}(x+m_i),$$ (9.84) where $`\varphi _a`$ are subject to an unusual constraint Eq. (9.82). Apart from the constraint, the problem is similar to earlier study of the $`r`$-root, where we take $`\varphi _a=m_a`$ to have double zeros at $`x=m_a`$. The maximum number of $`\varphi _a`$’s that can be chosen this way is, however, $`\stackrel{~}{n}_c1`$ because of the constraint. Let us choose $`r`$ out of $`\stackrel{~}{n}_c`$ to coincide with $`r`$ of $`m_i`$ ($`{}_{n_f}{}^{}C_{r}^{}`$ choices), e.g., $$\varphi _1=m_1,\mathrm{},\varphi _r=m_r,$$ (9.85) While the remaining $`\varphi _a`$ are $$\varphi _k=\frac{1}{\stackrel{~}{n}_cr}\underset{i=r+1}{\overset{n_f}{}}m_i+\delta \varphi _k,k=1,\mathrm{},\stackrel{~}{n}_cr.$$ (9.86) The fluctuations are subject to the constraint $$\underset{k=1}{\overset{\stackrel{~}{n}_cr}{}}\delta \varphi _k=0.$$ (9.87) The low-energy curve Eq. (9.84) then becomes $$\frac{y^2}{\mathrm{\Lambda }^{2(n_c\stackrel{~}{n}_c)}}=\underset{a=1}{\overset{r}{}}(xm_a)^2\underset{k=1}{\overset{\stackrel{~}{n}_cr}{}}\left(x+\frac{1}{\stackrel{~}{n}_cr}\underset{i=r+1}{\overset{n_f}{}}m_i\delta \varphi _k\right)^2+\frac{4}{\mathrm{\Lambda }^{n_c\stackrel{~}{n}_c}}\underset{i=1}{\overset{n_f}{}}(x+m_i).$$ (9.88) Now shift $`x`$ to $`x\frac{1}{\stackrel{~}{n}_cr}_{i=r+1}^{n_f}m_i`$, and we assume that the remaining fluctuations $`x,\delta \varphi _km`$ which will be justified a posteriori. The curve is further approximated as $`{\displaystyle \frac{y^2}{\mathrm{\Lambda }^{2(n_c\stackrel{~}{n}_c)}}}`$ $`=`$ $`{\displaystyle \underset{a=1}{\overset{r}{}}}\left({\displaystyle \frac{1}{\stackrel{~}{n}_cr}}{\displaystyle \underset{i=r+1}{\overset{n_f}{}}}m_i+m_a\right)^2{\displaystyle \underset{k=1}{\overset{\stackrel{~}{n}_cr}{}}}(x\delta \varphi _k)^2`$ (9.89) $`+{\displaystyle \frac{4}{\mathrm{\Lambda }^{n_c\stackrel{~}{n}_c}}}{\displaystyle \underset{i=1}{\overset{n_f}{}}}\left(m_i{\displaystyle \frac{1}{\stackrel{~}{n}_cr}}{\displaystyle \underset{i=r+1}{\overset{n_f}{}}}m_i\right).`$ Up to an overall constant, this is the curve of pure $`SU(\stackrel{~}{n}_cr)`$ Yang–Mills theories. The dynamical scale of this theory is $`\mathrm{\Lambda }_{\mathrm{𝑝𝑢𝑟𝑒}}^{2(\stackrel{~}{n}_cr)}`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Lambda }^{n_c\stackrel{~}{n}_c}}}{\displaystyle \frac{_{i=1}^{n_f}\left(m_i\frac{1}{\stackrel{~}{n}_cr}_{i=r+1}^{n_f}m_i\right)}{_{a=1}^r\left(\frac{1}{\stackrel{~}{n}_cr}_{i=r+1}^{n_f}m_i+m_a\right)^2}}`$ (9.90) $``$ $`m^{2(\stackrel{~}{n}_cr)}\left({\displaystyle \frac{m}{\mathrm{\Lambda }}}\right)^{n_c\stackrel{~}{n}_c}m^{2(\stackrel{~}{n}_cr)}.`$ Therefore the singularities of this curve are located at $`x,\delta \varphi _k\mathrm{\Lambda }_{\mathrm{𝑝𝑢𝑟𝑒}}m`$ and hence the approximation is justified. There are $`\stackrel{~}{n}_cr`$ such singularities (basically the Witten index) for each $`r`$, and hence the total number of vacua obtained in this subsection are $$𝒩_2=\underset{r=0}{\overset{\stackrel{~}{n}_c1}{}}(\stackrel{~}{n}_cr){}_{n_f}{}^{}C_{r}^{}.$$ (9.91) iv-b) Keeping all of “small” double zeros Another possibility is that all $`\stackrel{~}{n}_c`$ “small” zeros near $`x=0`$ are doubled, while one of the $`n_c\stackrel{~}{n}_c`$ “large” double zeros is split into two single zeros. The crucial difference from the previous case is that we keep one less “large” double zeros and hence fluctuations $`\varphi _1,\mathrm{},\varphi _{\stackrel{~}{n}_c}`$ are not subject to a constraint. Even though the low-energy curve $$\frac{y^2}{\mathrm{\Lambda }^{2(n_c\stackrel{~}{n}_c)}}=\underset{a=1}{\overset{\stackrel{~}{n}_c}{}}(x\varphi _a)^2+\frac{4}{\mathrm{\Lambda }^{n_c\stackrel{~}{n}_c}}\underset{i=1}{\overset{n_f}{}}(x+m_i),$$ (9.92) is the same as Eq. (9.84) in the previous subsection, $`\varphi _a`$ can all freely vary and hence can all be matched to the quark masses as in the cases i) above, as $$\varphi _1=m_1,\mathrm{},\varphi _{\stackrel{~}{n}_c}=m_{\stackrel{~}{n}_c},$$ (9.93) and there are $`{}_{n_f}{}^{}C_{\stackrel{~}{n}_c}^{}`$ choices. Recall that there are $`n_c\stackrel{~}{n}_c`$ “large” double zeros at the baryonic root and we can give up one of them; therefore there are actually $`(n_c\stackrel{~}{n}_c){}_{n_f}{}^{}C_{\stackrel{~}{n}_c}^{}`$ choices. v) Summary of the vacuum counting in $`SU(n_c)`$ theories The “nonbaryonic” branch roots Eq.(9.2), Eq.(9.29), yield, upon quark mass perturbation, $$(2n_cn_f)2^{n_f1}(n_c\stackrel{~}{n}_c){}_{n_f}{}^{}C_{\stackrel{~}{n}_c}^{}=𝒩_1(n_c\stackrel{~}{n}_c){}_{n_f}{}^{}C_{\stackrel{~}{n}_c}^{}$$ (9.94) vacua. The baryonic root, Eq.(9.53), Eq.(9.54), leads to $$\underset{r=0}{\overset{\stackrel{~}{n}_c1}{}}(\stackrel{~}{n}_cr){}_{n_f}{}^{}C_{r}^{}+(n_c\stackrel{~}{n}_c){}_{n_f}{}^{}C_{\stackrel{~}{n}_c}^{}=𝒩_2+(n_c\stackrel{~}{n}_c){}_{n_f}{}^{}C_{\stackrel{~}{n}_c}^{}$$ (9.95) vacua. Their sum coincides with the total number of $`N=1`$ vacua found from the semiclassical analysis as well as from the analyses at large $`\mu `$. ### 9.2 Perturbation around CFT points of the $`USp(2n_c)`$ Curve We start from the CFT points described by the Chebyshev polynomial, Eq.(6.26)-Eq.(6.29), and add generic quark masses $`m_i`$. i) Chebyshev point: Odd $`n_f`$ Let us take $`n_f`$ odd first. We first study the specific case of $`2n_c=n_f1`$. The more general cases will be discussed later on. The Chebyshev solution is obtained in the massless limit by setting all of $`\varphi _a=0`$: $$xy^2=\left[x^{n_c+1}\right]^24\mathrm{\Lambda }^2x^{2n_c+1}=x^{2n_c+1}(x4\mathrm{\Lambda }^2).$$ (9.96) The zero at $`x=0`$ is of degree $`2n_c`$, and there is another isolated zero at $`x=4\mathrm{\Lambda }^2`$. There is also a branch point at $`x=\mathrm{}`$. Under the perturbation by generic quark masses, we go back to the original curve. The only way that the curve can be arranged to have $`n_c`$ double zeros as $$xy^2=x(x4\mathrm{\Lambda }^2\beta )\underset{a=1}{\overset{n_c}{}}(x\alpha _a)^2$$ (9.97) is by assuming $$\alpha _am^2,\varphi _1^2m\mathrm{\Lambda },\varphi _2^2\mathrm{}\varphi _{n_c}^2m^2.$$ (9.98) By neglecting $`\beta ,x\mathrm{\Lambda }^2`$, we need to solve the equation $$\left[x\underset{a=1}{\overset{n_c}{}}(x\varphi _a^2)+2\mathrm{\Lambda }m_1\mathrm{}m_{n_f}\right]^24\mathrm{\Lambda }^2\underset{i=1}{\overset{n_f}{}}(x+m_i^2)=4\mathrm{\Lambda }^2x\underset{a=1}{\overset{n_c}{}}(x\alpha _a)^2.$$ (9.99) To address this question, we introduce a few preliminary facts. First of all, consider a generic polynomial $$\underset{i=1}{\overset{N}{}}(z\rho _i)=\underset{k=0}{\overset{N}{}}(1)^ks_k(\rho )z^{Nk}.$$ (9.100) The symmetric polynomials $`s_j(\rho )`$ are given by $`s_0(\rho )=1`$, $`s_1(\rho )=_{i=1}^N\rho _i`$, $`s_2(\rho )=_{i<j}\rho _i\rho _j`$, etc. This defines the notation $`s_k`$. Then Eq. (9.99) is written as $$\left[x\underset{k=0}{\overset{n_c}{}}(1)^ks_k(\varphi ^2)x^{n_ck}+2\mathrm{\Lambda }m_1\mathrm{}m_{n_f}\right]^24\mathrm{\Lambda }^2\underset{i=1}{\overset{n_f}{}}(x+m_i^2)=4\mathrm{\Lambda }^2x\underset{a=1}{\overset{n_c}{}}(x\alpha _a)^2.$$ (9.101) Note that $`s_k(\varphi ^2)=O(m^{2k1}\mathrm{\Lambda })`$ for $`k0`$ because one of the $`\varphi ^2`$’s is $`O(m\mathrm{\Lambda })`$. This allows us to neglect $`s_0(\varphi ^2)x^{n_c}=x^{n_c}`$ term in the sum, and the equation becomes $$\left[\underset{k=1}{\overset{n_c}{}}(1)^ks_k(\varphi ^2)x^{n_ck+1}+2\mathrm{\Lambda }m_1\mathrm{}m_{n_f}\right]^24\mathrm{\Lambda }^2\underset{i=1}{\overset{n_f}{}}(x+m_i^2)=4\mathrm{\Lambda }^2x\underset{a=1}{\overset{n_c}{}}(x\alpha _a)^2.$$ (9.102) Now rewrite it as $`{\displaystyle \underset{i=1}{\overset{n_f}{}}}(x+m_i^2)`$ $`=`$ $`x\left[{\displaystyle \underset{k=0}{\overset{n_c}{}}}(1)^ks_k(\alpha )x^{n_ck}\right]^2+\left[{\displaystyle \underset{k=1}{\overset{n_c}{}}}(1)^k{\displaystyle \frac{s_k(\varphi ^2)}{2\mathrm{\Lambda }}}x^{n_ck+1}+m_1\mathrm{}m_{n_f}\right]^2`$ $`=`$ $`x\left[{\displaystyle \underset{k=0}{\overset{n_c}{}}}(1)^ks_k(\alpha )x^{n_ck}\right]^2+\left[{\displaystyle \underset{k=0}{\overset{n_c1}{}}}(1)^k{\displaystyle \frac{s_{k+1}(\varphi ^2)}{2\mathrm{\Lambda }}}x^{n_ck}+m_1\mathrm{}m_{n_f}\right]^2.`$ Now consider the following polynomial $$F(x)=\underset{i=1}{\overset{n_f}{}}(\sqrt{x}+im_i)=\underset{k=0}{\overset{n_f}{}}i^ks_k(m)\sqrt{x}^{n_fk}.$$ (9.104) This polynomial can be divided into the “real” and “imaginary” parts (this is not strictly true because $`m`$’s are complex, but what is meant here is the division between terms of odd powers in $`i`$ and of even powers in $`i`$), $`F(x)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{(n_f1)/2}{}}}(1)^ks_{2k}(m)\sqrt{x}^{n_f2k}+{\displaystyle \underset{k=0}{\overset{(n_f1)/2}{}}}i(1)^ks_{2k+1}(m)\sqrt{x}^{n_f2k1}`$ (9.105) $`=`$ $`\sqrt{x}{\displaystyle \underset{k=0}{\overset{n_c}{}}}(1)^ks_{2k}(m)x^{n_ck}+i{\displaystyle \underset{k=0}{\overset{n_c}{}}}(1)^ks_{2k+1}(m)x^{n_ck},`$ where we used $`n_f=2n_c+1`$. Similarly consider the polynomial $$G(x)=\underset{i=1}{\overset{n_f}{}}(\sqrt{x}im_i)=\sqrt{x}\underset{k=0}{\overset{n_c}{}}(1)^ks_{2k}(m)x^{n_ck}i\underset{k=0}{\overset{n_c}{}}(1)^ks_{2k+1}(m)x^{n_ck}.$$ (9.106) From the definition, $$F(x)G(x)=\underset{i=1}{\overset{n_f}{}}(x+m_i^2).$$ (9.107) On the other hand, this product is also given by $$F(x)G(x)=\left[\sqrt{x}\underset{k=0}{\overset{n_c}{}}(1)^ks_{2k}(m)x^{n_ck}\right]^2+\left[\underset{k=0}{\overset{n_c}{}}(1)^{n_fk}s_{2k+1}(m)x^k\right]^2.$$ (9.108) This is precisely the same as Eq. (LABEL:eq:tosolve2) upon identifications $$s_k(\alpha )=s_{2k}(m),s_{k+1}(\varphi ^2)=2\mathrm{\Lambda }(1)^{n_c}s_{2k+1}(m).$$ (9.109) In the last identification, we used the fact that $`s_{2n_c+1}(m)=s_{n_f}(m)=m_1\mathrm{}m_{n_f}`$. This gives explicit solutions to the vacuum. Once we have this solution, however, we can obtain other $`2^{n_f1}1`$ solutions as follows. First note that the curve Eq. (6.24) is invariant under changing signs of even number of masses. Therefore, we can change signs of even number of masses from the solution (9.109). This gives $`2^{n_f1}`$ solutions in total agreeing with $`𝒩_1`$ in the large $`\mu `$ analysis. This solution therefore decomposes under $`U(n_f)`$ for the equal mass case as $`2^{n_f1}={}_{n_f}{}^{}C_{0}^{}+{}_{n_f}{}^{}C_{2}^{}+\mathrm{}{}_{n_f}{}^{}C_{n_f1}^{}`$, reminiscent of the spinor representation. The case of equal mass deserves further comments. As noted above, we can flip the signs of even number of quark masses, and therefore a general situation has $`2r`$ negative masses $`m`$ and $`n_f2r`$ positive masses $`m`$. Let us study the location of branch points in this situation. The curve studied above has the form $$y^2=4\mathrm{\Lambda }^2\underset{a=1}{\overset{n_c}{}}(x\alpha _a)^2$$ (9.110) near $`x0`$. Given the solutions Eq. (9.109), we can write $`2\mathrm{\Lambda }{\displaystyle \underset{a=1}{\overset{n_c}{}}}(x\alpha _a)=2\mathrm{\Lambda }{\displaystyle \underset{k=0}{\overset{n_c}{}}}(1)^ks_k(\alpha )x^{n_ck}`$ $`=2\mathrm{\Lambda }{\displaystyle \underset{k=0}{\overset{n_c}{}}}(1)^ks_{2k}(m)x^{n_ck}`$ $`=2\mathrm{\Lambda }{\displaystyle \frac{1}{2\sqrt{x}}}\left[{\displaystyle \underset{i=1}{\overset{n_f}{}}}(\sqrt{x}+im_i)+{\displaystyle \underset{i=1}{\overset{n_f}{}}}(\sqrt{x}im_i)\right]`$ $`=\mathrm{\Lambda }{\displaystyle \frac{1}{\sqrt{x}}}\left[(\sqrt{x}+im_i)^{n_f2r}(\sqrt{x}im_i)^{2r}+(\sqrt{x}im_i)^{n_f2r}(\sqrt{x}+im_i)^{2r}\right].`$ Note that $`1/\sqrt{x}`$ does not introduce a singularity at $`x=0`$. If $`4r<n_f`$, we can factor $`(x+m^2)^{2r}`$ from above and hence $`(x+m^2)^{4r}`$ from the curve. We interpret this factor as the emergence of $`SU(2r)`$ gauge theory. If $`4r>n_f`$, we can factor $`(x+m^2)^{n_f2r}`$ from above and hence $`(x+m^2)^{2n_f4r}`$ from the curve. We interpret this factor as the emergence of $`SU(n_f2r)=SU(2(n_cr)+1)`$ gauge theory. Combining both, we see that gauge groups up to $`SU((n_f1)/2)`$ are possible. This fact can also be understood from the Higgs branch picture. When quark masses are large and equal, one can cancel quark masses by the adjoint VEVS classically (squark singularity) and obtain $`U(k)`$ gauge theories $`k=0,1,\mathrm{},(n_f1)/2`$ depending on how many components of the adjoint field is used to cancel the quark masses. These are all infrared-free theories and the gauge fields survive the quantum effects. While smoothly decreasing the quark masses, the Higgs branch emanating from the squark singularity does not change due to the non-renormalization theorem and hence the $`U(k)`$ theories survive to the small mass studied using the Coulomb branch. The effective low-energy Lagrangian which describes physics around the squark singularity is therefore nothing but that of Argyres–Plesser–Seiberg for $`SU(n_c)`$ theories at the non-baryonic roots . We have shown (Section 8.1) that these theories indeed produce $`{}_{n_f}{}^{}C_{k}^{}`$ vacua upon $`\mu 0`$ and generic quark mass perturbation. Therefore the whole picture nicely fits together. On the other hand, the strictly massless case has the high singularity $`x^{2n_c}`$ and appears to give a new superconformal theory with a global symmetry $`SO(2n_f)`$. We do not know of a weakly coupled description of this singularity in terms of a local field theory. We have checked that this singularity indeed produces mutually non-local dyons for $`n_c=2`$ and $`n_f=5`$. Now we come back to the case of smaller $`n_f`$. Around the Chebyshev solution in the presence of quark masses, the curve is $`xy^2`$ $`=`$ $`\left[x{\displaystyle \underset{a=1}{\overset{(n_f1)/2}{}}}(x\varphi _a^2){\displaystyle \underset{k=1}{\overset{n_c(n_f1)/2}{}}}(x\varphi _k^2\kappa _k^2)+2\mathrm{\Lambda }^{2n_c+2n_f}m_1\mathrm{}m_{n_f}\right]^2`$ (9.112) $`4\mathrm{\Lambda }^{2(2n_c+2n_f)}{\displaystyle \underset{i=1}{\overset{n_f}{}}}(x+m_i^2),`$ where $`\varphi _k^2`$ are given by those for the Chebyshev polynomial. However, for the purpose of studying the behavior around $`x=0`$, both $`x`$ and the shifts $`\kappa _k^2`$ can be ignored relative to $`\varphi _k^2`$. We need to know $`_{k=1}^{n_c(n_f1)/2}(\varphi _k^2)`$. This can be calculated as $`{\displaystyle \underset{k=1}{\overset{n_c(n_f1)/2}{}}}(\varphi _k^2)`$ $`=`$ $`{\displaystyle \frac{2\mathrm{\Lambda }^{2n_c+1n_f}}{\sqrt{x}}}T_{2n_c+2n_f}\left({\displaystyle \frac{\sqrt{x}}{2\mathrm{\Lambda }}}\right)\|_{x0}`$ (9.113) $`=`$ $`\mathrm{\Lambda }^{2n_c+1n_f}\underset{t0}{lim}{\displaystyle \frac{1}{t}}\mathrm{cos}\left((2n_c+2n_f)\mathrm{arccos}t\right)`$ $`=`$ $`(1)^{n_c(n_f1)/2}(2n_c+2n_f)\mathrm{\Lambda }^{2n_c+1n_f}.`$ Then the curve is approximated as $`{\displaystyle \frac{xy^2}{\mathrm{\Lambda }^{2(2n_c+1n_f)}}}`$ $`=`$ $`\left[(1)^{n_c(n_f1)/2}(2n_c+2n_f)x{\displaystyle \underset{a=1}{\overset{(n_f1)/2}{}}}(x\varphi _a^2)+2\mathrm{\Lambda }m_1\mathrm{}m_{n_f}\right]^2`$ (9.114) $`4\mathrm{\Lambda }^2{\displaystyle \underset{i=1}{\overset{n_f}{}}}(x+m_i^2).`$ This is nothing but the curve of $`USp(2n_c^{})`$ theory with $`n_c^{}=(n_f1)/2`$ upon changing normalizations of $`x`$, $`y`$, $`\varphi _a^2`$. The rest of the analysis therefore follows exactly the same as in $`2n_c=n_f1`$ case. Even when other Chebyshev solutions obtained by $`Z_{2n_c+2n_f}`$ are used, they simply amount to the change of phase of $`\mathrm{\Lambda }`$ in the above approximate curve and the analysis remains the same. ii) Chebyshev point: even $`n_f`$ Consider now even $`n_f`$ cases. Again, let us study the specific case of $`n_f=2n_c`$ first. We shall come back to the more general cases later on. The Chebyshev solution is obtained in the massless limit by setting all but one of $`\varphi _a=0`$: $$xy^2=\left[x^{n_c}(x\varphi _{n_c}^2)\right]^24\mathrm{\Lambda }^4x^{2n_c}=x^{2n_c}(x\varphi _{n_c}^22\mathrm{\Lambda }^2)(x\varphi _{n_c}^2+2\mathrm{\Lambda }^2).$$ (9.115) We take $`\varphi _{n_c}^2=\pm 2\mathrm{\Lambda }^2`$ so that the zero at $`x=0`$ has degree $`2n_c`$. There is another isolated zero at $`x=\pm 4\mathrm{\Lambda }^2`$. There is also a branch point at $`x=\mathrm{}`$. We first consider the case $`\varphi _{n_c}^2=+2\mathrm{\Lambda }^2`$ and will come back to the case $`\varphi _{n_c}^2=2\mathrm{\Lambda }^2`$ later on. Under the perturbation by generic quark masses, we go back to the original curve. The only way that the curve $$xy^2=\left[x\underset{a=1}{\overset{n_c1}{}}(x\varphi _a^2)(x2\mathrm{\Lambda }^2\beta )+2\mathrm{\Lambda }^2m_1\mathrm{}m_{n_f}\right]^24\mathrm{\Lambda }^4\underset{i=1}{\overset{n_f}{}}(x+m_i^2)$$ (9.116) can be arranged to have $`n_c`$ double zeros as $$xy^2=x(x4\mathrm{\Lambda }^2\gamma )\underset{a=1}{\overset{n_c}{}}(x\alpha _a)^2,$$ (9.117) is by assuming $$\varphi _am^2,\alpha _{n_c}m\mathrm{\Lambda },\alpha _1^2\mathrm{}\alpha _{n_c1}^2m^2.$$ (9.118) By neglecting $`\beta ,\gamma ,x\mathrm{\Lambda }^2`$, we must solve the equation $`\left[2\mathrm{\Lambda }^2x{\displaystyle \underset{a=1}{\overset{n_c1}{}}}(x\varphi _a^2)+2\mathrm{\Lambda }^2m_1\mathrm{}m_{n_f}\right]^24\mathrm{\Lambda }^4{\displaystyle \underset{i=1}{\overset{n_f}{}}}(x+m_i^2)`$ (9.119) $`=4\mathrm{\Lambda }^2x{\displaystyle \underset{a=1}{\overset{n_c}{}}}(x\alpha _a)^2=4\mathrm{\Lambda }^2\alpha _{n_c}^2x{\displaystyle \underset{a=1}{\overset{n_c1}{}}}(x\alpha _a)^2.`$ It is interesting to note that this is the curve for the superconformal $`USp(n_f2)`$ theory with $`n_f`$ flavors with a special choice of $`g=1`$. This can be rewritten as $$\left[x\underset{k=0}{\overset{n_c1}{}}(1)^ks_k(\varphi ^2)x^{n_c1k}+m_1\mathrm{}m_{n_f}\right]^2\underset{i=1}{\overset{n_f}{}}(x+m_i^2)=\frac{\alpha _{n_c}^2}{\mathrm{\Lambda }^2}x\underset{a=1}{\overset{n_c1}{}}(x\alpha _a)^2.$$ (9.120) By moving terms, we find $`{\displaystyle \underset{i=1}{\overset{n_f}{}}}(x+m_i^2)`$ $`=`$ $`\left[{\displaystyle \underset{k=0}{\overset{n_c1}{}}}(1)^ks_k(\varphi ^2)x^{n_ck}+m_1\mathrm{}m_{n_f}\right]^2+{\displaystyle \frac{\alpha _{n_c}^2}{\mathrm{\Lambda }^2}}x\left[{\displaystyle \underset{k=0}{\overset{n_c1}{}}}(1)^ks_k(\alpha )x^{n_c1k}\right]^2.`$ Now consider the following polynomial $$F(x)=\underset{i=1}{\overset{n_f}{}}(\sqrt{x}+im_i)=\underset{k=0}{\overset{n_f}{}}i^ks_k(m)\sqrt{x}^{n_fk}.$$ (9.122) This polynomial can be divided into the “real” and “imaginary” parts (this is not strictly true because $`m`$’s are complex, but what is meant here is the division between terms of odd powers in $`i`$ and of even powers in $`i`$), $`F(x)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{n_f/2}{}}}(1)^ks_{2k}(m)\sqrt{x}^{n_f2k}+{\displaystyle \underset{k=0}{\overset{n_f/21}{}}}i(1)^ks_{2k+1}(m)\sqrt{x}^{n_f2k1}`$ (9.123) $`=`$ $`{\displaystyle \underset{k=0}{\overset{n_c}{}}}(1)^ks_{2k}(m)x^{n_ck}+i\sqrt{x}{\displaystyle \underset{k=0}{\overset{n_c1}{}}}(1)^ks_{2k+1}(m)x^{n_ck1},`$ where we used $`n_f=2n_c`$. Similarly consider the polynomial $$G(x)=\underset{i=1}{\overset{n_f}{}}(\sqrt{x}im_i)=\underset{k=0}{\overset{n_c}{}}(1)^ks_{2k}(m)x^{n_ck}i\sqrt{x}\underset{k=0}{\overset{n_c1}{}}(1)^ks_{2k+1}(m)x^{n_ck1}.$$ (9.124) From the definition, $$F(x)G(x)=\underset{i=1}{\overset{n_f}{}}(x+m_i^2).$$ (9.125) On the other hand, this product is also given by $$F(x)G(x)=\left[\underset{k=0}{\overset{n_c}{}}(1)^ks_{2k}(m)x^{n_ck}\right]^2+\left[\sqrt{x}\underset{k=0}{\overset{n_c1}{}}(1)^ks_{2k+1}(m)x^{n_ck1}\right]^2.$$ (9.126) This is precisely the same as Eq. (LABEL:eq:tosolve4) upon identifications $$\frac{1}{\mathrm{\Lambda }}\alpha _{n_c}s_k(\alpha )=s_{2k+1}(m),s_k(\varphi ^2)=(1)^{n_c}s_{2k}(m).$$ (9.127) The first equation gives $`\alpha _{n_c}=\mathrm{\Lambda }s_1(m)`$ by setting $`k=0`$ and using $`s_0=1`$. In the last identification, we used the fact that $`s_{2n_c}(m)=s_{n_f}(m)=m_1\mathrm{}m_{n_f}`$. Note that $`s_k(\varphi ^2)`$ excludes $`\varphi _{n_c}^2`$ and hence are different from the conventional gauge-invariant polynomials. This gives explicit solutions to the vacuum. Once we have this solution, however, we can obtain other $`2^{n_f1}1`$ solutions as follows. First note that the curve Eq. (6.24) is invariant under changing signs of even number of masses. Therefore, we can change signs of even number of masses from the solution (9.127). This gives $`2^{n_f1}`$ solutions in total agreeing with $`𝒩_1`$ in the large $`\mu `$ analysis. This solution therefore decomposes under $`U(n_f)`$ for the equal mass case as $`2^{n_f1}={}_{n_f}{}^{}C_{0}^{}+{}_{n_f}{}^{}C_{2}^{}+\mathrm{}{}_{n_f}{}^{}C_{n_f}^{}`$, i.e., to even-rank anti-symmetric tensors, reminiscent of the spinor representation. Now we come back to the other solution $`\varphi _{n_c}^2=2\mathrm{\Lambda }^2`$. This choice changes Eq. (LABEL:eq:tosolve4) to $`{\displaystyle \underset{i=1}{\overset{n_f}{}}}(x+m_i^2)`$ $`=`$ $`\left[+{\displaystyle \underset{k=0}{\overset{n_c1}{}}}(1)^ks_k(\varphi ^2)x^{n_ck}+m_1\mathrm{}m_{n_f}\right]^2{\displaystyle \frac{\alpha _{n_c}^2}{\mathrm{\Lambda }^2}}x\left[{\displaystyle \underset{k=0}{\overset{n_c1}{}}}(1)^ks_k(\alpha )x^{n_c1k}\right]^2.`$ This is again the curve for the superconformal $`USp(2n_f2)`$ theory with $`n_f`$ flavors with a different choice of $`g=+1`$. The sign changes can be absorved if we flip the sign of $`m_1`$ and change $`\alpha _{n_c}`$ to $`i\alpha _{n_c}`$. The sign flip of one of the quark masses implies that we have an odd number of minus signs in quark masses. This solution therefore decomposes under $`U(n_f)`$ for the equal mass case as $`2^{n_f1}={}_{n_f}{}^{}C_{1}^{}+{}_{n_f}{}^{}C_{3}^{}+\mathrm{}{}_{n_f}{}^{}C_{n_f1}^{}`$,, i.e., to odd-rank anti-symmetric tensors, reminiscent of the anti-spinor representation. Comments on the case of equal mass are in order. As noted above, we can flip the signs of even number of quark masses, and therefore a general situation has $`2r`$ negative masses $`m`$ and $`n_f2r`$ positive masses $`m`$. Let us study the location of branch points in this situation. The curve studied above has the form $$y^2=4\mathrm{\Lambda }^2\alpha _{n_c}^2\underset{a=1}{\overset{n_c1}{}}(x\alpha _a)^2$$ (9.129) near $`x0`$. Given the solutions Eq. (9.127), we can write $`2\mathrm{\Lambda }\alpha _{n_c}{\displaystyle \underset{a=1}{\overset{n_c1}{}}}(x\alpha _a)=2\mathrm{\Lambda }\alpha _{n_c}{\displaystyle \underset{k=0}{\overset{n_c1}{}}}(1)^ks_k(\alpha )x^{n_c1k}=2\mathrm{\Lambda }^2{\displaystyle \underset{k=0}{\overset{n_c1}{}}}(1)^ks_{2k+1}(m)x^{n_c1k}`$ $`=2\mathrm{\Lambda }^2{\displaystyle \frac{1}{2i}}\sqrt{x}\left[{\displaystyle \underset{i=1}{\overset{n_f}{}}}(\sqrt{x}+im_i){\displaystyle \underset{i=1}{\overset{n_f}{}}}(\sqrt{x}im_i)\right]`$ $`=i\mathrm{\Lambda }^2\sqrt{x}\left[(\sqrt{x}+im_i)^{n_f2r}(\sqrt{x}im_i)^{2r}(\sqrt{x}im_i)^{n_f2r}(\sqrt{x}+im_i)^{2r}\right].`$ If $`4rn_f`$, we can factor $`(x+m^2)^{2r}`$ from above and hence $`(x+m^2)^{4r}`$ from the curve. We interpret this factor as the emergence of $`SU(2r)`$ gauge theory. If $`4r>n_f`$, we can factor $`(x+m^2)^{n_f2r}`$ from above and hence $`(x+m^2)^{2n_f4r}`$ from the curve. We interpret this factor as the emergence of $`SU(n_f2r)=SU(2(n_cr))`$ gauge theory. This fact can also be understood from the Higgs branch picture. When quark masses are large and equal, one can cancel quark masses by the adjoint VEV classically (squark singularity) and obtain $`U(k)`$ gauge theories $`k=0,1,\mathrm{},n_f/2=n_c`$ depending on how many components of the adjoint field is used to cancel the quark masses. These are all infrared-free or scale-invariant theories and the gauge fields survive the quantum effects. While smoothly decreasing the quark masses, the Higgs branch emanating from the squark singularity does not change due to the non-renormalization theorem and hence the $`U(k)`$ theories survive to the small mass studied using the Coulomb branch. The effective low-energy Lagrangian which describes physics around the squark singularity is therefore the one given by Argyres–Plesser–Seiberg for $`SU(n_c)`$ theories at the non-baryonic roots . These theories indeed produce $`{}_{n_f}{}^{}C_{k}^{}`$ vacua upon $`\mu 0`$ and generic quark mass perturbation. Therefore the whole picture nicely fit together. On the other hand, the strictly massless case has the high singlarity $`x^{2n_c}`$ and appears to give a new superconformal theory with a global symmetry $`SO(2n_f)`$. We do not know of a weakly coupled description of this singularity in terms of a weakly coupled local field theory. We have checked that this singularity indeed produces mutually non-local dyons for $`n_c=2`$ and $`n_f=4.`$ Now we come back to the case of smaller $`n_f`$. Around the Chebyshev solution in the presence of quark masses, we study the curve $`xy^2`$ $`=`$ $`\left[x{\displaystyle \underset{a=1}{\overset{n_f/21}{}}}(x\varphi _a^2){\displaystyle \underset{k=1}{\overset{n_c+1n_f/2}{}}}(x\varphi _k^2\kappa _k^2)+2\mathrm{\Lambda }^{2n_c+2n_f}m_1\mathrm{}m_{n_f}\right]^2`$ (9.131) $`4\mathrm{\Lambda }^{2(2n_c+2n_f)}{\displaystyle \underset{i=1}{\overset{n_f}{}}}(x+m_i^2),`$ where $`\varphi _k^2`$ are given by those for the Chebyshev polynomial. However, for the purpose of studying the behavior around $`x=0`$, both $`x`$ and the shifts $`\kappa _k^2`$ can be ignored relative to $`\varphi _k^2`$. We need to know $`_{k=1}^{n_c+1n_f/2}(\varphi _k^2)`$. This can be calculated as $`{\displaystyle \underset{k=1}{\overset{n_c+1n_f/2}{}}}(\varphi _k^2)`$ $`=`$ $`2\mathrm{\Lambda }^{2n_c+2n_f}T_{2n_c+2n_f}\left({\displaystyle \frac{\sqrt{x}}{2\mathrm{\Lambda }}}\right)\|_{x0}`$ (9.132) $`=`$ $`2\mathrm{\Lambda }^{2n_c+2n_f}\mathrm{cos}\left[(2n_c+2n_f)\mathrm{arccos}{\displaystyle \frac{\sqrt{x}}{2\mathrm{\Lambda }}}\right]\|_{x0}`$ $`=`$ $`2(1)^{n_c+1n_f/2}\mathrm{\Lambda }^{2n_c+2n_f}.`$ Then the curve is approximated as $`{\displaystyle \frac{xy^2}{4\mathrm{\Lambda }^{2(2n_c+2n_f)}}}`$ $`=`$ $`\left[(1)^{n_c+1n_f/2}x{\displaystyle \underset{a=1}{\overset{n_f/21}{}}}(x\varphi _a^2)+m_1\mathrm{}m_{n_f}\right]^2`$ (9.133) $`{\displaystyle \underset{i=1}{\overset{n_f}{}}}(x+m_i^2),`$ This is nothing but the left hand side of Eq.(9.119) for the effective curve of $`USp(n_f)`$ theory upon changing normalizations of $`x`$, $`y`$, $`\varphi _a^2`$. The rest of the analysis therefore is exactly the same as in the $`2n_c=n_f`$ case. With other Chebyshev solutions obtained by $`Z_{2n_c+2n_f}`$, the relative sign between the two terms in the square bracket changes. This sign change corresponds to two different solutions $`\varphi _{n_c}^2=\pm 2\mathrm{\Lambda }^2`$ in the $`2n_c=n_f`$ case and hence they give decompositions into even-rank (odd-rank) anti-symmetric tensors under $`U(n_f)`$ for the equal mass perturbation, respectively. iii) Special (baryonic-like) point $`USp(2n_c)`$ theories also have special Higgs branch root (Eq.(6.30)-Eq.(6.32)), similar to the baryonic roots of the $`SU(n_c)`$ theories. This point is obtained by setting $`\varphi _1,\mathrm{},\varphi _{n_c\stackrel{~}{n}_c}0`$, and $`\varphi _{n_c\stackrel{~}{n}_c+1}=\mathrm{}=\varphi _{n_c}=0`$ in the original $`USp(2n_c)`$ curve $$xy^2=\left[x\underset{a=1}{\overset{n_c}{}}(x\varphi _a^2)+\mathrm{\Lambda }^{2n_c+2n_f}\underset{i=1}{\overset{n_f}{}}m_i\right]^2\mathrm{\Lambda }^{4n_c+42n_f}\underset{i=1}{\overset{n_f}{}}(x+m_i^2),$$ (9.134) (here and below, $`\stackrel{~}{n}_c=n_fn_c2`$), leading to $$y^2=x^{2\stackrel{~}{n}_c+1}\underset{k=1}{\overset{n_c\stackrel{~}{n}_c}{}}(x\mathrm{\Phi }_k)^24\mathrm{\Lambda }^{4n_c+42n_f}x^{n_f1}.$$ (9.135) Nonvanishing $`\mathrm{\Phi }`$’s are taken as $$(\mathrm{\Phi }_1^2,\mathrm{},\mathrm{\Phi }_k^2,\mathrm{},\mathrm{\Phi }_{n_c\stackrel{~}{n}_c}^2)=\mathrm{\Lambda }^2(\omega ,\mathrm{},\omega ^{2k1},\mathrm{},\omega ^{2(n_c\stackrel{~}{n}_c)1}),$$ (9.136) where $`\omega =e^{\pi i/(n_c\stackrel{~}{n}_c)}`$. Note that our $`\omega `$ is the square root of $`\omega `$ in Argyres–Plesser–Seiberg paper because of later convenience. Then the product $`_{k=1}^{n_c\stackrel{~}{n}_c}(x\mathrm{\Phi }_k)`$ can be rewritten as $`x^{n_c\stackrel{~}{n}_c}+\mathrm{\Lambda }^{2(n_c\stackrel{~}{n}_c)}`$, and the curve becomes $`y^2`$ $`=`$ $`x^{2\stackrel{~}{n}_c+1}\left[(x^{n_c\stackrel{~}{n}_c}+\mathrm{\Lambda }^{2(n_c\stackrel{~}{n}_c)})^24\mathrm{\Lambda }^{4n_c+42n_f}x^{n_c\stackrel{~}{n}_c}\right]`$ (9.137) $`=`$ $`x^{2\stackrel{~}{n}_c+1}(x^{n_c\stackrel{~}{n}_c}\mathrm{\Lambda }^{2(n_c\stackrel{~}{n}_c)})^2.`$ The double zeros of the factor in the parenthesis are at $$x=\mathrm{\Lambda }^2\omega ^2,\mathrm{\Lambda }^2\omega ^4,\mathrm{},\mathrm{\Lambda }^2\omega ^{2k},\mathrm{},\mathrm{\Lambda }^2\omega ^{2(n_c\stackrel{~}{n}_c)}=\mathrm{\Lambda }^2.$$ (9.138) One crucial difference from the $`SU(n_c)`$ case is that there is no choice but keep all of the “large” double zeros because the zeros at $`x=0`$ has an odd power $`2\stackrel{~}{n}_c+1`$ and hence leaves one of the zeros not doubled anyway. This explains why the separation between $`𝒩_1`$ and $`𝒩_2`$ works out nicely with the $`USp(2n_c)`$ theories, but one of the $`r`$-roots gets mixed up with the baryonic root for $`SU(n_c)`$ theories. Another crucial difference is that there is no constraint among $`\varphi _a`$’s. Therefore we can just fix “large” ones and study the low-energy curve. Using $$\underset{k=1}{\overset{n_c\stackrel{~}{n_c}}{}}(x\mathrm{\Lambda }^2\omega ^{2k1})=x^{n_c\stackrel{~}{n}_c}+\mathrm{\Lambda }^{2(n_c\stackrel{~}{n}_c)}\mathrm{\Lambda }^{2(n_c\stackrel{~}{n}_c)},$$ (9.139) the curve is approximated as $`xy^2`$ $`=`$ $`\left[x{\displaystyle \underset{a=1}{\overset{\stackrel{~}{n}_c}{}}}(x\varphi _a^2)\mathrm{\Lambda }^{2(n_c\stackrel{~}{n}_c)}+\mathrm{\Lambda }^{2n_c+2n_f}{\displaystyle \underset{i=1}{\overset{n_f}{}}}m_i\right]^2\mathrm{\Lambda }^{4n_c+42n_f}{\displaystyle \underset{i=1}{\overset{n_f}{}}}(x+m_i^2)`$ $`=`$ $`\mathrm{\Lambda }^{4(n_c\stackrel{~}{n}_c)}\left\{\left[x{\displaystyle \underset{a=1}{\overset{\stackrel{~}{n}_c}{}}}(x\varphi _a^2)+{\displaystyle \frac{1}{\mathrm{\Lambda }^{n_c\stackrel{~}{n}_c}}}{\displaystyle \underset{i=1}{\overset{n_f}{}}}m_i\right]^2{\displaystyle \frac{1}{\mathrm{\Lambda }^{2(n_c\stackrel{~}{n}_c)}}}{\displaystyle \underset{i=1}{\overset{n_f}{}}}(x+m_i^2)\right\}.`$ This curve clearly describes an infrared free $`USp(2\stackrel{~}{n}_c)`$ theory. Now we choose $`r`$ out of $`\stackrel{~}{n}_c`$ $`\varphi _a`$’s to match $`r`$ of the mass-squared ($`{}_{n_f}{}^{}C_{r}^{}`$ choices), e.g., $$\varphi _1^2=m_1^2,\mathrm{},\varphi _r^2=m_r^2,$$ (9.141) while the other $`\varphi _a^2`$ are still allowed to fluctuate but with magnitudes much less than $`m^2`$. Note that the absence of constraints among $`\varphi _a`$ allows us to have $`r\stackrel{~}{n}_c`$ while we had $`r<\stackrel{~}{n}_c`$ in $`SU(n_c)`$ theories. Then the low-energy curve (LABEL:eq:USpLEcurve) can further be approximated as $$xy^2=\mathrm{\Lambda }^{4(n_c\stackrel{~}{n}_c)}\left[x\underset{a=1}{\overset{\stackrel{~}{n}_cr}{}}(x\varphi _a^2)\underset{i=1}{\overset{r}{}}m_i^2+\frac{1}{\mathrm{\Lambda }^{n_c\stackrel{~}{n}_c}}\underset{i=1}{\overset{n_f}{}}m_i\right]^2\frac{1}{\mathrm{\Lambda }^{2(n_c\stackrel{~}{n}_c)}}\underset{i=1}{\overset{n_f}{}}m_i^2,$$ (9.142) which is nothing but the same as the curve of pure $`USp(2(\stackrel{~}{n}_cr+1))`$ Yang–Mills theories with the dynamical scale $$\mathrm{\Lambda }_{\mathrm{𝑝𝑢𝑟𝑒}}^{2(\stackrel{~}{n}_cr+1)}=\frac{1}{\mathrm{\Lambda }^{n_c\stackrel{~}{n}_c}}\frac{_{i=1}^{n_f}m_i}{_{i=1}^rm_i^2}m^{2(\stackrel{~}{n}_cr+1)}\left(\frac{m}{\mathrm{\Lambda }}\right)^{n_c\stackrel{~}{n}_c}m^{2(\stackrel{~}{n}_cr+1)}.$$ (9.143) This justifies the assumption of $`\varphi _a^2m^2`$. Since this is the curve of pure $`USp(2(\stackrel{~}{n}_cr+1))`$ Yang–Mills theories, it gives $`(\stackrel{~}{n}_cr+1)`$ vacua, and hence in total we find $$𝒩_2=\underset{r=0}{\overset{\stackrel{~}{n}_c}{}}(\stackrel{~}{n}_cr+1){}_{n_f}{}^{}C_{r}^{}.$$ (9.144) iv) Summary of the vacuum counting in $`USp(2n_c)`$ theories There are thus two groups of $`N=1`$ vacua predicted by the Seiberg-Witten curve in $`USp(2n_c)`$ theories. The Chebyshev point Eq.(6.27), Eq.(6.28), gives rise to $`𝒩_1=(2n_c+2n_f)2^{n_f1}`$ vacua upon mass perturbation , while the special point Eq.(6.30), Eq.(6.31), leads to $`𝒩_2=_{r=0}^{\stackrel{~}{n}_c}(\stackrel{~}{n}_cr+1){}_{n_f}{}^{}C_{r}^{}`$ vacua. Their sum coincides with the total number of $`N=1`$ vacua found from the semiclassical as well as from large $`\mu `$ analyses. Acknowledgment One of the authors (K.K.) thanks Lawrence Berkeley National Laboratory, University of California, Berkeley, and ITP, University of California Santa Barbara, for their warm hospitality. Part of the work was done during the workshop, “Supersymmetric Gauge Dynamics and String Theory” at ITP, UCSB, to which two of us (G.C. and K.K.) participated. The authors acknowledge useful discussions with Philip Argyres, Alex Buchel, S. Elitzur, Amit Giveon, Prem Kumar, Rita Pardini, Misha Shifman and Arkady Vainshtein. This research was supported in part by the National Science Foundation under Grant No. PHY-94-07194, PHY-95-14797, and in part by the Director, Office of Science, Office of High Energy and Nuclear Physics, Division of High Energy Physics of the U.S. Department of Energy under Contract DE-AC03-76SF00098. ## Appendix A $`SO(2N)USp(2N)=U(N)`$ The $`SO(2N)`$ generators are the most general pure imaginary anti-symmetric matrices. Break it down to the $`N\times N`$ blocks, and write them down as: $$\left(\begin{array}{cc}E& F\\ ^tF& D\end{array}\right),$$ (A.1) where $`D`$, $`E`$, $`F`$ are all pure imaginary $`N\times N`$ matrices, with the constraints $`{}_{}{}^{t}E=E`$, $`{}_{}{}^{t}D=D`$. The generators of $`USp(2N)`$ are given by $$\left(\begin{array}{cc}B& A\\ C& ^tB\end{array}\right),$$ (A.2) with the constraints, $`{}_{}{}^{t}A=A`$, $`{}_{}{}^{t}C=C`$, $`A^{}=C`$, $`B^{}=B`$. The way to compare them is to go to the bases of $`SO(2N)`$ where it naturally breaks to $`N+\overline{N}`$ under $`U(N)`$. This can be done by the following rotation, $`\left(\begin{array}{cc}1/\sqrt{2}& i/\sqrt{2}\\ i/\sqrt{2}& 1/\sqrt{2}\end{array}\right)\left(\begin{array}{cc}E& F\\ ^tF& D\end{array}\right)\left(\begin{array}{cc}1/\sqrt{2}& i/\sqrt{2}\\ i/\sqrt{2}& 1/\sqrt{2}\end{array}\right)`$ (A.12) $`={\displaystyle \frac{1}{2}}\left(\begin{array}{cc}(E+D)+i(F+^tF)& i(ED)+(F^tF)\\ i(ED)+(F^tF)& (E+D)i(F+^tF)\end{array}\right).`$ Since both $`E`$, $`D`$ are anti-symmetric, $`(E+D)`$ in the 1st block is the most general anti-symmetric imaginary matrix, while $`i(F+^tF)`$ is the most general symmetric real matrix. Their sum gives the most general hermitian matrix. Comparing to the $`USp(N)`$ generators, the off-diagonal blocks are completely symmetric for $`USp(N)`$ and completely anti-symmetric for $`SO(2N)`$, and hence there is no overlap. While the diagonal blocks are the most general $`N\times N`$ hermitian matrices, and overlap completely, hence $`SO(2N)USp(2N)=U(N)`$. ## Appendix B Semiclassical Monopole States The Jackiw–Rebbi zero mode has the form $$\psi _L^{(0)}=ib\eta (x);\psi _R^{(0)}=b\eta (x).$$ (B.1) in the chiral representation, where the commutation relations are the standard one: $$\{b^i,(b^j)^{}\}=\delta ^{ij}.$$ (B.2) Given the monopole state $`\|\mathrm{\Omega }`$, one can construct $`2^{n_f}`$ positive-norm states by acting various number of creation operators upon it, $$(b^{i_1})^{}(b^{i_2})^{}\mathrm{}(b^{i_k})^{}\|\mathrm{\Omega },$$ (B.3) which are all spinless bosons . In general $`SU(n_c)`$ theories ($`n_c3`$) with $`n_f`$ flavors the Noether currents are: $$J_\mu ^A=\overline{\psi }_D^i\gamma _\mu \lambda _{ij}^A\psi _D^j$$ (B.4) so that the charge operators are $$Q^A=(b^i)^{}\lambda _{ij}^Ab_j+\text{non zero modes}.$$ (B.5) The semiclassical monopole multiplets are formed by the state $`\|\mathrm{\Omega }`$ which is a singlet, the $`n_f`$ states $`(b^i)^{}\|\mathrm{\Omega }`$ belonging to a $`\underset{¯}{n}_f`$, the $`{}_{n_f}{}^{}C_{2}^{}`$ states $`(b^i)^{}(b^j)^{}\|\mathrm{\Omega }`$ which is the second rank antisymmetric tensor, etc. Although semi-classically $`2^{n_f}`$ states (B.3) are all degenerates, higher quantum effects lift such a degeneray in general, and only those belonging to the same irreducible representation will have the same mass (monopole multiplets). In the $`USp(2n_c)`$ gauge theories the fermionic part of the lagrangian is $$\underset{i=1}{\overset{2n_f}{}}[\overline{\psi }_a^ii\overline{\sigma }^\mu (𝒟_\mu )_{ab}\psi _b^i+\frac{1}{\sqrt{2}}\psi _a^i\varphi _{ab}\psi _c^iJ^{bc}+\mathrm{h}.\mathrm{c}.],$$ (B.6) where all fermions are pure left–handed. In this basis of fermions the $`SO(2n_f)`$ symmetry is manifest and the global symmetry current is simply $$J_\mu ^{ij}=\overline{\psi }_a^i\overline{\sigma }^\mu \psi _a^j(ij),$$ (B.7) their charges ($`SO(2n_f)`$ generators) are $$Q_\mu ^{ij}=\overline{\psi }_a^i\psi _a^j(ij).$$ (B.8) The zero mode operators in $`USp(2n_c)`$ theories have the form ($`\gamma ^i`$ are just the name of these operators, not gamma matrices) $$\psi ^i=\gamma ^i\eta (x)+\mathrm{},(i=1,2,\mathrm{}2n_f),$$ (B.9) where $$\gamma ^i=b^i+b_{}^{i}{}_{}{}^{};\gamma ^{n_f+i}=\frac{1}{i}(b^ib_{}^{i}{}_{}{}^{}),(i=1,2,\mathrm{}n_f).$$ (B.10) This particular form of the zero mode contribution reflects the fact that the fermion basis in which the standard Jackiw–Rebbi solution Eq. (B.1) applies and the one which transforms as an $`SO(2n_f)`$ vector, are related by: $$\psi _a^i=\frac{1}{\sqrt{2}}(\widehat{\psi }_a^{2i1}+\widehat{\psi }_a^{2i}),\psi _a^{n_f+i}=\frac{1}{i\sqrt{2}}(\widehat{\psi }_a^{2i1}\widehat{\psi }_a^{2i}),(i=1,2,\mathrm{}n_f).$$ (B.11) where $$\widehat{\psi }_a^{2i1}\psi _{La}^i\widehat{\psi }_a^{2i}\overline{\psi }_{Ra}^i,(i=1,\mathrm{},n_f).$$ (B.12) Note that $`\gamma ^i`$’s are all real. These relations show that $`\gamma ^i`$’s obey the Clifford algebra, $$\{\gamma ^i,\gamma ^j\}=2\delta _{ij}.$$ (B.13) By substituting Eq. (B.9) into Eq. (B.8) we find that the generators of $`SO(2n_f)`$ symmetry are (by renormalizing by a constant): $$\mathrm{\Sigma }_{ij}=\frac{1}{4i}[\gamma ^i,\gamma ^j],$$ (B.14) which obviously satisfies the standard $`SO(2n_f)`$ algebra. This shows that various monopole states: $$(b^{i_1})^{}(b^{i_2})^{}\mathrm{}(b^{i_k})^{}\|\mathrm{\Omega }.$$ (B.15) transform as spinor representations of $`SO(2n_f)`$. Furthermore one notes that states with odd or even numbers of creation operators have definite “chirality” with respect to $$\gamma ^{2N_f+1}=(i)^{N_f}\gamma ^1\gamma ^2\mathrm{}\gamma ^{2N_f}=\underset{i=1}{\overset{N_f}{}}(12b^ib^i)$$ (B.16) so that each of them transform independently. Each monopole state thus belongs to a spinor representation of definite chirality of the global $`SO(2n_f)`$ group. ## Appendix C Explicit Formulae for $`a_{Di}`$, $`a_i`$, $`\frac{a_{Di}}{u_j},`$ and $`\frac{a_i}{u_j}`$ $$y^2=\underset{k=1}{\overset{n_c}{}}(x\varphi _k)^2+4\mathrm{\Lambda }^{2n_cn_f}\underset{j=1}{\overset{n_f}{}}(x+m_j),SU(n_c),n_f2n_c2,$$ (C.1) and $$y^2=\underset{k=1}{\overset{n_c}{}}(x\varphi _k)^2+4\mathrm{\Lambda }\underset{j=1}{\overset{n_f}{}}\left(x+m_j+\frac{\mathrm{\Lambda }}{n_c}\right),SU(n_c),n_f=2n_c1,$$ (C.2) with $`\varphi _k`$ subject to the constraint $`_{k=1}^{n_c}\varphi _k=0`$, and $$xy^2=\left[x\underset{a=1}{\overset{n_c}{}}(x\varphi _a^2)^2+2\mathrm{\Lambda }^{2n_c+2n_f}m_1\mathrm{}m_{n_f}\right]^24\mathrm{\Lambda }^{2(2n_c+2n_f)}\underset{i=1}{\overset{n_f}{}}(x+m_i^2),USp(2n_c).$$ (C.3) In each case, these represent a genus $`g=n_c1`$ ($`g=n_c`$ for the $`USp(2n_c)`$ case) hypertorus, which are characterized by $`2g`$ homology cycles $`\alpha _i`$, $`\beta _i`$, $`i=1,2,\mathrm{},g`$. These cycles are taken in the doubly sheeted $`x`$ plane to surround two branch points of $`y`$, and such that they intersect pairwise, in the canonical way, $`(\alpha _i\beta _j)=\delta _{ij}.`$ $`a_{Di}/u_j,`$ and $`a_i/u_j`$ are given by the $`g\times 2g`$ period integrals of the holomorphic differentials (neglecting the normalization constant) - $$\frac{a_{Di}}{u_j}=_{\alpha _i}\frac{dxx^{j1}}{y};\frac{a_i}{u_j}=_{\beta _i}\frac{dxx^{j1}}{y};$$ (C.4) whereas $`a_{Di}`$, $`a_i`$ are given by the integrals of the meromorphic differential $`\lambda `$ (defined such as $`d\lambda /du_j=\frac{dxx^{j1}}{y};`$ $$a_{Di}=_{\alpha _i}\lambda ,a_i=_{\beta _i}\lambda ,$$ (C.5) where some additive terms proportional to the bare quark masses are neglected. ## Appendix D Absence of the Non-Baryonic Root with $`r=\stackrel{~}{n}_c=n_fn_c`$ In this Appendix we prove the absence of the non-baryonic root with $`r=\stackrel{~}{n}_c=n_fn_c`$ for $`SU(n_c)`$ gauge theory. The nonbaryonic branch root in question is characterized by the adjoint VEVS $$\text{diag}\varphi =(\underset{n_fn_c}{\underset{}{0,0,\mathrm{},0}},\varphi _1,\mathrm{}\varphi _{2n_cn_f}),\underset{a=1}{\overset{2n_cn_f}{}}\varphi _a=0:$$ (D.1) the curve has the form $`y^2`$ $`=`$ $`x^{2\stackrel{~}{n}_c}{\displaystyle \underset{a=1}{\overset{2n_cn_f}{}}}(x\varphi _a)^24\mathrm{\Lambda }^{2n_cn_f}x^{n_f}`$ (D.2) $`=`$ $`x^{2\stackrel{~}{n}_c}\left[{\displaystyle \underset{a=1}{\overset{N}{}}}(x\varphi _a)^24\mathrm{\Lambda }^Nx^N\right],N=2n_cn_f.`$ We prove below that this curve cannot be put (for whatever $`\{\varphi _a\}`$) in the form, $$y^2=x^{2\stackrel{~}{n}_c}\underset{i=1}{\overset{N1}{}}(x\alpha _i)^2(x\gamma )(x\delta ),\gamma \delta .$$ (D.3) 1. Theorem: The function $$F(x)=\underset{a=1}{\overset{N}{}}(x\varphi _a)^24x^N,\underset{a=1}{\overset{N}{}}\varphi _a0,$$ (D.4) $`x`$, $`\{\varphi \}`$ complex, cannot have exactly $`N1`$ double factors. 2. $`N=2,3,4`$ For $`N=2,3,4`$, we have checked explicitly that there are indeed no $`\varphi `$ configurations such that $`F(x)`$ has exactly $`N1`$ pairs of double factors. There are either $`N`$ pairs, as can be realized by taking $$\varphi _a=(\omega _0)^a,\omega _0=e^{2\pi i/N},$$ (D.5) or less than $`N1`$ pairs of double factors. 3. $`N=2n_cn_f`$ even In this case, $$F(x)=F_+(x)F_{}(x),F_\pm (x)=\underset{a=1}{\overset{N}{}}(x\varphi _a)\pm 2x^{N/2}.$$ (D.6) First of all, there cannot be any common factor between $`F_+(x)`$ and $`F_{}(x)`$. For if there is one, $`(xx^{})`$, $`F_+(x^{})=F_{}(x^{})=0,`$ hence $`x^{}=0`$. It means that there is an extra power of $`x^2`$ in front (an extra $`\varphi _a=0`$), which is not possible because $`_{a=1}^N\varphi _a0`$. Since there are no common factor in $`F_+(x)`$ and $`F_{}(x)`$, in order to get at least $`N1`$ double factors, one of $`F_+(x)`$ and $`F_{}(x)`$ must be fully doubled up, say: $$F_+(x)=\underset{a=1}{\overset{N}{}}(x\varphi _a)+2x^{N/2}=\underset{a=1}{\overset{N/2}{}}(x\alpha _a)^2.$$ (D.7) We wish to prove that in this case $`F_{}(x)`$ is a perfect square also. In order to show it, note that $`F(x)`$ is invariant under the transformation, $$x\omega x;\varphi _a\omega \varphi _a,$$ (D.8) where $`\omega =\mathrm{exp}2\pi i/N`$. Note that under this transformation, $`F_+(x)`$ and $`F_{}(x)`$ get exchanged: $$F_+(x)F_{}(x);F_{}(x)F_+(x).$$ (D.9) Assume now that a configuration $`\{\varphi \}`$ such that (D.7 ) holds was found. $`\alpha _i`$’s are functions of $`\{\varphi \}`$ : $$\alpha _i=\alpha _i(\{\varphi \}).$$ (D.10) According to (D.9), $`F_{}(x)`$ can be found by the $`\omega `$ transformation from $`F_+(x)`$: $$F_{}(x)=\underset{i=1}{\overset{N/2}{}}[x\omega ^1\alpha _i(\{\omega \varphi \})]^2.$$ (D.11) Thus we have proved that if $`F(x)`$ has at least $`N1`$ double factors, then it has $`N`$ of them. 4. General $`N=2n_cn_f`$ Assume that $`\{\varphi _a\}`$ ’s are found such that $$F(x)=\underset{a=1}{\overset{N}{}}(x\varphi _a)^24x^N\underset{A=1}{\overset{N1}{}}(x\alpha _A)^2(x\gamma )(x\delta ),\gamma \delta ,$$ (D.12) where $`\alpha _A`$’s are all different among each other and none of them coincides either with $`\gamma `$ or with $`\delta `$. The left hand side of Eq.(D.12) is invariant under the transformation Eq.(D.8), so must be also the right hand side. That is $`{\displaystyle \underset{A=1}{\overset{N1}{}}}(x\omega ^1\alpha _A(\{\omega \varphi \}))^2(x\omega ^1\gamma (\{\omega \varphi \}))(x\omega ^1\delta (\{\omega \varphi \}))`$ (D.13) $`=`$ $`{\displaystyle \underset{A=1}{\overset{N1}{}}}(x\alpha _A)^2(x\gamma )(x\delta );`$ it implies however $$\underset{A=1}{\overset{N1}{}}(x\omega ^1\alpha _A(\{\omega \varphi \}))=\underset{A=1}{\overset{N1}{}}(x\alpha _A),$$ (D.14) $$(x\omega ^1\gamma (\{\omega \varphi \}))(x\omega ^1\delta (\{\omega \varphi \}))=(x\gamma )(x\delta ).$$ (D.15) Now Eq.(D.14) and Eq.(D.15), which are equivalent to $$\underset{A=1}{\overset{N1}{}}(\omega x\alpha _A(\{\omega \varphi \}))=\omega ^{N1}\underset{A=1}{\overset{N1}{}}(x\alpha _A(\{\varphi \})),$$ (D.16) $$(\omega x\gamma (\{\omega \varphi \}))(\omega x\delta (\{\omega \varphi \}))=\omega ^2(x\gamma (\{\varphi \}))(x\delta (\{\varphi \})),$$ (D.17) imply that the polynomials $$H_1(x,\varphi )=\underset{A=1}{\overset{N1}{}}(x\alpha _A),$$ (D.18) $$H_2(x,\varphi )=(x\gamma )(x\delta ),$$ (D.19) of order $`N1`$ and $`2`$, are both homogeneous in $`(\{\varphi \},x)`$: namely, $$H_1(x,\varphi )=x^{N1}+\underset{i=1}{\overset{N1}{}}s_i(\varphi )x^{N1i};$$ (D.20) $$H_2(x,\varphi )=x^2+\underset{i=1}{\overset{2}{}}t_i(\varphi )x^{2i};$$ (D.21) where $`s_i(\varphi )`$ and $`t_i(\varphi )`$ satisfy $$s_i(\omega \varphi )=\omega ^is_i(\varphi );t_i(\omega \varphi )=\omega ^it_i(\varphi ).$$ (D.22) This is so because $`H_1,H_2`$, being polynomials in $`x`$ of order less than $`N`$, have each term in them transformed non trivially under the $`\omega `$ transformation. It follows now that $$F(x)=H_1(x)^2H_2(x)$$ (D.23) is a homogeneous expression in $`(\{\varphi \},x)`$ with nontrivial coefficients in the expansion in $`x`$, which contradicts the form of $`F(x)`$, Eq.(D.4). We have thus shown that the assumption Eq.(D.12) is impossible. ## Appendix E Monodromies in $`SU(3)`$ Theories with $`n_f=4`$ In this Appendix we briefly describe the analysis of monodromy transformation around various singularities for $`SU(3)`$ gauge theory with $`n_f=4`$. To study the monodromies one sets $`(u,v)`$ slightly off the singularity interested, and lets $`(u,v)`$ make a small circle around it in the parameter space (QMS). From the way the branch points move around and the branch cuts get entangled one easily finds the manodromy matrix for $`a_{D1},a_{D2},a_1,a_2`$. One must also study how the positions of the branch points and cuts are varied as one goes from one singularity to another. This allows one to determine the homology cycles defining various periods, $`a_{D1},a_{D2},a_1,a_2`$, in a globally consistent manner. The quantum numbers of the massless states at each singularity are found from the nonvanishing eigenvectors of the monodromy matrix thus obtained. In many cases, the movements of the branch points can be studied analytically as well: we illustrate below how such a check can be made, in some examples. The $`SU(3)`$ gauge theory with four flavor has $`17`$ vacua for generic quark masses and for a nonvanishing adjoint mass. They collapse to six vacua in the limit of equal quark masses, three singlets, two quartets and one sextet. For $`\mathrm{\Lambda }=2`$ and $`m=2^6`$ they are at: \| 1. \| $`(u,v)=(0,0)`$ \| sextet \| \| --- \| --- \| --- \| \| 2. \| $`(u,v)=(0.92,0.14)`$ \| singlet \| \| 3. \| $`(u,v)=(0.85,0.09)`$ \| singlet \| \| 4. \| $`(u,v)=(1.05,0.02)`$ \| quartet \| \| 5. \| $`(u,v)=(0.95,0.01)`$ \| quartet \| \| 6. \| $`(u,v)=(1.00,0.06)`$ \| singlet \| The branch points and cuts (dotted lines) are chosen as shown in Fig. 9. Let us analyze each singularity, starting from the singularity 2. Singularity 2. The branch points near this singularity are located as in (1.1) in Fig. 10 with $`x_2x_3`$ and $`x_5x_6`$ on the singularity. To determine the massless BPS states condensing on the singularity, we perform a small circle around the singularity itself in the parameter space. The branch points transform as in 1.2 (Fig. 10) therefore the monodromy matrix is: $$M_2=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 1& 0& 1& 0\\ 0& 1& 0& 1\end{array}\right).$$ (E.1) The eigenvectors, with unimodular eigenvalues, of the (transpose of the) monodromy matrix give the charges of the massless particles condensing in the singularity. In this case, we have: $$(n_{m1},n_{m2},n_{e1},n_{e2})=(1,0,0,0),(0,1,0,0),$$ (E.2) i.e., the two monopoles of the two abelian factors. Singularity 3. The $`x`$–plane is shown in 1.3 of Fig. 10. On the singularity: $`x_1x_2`$ and $`x_4x_5`$. The monodromy matrix is (see 1.4 Fig. 10): $$M_3=\left(\begin{array}{cccc}0& 0& 1& 0\\ 0& 0& 0& 1\\ 1& 0& 2& 0\\ 0& 1& 0& 2\end{array}\right).$$ (E.3) The charges of the condensing particles are: $$(n_{m1},n_{m2},n_{e1},n_{e2})=(0,1,0,1),(1,0,1,0),$$ (E.4) i.e. two dyons of the two abelian factors. By an appropriate redefinition of $`(n_{m1},n_{m2},n_{e1},n_{e2})`$ they become $$(n_{m1}^{^{}},n_{m2}^{^{}},n_{e1},n_{e2})=(0,1,0,0),(1,0,0,0):$$ (E.5) two monopoles of $`U(1)^2`$. Singularity 6. At the singularity $`6`$ the coalescing branch points are: $`x_1x_2`$ and $`x_5x_6`$. See (1.5), (1.6) of Figs. 11. The monodromy matrix is: $$M_6=\left(\begin{array}{cccc}0& 0& 1& 0\\ 0& 1& 0& 0\\ 1& 0& 2& 0\\ 0& 1& 0& 1\end{array}\right).$$ (E.6) The charges of the condensing states: $$(n_{m1},n_{m2},n_{e1},n_{e2})=(1,0,1,0),(0,1,0,0).$$ (E.7) Note that, these objects are mutually local. By an appropriate redefinition of $`(n_{m1},n_{m2},n_{e1},n_{e2})`$ they become again $$(n_{m1}^{^{\prime \prime }},n_{m2},n_{e1},n_{e2})=(0,1,0,0),(1,0,0,0):$$ (E.8) two monopoles of $`U(1)^2`$. Singularity 1. This is the sextet singularity. The branch points ($`x`$-plane) are in the positions depicted in the 1.7 (Fig. 11), with $`x_2x_3x_4x_5`$ exactly on the singularity. Performing a small circle around the singularity (in the QMS) the branch points move as indicated in Fig. 1.8. We find the monodromy matrix $$M_1=\left(\begin{array}{cccc}1& 2& 0& 2\\ 2& 2& 2& 1\\ 1& 0& 1& 0\\ 2& 1& 2& 0\end{array}\right).$$ (E.9) so the massless particles have the quantum numbers: $$(n_{m1},n_{m2},n_{e1},n_{e2})=(0,1,0,1),(1,2,2,0),(1,1,0,1).$$ (E.10) The first and the second of these states are relatively nonlocal, hence this is a conformally invariant vacuum. This was to be expected since this singularity corresponds to a class 3 conformal theory of Eguch et. al. (see the main text). Singularity 4, 5. For singularity 4, the $`x`$–plane looks like in 1.9 (Fig. 12) and the branch points rotates as in Fig. 1.10. The coincident branch points, on the singularity are: $`x_3x_4`$ and $`x_5x_6`$. As for the singularity 5, the $`x`$–plane looks like in 1.11 and the branch points rotates as in 1.12. The coincident branch points, on the singularity are: $`x_1x_2`$ and $`x_3x_4`$. The monodromy matrix at the singularity $`4`$ turns out to be $$M_4=\left(\begin{array}{cccc}3& 4& 4& 0\\ 0& 1& 0& 0\\ 4& 4& 5& 0\\ 4& 5& 4& 1\end{array}\right).$$ (E.11) so the massless particles have the quantum numbers: $$(n_{m1},n_{m2},n_{e1},n_{e2})=(1,0,1,0),(0,1,0,0).$$ (E.12) They are relatively local. Note that $$M_4=T^4A,T=\left(\begin{array}{cccc}0& 1& 1& 0\\ 0& 1& 0& 0\\ 1& 1& 2& 0\\ 1& 2& 1& 1\end{array}\right),A=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 1& 0\\ 0& 3& 0& 1\end{array}\right),$$ (E.13) with the matrix $`T`$ having the same charge eigenvectors (E.12) and $`A`$ representing a possible change of homology cycles. This shows that the this singularty correspond to a quartet of singularity. The monodromy matrix at the singularity $`5`$ is $$M_5=\left(\begin{array}{cccc}4& 4& 5& 0\\ 0& 1& 0& 0\\ 5& 4& 6& 0\\ 4& 5& 4& 1\end{array}\right).$$ (E.14) so the massless particles have the same quantum numbers as at $`4`$: $$(n_{m1},n_{m2},n_{e1},n_{e2})=(1,0,1,0),(0,1,0,0).$$ (E.15) They are again relatively local and the same as at the point 4. Analytic determination of the monodromy One can actually study the movement of the branch points analytically in many cases. For instance take the singularities $`4`$ or $`5`$ and condsider the double branch point at $`m=\frac{1}{64}`$. (At both singularities, one of the double branch point occurs at $`x=m`$.) The auxiliary curve $$y^2=4\left(x+\frac{1}{64}\right)^4+(vux+x^3)^2,$$ (E.16) can be rewritten as: $$y^2=4\left(x+\frac{1}{64}\right)^4+\left\{\left(x+\frac{1}{64}\right)f(x)(uu_0)x(vv_0)\right\}^2$$ (E.17) with $$f(x)=x^2\frac{1}{64}x64v_0.$$ (E.18) The positions of the singularities $`4`$ and $`5`$ are given by $$u_0=\frac{17149}{16384};v_0=\frac{u_0}{64}=\frac{17149}{1048576}$$ (E.19) (the vacuum $`4`$), and by $$u_0=\frac{15613}{16384};v_0=\frac{u_0}{64}=\frac{15613}{1048576}$$ (E.20) (the vacuum $`5`$). Now, shift $`x`$ by $$x+\frac{1}{64}x^{}$$ (E.21) and rewrite the curve as: $$y^2=4x_{}^{}{}_{}{}^{4}+\left\{x^{}f\left(x^{}\frac{1}{64}\right)(uu_0)\left(x^{}\frac{1}{64}\right)(vv_0)\right\}^2.$$ (E.22) For small but nonzero values of $`uu_0`$ and/or $`vv_0`$ of order $`ϵ`$, the second term of the right hand side has the form $$\{(c_0x^{}+c_1x_{}^{}{}_{}{}^{2}+\mathrm{})+ϵx^{}+ϵ\}^2.$$ (E.23) and the curve looks like $$y^2x_{}^{}{}_{}{}^{4}+(x^{}+ϵ)^2.$$ (E.24) Shifting further $`x^{}`$ as $`\stackrel{~}{x}=x^{}+ϵ`$ one gets $$y^2\stackrel{~}{x}^2+(\stackrel{~}{x}ϵ)^4.$$ (E.25) The approximately doublet zeros of the right hand side are then found from $$\stackrel{~}{x}^2+ϵ^44ϵ^3\stackrel{~}{x}+6ϵ^2\stackrel{~}{x}^24ϵ\stackrel{~}{x}^3+\stackrel{~}{x}^4=0$$ to be of order $`\stackrel{~}{x}\pm ϵ^2`$. The splitting of the double branch point $`\frac{1}{64}`$ is therefore given by: $$x^{}ϵ\pm ϵ^2,.^..x\frac{1}{64}+ϵ\pm ϵ^2.$$ (E.26) A small circular motion in QMS around $`u_0`$ and/or $`v_0`$ yields a convoluted circular movements of the split double zeros: their relative position makes a $`4\pi `$ rotation (this is relevant to the monodromy analysis) while their center of mass performs a single $`2\pi `$ rotation. This movement of the branch points has been confirmed by the numerical analysis (Fig. 1.10.). Now, consider the other double zero at $`x_01`$ (in the case of the singularity No. $`4`$). The curve Eq.(E.17) becomes $$y^2=\left(x+\frac{1}{64}\right)^2\left[4\left(x+\frac{1}{64}\right)^2+f(x)^2\right]$$ (E.27) at the singularity. The double zero at $`x_01`$ come from (as can be seen checked explicitly): $$f(x)+2\left(x+\frac{1}{64}\right)=(xx_0)^2.$$ (E.28) Near the singularity the curve looks like $$y^2=4\left(x+\frac{1}{64}\right)^4+\left[\left(x+\frac{1}{64}\right)f(x)ϵ\right]^2$$ (E.29) ($`uu_0ϵ,`$ $`vv_0ϵ.`$) Near $`x=x_01`$, the double zero is split by the presence of terms linear in $`ϵ`$: $$(xx_0)^22ϵ\left(x+\frac{1}{64}\right)f(x)0$$ (E.30) that is: $$xx_0\pm 2ϵ^{1/2}.$$ (E.31) In conclusion, a small $`2\pi `$ circle around the singularity implies that the two branch points (double zero splitted) simply exchange between themselves. Again these movements of the branch points have been confirmed by a numerical analysis (1.10 in (Fig. 12)).
warning/0005/cond-mat0005175.html	ar5iv	text	# Electron correlations in narrow energy bands: ground state energy and metal-insulator transition ## 1 Introduction Among the observed in narrow-bands materials metal-insulator transitions (MIT) the significant interest is attracted to the transitions from paramagnetic metal state to paramagnetic insulator state at increase of temperature which exhibit the systems NiS<sub>2-x</sub>Se<sub>x</sub> -, (V<sub>1-x</sub>Cr<sub>x</sub>)<sub>2</sub>O<sub>3</sub> and Y<sub>1-x</sub>Ca<sub>x</sub>TiO<sub>3</sub> ; in these systems the paramagnetic insulator - paramagnetic metal transitions under external pressure are observed also. There are reasons to believe that noted transitions are the consequences of electron-electron interactions and can be described within the framework of Hubbard model . Hubbard model is the simplest model describing MIT in materials with narrow energy bands. This model describes a single non-degenerate band of electrons with the local Coulomb interaction. The model Hamiltonian contains two energy parameters: the hopping integral of an electron from one site to another and the intraatomic Coulomb repulsion of two electrons of the opposite spins. This model is used intensively (for recent reviews see Refs. -) in order to describe the peculiarities of physical properties of narrow-band materials; in this connection two-pole approaches are attractive. The two-pole approaches in the Hubbard model and the Hubbard bands conception (being the consequence of two-pole approximation) have been useful for understanding of the peculiarities of electric and magnetic properties of narrow-band materials . However within the framework of two-pole approaches there are series of issues, in particular the problem of metal-insulator transition description -. In the present paper recently proposed two-pole approximation is used to study effects of electron correlations in the Hubbard model. The single particle Green function and energy spectrum are obtained. In the ground state analitical dependences of energy gap, polar states (doublons or holes) concentration and energy of the system on model parameters are found. Dependences for energy gap, polar states concentration on temperature are calculated. The obtained results are compared with corresponding results of other approximations and are used for the interpretation of some experimental data. In particular, the observable transitions from an insulating state to a metallic one at increase of bandwidth, and from a metallic state to an insulating one at increasing temperature are explained. ## 2 Single-particle Green function and energy spectrum The Hubbard Hamiltonian in terms of transition-operators of $`i`$-site from state $`\|l`$ to state $`\|k`$ $`X_i^{kl}`$ is written as $`H=H_0+H_1+H_1^{};`$ (1) $`H_0=\mu {\displaystyle \underset{i\sigma }{}}\left(X_i^\sigma +X_i^2\right)+U{\displaystyle \underset{i}{}}X_i^2,`$ (2) $`H_1={\displaystyle \underset{ij\sigma ,ij}{}}t_{ij}\left(X_i^{\sigma 0}X_j^{0\sigma }+X_i^{2\overline{\sigma }}X_j^{\overline{\sigma }2}\right),`$ (3) $`H_1^{}={\displaystyle \underset{ij,ij}{}}t_{ij}(X_i^{\overline{\sigma }0}X_j^{\sigma 2}X_i^{\sigma 0}X_j^{\overline{\sigma }2}+h.c.),`$ (4) where $`\mu `$ is the chemical potential, $`U`$ is the intra-atomic Coulomb repulsion, $`t_{ij}`$ is the nearest-neighbor hopping integral, $`X_i^k`$ is the operator of number of $`\|k`$-states on $`i`$-site; $`\sigma `$ denotes spin of an electron ($`\sigma =,`$) and $`\overline{\sigma }`$ denotes the projection of electron spin opposite to $`\sigma `$; $`H_0`$ describes system in the atomic limit, $`H_1`$ describes electron hoppings between single occupied sites and empty sites (holes) (the first sum in $`H_1`$ – processes forming “h-band”) and electron hoppings between doubly occupied sites (doublons) and single occupied sites (the second sum in $`H_1`$ – processes forming “d-band”). $`H_1^{}`$ describes “hybridization” between the “h-band” and “d-band” (the processes of pair creation and annihilation of holes and doublons). The single-particle Green function is written in $`X_i^{kl}`$-operators as $`a_p\|a_s^+=X_p^2\|X_s^2X_p^0\|X_s^2X_p^2\|X_s^0+X_p^0\|X_s^0.`$ (5) The functions $`X_p^2\|X_s^2`$ and $`X_p^0\|X_s^2`$ satisfy the equations $`(E+\mu U)X_p^2\|X_s^2`$ $`=`$ $`{\displaystyle \frac{\delta _{ps}}{2\pi }}X_p^{}+X_p^2+[X_p^2,H_1]_{}\|X_s^2`$ (6) $`+`$ $`[X_p^2,H_1^{}]_{}\|X_s^2,`$ $`(E+\mu )X_p^0\|X_s^2`$ $`=`$ $`[X_p^0,H_1]_{}\|X_s^2+[X_p^0,H_1^{}]_{}\|X_s^2,`$ with $`[A,B]_{}=ABBA`$. To obtain the closed system of equations we apply new two-pole approximation, proposed in work . Suppose in Eq. (6) that $`[X_p^0,H_1]_{}={\displaystyle \underset{j}{}}ϵ(pj)X_j^0,`$ $`[X_p^2,H_1]_{}={\displaystyle \underset{j}{}}\stackrel{~}{ϵ}(pj)X_j^2,`$ (7) where $`ϵ(pj)`$ and $`\stackrel{~}{ϵ}(pj)`$ are non-operator expressions which we calculate using the method of work . At electron concentration $`n`$=1 in a paramagnetic state we have $`ϵ(pj)=(12d)t_{pj},`$ $`\stackrel{~}{ϵ}(pj)=(12d)t_{pj},`$ (8) with $`d=X_p^2`$ being the concentration of doublons. Let us take into account the functions $`[X_p^2,H_1^{}]_{}\|X_s^2`$ and $`[X_p^0,H_1^{}]_{}\|X_s^2`$ in the mean-field approximation: $`[X_p^2,H_1^{}]_{}\|X_s^2`$ $`=`$ $`{\displaystyle \underset{i,ip}{}}t_{ip}[(X_p^{}+X_p^2)X_i^0\|X_s^2+X_p^{02}X_i^2\|X_s^2`$ (9) $``$ $`X_p^{}X_i^0\|X_s^2]{\displaystyle }_{i,ip}t_{ip}X_p^{}+X_p^2X_i^0\|X_s^2,`$ $`[X_p^0,H_1^{}]_{}\|X_s^2`$ $`=`$ $`{\displaystyle \underset{i,ip}{}}t_{ip}[(X_p^0+X_p^{})X_i^2\|X_s^2+X_p^{02}X_i^0\|X_s^2`$ $``$ $`X_p^{}X_i^2\|X_s^2]{\displaystyle }_{i,ip}t_{ip}X_p^0+X_p^{}X_i^2\|X_s^2;`$ in this way we neglect the processes describing the “inter-band” hoppings of electrons which are connected with spin turning over and “inter-band” hoppings with creation or annihilation of two electrons on the same site. So we obtain the closed system of equations $`(E\mu +U)X_p^2\|X_s^2{\displaystyle \underset{i}{}}\stackrel{~}{ϵ}(pi)X_i^2\|X_s^2+X_p^{}+X_p^2{\displaystyle \underset{i,ip}{}}t_{ip}X_i^2\|X_s^2`$ $`={\displaystyle \frac{X_p^2+X_p^{}}{2\pi }}\delta _{ps},`$ (10) $`(E\mu )X_p^0\|X_s^2{\displaystyle \underset{i}{}}ϵ(pi)X_i^0\|X_s^2+X_p^0+X_p^{}{\displaystyle \underset{i,ip}{}}t_{ip}X_i^2\|X_s^2=0.`$ After the Fourier transformation we obtain solutions of system of Eqs. (2): $`X_p^2\|X_s^2_𝐤`$ $`=`$ $`{\displaystyle \frac{X_p^2+X_p^{}}{2\pi }}\left({\displaystyle \frac{A_𝐤^1}{EE_h(𝐤)}}+{\displaystyle \frac{B_𝐤^1}{EE_d(𝐤)}}\right),`$ (11) $`A_𝐤^1`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1{\displaystyle \frac{Uϵ(𝐤)+\stackrel{~}{ϵ}(𝐤)}{E_d(𝐤)E_h(𝐤)}}\right),B_𝐤^1=1A_𝐤^1,`$ $`X_p^0\|X_s^2_𝐤`$ $`=`$ $`{\displaystyle \frac{X_p^2+X_p^{}X_p^0+X_p^{}}{2\pi }}`$ (12) $`\times {\displaystyle \frac{t(𝐤)}{E_d(𝐤)E_h(𝐤)}}\left({\displaystyle \frac{1}{EE_h(𝐤)}}{\displaystyle \frac{1}{EE_d(𝐤)}}\right).`$ Here $`t(𝐤)`$ is the hopping intergral in $`𝐤`$representation and $`E_h(𝐤)`$ $`=`$ $`\mu +{\displaystyle \frac{U}{2}}+{\displaystyle \frac{ϵ(𝐤)+\stackrel{~}{ϵ}(𝐤)}{2}}`$ (13) $``$ $`{\displaystyle \frac{1}{2}}\sqrt{[Uϵ(𝐤)+\stackrel{~}{ϵ}(𝐤)]^2+X_p^0+X_p^{}X_p^{}+X_p^2(t(𝐤))^2},`$ $`E_d(𝐤)`$ $`=`$ $`\mu +{\displaystyle \frac{U}{2}}+{\displaystyle \frac{ϵ(𝐤)+\stackrel{~}{ϵ}(𝐤)}{2}}`$ (14) $`+`$ $`{\displaystyle \frac{1}{2}}\sqrt{[Uϵ(𝐤)+\stackrel{~}{ϵ}(𝐤)]^2+X_p^0+X_p^{}X_p^{}+X_p^2(t(𝐤))^2}`$ are the energies of electron in lower (“hole”) and upper (“doublon”) subbands, respectively; $`ϵ(𝐤)`$ and $`\stackrel{~}{ϵ}(𝐤)`$ are the Fourier components of $`ϵ(pj)`$ and $`\stackrel{~}{ϵ}(pj)`$. Analogous procedure gives for functions $`X_p^2\|X_s^0`$ and $`X_p^0\|X_s^0`$ the following expressions: $`X_p^2\|X_s^0_𝐤`$ $`=`$ $`X_p^0\|X_s^2_𝐤,`$ $`X_p^0\|X_s^0_𝐤`$ $`=`$ $`{\displaystyle \frac{X_p^0+X_p^{}}{2\pi }}\left({\displaystyle \frac{A_𝐤^2}{EE_h(𝐤)}}{\displaystyle \frac{B_𝐤^2}{EE_d(𝐤)}}\right),`$ (15) $`A_𝐤^2`$ $`=`$ $`B_𝐤^1,B_𝐤^2=A_𝐤^1.`$ Finally, in k-representation single-particle Green function (5) we obtain $`a_p\|a_s^+_𝐤={\displaystyle \frac{1}{2\pi }}\left({\displaystyle \frac{A_𝐤}{EE_h(𝐤)}}+{\displaystyle \frac{B_𝐤}{EE_d(𝐤)}}\right),`$ (16) $`A_𝐤={\displaystyle \frac{1}{2}}\left(1{\displaystyle \frac{(C_1C_2)(Uϵ(𝐤)+\stackrel{~}{ϵ}(𝐤))+4t(𝐤)C_1C_2}{E_d(𝐤)E_h(𝐤)}}\right),`$ $`B_𝐤=1A_𝐤,`$ where $`C_1=X_p^0+X_p^{}`$, $`C_2=X_p^2+X_p^{}`$. In the important for an investigation of metal-insulator transition case $`n=1`$ in a paramagnetic state ($`X_p^{}=X_p^{}`$) single-particle Green function (16) has the form $`a_p\|a_s^+_𝐤={\displaystyle \frac{1}{2\pi }}\left({\displaystyle \frac{A_𝐤}{EE_h(𝐤)}}+{\displaystyle \frac{B_𝐤}{EE_d(𝐤)}}\right),`$ (17) $`A_𝐤={\displaystyle \frac{1}{2}}\left(1{\displaystyle \frac{t(𝐤)}{\sqrt{U^2+(t(𝐤))^2}}}\right),`$ $`B_𝐤=1A_𝐤,`$ where single-particle energy spectrum is $`E_h(𝐤)=(12d)t(𝐤){\displaystyle \frac{1}{2}}\sqrt{U^2+(t(𝐤))^2},`$ $`E_d(𝐤)=(12d)t(𝐤)+{\displaystyle \frac{1}{2}}\sqrt{U^2+(t(𝐤))^2}`$ (18) (here we took into account that $`\mu =\frac{U}{2}`$ for $`n=1`$). Single-particle Green function (17) and energy spectrum (2) are exact in the band and atomic limits. It is worthwile to note, that in distinction from the results of two-pole approximations of Hubbard and Ikeda, Larsen, Mattuck the energy spectrum (2) depends on polar states concentration (thus on temperature). In distinction from approximations based on ideology of Roth (in this connection see also Refs. -) the energy spectrum (2) describes metal-insulator transition. Energy spectrum which describes metal-insulator transition was earlier obtained in work . Expressions (2) differs from the respective expressions in work by presence of term $`\sqrt{U^2+t^2(𝐤)}`$ instead of $`\sqrt{U^2+4d^2t^2(𝐤)}`$. This leads to the series of distinctions between results of this work and results of work ($`d(U/w)`$–dependence, the condition of metal-insulator transition, etc); at the same time expression (2) depends on polar state concentration similarly to respective expression in work . ## 3 Energy gap and polar states concentration The energy gap (difference of energies between bottom of the upper and top of the lower Hubbard bands) is given by $`\mathrm{\Delta }E=E_d(w)E_h(w)=2w(12d)+\sqrt{U^2+w^2},`$ (19) (where $`w=z\|t\|`$ is the halfwidth of uncorrelated electron band, $`z`$ is the number of nearest neighbors to a site). Expression (19) describes the vanishing of the energy gap in the spectrum of paramagnetic insulator at critical value $`(\frac{U}{w})_c`$ when the halfbandwidth $`w`$ increase (under pressure). Dependence of $`\mathrm{\Delta }E`$ on temperature can lead to the transition from metallic to insulator state with increase of temperature (in this connection note the transition at increasing temperatures from the state of paramagnetic metal to the paramagnetic insulator state in the systems NiS<sub>2-x</sub>Se<sub>x</sub>, (V<sub>1-x</sub>Cr<sub>x</sub>)<sub>2</sub>O<sub>3</sub> and Y<sub>1-x</sub>Ca<sub>x</sub>TiO<sub>3</sub>). For the calculation of polar states concentration we use function (11). At $`T=0`$ and rectangular density of states the concentration of polar states is $`d={\displaystyle \frac{1}{4}}+{\displaystyle \frac{U}{8w}}\mathrm{ln}\left({\displaystyle \frac{14d}{34d}}\right)`$ (20) if $`(\frac{U}{w})(\frac{U}{w})_c`$ and $`d={\displaystyle \frac{1}{4}}+{\displaystyle \frac{U}{8w}}\mathrm{ln}\left({\displaystyle \frac{\sqrt{1+\frac{U}{w}^2}+1}{\sqrt{1+\frac{U}{w}^2}1}}\right)`$ (21) if $`(\frac{U}{w})>(\frac{U}{w})_c`$. At $`T=0`$ we have $`(\frac{U}{w})_c=1.672`$. The dependence $`d(\frac{U}{w})`$ given by Eqs.(20)-(21) is plotted on Fig. 1. One can see that in the point $`(\frac{U}{w})_c`$ the slope of $`d(\frac{U}{w})`$dependence changes; the concentration of doublons vanishes at $`\frac{U}{w}\mathrm{}`$. Our result for $`d(\frac{U}{w})`$ in region of MIT is in good agreement with result of papers obtained in the limit of infinite dimensions (Fig. 2). The parameter $`U`$ is normalized by averaged band energy in absence of correlation $`\epsilon _0`$. In Fig. 3 the dependences of polar states concentration on parameter $`\frac{U}{w}`$ at different temperatures are presented. Note the important difference (see Fig.4) of the dependence of $`d`$ on temperature from result of papers : we found that at any temperature polar states concentration increases monotonically with increasing temperature at the fixed value of $`\frac{U}{w}`$ when respective dependence in has a minimum. The dependence of $`\frac{\mathrm{\Delta }E}{U}`$ on parameter $`\frac{U}{w}`$ at zero temperature is plotted in Fig. 5. It is important to note that in the point of gap disappearence $`d0`$ in contrast to the previously obtained result . At increasing $`\frac{U}{w}`$ the energy gap width increases (the negative values of $`\mathrm{\Delta }E`$ correspond to the overlapping of the subbands). For comparison on Fig.5 results of approximation “Hubbard-I” is also plotted. In the point of energy gap vanishing $`(\frac{U}{w})_c=1.672`$ what is very close to result of “Hubbard-III” approximation . At increase of temperature in metallic state the overlapping of subbands decreases and temperature induced transition from metallic to insulating state can occur at some values of parameter $`\frac{U}{w}`$ (Fig. 6). The obtained results allows us to draw the $`(w/U,T)`$ phase diagram of the model (Fig. 7). This phase diagram can explain the experimentally observed transitions from metallic to insulating state with increase of temperature and from insulating to metallic state with increase of bandwidth (under external pressure) in paramagnetic state. ## 4 Ground state energy The ground state energy of the model $`{\displaystyle \frac{E_0}{N}}={\displaystyle \frac{1}{N}}{\displaystyle \underset{ij\sigma }{}}t_{ij}a_{i\sigma }^+a_{j\sigma }+Ud,`$ (22) calculated using single particle Green function (17) and expressions (20)-(21) for the concentration of polar states has the form: $`{\displaystyle \frac{E_0}{N}}={\displaystyle \frac{w}{2}}+{\displaystyle \frac{U}{4}}(1+3d){\displaystyle \frac{U^2}{2w}}{\displaystyle \frac{(14d)}{4(12d)^21}}`$ (23) if $`(\frac{U}{w})(\frac{U}{w})_c`$ and $`{\displaystyle \frac{E_0}{N}}={\displaystyle \frac{1}{2}}\sqrt{U^2+w^2}+2U({\displaystyle \frac{1}{4}}d)`$ (24) if $`(\frac{U}{w})>(\frac{U}{w})_c`$. In Fig. 8 the dependence of the ground state energy on parameter $`\frac{U}{w}`$ given by Eqs.(23)-(24) is compared with the exact result, found in one-dimensional case . The upper and lower bounds on ground state energy in one-dimensional case found in paper are also shown. Our result for the ground state energy in metallic state lies slightly lower than exact one and in insulator state fits the exact ground state energy very well. In Fig. 9 our plot of the ground state energy is compared with the best upper and lower bounds on ground state energy in infinite-dimensional case . In Fig. 10 we have the comparison with bounds on ground state energy for three-dimensional simple cubic lattice obtained in paper . In Figs. 8-10 the ground state energy per electron is normalized by averaged band energy in absence of correlation $`\epsilon _0`$; in considered case and rectangular density of states $`\epsilon _0=\frac{w}{2}`$. Figs. 8-10 show that our result present a good approximation for the ground state energy of the system. In Fig. 11 we plot our result for the kinetic part of ground state energy. This plot describes the same behavior of kinetic energy of electrons with change of correlation strength in paramagnetic state as respective result of work : in metallic state absolute value of kinetic energy decrease rapidly due to rapid decrease of doublon (hole) concentration. In insulating state absolute value of kinetic energy decrease slowly what in the approximation of effective Hamiltonian (obtained for the case $`\frac{t_{ij}}{U}1`$) is equivalent to the interaction of local magnetic moments. ## 5 Conclusions In this paper we have studied electron correlations in narrow energy bands using recently proposed approximation . We assume that state of the narrow-band system is paramagnetic insulator or paramagnetic metal. The single-particle Green function and energy spectrum dependent on model parameters and on polar states concentration (thus on temperature) have been found in paramagnetic state at half-filling ($`n=1`$). The obtained expression for energy gap allows to describe MIT at changes of bandwidth (under external pressure) or temperature. The comparison of calculated ground state energy with results of other approximations and the exact result found in one-dimensional case shows that the used method is a good approximation for the model under consideration. The obtained phase diagram of the model can explain the transitions from paramagnetic metal state to paramagnetic insulator state at increase of temperature and the paramagnetic insulator - paramagnetic metal transitions under external pressure observed in the systems NiS<sub>2-x</sub>Se<sub>x</sub>, (V<sub>1-x</sub>Cr<sub>x</sub>)<sub>2</sub>O<sub>3</sub> and Y<sub>1-x</sub>Ca<sub>x</sub>TiO<sub>3</sub>. It is worthwhile to note that approximation used in this paper can be generalized to describe effects of antiferromagnetic ordering. Such a generalization will be considered in subsequent paper. Figure captions Fig.1 The dependence of doublon concentration $`d`$ on $`U/w`$ at zero temperature. Fig.2 The comparison of $`d(U/w)`$ dependences: solid line - our result, dashed line - iterative-perturbative theory , circles - QMC method . Fig.3 The dependences of doublon concentration $`d`$ on $`U/w`$ at different temperatures: lower curve corresponds to $`kT/w=0.16`$, middle curve corresponds to $`kT/w=0.08`$, upper curve corresponds to $`kT/w=0`$. Fig.4 The dependences of doublon concentration $`d`$ on temperature at different $`U/w`$: values of $`U/w`$ from down to up are 3, 2, 1.5, 1, 0.5, 0 Fig.5 The dependences of energy gap width on $`U/w`$: “Hubbard-I” approximation (upper curve), our result (middle curve), approximation (lower curve). Fig.6 The dependences of energy gap width on temperature at different $`U/w`$: values of $`U/w`$ from down to up are 0.5, 1.2, 1.5 Fig.7 The obtained ($`kT,w/U`$)-phase diagram of the model. Fig.8 The comparison of ground state energies in one-dimensional case: dashed curves correspond to upper and lower bounds given by Langer and Mattis , upper solid curve corresponds to exact ground state (Lieb and Wu ) lower solid curve corresponds to result of this paper. Fig.9 The ground state energy found in this paper (upper curve), best upper (middle curve) and lower (lower curve) bounds on ground state energy in infinite-dimensional case. Fig.10 The upper (upper curve) and lower (lower curve) bounds on ground state energy in three-dimensional case and the ground state energy found in this paper (middle curve). Fig.11 The kinetic part of ground state energy as a function of $`U/w`$.
warning/0005/hep-ph0005034.html	ar5iv	text	# 1 Introduction ## 1 Introduction Lattice Gauge Theory has been applied to the study of the strong interaction in earnest for the last 20 years or so. In that time, it has grown from an fledgling, optimistic area of research to a well-developed, mature (and optimistic!) discipline. The reason for this optimism is that all the approximations involved in the technique are systematically improvable. This means that, given enough time on a computer powerful enough, we will have predictions for the QCD bound-state spectrum (for example) with arbitrary accuracy. Other approaches of studying the strong interaction, while they may have their own advantages, cannot make this claim. In this talk I will begin by overviewing the lattice method for obtaining hadronic spectral quantities via the calculation of $`n`$point functions. I will outline the caveats that exist with current lattice simulations as well as detailing some of the successes of the approach. The vexed question of performing “full” QCD calculations (i.e. without the quenched approximation) is discussed and I will explain why these simulations are even more difficult than they at first seemed. In Sec. 3, I will detail some recent results from the UKQCD collaboration, focusing on our recent full QCD work. Sec. 4 outlines a new and promising approach of obtaining physical results from the lattice by using the spectral function representation. This is an exciting area of research, and, if it reaches its potential, promises to become a standard approach in the future. I then summarise the main points raised in this talk in Sec. 5. ## 2 Overview of Lattice Gauge Theory ### 2.1 The Method I attempt here to describe the method normally used in lattice calculations of the hadronic spectrum of QCD. There are many excellent reviews of this topic which cover the approach in more detail. The conventional lattice approach to QCD spectrum calculations is performed in the Euclidean path integral framework. It requires the calculation of $`n`$-point correlation functions, $`G_n(t)`$, of hadronic interpolating operators, $`J`$, in a background sea of glue (and sea quarks in the case of full, i.e. “unquenched” calculations). A Monte Carlo approach is used to generate these background configurations with the appropriate Boltzmann factor $`e^𝒮`$ where $`𝒮`$ is the Euclidean action. There is an obvious analogy between this Lagrangian approach to Lattice Gauge Theory, and statistical mechanics: clearly the lattice path integral corresponds to the partition function of statistical mechanics. Note that in lattice simulations there is freedom to choose different quark fields in the interpolating operator $`J`$ compared to those in the Lagrangian. Hence we are able to distinguish between “valence” and “sea” quark masses, $`m^{val}`$ and $`m^{sea}`$. Because the QCD Lagrangian contains fermionic fields that appear quadratically, they can be integrated out analytically giving the usual determinant factor, $`det(D/+m)`$. Including this factor is a technical headache because it involves a huge increase in computational requirements. The usual way around this is to invoke the quenched approximation where two things are done: (i) the fermionic determinant is replaced by unity; and (ii) the gauge coupling, $`g`$, is rescaled so that the physical predictions for a particular test quantity (like the rho mass, for example) is in agreement with its experimental value. The extent to which the quenched approximation reproduces physical predictions for other quantities is a measure of its success. It is a remarkable fact that for many spectral quantities, the quenched predictions agree with the experimental (i.e. the “full” QCD) values to within 10% (see ). Despite the success of the quenched approximation, it is obviously essential to perform full QCD calculations in order to study the real world. Furthermore, for some quantities (such as the deconfinement temperate) the quenched approximation is poor, and for others (such as the $`\eta ^{}`$ mass) the quenched approximation fails completely. In the calculation of the two-point function, $`G_2(t)`$, if the exact operator, $`J^{exact}`$ for the hadron in question was used, then $`G_2(t)`$ would contain information on that hadron and no others. However, since $`J^{exact}`$ is not known, $`G_2(t)`$ receives contributions from all hadronic states which have non-zero overlap with $`J`$. It is straightforward to show that, in this case, $`G_2(t)`$ has the following form $$G_2(t)=\underset{i}{}Z_ie^{M_it}$$ where the sum is over the hadronic states $`i`$, and $`M_i`$ and $`Z_i`$ are the corresponding hadronic mass and overlap. Note that since the calculation is performed in Euclidean space-time, the excited states are exponentially suppressed with respect to the fundamental state (i.e. the exponentials have real arguments). This means that $`G_2(t)`$ asymptotes to the two-point function of the ground state hadron as $`t\mathrm{}`$. Obviously the parameters of the ground state, in particular the mass $`MM_0`$, can be extracted by fitting $`G_2(t)`$ to an exponential, $`e^{Mt}`$, for $`t`$ sufficiently large. ### 2.2 Caveats In this sub-section I explain the caveats that one must apply to any calculation using the method described above. The main point to make is that the calculations of the hadronic properties are performed with input parameter values which are not those of the real world. Specifically, these parameters are * valence quark mass(es), $`m^{val}`$, (typically $`m^{val}>\frac{1}{2}m^{strange}`$); * sea quark mass(es), $`m^{sea}`$, (typically $`m^{sea}>m^{strange}`$); * lattice volume, $`V`$, (typically $`V<(2fm)^3`$); * lattice spacing, $`a`$, or, through dimensional transmutation, $`g_0`$, (typically $`a>\mathrm{\hspace{0.33em}.05}fm`$); * number of dynamical fermion flavours, $`N_f`$. Quenching corresponds to $`N_f=0`$. (Typically, $`N_f=0`$ or $`2`$.) Therefore the mass, $`M`$, which was obtained using the procedure above is not the mass of the real world hadron, but the mass of the corresponding hadron in a world where the quark masses are the same as those input into the lattice calculation, the volume is the finite volume used in the simulation etc. So, strictly speaking, $`M`$ is a function of the above 5 input parameters. The final prediction of the real world value, $`M^{expt}`$, should be obtained by the following extrapolations: $`\sqrt{\sqrt{}}`$ $`m^{val}`$ $`\text{few}MeV`$ $`\sqrt{}`$ $`m^{sea}`$ $`\text{few}MeV`$ $`\sqrt{\sqrt{}}`$ $`V`$ $`\mathrm{}`$ $`\sqrt{\sqrt{}}`$ $`a`$ $`0`$ (1) $``$ $`N_f`$ $`\text{}2{\displaystyle \frac{1}{2}}\text{}\text{i.e. two light flavours for }u,d\text{ and one heavier for }s`$ Note that the limit $`m^{sea}=\mathrm{}`$ corresponds to $`N_f=0`$ so the $`m^{sea}`$ and $`N_f`$ extrapolations are not independent. While these extrapolations muddy the water significantly, many of them are theoretically well-understood and numerically under control. The number of $``$’s which appear next to the extrapolations is an indication of how well the extrapolation is under control. As an example of the quality of the extrapolations, Fig. 1 shows an extrapolation of the nucleon and vector meson masses as a function of the lattice spacing, $`a`$, taken from . In these plots, data points obtained with both the Wilson and various improved lattice actions (designed to have discretisation errors smaller than $`𝒪(a)`$) are shown. Fits to the relevant functional forms are included in the figure and the symbols on the left of the plot are the continuum extrapolations (which clearly agreement with each other). ### 2.3 Successes The caveats listed in the previous sub-section do not hinder the success of the Lattice technique as a means of obtaining accurate predictions from the strong interaction. To give an example of the current status of lattice calculations, Fig. 2 shows the hadronic spectrum obtained by the CP-PACS collaboration using the quenched approximation. There are two features to note. The error bars in the predictions are tiny ( $``$$`<`$ 3%) - which is a clear measure of the success of the lattice technique. Secondly, there is a small, but statistically significant discrepancy between the lattice predictions and the experimental numbers - which is a signal that unquenching is required in order to make further progress. While the quenched lattice calculations have clearly matured with precision estimates of many quantities of only a few percent, serious calculations involving full QCD have only recently begun. Typically the current errors in these calculations are several times that of equivalent quenched results. ## 3 Recent Dynamical Results from the UKQCD Collaboration In this section I review some of the recent results from the UKQCD Collaboration’s dynamical simulations. ### 3.1 Dependency on Sea Quark Mass As was outlined in Sec. 2.2 all lattice predictions are functions of $`V`$, $`a`$, $`N_f`$ etc. In this section, the functional dependency on $`m^{sea}`$ is discussed. Fig. 3 (taken from ) shows how the lattice spacing, obtained from the hadronic length scale, $`r_00.5fm`$, depends strongly on the sea quark mass, $`m^{sea}`$ (here expressed in terms of the hopping parameter $`\kappa ^{sea}`$). These calculations were performed at fixed bare coupling, $`g_0`$. This effect has important consequences. Simulations at a fixed bare coupling, $`g_0`$, and with several values of $`m^{sea}`$, correspond to different physical volumes, and, furthermore, different points on the continuum extrapolation $`a0`$. Thus finite volume and $`𝒪(a)`$ systematics are mixed in this case. UKQCD have performed two types of calculation. The first was at fixed $`g_0`$ for various values of $`m^{sea}`$ in order to calibrate this effect. We then performed more sophisticated simulations at several values of the parameter pairs $`(g_0,m^{sea})`$ which were chosen in order to maintain fixed lattice spacing $`a0.11fm`$ (and therefore also fixed volume). This second calculation utilised the “matching” technology of . In both these calculations, an improved action was used in order to reduce the effect of $`𝒪(a)`$ errors. ### 3.2 Results There is not space to present full details of UKQCD’s recent unquenched calculations. I discuss the results of only two quantities, and refer the reader to the original papers for full details. One of the benchmark quantities of lattice calculations is the static quark potential. Fig. 4 shows UKQCD’s result for this quantity in units of $`r_0`$ from . It can be noted that there is little immediate dependency on $`m^{sea}`$ in this quantity.<sup>1</sup><sup>1</sup>1A closer look however at the data points close to the origin shows a systematic effect which can be interpreted as different runnings of the coupling as function of $`m^{sea}`$. Fig. 5 shows the vector meson mass, $`M_V`$, plotted against the pseudoscalar mass squared, $`M_{PS}^2`$ from . Again, these quantities are expressed in units of $`r_0`$. Quenched data which corresponds to the same lattice spacing as the unquenched simulations are included as a comparison. The experimental points corresponding to the strange mesons are also plotted. In summary, the results of indicate that the effects of unquenching are small for these values of $`m^{sea}>m^{strange}`$. This motivates the need to move to more physical values of $`m^{sea}\text{few }GeV`$ in the future. ## 4 Lattice Spectral Functions As outlined in Sec. 2.1, the conventional method of determining ground state properties from lattice simulations is by fitting exponentials to the tails of $`n`$-point functions. There have been several attempts at developing other strategies for uncovering spectral quantities from lattice data. These all revolve around the spectral function (SF) $`\rho (s)`$ which can be defined through $$G_2(t)=_0^{\mathrm{}}K(t,s)\rho (s)𝑑s,$$ (2) where $`K(t,s)`$ is the kernel function - typically just $`e^{st}`$ for this work. The SF contains much richer information on the channel being considered than just the ground state parameters. It also has the advantage that theoretical input can be used to guide its form for large values of $`s`$ where perturbation theory is valid. In fact, this was the approach taken by who used (continuum) perturbation theory to derive a functional form for $`\rho ^{PT}(s)`$. This $`\rho ^{PT}(s)`$ was then used above a certain threshold of energy $`s_0`$ and a $`\delta `$ function used for the ground state following the approach of QCD Sum Rules. A very new and promising technique takes the marriage of lattice data and spectral functions one step further. This approach uses the lattice data itself to determine the SF by inverting eq.(2). This is a very numerically technical approach; it is in fact an “ill-posed” problem - $`G_2(t)`$ is known only at a small number of discrete values of $`t`$, whereas the aim is to determine $`\rho (s)`$ for a large number of values of $`s`$ (ideally for a continuous range in $`s`$). In a “maximum entropy method” is employed to overcome these hurdles. ## 5 Conclusion I have given a brief overview of the current state of play of lattice gauge theory calculations of the hadronic spectrum. The problems of extrapolating the input parameters of any lattice simulation to their physical value is emphasised. In particular I have shown that interpreting results from naive unquenched calculations can prove difficult due to the dependency of the lattice spacing on $`m^{sea}`$. A summary has been given of recent unquenched results from the UKQCD collaboration for the static quark potential and vector meson mass. Finally the new and interesting method of using spectral functions in the analysis of lattice data is discussed. ## Acknowledgements I am grateful to my colleagues in the UKQCD collaboration especially Joyce Garden, Balint Joo, Alan Irving and Ken Bowler. I acknowledge the support of the Particle Physics & Astronomy Research Council and the Royal Society.
warning/0005/cond-mat0005466.html	ar5iv	text	# Spin and charge excitations in incommensurate spin density waves ## Abstract Collective excitations both for spin- and charge-channels are investigated in incommensurate spin density wave (or stripe) states on two-dimensional Hubbard model. By random phase approximation, the dynamical susceptibility $`\chi (𝐪,\omega )`$ is calculated for full range of $`(𝐪,\omega )`$ with including all higher harmonics components. An intricate landscape of the spectra in $`\chi (𝐪,\omega )`$ is obtained. We discuss the anisotropy of the dispersion cones for spin wave excitations, and for the phason excitation related to the motion of the stripe line. Inelastic neutron experiments on Cr and its alloys and stripe states of underdoped cuprates are proposed. Recently a remarkable series of elastic neutron experiments has been performed on La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> . It reveals static incommensurate spin density wave (ISDW) structure in underdoped regions. The superconducting samples show that the modulation vector $`𝐐`$ are characterized by $`𝐐=2\pi (\frac{1}{2},\frac{1}{2}\pm \delta )`$ or $`2\pi (\frac{1}{2}\pm \delta ,\frac{1}{2})`$ in momentum space. The incommensurate modulation runs vertically in the $`a`$\- or $`b`$-axis of the CuO<sub>2</sub> plane. The static stripe structures have been observed also in superconducting (La, Nd)<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> and insulating La<sub>2-x</sub>Sr<sub>x</sub>NiO<sub>4</sub> . Since the establishment of static orders on these systems, the study of spin and charge dynamics; $`\chi (𝐪,\omega )`$ has been just started experimentally. Prior to these studies on cuprates, ISDW is well known in itinerant electron systems, such as a typical example Cr and its alloys for long time. While its static properties are fairly well understood, its dynamical properties remain largely unexplored . A part of reasons stems from the fact that the theoretical description is underdeveloped, thus the interplay between theory and experiment was not satisfactory. The intensive efforts for extracting spin and charge excitations in Cr and its alloys and also high $`T_\mathrm{c}`$ cuprates are strongly motivated by a hope that knowing the collective excitations or fluctuation spectrum in strongly correlated systems may lead us to a clue to understanding the mechanism of high $`T_\mathrm{c}`$ superconductivity. There exist a lot of theoretical works on dynamical properties, beginning from a seminal paper by Fedders and Martin to recent work by Fishman and Liu. The former first derives spin wave mode in the transverse spin excitation of itinerant electron systems. The latter investigates various transverse and longitudinal modes at the $`𝐐`$ position, based on a one-dimensional (1D) model within the random phase approximation (RPA). Their theory takes account of only fundamental order parameter $`\mathrm{\Delta }_𝐐`$, neglecting higher harmonics $`\mathrm{\Delta }_{2𝐐}`$, $`\mathrm{\Delta }_{3𝐐}`$, $`\mathrm{}`$ associated with the incommensurability $`\delta `$. By extending their work , we calculate the dynamical spin and charge susceptibilities $`\chi (𝐪,\omega )`$ in the full range of the two dimensional (2D) wave number $`𝐪`$ for entire Brillouin zone and the excitation energy $`\omega `$ up to the band width. We take account of all the possible higher harmonics. This will turn out to be extremely crucial in correctly evaluating these quantities. The multi-dimensional calculation here allows us to extract a wealth of information on stripe motions such as translation, or meandering, etc. and on the anisotropy of excitation cones. This kind of calculation has not been done before to our knowledge. We start with the Hubbard model on 2D square lattice: $`H=t_{i,j,\sigma }C_{i,\sigma }^{}C_{j,\sigma }+U_in_in_i`$. To calculate $`\chi (𝐪,\omega )`$, we first set up the incommensurate SDW ground state. Assuming a periodic spin and associated charge orderings, we introduce the order parameter $`n_{i,\sigma }=_le^{\mathrm{i}l𝐐𝐫_𝐢}n_{l𝐐,\sigma }`$ with $`\delta =1/N`$. In $`N`$-site periodic case, the Brillouin zone is reduced to $`1/N`$-area. and energy dispersion is split to $`N`$ bands. We write $`𝐤=𝐤_0+m𝐐`$ ($`m=0,1,\mathrm{},N1`$), where $`𝐤_0`$ is restricted within the reduced Brillouin zone. Then the Hamiltonian is reduced to $`H={\displaystyle \underset{𝐤_0,\sigma }{}}{\displaystyle \underset{m,n}{}}C_{𝐤_0+m𝐐,\sigma }^{}(\widehat{H}_{𝐤_0,\sigma })_{m,m^{}}C_{𝐤_0+m^{}𝐐,\sigma }`$ (1) $`={\displaystyle \underset{𝐤_0,\sigma ,\alpha }{}}E_{𝐤_0,\sigma ,\alpha }\gamma _{𝐤_0,\sigma ,\alpha }^{}\gamma _{𝐤_0,\sigma ,\alpha }.`$ (2) The $`N\times N`$ Hamiltonian matrix $`(\widehat{H}_{𝐤_0,\sigma })_{m,m^{}}=ϵ(𝐤_0+m𝐐)\delta _{m,m^{}}+Un_{(mm^{})𝐐,\sigma }`$ is diagonalized by a unitary transformation $`C_{𝐤_0+m𝐐,\sigma }=_\alpha u_{𝐤_0,\sigma ,\alpha ,m}\gamma _{𝐤,\sigma ,\alpha }`$. The calculation is iterated until all the order parameters satisfy the self-consistent condition $`n_{lQ\sigma }=N_k^1_{k_0,m,\alpha }u_{𝐤_0,\sigma ,\alpha ,m}^{}u_{𝐤_0,\sigma ,\alpha ,m+l}f(E_{𝐤_0,\sigma ,\alpha })`$. Here, $`N_k=_{𝐤_0,m}1`$. We construct the thermal Green function as $$g_\sigma (𝐫,𝐫^{},\mathrm{i}\omega _n)=\underset{𝐤_0,\alpha }{}\frac{u_{𝐤_0,\sigma ,\alpha }(𝐫)u_{𝐤_0,\sigma ,\alpha }^{}(𝐫^{})}{\mathrm{i}\omega _nE_{𝐤_0,\sigma ,\alpha }},$$ (3) with $`u_{𝐤_0,\sigma ,\alpha }(𝐫_i)=N_k^{1/2}_m\mathrm{e}^{\mathrm{i}(𝐤_0+m𝐐)𝐫_i}u_{𝐤_0,\sigma ,\alpha ,m}`$, and evaluate the dynamical susceptibilities; the spin longitudinal mode $`\chi _{zz}(𝐪,\omega )=S_z;S_z_{𝐪,\omega }`$, the transverse one $`\chi _{xx}(𝐪,\omega )=S_x;S_x_{𝐪,\omega }`$, and the charge susceptibility $`\chi _{nn}(𝐪,\omega )=n;n_{𝐪,\omega }`$. The Fourier transformation of Eq. (3) is given by $`g_\sigma (𝐤+l_1𝐐,𝐤+l_2𝐐,\mathrm{i}\omega _n)`$ (4) $`={\displaystyle \frac{1}{N_k}}{\displaystyle \underset{𝐫_1,𝐫_2}{}}e^{\mathrm{i}(𝐤+l_1𝐐)𝐫_1}e^{\mathrm{i}(𝐤+l_2𝐐)𝐫_2}g_\sigma (𝐫_1,𝐫_2,\mathrm{i}\omega _n)`$ (5) $`={\displaystyle \underset{\alpha }{}}{\displaystyle \frac{u_{𝐤,\sigma ,\alpha ,l_1}u_{𝐤,\sigma ,\alpha ,l_2}^{}}{\mathrm{i}\omega _nE_{𝐤,\sigma ,\alpha }}},`$ (6) where $`l_1,l_2=0,1,\mathrm{},N1`$. In the presence of the order parameter $`n_{l𝐐,\sigma }`$, the incoming- and outgoing-momentum of $`g_\sigma `$ can differ by $`l𝐐=(l_1l_2)𝐐`$. Then $`g_\sigma (𝐤+l_1𝐐,𝐤+l_2𝐐,\mathrm{i}\omega _n)`$ is an $`N\times N`$ matrix with indexes $`l_1`$ and $`l_2`$. The bare susceptibility is given by $`\chi _0^{\sigma \sigma ^{}}(𝐫_1,𝐫_2,\mathrm{i}\mathrm{\Omega }_n)=T_{\omega _n}g_\sigma (𝐫_1,𝐫_2,\mathrm{i}\omega _n+\mathrm{i}\mathrm{\Omega }_n)g_\sigma ^{}(𝐫_2,𝐫_1,\mathrm{i}\omega _n).`$ Its Fourier transformation $`\chi _0^{\sigma \sigma ^{}}(𝐤+l_1𝐐,𝐤+l_2𝐐,\mathrm{i}\omega _n)`$ is also $`N\times N`$ matrix . We use the analytic continuation $`\mathrm{i}\omega _n\omega +\mathrm{i}\eta `$. Typically, we use $`\eta =0.001t`$ in our numerical calculation. The RPA equation for $`\chi _{S_{}S_{}}=S_{};S_{}`$ is written as $`\chi _{S_{}S_{}}(𝐫_1,𝐫_3,\omega )=\chi _0^{}(𝐫_1,𝐫_3,\omega )`$ (7) $`+U{\displaystyle \underset{𝐫_2}{}}\chi _0^{}(𝐫_1,𝐫_2,\omega )\chi _{S_{}S_{}}(𝐫_2,𝐫_3,\omega ).`$ (8) After Fourier transformation to $`𝐤`$-space, Eq. (8) is reduced to a matrix equation of $`\chi _0^{}`$ and $`\chi _{S_{}S_{}}`$. By solving it, we obtain $`\chi _{xx}(𝐤+l_1𝐐,𝐤+l_2𝐐,\omega )`$. In a similar manner, we calculate $`n_{};n_{}`$ and $`n_{};n_{}`$, and obtain $`\chi _{zz}(𝐤+l_1𝐐,𝐤+l_2𝐐,\omega )`$ and $`\chi _{nn}(𝐤+l_1𝐐,𝐤+l_2𝐐,\omega )`$. The neutron scattering experiments observe the imaginary part of the dynamical susceptibility $`\chi ^{\prime \prime }(𝐪,\omega )=\mathrm{Im}\chi (𝐪,𝐪,\omega )`$. Since the signal is observed as spatial average, we consider the diagonal part of $`\chi (𝐤+l_1𝐐,𝐤+l_2𝐐,\omega )`$. According to standard linear response theory, the spatio-temporal oscillation pattern of a collective mode can be analyzed by $`\chi `$. In the presence of an infinitesimal external field $`h_O^{}(𝐤^{},\omega )`$ coupled to an operator $`O^{}(=S_x,S_z,n)`$, the response of the operator $`O`$ is given by $`\delta O(𝐤_0+l_1𝐐,\omega )=_{l_2}\chi _{OO^{}}(𝐤_0+l_1𝐐,𝐤_0+l_2𝐐,\omega )h_O^{}(𝐤_0+l_2𝐐,\omega ).`$ When the external field is a plane wave; $`h_O^{}(𝐫,t)=\overline{h}_O^{}\mathrm{e}^{\mathrm{i}(𝐪𝐫\omega t)}`$ with a small amplitude $`\overline{h}_O^{}`$, the response is given by $`\delta O(𝐫,t)={\displaystyle \underset{l}{}}\chi _{OO^{}}(𝐪+l𝐐,𝐪,\omega )\mathrm{e}^{\mathrm{i}(𝐪+l𝐐)𝐫}\overline{h}_O^{}\mathrm{e}^{𝐢\omega t}.`$ (9) (10) We consider the vertical stripe case for $`U=4t`$ and the hole density $`n_\mathrm{h}=1/8`$ as a representative case (The enrgy is scaled by $`t`$ from now on). As we do not include the nearest neighbor hopping $`t^{}`$, the lowest energy ground state is an SDW-gapped insulator with $`\delta =n_\mathrm{h}/2`$, i.e. $`N=16`$. The detailed ground state properties are reported previously . For our parameters, the single particle SDW gap $`E_g=0.41`$. In the ISDW state, the spatial profile is characterized by a distorted sinusoidal, or soliton form with a midgap band. The higher harmonics are determined as $`\|M_{l𝐐}/M_𝐐\|=`$0.08 ($`l=3`$), 0.005 ($`l=5`$), $`2.3\times 10^3`$ ($`l=7`$), where $`M_{l𝐐}=n_{l𝐐,}n_{l𝐐,}`$. In the limit of the half-filling $`(N\mathrm{})`$, the ratio $`\|M_{(2n+1)𝐐}/M_𝐐\|`$ increases and approaches $`(2n+1)^1`$, since the profile of the spin structure approaches square wave form . The structure factor has spots at $`(2n+1)𝐐`$ in the spin structure, and at $`2n𝐐`$ in the charge structure. These spots are observed by the elastic neutron scattering in Cr and its alloys . Their positions in the momentum space are shown in Fig. 1. Let us start with the excitation of the spin transverse mode. In Fig. 2, we show $`\chi _{xx}^{\prime \prime }(𝐪,\omega )`$ along paths A and B of Fig. 1. The gapless spin wave modes emanate not only $`𝐐`$, but also from $`3𝐐`$ and other odd harmonics. The ridge of $`\chi ^{\prime \prime }(𝐪,\omega )`$ shows singularity reflecting the dispersion of the collective mode. These modes at $`(2n+1)𝐐`$ have an identical dispersion relation, because every $`(2n+1)𝐐`$-modes couple each other in the RPA equation (8). The same dispersion pattern appears in each reduced Brillouin zone. But their intensities are different. With increasing $`l`$, the intensity is decreased as follows, $`\chi _{xx}^{\prime \prime }(l𝐐,0.025t)/\chi _{xx}^{\prime \prime }(𝐐,0.025t)=6.5\times 10^3`$ ($`l=3`$), $`2.7\times 10^5`$ ($`l=5`$), $`5.8\times 10^6`$ ($`l=7`$). It is found that the intensity ratio is obeyed $$\frac{\chi _{xx}^{\prime \prime }((2n+1)𝐐,\omega 0)}{\chi _{xx}^{\prime \prime }(𝐐,\omega 0)}\left[\frac{M_{(2n+1)𝐐}}{M_𝐐}\right]^2.$$ (11) With approaching the half-filling ($`N\mathrm{}`$) the above ratio is expected to increase and approach $`(2n+1)^2`$. Then, the higher harmonics spot at $`(2n+1)𝐐`$ may have enough intensity to be observed near half-fillings such as in Cr or underdoped cuprates. We analyze the oscillation pattern by Eq. (10). The response of $`\delta S_x(𝐫,t)`$ shows the same spin pattern as $`S_z`$ of the ground state at $`\omega 0`$ for $`(2n+1)𝐐`$. It means that the ground state spin structure is rotated as it is without modulation, since it is a Goldstone mode. The external field of the wave number $`(2n+1)𝐐`$ couples to $`M_{(2n+1)𝐐}`$ and makes the ground state spin structure rotate. Then, we can conclude that the spin transverse mode is a spin wave. With increasing $`\omega `$ along the dispersion curve, the spin wave oscillation shows the deviation from the spin pattern of the ground state, reflecting the wave number of the external field. The reconnection of the dispersion curve occurs at the reduced Brillouin zone boundary. Then, instead of the simple intersection of two dispersions, the small gap appears at $`\omega 0.3t`$ in Fig. 2. With increasing $`\omega `$, the intensity of $`\chi _{xx}^{\prime \prime }(𝐪,\omega )`$ decreases as $`1/\omega `$ along the dispersion curve, except for the weakened intensity at the gap position. For $`\omega >E_g`$, there exists other modes reflecting the fluctuation of the magnetic moment amplitude. It is a character of itinerant magnets and absent in localized spin magnets. Figure 3 shows the dispersion curve along the path A ($`q_y`$-direction) and B ($`q_x`$-direction) near $`𝐐`$ in Figs. 2 (a) and (b). The spin wave velocity $`v_{\mathrm{spin}}^x`$ ($`v_{\mathrm{spin}}^y`$) is defined by the slope of the $`q_x`$\- ($`q_y`$-) direction at $`\omega 0`$. The spin modulation parallel to the stripe (domain wall) corresponds to $`v_{\mathrm{spin}}^x`$. In this direction, the staggered spin moment has a constant amplitude. The modulation perpendicular to the stripe corresponds to $`v_{\mathrm{spin}}^y`$. In this direction, the spin moment is modulated and suppressed when it crosses the stripe region. In Fig. 3, $`v_{\mathrm{spin}}^x>v_{\mathrm{spin}}^y`$. It indicates that the spin modulation is easier for the direction perpendicular to the stripe. In other words, the effective exchange integral $`J`$ across the stripe becomes weaker than that parallel to the stripe. It is the first time to microscopically derive the anisotropy of the spin wave velocity. As $`U`$ increases, $`v_{\mathrm{spin}}`$ decreases. These results are reasonable in view of the correspondence between Hubbard and Heisenberg models: $`v_{\mathrm{spin}}Jt^2/U`$. We have done the same calculation for the diagonal stripe to confirm that the spin wave velocity is similar. The longitudinal spin mode $`\chi _{zz}^{\prime \prime }(𝐪,\omega )`$ is shown in Fig. 4 along paths A and B of Fig. 1. Low energy modes appears at $`(2n+1)𝐐`$. Along path A, the dispersion relation is continuous and repeated, touching at $`(2n+1)𝐐`$. Away from $`𝐐`$, the intensity decreases as $`\chi _{zz}^{\prime \prime }(l𝐐,0.01t)/\chi _{zz}^{\prime \prime }(𝐐,0.01t)=5.9\times 10^2`$ ($`l=3`$), $`6.9\times 10^4`$ ($`l=5`$), $`5.3\times 10^4`$ ($`l=7`$). These ratios are larger than those given by Eq. (11), which is, thus, not satisfied for $`\chi _{zz}`$. The intensity $`\chi _{zz}^{\prime \prime }(𝐪,\omega )\frac{1}{3}\chi _{xx}^{\prime \prime }(𝐪,\omega )`$ near $`𝐐`$ along each dispersion relation in our results. There, $`\chi _{zz}^{\prime \prime }(𝐪,\omega )1/\omega `$ with increasing $`\omega `$. The charge mode $`\chi _{nn}^{\prime \prime }(𝐪,\omega )`$ is shown in Fig. 5 along paths C and D of Fig. 1. The low energy mode appears at $`2n𝐐`$. It has the identical dispersion curve as that of $`\chi _{zz}`$, because $`S_z`$ and $`n`$ couple each other in the RPA equation. As in Fig. 4 (a), it is a continuous curve along path C and its intensity deceases away from $`2𝐐`$. The $`2𝐐`$ mode comes from even harmonics, as $`n_{2𝐐,\sigma }`$ describes the charge modulation in the ISDW. We analyze the oscillation pattern of $`S_z`$ and $`n`$ by Eq. (10). Along the dispersion curve of $`\chi _{zz}`$ of Fig. 4 ($`\chi _{nn}`$ of Fig. 5), the response is enhanced by the resonance with the external field coupled to $`S_z`$ ($`n`$). The response of $`\delta S_z(𝐫,t)`$ and $`\delta n(𝐫,t)`$ shows large amplitude near the stripe region. It means that this excitation is related to the motion of the stripe, i.e., phason mode. This analysis shows that the collective mode at $`\omega 0`$ is the translational motion, where whole stripes move together to the same direction. When the pinning (such as the energy difference between the site-centered stripe and the bond-centered stripe) is negligible, this translational mode is a Goldstone mode and gapless as in Figs 4 and 5. From the analysis of Eq. (10), we understand that the excitation along the $`q_y`$-direction is a compress mode and that along the $`q_x`$-direction is a meandering mode. In the compress mode, the inter-stripe distance is modulated periodically in the direction perpendicular to the stripe, with keeping the straight line shape. In the meandering mode, each stripe meanders along the stripe direction with keeping the same inter-stripe distance. From Fig. 3, we see $`v_{\mathrm{phason}}^x<v_{\mathrm{phason}}^y`$. It means the meandering motion is easier to occur compared with the compress mode in the vertical stripe case of this model. As $`N`$ increases, $`v_{\mathrm{phason}}^x`$ and $`v_{\mathrm{phason}}^y`$ decrease. It is because effective interaction between neighbor stripes is weak and each stripe can move more freely when the inter-stripe distance becomes long. This slow velocity of the longitudinal mode for large $`N`$ may be related to the un-identified Fincher-Burke mode observed in Cr. This identification deserves further experimental and theoretical studies. We are also interested in the silent position $`𝐪=𝐒=2\pi (\frac{1}{2}\pm \delta ,\frac{1}{2})`$, which is an equivalent position to $`𝐐`$ in the paramagnetic state above Néel temperature (see Fig.1). This silent mode is related to a critical scattering in Cr . There exists a large intensity in $`\chi _{xx}^{\prime \prime }(𝐒,\omega )`$, whose dispersion has a gap of the order $`E_g`$. In $`\chi _{zz}^{\prime \prime }(𝐪,\omega )`$, the excitation at $`𝐒`$ (but slightly shifted to lower $`q_x`$) is shifted to lower energy as $`N=81216`$. It almost touches $`\omega =0`$ for $`N=16`$. It may suggest that we are approaching the transition to another low energy ground state (such as diagonal stripe) for lower $`n_h`$ . So far we mention only the insulating stripe state. The corresponding metallic stripe state can be also stabilized by merely introducing the next nearest hopping $`t^{}`$. The lowest energy state is given by $`\delta n_\mathrm{h}`$, and the Fermi level situates in the so-called midgap band. Then, we set $`N=8`$ in the metallic case. Since $`v_{\mathrm{spin}}^y`$ increases and $`v_{\mathrm{spin}}^x`$ decreases, we find $`v_{\mathrm{spin}}^y>v_{\mathrm{spin}}^x`$ in $`\chi _{xx}^{\prime \prime }`$. There is low energy excitation also at $`𝐒`$. In $`\chi _{zz}^{\prime \prime }`$ and $`\chi _{nn}^{\prime \prime }`$, the excitation at $`𝐐`$ or $`2𝐐`$ has a gap. The low energy excitation appears along the line at $`q_x=\pi /2`$ which is presented by a line in Fig. 1. It is the 1D CDW or SDW fluctuation mode within the stripe line, and originated from the Fermi surface nesting $`2𝐤_{\mathrm{F1D}}`$ of the parallel 1D Fermi lines (Fig. 12(d) in Ref. ) of the stripe state . Since the 1D Fermi state has a gap near $`(\frac{\pi }{2},\frac{\pi }{2})`$, the intensity of $`\chi ^{\prime \prime }(2𝐤_{\mathrm{F1D}},\omega 0)`$ vanishes near $`(\frac{\pi }{2},\pi )`$. These low energy excitations are diffusive since $`E_g=0`$ in metallic state. In summary, we have investigated the dynamical susceptibilities of transverse and longitudinal spin channels and charge one for the whole space spanned by 2D wave vector $`𝐪`$ and the energy $`\omega `$, and identified several elementary excitations; the spin wave mode and the phason mode related to the motion of the stripe line. This allows us to construct the whole landscape of $`(𝐪,\omega )`$ space for the excitation spectra of various channels: The identical dispersion relation is replicated at every $`(2n+1)𝐐`$, which has the anisotropic excitation cones along $`q_x`$\- and $`q_y`$-directions. Our predictions about $`\chi (𝐪,\omega )`$ are directly testable by careful inelastic neutron experiments on Cr alloys and underdoped cuprates. We thank G. Shirane, Y. Endoh, K. Yamada, T. Fukuda, M. Wakimoto, M. Matsuda and M. Fujita for their useful discussions and information.
warning/0005/hep-th0005188.html	ar5iv	text	# 1 Introduction ## 1 Introduction The possibility of writing down a matter-only theory which contains all the interactions of nature was first considered by Heisenberg in his study of a unified field theory. Along a somewhat different line of thought, motivated by the model of dynamical symmetry breaking by Nambu and Jona-Lasinio, Bjorken proposed a fermionic theory without an elementary gauge field which might explain the masslessness of the photon as a consequence of a symmetry breaking. He actually found his model equivalent to the standard QED, and the equivalence was further discussed by Bialynicki-Birula, Lurié and Macfarlane, and Guralnik among others. Similarly, the equivalence of the Nambu-Jona-Lasinio model with a manifestly renormalizable Yukawa theory was shown by Eguchi. A clear explanation of such equivalence between a manifestly renormalizable theory and an apparently non-renormalizable theory was only relatively recently given in the work of Hasenfratz et al. from the viewpoint of the Wilsonian renormalization group. It is this viewpoint which was extended to the equivalence between QED and the Bjorken model in the previous paper by the author. The purpose of the present paper is two-fold. First we improve the explanation of the equivalence between a manifestly renormalizable model and an apparently non-renormalizable model. Using the $`\frac{1}{N}`$ expansions to leading order, we give a simple yet well defined procedure for building a non-renormalizable model which is equivalent to the original renormalizable model. Though the equivalence is guaranteed only to leading order in $`\frac{1}{N}`$, the possibility of extending the equivalence beyond the leading order is supported by the general renormalization group argument. Second, as a concrete example, we show how to use the above procedure to construct a matter-only lagrangian equivalent to a chiral QED which has a much richer structure than the models of QED discussed in the previous paper. The paper is organized as follows. In sect. 2 we give a simpler explanation of the equivalence of the Bjorken model with QED than was given in the previous paper. Using this example, we construct a simple procedure for constructing a matter-only theory equivalent to the original model to leading order in $`\frac{1}{N}`$. In sect. 3 we summarize the relevant properties of a chiral QED with two flavors of chiral fermions and two complex scalar fields. In sect. 4 we follow the procedure given in sect. 2 to construct a matter-only theory equivalent to the chiral QED of sect. 3. We discuss the realization of chiral anomaly in the matter-only model in sect. 5 before we conclude the paper in sect. 6. Three appendices are given for completeness. We work in the four dimensional euclidean space throughout and use the same convention for the spinors as in ref. . ## 2 Equivalence revisited The equivalence of QED with a matter-only theory was discussed from the Wilsonian RG (renormalization group) viewpoint in the previous paper . In this section we wish to give an improved explanation of the equivalence. The standard QED, which is manifestly renormalizable by perturbation theory, is defined by the lagrangian $`_{QED}`$ $`=`$ $`{\displaystyle \frac{1}{4e_0^2}}F_{\mu \nu }^2+{\displaystyle \frac{1}{2e_0^2\xi _0}}(_\mu A_\mu )^2+{\displaystyle \frac{m_0^2}{2e_0^2}}A_\mu ^2+{\displaystyle \frac{\lambda }{N(4\pi )^2}}{\displaystyle \frac{(A_\mu ^2)^2}{8}}`$ (1) $`+\overline{\psi }^I\left({\displaystyle \frac{1}{i}}\text{/}+{\displaystyle \frac{1}{\sqrt{N}}}A\text{/}+iM\right)\psi ^I`$ where $`I=1,\mathrm{},N`$. We regularize the theory with a momentum cutoff $`\mathrm{\Lambda }`$. To leading order in $`\frac{1}{N}`$ the theory is renormalized by $`{\displaystyle \frac{(4\pi )^2}{e^2}}`$ $``$ $`{\displaystyle \frac{4}{3}}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }_0^2}{\mu ^2}}={\displaystyle \frac{(4\pi )^2}{e_0^2}}+{\displaystyle \frac{4}{3}}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{\mu ^2}}1`$ (2) $`{\displaystyle \frac{m_\gamma ^2}{e^2}}`$ $`=`$ $`{\displaystyle \frac{m_0^2}{e_0^2}}{\displaystyle \frac{2}{(4\pi )^2}}(\mathrm{\Lambda }^2M^2)`$ (3) $`{\displaystyle \frac{(4\pi )^2}{e^2\xi }}`$ $`=`$ $`{\displaystyle \frac{(4\pi )^2}{e_0^2\xi _0}}+{\displaystyle \frac{1}{3}}`$ (4) where $`m_\gamma `$ is a finite photon mass, and $`\mu `$ is an arbitrary low energy renormalization scale. To satisfy the Ward identity we must choose $$\lambda =\frac{4}{3}$$ (5) to leading order in $`\frac{1}{N}`$ so that the four-point proper vertex of the photon field vanishes at zero external momenta. The momentum cutoff does not allow the shift of loop momenta, and it does not respect the gauge invariance of the lagrangian. Hence, we need to introduce the self-coupling $`\lambda `$ to enforce the Ward identity. Now, let us try to construct a matter-only model which is equivalent to the above, to leading order in $`\frac{1}{N}`$. There is no unique choice for the lagrangian, but one straightforward choice is the following: $$_{Bj}=\overline{\psi }^I\left(\frac{1}{i}\text{/}+iM\right)\psi ^I\frac{1}{2Nv^2}\left(\overline{\psi }^I\gamma _\mu \psi ^I\right)^2+\mathrm{\Delta }$$ (6) where $$v^2\frac{m_0^2}{e_0^2}=\frac{2}{(4\pi )^2}(\mathrm{\Lambda }^2M^2)+\frac{m_\gamma ^2}{e^2}$$ (7) Here, the counterterms $`\mathrm{\Delta }`$ are given by $$\mathrm{\Delta }=\frac{1}{4e_0^2}(_\mu B_\nu _\mu B_\mu )^2+\frac{1}{2e_0^2\xi _0}(_\mu B_\mu )^2+\frac{\lambda }{N(4\pi )^2}\frac{(B_\mu ^2)^2}{8}$$ (8) where $$B_\mu \frac{(4\pi )^2}{2\mathrm{\Lambda }^2}\frac{1}{\sqrt{N}}\overline{\psi }^I\gamma _\mu \psi ^I$$ (9) We can understand the equivalence between $`_{QED}`$ and $`_{Bj}`$ from the viewpoint of the Wilsonian RG. We notice that the lagrangian (6) has two distinct critical points. One is trivial: $`M=0`$ and $`v`$ arbitrary (except for $`v_c`$ given below). The other parameters $`e_0^2,\xi _0,\lambda `$ are also arbitrary. This critical point, which describes $`N`$ free massless fermions, has codimension $`1`$ in theory space, and the fermion mass is the sole relevant parameter. There is another non-trivial critical point, however, which is given by $`M=0`$ and $`v^2=v_c^2`$ (where $`v_c^2=\frac{2\mathrm{\Lambda }^2}{(4\pi )^2}`$ to leading order in $`\frac{1}{N}`$). The remaining parameters are arbitrary. This critical point corresponds to $`N`$ free massless fermions and one free massless vector. The criticality has codimension $`2`$ in theory space, and the two relevant parameters are the fermion and photon masses. In finding the non-trivial critical point of the matter-only theory, the $`\frac{1}{N}`$ expansions are helpful. As long as we use perturbation expansions in powers of the four-fermi interaction, we cannot detect the second critical point. The Wilsonian RG tells us that arbitrary theories which are almost critical are completely characterized by the relevant and marginally irrelevant parameters. The corrections are suppressed by the negative powers of the cutoff. The critical point of $`_{QED}`$ at $`M=0`$, $`\frac{m_0^2}{e_0^2}=v_c^2`$ and that of $`_{Bj}`$ at $`M=0`$, $`v^2=v_c^2`$ describe the same criticality; both run to the same fixed point under the RG. Therefore, the two lagrangians must define the same theory. The form of the matter-only lagrangian (6) is by no means unique. Any lagrangian in the neighborhood of the non-trivial critical point will do as long as it has enough degrees of freedom to allow for two marginally irrelevant parameters $`e_0^2,\lambda `$. In ref. we have given a different but equivalent lagrangian. The procedure to get $`_{Bj}`$ out of $`_{QED}`$ can be made systematic. We first extract the gaussian part of $`_{QED}`$ quadratic in the gauge field with no derivatives: $$_{gauss}=\frac{v^2}{2}A_\mu ^2+A_\mu \frac{1}{\sqrt{N}}\overline{\psi }^I\gamma _\mu \psi ^I$$ (10) where $`v^2\frac{m_0^2}{e_0^2}`$. The equation of motion for this lagrangian gives $$A_\mu =B_\mu ^{}\frac{1}{\sqrt{N}v^2}\overline{\psi }^I\gamma _\mu \psi ^I$$ (11) Now we can construct a matter-only lagrangian by substituting the above interpolating field $`B_\mu ^{}`$ into $`A_\mu `$ in the remaining part of the lagrangian $`_{QED}`$: $`_{matter}=\overline{\psi }^I\left({\displaystyle \frac{1}{i}}\text{/}+iM\right)\psi ^I+{\displaystyle \frac{v^2}{2}}\left(A_\mu B_\mu ^{}\right)^2{\displaystyle \frac{v^2}{2}}B_{\mu }^{}{}_{}{}^{2}`$ $`+{\displaystyle \frac{1}{4e_0^2}}(_\mu B_\nu ^{}_\mu B_\mu ^{})^2+{\displaystyle \frac{1}{2e_0^2\xi _0}}(_\mu B_\mu ^{})^2+{\displaystyle \frac{\lambda }{N(4\pi )^2}}{\displaystyle \frac{(B_{\mu }^{}{}_{}{}^{2})^2}{8}}`$ (12) Note that the above lagrangian $`_{matter}`$ is quadratic in $`A_\mu `$ which plays the role of an auxiliary field. The equation of motion for $`A_\mu `$ is given by Eq. (11), and to leading order in $`\frac{1}{N}`$ we can substitute the equation of motion inside the correlation functions. This is because the proper two-point function of $`A_\mu `$ and $`B_\nu ^{}`$ is given by the diagram in Fig. 1 to leading order in $`\frac{1}{N}`$, and the proper part of the correlation function of $`\frac{1}{n!}\left(B_{\mu }^{}{}_{}{}^{2}/2\right)^n`$ with $`2n`$ $`A_\mu `$ fields is given by $`1`$ (except for the tensorial factor) to leading order in $`\frac{1}{N}`$. (Fig. 2) The errors of this substitution are suppressed by negative powers of the cutoff. Integrating over the auxiliary field $`A_\mu `$, we obtain the lagrangian $`_{Bj}`$ (6). Except for the mass term $`B_\mu ^2`$, we can replace $`B_\mu ^{}`$ by $`B_\mu `$ since the two differ only by a normalization which is unity up to $`\frac{\mu ^2}{\mathrm{\Lambda }^2}`$. Thus, to leading order in $`\frac{1}{N}`$, the two lagrangians $`_{QED}`$ (1) and $`_{Bj}`$ (6) are equivalent. While the equivalence between (1) and (6) is strictly valid only to leading order in $`\frac{1}{N}`$, the equivalence can be generalized beyond the leading order by choosing appropriate values for the parameters of (6), as is assured by the general Wilsonian RG argument. Finally we note that the lagrangian of the matter-only theory simplifies somewhat if we raise the momentum cutoff $`\mathrm{\Lambda }`$ to the Landau scale $`\mathrm{\Lambda }_0`$ defined by (2). Then the bare charge $`e_0`$ diverges, and the lagrangian (6) becomes $$_{Bj}^{}=\overline{\psi }^I\left(\frac{1}{i}\text{/}+iM\right)\psi ^I\frac{1}{2Nv^2}\left(\overline{\psi }^I\gamma _\mu \psi ^I\right)^2+\frac{\lambda }{N(4\pi )^2}\frac{(B_\mu ^2)^2}{8}$$ (13) This choice of the cutoff is known as the “compositeness condition” in the literature. We hope that the reader is convinced that any gauge theory (at least as long as it is abelian) can be rewritten as a matter-only theory. For completeness we also derive the equivalence between the manifestly renormalizable Yukawa model with the Nambu-Jona-Lasinio model along the above line in Appendix A. ## 3 Chiral QED with a momentum cutoff We consider a QED with chiral fermions in the Higgs phase.<sup>1</sup><sup>1</sup>1We do not know of any appropriate reference or textbook to point to. Contrary to the examples considered in ref. which had an explicit photon mass, here the photon mass results dynamically from the Higgs mechanism. For our purposes it is important to use a momentum cutoff regularization. As we have seen in the previous section, composite interpolating fields such as $`B_\mu `$ (9) play an important role in constructing matter-only theories, and their definitions need a cutoff explicitly. With a momentum cutoff $`\mathrm{\Lambda }`$, the theory is defined by the following lagrangian $`_{chiral}={\displaystyle \frac{1}{4e_0^2}}F_{\mu \nu }^2+{\displaystyle \frac{1}{2e_0^2\xi _0}}(_\mu A_\mu )^2+{\displaystyle \frac{m_0^2}{2e_0^2}}A_\mu ^2`$ $`+\overline{u}^I\left({\displaystyle \frac{1}{i}}\text{/}+{\displaystyle \frac{1}{\sqrt{N}}}A\text{/}{\displaystyle \frac{1+\gamma _5}{2}}\right)u^I+\overline{d}^I\left({\displaystyle \frac{1}{i}}\text{/}{\displaystyle \frac{1}{\sqrt{N}}}A\text{/}{\displaystyle \frac{1+\gamma _5}{2}}\right)d^I`$ $`+{\displaystyle \frac{1}{g_{u,0}^2}}\left(_\mu {\displaystyle \frac{i}{\sqrt{N}}}A_\mu \right)\varphi _u^{}\left(_\mu +{\displaystyle \frac{i}{\sqrt{N}}}A_\mu \right)\varphi _u`$ $`+{\displaystyle \frac{1}{g_{d,0}^2}}\left(_\mu +{\displaystyle \frac{i}{\sqrt{N}}}A_\mu \right)\varphi _d^{}\left(_\mu {\displaystyle \frac{i}{\sqrt{N}}}A_\mu \right)\varphi _d`$ $`+{\displaystyle \frac{i}{\sqrt{N}}}\left(\varphi _u\overline{u}_R^Iu_L^I+\varphi _u^{}\overline{u}_L^Iu_R^I\right)+{\displaystyle \frac{i}{\sqrt{N}}}\left(\varphi _d\overline{d}_R^Id_L^I+\varphi _d^{}\overline{d}_L^Id_R^I\right)`$ (14) $`+{\displaystyle \frac{1}{N}}{\displaystyle \frac{\lambda _{u,0}}{4g_{u,0}^4}}\left(\varphi _u^{}\varphi _u\right)^2+{\displaystyle \frac{M_{u,0}^2}{g_{u,0}^2}}\varphi _u^{}\varphi _u+{\displaystyle \frac{1}{N}}{\displaystyle \frac{\lambda _{d,0}}{4g_{d,0}^4}}\left(\varphi _d^{}\varphi _d\right)^2+{\displaystyle \frac{M_{d,0}^2}{g_{d,0}^2}}\varphi _d^{}\varphi _d`$ $`+{\displaystyle \frac{1}{N}}{\displaystyle \frac{\stackrel{~}{\lambda }}{g_{u,0}^2g_{d,0}^2}}\left(\varphi _u^{}\varphi _u\right)\left(\varphi _d^{}\varphi _d\right)+\mathrm{\Delta }`$ where $`\mathrm{\Delta }`$ denotes the counterterms to be given below to compensate for the gauge non-invariance of the momentum cutoff $`\mathrm{\Lambda }`$. The lagrangian is invariant under the following global chiral transformations: $`u_R^I`$ $``$ $`\text{e}^{i\theta _u}u_R^I,\varphi _u\text{e}^{i\theta _u}\varphi _u`$ (15) $`d_R^I`$ $``$ $`\text{e}^{i\theta _d}d_R^I,\varphi _d\text{e}^{i\theta _d}\varphi _d`$ (16) The off-diagonal part $`\theta _u=\theta _d`$ can be gauged thanks to the anomaly cancellation between the two flavors $`u,d`$. We note that the above global invariance properties are preserved by the momentum cutoff regularization. Though a non-vanishing $`\stackrel{~}{\lambda }`$ is allowed by the above chiral symmetry, we take $`\stackrel{~}{\lambda }=0`$ and ignore the mixing between $`\varphi _u`$ and $`\varphi _d`$ in the rest of the paper. This is solely to simplify the calculations. We now consider the Higgs phase in which both $`u`$ and $`d`$ acquire a mass. Denoting the fermion mass by $`m_i`$ for $`i=u,d`$, we obtain the following relation to leading order in $`\frac{1}{N}`$: $$\frac{1}{2}\lambda _{i,0}\frac{m_i^2}{g_{i,0}^2}+M_{u,0}^2=\frac{2}{(4\pi )^2}g_{i,0}^2\left(\mathrm{\Lambda }^2m_i^2\mathrm{ln}\frac{\mathrm{\Lambda }^2}{m_i^2}\right)$$ (17) Thus, for the theory to be in the Higgs phase we must choose $$M_{i,0}^2<\frac{2}{(4\pi )^2}g_{i,0}^2\mathrm{\Lambda }^2$$ (18) With the following shifts of fields the fields $`\rho _i,\phi _i`$ have vanishing expectation values: $$\varphi _i=\sqrt{N}m_i+\frac{g_{i,0}}{\sqrt{2}}\left(\rho _i+i\phi _i\right)$$ (19) To compensate for the non-gauge invariance of the momentum-cutoff regularization we introduce the following counterterms $`\mathrm{\Delta }`$: $$\mathrm{\Delta }=C\frac{1}{8(4\pi )^2N}\left(A_\mu ^2\right)^2+\frac{1}{(4\pi )^2N}A_\mu ^2\left(a_u\varphi _u^{}\varphi _u+a_d\varphi _d^{}\varphi _d\right)$$ (20) To satisfy the Ward identities (which are summarized in Appendix C), we must make the following choice to leading order in $`\frac{1}{N}`$: $`{\displaystyle \frac{m_0^2}{e_0^2}}`$ $`=`$ $`{\displaystyle \frac{2\mathrm{\Lambda }^2}{(4\pi )^2}}`$ (21) $`C`$ $`=`$ $`{\displaystyle \frac{4}{3}},a_u=a_d={\displaystyle \frac{1}{2}}`$ (22) Unlike the case of QED which has the photon mass as a free parameter, the mass $`m_0^2`$ is uniquely determined by the Ward identity. To leading order in $`\frac{1}{N}`$ the renormalization is done as follows: $`e`$ $`=`$ $`\sqrt{Z_3}e_0`$ (23) $`g_i`$ $`=`$ $`\sqrt{Z_i}g_{i,0}(i=u,d)`$ (24) $`{\displaystyle \frac{\lambda _i}{4g_i^4}}`$ $`=`$ $`{\displaystyle \frac{\lambda _{i,0}}{4g_{i,0}^4}}+{\displaystyle \frac{1}{(4\pi )^2}}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{\mu ^2}}`$ (25) where the renormalization constants are defined by $`Z_3`$ $``$ $`1{\displaystyle \frac{4}{3}}{\displaystyle \frac{e^2}{(4\pi )^2}}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{\mu ^2}}={\displaystyle \frac{1}{1+\frac{4}{3}\frac{e_0^2}{(4\pi )^2}\mathrm{ln}\frac{\mathrm{\Lambda }^2}{\mu ^2}}}`$ (26) $`Z_i`$ $``$ $`1{\displaystyle \frac{g_i^2}{(4\pi )^2}}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{\mu ^2}}={\displaystyle \frac{1}{1+\frac{g_{i,0}^2}{(4\pi )^2}\mathrm{ln}\frac{\mathrm{\Lambda }^2}{\mu ^2}}}`$ (27) The fermion masses $`m_i(i=u,d)`$ are unrenormalized to leading order in $`\frac{1}{N}`$. The products $`e_0A_\mu `$, $`g_{i,0}\rho _i`$, $`g_{i,0}\phi _i`$ are also left unrenormalized. We summarize the results for the renormalized correlation functions in Appendix B. ## 4 Matter-only model The lagrangian of the matter-only model can be obtained following the procedure given in sect. 2. The gaussian part of the lagrangian is given by $`_{gauss}={\displaystyle \frac{v^2}{2}}\left(A_\mu ^22A_\mu B_\mu ^{}\right)`$ $`+V_u^2\left(\varphi _u^{}\varphi _u\mathrm{\Phi }_{}^{}{}_{u}{}^{}\varphi _u\varphi _u^{}\mathrm{\Phi }_u^{}\right)+V_d^2\left(\varphi _d^{}\varphi _d\mathrm{\Phi }_{}^{}{}_{d}{}^{}\varphi _d\varphi _d^{}\mathrm{\Phi }_d^{}\right)`$ (28) where $`v^2`$ $``$ $`{\displaystyle \frac{m_0^2}{e_0^2}}={\displaystyle \frac{2}{(4\pi )^2}}\mathrm{\Lambda }^2`$ (29) $`V_u^2`$ $``$ $`{\displaystyle \frac{M_{u,0}^2}{g_{u,0}^2}}={\displaystyle \frac{2}{(4\pi )^2}}\mathrm{\Lambda }^22m_u^2\left({\displaystyle \frac{\lambda _{u,0}}{4g_{u,0}^4}}+{\displaystyle \frac{1}{(4\pi )^2}}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{m_u^2}}\right)`$ (30) $`V_d^2`$ $``$ $`{\displaystyle \frac{M_{d,0}^2}{g_{d,0}^2}}={\displaystyle \frac{2}{(4\pi )^2}}\mathrm{\Lambda }^22m_d^2\left({\displaystyle \frac{\lambda _{d,0}}{4g_{d,0}^4}}+{\displaystyle \frac{1}{(4\pi )^2}}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{m_d^2}}\right)`$ (31) and $`B_\mu ^{}`$ $``$ $`{\displaystyle \frac{1}{v^2}}{\displaystyle \frac{1}{\sqrt{N}}}\left(\overline{u}_R^I\gamma _\mu u_R^I\overline{d}_R^I\gamma _\mu d_R^I\right)`$ $`\mathrm{\Phi }_u^{}`$ $``$ $`{\displaystyle \frac{1}{V_u^2}}{\displaystyle \frac{1}{\sqrt{N}}}i\overline{u}_L^Iu_R^I\mathrm{\Phi }_d^{}{\displaystyle \frac{1}{V_d^2}}{\displaystyle \frac{1}{\sqrt{N}}}i\overline{d}_L^Id_R^I`$ (32) Hence, we obtain the following matter-only lagrangian: $`_{matter}=\overline{u}^I{\displaystyle \frac{1}{i}}\text{/}u^I+\overline{d}^I{\displaystyle \frac{1}{i}}\text{/}d^I{\displaystyle \frac{1}{2Nv^2}}\left(\overline{u}_R^I\gamma _\mu u_R^I\overline{d}_R^I\gamma _\mu d_R\right)^2`$ $`+{\displaystyle \frac{1}{NV_u^2}}\left(\overline{u}_L^Iu_R^I\right)\left(\overline{u}_R^Ju_L^J\right)+{\displaystyle \frac{1}{NV_d^2}}\left(\overline{d}_L^Id_R^I\right)\left(\overline{d}_R^Jd_L^J\right)+\mathrm{\Delta }`$ (33) The counterterms are given by $`\mathrm{\Delta }={\displaystyle \frac{1}{e_0^2}}{\displaystyle \frac{1}{4}}\left(_\mu B_\nu _\nu B_\mu \right)^2+{\displaystyle \frac{1}{\xi _0e_0^2}}{\displaystyle \frac{1}{2}}(_\mu B_\mu )^2+{\displaystyle \frac{C}{8(4\pi )^2N}}\left(B_\mu ^2\right)^2`$ $`+{\displaystyle \frac{1}{g_{u,0}^2}}\left(_\mu i{\displaystyle \frac{1}{\sqrt{N}}}B_\mu \right)\mathrm{\Phi }_u^{}\left(_\mu +i{\displaystyle \frac{1}{\sqrt{N}}}B_\mu \right)\mathrm{\Phi }_u`$ $`+{\displaystyle \frac{1}{g_{d,0}^2}}\left(_\mu i{\displaystyle \frac{1}{\sqrt{N}}}B_\mu \right)\mathrm{\Phi }_d^{}\left(_\mu +i{\displaystyle \frac{1}{\sqrt{N}}}B_\mu \right)\mathrm{\Phi }_d`$ (34) $`+{\displaystyle \frac{B_\mu ^2}{N(4\pi )^2}}\left(a_u\mathrm{\Phi }_u^{}\mathrm{\Phi }_u+a_d\mathrm{\Phi }_d^{}\mathrm{\Phi }_d\right)+{\displaystyle \frac{1}{N}}{\displaystyle \frac{\lambda _{u,0}}{4g_{u,0}^4}}\left(\mathrm{\Phi }_u^{}\mathrm{\Phi }_u\right)^2+{\displaystyle \frac{1}{N}}{\displaystyle \frac{\lambda _{d,0}}{4g_{d,0}^4}}\left(\mathrm{\Phi }_d^{}\mathrm{\Phi }_d\right)^2`$ where $`B_\mu `$ $``$ $`{\displaystyle \frac{(4\pi )^2}{2\mathrm{\Lambda }^2}}{\displaystyle \frac{1}{\sqrt{N}}}\left(\overline{u}_R^I\gamma _\mu u_R^I\overline{d}_R^I\gamma _\mu d_R^I\right)`$ (35) $`\mathrm{\Phi }_u`$ $``$ $`{\displaystyle \frac{(4\pi )^2}{2\mathrm{\Lambda }^2}}{\displaystyle \frac{1}{\sqrt{N}}}i\overline{u}_L^Iu_R^I,\mathrm{\Phi }_d{\displaystyle \frac{(4\pi )^2}{2\mathrm{\Lambda }^2}}{\displaystyle \frac{1}{\sqrt{N}}}i\overline{d}_L^Id_R^I`$ (36) If we had considered a non-vanishing $`\stackrel{~}{\lambda }`$ in the previous section, we would have obtained a term proportional to $`\|\mathrm{\Phi }_u\|^2\|\mathrm{\Phi }_d\|^2`$ which is allowed by the global symmetry of the lagrangian. The counterterms $`\mathrm{\Delta }`$ are essential for the complete equivalence to the original chiral QED. Without $`\mathrm{\Delta }`$, the matter-only lagrangian would depend only on the cutoff $`\mathrm{\Lambda }`$ and the three parameters $`v^2,V_u^2,V_d^2`$, and in no way the theory would be equivalent to the chiral QED which has more number of marginal parameters. ## 5 Chiral anomaly The chiral QED has two massless scalar modes $`\phi _u,\phi _d`$, and one is used for the Higgs mechanism to give the photon a mass, but the other is left. The remaining massless mode couples to two photons, and the renormalized vertex is given by (using the results in Appendix B) $`\left({\displaystyle \frac{g_d}{m_d}}\phi _u+{\displaystyle \frac{g_u}{m_u}}\phi _d\right)((k+l))A_\alpha (k)A_\beta (l)`$ $`{\displaystyle \frac{1}{\sqrt{2N}}}{\displaystyle \frac{i}{(4\pi )^2}}{\displaystyle \frac{8}{3}}{\displaystyle \frac{g_ug_d}{m_um_d}}ϵ_{\alpha \beta \mu \nu }k_\mu l_\nu `$ (37) Because of the equivalence between the chiral QED and the matter-only theory of the previous section, the same result is obtained for the correlation of the corresponding interpolating fields. In other words chiral anomaly is correctly reproduced by the matter-only theory. ## 6 Conclusion The essence of the present and previous papers is that any field theory can be renormalized by fine tuning the relevant parameters. This is what the Wilsonian renormalization group tells us. Hence, given an arbitrary lagrangian, whether it is manifestly renormalizable or not, if we can find a critical point, we get a renormalizable theory by finely adjusting the relevant parameters so that the physical mass is much smaller than the cutoff. The resulting theory depends only on the relevant and marginal degrees of freedom. We have elaborated on this expectation using concrete examples of abelian gauge theories in this and previous papers. We have given a very simple procedure for constructing a matter-only lagrangian which is equivalent to the original manifestly renormalizable gauge theory. We have used the $`\frac{1}{N}`$ expansions to locate non-trivial critical points which would be missed if we used perturbation expansions in non-renormalizable interactions. There has been some hope that the theories without elementary gauge fields or scalar fields would be more tightly constrained. From the Wilsonian RG viewpoint, this hope is not well founded. Unless the symmetry of the theory is enhanced by imposing the vanishing of a relevant or marginally irrelevant parameter, there is no justification for the relationship among the free parameters of the theory. Writing a matter-only theory in a particular form has, at best, as much significance as imposing an arbitrary relation among the parameters of the theory. The author is aware that much of what is written in this paper may sound obvious to those versed in the Wilsonian RG. Since there is no reference to point to, however, the author hopes that this paper has its merit in describing a proper and modern way of looking at this old subject. The author thanks Prof. K. Akama for informing him of the references to early works on the subject. This work was supported in part by the Grant-In-Aid for Scientific Research (No. 11640279) from the Ministry of Education, Science, and Culture, Japan. ## Appendix A The Yukawa model vs. the NJL model The perturbatively renormalizable Yukawa model, which is invariant under a chiral $`U(1)`$, is defined by $`_Y`$ $`=`$ $`_\mu \varphi ^{}_\mu \varphi +{\displaystyle \frac{\lambda _0}{4N}}\left(\varphi ^{}\varphi \right)^2+m_0^2\varphi ^{}\varphi `$ (38) $`+\overline{\psi }^I{\displaystyle \frac{1}{i}}\text{/}\psi ^I+i{\displaystyle \frac{g_0}{\sqrt{N}}}\left(\varphi \overline{\psi }_R^I\psi _L^I+\varphi ^{}\overline{\psi }_L^I\psi _R^I\right)`$ with a momentum cutoff $`\mathrm{\Lambda }`$. By rescaling $`\varphi `$ by $`\frac{1}{g_0}\varphi `$, we rewrite the above as $`_Y`$ $`=`$ $`{\displaystyle \frac{1}{g_0^2}}_\mu \varphi ^{}_\mu \varphi +{\displaystyle \frac{\lambda _0}{4g_0^4N}}\left(\varphi ^{}\varphi \right)^2+{\displaystyle \frac{m_0^2}{g_0^2}}\varphi ^{}\varphi `$ (39) $`+\overline{\psi }^I{\displaystyle \frac{1}{i}}\text{/}\psi ^I+i{\displaystyle \frac{1}{\sqrt{N}}}\left(\varphi \overline{\psi }_R^I\psi _L^I+\varphi ^{}\overline{\psi }_L^I\psi _R^I\right)`$ To leading order in $`\frac{1}{N}`$, the renormalized parameters are given by $`{\displaystyle \frac{1}{g^2}}`$ $``$ $`{\displaystyle \frac{1}{(4\pi )^2}}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }_0^2}{\mu ^2}}={\displaystyle \frac{1}{g_0^2}}+{\displaystyle \frac{1}{(4\pi )^2}}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{\mu ^2}}`$ (40) $`{\displaystyle \frac{m^2}{g^2}}`$ $`=`$ $`{\displaystyle \frac{m_0^2}{g_0^2}}{\displaystyle \frac{2}{(4\pi )^2}}\mathrm{\Lambda }^2`$ (41) $`{\displaystyle \frac{\lambda }{4g^4}}`$ $`=`$ $`{\displaystyle \frac{\lambda _0}{4g_0^4}}+{\displaystyle \frac{1}{(4\pi )^2}}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{\mu ^2}}`$ (42) where $`\mu `$ is an arbitrary low energy renormalization scale. The theory is in the symmetric (broken) phase if $$\frac{m_0^2}{g_0^2}>(<)\frac{2}{(4\pi )^2}\mathrm{\Lambda }^2$$ (43) Following the procedure given in sect. 2, we obtain an equivalent matter-only lagrangian of the generalized Nambu-Jona-Lasinio model as $`_{NJL}`$ $`=`$ $`\overline{\psi }^I{\displaystyle \frac{1}{i}}\text{/}\psi ^I+{\displaystyle \frac{1}{Nv^2}}\left(\overline{\psi }_R^I\psi _L^I\right)\left(\overline{\psi }_L^J\psi _R^J\right)`$ (44) $`+{\displaystyle \frac{1}{g_0^2}}_\mu \mathrm{\Phi }^{}_\mu \mathrm{\Phi }+{\displaystyle \frac{1}{N}}{\displaystyle \frac{\lambda _0}{4g_0^4}}\left(\mathrm{\Phi }^{}\mathrm{\Phi }\right)^2,`$ where $`v^2\frac{m_0^2}{g_0^2}=\frac{2}{(4\pi )^2}\mathrm{\Lambda }^2+\frac{m^2}{g^2}`$, and $$\mathrm{\Phi }\frac{i}{\sqrt{N}}\frac{(4\pi )^2}{2\mathrm{\Lambda }^2}\overline{\psi }_L^I\psi _R^I,\mathrm{\Phi }^{}\frac{i}{\sqrt{N}}\frac{(4\pi )^2}{2\mathrm{\Lambda }^2}\overline{\psi }_R^I\psi _L^I$$ (45) If we raise the cutoff $`\mathrm{\Lambda }`$ to the Landau scale $`\mathrm{\Lambda }_0`$, we get $`\frac{1}{g_0^2}0`$, and $`\frac{\lambda _0}{4g_0^4}\frac{\lambda }{4g^4}\frac{1}{g^2}`$. Hence, we obtain a simpler lagrangian $$_{NJL}^{}=\overline{\psi }^I\frac{1}{i}\text{/}\psi ^I+\frac{1}{Nv^2}\left(\overline{\psi }_R^I\psi _L^I\right)\left(\overline{\psi }_L^J\psi _R^J\right)+\frac{1}{N}\left(\frac{\lambda }{4g^4}\frac{1}{g^2}\right)\left(\mathrm{\Phi }^{}\mathrm{\Phi }\right)^2$$ (46) For an original discussion, see ref. . ## Appendix B The renormalized two-, three-, and four-point functions in chiral QED For definiteness we will summarize the results of the lowest order calculations in $`\frac{1}{N}`$, ignoring the contributions suppressed by negative powers of the cutoff. We only list those necessary for verifying the Ward identities. ### B.1 two-point functions For small external momenta, we obtain the following approximate results for the proper two-point functions: $`\mathrm{\Pi }_{A_\mu A_\nu }(k)`$ $``$ $`\left[{\displaystyle \frac{1}{e^2}}+{\displaystyle \frac{1}{(4\pi )^2}}\left({\displaystyle \frac{2}{3}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_u^2}}+{\displaystyle \frac{2}{3}}\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_d^2}}{\displaystyle \frac{5}{3}}\right)\right]\left(k^2\delta _{\mu \nu }k_\mu k_\nu \right)`$ $`+{\displaystyle \frac{1}{e^2\xi }}k_\mu k_\nu +\delta _{\mu \nu }2[{\displaystyle \frac{m_u^2}{g_u^2}}+{\displaystyle \frac{m_d^2}{g_d^2}}`$ $`+{\displaystyle \frac{1}{(4\pi )^2}}\{m_u^2(\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_u^2}}1)+m_d^2(\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_d^2}}1)\}]`$ where $`\frac{1}{\xi e^2}=\frac{1}{\xi _0e_0^2}\frac{1}{(4\pi )^2}\frac{1}{3}`$. $`\mathrm{\Pi }_{\phi _u\phi _u}(k)`$ $``$ $`k^2\left[1+{\displaystyle \frac{g_u^2}{(4\pi )^2}}\left(\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_u^2}}1\right)\right]`$ (48) $`\mathrm{\Pi }_{\phi _uA_\nu }(k)`$ $``$ $`ik_\nu \sqrt{2}{\displaystyle \frac{m_u}{g_u}}\left[1+{\displaystyle \frac{g_u^2}{(4\pi )^2}}\left(\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_u^2}}1\right)\right]`$ (49) $`\mathrm{\Pi }_{\rho _u\rho _u}(p)`$ $``$ $`p^2\left[1+{\displaystyle \frac{g_u^2}{(4\pi )^2}}\left(\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_u^2}}{\displaystyle \frac{5}{3}}\right)\right]`$ (50) $`+\lambda _u{\displaystyle \frac{m_u^2}{g_u^2}}+4m_u^2{\displaystyle \frac{g_u^2}{(4\pi )^2}}\left(\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_u^2}}1\right)`$ ### B.2 three-point functions For small external momenta, we find $`\mathrm{\Pi }_{\rho _uA_\mu \phi _u}(p,k,(p+k)){\displaystyle \frac{i}{\sqrt{N}}}[(2p+k)_\mu `$ $`+{\displaystyle \frac{g_u^2}{(4\pi )^2}}\{k_\mu (\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_u^2}}3)+p_\mu (2\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_u^2}}4)\}]`$ (51) $`\mathrm{\Pi }_{\rho _u\phi _u\phi _u}(p,k,(p+k)){\displaystyle \frac{1}{\sqrt{2N}}}{\displaystyle \frac{g_u}{m_u}}[\lambda _u{\displaystyle \frac{m_u^2}{g_u^2}}`$ $`+{\displaystyle \frac{g_u^2}{(4\pi )^2}}\{4m_u^2(\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_u^2}}1)2(k^2+kp+{\displaystyle \frac{p^2}{3}})\}]`$ (52) $`\mathrm{\Pi }_{A_\mu A_\nu \rho _u}(k,l,(k+l))`$ $`{\displaystyle \frac{2\sqrt{2}}{\sqrt{N}}}{\displaystyle \frac{m_u}{g_u}}\delta _{\mu \nu }\left[1+{\displaystyle \frac{g_u^2}{(4\pi )^2}}\left(\mathrm{ln}{\displaystyle \frac{\mu ^2}{m_u^2}}2\right)\right]`$ (53) $`\mathrm{\Pi }_{A_\mu A_\nu \phi _u}(k,l,(k+l)){\displaystyle \frac{1}{\sqrt{2N}}}{\displaystyle \frac{i}{(4\pi )^2}}{\displaystyle \frac{4}{3}}{\displaystyle \frac{g_u}{m_u}}ϵ_{\alpha \beta \mu \nu }k_\mu l_\nu `$ (54) Note $`\mathrm{\Pi }_{AAA}`$ vanishes due to the CP invariance. ### B.3 four-point functions For small external momenta we find $`\mathrm{\Pi }_{A_\alpha A_\beta A_\gamma A_\delta }`$ $``$ $`{\displaystyle \frac{1}{N}}{\displaystyle \frac{1}{(4\pi )^2}}{\displaystyle \frac{4}{3}}(\delta _{\alpha \beta }\delta _{\gamma \delta }+\text{permutations})`$ (55) $`\mathrm{\Pi }_{\phi _uA_\alpha A_\beta A_\gamma }(k,\mathrm{})`$ $``$ $`{\displaystyle \frac{i}{N}}{\displaystyle \frac{g_u}{m_u}}{\displaystyle \frac{1}{(4\pi )^2}}{\displaystyle \frac{\sqrt{2}}{3}}\left(\delta _{\alpha \beta }k_\gamma +\delta _{\alpha \gamma }k_\beta +\delta _{\beta \gamma }k_\alpha \right)`$ (56) ## Appendix C Ward identities The following Ward identities must be satisfied. $`ik_\mu \mathrm{\Pi }_{A_\mu A_\nu }(k)+\sqrt{2}{\displaystyle \frac{m_u}{g_u}}\mathrm{\Pi }_{\phi _uA_\nu }(k)\sqrt{2}{\displaystyle \frac{m_d}{g_d}}\mathrm{\Pi }_{\phi _dA_\nu }(k)=0`$ (57) $`ik_\mu \mathrm{\Pi }_{A_\mu \phi _u}(k)+\sqrt{2}{\displaystyle \frac{m_u}{g_u}}\mathrm{\Pi }_{\phi _u\phi _u}(k)\sqrt{2}{\displaystyle \frac{m_d}{g_d}}\mathrm{\Pi }_{\phi _d\phi _d}(k)=0`$ (58) $`ik_\mu \mathrm{\Pi }_{A_\mu \rho _u\phi _u}(k,p,(k+p))+\sqrt{2}{\displaystyle \frac{m_u}{g_u}}\mathrm{\Pi }_{\phi _u\rho _u\phi _u}(k,p,(k+p))`$ $`={\displaystyle \frac{1}{\sqrt{N}}}\mathrm{\Pi }_{\phi _u\phi _u}(p+k)+{\displaystyle \frac{1}{\sqrt{N}}}\mathrm{\Pi }_{\rho _u\rho _u}(p)`$ (59) $`ik_\mu \mathrm{\Pi }_{A_\mu \rho _uA_\nu }(k,p,(p+k))+\sqrt{2}{\displaystyle \frac{m_u}{g_u}}\mathrm{\Pi }_{\phi _u\rho _uA_\nu }(k,p,(p+k))`$ $`={\displaystyle \frac{1}{\sqrt{N}}}\mathrm{\Pi }_{\varphi _uA_\nu }(k)`$ (60) $`i(k_\mu +l_\mu )\mathrm{\Pi }_{A_\mu A_\alpha A_\beta }((k+l),k,l)`$ (61) $`=\sqrt{2}{\displaystyle \frac{m_u}{g_u}}\mathrm{\Pi }_{\phi _uA_\alpha A_\beta }((k+l),k,l)+\sqrt{2}{\displaystyle \frac{m_d}{g_d}}\mathrm{\Pi }_{\phi _dA_\alpha A_\beta }((k+l),k,l)`$ $`ik_\alpha \mathrm{\Pi }_{A_\alpha A_\beta A_\gamma A_\delta }(k,\mathrm{})`$ $`=\sqrt{2}{\displaystyle \frac{m_u}{g_u}}\mathrm{\Pi }_{\phi _uA_\beta A_\gamma A_\delta }(k,\mathrm{})+\sqrt{2}{\displaystyle \frac{m_d}{g_d}}\mathrm{\Pi }_{\phi _dA_\beta A_\gamma A_\delta }(k,\mathrm{})`$ (62)
warning/0005/hep-ph0005083.html	ar5iv	text	# Nonperturbative Renormalization and the QCD Vacuum ## I Introduction There is a long history of using mean field techniques to explore dynamical chiral symmetry breaking in Quantum Chromodynamics. Chiral symmetry breaking is of fundamental importance because it is directly related to the vacuum structure of QCD. It also has important practical consequences, such as making chiral perturbation theory a viable approximation to low energy QCD. Mean field equations are typically derived by analyzing a truncation of the Schwinger-Dyson equations or, when dealing with a Hamiltonian, by performing a Bogoliubov-Valatin transformation on the particle basis. An added complexity with Hamiltonian approaches has to do with renormalization. Unlike covariant approaches to field theory which generate a finite number of counterterms, the Hamiltonian formalism necessarily involves a noncovariant truncation of the theory and hence may generate noncanonical counterterms. As a result it may not be possible to remove the ultraviolet cutoff. In this sense, an effective field theory approach to Hamiltonian-based renormalization becomes natural. The problem of renormalization in models derived from a QCD Hamiltonian is often ignored. An exception was the work of Adler and Davis who have noted that, when using the BCS Ansatz, consistency with the Ward identities imposes a definite counterterm structure on the gap equation. They also noted that calculations should be performed with the full Hamiltonian (as opposed to the normal ordered Hamiltonian). Although the quark sector of the vacuum has been studied relatively well, the gluonic sector has largely been ignored. However, the gluonic vacuum is also relevant to studies of chiral symmetry breaking because the quark and gluonic gap equations are coupled and, more importantly, the gluon propagator is related to the quark-antiquark confining interaction. Furthermore, the gluonic vacuum determines the properties of the gluonic quasiparticles which in turn may be used to construct a Fock space expansion of hadrons such as glueballs or hybrids. In this paper we derive and renormalize the gap equation for the purely gluonic sector of QCD. This work extends three previous papers. The first of these presented the unrenormalized gluonic gap equation and calculated the gluon condensate and glueball spectrum. The other two described how to renormalize Hamiltonian QCD using the similarity or flow equation methods for evolving the renormalization scale. Both methods are necessarily perturbative, which means that gap equations may be derived with one (or more) loop corrections. However, since the phenomena being studied are nonperturbative, it is preferable to construct a nonperturbative renormalization scheme from the start. Unfortunately, this means abandoning the elegant flow equation methodology which is so well suited to Hamiltonian-based perturbation theory. We further extend the previous work by showing how to obtain the long range interaction between color sources. This is done in the BCS Ansatz by simultaneously considering the single gluon spectral function and the non-Abelian Coulomb potential. This subject was previously considered in Ref. and (we discuss the various approaches below). We derive a self consistent gap equation for the gluon spectral function. Self consistency arises because the non-Abelian potential appearing in this equation depends on the gluon spectral function. The solution to the gap equation gives rise to an infrared enhancement in the effective non-Abelian potential which may be identified with linear confinement. In the following section we give a brief review of QCD in Coulomb gauge and present our regularization scheme. We then discuss the renormalization scheme adopted, the gap equation, and the effective interaction emerging from the self consistent solution. We conclude and discuss future applications in Section III. ## II QCD in Coulomb Gauge Understanding the properties of soft gluons is one of the major challenges of hadronic physics. It is natural to study soft gluons using the Hamiltonian formulation of QCD in a physical gauge such as the Coulomb gauge. This has many advantages: it is closest in spirit to quantum mechanical models of QCD, all degrees of freedom are physical, no additional constraints are need on the Fock space, the norm is positive definite, spurious retardation effects are minimized, and confinement may be rigorously identified with the non-Abelian Coulomb interaction in the heavy quark limit . Finally this is also the natural choice for the study of nonrelativistic bound states of a few constituent degrees of freedom. It is of relevance in the gluonic sector since one would expect the gluonic quasiparticles to be heavy due to strong confining interactions. The pure gauge QCD Hamiltonian in Coulomb gauge may be written as $`H=H_g+H_C`$, where the terms are given by, $`H_g`$ $`=`$ $`\mathrm{Tr}{\displaystyle d^3x\left[𝒥^1𝚷𝒥𝚷+𝐁𝐁\right]}`$ (1) $`H_C`$ $`=`$ $`{\displaystyle \frac{1}{2}}g^2{\displaystyle d^3xd^3y𝒥^1\rho ^a(𝐱)K^{ab}(𝐱,𝐲)𝒥\rho ^b(𝐲)}`$ (2) The second term is the instantaneous non-Abelian Coulomb interaction for QCD. The kernel $`K`$ is given by $$K^{ab}(𝐱,𝐲)=𝐱,a\|(𝐃)^1(^2)(𝐃)^1\|𝐲,b$$ (3) and the color charge density $`\rho ^a`$ is given by $$\rho ^a(𝐱)=f^{abc}𝐀^b(𝐱)𝚷^c(𝐱).$$ (4) When expanded in powers of the strong coupling, the Coulomb kernel contains infinitely many terms arising from the inverse of the adjoint covariant derivative: $$𝐃^{ab}=\delta ^{ab}gf^{abc}𝐀^c.$$ (5) At lowest order in the coupling the kernel is given by $`K^{ab}(𝐱,𝐲)=\delta ^{ab}/(4\pi \|𝐱𝐲\|)`$ as in QED. The Faddeev-Popov determinant appears in the kinetic energy and Coulomb interaction terms due to the curvature of the gauge manifold and is given by $`𝒥=\mathrm{det}\left[𝐃\right]`$. Lastly, the components of the non-Abelian magnetic field are given by $$B_i^a=ϵ_{ijk}\left(_jA_k^a+\frac{g}{2}f^{abc}A_j^bA_k^c\right).$$ (6) It should be stressed that this is a bare Hamiltonian. Regularization leads to counterterms which in practice will depend on the approximation schemes used to diagonalize the Hamiltonian. As stated above, our goal is to study the QCD vacuum with the aid of the BCS Ansatz. In the past we have chosen to regulate the Hamiltonian by restricting its matrix elements in the basis of eigenstates of the bare kinetic energy to a band-diagonal form . Although an elegant formulation of the renormalization procedure is possible with this choice of basis, it gives rise to lengthy expressions because all calculations must be performed in the particle basis. Furthermore, such a scheme does not seem to be particularly relevant when dealing with nonperturbative renormalization since there is no particular advantage in using the bare basis. Hence, it is preferable to employ the field basis and we introduce a field-based regulator which is analogous to Schwinger’s point splitting. This consists of smearing field operators over a small spatial region: $$\stackrel{~}{\mathrm{A}}_i^b(𝐱)=d^3y\frac{\mathrm{\Lambda }^3}{(2\pi )^{3/2}}\mathrm{A}_i^b(𝐲)\mathrm{e}^{(𝐱𝐲)^2\frac{\mathrm{\Lambda }^2}{2}}$$ (7) Here $`\mathrm{\Lambda }`$ is the UV cutoff and the fields are effectively smeared over a distance $`O(1/\mathrm{\Lambda })`$. The mode expansions are $`\stackrel{~}{\mathrm{A}}_i^b(𝐱)`$ $`=`$ $`{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{1}{\sqrt{2\omega _0(k)}}\left(a_i^b(𝐤)+a_i^b(𝐤)\right)\mathrm{e}^{i𝐤𝐱\frac{k^2}{2\mathrm{\Lambda }^2}}}`$ (8) $`\stackrel{~}{\mathrm{\Pi }}_i^b(𝐱)`$ $`=`$ $`{\displaystyle \frac{d^3k}{(2\pi )^3}i\sqrt{\frac{\omega _0(k)}{2}}\left(a_i^b(𝐤)a_i^b(𝐤)\right)\mathrm{e}^{i𝐤𝐱\frac{k^2}{2\mathrm{\Lambda }^2}}}`$ (9) where $`\omega _0(k)=\|𝐤\|`$ in the perturbative vacuum. We shall subsequently drop the tildes. Contractions of the field operators which are needed below are given as follows $`\mathrm{A}_i^a(𝐱)\mathrm{A}_j^b(𝐲)`$ $`=`$ $`\delta ^{ab}{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{D_{ij}(k)}{2\omega _0(k)}\mathrm{e}^{i𝐤(𝐱𝐲)\frac{k^2}{\mathrm{\Lambda }^2}}}`$ (10) $`\mathrm{\Pi }_i^a(𝐱)\mathrm{\Pi }_j^b(𝐲)`$ $`=`$ $`\delta ^{ab}{\displaystyle \frac{d^3k}{(2\pi )^3}D_{ij}(k)\frac{\omega _0(k)}{2}\mathrm{e}^{i𝐤(𝐱𝐲)\frac{k^2}{\mathrm{\Lambda }^2}}}`$ (11) $`\mathrm{A}_i^a(𝐱)\mathrm{\Pi }_j^b(𝐲)`$ $`=`$ $`i\delta ^{ab}{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{D_{ij}(k)}{2}\mathrm{e}^{i𝐤(𝐱𝐲)\frac{k^2}{\mathrm{\Lambda }^2}}}`$ (12) where $`D_{ij}(k)=\delta _{ij}\widehat{k}_i\widehat{k}_j`$ is the transverse delta function associated with Coulomb gauge. ### A Nonperturbative Renormalization In order to reduce the cutoff dependence induced by regularization we allow the couplings to be $`\mathrm{\Lambda }`$-dependent and add an infinite set of counterterms to the Hamiltonian. The counterterms are organized in powers of the cutoff, $$\delta H\underset{n=2}{}\frac{c_n(\mathrm{\Lambda })}{\mathrm{\Lambda }^n}𝒪^{(n)}.$$ (13) As long at the cutoff only affects the operator products at short relative distances, the counterterms $`𝒪`$ must be local. Thus they may be classified according to their canonical dimension, $`n+1`$. Furthermore they have to preserve the unbroken symmetries i.e. be rotationally invariant, color, spin, and flavor singlets. Notice that the series starts at order $`\mathrm{\Lambda }^2`$. This is because the lowest dimension operator satisfying the above restrictions is $`𝑑\mathrm{𝐱𝐀}^2(𝐱)`$ which is of dimension $`1`$, ($`n=2`$). The needed marginal and relevant operators contributing to the gluon sector of the Hamiltonian are $`\delta H`$ $`=`$ $`{\displaystyle \frac{M^2(\mathrm{\Lambda })}{2}}{\displaystyle d^3xA^2}+{\displaystyle \frac{c_0(\mathrm{\Lambda })}{2}}{\displaystyle d^3xA^2A}`$ (15) $`+{\displaystyle \frac{Z^1(\mathrm{\Lambda })1}{2}}{\displaystyle d^3x\mathrm{\Pi }^2}+{\displaystyle \frac{Z(\mathrm{\Lambda })1}{2}}{\displaystyle d^3x𝐁^2}`$ We have established contact with more traditional notation by defining $`M^2(\mathrm{\Lambda })=\mathrm{\Lambda }^2c_2`$ where $`c_2`$ is dimensionless. Operators of dimension three and five do not occur. Thus the next set of relevant operators are of dimension six: $`f_{abc}f_{cde}A^{ia}\mathrm{\Pi }^{ib}A^{jd}\mathrm{\Pi }^{je}`$, $`A\mathrm{\Pi }A\mathrm{\Pi }`$, $`A^4A`$, $`\mathrm{\Pi }^2\mathrm{\Pi }`$, $`f_{abc}\mathrm{\Pi }^{ia}A^b\mathrm{\Pi }^{ic}`$, combinations of six gluon fields and gradients, and so on. It is important to recall that in an effective field theory the cutoff is not removed. Rather the various coefficients are determined by requiring that observables are accurate to a given order in $`p/\mathrm{\Lambda }`$, where $`p`$ is some characteristic, measurable momentum scale. In our case the $`c_n(\mathrm{\Lambda })`$ should be tuned to reproduce experimental data calculated in some nonperturbative scheme. Here we are referring to a number of standard many-body calculational schemes such as the BCS vacuum Ansatz, the Tamm-Dancoff truncation, or the random phase approximation. Transverse gluons do not contribute to the counterterms because we employ a BCS Ansatz in this work. Thus at this stage of the calculation, the $`\mathrm{\Lambda }`$ dependence of the counterterms will not match that of perturbative QCD. Nevertheless couplings to the transverse gluons with momenta below the cutoff are still present in the Hamiltonian and lead to important effects which (if the valence sector dominates) could be taken into account in perturbation theory around quasiparticle bound states. Removal of the cutoff without taking into account the effects from transverse gluons would clearly be incorrect. This will be seen explicitly below when we compare, in the weak coupling limit, the counterterms calculated with the BCS Ansatz with the results coming from perturbation theory. ### B The Gap Equation There are many equivalent formulations of the BCS approach. The one appropriate to many-body physics is the Bogoliubov-Valatin (BV) canonical transformation on the particle operators. In our case this may be written as $`a_i^b(𝐤)`$ $`=`$ $`c(k)\alpha _i^b(𝐤)+s(k)\alpha _i^b(𝐤)`$ (16) $`a_i^b(𝐤)`$ $`=`$ $`c(k)\alpha _i^b(𝐤)+s(k)\alpha _i^b(𝐤)`$ (17) where the rotation is parameterized in terms of an unknown gap function, $`\omega `$, as $`c(k)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\sqrt{{\displaystyle \frac{\omega _0(k)}{\omega (k)}}}+\sqrt{{\displaystyle \frac{\omega (k)}{\omega _0(k)}}}\right)`$ (18) $`s(k)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\sqrt{{\displaystyle \frac{\omega _0(k)}{\omega (k)}}}\sqrt{{\displaystyle \frac{\omega (k)}{\omega _0(k)}}}\right).`$ (19) Because we have employed a field-based regulator, the effect of the BV transformation is to simply replace $`\omega _0(k)=k`$ with $`\omega _0=\omega (k)`$ in the mode expansions of Eq. (9). Thus all matrix elements evaluated in the BCS vacuum are simply related to those evaluated in the perturbative vacuum with the same replacement. The gap equation may now be obtained from the one-body portion of the QCD Hamiltonian. Normal ordering with respect to the BCS vacuum yields $$H_{1b}=\frac{d^3q}{(2\pi )^3}\mathrm{e}^{\frac{k^2}{\mathrm{\Lambda }^2}}[E(q)\alpha _i^b(𝐪)\alpha _i^b(𝐪)+G(q)(\alpha _i^b(𝐪)\alpha _i^b(𝐪)+H.c.)]$$ (20) where $`E(q)`$ $`=`$ $`{\displaystyle \frac{1}{2\omega (q)}}[{\displaystyle \frac{\omega ^2(q)}{Z}}+Zq^2+M^2(\mathrm{\Lambda })+q^2c_0(\mathrm{\Lambda })`$ (23) $`+{\displaystyle \frac{N_c}{4}}{\displaystyle \frac{d^3k}{(2\pi )^3}\mathrm{e}^{\frac{k^2}{\mathrm{\Lambda }^2}}\stackrel{~}{V}(𝐤+𝐪)(1+(\widehat{k}\widehat{q})^2)\frac{\omega ^2(k)+\omega ^2(q)}{\omega (k)}}`$ $`+\pi \alpha _s(\mathrm{\Lambda })N_c{\displaystyle }{\displaystyle \frac{d^3k}{(2\pi )^3}}\mathrm{e}^{\frac{k^2}{\mathrm{\Lambda }^2}}{\displaystyle \frac{(3(\widehat{k}\widehat{q})^2)}{\omega (k)}}]`$ and $`G(q)`$ $`=`$ $`{\displaystyle \frac{1}{4\omega (k)}}[{\displaystyle \frac{\omega ^2(q)}{Z}}+Zq^2+M^2(\mathrm{\Lambda })+c_0(\mathrm{\Lambda })q^2`$ (24) $`+`$ $`{\displaystyle \frac{N_c}{4}}{\displaystyle \frac{d^3k}{(2\pi )^3}\stackrel{~}{V}(𝐤+𝐪)\mathrm{e}^{\frac{k^2}{\mathrm{\Lambda }^2}}(1+(\widehat{k}\widehat{q})^2)\frac{\omega ^2(k)\omega ^2(q)}{\omega (k)}}`$ (25) $`+`$ $`\pi \alpha _s(\mathrm{\Lambda })N_c{\displaystyle }{\displaystyle \frac{d^3k}{(2\pi )^3}}\mathrm{e}^{\frac{k^2}{\mathrm{\Lambda }^2}}{\displaystyle \frac{(3(\widehat{k}\widehat{q})^2)}{\omega (k)}}]`$ (26) The spectral function $`\omega (k)`$ is obtained from solving the gap equation $$G(q)=0$$ (27) which demands that the quasiparticle, BCS vacuum decouples from states with pairs of gluons –a feature reminiscent of the constituent quark model. The static potential in these expressions is the expectation value of the non-Abelian Coulomb kernel in the BCS vacuum state, $$\stackrel{~}{V}(𝐪)=d^3r\mathrm{e}^{i𝐪𝐫}BCS\|K^{ab}(𝐫,\mathrm{𝟎})\|BCS.$$ (28) At lowest order in the coupling and in the perturbative vacuum this potential is simply $`4\pi \alpha _s(\mathrm{\Lambda })/q^2`$. It is possible to show that the non-Abelian Coulomb term describes the complete quark-antiquark interaction in the heavy quark limit. Thus one may use lattice data to determine that the potential takes on the familiar Coulomb+linear form. It is natural to assume that, with the appropriate color structure, this form holds for gluonic color sources as well (in fact this is confirmed by lattice calculations ). Thus, in previous calculations we have simply replaced the expectation of the non-Abelian kernel with a linear potential. However, one of our current goals is to demonstrate that the BCS Ansatz and linear confinement can be described consistently by the same formalism. The central idea is that the gap equation contains a kernel which is itself the expectation value of the non-Abelian Coulomb interaction in the BCS vacuum. Thus the kernel is functionally dependent on the gap function and it is possible to obtain both the gap function and the effective potential. Since it is known that the non-Abelian Coulomb interaction gives rise to confinement in the heavy quark limit, one can hope that this procedure will yield a linear potential at large distances ,. In an attempt to gain some insight into this hypothesis, we consider perturbative corrections to the effective potential. As stated above, the non-Abelian Coulomb kernel is an operator which depends on the transverse degrees of freedom of the gauge field (cf. Eq. (3)). It can be shown that the $`\beta `$ function can be obtained from loop corrections to this kernel . One finds that at $`O(\alpha _s)`$, there are three contributions to $`\beta =\beta _0/(2\pi )\alpha _s^2`$: $`\beta _0=\beta _c+\beta _g+\beta _q=+4N_c+(1)+(2n_f/3)`$. The first term, $`\beta _c`$ comes directly from the expansion of the non-Abelian Coulomb potential. The second and third terms are due to mixing with two-gluon and two-quark intermediate states respectively. In the following we consider the perturbative corrections to the Coulomb kernel since it gives rise to the majority of the pure gauge $`\beta `$ function. The contribution from the transverse gluons in the expansion of the Coulomb kernel may be computed with the expansion $$[\frac{1}{\mathbf{}𝐃}(\mathbf{}^2)\frac{1}{\mathbf{}𝐃}]_{ab}=\frac{1}{\mathbf{}^2}\left[\delta _{ab}+2g_1\frac{1}{\mathbf{}^2}f_{abc}𝐀^c\mathbf{}+3g_2\frac{1}{\mathbf{}^2}f_{acd}𝐀^d\mathbf{}\frac{1}{\mathbf{}^2}f_{cbe}𝐀^e\mathbf{}+\mathrm{}\right]$$ (29) with the cutoff dependence of the couplings, $`g_ig(\mathrm{\Lambda })Z_i(\mathrm{\Lambda })`$ to be determined. To order $`g_i^2`$, the Coulomb interaction is given by $$V_c(q)=\frac{4\pi \alpha _s}{q^2}\left[1+\frac{3Z_2^2(\mathrm{\Lambda })\alpha _s(\mathrm{\Lambda })\beta _c}{16\pi ^2}I(q,\mathrm{\Lambda };\omega _0)\right]$$ (30) where $`\alpha _s(\mathrm{\Lambda })=g^2(\mathrm{\Lambda })/4\pi `$ and $$I(q,\mathrm{\Lambda };\omega _0)=d^3𝐤\frac{(1(\widehat{𝐤}\widehat{𝐪})^2)\mathrm{e}^{k^2/\mathrm{\Lambda }^2}}{\omega _0(k)(𝐪𝐤)^2}.$$ (31) In the BCS vacuum the same expression arises except that $`\omega _0`$ is replaced with the unknown function, $`\omega `$. Thus, as stated above, the solution to the gap equation depends on a potential which itself is a function of the gluon dispersion relation. It therefore becomes possible to obtain the gluon gap function and the nonperturbative potential in a self consistent fashion. We pursue this scenario by defining a nonperturbative coupling via the relation $$V_c(q)\frac{4\pi \stackrel{~}{\alpha }_{eff}(q^2)}{q^2}.$$ (32) Furthermore, we take part of the rainbow-ladder approximation to $`V_c`$ by summing all diagrams proportional to $`I^n`$. This leads to an expression familiar from the perturbative leading log approximation, $$\stackrel{~}{\alpha }_{eff}(q^2)=\frac{\alpha _s(\mathrm{\Lambda })}{1\frac{3Z_2^2(\mathrm{\Lambda })\alpha _s(\mathrm{\Lambda })\beta _c}{16\pi ^2}I(q,\mathrm{\Lambda };\omega )}.$$ (33) Although this potential may be inserted directly into the gap equation we have found it useful to approximate it with the aid of the following substitution $$I[\omega ]\frac{4\pi }{3}\mathrm{log}\left[\frac{\mathrm{\Lambda }^2}{\omega (q)^2}\right].$$ (34) As shown in Fig. 1, this form is accurate to roughly 10% for $`q/\mathrm{\Lambda }<0.1`$. With this expression for $`V_c`$, the gap equation is given by $`{\displaystyle \frac{\omega (q)^2}{Z}}`$ $`=`$ $`Zq^2+M^2(\mathrm{\Lambda })+c_0(\mathrm{\Lambda })q^2+\pi \alpha _s(\mathrm{\Lambda })N_c{\displaystyle \frac{d^3k}{(2\pi )^3}\mathrm{e}^{\frac{k^2}{\mathrm{\Lambda }^2}}\frac{(3(\widehat{k}\widehat{q})^2)}{\omega (k)}}`$ (35) $`+`$ $`{\displaystyle \frac{N_c}{4}}{\displaystyle \frac{d^3k}{(2\pi )^3}\mathrm{e}^{\frac{k^2}{\mathrm{\Lambda }^2}}\stackrel{~}{V}_c(𝐤+𝐪)\left(1+(\widehat{q}\widehat{k})^2\right)\frac{\omega ^2(k)\omega ^2(q)}{\omega (k)}}`$ (36) Many of the terms of higher order in $`1/\mathrm{\Lambda }`$ also contribute to the gap equation. Additional terms include $`c_2(\mathrm{\Lambda })q^2\omega (q)^2/\mathrm{\Lambda }^2`$, $`d_2(\mathrm{\Lambda })q^4/\mathrm{\Lambda }^2`$, and $`e_2(\mathrm{\Lambda })/\mathrm{\Lambda }^2N_c/3\mathrm{exp}(k^2/\mathrm{\Lambda }^2)(\omega (k)^2\omega (q)^2)/\omega (k)`$. As expected, these only affect the solution at short range, thus the long range behavior of the effective potential is completely specified by the nonperturbative model (the BCS Ansatz in this case). Nevertheless, one may get an indication about the viability of the low energy model by studying the coefficients, $`c_2`$, $`d_2`$, etc – a strong cutoff dependence at low momenta would indicate that the nonperturbative model is inadequate and that additional counterterms are therefore required. At this point the standard procedure would be to choose values for the cutoff and coefficients, solve the gap equation to obtain $`\omega `$, use the solution to obtain, say, glueball masses<sup></sup><sup></sup>*This procedure is followed, without renormalization, in Ref. ., and vary the coefficients so that the predictions agree with experimental (or lattice) data. However, in this study we choose to perform a simpler analysis and fit the derived form of the effective potential, $`V_{eff}=\stackrel{~}{V}(\omega )`$ to the Wilson loop lattice potential. This is similar in spirit to a lattice renormalization procedure advocated by Lepage and Mackenzie. Dimension six operators are neglected so that the couplings to be varied are the five coefficients $`Z`$, $`c_2`$, $`c_0`$, $`Z_2`$, and $`\alpha _s`$. First we note that not all of the coefficients are independent since one requires that the effective potential, and therefore $`\omega `$, is $`\mathrm{\Lambda }`$ independent (to leading order in $`p/\mathrm{\Lambda }`$). Multiplying the gap equation by $`Z`$ and introducing $`\stackrel{~}{\alpha _s}=Z\alpha _s`$, $`\stackrel{~}{c}_2=Z(c_2+L)`$, $`\stackrel{~}{c}_0=Z^2+Zc_0`$ leads to $$\omega ^2(q)=\stackrel{~}{c}_0(\mathrm{\Lambda })q^2+\mathrm{\Lambda }^2\stackrel{~}{c}_2(\mathrm{\Lambda })+\frac{N_c}{(4\pi )^2}_0^{\mathrm{}}𝑑kk^2d(\widehat{q}\widehat{k})\stackrel{~}{V}_c(𝐤+𝐪)\left(1+(\widehat{q}\widehat{k})^2\right)\frac{\omega ^2(k)\omega ^2(q)}{\omega (k)}$$ (37) where $$\stackrel{~}{V}_c(q)=\frac{4\pi \stackrel{~}{\alpha }_s(\mathrm{\Lambda })}{q^2\left(1\frac{\stackrel{~}{\alpha }_s(\mathrm{\Lambda })\beta _c}{4\pi }\mathrm{log}\left[\frac{\mathrm{\Lambda }^2}{\omega ^2(q)}\right]\right)}$$ (38) and $`L`$ is defined to be the constant, $`\mathrm{\Lambda }`$-dependent contribution from the third term on the right hand side of Eq. (36) (this term comes from the $`A^4`$ operator in the Hamiltonian). In obtaining Eq. (37) we have further chosen $`Z=Z_2^2`$ which is required to make the effective interaction cutoff independent. We now define a QCD scale, $`\mu `$, as $$\mu ^2\mathrm{\Lambda }^2\mathrm{e}^{\frac{4\pi }{\stackrel{~}{\alpha }_s(\mathrm{\Lambda })\beta _c}}$$ (39) and set $`m^2(q)\omega ^2(q)q^2`$. The expression for the running coupling becomes cutoff independent (as long as $`m(q)`$ is $`\mathrm{\Lambda }`$-independent) and is given by $$\alpha _{eff}(q^2)=\frac{4\pi }{\beta _c\mathrm{log}(\frac{q^2+m(q)^2}{\mu ^2})},\stackrel{~}{V}_c(q)=\frac{4\pi \alpha _{eff}(q^2)}{q^2}.$$ (40) We note that if $`m(0)=\mu `$ the long range behavior of the effective potential corresponds to a linear potential in position space. Equating the coefficient of the leading $`1/q^4`$ behavior with the standard form $`6\pi b/q^4`$ ($`b`$ is the string tension), yields $`\mu ^2=3(1+m^{})\beta _cb/(8\pi )`$ where $`m^{}=dm^2(0)/dq^2`$. The two remaining coefficients in the gap equations can thus be fixed by demanding that the effective potential represents linear confinement with the right slope. In numerical computations, instead of fixing $`m^{}`$ we have chosen to fix $`m^2(q)`$ at a finite value of $`q^2`$, typically $`q^2=1\text{ GeV}^2`$, which is numerically easier. The two renormalization conditions, one for $`m(0)`$ and one for $`m^{}`$ (or in general for $`m^2(q_0)`$ and $`m^2(q_1)`$), can be implemented directly into the gap equation by two subtractions, $`\omega ^2(q)`$ $`=`$ $`q^2{\displaystyle \frac{\omega _1^2\omega _0^2}{q_1^2q_0^2}}+{\displaystyle \frac{q_1^2\omega _0^2q_0^2\omega _1^2}{q_1^2q_0^2}}+{\displaystyle _0^{\mathrm{}}}𝑑k𝒱(k,q){\displaystyle \frac{\omega ^2(k)\omega ^2(q)}{\omega (k)}}`$ (41) $`+`$ $`{\displaystyle \frac{q^2q_1^2}{q_1^2q_0^2}}{\displaystyle _0^{\mathrm{}}}𝑑k𝒱(k,q_0){\displaystyle \frac{\omega ^2(k)\omega _0^2}{\omega (k)}}+{\displaystyle \frac{q_0^2q^2}{q_1^2q_0^2}}{\displaystyle _0^{\mathrm{}}}𝑑k𝒱(k,q_1){\displaystyle \frac{\omega ^2(k)\omega _1^2}{\omega (k)}}`$ (42) where $`\omega _0=\omega (q_0)`$, $`\omega _1=\omega (q_1)`$ and $$𝒱(k,q)=\frac{N_c}{(4\pi )^2}_1^1d(\widehat{q}\widehat{k})k^2\mathrm{e}^{\frac{k^2}{\mathrm{\Lambda }^2}}\stackrel{~}{V}_c(𝐤+𝐪)(1+(\widehat{q}\widehat{k})^2)$$ (43) The original couplings may be reconstructed as follows: $$\mathrm{\Lambda }^2\stackrel{~}{c}_2=\frac{q_1^2\omega _0^2q_0^2\omega _1^2}{q_1^2q_0^2}\frac{q_1^2}{q_1^2q_0^2}_0^{\mathrm{}}𝑑k𝒱(k,q_0)\frac{\omega ^2(k)\omega _0^2}{\omega (k)}+\frac{q_0^2}{q_1^2q_0^2}_0^{\mathrm{}}𝑑k𝒱(k,q_1)\frac{\omega ^2(k)\omega _1^2}{\omega (k)}$$ (44) $$\stackrel{~}{c}_0=\frac{\omega _1^2\omega _0^2}{q_1^2q_0^2}\frac{1}{q_1^2q_0^2}_0^{\mathrm{}}𝑑k𝒱(k,q_1)\frac{\omega ^2(k)\omega _1^2}{\omega (k)}+\frac{1}{q_1^2q_0^2}_0^{\mathrm{}}𝑑k𝒱(k,q_0)\frac{\omega ^2(k)\omega _0^2}{\omega (k)}$$ (45) The single gluon energy is given by $$E(q)=\omega (q)\left[Z^1+_0^{\mathrm{}}𝑑k\frac{𝒱(k,q)}{\omega (k)}\right]$$ (46) If a linear potential is to be recovered (as it must because we are fitting to lattice data), one must have $`\omega _0=\mu `$ at $`q_0=0`$, (see Eq. (40)). Thus the QCD scale $`\mu `$ must be related to the string tension. As discussed above, instead of taking the limit $`q_1q_0`$ and fixing $`d\omega ^2(0)/dq^2`$ to reproduce the magnitude of string tension we keep $`q_1`$ finite, choose $`q_1=1\text{ GeV}`$ and fit $`\omega _1`$ to reproduce the strength of the confining potential at low momentum. In our computations the one dimensional nonlinear integral (after numerical evaluation of the angular integrals) equation was solved numerically using a modified Levenberg–Marquardt algorithm. We note that the equations are extremely sensitive to the proper treatment of the IR singularity (which, however, is integrable) and that some solutions which have appeared in the literature are misleading. The results for $`\alpha _{eff}(q)=q^2V_{eff}(q)/4\pi `$ are shown in Fig. 2. We find that $`m^{}0`$ (typically $`m^{}=0.01`$) thus the QCD scale $`\mu `$ can be directly calculated from the physical string tension. Taking $`b=0.18\text{ GeV}^2`$ yields $`\mu ^2=0.2579\text{ GeV}^2`$. Finally, for large momenta the running of the effective coupling is given by $$\alpha _{eff}^{UV}(q)=\frac{4\pi }{\beta _c\mathrm{log}\left(\frac{q^2}{0.2579GeV^2}\right)}.$$ (47) Over the a range of momenta the full solution to the effective coupling (or potential) can be well approximated by a running coupling version of the Cornell potential $$\alpha _{eff}(q)=\frac{4\pi }{\beta _c\mathrm{log}\left[\frac{q^2}{\mu ^2}+\left(\frac{\mu _0}{\mu }\right)^2\right]}+\frac{3}{2}\frac{b}{q^2}=\frac{\alpha _{eff}^{UV}(\mathrm{\Lambda })}{\left(1+\frac{\alpha _{eff}^{UV}(\mathrm{\Lambda })\beta _c}{4\pi }\mathrm{log}\frac{p^2+\mu _0^2}{\mathrm{\Lambda }^2}\right)}+\frac{3}{2}\frac{b}{q^2}.$$ (48) Here $`\mu _0`$ is a parameter which allows interpolation between the infrared and ultraviolet limits of the potential. We find $`(\mu _0/\mu )^210`$. Essentially identical results are obtained upon comparison with the standard Coulomb+linear form of the Cornell potential. That the form of the derived effective potential so nearly reproduces linear confinement at long range is a tantalizing indication that our goal of deriving the confinement potential in a self-consistent manner from the BCS vacuum Ansatz may be achieved. The derived gluon dispersion relation $`\omega (q)`$ is shown in Fig. 3. The doubly subtracted form of the gap equation removes $`\mathrm{\Lambda }^2`$ and log$`(\mathrm{\Lambda })`$ dependence from the gap equation and leaves subleading terms of order $`q^2/\mathrm{\Lambda }^2`$. This is reflected in the cutoff independence evident in Fig. 3 at low momentum. In Fig. 3 we also compare the full solution of the gap equation for $`\mathrm{\Lambda }=10`$ with the approximate one obtained using the “Coulomb+linear” potential with the parameters fitted to the self consistent solution as discussed above. Both of them lead to $`m^{}0`$ and both are renormalized at the same point, $`q_1=1`$ GeV, however the intermediate $`q`$ dependence is somewhat different. As seen from Fig. 2, this small difference does not alter the behavior of the effective interaction. One may expect that $`\omega (q)`$ should approach the free solution, $`\omega _0(q)=q`$ at large momenta. However, this need not be true, both in perturbation theory, where logarithms can ruin the asymptotic behavior, and nonperturbatively. This is because the counterterms $`ZB^2`$ and $`c_0AA`$ give rise to the term $`\stackrel{~}{c}_0q^2`$ in the gap equation. Renormalizing $`\omega (q)`$ at a low momentum scales $`q_0,q_1<<\mathrm{\Lambda }`$ will result in a logarithmic dependence on $`\mathrm{\Lambda }`$ in $`\stackrel{~}{c}_0`$. Thus the effective gluon mass $`m^2(q)=\omega ^2(q)q^2`$ will remain proportional to $`q^2`$ at large momenta. As mentioned above, this behavior is also true in perturbation theory. In the following we examine the small coupling limit of the counterterm coefficients to establish contact with perturbative QCD. Evaluating Eqns 37 and 38 to $`O\left(\alpha _s=\stackrel{~}{\alpha }_s(\mathrm{\Lambda })\right)`$ yields $$\stackrel{~}{c}_2=\frac{\alpha _sN_C}{4\pi }\left[\frac{4}{3}\frac{8}{3}\right]=\frac{\alpha _s\beta _0}{N_C}$$ (49) where the first contribution comes from the vacuum expectation value of the Coulomb potential and the second from the $`A^4`$ term in the Hamiltonian, and $$\stackrel{~}{c}_0=1+\frac{\alpha _sN_c}{4\pi }\frac{8}{15}\frac{q_1^2+q_0^2}{q_1^2q_0^2}\mathrm{ln}\left(\frac{q_0}{q_1}\right)+\frac{\alpha _sN_c}{4\pi }\frac{8}{15}\mathrm{ln}\left(\frac{\mathrm{\Lambda }^2}{q_0q_1}\right)=\frac{\alpha _sN_c}{4\pi }\frac{8}{15}\mathrm{ln}\mathrm{\Lambda }^2+\text{finite}.$$ (50) In perturbation theory, to $`O(\alpha _s)`$, $`\mathrm{\Lambda }`$-dependence of $`Z`$ and $`c_0`$ follows from Eq. (50) and the requirement that the gluon energy (an observable in perturbation theory) is independent of $`\mathrm{\Lambda }`$. Thus from Eq. (46) it follows that $$E(q)=q\left[Z^1+\frac{\alpha _sN_c}{4\pi }\frac{4}{3}\mathrm{ln}\frac{\mathrm{\Lambda }^2}{q^2}\right]$$ (51) which leads to $$Z=\frac{\alpha _sN_c}{4\pi }\frac{4}{3}\mathrm{ln}\mathrm{\Lambda }^2+\text{finite}$$ (52) and finally $$c_0=\frac{\alpha _sN_c}{4\pi }\frac{32}{15}\mathrm{ln}\mathrm{\Lambda }^2$$ (53) These expressions agree with the corresponding expressions in a full perturbative Hamiltonian QCD calculation. Note, however, that in Ref. it has been shown that $`c_0=0`$ to $`O(\alpha _s)`$. This arises because the self energy contribution from a gluon loop with transverse gluon exchange cancels the contribution from the Coulomb and the $`A^4`$ vacuum expectation values Note that the UV cutoff $`\mathrm{\Lambda }`$ used in Ref. is by a factor of $`\sqrt{2}`$ larger then the one used here.. The BCS calculation only contains the latter two terms. In a nonperturbative calculation based on the BCS Ansatz transverse gluons are confined and the single gluon energy of Eq. (46) becomes IR singular. As a result, transverse gluons only affect observables in bound state calculations of color singlet states. In Fig.4 we show the behavior of the counterterms, $`\stackrel{~}{c}_2`$ and $`\stackrel{~}{c}_0`$ for small couplings. The points represent the full expressions obtained from Eqns 37 and 38. The lines are the perturbative results given above with the finite portions fit to the full results. As expected the agreement is very good for small coupling. Fig. 4a shows a dramatic dependence of $`\stackrel{~}{c_2}`$ on the value of the UV cutoff. This dependence originates with the subleading $`1/q^4`$ behavior of the effective potential which makes a log$`(\mathrm{\Lambda })`$ contribution to $`\stackrel{~}{c}_2`$. The solid lines in Fig. 4a incorporate this behavior. The logarithmic term dominates in the small coupling limit and leads to the large relative splitting seen in the figure. ### C Relationship to Other Approaches The Schwinger-Dyson formalism if often used to discuss the issues presented here. Advantages are that it is covariant and that calculations beyond the rainbow approximation are possible. However, the resulting integral equations are four dimensional and are therefore difficult to solve. They are also normally evaluated in Euclidean space which introduces interpretational difficulties, especially when dealing with confinement. Furthermore, the connection to the parton model is unclear since it is difficult to continue wavefunctions to Minkowski space. In the Hamiltonian framework color nonsinglets decouple from the physical spectrum as in the Schwinger-Dyson approach, but hadronic wavefunctions are well-defined, and thus a parton picture of high energy QCD is possible. We also note that mean field formalisms of Coulomb gauge QCD have been extensively studied in Refs. and . The numerical solution presented in Ref. are however doubtful, as they do not have the proper large momentum behavior. In Ref. the Coulomb potential was studied with the aid of a model of the gluon dispersion function (the gap equation was not solved). A calculation with similar goals to those expressed here has recently appeared. This paper employs the flow equation renormalization methodology to derive a gluonic gap equation which is similar in form to Eq. (36). Differences are that convergence factors have been inserted by hand into two of their loop integrals, a phenomenological potential has been used for the interaction kernel, and the perturbative expression for the gluon mass counterterm has been used. Their calculation also includes transverse gluon contributions at one loop, which is incorrect since the low momentum transverse gluon propagator is strongly modified by the confining interaction. As stated above, the implication of this is that any propagator of color nonsinglet states is IR divergent and that transverse gluons should only be incorporates into color singlet bound state calculations. We stress our belief that the interaction can, and should, be derived in a self-consistent fashion and that the use of perturbative renormalization is not a viable approach to nonperturbative problems. A related problem is seen in their Figure 2 which shows $`\omega (q)q`$ approaching zero for $`q>4`$ GeV. This is not correct, $`\omega `$ will in fact not approach $`q`$ at large momentum, but some multiple of $`q`$. Curtis and Pennington have noted that the gap equation derived in the Schwinger-Dyson formalism is not multiplicatively renormalizable. This means that it is not possible to exchange the $`\mathrm{\Lambda }`$-dependence of the gap solution with a renormalization point dependence. As stated above, multiplicative renormalizability is not a concern in the effective field theory. An obvious manifestation of the difference is the appearance of noncanonical counterterms like the gluon mass term. In our approach the cutoff has to be kept finite, otherwise neglected contributions, i.e. from transverse gluons, would cause new divergences to appear. Instead we keep $`\mathrm{\Lambda }`$ finite and use higher dimensional local operators to account for short distance contributions from transverse gluons. In this way, the neglected contributions are kept finite (below the cutoff) and can be evaluated by expanding the Fock space used to diagonalize the Hamiltonian. ## III Conclusions We have presented a computational scheme for evaluating vacuum and quasiparticle properties of QCD. The scheme is based on the QCD Hamiltonian in Coulomb gauge, a field-based regularization, and the effective field theory approach to renormalization. The use of the Hamiltonian is advantageous because all degrees of freedom are physical and a transparent picture of confinement emerges – confinement is generated by the instantaneous non-Abelian Coulomb interaction. We found it useful to regulate the theory by field-smearing because this greatly reduced the complexity of the computations. Field-based calculations have forced us to adopt a field-based renormalization procedure and we have found that the effective field theory approach is ideal for our work. A gap equation was derived from the QCD Hamiltonian using the Bogoliubov-Valatin transformation. The gap equation necessarily involves the Coulomb interaction kernel which should also be evaluated in the BCS vacuum. Thus self-consistent gap and effective instantaneous color source interaction equations may be derived. With the aid of an approximation to the matrix element of the full kernel in the BCS vacuum, we have found that it is possible to obtain an effective potential which agrees well with the (Abelian) Coulomb and linear regimes of the Wilson loop measurement of the heavy source nonperturbative potential. Furthermore, the derived effective potential is nearly cutoff independent, indicating that a good low energy description has been found. Encouraged by this success, we look forward to repeating the calculation with the exact BCS interaction kernel. We fully expect this calculation to be equally successful. It will also be of interest to examine the consistency of our results with those derived in the quark sector. Since the quark and gluon vacuum sectors couple at one loop, an examination of the coupled gap equations should prove illuminating. Finally the efficacy of the methodology presented here can be tested by calculating glueball and meson spectra. ###### Acknowledgements. ES acknowledges support from the DOE under grant DE-FG02-96ER40944 and DOE contract DE-AC05-84ER40150 under which the Southeastern Universities Research Association operates the Thomas Jefferson National Accelerator Facility. AS acknowledges support from the DOE under grant DE-FG02-87ER40365.
warning/0005/astro-ph0005104.html	ar5iv	text	# Newtonian limit of Induced Gravity ## 1 Introduction Extended theories of gravity have become a sort of paradigm in the study of gravitational interaction since several motivations push for enlarging the traditional scheme of Einstein general relativity. Such issues come, essentially, from cosmology and quantum field theory. In the first case, it is well known that higher–derivative theories and scalar–tensor theories furnish inflationary cosmological solutions capable, in principle, of solving the shortcomings of standard cosmological model . Besides, they have relevant features also from the quantum cosmology point of view since they give interesting solutions to the initial condition problem, at least in the restricted context of minisuperspaces . In the second case, every unification scheme as superstrings, supergravity or great unified theories, takes into account effective actions where nonminimal coupling to the geometry or higher order terms in the curvature invariants come out. Such contributions are due to one–loop or higher–loop corrections in the high curvature regimes near the full (not yet available) quantum gravity regime . However, in the weak–limit approximation, all these classes of theories should be expected to reproduce the Einstein general relativity which, in any case, is experimentally tested only in this limit . This fact is matter of debate since several relativistic theories do not reproduce Einstein results at the Newtonian approximation but, in some sense, generalize them. In fact, as it was firstly noticed by Stelle , a $`R^2`$–theory gives rise to Yukawa–like corrections to the Newtonian potential which could have interesting physical consequences. For example, some authors claim for explaining the flat rotation curves of galaxies by using such terms . Others have shown that a conformal theory of gravity is nothing else but a fourth–order theory containing such terms in the Newtonian limit and, by invoking these results, it could be possible to explain the missing matter problem without dark matter. Besides, indications of an apparent, anomalous, long–range acceleration revealed from the data analysis of Pioneer 10/11, Galileo, and Ulysses spacecrafts could be framed in a general theoretical scheme by taking into account Yukawa–like or higher order corrections to the Newtonian potential . In general, any relativistic theory of gravitation can yield corrections to the Newton potential (see for example ) which in the post-Newtonian (PPN) formalism, could furnish tests for the same theory . Futhermore the newborn gravitational lensing astronomy is giving rise to additional tests for general relativity over small, large, and very large scales which very soon will provide direct measurements for the variation of Newton coupling $`G_N`$ , the potential of galaxies and clusters of galaxies and several other features of gravitating systems. Such data will be very likely capable of confirming or ruling out the physical consistency of extended theories of gravity. In this paper, we want to discuss the Newtonian limit of a class of scalar–tensor theories of gravity, the induced gravity theories, which are invoked in order to solve several problems in cosmology (e.g. extended inflation ) and in fundamental physics (e.g. the tree–level–string–dilaton effective action can be recast as an induced gravity theory ). In particular, we show that the Newtonian limit of an induced gravity theory depends on the parameters of the nonminimal coupling between scalar field and Ricci scalar and the self–interaction potential of such a scalar field. Furthermore, a quadratic correction of the Newtonian potential strictly depends on the presence of the scalar–field potential which acts, in the low energy limit, as a cosmological constant (Sec.2). In Sec.3, we obtain a static, spherically symmetric, asymptotically flat exact solution which is the generalization, for the nonminimally coupled theory with coupling of the form $`F(\varphi )=\xi \varphi ^2`$, of the well–known Janis–Newman–Winicour solution to the Einstein–massless scalar equations . The deflection angle and some microlensing observables are deduced in Sec.4. Conclusions are drawn in Sec.5. ## 2 The field equations and the linearized solutions The most general action for a theory of gravity where a scalar–field is nonminimally coupled to the geometry is, in four dimensions, $$𝒜=d^4x\sqrt{g}\left[F(\varphi )R+\frac{1}{2}g^{\mu \nu }\varphi _{;\mu }\varphi _{;\nu }V(\varphi )+_m\right],$$ (2.1) where $`F(\varphi )`$, $`V(\varphi )`$ are generic functions of $`\varphi `$ and $`_m`$ is the ordinary matter Lagrangian density ,. We obtain the field equations by varying the action with respect to the metric tensor field $`g_{\mu \nu }`$ $$G_{\mu \nu }=T_{\mu \nu }^{(eff)},$$ (2.2) where $`G_{\mu \nu }`$ is the Einstein tensor, $$G_{\mu \nu }=R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R,$$ (2.3) and $`T_{\mu \nu }^{(eff)}`$, whose expression is found to be $`T_{\mu \nu }^{(eff)}={\displaystyle \frac{1}{F(\varphi )}}\{{\displaystyle \frac{1}{2}}\varphi _{;\mu }\varphi _{;\nu }+{\displaystyle \frac{1}{4}}g_{\mu \nu }g^{\alpha \beta }\varphi _{;\alpha }\varphi _{;\beta }`$ $``$ $`{\displaystyle \frac{1}{2}}g_{\mu \nu }V(\varphi )+`$ $`g_{\mu \nu }\mathrm{}F(\varphi )`$ $`+`$ $`F(\varphi )_{;\mu \nu }4\pi \stackrel{~}{G}T_{\mu \nu }\},`$ (2.4) is the effective stress-energy tensor containing terms of nonminimal coupling, kinetic terms, potential of the scalar field $`\varphi `$ and the usual stress-energy tensor of matter, $`T_{\mu \nu }`$, calculated from $`_m`$. $`\stackrel{~}{G}`$ is a dimensional, strictly positive, constant. In our units, with $`c=1`$, Einstein general relativity is obtained when the scalar field coupling $`F(\varphi )`$ is a constant and $`\stackrel{~}{G}`$ reduces to the Newton gravitational constant $`G_N`$ . Equations for the scalar field are found by varying the action with respect to the same field $$\mathrm{}\varphi RF^{}(\varphi )+V^{}(\varphi )=0,$$ (2.5) where $`F^{}(\varphi )=dF(\varphi )/d\varphi `$, $`V^{}(\varphi )=dV(\varphi )/d\varphi `$ and $`\mathrm{}\varphi g^{\mu \nu }\varphi _{;\mu }\varphi _{;\nu }`$. This Klein-Gordon equation can also be obtained from the contracted Bianchi identity . Field equations derived from the action $`(\text{2.1})`$ can be recast in a Brans–Dicke equivalent form by choosing $$\phi =F(\varphi ),\omega \left(\phi \right)=\frac{F(\varphi )}{2F^{}(\varphi )^2},W\left(\phi \right)=V\left(\varphi (\phi )\right).$$ (2.6) Action (2.1) now reads $$𝒜=d^4x\sqrt{g}\left[\phi R+\frac{\omega (\phi )}{\phi }g^{\mu \nu }\phi _{;\mu }\phi _{;\nu }W(\phi )+_m\right],$$ (2.7) that is nothing else but the extension of the original Brans–Dicke proposal (where $`\omega `$ is a constant), with $`\omega =\omega (\phi )`$ plus a potential term $`W(\phi )`$. By varying the new action with respect to $`g_{\mu \nu }`$ and the new scalar field $`\phi `$, we obtain again the field equations. The effective stress energy tensor in Eq.(2.2) and the Klein–Gordon equation (2.5), in the same units as before, become $`T_{\mu \nu }^{(eff)}={\displaystyle \frac{4\pi \stackrel{~}{G}}{\phi }}T_{\mu \nu }{\displaystyle \frac{\omega (\phi )}{\phi ^2}}\left(\phi _{;\mu }\phi _{;\nu }{\displaystyle \frac{1}{2}}g_{\mu \nu }g^{\alpha \beta }\phi _{;\alpha }\phi _{;\beta }\right)`$ $`+`$ $`+{\displaystyle \frac{1}{\phi }}\left(\phi _{;\mu \nu }g_{\mu \nu }\mathrm{}\phi \right)`$ $``$ $`{\displaystyle \frac{1}{2\phi }}g_{\mu \nu }W\left(\phi \right),`$ (2.8) $$2\omega \left(\phi \right)\mathrm{}\phi \frac{\omega \left(\phi \right)}{\phi }g^{\alpha \beta }\phi _{;\alpha }\phi _{;\beta }\frac{d\omega \left(\phi \right)}{d\phi }g^{\alpha \beta }\phi _{;\alpha }\phi _{;\beta }\phi R+\phi W^{}\left(\phi \right)=0.$$ (2.9) This last equation is usually rewritten eliminating the scalar curvature term $`R`$ with the help of Eq.(2), so that one obtains $$\mathrm{}\phi =\frac{1}{32\omega (\phi )}\left(4\pi \stackrel{~}{G}T2W\left(\phi \right)+\phi W^{}\left(\phi \right)\frac{d\omega \left(\phi \right)}{d\phi }g^{\alpha \beta }\phi _{;\alpha }\phi _{;\beta }\right).$$ (2.10) The minus sign in the denominator comes from the sign chosen in our action (2.7). Our aim now is to study the linearized equations derived from the action (2.1) or, equivalently, from (2.7). Before starting, we need a choice for the up to now arbitrary functions $`F(\varphi )`$ and $`V(\varphi )`$. A rather general choice is given by $$F(\varphi )=\xi \varphi ^m,$$ (2.11) $$V(\varphi )=\lambda \varphi ^n,$$ (2.12) where $`\xi `$ is a coupling constant, $`\lambda `$ gives the self–interaction potential strength, $`m`$ and $`n`$ are arbitrary, for the moment, parameters. This choice is in agreement with the existence of a Noether symmetry in the action (2.1) as discussed in ,. Furthermore, several scalar–tensor physical theories (e.g. induced gravity) admit such a form for $`F(\varphi )`$ and $`V(\varphi )`$. In order to recover the Newtonian limit, we write, as usual, the metric tensor as $$g_{\mu \nu }=\eta _{\mu \nu }+h_{\mu \nu },$$ (2.13) where $`\eta _{\mu \nu }`$ is the Minkoskwi metric and $`h_{\mu \nu }`$ is a small correction to it. In the same way, we define the scalar field $`\psi `$ as a perturbation, of the same order of the components of $`h_{\mu \nu }`$, of the original field $`\varphi `$, that is $$\varphi =\phi _0+\psi ,$$ (2.14) where $`\phi _0`$ is a constant of order unit. It is clear that for $`\phi _0=1`$ and $`\psi =0`$ Einstein general relativity is recovered. To write in an appropriate form the Einstein tensor $`G_{\mu \nu }`$, we define the auxiliary fields $$\overline{h}_{\mu \nu }h_{\mu \nu }\frac{1}{2}\eta _{\mu \nu }h,$$ (2.15) and $$\sigma _\alpha \overline{h}_{\alpha \beta ,\gamma }\eta ^{\beta \gamma }.$$ (2.16) Given these definitions, to the first order in $`h_{\mu \nu }`$, we obtain $$G_{\mu \nu }=\frac{1}{2}\left\{\mathrm{}_\eta \overline{h}_{\mu \nu }+\eta _{\mu \nu }\sigma _{\alpha ,\beta }\eta ^{\alpha \beta }\sigma _{\mu ,\nu }\sigma _{\nu ,\mu }\right\},$$ (2.17) where $`\mathrm{}_\eta \eta ^{\mu \nu }\varphi _{,\mu }\varphi _{,\nu }`$. We have not fixed the gauge yet. We pass now to the right hand side of Eq.(2.2), namely to the effective stress energy tensor.Up to the second order in $`\psi `$, the coupling function $`F(\varphi )`$ and the potential $`V(\varphi )`$, by using Eqs.(2.11) and (2.12), become $$F(\varphi )\xi \left(\phi _0^m+m\phi _0^{m1}\psi +\frac{m(m1)}{2}\phi _0^{m2}\psi ^2\right),$$ (2.18) $$V(\varphi )\lambda \left(\phi _0^n+n\phi _0^{n1}\psi +\frac{n(n1)}{2}\phi _0^{n2}\psi ^2\right).$$ (2.19) To the first order, the effective stress–energy tensor becomes $$\stackrel{~}{T}_{\mu \nu }=m\phi _0^{2m1}\eta _{\mu \nu }\mathrm{}_\eta \psi +m\phi _0^{2m1}\psi _{,\mu \nu }\frac{\lambda \phi _0^{m+n}}{2\xi }\eta _{\mu \nu }(4\pi \stackrel{~}{G})\frac{\phi _0^m}{\xi }T_{\mu \nu },$$ (2.20) and then the field equations are $$\frac{1}{2}\left\{\mathrm{}_\eta \overline{h}_{\mu \nu }+\eta _{\mu \nu }\sigma _{\alpha ,\beta }\eta ^{\alpha \beta }\sigma _{\mu ,\nu }\sigma _{\nu ,\mu }\right\}=m\phi _0^{2m1}\eta _{\mu \nu }\mathrm{}_\eta \psi m\phi _0^{2m1}\psi _{,\mu \nu }+$$ (2.21) $$+\frac{\lambda \phi _0^{m+n}}{2\xi }\eta _{\mu \nu }+(4\pi \stackrel{~}{G})\frac{\phi _0^m}{\xi }T_{\mu \nu }.$$ We can eliminate the term proportional to $`\psi _{,\mu \nu }`$ by choosing an appropriate gauge. In fact, by writing the auxiliary field $`\sigma _\alpha `$, given by Eq. (2.16), as $$\sigma _\alpha =m\phi _0^{2m1}\psi _{,\alpha },$$ (2.22) field equations (2.21) read $$\mathrm{}_\eta \overline{h}_{\mu \nu }m\phi _0^{2m1}\eta _{\mu \nu }\mathrm{}_\eta \psi \frac{\lambda \phi _0^{m+n}}{\xi }\eta _{\mu \nu }+(8\pi \stackrel{~}{G})\frac{\phi _0^m}{\xi }T_{\mu \nu }$$ (2.23) By defining the auxiliary field with components $`\stackrel{~}{h}_{\mu \nu }`$ as $$\stackrel{~}{h}_{\mu \nu }\overline{h}_{\mu \nu }m\phi _0^{2m1}\eta _{\mu \nu }\psi ,$$ (2.24) the field equations take the simpler form $$\mathrm{}_\eta \stackrel{~}{h}_{\mu \nu }=\frac{\lambda \phi _0^{m+n}}{\xi }\eta _{\mu \nu }+(8\pi \stackrel{~}{G})\frac{\phi _0^m}{\xi }T_{\mu \nu }$$ (2.25) The original perturbation field $`h_{\mu \nu }`$ can be written in terms of the new field as (with $`\stackrel{~}{h}\eta ^{\mu \nu }\stackrel{~}{h}_{\mu \nu }`$) $$h_{\mu \nu }=\stackrel{~}{h}_{\mu \nu }\frac{1}{2}\eta _{\mu \nu }\stackrel{~}{h}m\phi _0^{2m1}\eta _{\mu \nu }\psi .$$ (2.26) We turn now to the Klein-Gordon Eq.(2.5). $`\mathrm{}_\eta \psi `$ can be written, from the linearized Klein- Gordon equation, in terms of the matter stress energy tensor and of the potential term. If we calculate the scalar invariant of curvature $`R=g^{\mu \nu }R_{\mu \nu }`$ from Eq.(2), we find $$\mathrm{}\varphi +\frac{F^{}(\varphi )}{F(\varphi )}\left(\frac{1}{2}g^{\alpha \beta }\varphi _{;\alpha }\varphi _{;\beta }2V(\varphi )3\mathrm{}F(\varphi )4\pi \stackrel{~}{G}T\right)+V^{}(\varphi )=0,$$ (2.27) and, to the first order, it reads $$\mathrm{}_\eta \psi +\frac{\lambda (n2m)(n1)\phi _0^{n2}}{13\xi m^2\phi _0^{m2}}\psi =\frac{\lambda (2mn)\phi _0^{n1}}{13\xi m^2\phi _0^{m2}}+\frac{4\pi \stackrel{~}{G}m^2}{(13\xi m^2\phi _0^{m2})\phi _0}T.$$ (2.28) We work in the weak-field and slow motion limits, namely we assume that the matter stress-energy tensor $`T_{\mu \nu }`$ is dominated by the mass density term and we neglect time derivatives with respect to the space derivatives, so that $`\mathrm{}_\eta \mathrm{\Delta }`$, where $`\mathrm{\Delta }`$ is the ordinary Laplacian operator in flat spacetime. The linearized field equations (2.25) and (2.28) have, for point–like distribution of matter<sup>1</sup><sup>1</sup>1 To be precise, we can define a Schwarzschild mass of the form $$M=(2T_0^0T_\mu ^\mu )\sqrt{g}d^3x.$$ , which is $`\rho (r)=M\delta (r)`$, the following solutions: for $`n2m`$, $`n1`$, we get $`h_{00}`$ $``$ $`\left[(4\pi \stackrel{~}{G}){\displaystyle \frac{\phi _0^m}{\xi }}\right]{\displaystyle \frac{M}{r}}\left[{\displaystyle \frac{4\pi \lambda \phi _0^{m+n}}{\xi }}\right]r^2\left[(4\pi \stackrel{~}{G}){\displaystyle \frac{m^2\phi _0^{2m2}M}{13\xi m^2\phi _0^{m2}}}\right]{\displaystyle \frac{e^{pr}}{r}}+`$ (2.29) $`\left[{\displaystyle \frac{4\pi m\phi _0^{2m}}{n1}}\right]\mathrm{cosh}(pr),`$ $`h_{il}`$ $``$ $`\delta _{il}\left\{\left[(4\pi \stackrel{~}{G}){\displaystyle \frac{\phi _0^m}{\xi }}\right]{\displaystyle \frac{M}{r}}+\left[{\displaystyle \frac{4\pi \lambda \phi _0^{m+n}}{\xi }}\right]r^2+\left[(4\pi \stackrel{~}{G}){\displaystyle \frac{m^2\phi _0^{2m2}M}{13\xi m^2\phi _0^{m2}}}\right]{\displaystyle \frac{e^{pr}}{r}}\right\}+`$ (2.30) $`\delta _{il}\left[{\displaystyle \frac{4\pi \phi _0^{2m}m}{n1}}\right]\mathrm{cosh}(pr),`$ $`\psi `$ $``$ $`\left[(4\pi \stackrel{~}{G}){\displaystyle \frac{mM}{13\xi m^2\phi _0}}\right]{\displaystyle \frac{e^{pr}}{r}}\left[{\displaystyle \frac{4\pi \phi _0}{n1}}\right]\mathrm{cosh}(pr),`$ (2.31) where the parameter $`p`$ is given by $$p^2=\frac{\lambda (n2m)(n1)\phi _0^{m2}}{13\xi m^2\phi _0^{m2}}.$$ (2.32) For $`n=2m`$, we obtain $`h_{00}`$ $``$ $`\left[{\displaystyle \frac{(4\pi \stackrel{~}{G})\phi _0^m(14\xi m^2\phi _0^{m2})}{\xi (13\xi m^2\phi _0^{m2})}}\right]{\displaystyle \frac{M}{r}}\left[{\displaystyle \frac{4\pi \lambda \phi _0^{m+n}}{\xi }}\right]r^2\mathrm{\Lambda },`$ (2.33) $`h_{il}`$ $``$ $`\delta _{il}\left\{\left[{\displaystyle \frac{(4\pi \stackrel{~}{G})\phi _0^m(12\xi m^2\phi _0^{m2})}{\xi (13\xi m^2\phi _0^{m2})}}\right]{\displaystyle \frac{M}{r}}+\left[{\displaystyle \frac{4\pi \lambda \phi _0^{m+n}}{\xi }}\right]r^2+\mathrm{\Lambda }\right\},`$ (2.34) $`\psi `$ $``$ $`\left[{\displaystyle \frac{(4\pi \stackrel{~}{G})m}{(13\xi m^2\phi _0^{m2})\phi _0}}\right]{\displaystyle \frac{M}{r}}+\psi _0`$ (2.35) where $`\mathrm{\Lambda }=m\phi _0^{2m1}\psi _0`$ and $`\psi _0`$ are arbitrary integration constants. Let us note that the values $`m=1`$, $`n=2`$ and $`m=2`$, $`n=4`$ correspond to the well known couplings and potentials, i.e. $`F\varphi `$, $`V\varphi ^2`$ and $`F\varphi ^2`$, $`V\varphi ^4`$, respectively. Finally for $`n=1`$, we obtain $$h_{00}\left[\frac{(4\pi \stackrel{~}{G})\phi _0^m(14\xi m^2\phi _0^{m2})}{\xi (13\xi m^2\phi _0^{m2})}\right]\frac{M}{r}\left[\frac{4\pi \lambda \phi _0^{m+n}[1\xi m(m+1)\phi _0^{m2}]}{\xi (13\xi m^2\phi _0^{m2})}\right]r^2,$$ (2.36) $$h_{il}\delta _{il}\left\{\left[\frac{(4\pi \stackrel{~}{G})\phi _0^m(12\xi m^2\phi _0^{m2})}{\xi (13\xi m^2\phi _0^{m2})}\right]\frac{M}{r}+\left[\frac{4\pi \lambda \phi _0^{m+n}[1\xi m(m+1)\phi _0^{m2}]}{\xi (13\xi m^2\phi _0^{m2})}\right]r^2\right\},$$ (2.37) $$\psi \left[\frac{(4\pi \stackrel{~}{G})m}{(13\xi m^2)\phi _0}\right]\frac{M}{r}+\left[\frac{4\pi \lambda (2m1)\phi _0^{n1}}{13\xi m^2\phi _0^{m2}}\right]r^2.$$ (2.38) If we demand the $`(0,0)`$–component of the field Eq.(2.25), when $`\lambda =0`$, to read as the usual Poisson equation (that is nothing else but a definition of the mass) $$\mathrm{\Delta }\mathrm{\Phi }=4\pi G_N\rho ,$$ (2.39) where $`\mathrm{\Phi }`$, linked with the metric tensor by the relation $`h_{00}=2\mathrm{\Phi }`$, is the Newtonian potential, we have to put $$G_N=\frac{\phi _0^m}{2\xi }\left(\frac{14\xi m^2\phi _0^{m2}}{13\xi m^2\phi _0^{m2}}\right)\stackrel{~}{G}.$$ (2.40) We may now rewrite the nonzero components of $`h_{\mu \nu }`$ and the scalar perturbed field. Let us take into account, for example, Eqs. (2.33)–(2.35). We get $$h_{00}\frac{2G_NM}{r}\frac{4\pi \lambda \phi _0^{m+n}}{\xi }r^2\phi _0^{2m1}\psi _0$$ (2.41) $$h_{il}\delta _{il}\left\{\frac{2G_NM}{r}\left(\frac{12\xi m^2\phi _0^{m2}}{14\xi m^2\phi _0^{m2}}\right)+\frac{4\pi \lambda \phi _0^{m+n}}{\xi }r^2+m\phi _0^{2m1}\psi _0\right\},$$ (2.42) and $$\psi =\frac{2G_NM}{r}\left(\frac{\xi m\phi _0^{m1}}{14\xi m^2}\right)+\psi _0,$$ (2.43) where the Newton constant explicitly appears. Similar considerations hold in the other cases. What we have obtained are solutions of the linearized field equations, starting from the action of a scalar field nonminimally coupled to the geometry, and minimally coupled to the ordinary matter. Such solutions depend on the parameters which characterize the theory: $`\xi ,m,n,\lambda `$. The results of Einstein general relativity are obtained for $`F(\varphi )=F_0`$, with $`F_0`$ negatively–defined due to the sign choice in the action (2.1). As we can easily see from above, in particular from Eqs.(2.41)–(2.43), we have the usual Newtonian potential and a sort of cosmological term ruled by $`\lambda `$ which, from the Poisson equation, gives a quadratic contribution. We consider now the Brans-Dicke-like action (2.7) where $`\omega =\omega \left(\phi \right)`$. It is actually simple to see, from the field Eqs.(2) and (2.10), that, if we want to limit ourselves to the linear approximation, we may consider as well $`\omega =`$constant. The link with the results that follows from the action (2.7) are given by the transformation laws (2.6), that is $$\omega \left(\phi (\varphi )\right)=\frac{1}{2\xi m^2}\varphi ^{2m}.$$ (2.44) The potential term $`W\left(\phi \right)`$ in the linear approximation, behaves as $`V(\varphi )`$. Results in the approaches given by the actions (2.1) and (2.7) are completely equivalent. ## 3 An exact solution We are going now to consider a particular choice for the coupling field $`F(\varphi )`$ and the potential $`V(\varphi )`$ and to obtain an exact solution for the field equations in the case of the external field of a distribution of matter endowed with spherical symmetry. Namely, we consider $$F(\varphi )=\xi \varphi ^2,$$ (3.1) $$V(\varphi )=0.$$ (3.2) The physical relevance of the above assumption is widely discussed in literature (see e.g. and references therein). Here we want to show that it is possible to find an exact solution of the same form of those (more general) discussed in previous section. This is nothing else but a generalization to the nonminimal coupling case of the well known solution of Janis–Newman–Winicour ,. The exact solution that we will find is linked, of course, with the one obtained in Brans-Dicke theory with $`\omega =`$constant (from (2.44) it follows that it is just for the value $`m`$=2, that we find $`\omega =`$constant). We look for a solution for a static scalar field, so that, given the symmetry of the problem, the Birkoff theorem holds, and we can write the line element as $$ds^2=a(r)dt^2b(r)dr^2c(r)r^2d\mathrm{\Omega }^2,$$ (3.3) where $`a(r),b(r)`$ and $`c(r)`$ are strictly positive functions of the radial coordinate $`r`$, $`d\mathrm{\Omega }^2`$ is the usual spherical element. We can write down the Einstein field equations (2.2) as $$\frac{a}{r^2b}+\frac{a}{r^2c}+\frac{ab^{}}{rb^2}\frac{3ac^{}}{rbc}+\frac{ab^{}c^{}}{2b^2c}+\frac{ac^2}{4bc^2}\frac{ac^{\prime \prime }}{bc}=$$ $$=\frac{4a\varphi ^{}}{rb\varphi }\frac{ab^{}\varphi ^{}}{b^2\varphi }+\frac{2ac^{}\varphi ^{}}{bc\varphi }+\frac{2a\varphi ^2}{b\varphi ^2}\frac{a\varphi ^2}{4\xi b\varphi ^2}+\frac{2a\varphi ^{\prime \prime }}{b\varphi },$$ (3.4) $$\frac{1}{r^2}\frac{b}{r^2c}+\frac{a^{}}{ra}+\frac{c^{}}{rc}+\frac{a^{}c^{}}{2ac}+\frac{c^2}{4c^2}=\frac{4\varphi ^{}}{r\varphi }\frac{a^{}\varphi ^{}}{a\varphi }\frac{2c^{}\varphi ^{}}{c\varphi }\frac{\varphi ^2}{4\xi \varphi ^2},$$ (3.5) $$\frac{rca^{}}{2ab}\frac{r^2ca^2}{4a^2b}\frac{rcb^{}}{2b^2}\frac{r^2ca^{}b^{}}{4ab^2}+\frac{rc^{}}{b}+\frac{r^2a^{}c^{}}{4ab}\frac{r^2b^{}c^{}}{4b^2}\frac{r^2c^2}{4bc}+\frac{r^2ca^{\prime \prime }}{2ab}+\frac{r^2c^{\prime \prime }}{2b}=$$ $$=\frac{2rc\varphi ^{}}{b\varphi }\frac{r^2ca^{}\varphi ^{}}{ab\varphi }+\frac{r^2cb^{}\varphi ^{}}{b^2\varphi }\frac{r^2c^{}\varphi ^{}}{b\varphi }\frac{2r^2c\varphi ^2}{b\varphi ^2}+\frac{r^2c\varphi ^2}{4\xi b\varphi ^2}\frac{2r^2c\varphi ^{\prime \prime }}{b\varphi },$$ (3.6) where the prime now indicates the derivative with respect to $`r`$. The Klein-Gordon equation is $$\left(112\xi \right)\left(\mathrm{}\varphi \frac{\varphi ^2}{b\varphi }\right)=0.$$ (3.7) The case $`\xi =1/12`$ is the conformal case. Of course, the usual Schwarzschild solution of Einstein general relativity is recovered for $`\varphi =`$constant. We look for a solution of the form $$ds^2=\left(1\frac{\chi }{r}\right)^\alpha dt^2\left(1\frac{\chi }{r}\right)^\beta dr^2\left(1\frac{\chi }{r}\right)^\nu r^2d\mathrm{\Omega }^2,$$ (3.8) $$\varphi \left(r\right)=\left(1\frac{\chi }{r}\right)^\delta ,$$ (3.9) where $`\alpha ,\beta `$ and $`\nu `$ are the parameters that specify the solution (and we expect to find $`\alpha =\beta =1`$ and $`\nu =0`$ as a particular solution), $`\chi `$ is a constant related to the mass of the system which generates the gravitational field. Eq.(3.7) is satisfied when the algebraic relation $$4\delta +\alpha \beta +2\nu 2=0,$$ (3.10) holds. By studying Einstein equations, we find that a solution of the form (3.8) does exist for $$\alpha =2\delta +\epsilon ,$$ (3.11) $$\beta =2\delta \epsilon ,$$ (3.12) $$\gamma =2\delta \epsilon +1,$$ (3.13) where $`\epsilon `$ is an auxiliary parameter which we discuss below. The link between them and the coupling constant $`\xi `$ is the quadratic algebraic equation $$\delta ^2\left(112\xi \right)+\xi \left(1\epsilon ^2\right)=0.$$ (3.14) We notice that, in order to specify completely the solution, we have to give, besides $`\chi `$, the value of two of the three parameters $`\delta ,\epsilon `$ and $`\xi `$. It is only to recover the Schwarzschild solution that we need to specify $`\epsilon ^2=1`$. We have also to note that, in the weak field limit, that we are now going to study, the parameters $`\delta `$ and $`\epsilon `$ can be expressed in term of $`\xi `$. We may study the range of applicability of the found solution writing the Ricci scalar of curvature whose expression is given by $$R=\left[\frac{1\epsilon ^2}{112\xi }\right]\left[\frac{\chi ^2\left(1\frac{\chi }{r}\right)^{2\delta +\epsilon }}{2r^4\left(1\frac{\chi }{r}\right)^2}\right].$$ (3.15) We find, beside the usual singularity in $`r=0`$, a singularity in $`r=\chi `$, which is a null surface (see also ). Actually $`\chi `$ can take both positive and negative values, so that the solution is defined respectively for $`r>\chi `$ and $`r>0`$. We now expand the solution found at the first order in $`\chi /r`$ and we identify it with the result of the linearized field equations for a point-like distribution of matter, with mass $`M`$. We demand the Poisson equation to hold, that is, we are taking into account a particular case of what we considered in the last section. We find $$\chi =2\epsilon G_NM,$$ (3.16) and from Eq.(2.40), we have (for $`\phi _0=1,`$ and $`m=2`$) $$G_N=\frac{1}{2\xi }\left(\frac{116\xi }{112\xi }\right)\stackrel{~}{G}.$$ (3.17) $`\stackrel{~}{G}`$, as above, is a dimensional strictly positive constant. Moreover we find, from the $`(0,0)`$–component of the field Eq. (2.25) and Eq.(2.28) for the scalar perturbation field $`\psi `$, where we put $`m=2`$ and $`\lambda =0`$, a second equation for $`\delta ,\epsilon `$ and $`\xi `$, $$\left(116\xi \right)\delta =2\left(\epsilon 2\delta \right)\xi .$$ (3.18) Together with Eq.(3.16), this equation allows us to express the solution as a function of the parameter $`\xi `$ alone. It results $$\epsilon =\sqrt{\frac{112\xi }{116\xi }},$$ (3.19) and $$\delta =\frac{2\epsilon \xi }{112\xi }=\frac{2\xi }{\sqrt{(112\xi )(116\xi )}}.$$ (3.20) As it can be seen, the allowed range for $`\xi `$ is $`\xi <1/16`$, $`\xi >1/12`$ (from which one should exclude the value $`\xi =0`$, the value that would give $`\delta =0`$ and $`\epsilon ^2=1`$). The role of the parameter $`\epsilon `$ is now clarified: it tells us how much the solution differs from the usual Schwarzschild solution of general relativity since we are considering the nonminimal coupling $`F(\varphi )=\xi \varphi ^2`$. The linearized solution can then be written as $`ds^2`$ $`=`$ $`\left(1{\displaystyle \frac{2G_NM}{r}}\right)dt^2\left[1{\displaystyle \frac{2G_NM}{r}}\left({\displaystyle \frac{18\xi }{116\xi }}\right)\right]dr^2+`$ (3.21) $``$ $`\left[1{\displaystyle \frac{2G_NM}{r}}\left({\displaystyle \frac{18\xi }{116\xi }}\sqrt{{\displaystyle \frac{112\xi }{116\xi }}}\right)\right]r^2d\mathrm{\Omega }^2,`$ $$\varphi \left(r\right)=1\frac{2G_NM}{r}\left(\frac{2\xi }{116\xi }\right),$$ (3.22) which, by a linear transformation for the radial coordinate, can be put in the same form of the solution (2.41), (2.42) and (2.43) with the choice $`m=2`$ and $`\lambda =0`$. In order to recover the corresponding Brans-Dicke solution, with $`\omega =`$constant and $`\varphi \phi =\xi \varphi ^2`$, let us introduce the coordinate transformation $$\left(\frac{dr}{dr^{}}\right)^2=\left(\frac{r}{r^{}}\right)^2\frac{c(r)}{b(r)}=\left(\frac{r}{r^{}}\right)^2\left(1\frac{\chi }{r}\right)^{\nu \beta }=\left(\frac{r}{r^{}}\right)^2\left(1\frac{\chi }{r}\right),$$ (3.23) where, as already noted, $`\beta =1`$ and $`\nu =0`$. The new variable $$r=r^{}\left(1+\frac{\chi }{4r^{}}\right)^2,$$ (3.24) is consistent with (3.23). The line element in the new coordinate system is $$ds^2=\left(\frac{1\chi /r^{}}{1+\chi /r^{}}\right)^{2\alpha }dt^2\left(1+\frac{\chi }{r^{}}\right)^4\left(\frac{1\chi /r^{}}{1+\chi /r^{}}\right)^{2\beta +2}\left(dr^2+r^2d\mathrm{\Omega }^2\right),$$ (3.25) while the scalar field is $$\varphi \left(r^{}\right)=\left(\frac{1\chi /r^{}}{1+\chi /r^{}}\right)^{2\delta }.$$ (3.26) The Brans-Dicke solution is then immediately found (see also for the asymptotic discussion). The line element is the same as (3.25), whereas the scalar field is equal to $$\phi \left(r^{}\right)=\frac{1}{8\omega }\left(\frac{1\chi /r^{}}{1+\chi /r^{}}\right)^{4\delta }.$$ (3.27) If we introduce the new parameters $`\epsilon ^{}`$, $`\delta ^{}`$ and $`\nu ^{}`$ $$2\alpha \alpha ^{}\frac{2}{\epsilon ^{}},2\beta +2\beta ^{}\frac{2}{\epsilon ^{}}\left(\epsilon ^{}\delta ^{}1\right),2\nu \nu ^{}4\delta ,$$ (3.28) we find $$\nu ^{}=\frac{\delta ^{}}{\epsilon ^{}}.$$ (3.29) As it has been done for $`\xi `$, it is possible to express everything as a function of the parameter $`\omega `$. It is worth noticing that the above solutions are equivalent, up to a conformal transformation, to the Schwarzschild-like solution obtained in the contest of minimally coupled theory of gravitation with a scalar field given by the action $$𝒜=d^4x\sqrt{\stackrel{~}{g}}\left[k\stackrel{~}{R}+\frac{1}{2}\stackrel{~}{g}^{\mu \nu }\stackrel{~}{\varphi }_{;\mu }\stackrel{~}{\varphi }_{;\nu }\right],$$ (3.30) where $`k`$ is a dimensional, strictly negative, constant that fixes units and $`\sqrt{\stackrel{~}{g}}`$ is the square root of the determinant of conformally transformed metric. This solution is known as the Janis–Newman–Winicour one of the Einstein–massless scalar field theory . However, we have to stress that, when we look for the linearized equations, such a conformal transformation has to be performed, if we want to express the parameters that specify the solution as a function of the physical quantities of the system. For example, to define the mass, we must choose the coupling with the ordinary matter. That is, we must decide whether the Jordan frame description assumed for the action (2.1) or (2.7) is the physical one, or if, by coupling minimally to the scalar field also the Lagrangian for ordinary matter in action (3.30), the Einstein frame description is the physical one. Of course, in the two case we get different results. For a detailed discussion, see e.g. . Furthermore, deviations from the usual expression for the bending angle of light in Einstein general relativity, are already present at first order for solutions in the Jordan frame, but only at second order in the Einstein frame . In this sense the conformal equivalence of Schwarzschild-like solutions in the two frames is broken when we introduce ordinary matter. ## 4 Deflection Angle and Microlensing Observables In the previous sections, we dealt with the problem of the determination of the geometry of a static and spherically symmetric spacetime in a given nonminimally coupled theory of gravity. From this solution it is possible to calculate, in a straightforward way, the deflection of light. We start with a line element written in the form $$ds^2=a\left(r\right)dt^2b\left(r\right)\left(dr^2+r^2d\mathrm{\Omega }^2\right),$$ (4.1) from which the geodesic equation of motion for photons is found to be $$u^2\left(u_{\widehat{\varphi }\widehat{\varphi }}+u\right)+\frac{1}{2}\left(b^1b^{}a^1a^{}\right)\left(u^2+u_{\widehat{\varphi }}^2\right)=0,$$ (4.2) where $`u1/r`$ and $`u_{\widehat{\varphi }}du/d\widehat{\varphi }`$ (here $`\widehat{\varphi }`$ is the azimuthal angle). If we consider this equation in the same limit where our solution (2.41), (2.42) is valid, we get (putting $`\psi _0=0`$) $$a\left(r\right)=12\frac{G_NM}{r}\frac{4\pi \lambda \phi _0^{m+n}}{\xi }r^2,$$ (4.3) $$b\left(r\right)=1+2\frac{G_NM}{r}\gamma +\frac{4\pi \lambda \phi _0^{m+n}}{\xi }r^2,$$ (4.4) where $`\lambda `$ and $`\gamma =\gamma (\xi ,m)`$, which is given by $$\gamma =\frac{12\xi m^2\phi _0^{m2}}{14\xi m^2\phi _0^{m2}},$$ (4.5) parametrize the scalar–tensor theory. We can recast Eq.(4.2) as $$u_{\varphi \varphi }+u=G_NM\left(1+\gamma \right)\left(u^2+u_\varphi ^2\right).$$ (4.6) It is worth noticing that the potential term $`\lambda `$ does not appear in this equation, so that the deflection angle, $`\widehat{\alpha }`$, does not depend from it either (at least in the approximation order which we have considered). We have $$\widehat{\alpha }=2\frac{G_NM}{r_0c^2}\left(1+\gamma \right),$$ (4.7) where, as usual, $`r_0`$ is the minimum distance from the deflector, and, to this order, it is nothing else but an impact parameter. We have reintroduced the dimensional constant $`c`$. The result is similar to that given in where also the singular isothermal sphere model has been studied. These arguments specify how deflection of light depends on the parameters $`\xi `$ and $`m`$ of our scalar–tensor theory. This give us the chance of looking at the well known effect of gravitational lensing, as a possible test for these theories. We are now going to give the expression for some quantities relevant for this phenomenon. We consider the simplest situation, that is, the so called Schwarzschild lens. With this geometry usually two images are formed, but when they are no longer separables, and the source brightness is magnified, one speaks about microlensing. In this context the relevant quantity is the Einstein radius, $`R_E`$, that is given by $$R_E=\sqrt{\frac{2G_NM}{c^2}\left(1+\gamma \right)\frac{D_lD_{ls}}{D_s}},$$ (4.8) where $`D_l`$, $`D_{ls}`$ and $`D_s`$ stand respectively for the distances between observer and lens, lens and source, observer and source. The Einstein angle, $`\theta _E`$ is given by $`R_ED_l`$. The maximum amplification is equal to $$A_{max}=\frac{r_0^2+2R_E^2}{r_0\sqrt{r_0^2+4R_E^2}}.$$ (4.9) Clearly, also $`A_{max}`$ depends on the parameters of the given scalar–tensor theory and, in general, all the quantities of gravitational lensing can be parametrized by them. In conclusion, beside classical tests , gravitational lensing could be useful, in general, to test gravity. On the other hand, anomalies in lensing features could be explained by enlarging the available set of gravitational theories. It is interesting to observe that Eq. (4.5) allows to determine the relation between $`\xi `$ and $`m`$. In fact, inverting (4.5) one gets $$\xi m^2=\frac{1\gamma }{2(12\gamma )},$$ (4.10) considering, for simplicity, $`\phi _0=1`$. A measurement of $`\gamma `$, i.e. the deviation from the expected value of the gravitational mass $`G_NM`$, estimated for example by Keplerian motions, could roughly indicate the form of the relativistic theory of gravity which holds at a given scale. On the other hand, the result strictly depends on the accuracy by which $`\widehat{\alpha }`$ is measured. ## 5 Discussion and Conclusions The Newtonian limit of induced–gravity theories, where the scalar field is nonminimally coupled to the geometry, strictly depends on the parameters of the coupling and the self–interaction potential as it has to be by a straightforward PPN parametrization . In this paper, we have constructed the weak limit solution of induced–gravity assuming, as it is quite natural to do, power law couplings and potentials. As we have seen, the role of the self–interaction potential is essential to obtain corrections to the Newtonian potential. Such corrections are, in any case, constant, quadratic or Yukawa–like as for other generalized theories of gravity , so also induced gravity could be a candidate to solve the problems of flat rotation curves of spiral galaxies \[see e.g. the solution (2.29)\]. Essentially the corrections depends on the strength of the coupling and the “mass” of the scalar field $`\varphi `$ given by $`\lambda `$. Besides, we have always scale lengths where Einstein general relativity can be recovered. This fact could account why measurements inside the Solar System confirm such theory while outside of it there are probable deviations . Furthermore, it is possible to show that the induced-gravity picture and the Brans-Dicke picture are completely equivalent in the Newtonian limit as it trivially has to be. The so called Janis–Newman–Winicour solution of minimally coupled scalar–tensor gravity can be easily extended to nonminimally coupled theories (at least without self–interaction potential) and it appears as a particular case of the general Newtonian solution (2.41)–(2.43) discussed above. Finally all the lensing observables, first of all the deflection angle, are affected if we introduce a nonminimal coupling and, consequently, gravitational lensing could constitute a further tool to investigate relativistic theories of gravity. Vice versa, relativistic theories could explain anomalous effects in gravitational lensing. In the case which we have analysed, it is interesting to stress the fact that the self–interaction potential of the scalar field does not intervene while deviations and then the “parametrization” strictly depend on the coupling.
warning/0005/hep-ph0005324.html	ar5iv	text	# A general reduction method for one-loop N-point integralsTalk presented at Loops and Legs in Quantum Field Theory, April 2000, Bastei, Germany ## 1 Motivation The study of processes with multi–particle/jet final states is of major importance not only to test perturbative QCD, but also to obtain a precise estimate of the background for search experiments. At the Tevatron, multi–jet cross sections with up to six jets have been measured and compared to tree level QCD predictions . However, leading order predictions are very unstable with respect to variations of the renormalization and factorization scales, such that the NLO corrections are necessary to obtain reliable predictions. The calculation of e.g. a $`2N2`$ parton process at NLO requires the knowledge of $`N`$-point scalar and tensor one-loop integrals with massless propagators and external legs, for $`e^+e^{}Njets`$ one needs $`(N+1)`$-point integrals with one off-shell external leg. These integrals have to be further processed by reducing the tensor integrals to scalar integrals and, for $`N5`$, reducing the scalar $`N`$-point integrals to known integrals with less propagators. Of course, reduction methods have been worked out and successfully applied before , but mostly for integrals with massive propagators, as needed e.g. for electroweak radiative corrections. In contrast, when dealing with massless partons, the presence of infrared divergences requires a different approach. For this case, reduction methods have been worked out which are valid for up to 6 external legs. How to extend these results to more than 6 external legs was not clear until recently due to problems stemming from the inversion of kinematical matrices which were singular if the external momenta were kept in four dimensions. On the other hand, the use of helicity techniques requires 4-dimensional external momenta. The main ideas how to overcome this problem will be outlined in the following, leading to reduction formulas for scalar and tensor integrals valid for an arbitrary number of external legs. ## 2 Reduction of scalar integrals Consider the scalar integral $`I_N^n`$ $`=`$ $`{\displaystyle \frac{d^nk}{i\pi ^{n/2}}\frac{1}{_{l=1}^Nq_l^2}}`$ $`q_l`$ $`=`$ $`kr_l,r_l={\displaystyle \underset{j=1}{\overset{l}{}}}p_j`$ in $`n=42ϵ`$ dimensions with $`N`$ massless propagators, depicted in Fig. 1. The corresponding expression in Feynman parameter space reads $`I_N^n`$ $`=`$ $`(1)^N\mathrm{\Gamma }(Nn/2)`$ $`\times {\displaystyle _0^{\mathrm{}}}d^Nz{\displaystyle \frac{\delta (1_{l=1}^Nz_l)}{(zSz)^{Nn/2}}}.`$ The kinematic information is contained in the matrix $`S`$ which is related to the Gram matrix $`G`$ by $`S_{kl}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(r_lr_k)^2={\displaystyle \frac{1}{2}}(G_{kl}r_l^2r_k^2)`$ (1) $`G_{kl}`$ $`=`$ $`2r_kr_l,k,l=1,\mathrm{},N.`$ The external momenta are kept in four dimensions; momentum conservation implies $`r_N=0`$. The integral $`I_N^n`$ can be split into a singular part with less propagators and a finite, but formally higher dimensional integral as $`I_N^n`$ $`=`$ $`I_{div}+I_{fin}`$ $`I_{div}`$ $`=`$ $`{\displaystyle \frac{d^nk}{i\pi ^{n/2}}\frac{_{l=1}^Nb_lq_l^2}{_{l=1}^Nq_l^2}}`$ (2) $`I_{fin}`$ $`=`$ $`{\displaystyle \frac{d^nk}{i\pi ^{n/2}}\frac{\left[1_{l=1}^Nb_lq_l^2\right]}{_{l=1}^Nq_l^2}}`$ (3) $`=`$ $`\left({\displaystyle \underset{l=1}{\overset{N}{}}}b_l\right)(Nn1)I_N^{n+2},`$ where the coefficients $`b_l`$ have to fulfill the following equations: $`{\displaystyle \underset{l=1}{\overset{N}{}}}S_{kl}b_l`$ $`=`$ $`{\displaystyle \frac{1}{2}}`$ $`{\displaystyle \underset{l=1}{\overset{N1}{}}}G_{kl}b_l`$ $`=`$ $`r_k^2{\displaystyle \underset{l=1}{\overset{N}{}}}b_l,{\displaystyle \underset{l=1}{\overset{N1}{}}}r_l^2b_l=1`$ (4) The integral $`I_{div}`$ in (2) is a sum over reduced integrals where one propagator could be cancelled. Our aim is to show that the higher dimensional integrals $`I_N^{n+2}`$ in (3), which are very hard to calculate for $`N>4`$ and hinder a reduction, drop out for $`N>4`$. To this end we have to solve Eqs. (2) for the parameters $`b_l`$. Note that for $`N=5`$, expression (3) is of order $`ϵ`$, which means that $`I_{fin}`$ can be dropped right away in that case because $`(n+2)`$-dimensional integrals are finite for $`N4`$. Since only four 4-dimensional external vectors $`r_i^\mu `$ can be linearly independent in Minkowski space, rank($`S`$)=min($`N`$,6) and rank($`G`$)=min($`N`$-1,4). Hence the case $`N=6`$ is special since $`S`$ is still invertible for $`N=6`$, whereas $`G`$ is not. Therefore Eqs.(2) can still be solved by inverting $`S`$ for $`N=6`$, whereas for $`N>6`$ both, $`G`$ and $`S`$ have a vanishing determinant. In this case, the solution to (2) can be constructed in terms of the so-called pseudoinverse $`H`$ to $`G`$. It is defined by the properties $`HGH=H`$, $`GHG=G`$. Given a pseudoinverse $`H`$ a solution to $`Gx=y`$ exists if and only if $`y=GHy`$. Introducing a basis $`E_{l=1,\mathrm{},4}^\mu `$ of Minkowski space and defining the matrices $`R`$ and $`\stackrel{~}{G}`$ by $$r_j^\mu =\underset{l=1}{\overset{4}{}}R_{lj}E_l^\mu ,\stackrel{~}{G}_{jk}=2E_jE_k,$$ (5) the uniquely defined<sup>1</sup><sup>1</sup>1In general $`H`$ is not unique. In our case, the uniqueness follows from the symmetry of $`G`$. pseudoinverse is given by $$H=R^T(RR^T)^1\stackrel{~}{G}^1(RR^T)^1R.$$ (6) Now Eq. (2) can be solved for general $`N`$ . Doing so one finds $$\underset{l=1}{\overset{N}{}}b_l=0\text{ for }N6$$ such that the contribution (3) completely vanishes for all $`N6`$. This means that all scalar $`N`$point integrals, for arbitrary $`N`$, can – by recursion – be expressed in terms of box integrals. Explicitly, one finds the following reduction formulas: For $`N6`$, i.e. $`det(S)0`$: $`I_N^n`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{k,l=1}{\overset{N}{}}}S_{kl}^1I_{N1,l}^n`$ (7) $`+(Nn1){\displaystyle \frac{det(G)}{2^Ndet(S)}}I_N^{n+2}.`$ For $`N>6`$, i.e. $`det(S)=0`$: $`I_N^n`$ $`=`$ $`{\displaystyle \frac{1}{vKv}}{\displaystyle \underset{l=1}{\overset{N1}{}}}\left((Kv)_l+{\displaystyle \underset{k=1}{\overset{N6}{}}}\beta _kU_l^{(k)}\right)`$ (8) $`\left(I_{N1,N}^nI_{N1,l}^n\right).`$ The reduced integrals $`I_{N1,l}`$ are defined as $`I_{N1,l}^n`$ $`=`$ $`{\displaystyle \frac{d^nk}{i\pi ^{n/2}}\frac{q_l^2}{_{k=1}^Nq_k^2}}`$ $`\text{and }K`$ $`=`$ $`1_{N1}HG,v_l=r_l^2.`$ The $`(N6)\beta _k`$ in (8) are free parameters to construct linear combinations of the vectors $`\{U^{(k)}\}(k=1,\mathrm{},N6)`$ which, together with $`U^{(N5)}`$, form a basis for the kernel of $`G`$. We chose $$U^{(N5)}=Kv/(vKv)$$ parallel to $`v`$ and the others orthogonal, $$vU^{(k=1,\mathrm{},N6)}=0.$$ Obviously one can choose the $`\beta `$’s such that $`(N6)`$ $`b`$’s from the set $`\{b_1,\mathrm{},b_{N1}\}`$ are zero. Doing so one observes that $`I_N^n`$ can be expressed by only 6 $`(N1)`$-point graphs for arbitrary $`N6`$. This is the generalization to massless kinematics of a similar result in the infrared finite case which has been derived for integer dimensions using 4-dimensional Schouten identities. ## 3 Reduction of tensor integrals Tensor integrals can also be split into a part containing reduced integrals ($`K_N^{\mathrm{}}`$) and a part containing higher dimensional objects ($`J_N^{\mathrm{}}`$). Using momentum conservation ($`r_N=0`$) one can write: $`I_N^{\mu _1\mathrm{}\mu _L}`$ $`=`$ $`{\displaystyle \frac{d^nk}{i\pi ^{n/2}}\frac{k^{\mu _1}\mathrm{}k^{\mu _L}}{_{l=1}^Nq_l^2}}`$ $`=`$ $`K_N^{\mu _1\mathrm{}\mu _L}+J_N^{\mu _1\mathrm{}\mu _L}`$ $`K_N^{\mu _1\mathrm{}\mu _L}`$ $`=`$ $`{\displaystyle \underset{l=0}{\overset{L}{}}}{\displaystyle \underset{j_1,\mathrm{},j_l=1}{\overset{N1}{}}}(\{r_j\})_{j_1,\mathrm{},j_l}^{\mu _1\mathrm{}\mu _L}`$ $`\times I_{Nl,Nj_1,\mathrm{},Nj_l}^n`$ $`J_N^{\mu _1\mathrm{}\mu _L}`$ $`=`$ $`{\displaystyle \underset{l=1}{\overset{[L/2]}{}}}{\displaystyle \underset{j_1,\mathrm{},j_{L2l}=1}{\overset{N1}{}}}𝒢(\{r_j\})_{j_1,\mathrm{},j_{L2l}}^{\mu _1\mathrm{}\mu _L}`$ (9) $`\times I_N^{n+2l}(j_1,\mathrm{},j_{L2l})`$ where $`I_{Nl,Nj_1,\mathrm{},Nj_l}^n={\displaystyle \frac{d^nk}{i\pi ^{n/2}}\frac{\underset{m=1}{\overset{l}{}}(q_N^2q_{j_m}^2)}{_{i=1}^Nq_i^2}}`$ are differences of reduced integrals and $$I_N^{n+2l}(j_1,\mathrm{},j_{L2l})$$ are higher dimensional integrals with Feynman parameters $`z_j`$ in the numerator. The Lorentz tensors $`,𝒢`$ present in the above expressions can be written in terms of the objects ($`v_l=r_l^2`$) $`^{\mu \nu }`$ $`=`$ $`r^\mu Hr^\nu `$ $`𝒦_l^\mu `$ $`=`$ $`(r^\mu H)_l`$ $`𝒲^\mu `$ $`=`$ $`𝒦^\mu v`$ (10) Note that the pseudoinverse $`H`$ to the Gram matrix $`G`$ also appears here. The important feature is that one always has $$𝒢g^{\mu \nu }/2r^\mu Hr^\nu .$$ If now the Lorentz vectors $`r_j^\mu (j=1,\mathrm{},N1)`$ span Minkowski space, which is generically the case for $`N5`$, one immediately derives, using the explicit expresssion of $`H`$ given above, that $$r^\mu Hr^\nu =(g^{\mu \nu }/2)^{4\mathrm{dim}}.$$ This means $$g^{\mu \nu }/2r^\mu Hr^\nu =𝒪(ϵ)$$ and thus the higher dimensional objects $`I_N^{n+2l}`$ in Eq. (9) can be dropped for $`N5`$. Consequently it follows from Eq. (9) that for the reduction of an $`N`$-point integral only a small number of higher dimensional integrals is needed, i.e. one has to know $`J_N^{\mu _1\mathrm{}\mu _L}`$ only for $`2LN=2,3,4`$. These can be calculated once and forever, the explicit expressions are $`J_{N=2,3,4}^{\mu _1\mu _2}=\left(g^{\mu _1\mu _2}/2^{\mu _1\mu _2}\right)I_N^{n+2}`$ $`J_{N=3,4}^{\mu _1\mu _2\mu _3}=\left[\left(g/2\right)^{}𝒲^{}\right]^{\{\mu _1\mu _2\mu _3\}}I_N^{n+2}`$ $`\left[\left(g/2\right)^{}𝒦_l^{}\right]^{\{\mu _1\mu _2\mu _3\}}I_{N1,Nl}^{n+2}`$ $`J_4^{\mu _1\mu _2\mu _3\mu _4}=[(g/2)^{}`$ $`\times (g/2)^{}]^{\{\mu _1\mu _2\mu _3\mu _4\}}I_4^{n+4}`$ $`\left[\left(g/2\right)^{}𝒲^{}𝒲^{}\right]^{\{\mu _1\mu _2\mu _3\mu _4\}}I_4^{n+2}`$ $`{\displaystyle \frac{1}{2}}\left[\left(g/2\right)^{}𝒲^{}𝒦_l^{}\right]^{\{\mu _1\mu _2\mu _3\mu _4\}}I_{3,4l}^{n+2}`$ $`\left[\left(g/2\right)^{}_\nu ^{}𝒦_l^{}\right]^{\{\mu _1\mu _2\mu _3\mu _4\}}I_{3,4l}^{\nu ,n+2}.`$ The indexed bracket is an abbreviation for the sum of all distinguishable distributions of Lorentz indices to the objects inside the bracket, e.g. $`\left[X^{}Y^{}Y^{}\right]^{\{\mu _1\mu _2\mu _3\}}=X^{\mu _1}Y^{\mu _2}Y^{\mu _3}`$ $`+X^{\mu _2}Y^{\mu _3}Y^{\mu _1}+X^{\mu _3}Y^{\mu _1}Y^{\mu _2}.`$ For the $`n`$-dimensional part, $`K_N^{\mu _1\mathrm{}\mu _L}`$, a recursion relation holds. Skipping all the details (see ) one finds $`K_N^{\mu _1\mathrm{}\mu _L}`$ $`=`$ $`{\displaystyle \frac{1}{L}}\left[𝒲^{}K_N^{\{L1\mathrm{dots}\}}\right]^{\{\mu _1\mathrm{}\mu _L\}}`$ (11) $`+`$ $`{\displaystyle \frac{2^{(L1)}}{L!}}\left[𝒦_l^{}_{\nu _1}^{}\mathrm{}_{\nu _{L1}}^{}\right]^{\{\mu _1\mathrm{}\mu _L\}}`$ $`\times I_{N1,Nl}^{\nu _1\mathrm{}\nu _{L1}}.`$ Putting everything together, the final recursion formula for $`N`$-point tensor integrals reads $`I_N^{\mu _1\mathrm{}\mu _L}`$ $`=`$ $`{\displaystyle \frac{1}{L}}\left[𝒲^{}(I_NJ_N)^{\{L1\mathrm{dots}\}}\right]^{\{\mu _1\mathrm{}\mu _L\}}`$ (12) $`+{\displaystyle \frac{2^{(L1)}}{L!}}\left[𝒦_l^{}_{\nu _1}^{}\mathrm{}_{\nu _{L1}}^{}\right]^{\{\mu _1\mathrm{}\mu _L\}}`$ $`\times I_{N1,Nl}^{\nu _1\mathrm{}\nu _{L1}}+J_N^{\mu _1\mathrm{}\mu _L}.`$ This form is not immediately suited for iteration, as some reduced integrals do not contain the trivial propagator $`k^2`$ anymore. Therefore one has to shift the loop momenta of the respective integrals by $`kk+r_l`$ before iteration, leading to $`I_{N1,Nl}^{\mu _1\mathrm{}\mu _P}(R)=`$ $`I_{N1}^{\mu _1\mathrm{}\mu _P}(\widehat{R}_{[N]})I_{N1}^{\mu _1\mathrm{}\mu _P}(\widehat{R}_{[l]})`$ $`={\displaystyle \underset{k=0}{\overset{P}{}}}\left[r_l^{(k)}I_{N1}^{\{Pk\mathrm{dots}\}}(R_{[l]})\right]^{\{\mu _1\mathrm{}\mu _P\}}`$ $`I_{N1}^{\mu _1\mathrm{}\mu _P}(\widehat{R}_{[l]}).`$ (13) The argument vectors are $`R`$ $`=`$ $`(r_1,\mathrm{},r_N)`$ $`\widehat{R}_{[k]}`$ $`=`$ $`(r_1,\mathrm{},\widehat{r}_k,\mathrm{},r_N)`$ $`R_{[l]}`$ $`=`$ $`(r_{l+1}r_l,r_{l+2}r_l,\mathrm{},r_{N2+l}r_l,0)`$ where $`\widehat{r}`$ means that the respective vector is missing. The vector indices are understood to be taken cyclically symmetric with periodicity $`N`$, i.e. $`r_{N+l}=r_l`$ for $`l\{1,\mathrm{},N\}`$. As we assume $`r_N=0`$, the integrals on the right-hand side of (3) have again at least one trivial propagator and thus are suited for a further reduction step. As a simple example, consider the rank two 3-point integral $`I_3^{\mu _1\mu _2}(R)`$. After the first reduction step and subsequent shift, one obtains $`I_3^{\mu _1\mu _2}(R)={\displaystyle \frac{1}{2}}\left(𝒲^{\mu _1}I_3^{\mu _2}(R)+𝒲^{\mu _2}I_3^{\mu _1}(R)\right)`$ $`+{\displaystyle \underset{l=1}{\overset{2}{}}}\left(𝒦_l^{\mu _1}_\nu ^{\mu _2}+𝒦_l^{\mu _2}_\nu ^{\mu _1}\right)`$ $`\times \left(r_l^\nu I_2^n(R_{[l]})+I_2^\nu (R_{[l]})I_2^\nu (\widehat{R}_{[l]})\right)`$ $`(g^{\mu _1\mu _2}/2^{\mu _1\mu _2})I_3^{n+2}(R)`$ The recursive structure together with the shift can be implemented in an algebraic manipulation program in a straightforward manner. The recursion stops if all tensor integrals are transformed into expressions containing scalar integrals and Lorentz tensors composed of the objects given in Eq. (10). Starting with an $`N`$-point tensor integral one finds a linear combination of scalar integrals with $`KN`$ legs. To these scalar integrals the reduction formulas (7),(8) can be applied. In this way our reduction formalism solves the problem of calculating massless $`N`$-point Feynman diagrams for arbitrary $`N`$. The generalization to the case where some of the propagators are massive will work similarly. ## 4 Summary/Outlook In conclusion, massless scalar and tensor one-loop integrals for graphs with an arbitrary number of external legs ($`N>4`$) can recursively be reduced to linear combinations of (known) box integrals. To deal with the problem of singular Gram matrices we introduced the concept of the pseudoinverse, which renders the whole formalism mathematically well defined for arbitrary $`N`$. The recursive structure can easily be implemented in algebraic manipulation programs. This allows in principle for the NLO calculation of multi–jet cross sections with a large number of final states. Applications relevant for the near future are e.g. the NLO calculations for $`ppb\overline{b}b\overline{b}`$ or $`pp4jets/photons`$. The most complicated analytical structure in such a calculation is the scalar 6-point function since, as shown above, no dangerous higher-dimensional objects have to be computed. With the given methods, an explicit analytical expression for the scalar 6-point function has been derived and presented in . ## Acknowledgments We would like to thank the organizers for their work. It was a pleasure for us to take part in the Loops and Legs in Quantum Field Theory 2000 conference. This work was supported by the EU Fourth Training Programme ”Training and Mobility of Researchers”, Network ”Quantum Chromodynamics and the Deep Structure of Elementary Particles”, contract FMRX–CT98–0194 (DG 12 - MIHT).
warning/0005/cond-mat0005211.html	ar5iv	text	# Commensurate and incommensurate correlations in Haldane gap antiferromagnets ## I Introduction The $`S=1`$ bilinear-biquadratic chain $$H=\underset{i}{}\left[\mathrm{cos}\theta 𝐒_i𝐒_{i+1}+\mathrm{sin}\theta (𝐒_i𝐒_{i+1})^2\right]$$ (1) is one of the prototype models for the physics of Haldane gap antiferromagnets. Its zero temperature phase diagram has been the subject of intensive studies in recent years.<sup>-</sup> By now it is well-established that the energy gap (Haldane gap) persists in a wide range $`\pi /4<\theta <\pi /4`$ around the conventional Heisenberg point $`\theta =0`$. The model is gapless at the special points $`\theta =\pm \pi /4`$ beyond which other phases with qualitatively different physical properties appear. Although the gap and the hidden (string) order characterizing the Haldane-gap systems persist for the whole Haldane phase $`\pi /4<\theta <\pi /4`$, one can divide this interval into (at least) two, somewhat different subphases. These subphases are separated by the so called valence-bond-solid (VBS) point $`\theta _{\mathrm{vbs}}=\mathrm{tan}^11/30.1024\pi `$ where the ground state properties can be obtained exactly. The two subphases differ in the form of the long distance asymptotics of the two-point correlation function $`S_i^zS_{i+n}^z`$. In the ”commensurate” Haldane phase (C-phase) for $`\pi /4<\theta <\theta _{\mathrm{vbs}}`$ the leading behavior is expected to be $$S_i^zS_{i+n}^z(1)^n\frac{e^{n/\xi }}{n^{1/2}}\text{for }n\mathrm{},$$ (2) while in the ”incommensurate” Haldane phase (IC-phase) for $`\theta _{\mathrm{vbs}}<\theta <\pi /4`$ this was predicted to take the form $$S_i^zS_{i+n}^z\frac{e^{n/\xi }}{n^{1/2}}\mathrm{cos}(qn+\varphi )\text{for }n\mathrm{}$$ (3) where $`q=q(\theta )(\pi ,2\pi /3)`$ is a $`\theta `$-dependent incommensurate wavenumber, $`\varphi `$ is a phase shift, and $`\xi =\xi (\theta )`$ is the correlation length. At the ”C-IC transition point” (also called the ”disorder point”) $`\theta _{\mathrm{cic}}=\theta _{\mathrm{vbs}}`$ the correlation function is known rigorously and it is purely exponential without any algebraic prefactor $$S_i^zS_{i+n}^z=\frac{4}{3}()^ne^{n/\xi }\frac{2}{3}\delta _{n0},$$ (4) with $`\xi (\theta _{\mathrm{vbs}})=1/\mathrm{ln}30.9102`$. Correlation functions of some other operators can also be studied rigorously at the VBS point. In particular, $`(S^z)_i^2(S^z)_{i+n}^2`$ has a similar purely exponential decay with $`\xi =1/\mathrm{ln}3`$ at $`\theta =\theta _{\mathrm{vbs}}`$. The commensurate-incommensurate (C-IC) transition of the spin-1 chain bears much similarity to C-IC transitions found in other models, e.g., in the anisotropic 2D Ising model on the triangular lattice at finite temperature, where it can be analyzed rigorously using the exact solution. Other, higher dimensional, not exactly solvable models with disorder points were investigated using an RPA approach to the susceptibility in Ref. IV. In general, C-IC transitions can be divided into two categories (two kinds). In this classification scheme the transition at the VBS point of the $`S=1`$ chain is a C-IC transition of the first kind, with the property that the incommensurate wavenumber $`q`$ in the IC regime is parameter dependent. For C-IC transitions of the first kind the correlation length is predicted to behave on the C and IC sides of the disorder point $`\theta _{\mathrm{cic}}`$ as $$\frac{d\xi }{d\theta }\|_\mathrm{C}=\mathrm{},\frac{d\xi }{d\theta }\|_{\mathrm{IC}}=\mathrm{finite}$$ (5) with $`\xi (\theta _{\mathrm{cic}})0`$ at the transition point. The characteristic wavenumber is expected to vary on the IC side as $$q(\theta )q(\theta _{\mathrm{cic}})\|\theta \theta _{\mathrm{cic}}\|^{1/2}.$$ (6) Moreover, the RPA theory also predicts the change of the asymptotic form of the correlation functions at the disorder point: the algebraic prefactor is $`n^{(D1)/2}`$ in $`D`$ dimensions except at the disorder point where $`DD^{}`$, $`D^{}<D`$ should be taken, reflecting a ”dimensional reduction”. In our case $`D=1+1=2`$, $`D^{}=1`$ as is seen in Eqs. (2-4). These general features of the C-IC transitions of the first kind have been tested numerically for the spin-1 chain, and except for some observed deviation from the predicted form in Eq. (3) slightly above $`\theta _{\mathrm{vbs}}`$, all were justified. Note that around the VBS point finite size corrections are very small, thus the numerical results (exact diagonalization and DMRG) are extremely precise. Within the IC subphase some other special points can be defined. The first one is $`\theta _{\mathrm{disp}}0.1210\pi `$, where the second derivative of the magnon dispersion at $`k=\pi `$ vanishes, and the dispersion becomes quartic. When the chain is subject to a uniform magnetic field the magnetization-vs-field curve has an anomalous, non-square-root-like singularity at this point. For $`\theta >\theta _{\mathrm{disp}}`$, the second derivative at $`k=\pi `$ is negative and the gap takes its minimum value for a momentum different from $`k=\pi `$. In this region a strong enough field causes the Haldane gap to collapse into a two-component Luttinger liquid (LL) phase instead of a conventional one-component LL phase. The next special point is $`\theta _{\mathrm{max}}0.123\pi `$ where the Haldane gap takes its maximum value. Note that this point has less physical relevance, since it strongly depends on the actual definition of the Hamiltonian; a $`\theta `$-dependent rescaling of the model can easily change its location. Finally one can define the Lifsitz point $`\theta _{\mathrm{Lifs}}0.1314\pi `$, where the incommensurability manifests itself in the structure factor. This is different from the disordered point in massive models due to the finite linewidth of the peaks. Figure 1 shows these special points. In both the C and IC Haldane phases the energy spectrum consists of a discrete triplet branch of ”magnons” separated by finite gaps from the singlet ground state, and also from a continuum of higher lying ”multi-magnon” excitations in a wide momentum range. This one-magnon branch is clearly discernible around the edges of the Brillouin zone $`k=\pm \pi `$. However, it merges into the continuum and vanishes due to magnon scattering processes in an extended range around $`k=0`$. The energy spectrum, especially in the Heisenberg point and the VBS point was studied numerically by many authors. All concluded that the magnon-magnon interactions are rather weak, and bound states do not play a role at low energies. Many properties of the system, from ground-state correlation functions to the onset of magnetization in uniform fields can be extremely well approximated using a massive relativistic free boson theory. Such an approximate theory can be derived directly from the non-linear $`\sigma `$-model (NL$`\sigma `$M) description of the spin chain at the Heisenberg point, or from the Majorana fermion representation of the integrable Takhtajan-Babujian model in the vicinity of $`\theta =\pi /4`$. However, all these microscopic theories utilize the a priori assumption that the important low-energy fluctuations are at momenta $`k=0`$ and $`\pi `$, and thus they cannot account for the C-IC transition, nor can they give any reliable description of the IC regime. Although we do not have a rigorous microscopic justification all results we possess are consistent with the assumption that the elementary excitations are essentially free bosons in the IC regime, too. Thus the aim of the present paper is to extend the free boson description to the whole Haldane phase, and give a general theory which is capable of accounting for the C-IC transition in simple terms. In the lack of a detailed microscopic formulation, however, this theory, at present, is only phenomenological. ## II Effective theory - continuum version A continuum field theory to describe Heisenberg antiferromagnetic chains ($`\theta =0`$) with integer $`S`$ was developed by Haldane. This is the nonlinear $`\sigma `$-model (NL$`\sigma `$M) – without topological terms – which can be derived in the large $`S`$ limit, but whose implications are believed to hold for $`S=1`$ too. The NL$`\sigma `$M is defined by the Lagrangian $$=\frac{1}{2g}\left[\frac{1}{v}(_t\stackrel{}{\varphi })^2v(_x\stackrel{}{\varphi })^2\right],$$ (7) where $`\stackrel{}{\varphi }(x,t)`$ is a vector field with unit length $`\stackrel{}{\varphi }^2=1`$, $`v`$ is a nonuniversal constant (velocity) setting the energy scale, and $`g=2/S`$ is the coupling constant. The associated Hamiltonian is $$=\frac{v}{2}\left[g\stackrel{}{l}^2+\frac{1}{g}(_x\stackrel{}{\varphi })^2\right],$$ (8) where the momentum canonically conjugate to $`\stackrel{}{\varphi }`$ turns out to be $$\stackrel{}{l}(x,t)\frac{1}{gv}\left[\stackrel{}{\varphi }\times \frac{\stackrel{}{\varphi }}{t}\right](x,t).$$ (9) The spin operator $`\stackrel{}{S}_n`$ can be expressed in terms of $`\stackrel{}{\varphi }`$ and $`\stackrel{}{l}`$ as $$\stackrel{}{S}_n(t)=(1)^nS\stackrel{}{\varphi }(x_n,t)+\delta x\stackrel{}{l}(x_n,t)$$ (10) where $`x_n=n\delta x`$, $`\delta x`$ being the lattice constant (usually set to unity), and $`S=1`$ in the present case. At the mean field level the NL$`\sigma `$M can be well approximated by a massive, essentially free vector-boson theory $$=\frac{1}{v}(_t\stackrel{}{\varphi })^2v(_x\stackrel{}{\varphi })^2m^2\stackrel{}{\varphi }^2,$$ (11) now without any constraint on the $`\stackrel{}{\varphi }`$ field. The boson field $`\stackrel{}{\varphi }`$ varies smoothly on the scale of the lattice constant in the commensurate regime (the pure Heisenberg model, which the theory applies for, is in the C phase) and thus higher derivatives, neglected in Eq. (11), are indeed small. The Lagrangian gives rise to a relativistic dispersion $`\omega (k)=\sqrt{m^2+v^2k^2/2m}`$ which takes its minimum at $`k=0`$. \[Note the factor $`(1)^n`$ introduced in Eq. (10) shifting all momenta by $`\pi `$.\] The NL$`\sigma `$M summarized above gives a good description of the low-energy behavior of the $`S=1`$ Heisenberg chain, but is unable to describe the C-IC transition in the bilinear-biquadratic chain. The key feature missing is that in the vicinity of our special points the shape of the magnon dispersion changes drastically, and the minimum at $`k=\pi `$ ($`k=0`$ in the momentum-shifted boson language) splits. Obviously, such an effect can never be obtained from a relativistic field theory such as Eq. (11) and we need to consider a non-relativistic model. The simplest continuum Lagrangian with this property is $``$ $`=`$ $`(_t\stackrel{}{\varphi })^2a(_x\stackrel{}{\varphi })^2b(_x^2\stackrel{}{\varphi })^2m^2\stackrel{}{\varphi }^2.`$ (12) At this level $`a`$, $`b`$, and $`m`$ are $`\theta `$-dependent phenomenological parameters. For stability $`b`$ is supposed to be positive, and $`a`$ is assumed to change sign at the special point $`\theta _{\mathrm{disp}}`$. \[Note that in the relativistic case $`b=0`$, and $`a`$ is in fact $`v^2`$, where $`v`$ is the ”speed of light”. Formally taking $`a<0`$, which we will consider in the sequel, implies an imaginary $`v`$. This only means, however, that one should redefine the physical (quasi)particles to live around the two new minima of the dispersion. In fact, one should introduce two new particles, now with $`v=`$real, for the two minima. This immediately leads to a two-particle (two-band) description which was analyzed in detail in Ref. IV for the massless case (above the critical magnetic field) and was shown to result a two-component (two-band) Luttinger liquid there. Here we shortcut these complications by simply assuming that the two particles are two chiral components ($`k<0`$ and $`k>0`$) of the same boson with $`a<0`$.\] Since our phenomenological model is not derived directly from the microscopic Hamiltonian in Eq. (1) the connection between the boson field and the original spin variables is not known rigorously. In principle $`\stackrel{}{S}_n`$ can be expanded in powers of the field $`\stackrel{}{\varphi }`$ and its derivatives, and any term permitted by the symmetries of the model can appear. Being guided by the NL$`\sigma `$M description at the Heisenberg point, in the following we will assume that the first two terms (linear and quadratic in $`\stackrel{}{\varphi }`$) are $$\stackrel{}{S}_n(t)=g_\varphi (1)^n\stackrel{}{\varphi }(x_n,t)+g_l\stackrel{}{l}(x_n,t)+𝒪(\|\stackrel{}{\varphi }\|^3),$$ (13) where $`g_\varphi `$ and $`g_l`$ are unknown constants. Note that the $`\stackrel{}{l}`$ term, defined by Eq. (9), is the most relevant 2-boson term which is even under parity and odd under time reversal as it should be. Although, we cannot exclude the possibility that up to $`𝒪(\stackrel{}{\varphi }^2)`$ some other terms with higher derivatives also appear, such extra derivatives, due to the finite mass gap, would not modify the long distance asymptotics of $`\stackrel{}{S}_0\stackrel{}{S}_n`$. Equation (13) shows that in the C phase where spin-spin correlations are antiferromagnetic $`\varphi `$ varies smoothly, and higher order derivatives in the Langrangian have minor role. Around the C-CI transition point, and in the IC-phase, however, this is no longer the case: as $`\stackrel{}{S}_n`$ picks up incommensuration, $`\stackrel{}{\varphi }(x)`$ must do so, as well. It no longer varies smoothly, and higher derivatives in the Lagrangian cannot be neglected. This is the effect which we try to take into consideration by the $`b`$ term in Eq. (12). The asymptotic behavior of the correlation function $`S_i^zS_{i+n}^z`$ is determined by the one-boson term; the two-boson term and in general any multi-boson terms only constitute minute corrections which decay at least twice as rapidly. Thus, up to the two-boson term in Eq. (13) the spin-spin correlation function has the behavior $$S_n^z(t)S_0^z(0)=g_\varphi ^2()^nG_\varphi (x_n,t)+g_l^2G_l(x_n,t),$$ (14) with $`G_\varphi (x,t)`$ $``$ $`\varphi ^z(x,t)\varphi ^z(0,0),`$ (15) $`G_l(x,t)`$ $``$ $`l^z(x,t)l^z(0,0).`$ (16) These are the quantities we will now calculate. The Euler-Lagrange equation associated with our Lagrangian gives a generalized Klein-Gordon equation for each component $`\alpha =x,y,z`$ of $`\stackrel{}{\varphi }`$ $$_t^2\varphi ^\alpha a_x^2\varphi ^\alpha +b_x^4\varphi ^\alpha +m^2\varphi ^\alpha =0,$$ (17) whose Green’s function is $$G_\varphi (\omega ,k)=\frac{1}{\omega ^2ak^2bk^4m^2+i\epsilon }.$$ (18) This defines the non-relativistic dispersion $$\omega (k)=\sqrt{m^2+ak^2+bk^4},$$ (19) which reduces to the relativistic dispersion of Ref. IV when $`b=0`$. The generalized theory can be quantized identically to the relativistic Klein-Gordon theory. The field $`\varphi ^\alpha `$, $`\alpha =x,y,z`$, has the following mode expansion $$\varphi ^\alpha (x,t)=\frac{dk}{\sqrt{4\pi \omega (k)}}\left[d_k^\alpha e^{i𝐊𝐗}+d_k^\alpha e^{i𝐊𝐗}\right]$$ (20) where $`𝐊𝐗=\omega tkx`$, and the normalization is $`[d_k^\beta ,d_k^{}^\alpha ]=\delta _{\beta \alpha }\delta (kk^{})`$. Using the mode expansion the equal time expressions $`G_\varphi (x)`$ and $`G_l(x)`$ can be easily reduced to Fourier transforms $`G_\varphi (x)`$ $`=`$ $`{\displaystyle \frac{dk}{4\pi }\frac{e^{ikx}}{\omega (k)}},`$ (21) $`G_l(x)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{dk^{}}{4\pi }\omega (k^{})e^{ik^{}x}\frac{dk}{4\pi }\frac{e^{ikx}}{\omega (k)}}{\displaystyle \frac{1}{2}}\delta ^2(x).`$ (23) Note that the integral determining $`G_\varphi `$ also appears as a multiplicative factor in $`G_l`$. Since the asymptotic behavior is determined by $`G_\varphi `$, we start our analysis with this. The evaluation of the Fourier transform starts with locating the zeros of $`\omega (k)`$ in the complex plane. The four zeros are given by $$k=\pm \left[\frac{1}{2b}\left(a\pm \sqrt{a^24m^2b}\right)\right]^{\frac{1}{2}}$$ (24) and depend on the single parameter $`\theta `$, via $`a`$, $`b`$ and $`m`$. Since our phenomenological model is not derived directly from the microscopic one, we have to make several assumptions about the $`\theta `$-dependence of $`a`$, $`b`$ and $`m`$. We will assume that this dependence is smooth (analytical), $`m(\theta )`$ and $`b(\theta )`$ are nonnegative whereas $`a`$ decreases with increasing $`\theta `$ and changes sign at $`\theta _{\mathrm{disp}}`$. We introduce the discriminant $$D(\theta )=a^24m^2b,$$ (25) which is hence an analytic function of $`\theta `$, and we suppose that it is positive for $`\theta <\theta _{\mathrm{vbs}}`$, vanishes at the VBS point, and it is negative from $`\theta _{\mathrm{vbs}}`$ to $`\theta =\pi /4`$ where it vanishes again. As we show below, under the above hypotheses the phenomenological model provides the expected asymptotic behavior of the correlation function in all parts of the Haldane phase. $`𝜽\mathbf{=}𝜽_{\mathrm{𝐯𝐛𝐬}}`$: At this point $`a>0`$ and $`D=0`$. The expression under the square-root in $`\omega (k)`$ is a complete square, thus $`\omega `$ becomes quadratic in $`k`$ and has two purely imaginary zeros, $`\pm i\sqrt{2m^2/a}=\pm i\sqrt{a/2b}`$ \[see Fig. 2(b)\]. The contour of integration can be closed in the upper half plane and the result is $$G_\varphi (x)=\frac{1}{2\sqrt{2a}}e^{\sqrt{\frac{a}{2b}}\|x\|}.$$ (26) Because of the purely exponential decay, this indeed corresponds to $`\theta =\theta _{\mathrm{vbs}}`$. $`𝜽\mathbf{<}𝜽_{\mathrm{𝐯𝐛𝐬}}`$: In this region $`a>0`$ and $`D>0`$. We may suppose $`b>0`$; the case when $`b=0`$ can be obtained by continuity. We get four purely imaginary zeros $`\pm iv_\pm `$ where $$v_\pm =\left[\frac{1}{2b}\left(a\pm \sqrt{D}\right)\right]^{\frac{1}{2}}.$$ (27) Now $$\omega (k)=\sqrt{b}(k^2+v_+^2)^{\frac{1}{2}}(k^2+v_{}^2)^{\frac{1}{2}}$$ (28) is single-valued in the complex plane with two cuts, one between $`iv_+`$ and $`iv_{}`$ and another one between $`iv_{}`$ and $`iv_+`$ \[see Fig. 2(a)\]. (We use the convention that $`z^{1/2}`$ has a cut along the negative real axis.) The contour of integration can be closed in the upper half-plane and drawn onto the upper cut, giving $$G_\varphi (x)=\frac{1}{2\pi \sqrt{b}}_v_{}^{v_+}\frac{e^{t\|x\|}}{\left[(t^2v_{}^2)(v_+^2t^2)\right]^{\frac{1}{2}}}𝑑t.$$ (29) This function is positive for all $`x`$, so with the factor $`(1)^n`$ introduced in Eq. (13) we obtain the expected antiferromagnetic modulation. As $`\|x\|`$ goes to infinity, the main contribution to the integral is coming from the vicinity of the end point $`v_{}`$, and it is legitimate to expand the integrand around this point. For $`\|x\|>\|D\|^{1/2}`$ this yields $$G_\varphi (x)=\frac{e^{v_{}\|x\|}}{2[2\pi v_{}\sqrt{D}]^{\frac{1}{2}}}\left[\|x\|^{\frac{1}{2}}+O(\|x\|^{\frac{3}{2}})\right].$$ (30) Together with Eq. (13), we find the expected asymptotic form Eq. (2) of the spin-spin correlation function with a correlation length $`\xi =1/v_{}`$ which is continuous at $`\theta =\theta _{\mathrm{vbs}}`$. $`𝜽\mathbf{>}𝜽_{\mathrm{𝐯𝐛𝐬}}`$: Here D becomes negative and $`\omega `$ has four complex zeros $`k_j=\pm u\pm iv`$ ($`j=1,\mathrm{},4`$) with $$u=\left[\frac{m}{2\sqrt{b}}\frac{a}{4b}\right]^{\frac{1}{2}}v=\left[\frac{m}{2\sqrt{b}}+\frac{a}{4b}\right]^{\frac{1}{2}}.$$ (31) If we write $`\omega `$ in the form $$\omega (k)=\sqrt{b}\underset{j=1}{\overset{4}{}}(kk_j)^{\frac{1}{2}}$$ (32) we see that it is single-valued on the complex plane with a cut between $`uiv`$ and $`uiv`$ and another cut between $`u+iv`$ and $`u+iv`$ \[see Fig. 2(c)\]. For $`k`$ real we get back the original positive function. The integration can be carried out along a contour which starts at $`u+i\mathrm{}`$, goes vertically down to $`u+iv`$, passes below the upper cut and goes vertically to $`u+i\mathrm{}`$. Next, we replace the integral below the cut by an integral going above the cut and in the opposite sense, and this latter by the sum of the two integrals along the vertical half-lines. These are complex conjugate to each other, so finally we obtain $`G_\varphi `$ $`(x)=`$ (34) $`{\displaystyle \frac{1}{\pi \sqrt{b}}}{\displaystyle _v^{\mathrm{}}}{\displaystyle \frac{e^{\|x\|t}\mathrm{cos}[u\|x\|\phi (t)]dt}{\sqrt{t^2v^2}[4u^2+(tv)^2]^{\frac{1}{4}}[4u^2+(t+v)^2]^{\frac{1}{4}}}},`$ where $$\phi (t)=\frac{1}{2}\left(\mathrm{arctan}\frac{tv}{2u}+\mathrm{arctan}\frac{t+v}{2u}\right).$$ (35) Now $`G_\varphi (x)`$ changes sign periodically, and we can identify $`\pi u`$ with the wave number $`q`$ of the incommensurate oscillation. At the VBS point $`u=0`$, and as it is shown by Eq. (31), the assumed analyticity of $`m`$, $`a`$ and $`b`$ assures that it has a square-root-type singularity above the VBS point in accordance with Eq. (6). The large-$`\|x\|`$ asymptotics of $`G_\varphi (x)`$ can be obtained by Watson’s lemma or by a direct expansion, $`G_\varphi `$ $`(x)=`$ (37) $`{\displaystyle \frac{e^{v\|x\|}\mathrm{cos}(u\|x\|\frac{1}{2}\mathrm{arctan}\frac{v}{u})}{(2\pi )^{\frac{1}{2}}\left(\frac{mD}{\sqrt{b}}\right)^{\frac{1}{4}}}}\left[\|x\|^{\frac{1}{2}}+O(\|x\|^{\frac{3}{2}})\right].`$ Formulas (34) and (37) apply for any $`\theta `$ between $`\theta _{\mathrm{vbs}}`$ and $`\pi /4`$, including $`\theta _{\mathrm{disp}}`$ which is a symmetry point of the domain of incommensurate oscillations ($`u=v`$). As $`\theta `$ approaches $`\pi /4`$, $`D`$ goes to zero and the zeros of $`\omega (k)`$ tend to the real axis. Thus, $`v`$ goes to zero and the correlation length diverges. This is what we expect at the boundary of the Haldane phase $`\theta =\pi /4`$ where the gap disappears. At this special point the ground state has a tripled periodicity, implying $`u(\theta =\pi /4)=2\pi /3`$. For any fixed $`x`$, $`G_\varphi (x)`$ depends analytically on $`\theta `$ inside the whole Haldane phase. This can be seen from the original form Eq. (LABEL:phiphi) of $`G_\varphi (x)`$ by inserting the original, non-factorized expression Eq. (19) for $`\omega (k)`$. A proof can be found in \[IV\]. The argument makes use of the continuity of the integrand in $`k`$ real, its (supposed) analyticity for any fixed $`k`$ as a function of $`\theta `$ in suitable complex domains, and the uniform convergence of the integral for $`\theta `$ in any of these domains. It is interesting to examine another kind of asymptotics, valid in a close neighborhood of $`\theta _{\mathrm{vbs}}`$, when $`\|D\|/a^21`$. In this case Eq. (29) reduces to the form $`G_\varphi (x)`$ $``$ $`{\displaystyle \frac{e^{\sqrt{\frac{a}{2b}}\|x\|}}{\sqrt{8a}}}{\displaystyle \frac{1}{\pi }}{\displaystyle _0^\pi }\mathrm{cosh}\left[\sqrt{{\displaystyle \frac{D}{8ab}}}\|x\|\mathrm{sin}\alpha \right]𝑑\alpha `$ (38) $``$ $`{\displaystyle \frac{e^{\sqrt{\frac{a}{2b}}\|x\|}}{\sqrt{8a}}}I_0\left(\sqrt{{\displaystyle \frac{D}{8ab}}}\|x\|\right)`$ (39) where $`I_0`$ is the zeroth order Bessel function. We can arrive at the same equation from Eq. (34), by changing the contour of integration (integrating around the upper cut). Now analyticity at the VBS point is manifest, because in the expansion of the hyperbolic cosine about zero only the even powers of $`\sqrt{D}`$ appear. Equation (39) shows that the crossover to the decay with the $`\|x\|^{\frac{1}{2}}`$ prefactor sets in at the characteristic distance $`\|x\|D^{\frac{1}{2}}`$, which diverges at the VBS point. This explains the numerical difficulties verifying the expected asymptotic behavior very close to the VBS point. At the VBS point there is an infinite jump in the derivative of the correlation length, as predicted by Eq. (5). Indeed, Eq. (27) yields $$\frac{d\xi }{d\theta }\|_{\theta _{\mathrm{vbs}}0}\frac{D^{}}{4\sqrt{2a}m}\|_{\theta _{\mathrm{vbs}}}D^{\frac{1}{2}}(\theta _{\mathrm{vbs}}\theta )^{\frac{1}{2}}$$ (40) because $`D^{}(\theta _{\mathrm{vbs}})<0`$. On the other hand, from Eq. (31) $$\frac{d\xi }{d\theta }\|_{\theta _{\mathrm{vbs}}+0}=\frac{1}{2\sqrt{a}}\left(\frac{m^{}}{2m}+\frac{a^{}}{2a}+\frac{m^2}{a^2}\right)\|_{\theta _{\mathrm{vbs}}},$$ (41) which is finite. The singularity of the correlation length at $`\theta _{\mathrm{vbs}}`$ is in no contradiction with the analyticity of $`G(x)`$ at a fixed $`x`$. Indeed, the divergence of the derivative of $`\xi `$ was extracted from the single-exponential asymptotic form Eq. (30) which, again, is valid only for $`\|x\|>(v_+v_{})^1D^{\frac{1}{2}}`$. To see the role of the two-boson term $`G_l`$ in the correlation functions, we can use the identity (after proper regularization) $$\frac{dk}{4\pi }\omega (k)e^{ikx}=\omega ^2\left(i\frac{}{x}\right)\frac{dk}{4\pi }\frac{e^{ikx}}{\omega (k)}.$$ (42) where $`\omega ^2(i/x)`$ is a shorthand for $`m^2a^2/x^2+b^4/x^4`$ in the present case. With this $$G_l(x)=G_\varphi (x)\omega ^2\left(i\frac{}{x}\right)G_\varphi (x),$$ (43) where we have neglected the singular, delta-function term of Eq. (LABEL:phiphi). This term can be evaluated directly for large $`x`$, knowing the asymptotic form of $`G_\varphi `$ in the different regimes. Using Eqs. (26), (30) and (37) we obtain $`G_l(x)\{\begin{array}{cc}{\displaystyle \frac{e^{2v_{}x}}{x^2}}\hfill & \text{if }\theta <\theta _{\mathrm{vbs}}\hfill \\ 0\hfill & \text{if }\theta =\theta _{\mathrm{vbs}}\hfill \\ {\displaystyle \frac{e^{2vx}}{x^2}}[c_1+c_2\mathrm{cos}(2ux+\alpha )]\hfill & \text{if }\theta >\theta _{\mathrm{vbs}},\hfill \end{array}`$ (47) where the constants $`c_1,c_2`$ and $`\alpha `$ can be expressed straightforwardly with $`a,b`$ and $`m`$. It is interesting to remark that $`G_l`$ is exactly zero for any $`x>0`$ at the disorder point, i.e., the 2-boson processes do not contribute to the equal time correlation function there. This can be easily verified by calculating $`\omega ^2(i/x)G_\varphi (x)`$ using Eq. (26) and the fact that $`D=0`$ at the VBS point. Here $`S_n^zS_0^z`$ only contains the $`G_\varphi `$ term in full accordance with the exact solution. For other values of $`\theta `$, $`G_l`$ decays twice as rapidly as $`G_\varphi `$. It is also of interest to see what predictions our simple field theory gives for the analytic properties of the ground state energy density and the gap. Using the mode expansion in Eq. (20), the Hamiltonian can be written as $$H=_\mathrm{\Lambda }^\mathrm{\Lambda }𝑑k\omega (k)\left[d_k^{}d_k+\frac{1}{2}\right],$$ (48) where $`\mathrm{\Lambda }`$ is an appropriate UV momentum cutoff, proportional to the inverse of the lattice constant. From this the ground state energy is $$E=\frac{1}{2}_\mathrm{\Lambda }^\mathrm{\Lambda }𝑑k\omega (k).$$ (49) The ground state energy depends analytically on $`a`$, $`b`$ and $`m`$ whenever the zeros of $`\omega (k)`$ are not on the real axis. Together with the supposed analyticity of $`a(\theta )`$, etc., this means that $`E(\theta )`$ is also analytic inside the whole Haldane phase. A straightforward expansion around $`\theta _{\mathrm{vbs}}`$ yields $`\underset{\mathrm{\Lambda }\mathrm{}}{lim}`$ $`\left[E(\theta )E(\theta _{\mathrm{vbs}})\right]=`$ (51) $`{\displaystyle \frac{8\pi a^{3/2}}{\sqrt{2}b}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{2^{6n}(4n4)!}{(2n2)!n!(n1)!}}\left({\displaystyle \frac{D}{a^2}}\right)^n,`$ which is convergent if $`\|D\|/a^2<1`$. We notice that in general $`E(\theta )`$ can be expressed in a closed form in terms of elliptic integrals of the first and second kind. The energy gap of the model is by definition $$\mathrm{\Delta }=\mathrm{min}_k\omega (k).$$ (52) This is obviously analytic at $`\theta _{\mathrm{cic}}=\theta _{\mathrm{vbs}}`$ but has a singularity at $`\theta _{\mathrm{disp}}`$, where the minima of $`\omega (k)`$ move away from $`k=0`$ as the parameter $`a=0`$ changes sign. While for $`\theta \theta _{\mathrm{disp}}`$ the minimum is taken at $`k=0`$, for $`\theta >\theta _{\mathrm{disp}}`$ it is taken at $`k=\pm \sqrt{a/2b}`$. The gap $`\mathrm{\Delta }`$ and its derivatives with respect to $`\theta `$ on the two sides of $`\theta _{\mathrm{disp}}`$ turn out ot be, resp., $`\mathrm{\Delta }(\theta _{\mathrm{disp}}0)`$ $`=`$ $`m,\mathrm{\Delta }(\theta _{\mathrm{disp}}+0)=m,`$ (53) $`\mathrm{\Delta }^{}(\theta _{\mathrm{disp}}0)`$ $`=`$ $`m^{},\mathrm{\Delta }^{}(\theta _{\mathrm{disp}}+0)=m^{},`$ (55) $`\mathrm{\Delta }^{\prime \prime }(\theta _{\mathrm{disp}}0)`$ $`=`$ $`m^{\prime \prime },\mathrm{\Delta }^{\prime \prime }(\theta _{\mathrm{disp}}+0)=m^{\prime \prime }{\displaystyle \frac{a^2}{4bm}}\|_{\theta _{\mathrm{disp}}}.`$ (56) We see that there is a discontinuity in the second derivative. This behavior of the effective theory seems consistent with the numerical results shown in Fig. 9 of Ref. \[IV\]. ## III Effective theory - lattice version The effective theory presented above is capable of providing a complete qualitative description of the C-IC transition of the spin-1 bilinear-biquadratic model in accordance with the available numerical data. However, in order to give quantitative predictions, too, the theory needs some refinement. We have seen earlier that our simple theory identifies the disorder point $`\theta _{\mathrm{cic}}`$ with the point where the discriminant $`D`$ defined by Eq. (25) vanishes. It is easy to verify that the second and fourth derivatives of $`\omega (k)`$ at $`k=0`$ are, resp., $$\frac{^2\omega }{k^2}(k=0)=\frac{a}{m},\frac{^4\omega }{k^4}(k=0)=\frac{3D}{m^3},$$ (57) thus in the above theory the fourth derivative vanishes at the disorder point. In contrast with this, Golinelli et al. measured numerically the second and fourth derivatives at the known disorder point $`\theta _{\mathrm{vbs}}`$ and found $$\frac{^2\omega }{k^2}(k=0)=0.9778(1),\frac{^4\omega }{k^4}(k=0)=1.202(1).$$ (58) Note that the fourth derivative is only zero far inside the IC regime, which seems inconsistent with the above theory. One step to improve the theory is to realize that the model is defined on a lattice, and thus the dispersion $`\omega (k)`$ must be a $`2\pi `$-periodic function of $`k`$ (from now on the lattice constant is set $`\delta x=1`$). This can be incorporated into the Lagrangian in Eq. (12) by the standard replacement $`_x\varphi [\varphi (n+\delta x)\varphi (n)]/\delta x`$, leading to the substitution $`k^22[1\mathrm{cos}(k)]`$ in the dispersion $`\omega (k)`$ $`=`$ (60) $`\sqrt{m^2+2a+6b(2a+8b)\mathrm{cos}(k)+2b\mathrm{cos}(2k)}.`$ The Green’s function now reads $$G_\varphi (n)=_\pi ^\pi \frac{dk}{4\pi }\frac{1}{\omega (k)}e^{ikn}.$$ (61) The condition that the expression under the square root in $`\omega (k)`$ is a complete square is again $`D(\theta )=0`$ with $`D`$ defined in Eq. (25). When this is satisfied the dispersion simplifies to $$\omega _{\mathrm{vbs}}(k)=m_{\mathrm{vbs}}+\frac{a_{\mathrm{vbs}}}{m_{\mathrm{vbs}}}[1\mathrm{cos}(k)],$$ (62) and $`1/\omega _{\mathrm{vbs}}(k)`$ has poles instead of branch cuts. \[One pole within the Brillouin zone $`\pi <\mathrm{Re}(k)\pi `$ with $`0<\mathrm{Im}(k)`$.\] Now the second and fourth derivatives at $`k=0`$ are \[cf. Eq. (57)\] $$\frac{^2\omega }{k^2}(k=0)=\frac{a}{m},\frac{^4\omega }{k^4}(k=0)=\frac{3D}{m^3}\frac{a}{m}.$$ (63) At the disorder point we find $$\frac{^4\omega }{k^4}(k=0)=\frac{a_{\mathrm{vbs}}}{m_{\mathrm{vbs}}},$$ (64) which is nonzero. The correlation length $`\xi =1/\mathrm{Im}(k)`$ at the disordered point is determined by the position of the pole, i.e., by the solution of the transcendental equation $$1+\frac{a_{\mathrm{vbs}}}{m_{\mathrm{vbs}}^2}[1\mathrm{cos}(k)]=0.$$ (65) Working the other way around, knowing that at the VBS (C-IC) point $`\xi _{\mathrm{vbs}}=1/\mathrm{ln}3`$, Eq. (65) gives $`a_{\mathrm{vbs}}=3m_{\mathrm{vbs}}^2/2`$ and thus the dispersion $$\omega _{\mathrm{vbs}}=m_{\mathrm{vbs}}\left[\frac{5}{2}\frac{3}{2}\mathrm{cos}(k)\right].$$ (66) With this expression the lattice Green’s function defined in Eq. (61) reads $$G_\varphi (n)=\frac{1}{4m_{\mathrm{vbs}}}e^{\|n\|\mathrm{ln}3},$$ (67) which should be conferred to Eq. (4). The functional form of the dispersion relation in Eq. (66) is exactly the same as the one appearing in the single mode approximation of the VBS model. There, one derives an upper bound for the gap $`\mathrm{\Delta }_{\mathrm{vbs}}=m_{\mathrm{vbs}}=\sqrt{40}/9`$, whereas here we should use the phenomenological (numerical) value $`\mathrm{\Delta }_{\mathrm{vbs}}=0.664314`$ in Eq. (66). This, together with Eq. (64) yields the value $`^4\omega /k^4=3m_{\mathrm{vbs}}/21.0`$, which is rather close to the numerical estimate in Eq. (58). The split of the double root in the vicinity of the disorder point can be analyzed similarly to the continuum theory. We do not go into details here, but emphasize that the critical exponents characterizing the behavior of $`\xi (\theta )`$ and $`q(\theta )`$ at $`\theta _{\mathrm{vbs}}`$, and the type of the singularity of the gap at $`\theta _{\mathrm{disp}}`$ remain the same. Similarly, we find that the ground state energy is analytic everywhere. One can wonder about the possible consequences of keeping higher order terms in the continuum Lagrangian Eq. (12) or in the improved lattice version. If the theory remains free the only effect is to bring about additional branch cuts or poles in $`1/\omega (k)`$. If the higher order terms are small the additional branching points are far in $`k`$ space, and the C-IC transition remains intact. The long distance asymptotics of the correlation functions do not change. Since at the VBS point the exact correlation function in Eq. (4) only contains a single exponential term, the free boson approximation does not allow any higher order spatial derivatives in the effective Lagrangian there. Although the dispersion $`\omega _{\mathrm{vbs}}`$ in Eq. (66) gives rise to the exponential term in Eq. (4), it misses the $`\delta _{n,0}`$ contribution. This is, however, another artifact of the continuum approach, which necessarily neglects some important short distance details. In fact, we should recall that in a quantum spin liquid with short-range valence bond ground state the elementary excitations (bosons) are physically triplet bonds living between lattice sites, rather than on the sites themselves. As was argued in Ref. \[IV\], $`S_i^z`$ acting on the VBS ground state produces the linear combination of two states, one of which containing a boson at site $`i1/2`$, the other a boson at site $`i+1/2`$. Hence the one-boson term of $`S_iS_{i+n}`$ is in fact $$S_n^zS_0^z=g_\varphi ^2()^n[G_\varphi (n1)+2G_\varphi (n)+G_\varphi (n+1)].$$ (68) Using Eq. (67) this leads directly to Eq. (4), including the $`\delta `$-function piece, if $`g_\varphi ^2=m_{\mathrm{vbs}}`$. ## IV Summary and discussion In summary, we proposed a simple effective field theory to describe the commensurate-incommensurate transition in the Haldane phase of the spin-1 bilinear-biquadratic chain. The theory is capable of reproducing many features of this transition previously seen in the numerical studies. Moreover it also has some new predictions. The effective theory predicts that the C-IC transition at $`\theta _{\mathrm{vbs}}`$ is not a phase transition in the conventional sense, since the ground state energy remains an analytic function of the control parameter $`\theta `$. The only singularity occurring is in the correlation length. We should emphasize, however, that the correlation function itself remains analytic as a function of $`\theta `$ for any fixed distance, unlike in conventional phase transitions. There is another point $`\theta _{\mathrm{disp}}`$ close to the disorder point where another quantity becomes singular. This is the energy gap (Haldane gap) whose second derivative produces a jump. This singularity in the singlet-triplet gap becomes important when a high enough magnetic field is applied, producing a crossing between these levels, and thus leading to the collapse of the gap. At the critical field, as $`\theta `$ is varied a real phase transition takes place at $`\theta _{\mathrm{disp}}`$, thus this point is the endpoint of a phase transition line on the magnetic field vs $`\theta `$ plain separating two Luttinger liquid type phases. In a technical sense the C-IC transition is a consequence of an accidental degeneracy of roots of the dispersion relation. We have shown in the free boson approach that this degeneracy causes the ”dimensional reduction” of the correlation function, and makes it to be a pure exponential at the VBS point. We also found that the two-magnon contributions, and presumably any higher order, multi-magnon contributions too, vanish exactly at this point. We derived a formula valid in the vicinity of the VBS point showing how the pure exponential decay emerges from the standard form with algebraic prefactors. In particular we found that there is a crossover between a pure exponential decay and a decay containing algebraic prefactors. The characteristic distance of this crossover tends to infinity as the VBS point is approached. The spin-1 bilinear-biquadratic model studied in this paper is not the only model which produces a C-IC transition. Another example is the spin-1/2 chain with next-nearest-neighbor interaction (this can also be visualized as a two-leg zig-zag ladder). By now it is well established that the Majumdar-Ghosh point of this model is a disorder point where a C-IC transition of the first kind occurs. On the same footing as described here, it seems possible to develop an effective theory which supposes that the elementary excitations (spin-1/2 solitons in that case) are essentially free particles with a non-relativistic dispersion. Care should, however, be taken on the facts that solitons are always created in pairs and that they are spin-1/2 particles. Beside the $`S=1/2`$ case, the appearence of disorder points has been demonstrated in other $`S1`$ frustrated Heisenberg chains too. Another interesting quasi-one-dimensional system where a commensurate-incommensurate transition have been reported in a numerical investigation is the SU(2)$`\times `$SU(2) symmetric coupled spin-orbit model. The elaboration and testing of effective theories, similar to the one described in this paper, for these models could be a possible direction of future research. We thank L. Balents and J. Sólyom for valuable discussions. This work was financially supported by the Hungarian Scientific Research Found (OTKA) under grant Nos. 30173, 30543, and F31949.
warning/0005/math0005074.html	ar5iv	text	# LOCALLY SASAKIAN MANIFOLDS 11footnote 1Research supported by Komitet Badań Naukowych (Grant nr 2 P03B 017 12) Michał Godliński, Wojciech Kopczyński22footnote 2e-mail: kopcz@fuw.edu.pl , Paweł Nurowski33footnote 3e-mail: nurowski@fuw.edu.pl Instytut Fizyki Teoretycznej Uniwersytet Warszawski, ul. Hoża 69, Warszawa, Poland ## 1 Introduction The Sasakian structure, wich is defined on an odd dimensional manifold is, in a sense, the closest possible analog of the Kähler geometry of even dimension. It was introduced by S. Sasaki in 1960, who considered it as a special kind of contact geometries. Sasakian structure consists, in paraticular, of the contact 1-form $`\eta `$ and the Riemannian metric $`g`$. The differential of $`\eta `$ defines a 2-form, which constitutes an analog of the fundamental form of Kähler geometry. Sasakian geometry was primarily studied as a substructure within the catagory of contact structures. A review of this approach can be found in . In this letter we exploit the analogy between Sasakian and Kähler geometry. We show that a well known fact that a Kähler geometry can be locally generated by a Kähler potential has its Sasakian counterpart. This result may be of some use in constructing a vast family of examples of Sasakian and Sasakian-Einstein structures. The Sasakian and Sasakian-Einstein structures appear in physics in the context of string theory. It turns out that a metric cone $`(C(𝒮)=𝐑_+\times 𝒮,\overline{g}=\mathrm{d}r^2+r^2g)`$ over a Sasakian-Einstein manifold $`(𝒮,g)`$ is Kähler and Ricci flat, i.e. it constitutes a Calabi-Yau manifold. Moreover, the Sasaki-Einstein manifolds in dimensions $`2k+1`$ and Sasakian manifolds with three Sasakian structures in dimesnion $`4k+3`$ are related to the Maldacena conjecture . It turns out that they are one of very few structures which can serve as a compact factor $`𝒮`$ in (anti-de-Sitter)$`\times 𝒮`$ background for classical field theories which, via the Maldacena conjecture, correspond to the large $`N`$ limit of certain quantum conformal field theories. A formal definition of a Sasaki manifold is as follows. Definition 1 Let $`𝒮`$ be a $`(2k+1)`$-dimensional manifold equipped with a structure $`(\varphi ,\xi ,\eta ,g)`$ such that: (i) $`\varphi `$ is a (1,1) tensor field, (ii) $`\xi `$ is a vector field, (iii) $`\eta `$ is a field of an 1-form, (iv) $`g`$ is a Riemannian metric. Assume, in addition, that for any vector fields $`X`$ and $`Y`$ on $`𝒮`$, $`(\varphi ,\xi ,\eta ,g)`$ satisfy the following algebraic conditions: * $`\varphi ^2X=X+\eta (X)\xi `$, * $`\eta (\xi )=1`$, * $`g(\varphi X,\varphi Y)=g(X,Y)\eta (X)\eta (Y)`$, * $`g(\xi ,X)=\eta (X)`$, and the folllowing differential conditions: * $`N_\varphi +\mathrm{d}\eta \xi =0`$, where $`N_\varphi (X,Y)=[\varphi X,\varphi Y]+\varphi ^2[X,Y]\varphi [\varphi X,Y]\varphi [X,\varphi Y]`$ is the Nijenhuis tensor for $`\varphi `$, * $`\mathrm{d}\eta (X,Y)=g(\varphi X,Y)`$. Then $`𝒮`$ is called a Sasakian manifold. Example 1 A standard example of a Sasaki manifold is an odd dimensional sphere $$𝐒^{2k+1}=\{𝐂^{k+1}(z^1,\mathrm{},z^{k+1}):\|z^1\|^2+\mathrm{}+\|z^{k+1}\|^2=1\}𝐂^{k+1},$$ viewed as a submanifold of $`𝐂^{k+1}`$. Let $`J`$ be the standard complex structure on $`𝐂^{k+1}`$, $`\stackrel{~}{g}`$ the standard flat metric on $`𝐂^{k+1}𝐑^{2k+2}`$, and $`n`$ be the unit normal to the sphere. The vector field $`\xi `$ on $`𝐒^{2k+1}`$ is defined by $`\xi =Jn`$. If $`X`$ is tangent vector to the sphere then $`JX`$ uniquely decomposes onto the part paralell to $`n`$ and the part tangent to the sphere. Denote this decomposition by $`JX=\eta (X)n+\varphi X`$. This defines 1-form $`\eta `$ and tensor field $`\varphi `$ on $`𝐒^{2k+1}`$. Denoting the restriction of $`\stackrel{~}{g}`$ to $`𝐒^{2k+1}`$ by $`g`$ we obtain $`(\varphi ,\xi ,\eta ,g)`$ structure on $`𝐒^{2k+1}`$. It is a matter of checking that this structure equips $`𝐒^{2k+1}`$ with a structure of a Sasakian-Einstein manifold. This construction is, in a certain sense, a Sasakian counterpart of the Fubini-Study Kähler structure on $`\mathrm{𝐂𝐏}^k`$. Notation We adapt the following notation: $`K_{,l}`$ denotes the partial derivative of a function $`K`$ with respect to the coordinate $`z^l`$. Complex conjugate of an indexed quantity, e.g. $`a_j^i`$, is usually denoted by a bar over it, i.e. $`\overline{a_j^i}`$. Our notation is: $`\overline{a_j^i}=\overline{a}_{\overline{j}}^{\overline{i}}`$. Symmetrized tensor products of 1-forms $`\eta `$ and $`\lambda `$ is denoted by $`\eta \lambda =\frac{1}{2}(\eta \lambda +\lambda \eta )`$. The main result The pouropuse of this letter is to prove the following theorem, which locally characterizes all Sasakian and Sasakian-Einstein manifolds. Theorem. Let $`𝒰`$ be an open set of $`𝐂^k\times 𝐑`$ and let $`(z^1,z^2,\mathrm{},z^k,x)`$ be Cartesian coordinates in $`𝒰`$. Consider: * a vector field $`\xi =_x`$ * a real-valued function $`K`$ on $`𝒰`$ such that $`\xi (K)=0`$ * an 1-form $$\eta =\mathrm{d}x+i\underset{m=1}{\overset{k}{}}(K_{,m}\mathrm{d}z^m)i\underset{\overline{m}=1}{\overset{k}{}}(K_{,\overline{m}}\mathrm{d}\overline{z}^{\overline{m}})$$ * a bilinear form $$g=\eta ^2+2\underset{m,\overline{k}=1}{\overset{k}{}}K_{,m\overline{k}}\mathrm{d}z^m\mathrm{d}\overline{z}^{\overline{k}}$$ * a tensor field $$\varphi =i\underset{m=1}{\overset{k}{}}[(_miK_{,m}_x)\mathrm{d}z^m]+i\underset{\overline{m}=1}{\overset{k}{}}(_{\overline{m}}+iK_{,\overline{m}}_x)\mathrm{d}\overline{z}^{\overline{m}}].$$ I) If the function $`K`$ is choosen in such a way that the bilinear form $`g`$ has positive definite signature then $`𝒰`$ equipped with the structure $`(\varphi ,\xi ,\eta ,g)`$ is a Sasakian manifold. Moreover, every Sasakian manifold can locally be generated by such a function $`K`$. II) The above Sasakian structure satisfies Einstein equation $`Ric(g)=\lambda g`$ if and only if $`\lambda =2k`$ and the function $`K`$ satsifies $$[\mathrm{log}det(K,_{i\overline{j}})],_{m\overline{n}}=2(k+1)K,_{m\overline{n}}.$$ ## 2 Almost conatact versus Sasakian manifolds Definition 2 Consider $`(2k+1)`$-dimensional manifold $`𝒮`$ equipped with a structure consisting of a (1,1) tensor field $`\varphi `$, a vector field $`\xi `$ and a field of an 1-form $`\eta `$. Assume, in addition, that for every vector field $`X`$ on $`𝒮`$ $`(\xi ,\eta ,\varphi )`$ satisfy the following algebraic conditions: * $`\varphi ^2X=X+\eta (X)\xi `$, * $`\eta (\xi )=1`$. Then $`𝒮`$ is called almost contact manifold. If, in addition, an almost contact manifold $`(𝒮,(\xi ,\eta ,\varphi ))`$ is equipped with Riemannian metric $`g`$ such that for every veector fields $`X`$ and $`Y`$ on $`𝒮`$ we have * $`g(\varphi X,\varphi Y)=g(X,Y)\eta (X)\eta (Y)`$, * $`g(\xi ,X)=\eta (X)`$, then the almost contact manifold is called almost contact metric manifold. Note that every Sasakian manifold is an almost contact metric manifold. Let T$`{}_{}{}^{𝐂}𝒮`$ be the complexification of the tangent bundle of an almost contact manifold $`𝒮`$. The almost contact structure $`(\xi ,\eta ,\varphi )`$ on $`𝒮`$ defines the decomposition $$\mathrm{T}^𝐂𝒮=𝐂\xi N\overline{N},$$ where $`𝐂\xi `$, $`N`$ and $`\overline{N}`$ are eigenspaces of $`\varphi `$ with eigenvalues 0, $`i`$ and $`i`$, respectively. We say that a vector subbundle $`Z`$ of T$`{}_{}{}^{𝐂}𝒮`$ is involutive if and only if $`[\mathrm{\Gamma }(Z),\mathrm{\Gamma }(Z)]\mathrm{\Gamma }(Z)`$, where $`\mathrm{\Gamma }(Z)`$ denotes the set of all sections of $`Z`$. Lemma For an almost contact structure the condition $`N_\varphi +\mathrm{d}\eta \xi =0`$ is satisfied if and only if the bundle $`N`$ is involutive, $`[\mathrm{\Gamma }(N),\mathrm{\Gamma }(N)]\mathrm{\Gamma }(N)`$, and $`[\xi ,\mathrm{\Gamma }(N)]\mathrm{\Gamma }(N)`$. Proof. Let $`X,Y\mathrm{\Gamma }(N)`$. Making use of the eigenvalue property of $`\varphi `$ and of property (1) of Definition 2, we get the following expressions for the Nijenhuis tensor of $`\varphi `$: $$N_\varphi (X,Y)=2[X,Y]+2i\varphi ([X,Y])+\eta ([X,Y])\xi ,$$ $$N_\varphi (X,\overline{Y})=\eta ([X,\overline{Y}])\xi ,$$ $$N_\varphi (X,\xi )=[X,\xi ]+i\varphi ([X,\xi ])+\eta ([X,\xi ])\xi .$$ Observe that the last component of the above formulae is the action of $`\mathrm{d}\eta \xi `$ on $`(X,Y)`$, $`(X,\overline{Y})`$ and $`(X,\xi )`$, respectively. Therefore $`N_\varphi +\mathrm{d}\eta \xi =0`$ if and only if $$\varphi ([X,Y])=i[X,Y]$$ and $$\varphi ([X,\xi ])=i[X,\xi ].$$ This finishes the proof. Corollary For an almost contact structure satisfying $`N_\varphi +\mathrm{d}\eta \xi =0`$ (in particular for a Sasaki structure) the bundle $`𝐂\xi N`$ is involutive. ## 3 Sasakian geometry in a null frame Let $`(𝒮,(\xi ,\eta ,\varphi ,g))`$ be a Sasakian manifold of dimension $`2k+1`$. The algebraic conditions (1)-(4) of Definition 1 imply an existence of a local basis $`(\xi ,m_i,\overline{m}_{\overline{i}})`$, $`i,\overline{i}=1,2,\mathrm{},k`$, of complex-valued vector fields on $`𝒮`$, with a cobasis $`(\eta ,\mu ^i,\overline{\mu }^{\overline{i}})`$, such that $$g=\eta ^2+2\underset{l=\overline{l},l=1}{\overset{k}{}}\mu ^l\overline{\mu }^{\overline{l}},$$ (1) $$\varphi =i\underset{\overline{j}=1}{\overset{k}{}}(\overline{m}_{\overline{j}}\overline{\mu }^{\overline{j}})i\underset{j=1}{\overset{k}{}}(m_j\mu ^j).$$ (2) Since $`(𝒮,(\xi ,\eta ,\varphi ,g))`$ is Sasakian then its bundle $`𝐂\xi N`$ is involutive. This is equivalent to the condition that the forms $`\mu ^1,\mu ^2,\mathrm{},\mu ^k`$ generate a closed differential ideal i.e. $$\mathrm{d}\mu ^i\mu ^1\mu ^2\mathrm{}\mu ^k=0i=1,2,\mathrm{},k.$$ (3) Condition (6) of Definition 1 of a Sasakian manifold in this basis reads $$\mathrm{d}\eta =2i\underset{l=\overline{l},l=1}{\overset{k}{}}\mu ^l\overline{\mu }^{\overline{l}}.$$ (4) Thus, the fact that the manifold is Sasakian necessarily implies the existence of a local basis $`(\xi ,m_i,\overline{m}_{\overline{i}})`$ with a dual basis $`(\eta ,\mu ^i,\overline{\mu }^{\overline{i}})`$ such that (1)-(4) holds. The converse is also true: if a real vector field $`\xi `$ on a manifold $`𝒮`$ can be suplemented by $`k`$ complex-valued vector fields $`m_i`$ such that $`(\xi ,m_i,\overline{m}_{\overline{i}})`$ and $`(\eta ,\mu ^i,\overline{\mu }^{\overline{i}})`$ form a mutually dual basis for T$`{}_{}{}^{𝐂}𝒮`$ and T$`{}_{}{}^{𝐂}𝒮`$, respectively, satisfying (3)-(4), then the structure $`(\xi ,\eta ,\varphi ,g)`$ defined by $`\xi `$, $`\eta `$ and $`g`$, $`\varphi `$ of (1)-(2) is a Sasakian manifold. This fact can be seen by noting that condition (3) is equivalent to the existence of complex-valued functions $`a_{ijk}`$, $`b_{i\overline{j}k}`$ and $`c_{ij}`$ such that $$\mathrm{d}\mu ^i=\underset{j,n=1}{\overset{k}{}}a_{ijn}\mu ^j\mu ^k+\underset{\overline{j},n=1}{\overset{k}{}}b_{i\overline{j}n}\overline{\mu }^{\overline{j}}\mu ^k+\underset{j}{\overset{k}{}}c_{ij}^j\eta .$$ (5) The dual conditions to conditions (4)-(5) imply that $`N`$ is involutive and that $`[\xi ,\mathrm{\Gamma }(N)]\mathrm{\Gamma }(N)`$. These, when compared with Lemma of Corollary 1, imply condition (7), which is the only condition from Definition 1 which, a’priori, was not assumed for $`(\xi ,\eta ,\varphi ,g)`$. This proves the following Proposition. Proposition 1 (I) (Local version) Let $`(\xi ,\eta ,\varphi ,g)`$ be a Sasakian structure on a manifold $`𝒮`$ of dimension $`2k+1`$. Then there exisists a local basis $`(\xi ,m_i,\overline{m}_{\overline{i}})`$, $`i,\overline{i}=1,2,\mathrm{}k`$ of T$`{}_{}{}^{𝐂}𝒮`$ with dual basis $`(\eta ,\mu ^i,\overline{\mu }^{\overline{i}})`$ such that $$g=\eta ^2+2\underset{l=\overline{l},l=1}{\overset{k}{}}\mu ^l\overline{\mu }^{\overline{l}}$$ $$\varphi =i\underset{\overline{j}=1}{\overset{k}{}}(\overline{m}_{\overline{j}}\overline{\mu }^{\overline{j}})i\underset{j=1}{\overset{k}{}}(m_j\mu ^j),$$ $$\mathrm{d}\mu ^i\mu ^1\mu ^2\mathrm{}\mu ^k=0i=1,2,\mathrm{},k,$$ $$\mathrm{d}\eta =2i\underset{l=\overline{l},l=1}{\overset{k}{}}\mu ^l\overline{\mu }^{\overline{l}}.$$ (II) (Global version) Every almost contact metric structure which satisfies condition (6) of Definition 1 is Sasakian if and only if its canonical decomposition $`\mathrm{T}^𝐂𝒮=𝐂\xi N\overline{N},`$ consists of involutive $`𝐂\xi N`$ part. We close this section with a quick application of part (I) of Proposition 1. It is well known that a vector field $`\xi `$ on a Sasakian manifold $`(𝒮,(\xi ,\eta ,\varphi ,g))`$ is a Killing vector field. This in particular means that the Lie derivatives $`_\xi `$ of $`g`$ and $`\eta `$ vanish. The second of these facts is an immediate consequence of (4). To calculate $`_\xi g`$ one uses (4) and (5). After some work one shows that vanishing of $`_\xi g`$ is equivalent to $`c_{ij}+\overline{c_{ji}}=0`$ where $`c_{ij}`$ are functions appearing in (5). On the other hand, these equations are automatically implied by application of d on both sides of equation (4). ## 4 Analogue of the Kähler potential We pass to a construction of local coordinates on a Sasakian manifold $`(𝒮,(\xi ,\eta ,\varphi ,g))`$. We assume that all the fields defining the Sasakian structure are smooth on $`𝒮`$. In a considered region of $`𝒮`$, we chose a local frame $`(\xi ,m_i,\overline{m}_{\overline{i}})`$ of Proposition 1. Now, the fact that $`\xi `$ is a Killing vector field on $`𝒮`$ together with the complex version of the Fröbenius theorem, assures that condition (3) is equivalent to an existence of complex-valued functions $`f_j^i`$ and $`z^i`$, $`i,j=1,2,\mathrm{}k`$ such that $$\mu ^i=f_j^i\mathrm{d}z^j.$$ (6) Since the forms $`(\mu ^i,\overline{\mu }^{\overline{i}})`$ form a part of the basis on the considered region of $`𝒮`$ then we also have $$\mathrm{d}z^1\mathrm{d}z^2\mathrm{}\mathrm{d}z^k\mathrm{d}\overline{z}^1\mathrm{d}\overline{z}^2\mathrm{}\mathrm{d}\overline{z}^k0.$$ (7) For the basis $`(\xi ,_{z^1},\mathrm{},_{z^k},_{\overline{z}^1},\mathrm{},_{\overline{z}^k})`$ and its dual $`(\eta ,\text{d}z^1,\mathrm{},\text{d}z^k,\text{d}\overline{z}^1,\mathrm{},\text{d}\overline{z}^k)`$ the Maurer-Cartan relations for $`\text{d}(\text{d}z^i)=0=\text{d}(\text{d}\overline{z}^{\overline{i}})`$ readilly show that $`[\xi ,/z^i]=0=[\xi ,/\overline{z}^{\overline{i}}]`$. Therefore, there exists a real coordinate $`x`$ complementary to $`z^1,\mathrm{},z^k,\overline{z}^1,\mathrm{},\overline{z}^k`$ such that $$\xi =_x$$ (8) and the form $`\eta `$ reads $$\eta =\mathrm{d}x+p_j\mathrm{d}z^j+\overline{p}_{\overline{j}}\mathrm{d}\overline{z}^{\overline{j}}.$$ (9) Comparing this with the fact that $`\xi `$ preserves $`\eta `$ leads to the conclusion that the functions $`p_i`$ are independent of coordinate $`x`$, $`p_i/x=0`$. Condition (4) is now equivalent to the follwing two conditions for the differentails of functions $`p_i`$: $$p_{i,j}p_{j,i}=0$$ (10) and $$p_{j,\overline{i}}\overline{p}_{\overline{i},j}=2i\underset{l=\overline{l},l=1}{\overset{k}{}}f_j^l\overline{f}_{\overline{i}}^{\overline{l}}.$$ (11) In a simply connected region of $`𝒮`$ equation (10) guarantees an existence of a complex-valued function $`V`$ such that $$p_i=\frac{V}{z^i}.$$ (12) Since $`p_i`$ is independent of $`x`$ it is enough to consider functions $`V`$ such that $`V/x=0`$. Inserting so determined $`p_i`$ to equation (11) we show that now (11) is equivalent to $$K_{,j\overline{i}}=\underset{l=\overline{l},l=1}{\overset{k}{}}f_j^l\overline{f}_{\overline{i}}^{\overline{l}},$$ (13) where we have introduce Im$`V=K`$ and Re$`V=L`$. Finally we note that now $$\eta =\mathrm{d}(x+L)+i\underset{j=1}{\overset{k}{}}K_{,j}\mathrm{d}z^ji\underset{\overline{j}=1}{\overset{k}{}}K_{,\overline{j}}\mathrm{d}\overline{z}^{\overline{j}},$$ so redefining the $`x`$ coordinate by $`xx+L`$ we simplify $`\eta `$ to the form $`\eta =\mathrm{d}x+i_{j=1}^kK_j\mathrm{d}z^ji_{\overline{j}=1}^kK_{\overline{j}}\mathrm{d}\overline{z}^{\overline{j}}`$. Using (13) we can eliminate functions $`f_j^i`$ from formulae defining our Sasakian structure. Indeed, $$g=\eta ^2+2\underset{l=\overline{l},l=1}{\overset{k}{}}\mu ^l\overline{\mu }^{\overline{l}}=$$ $$\eta ^2+2\underset{l=\overline{l},l=1}{\overset{k}{}}\underset{j,\overline{i}=1}{\overset{k}{}}f_j^l\overline{f}_{\overline{i}}^{\overline{l}}\mathrm{d}z^j\mathrm{d}\overline{z}^{\overline{i}}=$$ $$\eta ^2+2\underset{j,\overline{i}=1}{\overset{k}{}}K_{,j\overline{i}}\mathrm{d}z^j\mathrm{d}\overline{z}^{\overline{i}}.$$ In this way we obtain the following theorem. Theorem 1. Let $`𝒰`$ be an open set of $`𝐂^k\times 𝐑`$ and let $`(z^1,z^2,\mathrm{},z^k,x)`$ be Cartesian coordinates in $`𝒰`$. Consider: * a vector field $`\xi =_x`$ * a real-valued function $`K`$ on $`𝒰`$ such that $`\xi (K)=0`$ * an 1-form $$\eta =\mathrm{d}x+i\underset{m=1}{\overset{k}{}}(K_{,m}\mathrm{d}z^m)i\underset{\overline{m}=1}{\overset{k}{}}(K_{,\overline{m}}\mathrm{d}\overline{z}^{\overline{m}})$$ * a bilinear form $$g=\eta ^2+2\underset{m,\overline{k}=1}{\overset{k}{}}K_{,m\overline{k}}\mathrm{d}z^m\mathrm{d}\overline{z}^{\overline{k}}$$ * a tensor field $$\varphi =i\underset{m=1}{\overset{k}{}}[(_miK_{,m}_x)\mathrm{d}z^m]+i\underset{\overline{m}=1}{\overset{k}{}}(_{\overline{m}}+iK_{,\overline{m}}_x)\mathrm{d}\overline{z}^{\overline{m}}].$$ If the function $`K`$ is choosen in such a way that the bilinear form $`g`$ has positive definite signature then $`𝒰`$ equipped with the structure $`(\varphi ,\xi ,\eta ,g)`$ is a Sasakian manifold. Moreover, every Sasakian manifold can locally be generated by such a function $`K`$. The function $`K`$ appearing in the above theorem is a Sasakian analogue of the Kähler potential generating Kähler geometries. We call it a Sasakian potential. We close this section with a remark that several Sasakian potentials may generate the same Sasakian structure. This is evident if one notes that the following transformations $$KK+f(z^j)+\overline{f}(\overline{z}^{\overline{j}})xx+i\overline{f}(\overline{z}^{\overline{j}})if(z^j),$$ (14) with $`f`$ beeing a holomorphic function of $`z^j`$s, do not change the Sasakian structure of Theorem 1. Thus, transformations (14) are the gauge transformations for the Sasakian potential. ## 5 Locally Sasakian-Einstein structures In this section we calculate the Ricci tensor for the Sasakian metric $`g`$ generated in a region $`𝒰`$ by the Sasakian potential $`K`$ of Theorem 1. We also derive the equation which the Sasakian potential has to obey for Sasakian metric to satisfy Einstein equations $`Ric(g)=\lambda g`$. In this section we use the Einstein summation convention. Let $`𝒰`$ be a simply connected region of of $`𝐂^k\times 𝐑`$ as in Theorem 1. Consider a Sasakian structure defined in this Theorem by the Sasakian potential $`K`$. In the holonomic cobasis $$(\mathrm{d}y^\mu )=(\mathrm{d}x,\mathrm{d}z^j,\mathrm{d}\overline{z}^{\overline{j}})$$ the covariant components of the Sasakian metric read $$g_{\mu \nu }=\left(\begin{array}{ccc}1& iK_{,j}& iK_{,\overline{j}}\\ iK_{,i}& K_{,i}K_{,j}& K_{,i\overline{j}}+K_{,i}K_{,\overline{j}}\\ iK_{,\overline{i}}& K_{,\overline{i}j}+K_{,\overline{i}}K_{,j}& K_{,\overline{i}}K_{,\overline{j}}\end{array}\right).$$ (15) The contravariant components of the metric read $$g^{\mu \nu }=\left(\begin{array}{ccc}1+2K_{,i}K_{,\overline{j}}\kappa ^{i\overline{j}}& iK_{,\overline{j}}\kappa ^{i\overline{j}}& iK_{,j}\kappa ^{j\overline{i}}\\ iK_{,\overline{i}}\kappa ^{j\overline{i}}& 0& \kappa ^{j\overline{i}}\\ iK_{,i}\kappa ^{i\overline{j}}& \kappa ^{i\overline{j}}& 0\end{array}\right),$$ (16) where $$\kappa ^{j\overline{l}}K_{,i\overline{l}}=\delta _i^j\kappa ^{l\overline{j}}K_{,l\overline{i}}=\delta _{\overline{i}}^{\overline{j}}\overline{\kappa ^{l\overline{j}}}=\kappa ^{j\overline{l}}=\kappa ^{\overline{l}j}.$$ (17) The connection 1-forms $`\mathrm{\Gamma }_{\mu \nu }=\frac{1}{2}(g_{\mu \nu ,\rho }+g_{\rho \mu ,\nu }g_{\nu \rho ,\mu })\mathrm{d}y^\rho `$ read $$\mathrm{\Gamma }_{xx}=0\mathrm{\Gamma }_{xi}=iK_{,ij}\mathrm{d}z^j\mathrm{\Gamma }_{x\overline{i}}=\overline{\mathrm{\Gamma }_{xi}}\mathrm{\Gamma }_{ix}=iK_{,i\overline{j}}\mathrm{d}\overline{z}^{\overline{j}}\mathrm{\Gamma }_{\overline{i}x}=\overline{\mathrm{\Gamma }_{ix}}$$ $$\mathrm{\Gamma }_{ij}=K_{,i}K_{,jl}\mathrm{d}z^lK_{,j}K_{,i\overline{l}}\mathrm{d}\overline{z}^{\overline{l}}\mathrm{\Gamma }_{\overline{i}\overline{j}}=\overline{\mathrm{\Gamma }_{ij}}$$ (18) $$\mathrm{\Gamma }_{i\overline{j}}=iK_{,i\overline{j}}\mathrm{d}xK_{,l}K_{,i\overline{j}}\mathrm{d}z^l+(K_{,i\overline{j}\overline{l}}+K_{,i}K_{,\overline{j}\overline{l}}+K_{,\overline{j}}K_{,i\overline{l}}+K_{,\overline{l}}K_{,i\overline{j}})\mathrm{d}\overline{z}^{\overline{l}}\mathrm{\Gamma }_{,\overline{i}j}=\overline{\mathrm{\Gamma }_{i\overline{j}}}.$$ It is convenient to introduce the following functions: $$C_{jm}^i=\kappa ^{i\overline{l}}K_{,\overline{l}jm}B_{jm}^i=C_{jm}^i+\delta _m^iK_{,j}+\delta _j^iK_{,m}A_{jm}=C_{jm}^lK_{,l}+2K_{,j}K_{,m}K_{,jm}$$ $$C_{\overline{j}\overline{m}}^{\overline{i}}=\overline{C_{jm}^i}B_{\overline{j}\overline{m}}^{\overline{i}}=\overline{B_{jm}^i}A_{\overline{j}\overline{m}}=\overline{A_{jm}}.$$ Then the connection 1-forms $`\mathrm{\Gamma }_\nu ^\mu `$ read $$\mathrm{\Gamma }_x^x=\mathrm{d}K\mathrm{\Gamma }_j^x=K_{,j}\mathrm{d}xiA_{jm}\mathrm{d}z^m\mathrm{\Gamma }_x^i=i\mathrm{d}z^i\mathrm{\Gamma }_{\overline{j}}^x=\overline{\mathrm{\Gamma }_j^x}\mathrm{\Gamma }_x^{\overline{i}}=\overline{\mathrm{\Gamma }_x^i}$$ $$\mathrm{\Gamma }_j^i=i\delta _j^i\mathrm{d}x\delta _j^iK_{,\overline{l}}\mathrm{d}\overline{z}^{\overline{l}}+B_{jl}^i\mathrm{d}z^l\mathrm{\Gamma }_{\overline{j}}^i=K_{,\overline{j}}\mathrm{d}z^i\mathrm{\Gamma }_{\overline{j}}^{\overline{i}}=\overline{\mathrm{\Gamma }_j^i}\mathrm{\Gamma }_j^{\overline{i}}=\overline{\mathrm{\Gamma }_{\overline{j}}^i}.$$ The curvature 2-forms $`\mathrm{\Omega }_\nu ^\mu =\frac{1}{2}R_{\nu \rho \sigma }^\mu \mathrm{d}y^\rho \mathrm{d}y^\sigma =\mathrm{d}\mathrm{\Gamma }_\nu ^\mu +\mathrm{\Gamma }_\rho ^\mu \mathrm{\Gamma }_\nu ^\rho `$ read $$\mathrm{\Omega }_x^x=iK_{,l}\mathrm{d}x\mathrm{d}z^liK_{,\overline{l}}\mathrm{d}x\mathrm{d}\overline{z}^{\overline{l}}$$ $$\mathrm{\Omega }_j^x=K_{,j}K_{,l}\mathrm{d}x\mathrm{d}z^l+(K_{,j\overline{l}}+K_{,j}K_{,\overline{l}})\mathrm{d}x\mathrm{d}\overline{z}^{\overline{l}}+iA_{jl,\overline{m}}\mathrm{d}z^l\mathrm{d}\overline{z}^{\overline{m}}$$ $$\mathrm{\Omega }_x^i=\delta _j^i\mathrm{d}x\mathrm{d}z^j+i\delta _j^iK_{,l}\mathrm{d}z^j\mathrm{d}z^l+iK_{,\overline{l}}\mathrm{d}\overline{z}^{\overline{l}}\mathrm{d}z^i$$ $$\mathrm{\Omega }_j^i=i\delta _l^iK_{,j}\mathrm{d}x\mathrm{d}z^l+(B_{mn}^iB_{jl}^mB_{jn,l}^i+A_{jn}\delta _l^i)\mathrm{d}z^n\mathrm{d}z^l+(K_{,j}K_{,\overline{l}}\delta _n^i\delta _j^iK_{,n\overline{l}}B_{jn,\overline{l}}^i)\mathrm{d}z^n\mathrm{d}\overline{z}^{\overline{l}}$$ $$\mathrm{\Omega }_{\overline{j}}^i=i\delta _l^iK_{,\overline{j}}\mathrm{d}x\mathrm{d}z^l+(K_{,\overline{j}l}+K_{,\overline{j}}K_{,l})\delta _n^i\mathrm{d}z^n\mathrm{d}z^l\delta _n^iK_{,\overline{j}}K_{\overline{l}}\mathrm{d}z^n\mathrm{d}\overline{z}^{\overline{l}}$$ $$\mathrm{\Omega }_{\overline{j}}^x=\overline{\mathrm{\Omega }_j^x}\mathrm{\Omega }_x^{\overline{i}}=\overline{\mathrm{\Omega }_x^i}\mathrm{\Omega }_{\overline{j}}^{\overline{i}}=\overline{\mathrm{\Omega }_j^i}\mathrm{\Omega }_j^{\overline{i}}=\overline{\mathrm{\Omega }_{\overline{j}}^i}.$$ The Ricci tensor $`R_{\nu \sigma }=R_{\nu \mu \sigma }^\mu `$ componenets read $$R_{xx}=2kR_{xj}=2ikK_{,j}R_{ij}=2kK_{,i}K_{,j}R_{\overline{i}j}=2kK_{,\overline{i}}K_{,j}2K_{,\overline{i}j}C_{\overline{m}\overline{i},j}^{\overline{m}}$$ $$R_{x\overline{j}}=\overline{R_{xj}}R_{\overline{i}\overline{j}}=\overline{R_{ij}}.$$ Now, the Einstein equations $`Ric(g)=\lambda g`$, which are nontrivial only for the components $`R_{xx}`$ and $`R_{\overline{i}j}`$ become $$\lambda =2k(\kappa ^{\overline{m}l}K_{,\overline{m}l\overline{i}})_{,j}=2(k+1)K_{,\overline{i}j}.$$ Since the matrix $`(\kappa ^{\overline{m}l})`$ is the inverse of $`(K_{,\overline{i}j})`$ then the left hand side of the second equations above is $$(\kappa ^{\overline{m}l}K_{,\overline{m}l\overline{i}})_{,j}=\mathrm{log}(det(K_{,m\overline{n}}))_{,\overline{i}j},$$ see e.g. . Thus we arrive to the following theorem. Theorem 2 Any Sasakian manifold of dimension $`(2k+1)`$ can be locally represented by the Sasakian potential $`K`$ of Theorem 1. In the region where the potential is well defined the manifold satisfies Einstein equations $`Ric(g)=\lambda g`$ if and only if the cosmological constant $$\lambda =2k$$ and the potential satisfies $$\mathrm{log}(det(K_{,m\overline{n}}))_{,\overline{i}j}=2(k+1)K_{,\overline{i}j}.$$ (19) Surprisingly equation (19) is the same as the Einstein condition $`Ric(h)=2(k+1)h`$ for the Kähler metric $`h=2K_{,\overline{i}j}\mathrm{d}\overline{z}^{\overline{i}}\mathrm{d}z^j`$ in dimension $`2k`$. Thus we have the following Corollary. Corollary Every Sasakian-Einstein manifold in dimension $`(2k+1)`$ is locally in one to one correspondence with a Kähler-Einstein manifold in dimension $`2k`$ whose cosmological constant $`\lambda =2(k+1)`$. The correspondence is obtained by identifying the Kähler potential for the Kähler-Einstein manifold with the Sasaki potential for the Sasaki-Einstein manifold. Examples 1). Sasakian potential for the sphere $`𝐒^{2k+1}`$. Consider a function $$𝒦=\frac{1}{2}\mathrm{log}(z^1\overline{z}^1+\mathrm{}+z^{k+1}\overline{z}^{k+1})$$ defined on $`𝐂^{k+1}\{0\}`$. Let $$N=i(𝒦_{,j}\text{d}z^j𝒦_{,\overline{j}}\text{d}\overline{z}^{\overline{j}}),$$ $$H=2𝒦_{,i\overline{j}}\text{d}z^i\text{d}\overline{z}^{\overline{j}}$$ and $$G=N^2+H.$$ The tensor fields $`N`$ and $`G`$ restrict to a sphere $$𝐒^{2k+1}=\{(z^1,\mathrm{},z^{k+1})𝐂^{k+1}\{0\}\|z^1\overline{z}^1+\mathrm{}+z^{k+1}\overline{z}^{k+1}=1\}.$$ Denote these restrictions by $`\eta `$ and $`g`$, respectively. Then the 1-form $`\eta `$ and the Riemannian metric $`g`$ define a Sasakian-Einstein structure on $`𝐒^{2k+1}`$. This structure coincides with the one defined in Example 1 of Section 1. To see this, recal the Hopf fibration $`𝐔(1)𝐒^{2k+1}\mathrm{𝐂𝐏}^k`$ with the action of $`\mathrm{e}^{i\varphi }𝐔(1)`$ on $`(z^1,\mathrm{},z^{k+1})𝐒^{2k+1}`$ defined by $`\mathrm{e}^{i\varphi }(z^1,\mathrm{},z^{k+1})=(\mathrm{e}^{i\varphi }z^1,\mathrm{},\mathrm{e}^{i\varphi }z^{k+1})`$. The canonical projection is given by $`𝐒^{2k+1}(z^1,\mathrm{},z^{k+1})\mathrm{dir}(z^1,\mathrm{},z^{k+1})\mathrm{𝐂𝐏}^k`$. The sphere is covered by $`k+1`$ charts $$U_j=V_j\times 𝐔(1),$$ where $$V_j=\{\mathrm{dir}(z^1,\mathrm{},z^{k+1})\|(z^1,\mathrm{},z^{k+1})𝐒^{2k+1}\mathrm{and}z^j0\}.$$ The local coordinates on each $`U_j`$ are $$(\xi ^{ij}=\frac{z^i}{z^j},\varphi _j=\frac{i}{2}\mathrm{log}\frac{z^j}{\overline{z}^j}),i=1,\mathrm{},k+1,ij.$$ Then on each chart $`U_j`$ the form $`\eta _j=\eta \|_{U_j}`$ reads $$\eta _j=\text{d}\varphi _j+\frac{i}{2}\frac{_{i=1,ij}^{k+1}(\overline{\xi }^{ij}\text{d}\xi ^{ij}\xi ^{ij}\text{d}\overline{\xi }^{ij})}{1+_{i=1,ij}^{k+1}\|\xi ^{ij}\|^2}.$$ The metric $`g`$ restricted to $`U_j`$ is $$g_j=(\eta _j)^2+\frac{(1+_{i=1,ij}^{k+1}\|\xi ^{ij}\|^2)(_{i=1,ij}^{k+1}\|\text{d}\xi ^{ij}\|^2)\|_{i=1,ij}^{k+1}(\xi ^{ij}\text{d}\overline{\xi }^{ij})\|^2}{(1+_{i=1,ij}^{k+1}\|\xi ^{ij}\|^2)^2}.$$ Now, observe that on each $`U_j`$ the structure $`(g_j,\eta _j)`$ may be obtained by means of Theorems 1 and 2 choosing a Sasakian potential $$K^j=\frac{1}{2}\mathrm{log}(1+\underset{i=1,ij}{\overset{k+1}{}}\|\xi ^{ij}\|^2)$$ on the corresponding $`V_j`$. It is easy to check that $`K^j`$ satisfies equation (19) on $`V_j`$. Thus, Theorem 2 assures that the Sasakian structure generated by $`(g_j,\eta _j)`$ is Sasakian-Einstein. Easy, but lenghty, calculation shows that the Weyl tensor of $`g_j`$ vanishes identically on $`U_j`$. This proves that $`(U_j,g_j)`$ is locally isometric to a standard Riemannian structure on $`𝐒^{2k+1}`$. Since $`(g_j,\eta _j)`$ originate from the global structure $`(g,\eta )`$ then this global Sasakian structure must coincide with the standard Sasakian structure of Example 1. Note also that $`h_j=g_j(\eta _j)^2`$ projects to $`V_j`$ and patched together defines the Fubini-Study metric on $`\mathrm{𝐂𝐏}^k`$. In this sense the standard Sasakian structure on $`𝐒^{2k+1}`$ described in Example 1 is the analogue of the Fubini-Study Kähler structure on $`\mathrm{𝐂𝐏}^k`$. 2.) Sasakian-Einstein structure on $`𝐂^q\times 𝐂^n\times 𝐑`$. Consider a function $$K=\frac{1}{q+n+1}[\underset{i=1}{\overset{q}{}}\mathrm{log}(1+\|v^i\|^2)]+\frac{n+1}{2(q+n+1)}\mathrm{log}(1+\underset{I=1}{\overset{n}{}}\|w^I\|^2)$$ defined on $`𝐂^q\times 𝐂^n`$, with coordinates $`(z^\mu )=(v^i,w^I)`$. It is easy to check that $$[\mathrm{log}det(K_{,\mu \overline{\nu }})]_{,\rho \overline{\sigma }}=2(q+n+1)K_{,\rho \overline{\sigma }}.$$ Thus, via Theorems 1 and 2, such $`K`$ generates a Sasakian-Einstein structure on $`𝐂^q\times 𝐂^n\times 𝐑^1.`$ 3) Locally Sasakian-Einstein structures in dimension 5. If $`k=2`$ then, modulo the gauge (14), equation (19) may be integrated to the form $$K_{,1\overline{1}}K_{,2\overline{2}}K_{,1\overline{2}}K_{,2\overline{1}}=\mathrm{e}^{6K}.$$ This is a well known equation describing the gravitational instantons in four dimensions . Examples of the Kähler-Einstein metrics generated by its solutions can be found e.g. in . Via Theorems 1 and 2, each of these Kähler-Einstein structures defines a nontrivial Sasakian-Einstein manifold in dimension 5.
warning/0005/math0005137.html	ar5iv	text	# Homology for irregular connections ## 0. Introduction Consider the following formulas, culled, one may imagine, from a textbook on calculus: $`\sqrt{\pi }={\displaystyle _{\mathrm{}}^{\mathrm{}}}e^{t^2}𝑑t`$ $`(e^{2\pi is}1)\mathrm{\Gamma }(s)=(e^{2\pi is}1){\displaystyle _0^{\mathrm{}}}e^tt^s{\displaystyle \frac{dt}{t}}`$ Gamma function $`J_n(z)={\displaystyle \frac{1}{2\pi i}}{\displaystyle _{\{\|u\|=ϵ\}}}\mathrm{exp}({\displaystyle \frac{z}{2}}(u{\displaystyle \frac{1}{u}})){\displaystyle \frac{du}{u^{n+1}}}`$ Bessel function. These are a few familiar examples of periods associated to connections with irregular singular points on Riemann surfaces. Curiously, though of course such integrals have been studied for 200 years or so, and mathematicians in recent years have developed a powerful duality theory for holonomic $`𝒟`$-modules (for dimension $`1`$, which is the only case we will consider, cf. ,chap. IV, and ), it is not easy from the literature to interpret such integrals as periods arising from a duality between homological cycles and differential forms. A homological duality of this sort is well understood for differential equations with regular singular points. Our purpose in this note is to develop a similar theory in the irregular case. Of course, most of the “heavy lifting” was done by Malgrange op. cit. We hope, in reinterpreting his theory, to better understand relations between irregular connections and wildly ramified $`\mathrm{}`$-adic sheaves. There are striking relations between $`ϵ`$-factors for $`\mathrm{}`$-adic sheaves on curves over finite fields and determinants of irregular periods which merit further study. Finally, relations between irregular connections and the arithmetic theory of motives remain mysterious. Let $`X`$ be a smooth, compact, connected algebraic curve (Riemann surface) over $``$. Let $`D=\{x_1,\mathrm{},x_n\}X`$ be a non-empty, finite set of points (which we also think of as a reduced effective divisor), and write $`U:=XD\stackrel{j}{}X`$. Let $`E`$ be a vector bundle on $`X`$, and suppose given a connection with meromorphic poles on $`D`$ $$:EE\omega (D).$$ Here $`\omega `$ is the sheaf of holomorphic $`1`$-forms on $`X`$ and $`D`$ refers to meromorphic poles on $`D`$. Unless otherwise indicated, we work throughout in the analytic topology. The de Rham cohomology $`H_{DR}^{}(XD;E,)`$ is the cohomology of the complex of sections (0.1) $$\mathrm{\Gamma }(X,E(D))\stackrel{}{}\mathrm{\Gamma }(X,E\omega (D))$$ placed in degrees $`0`$ and $`1`$. These cohomology groups are finite dimensional , Proposition 6.20, (i). Let $`E^{}`$ be the dual bundle, and let $`^{}`$ be the dual connection, so (0.2) $$de,f^{}=(e),f^{}+e,^{}(f^{})$$ Define $`=\mathrm{ker}()`$, and $`^{}=\mathrm{ker}(^{})`$ to be the corresponding local systems of flat sections on $`U`$. We want to define homology with values in these local systems, or more precisely with values in associated cosheaves on $`X`$. For $`xXD`$, $`_x`$ will denote the stalk of $``$ at $`x`$. Define the co-stalk at $`0D`$ (0.3) $$_0:=_x/(1\sigma )_x$$ where $`x0`$ is a nearby point, and $`\sigma `$ is the local monodromy about $`0`$. We write $`𝒞_n=𝒞_n(E,)`$ for the group of $`n`$-chains with values in $``$ and rapid decay near $`0`$. Write $`\mathrm{\Delta }^n`$ for the $`n`$-simplex and $`b\mathrm{\Delta }^n`$ for its barycenter. Thus, $`𝒞_n(E,)`$ is spanned by elements $`cϵ`$ with $`c:\mathrm{\Delta }^nX`$ and $`ϵ_{c(b)}`$, where $`b\mathrm{\Delta }^n`$ is the barycenter. We assume $`c^1(0)=\text{union of faces}\mathrm{\Delta }^n`$ and that $`ϵ`$ has rapid decay near $`D`$. This is no condition if $`Dc(\mathrm{\Delta }^n)=\mathrm{}`$. If $`0Dc(\mathrm{\Delta }^n)`$, we take $`e_i`$ a basis for $`E`$ near $`0`$ and write $`ϵ=f_ic^{}(e_i)`$. Let $`z`$ be a local parameter at $`0`$ on $`\mathrm{\Delta }`$. We require that for all $`N`$, constants $`C_N>0`$ exist with $`\|f_i(z)\|C_N\|z\|^N`$ on $`\mathrm{\Delta }^nc^1(0)`$. Note that if $``$ has logarithmic poles in one point, then rapid decay implies vanishing. Thus in this case, we deal with the sheaf $`j_!`$, where $`j:XDX`$. There is a natural boundary map (0.4) $$:𝒞_n(E,)𝒞_{n1}(E,);(cϵ)=(1)^jc_jϵ_j$$ where $`c_j`$ are the faces of $`c`$. Note if $`b_j`$ is the barycenter of the $`j`$-th face and $`c(b_j)0`$, $`c`$ determines a path from $`c(b)`$ to $`c(b_j)`$ which is canonical upto homotopy on $`\mathrm{\Delta }\{0\}`$. (As a representative, one can take $`c[b_j,b]`$, the image of the straight line from $`b`$ to $`b_j`$. By assumption, $`c^1(0)`$ is a union of faces, so it does not meet the line.) Thus $`ϵ_{c(b)}`$ determines $`ϵ_j_{c(b_j)}`$. Similarly for $`0D`$, if $`c(b_j)=0`$ there is corresponding to $`ϵ`$ a unique $`ϵ_j_0`$ because we have taken coinvariants. If $`c:\mathrm{\Delta }^nD`$ is a constant simplex, there is no rapid decay condition. It is straightforward to compute that $`=0`$. Consider $`cϵ`$. If $`c(b)=0`$, where $`b\mathrm{\Delta }^2`$ is the barycentre, then $`c(\mathrm{\Delta }^2)=0`$ and $`ϵ=ϵ_i=(ϵ_i)_j_0`$ for all $`i`$ and $`j`$ involved, thus the condition is trivially fulfilled. If not, and some $`c(b_i)=0`$, then $`(ϵ_j)_i=(ϵ_i)_j_0`$ for all $`j`$, and if all $`c(b_i)0`$, then one has by unique analytic continuation in $`c(\mathrm{\Delta }^2)`$ the relation $`(ϵ_i)_j=(ϵ_j)_i_{\mathrm{edge}_{ij}}`$ for all $`i,j`$, if edge$`{}_{ij}{}^{}0`$, else in $`_0`$. We define (0.5) $$H_{}(X,D;E^{},^{}):=H_{}\left(𝒞_{}(X;E^{},^{})/𝒞_{}(D;E^{},^{})\right).$$ (The growth condition means this depends on more than just the topological sheaf $`^{}`$, so we keep $`E^{},^{}`$ in the notation.) We now define a pairing (0.6) $$(,):H_{DR}^{}(XD;E,)\times H_{}(X,D;E^{},^{});=0,1$$ by integrating over chains in the following manner. For $`=0`$, then $`H_0(X,D;E^{},^{})`$ is generated by sections of the dual local system $`^{}`$ in points $`X`$ while $`H_{DR}^0(XD;E,)`$ is generated by global flat sections in $``$ with moderate growth. So one can pair them. For $`=1`$, since $`D\mathrm{}`$, then $$H_{DR}^1(XD;E,)=H^0(X,\omega E(D))/H^0(X,E(D)),$$ and since classes $`cϵ`$ generating $`H_0(X,D;E^{},^{})`$ have rapid decay, the integral $`_c<f_ic^{}(e_i),\alpha >`$ is convergent, where $`\alpha H^0(X,\omega E(D))`$ and $`<>`$ is the duality between $`E^{}`$ and $`E`$. The rest of the note is devoted to the proof of the following theorem. ###### Theorem 0.1. The process of integrating forms over chains is compatible with homological and cohomological equivalences and defines a perfect pairing of finite dimensional vector complex spaces $$(,):H_{DR}^{}(XD;E,)\times H_{}(X,D;E^{},^{});=0,1.$$ ###### Example 0.2. (i). If $``$ has regular singular points, there are no rapidly decaying flat sections, so $`H_{}(X,D;E^{},^{})H_{}(XD;^{})`$. Also, $`H_{DR}^{}(XD;E,)H^{}(U,)`$ (cf. , Théorème 6.2), and the theorem becomes the classical duality between homology and cohomology. (ii). Suppose $`X=^1`$, $`D=\{0,\mathrm{}\}`$. Let $`E=𝒪_^1`$ with connection $`(1)=dt+s\frac{dt}{t}`$, for some $`s\{0,1,2,\mathrm{}\}`$. Then $`E_U=𝒪_U`$ is the trivial local system spanned by $`e^tt^s`$, so $`^{}E_U^{}=𝒪_U`$ is spanned by $`e^tt^s`$. We consider the pairing $`H_{DR}^1\times H_1`$ from theorem 0.1. Note first that $`H_{DR}^1`$ has dimension $`1`$, spanned by $`\frac{dt}{t}`$. This can either be checked directly from (0.1), using $$(t^p)=((p+s)t^{p1}t^p)dt,$$ or by showing the de Rham cohomology is isomorphic to the hypercohomology of the complex $`𝒪_^1\stackrel{}{}\omega ((0)+2(\mathrm{}))`$, which is easily computed. To compute $`H_1(X,D;E^{},^{})`$, the singularity at $`0`$ is regular, so there are no non-constant, rapidly decaying chains at $`0`$. The section $`ϵ^{}:=e^tt^s`$ of $`^{}`$ is rapidly decaying on the positive real axis near $`\mathrm{}`$, so the chain $`cϵ^{}`$ in fig. $`1`$ above represents a $`1`$-cycle. We have $$(ce^tt^s,\frac{dt}{t})=(e^{2\pi is}1)_0^{\mathrm{}}e^tt^s\frac{dt}{t}$$ which is a variant of Hankel’s formula (see , p. 245). (iii). Let $`X,D,E`$ be as in (ii), but take $`(1)=\frac{1}{2}(d(zu)d(\frac{z}{u}))`$ for some $`z\{0\}`$. Here the connection has pole order $`2`$ at $`0`$ and $`\mathrm{}`$ and it has trivial monodromy. Arguing as above, one computes $`dimH_{DR}^1=2`$, generated by $`u^pdu,p`$, with relations $`u^pdu=\frac{2p}{z}u^{p1}duu^{p2}du`$. The Gauß-Manin connection on this group is $$_{GM}(u^pdu)=\frac{1}{2}(u^{p+1}u^{p1})dudz.$$ Assume $`\text{Im}(z)>0`$. Then the vector space $`H_1(^1,\{0,\mathrm{}\};E^{},^{})`$ is generated by $$\{\|u\|=1\}\mathrm{exp}(\frac{1}{2}z(u\frac{1}{u})),\text{and}[0,i\mathrm{}]\mathrm{exp}(\frac{1}{2}z(u\frac{1}{u})).$$ (If $`\text{Im}(z)0`$, then the second path must be modified.) The integrals $$J_n(z):=_{\{\|u\|=1\}}\mathrm{exp}(\frac{1}{2}z(u\frac{1}{u}))\frac{du}{u^{n+1}};$$ $$H_n(z):=_0^i\mathrm{}\mathrm{exp}(\frac{1}{2}z(u\frac{1}{u}))\frac{du}{u^{n+1}}$$ are periods and satisfy the Bessel differential equation $$z^2\frac{d^2y}{dz^2}+z\frac{dy}{dz}+(z^2n^2)y=0$$ The function $`J_n`$ is entire. To show that $`H_n`$ is linearly independent of $`J_n`$, it will then be sufficient to show that $`H_n`$ is unbounded on the positive part of the imaginary axis $`\text{Re}(z)=0`$ as $`z0`$. Making the coordinate change $`v=\frac{1}{u}`$, and replacing $`y`$ by $`\frac{1}{2}y`$ one is led to show that $`E_n(y)=_0^{\mathrm{}}\mathrm{exp}(y(v+\frac{1}{v}))\frac{dv}{v^{n+1}}`$ is unbounded for $`y>0,y0`$. Writing $`E_n(v)=_0^1+_1^{\mathrm{}}`$, and making the change of variable $`v\frac{1}{v}`$ in the integal $`_0^1`$, one obtains $$E_n(y)=_1^{\mathrm{}}\mathrm{exp}(y(v+\frac{1}{v}))(\frac{1}{v^{n+1}}+v^{n1})𝑑v$$ $$_1^{\mathrm{}}\mathrm{exp}(2yv)(\frac{1}{v^{n+1}}+v^{n1})𝑑v.$$ For $`\|n\|1`$, then this expression is $`_1^{\mathrm{}}\mathrm{exp}(2yv)𝑑v`$ which is obviously unbounded. For $`n=0`$, one has $$E_0(y)2_1^{\mathrm{}}\mathrm{exp}(2yv)\frac{dv}{v}$$ $$2_{2y}^{\mathrm{}}\mathrm{exp}(v)\frac{dv}{v}$$ $$2_{2y}^1\mathrm{exp}(v)\frac{dv}{v},$$ where in the last inequality, we have assumed that $`2y1`$. This last integral is, upto something bounded, equal to $`2_{2y}^1\frac{dv}{v}=2\mathrm{log}(2y)`$, which is unbounded, as $`y>0,y0`$. Usually, for integers $`n`$, one considers $`J_n`$ as one standard solution, but not $`H_n`$ (see , p.371). Finally, to get Bessel functions for non-integral values of $`n`$, one may consider the connection $`(1)=\frac{1}{2}(d(zu)d(\frac{z}{u}))n\frac{du}{u}`$. ## 1. Chains Let $`D=\{x_1,\mathrm{},x_n\}`$ be as above, and let $`\mathrm{\Delta }_i`$ be a small disk about $`x_i`$ for each $`i`$. Let $`\delta _i`$ be the boundary circle. Define (1.1) $$\begin{array}{c}H_{}(\mathrm{\Delta }_i,\delta _i\{x_i\};E,)\hfill \\ \hfill =H_{}(𝒞_{}(\mathrm{\Delta }_i;E,)/(𝒞_{}(\delta _i;E,)+𝒞_{}(\{x_i\};E,))\end{array}$$ (Note, for a set like $`\delta _i`$ which is closed and disjoint from $`D`$, our chains coincide with the usual topological chains with values in the local system $``$. The group $`𝒞_{}(\{x_i\};E,)`$ consists of constant chains $`c:\mathrm{\Delta }^n\{x_i\}`$ with values in $$_{x_i}:=_x/(1\mu _i)_x$$ for some $`x`$ near $`x_i`$ as in (0.3), where $`\mu _i`$ is the local monodromy around $`x_i`$.) In the following theorem, $`H_{}(U,)`$ is the standard homology associated to the local system on $`U=XD`$. ###### Theorem 1.1. With notation as above, there is a long exact sequence (1.2) $$\begin{array}{c}0H_1(U,)H_1(X,D;E,)_iH_1(\mathrm{\Delta }_i,\delta _i\{x_i\};E,)\hfill \\ \hfill H_0(U,)H_0(X,D;E,)0\end{array}$$ ###### Proof. Let $`𝒞_{}:=𝒞_{}(X;E,)/𝒞_{}(D;E,)`$ be the complex calculating $`H_{}(X,D;E,)`$, and let $$𝒞_{}(U)𝒞_{}$$ be the subcomplex calculating $`H_{}(U,)`$, i.e. the subcomplex of chains whose support is disjoint from $`D`$. Of course, one has $`𝒞_{}(U;E,)=𝒞_{}(U;)`$, which justifies the notation. Write $`=𝒞_{}/𝒞_{}(U)`$. There is an evident map of complexes (1.3) $$\psi :_i𝒞_{}(\mathrm{\Delta }_i,\delta _i\{x_i\};E,)$$ which must be shown to be a quasi-isomorphism. Let $$(i)=\psi (𝒞_{}(\mathrm{\Delta }_i,\delta _i\{x_i\};E,))=$$ $$𝒞_{}(\mathrm{\Delta }_i,\delta _i\{x_i\};E,)/𝒞_{}(\mathrm{\Delta }_i\{x_i\};).$$ Obviously the map $`\alpha :_i(i)`$ is an inclusion. We claim first that $`\alpha `$ is a quasi-isomorphism. To see this, note that all these complexes admit subdivision maps subd which are homotopic to the identity. Given a chain $`c`$, there exists an $`N`$ such that $`\text{subd}^N(c)(i)`$. Taking $`c`$ with $`c=0`$, it follows that $`H_{}((i))`$ surjects onto $`H_{}()`$. If $`\alpha (x)=y`$, we choose $`N`$ such that $`\text{subd}^N(y)=\alpha (z)`$. Since $`\alpha `$ is injective and commutes with subd, it follows that $`\alpha `$ is injective on homology as well, so $`\alpha `$ is a quasi-isomorphism. It remains to show the surjective map of complexes $$\beta :𝒞_{}(\mathrm{\Delta }_i,\delta _i\{x_i\};E,)(i)$$ is a quasi-isomorphism. The kernel of $`\beta `$ is $$𝒞_{}(\mathrm{\Delta }_i\{x_i\};)/𝒞_{}(\delta _i;),$$ which is acyclic as $`\delta _i\mathrm{\Delta }_i\{x_i\}`$ admits an evident homotopy retract. The next point is to show (1.4) $$H_{}(\mathrm{\Delta }_i,\delta _i\{x_i\};E,)=(0);i=0,2.$$ The assertion for $`H_0`$ is easy because any point $`y`$ in $`\mathrm{\Delta }_i\{x_i\}`$ can be attached to $`\delta _i`$ by a radial path $`r`$ not passing through $`x_i`$. Then $`ϵ_y`$ extends uniquely to $`ϵ`$ on $`r`$ and $`(rϵ)yϵmod\text{chains on }\delta _i`$. Vanishing in (1.4) when $`i=2`$ will be proved in a sequence of lemmas. For convenience we drop the subscript $`i`$ and replace $`x_i`$ with $`0`$. ###### Lemma 1.2. Let $`\mathrm{}\mathrm{\Delta }`$ be a radial line meeting $`\delta `$ at $`p`$. Let $`_{\mathrm{}}`$ be the space of sections of the local system along $`\mathrm{}\{0\}`$ with rapid decay at $`0`$. Then $$H_{}(\mathrm{},\{0,p\};E,)\{\begin{array}{cc}0\hfill & 0\hfill \\ _{\mathrm{}}\hfill & =1\hfill \end{array}$$ ###### proof of lemma. Let $`𝒞_{}(\mathrm{})`$ be the complex of chains calculating this homology, and let $`𝒞_{}(\mathrm{}\{0\})𝒞_{}(\mathrm{})`$ be the subcomplex of chains not meeting $`0`$. Then $`𝒞_{}(\mathrm{}\{0\})`$ is contractible, and $$𝒞_{}(\mathrm{})/𝒞_{}(\mathrm{}\{0\})(_{}(\mathrm{})/_{}(\mathrm{}\{0\}))_{\mathrm{}}$$ where $`_{}`$ denotes classical topological chains. The result follows. ∎ One knows from the theory of irregular connections in $`dim1`$ that $`\mathrm{\Delta }\{0\}`$ can be covered by open sectors $`V\mathrm{\Delta }`$ such than (1.5) $$E,\|_V_i(L_iM_i)$$ where $`L_i`$ is rank $`1`$ and $`M_i`$ has a regular singular point. Let $`WV\{0\}`$ be a smaller closed sector with outer boundary $`\delta _W=\delta W`$ and radial sides $`\mathrm{}_1,\mathrm{}_2`$. Recall the Stokes lines are radial lines where the horizontal sections of the $`L_i`$ shift from rapid decay to rapid growth. We assume $`W`$ contains at most one Stokes line, and that $`\mathrm{}_1,\mathrm{}_2`$ are not Stokes lines. Writing $`W=W_1W_2`$, where $`W_i`$ are even smaller sectors, each of which containing the Stokes line if there is one, one may think of the following lemma as a Mayer-Vietoris sequence. ###### Lemma 1.3. With notation as above, Let $`w`$ be a basepoint in the interior of $`W`$. then $$H_{}(W,\delta _W\{0\};E,)\{\begin{array}{cc}0\hfill & 1\hfill \\ _\mathrm{}_1+_\mathrm{}_2_w\hfill & =1\hfill \end{array}$$ ###### proof of lemma. One has $$_iH_1(\mathrm{}_i,\{0,p_i\};E,)H_1(W,\delta _W\{0\};E,)$$ and of course the assertion of the lemma is that this coincides with $`_\mathrm{}_1_\mathrm{}_2_\mathrm{}_1+_\mathrm{}_2`$. To check this, by (1.5) one is reduced to the case $`E=LM`$ where $`L`$ has rank $`1`$ and $`M`$ has regular singular points. If $`W`$ does not contain a Stokes line for $`L`$ then $`_\mathrm{}_1=_\mathrm{}_2=_\mathrm{}_1+_\mathrm{}_2`$, and the argument is exactly as in lemma 1.2. Suppose $`W`$ contains a Stokes line for $`L`$. Then (say) $`_\mathrm{}_1=_w`$ and $`_\mathrm{}_2=(0)`$. Let $`𝒞_{}(W)`$ be the complex of chains calculating the desired homology, and let $`𝒞_{}(W\{0\})𝒞_{}(W)`$ be the chains not meeting $`0`$. As in the previous lemma, $`𝒞_{}(W\{0\})`$ is acyclic. We claim the map $$𝒞_{}(\mathrm{}_1)𝒞_{}(W)/𝒞_{}(W\{0\})$$ is a quasi-isomorphism. If we choose an angular coordinate $`\theta `$ such that $$\mathrm{}_1:\theta =0;\text{Stokes}:\theta =a>0;\mathrm{}_2:\theta =b>a,$$ then rotation $`re^{i\theta }re^{(1t)i\theta }`$ provides a homotopy contraction of the inclusion of $`\mathrm{}_1W`$. This homotopy contraction preserves the condition of rapid decay, proving the lemma. ∎ Let $`\pi _d:\mathrm{\Delta }\mathrm{\Delta }`$ be the ramified cover of degree $`d`$ obtained by taking the $`d`$-th root of a parameter at $`0`$. By the theory of formal connections , one has, for suitable $`d`$, a decomposition as in (1.5) for the formal completion of the pullback $`\widehat{\pi _d^{}E}_iL_iM_i`$. Let $`m_i`$ be the degree of the pole of the connection on $`L_i`$ when we identify $`L_i\widehat{𝒪}`$, i.e. $`_{L_i}(1)=g_i(z)dz`$ for a local parameter $`z`$, and $`m_i`$ is the order of pole of $`g_i`$. ###### Lemma 1.4. We have $$dimH_p(\mathrm{\Delta },\delta \{0\};E,)=\{\begin{array}{cc}0\hfill & p1\hfill \\ \frac{1}{d}_{m_i2}(m_i1)dim(M_i)\hfill & p=1.\hfill \end{array}$$ ###### proof of lemma. Assume first that we have a decomposition of the type (1.5) on $`\widehat{E}`$ itself, i.e. that no pullback $`\pi _d^{}`$ is necessary. We write $`\mathrm{\Delta }`$ as a union of closed sectors $`W_0,\mathrm{},W_{N1}`$ where $`W_i`$ has radial boundary lines $`\mathrm{}_i`$ and $`\mathrm{}_{i+1}`$. We assume each $`W_i`$ has at most one Stokes line. Using excision together with the previous lemmas we get (1.6) $$\begin{array}{c}0H_2(\mathrm{\Delta },\delta \{0\};E,)_{i=0}^{N1}H_1(\mathrm{}_i,\{p_i,0\};E,)\hfill \\ \hfill \stackrel{\nu }{}_{i=0}^{N1}H_1(W_i,\delta _{W_i}\{0\};E,)H_1(\mathrm{\Delta },\delta \{0\};E,)0\end{array}$$ By lemma 1.3, the map $`\nu `$ above is given by $$\nu (e_0,\mathrm{},e_{N1})=(e_0e_1,e_1e_2,\mathrm{},e_{N1}e_0).$$ An element in the kernel of $`\nu `$ is thus a section $`e`$ of $`\|_{\mathrm{\Delta }\{0\}}`$ which has rapid decay along each $`\mathrm{}_i`$. Since each $`W_i`$ contains at most one Stokes line, such an $`e`$ would necessarily have rapid decay on every sector and thus would be trivial. This proves vanishing for $`H_2(\mathrm{\Delta },\delta ;E,)`$. Finally, to compute the dimension of $`H_1`$, note that if $`L_i`$ has a connection with pole of order $`m_i`$, then it has a horizontal section of the form $`e^f`$, where $`f`$ has a pole of order $`m_i1`$. (The connection is $`1df`$.) Suppose $`f=az^{1m_i}+\mathrm{}`$. Stokes lines for this factor are radial lines where $`az^{1m_i}`$ is pure imaginary. Thus, there are $`2(m_i1)`$ Stokes lines for this factor. Consider one of the Stokes lines, and suppose it lies in $`W_k`$. If the real part of $`az^{1m_i}`$ changes from negative to positive as we rotate clockwise through this line, say we are in case $`+`$, otherwise we are in case $``$. We have (1.7) $$dim(_\mathrm{}_k+_{\mathrm{}_{k+1}})dim_\mathrm{}_k=\{\begin{array}{cc}0\hfill & \text{case }+\hfill \\ dim(M_i)\hfill & \text{case },\hfill \end{array}$$ since the two cases alternate, we get a contribution of $`(m_i1)dim(M_i)`$. If $`m_i1`$ there are no rapidly decaying sections, so that case can be ignored. Summing over $`i`$ with $`m_i2`$ gives the desired result. Finally, we must consider the general case when the decomposition (1.5) is only available on $`\widehat{\pi _d^{}E}`$ for some $`d2`$. By a trace argument, vanishing of the homology upstairs, i.e. for $`\widehat{\pi _d^{}E}`$, in degrees $`1`$ implies vanishing downstairs. Since $`\pi _d:\mathrm{\Delta }\{0\}\mathrm{\Delta }\{0\}`$ is unramified, an Euler characteristic argument (or, more concretely, just cutting into small sectors over which the covering splits) shows that the Euler characteristic multiplies by $`d`$ under pullback, proving the lemma. ∎ In particular, we have now completed the proof of theorem 1.1. ∎ ## 2. de Rham Cohomology In this section, using differential forms, we construct the dual sequence to the homology sequence from theorem 1.1. (More precisely, we continue to work with $`E,`$, so the sequence we construct will be dual to the homology sequence with coefficients in $`E^{},^{}`$). Consider the diagram of complexes (2.1) $$\begin{array}{ccccccccc}0& & E(D)& & j_{}E_U& & j_{}E_U/E(D)& & 0\\ & & _{\mathrm{mero}}& & _{\mathrm{an}}& & _{\mathrm{an}/\mathrm{mero}}& & \\ 0& & E(D)\omega & & j_{}E_U\omega & & (j_{}E_U/E(D))\omega & & 0\end{array}$$ A result of Malgrange is that $`_{\mathrm{an}/\mathrm{mero}}`$ is surjective. Define $`N:=_iN_i=\mathrm{ker}(_{\mathrm{an}/\mathrm{mero}})`$. Since none of these sheaves has higher cohomology (by assumption $`D\mathrm{}`$) we get a 5-term exact sequence by taking global sections and applying the serpent lemma: (2.2) $$\begin{array}{c}0H_{DR}^0(U;E,)H^0(U,)N\hfill \\ \hfill H_{DR}^1(U;E,)H^1(U,)0\end{array}$$ ###### Theorem 2.1. Integration of forms over chains defines a perfect pairing between the exact sequence (2.2) and the exact sequence from theorem 1.1: (2.3) $$\begin{array}{c}0H_1(U,^{})H_1(X,D;E^{},^{})_iH_1(\mathrm{\Delta }_i,\delta _i\{x_i\};E^{},^{})\hfill \\ \hfill H_0(U,^{})H_0(X,D;E^{},^{})0.\end{array}$$ ###### Proof. To establish the existence of a pairing, note that if $`cϵ^{}`$ is a rapidly decaying chain and $`\eta `$ is a form of the same degree with moderate growth, then elementary estimates show $`_cϵ^{},\eta `$ is well defined. Suppose $`c:\mathrm{\Delta }^nX`$ and write $`\mathrm{\Delta }^n=lim_{t0}\mathrm{\Delta }_t^n`$ where $`\mathrm{\Delta }_t^n`$ denotes $`\mathrm{\Delta }^n\text{tubular neighborhood of radius }t\text{ around }\mathrm{\Delta }^n`$. Let $`c_t=c\|_{\mathrm{\Delta }_t^n}`$ and suppose $`\eta =d\tau `$ where $`\tau `$ has moderate growth also. Then (2.4) $$_cϵ^{},\eta =\underset{t0}{lim}_{c_t}ϵ^{},d\tau =\underset{t0}{lim}_{c_t}ϵ^{},\tau =_cϵ^{},\tau .$$ Note $`c`$ may include simplices mapping to $`D`$. Our definition (0.5) of $`𝒞_{}(X,D;E^{},^{})`$ factors these chains out. Thus, we do get a pairing of complexes. Of course, chains away from $`D`$ integrate with forms with possible essential singularities on $`D`$. To complete the description of the pairing, we must indicate a pairing (2.5) $$(,):N_i\times H_1(\mathrm{\Delta }_i,\delta _i\{x_i\};E^{},^{}).$$ To simplify notation we will drop the subscript $`i`$ and take $`x_i=0`$. An element in $`H_1`$ can be represented in the form $`ϵ^{}c`$ where $`c`$ is a radial path. Let $`c\delta =\{p\}`$. Given $`nN`$, choose a sector $`W`$ containing $`c`$ on which $``$ has a basis $`ϵ_i`$. By assumption, we can represent $`n=a_iϵ_i`$ with $`a_i`$ analytic on the open sector, such that (2.6) $$(a_iϵ_i)=ϵ_ida_i=e_i\eta _i$$ where $`e_i`$ from a basis of $`E`$ in a neighborhood of $`0`$ and $`\eta _i`$ are meromorphic $`1`$-forms at $`0`$. then by definition (2.7) $$(ϵ^{}c,n):=_c\underset{i}{}ϵ^{},e_i\eta _i\underset{i}{}ϵ^{},ϵ_ia_i(p).$$ The pairing is taken to be trivial on chains which do not contain $`0`$. If $`s`$ is a $`2`$-chain bounding two radial segments $`c`$ and $`c^{}`$ and a path along $`\delta `$ from $`p`$ to $`p^{}`$. Then Cauchy’s theorem (together with a limiting argument at $`0`$) gives (2.8) $$\begin{array}{c}0=_c\underset{i}{}ϵ^{},e_i\eta _i_c^{}\underset{i}{}ϵ^{},e_i\eta _i+_p^p^{}\underset{i}{}ϵ^{},ϵ_ida_i\hfill \\ \hfill =(ϵ^{}c,n)(ϵ^{}c^{},n).\end{array}$$ Similar arguments show the pairing independent of the choice of the radius of the disk. Also, if $`a_iϵ_i=b_ie_i`$ with $`b_i`$ meromorphic at $`0`$, then (2.9) $$\begin{array}{c}(ϵ^{}c,n)=_cϵ^{},e_i\eta _i\underset{i}{}ϵ^{},ϵ_ia_i(p)\hfill \\ \hfill =_cdϵ^{},b_ie_i\underset{i}{}ϵ^{},ϵ_ia_i(p)\\ \hfill =ϵ^{},b_ie_i(p)\underset{i}{}ϵ^{},ϵ_ia_i(p)=0\end{array}$$ It follows that the pairing is well defined. ###### Lemma 2.2. The diagrams $$\begin{array}{ccc}H_1(X,D;E^{},^{})& & H_1(\mathrm{\Delta }_i,\delta _i;E^{},^{})\\ \times & & \times \\ H^1(XD;E,)& & N_i\\ & & \text{}\\ & & \end{array}$$ and $$\begin{array}{ccc}H_1(\mathrm{\Delta }_i,\delta _i;E^{},^{})& & H_0(U,^{})\\ \times & & \times \\ N_i& & H^0(U,)\\ & & \text{}\\ & & \end{array}$$ commute. ###### proof of lemma. Consider the top square. The top arrow is excision, replacing a chain with the part of it lying in the disks $`\mathrm{\Delta }_i`$. The bottom arrow maps an $`n`$ as above in some $`N_i`$ to $`e_j\eta _j=ϵ_jda_j`$. Along $`c`$ outside the disks $`e_j\eta _j`$ is exact; its integral along the chain is a sum of terms of the form $`_iϵ^{},ϵ_ja_j(p_i)`$ where $`p_ic\delta _i`$. For the part of the chain inside the $`\mathrm{\Delta }_i`$ of course we must take $`_{c\mathrm{\Delta }_i}ϵ^{},e_j\eta _j`$. Combining these terms with appropriate signs yields the desired compatibility. For the bottom square, the top arrow associates to a relative chain on $`\mathrm{\Delta }_i`$ its boundary on $`\delta _iU`$. The bottom arrow associates to a horizontal section $`ϵ`$ on $`U`$ the corresponding element in $`N`$. Note here the $`a_j`$ will be constant so in the pairing with $`N`$ only the term $`ϵ^{},ϵ_ja_j(p)`$ survives. The assertion of the lemma follows. ∎ Returning to the proof of the theorem, we see it reduces to a purely local statement for a connection on a disk. In the following lemma, we modify notation, writing $`N`$ to denote the corresponding group for a connection on a disk $`\mathrm{\Delta }`$ with a meromorphic singularity at $`0`$. ###### Lemma 2.3. The pairing $$(,):N\times H_1(\mathrm{\Delta },\delta ;E^{},^{})$$ is nondegenerate on the left, i.e. $`(ϵ^{}c,n)=0`$ for all relative $`1`$-cycles implies $`n=0`$. ###### proof of lemma. We work in a sector and we suppose the basis $`ϵ_i`$ taken in the usual way compatible (in the sector) with the decomposition into a direct sum of rank $`1`$ irregular connections tensor regular singular point connections. Let $`ϵ_i^{}`$ be the dual basis. Fix an $`i`$ and suppose first $`ϵ_i`$ and $`ϵ_i^{}`$ both have moderate growth. We claim $`a_i`$ has moderate growth. For this it suffices to show $`da_i`$ has moderate growth. But (2.10) $$da_i=(n),ϵ_i^{}=\underset{j}{}e_j,ϵ_i^{}\eta _j.$$ This has moderate growth because, $`e_j,ϵ_i^{},\text{and }\eta _j`$ all do. Now assume $`(ϵ^{}c,n)=0`$ for all $`ϵ^{}cH_1`$. Fix an $`i`$ and assume $`ϵ_i^{}`$ is rapidly decreasing in our sector. Let $`c`$ be a radius in the sector with endpoint $`p`$. We can find (cf. , chap. IV, p.53-56) a basis $`t_i`$ of $`E`$ on the sector with moderate growth and such that $`t_i=\psi _iϵ_i`$, so $`t_i^{}=\psi _i^1ϵ_i^{}`$. We are interested in the growth of $`a_iϵ_i`$ along $`c`$. We have (2.11) $$a_i(p)ϵ_i(p)=\left(_c\underset{j}{}ϵ_i^{},e_j\eta _j\right)ϵ_i(p)=\left(_c\psi _i\underset{j}{}t_i^{},e_j\eta _j\right)\psi _i(p)^1t_i(p).$$ Asymptotically, taking $`y`$ the parameter along $`c`$, $`\psi _i(y)\mathrm{exp}(ky^N)`$ as $`y0`$ for some $`k>0`$ and some $`N1`$. We need to know the integral (2.12) $$\mathrm{exp}(kp^N)_0^py^M\mathrm{exp}(ky^N)𝑑y$$ has moderate growth as $`p0`$. Changing variables, so $`x=y^1,q=p^1,u=xq`$, this becomes (2.13) $$\begin{array}{c}_0^{\mathrm{}}(u+q)^{M2}\mathrm{exp}(q^N(u+q)^N)𝑑u\hfill \\ \hfill =_0^{\mathrm{}}(u+q)^{M2}\mathrm{exp}(u^Nqf(u,q))𝑑u,\end{array}$$ where $`f`$ is a sum of monomials in $`q`$ and $`u`$ with positive coefficients. Clearly this has at worst polynomial growth as $`q\mathrm{}`$ as desired. Finally, assume $`ϵ_i^{}`$ is rapidly increasing and $`ϵ_i`$ is rapidly decreasing. We have as above (2.14) $$\underset{i}{}e_j\eta _j=\underset{j}{}ϵ_jda_j=\underset{j}{}\psi _j^1t_jda_j$$ In particular, $`\psi _i^1da_i`$ has moderate growth. This implies $`a_iϵ_i=a_i\psi _i^1t_i`$ has moderate growth as well. Indeed, changing notation, this amounts to the assertion that if $`g`$ is rapidly decreasing and $`g\frac{df}{dz}`$ has moderate growth, then $`gf`$ has moderate growth. Fix a point $`p_0`$ with $`0<p<p_0`$. the mean value theorem says there exists an $`r`$ with $`prp_0`$ such that $$g(p)f(p)=g(p)(f(p_0)+(pp_0)f^{}(r))$$ Suppose $`\|f^{}(q)g(q)\|<<q^N`$. We get $$\|g(p)f(p)\|<<\|g(p)f^{}(r)\|\|g(r)f^{}(r)\|<<r^Np^N$$ proving moderate growth. We conclude that our representation for $`n`$ has moderate growth, and hence it is zero in $`N`$. It follows that the pairing $`N\times H_1`$ is nondegenerate on the left. ∎ Returning to the global situation, we have now $$dimN_idimH_1(\mathrm{\Delta }_i,\delta _i;E^{},^{}),$$ and to finish the proof of the theorem, it will suffice to show these dimensions are equal. ###### Lemma 2.4. With notation as above, $`dimN_i=dimH_1(\mathrm{\Delta }_i,\delta _i;E^{},^{})`$. ###### proof of lemma. It will suffice to compute the difference of the two Euler characteristics (2.15) $$\chi (U,)\chi _{DR}(U;E,).$$ It is straightforward to show this difference is invariant if $`U`$ is replaced by a smaller Zariski open set, and that the Euler characteristics are multiplied by the degree in a finite étale covering $`VU`$. Using lemma 1.4, we reduce to the case where formally locally at each $`x_iD`$ we have $`E\widehat{K}_{x_i}_jL_{ij}M_{ij}`$ with $`L_{ij}`$ rank $`1`$ and $`M_{ij}`$ at worst regular singular. (Here $`\widehat{K}_{x_i}`$ is the Laurent power series field at $`x_i`$). Let $`m_{ij}`$ be the degree of the pole for the connection on $`L_{ij}`$. Then one can find coherent sheaves $$F_2F_1E(D)$$ such that $$F_1/F_2_{ij}M_{ij}/M_{ij}(m_{ij}x_i)$$ $$E^{}H^0(F_2);(F_2)F_1\omega $$ $$H^0(F_1\omega )H_{DR}^1(U;E,)$$ It follows that, writing $`g=\text{genus}(X)`$ (2.16) $$\begin{array}{c}\chi _{DR}(U;E,)\hfill \\ \hfill =\chi (F_2)\chi (F_1\omega )=\text{rk}(E)(2g2)\underset{ij}{}m_{ij}dim(M_{ij}).\end{array}$$ Since (2.17) $$\chi (U,)=\text{rk}(E)(2g2+n),$$ (which is proven algebraically as above, replacing $``$ by the regular connection associated to $``$) it follows that $$\chi _{DR}(U;E,)\chi (U,)=\underset{ij}{}(m_{ij}1)dim(M_{ij})$$ Referring to lemma 1.4, we see that this is the desired formula. ∎ This completes the proof of the theorem. ∎
warning/0005/hep-th0005265.html	ar5iv	text	# 1 Introduction ## 1 Introduction The abelian Born-Infeld (BI) action is known to be the low-energy part of any (gauge-fixed, world-volume) D-brane effective action. The BI action is merely dependent upon the abelian vector field strength, being independent upon spacetime derivatives of the field strength by definition. The supersymmetetric BI actions with linearly realized N=1 or N=2 supersymmetry in four dimensions describe the low-energy (world-volume) dynamics of a single D3-brane propagating in four or six dimensions, respectively. <sup>3</sup><sup>3</sup>3See ref. and references therein for a review. The D3-brane action is supposed to have the Goldstone-Maxwell interpretation associated with partial (1/2) spontaneous supersymmetry breaking, with the Goldstone fields in a (Maxwell) vector supermultiplet with respect to unbroken (linearly realized) supersymmetry. Hence, a supersymmetric BI action should be the low-energy part of the corresponding Goldstone-Maxwell action. The unbroken N=1 and N=2 supersymmetry can be made manifest in superspace. The N=1 manifestly supersymmetric BI action was first formulated in ref. , while its Goldstone-Maxwell interpretation was later established in ref. . The N=2 manifestly supersymmetric BI action was first formulated in ref. , whereas its relation to (yet to be fully determined) N=2 Goldstone-Maxwell action is briefly discussed in sect. 2, see also ref. for more. As was pointed out by Witten , there is a non-abelian gauge symmetry enhancement when $`N`$ parallel D3-branes coincide. A supersymmetric abelian BI action is then supposed to be replaced by a Non-abelian Born-Infeld (NBI) action where the world-volume fields are valued in the Lie algebra of $`U(N)`$. Both abelian and non-abelian BI actions are the effective actions, defined modulo local field redefinitions. Nevertheless, the bosonic abelian BI action (with tension $`T_3`$) $$S_{\mathrm{BI}}=T_3d^4x\sqrt{det\left(\eta _{\mu \nu }+2\pi \alpha ^{}F_{\mu \nu }\right)}$$ $`(1)`$ is unambiguous, being only dependent of the abelian field strength $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu `$ but not of spacetime derivatives of it $`(F)`$. In contrast, a bosonic NBI action is not well-defined, while there are two principal sources for ambiguities . The first type of non-abelian ambiguities is related to the obvious fact that the terms dependent of the gauge-covariant derivatives of the non-abelian field strength cannot be unambiguously separated from the $`F`$-dependent commutators since $`D_\mu ,D_\nu F_{\lambda \rho }=F_{\mu \nu },F_{\lambda \rho }`$. Any concrete proposal for an NBI action has to specify an order of the $`F`$-matrices and, hence, it may effectively include some of the $`DF`$-dependent terms, even if they do not explicitly appear in the action. Though the full (abelian or non-abelian) D3-brane effective action certainly includes the derivative-dependent terms, it does not make much sense to keep some of them while ignoring other possible terms. Perhaps, the best one can do with a bosonic NBI action is to define it for almost covariantly constant gauge fields with almost commuting field strengths, which does not seem to be very illuminating. The second (related) type of ambiguities is connected to the trace operation over the gauge group. For example, when using the abelian identity $`(2\pi \alpha ^{}=b)`$ $$det\left(\eta _{\mu \nu }+bF_{\mu \nu }\right)=1+\frac{b^2}{2}F^2\frac{b^4}{16}(F\stackrel{~}{F})^2,\stackrel{~}{F}^{\mu \nu }=\frac{1}{2}\epsilon ^{\mu \nu \lambda \rho }F_{\lambda \rho },$$ $`(2)`$ one gets two natural candidates for the bosonic NBI action, $$S_{(\mathrm{a})}=T_3d^4x\sqrt{1+\frac{b^2}{2}\mathrm{tr}(F^2)\frac{b^4}{16}\mathrm{tr}(F\stackrel{~}{F})^2}$$ $`(3a)`$ and $$S_{(\mathrm{b})}=T_3d^4x\mathrm{Str}\sqrt{det(\eta _{\mu \nu }+bF_{\mu \nu })},$$ $`(3b)`$ where $`F_{\mu \nu }=F_{\mu \nu }^at_a`$, $`\{t_a\}`$ are the hermitian generators of the gauge group, $`t_a,t_b=if_{ab}^ct_c`$, $`\mathrm{tr}(t_at_b)=\delta _{ab}`$, and $`\mathrm{Str}`$ is the symmetrized trace, $$\mathrm{Str}\left(t_{a_1}\mathrm{}t_{a_k}\right)=\frac{1}{k!}\underset{\genfrac{}{}{0pt}{}{\mathrm{permutations}}{\pi }}{}\mathrm{tr}\left(t_{\pi (a_1)}\mathrm{}t_{\pi (a_k)}\right).$$ $`(4)`$ The $`F`$-matrices effectively commute under the symmetrized trace, so that the formal definition (2) of the determinant still applies in eq. (3b). It is not difficult to verify that the equations of motion in the NBI theory (3b) on self-dual (Euclidean) configurations $`(F_{\mu \nu }=\stackrel{~}{F}_{\mu \nu })`$ coincide with the ordinary Yang-Mills equations, so that they have the same BPS solutions , though the existence of a BPS bound is not obvious in the non-abelian case. Away from self-dual configurations the action (3a) is much simpler than (3b), while it is also known to admit solitionic (glueball) solutions . The gauge-invariant actions (3a) and (3b) are obviously different, so that further resolution requirements are needed. Some extra conditions are provided by string theory, because the BI action is well-known to represent the effective action of slowly varying gauge fields in open string theory. The most basic requirement of string theory is the overall single trace of the non-abelian gauge field strength products . The overall symmetrized trace advocated by Tseytlin is a stronger condition based on the observation that it reproduces the $`F^4`$-terms in the non-abelian effective action of open superstrings in ten dimensions . In this Letter I show that adding supersymmetry unexpectedly gives rise to some more constraints on supersymmetric NBI actions in four dimensions, which are not apparent in the bosonic case. At first sight, it seems to be straightforward to supersymmetrize any NBI action, so that supersymmetry would not add anything new towards its intrinsic definition. However, in fact, supersymmetry does tell us something more about the BI actions. For example, linearly realized supersymmetry apparently prefers the parametrization of the abelian BI actions in terms of the (anti)self-dual combinations, $`F^\pm =\frac{1}{2}(F\pm \stackrel{~}{F})`$, rather than in terms of the naively expected tensors $`F`$ and $`\stackrel{~}{F}`$. More importantly, it is the spontaneously broken (non-linearly realized) supersymmetry on top of the unbroken (linearly realized) supersymmetry that is responsible for the complicated non-linear structure of the D3-brane action to be considered as the Goldstone-Maxwell action. Hence, the non-linear structure of the supersymmetric abelian BI actions should also be dictated by similar features. Though a Goldstone interpretation of the supersymmetric NBI actions is far from being obvious, if any, the well-established Goldstone form of the N=1 supersymmetric abelian BI action in superspace gives us the natural starting point for a construction of its generalizations, either non-abelian or with extended (linearly realized) supersymmetry. The paper is organized as follows. In sect. 2 the N=1 and N=2 supersymmetric abelian BI actions are reviewed by emphasizing their relation to the Goldstone-Maxwell actions. In sect. 3 the non-abelian generalizations of the BI actions are proposed in N=1 and N=2 superspace. Sect. 4 is my conclusion. ## 2 N=1 and N=2 abelian BI actions The abelian bosonic BI Lagrangian $`_{\mathrm{BI}}(F)`$ can be thought of as the unique non-linear generalization of the Maxwell Lagrangian, $`\frac{1}{4}F^2`$, under the conditions of preservation of causality, positivity of energy, and electric-magnetic duality. In particular, the duality invariance of an abelian Lagrangian $`(F)`$ amounts to the constraint $$G_{\mu \nu }\stackrel{~}{G}^{\mu \nu }+F_{\mu \nu }\stackrel{~}{F}^{\mu \nu }=0,\mathrm{where}\stackrel{~}{G}^{\mu \nu }=\frac{1}{2}\epsilon ^{\mu \nu \lambda \rho }G_{\lambda \rho }=2\frac{}{F_{\mu \nu }}.$$ $`(5)`$ Supersymmetry is known to be consistent with all these physical properties, so that the supersymmetric abelian BI actions enjoy similar features. The N=1 supersymmetric abelian BI action reads <sup>4</sup><sup>4</sup>4The deformation parameter $`b`$ is set to be one. The dependence upon $`b`$ can be easily restored for dimensional reasons. I also ignore $`T_3`$ for simplicity. $$S_{1\mathrm{B}\mathrm{I}}=\frac{1}{2}(d^4xd^2\theta W^2+\mathrm{h}.\mathrm{c}.)+d^4xd^4\theta Y(K,\overline{K})W^2\overline{W}^2,$$ $`(6)`$ where the structure function $`Y`$ is given by $$Y(K,\overline{K})=\frac{1}{1\frac{1}{2}(K+\overline{K})+\sqrt{1(K+\overline{K})+\frac{1}{4}(K\overline{K})^2}},$$ $`(7)`$ and $$K=\frac{1}{2}D^2W^2,D^2=D^\alpha D_\alpha ,W^2=W^\alpha W_\alpha ,\overline{W}^2=\overline{W}_\stackrel{_{{}_{}{}^{}}}{\alpha }\overline{W}^\stackrel{_{{}_{}{}^{}}}{\alpha },$$ $`(8)`$ in terms of the abelian N=1 chiral spinor superfield strength $`W_\alpha `$, $`\alpha =1,2`$, satisfying the off-shell superspace constraints (N=1 Bianchi identities) $$\overline{D}_\stackrel{_{{}_{}{}^{}}}{\alpha }W_\alpha =0,D^\alpha W_\alpha =\overline{D}_\stackrel{_{{}_{}{}^{}}}{\alpha }\overline{W}^\stackrel{_{{}_{}{}^{}}}{\alpha }.$$ $`(9)`$ The action (6) can be rewritten in the ‘non-linear sigma-model’ form $$S_{1\mathrm{B}\mathrm{I}}=d^4xd^2\theta \mathrm{\Phi }+\mathrm{h}.\mathrm{c}.,$$ $`(10)`$ where the N=1 chiral Lagrangian $`\mathrm{\Phi }`$ is the perturbative solution to the non-linear superfield constraint $$\mathrm{\Phi }=\frac{1}{2}\mathrm{\Phi }\overline{D}^2\overline{\mathrm{\Phi }}+\frac{1}{2}W^2.$$ $`(11)`$ It is worth mentioning that the constraint (11) is Gaussian in $`\mathrm{\Phi }`$, while its perturbative solution is unambiguously constructed in superspace by iterations. In fact, the simple constraint (11) is most useful in proving the invariance of the action $`S_{1\mathrm{B}\mathrm{I}}`$ under the second (non-linearly realized or spontaneously broken) supersymmetry with the rigid anticommuting spinor parameter $`\eta ^\alpha `$ , $$\delta _2\mathrm{\Phi }=\eta ^\alpha W_\alpha ,\delta _2W_\alpha =\eta _\alpha \left(1\frac{1}{2}\overline{D}^2\overline{\mathrm{\Phi }}\right)+i\overline{\eta }^\stackrel{_{{}_{}{}^{}}}{\alpha }_{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }}\mathrm{\Phi },$$ $`(12)`$ where the second equation follows from the first one after the use of eqs. (9) and (11). The constraint (11) generating the full action (6) is also quite useful in proving the electric-magnetic self-duality of $`S_{1\mathrm{B}\mathrm{I}}`$. The duality invariance amounts to another non-local constraint $$d^4xd^2\theta (W^2+M^2)=d^4xd^2\overline{\theta }(\overline{W}^2+\overline{M}^2),\frac{i}{2}M_\alpha =\frac{\delta S_{1\mathrm{B}\mathrm{I}}}{\delta W^\alpha },$$ $`(13)`$ which is the straightforward N=1 generalization of eq. (5). Thus, $`S_{1\mathrm{B}\mathrm{I}}=S_{\mathrm{GM}:\mathrm{N}=2/\mathrm{N}=1}`$. Similarly, the N=2 supersymmetric Abelian BI action in N=2 superspace reads $$S_{2\mathrm{B}\mathrm{I}}=\frac{1}{2}d^4xd^4\theta 𝒲^2+\frac{1}{4}d^4xd^8\theta 𝒴(𝒦,\overline{𝒦})𝒲^2\overline{𝒲}^2,$$ $`(14)`$ with the same structure function $$𝒴(𝒦,\overline{𝒦})=\frac{1}{1\frac{1}{2}(𝒦+\overline{𝒦})+\sqrt{1(𝒦+\overline{𝒦})+\frac{1}{4}(𝒦\overline{𝒦})^2}},$$ $`(15)`$ but $$𝒦=D^4𝒲^2,D^4=\underset{i,\alpha }{}D_\alpha ^i=\frac{1}{12}D_{ij}D^{ij},D_{ij}=D_i^\alpha D_{\alpha j}=D_{ji},i=1,2,$$ $`(16)`$ in terms of the N=2 restricted chiral gauge superfield strength $`𝒲`$ satisfying the off-shell constraints (N=2 Bianchi identities) $$\overline{D}_i^\stackrel{_{{}_{}{}^{}}}{\alpha }𝒲=0,D^4𝒲=\mathrm{}\overline{𝒲}.$$ $`(17)`$ The action (14) can be rewritten (modulo $`W`$-dependent terms) in the ‘non-linear sigma-model’ form $$S_{2\mathrm{B}\mathrm{I}}=\frac{1}{4}d^4xd^4\theta 𝒳+\mathrm{h}.\mathrm{c}.,$$ $`(18)`$ whose N=2 chiral Lagrangian $`𝒳`$ satisfies the non-linear N=2 superfield constraint $$𝒳=\frac{1}{4}𝒳\overline{D}^4\overline{𝒳}+𝒲^2.$$ $`(19)`$ Similarly to the N=1 abelian BI action, the non-linear constraint (19) gives us the convenient way of handling the complicated N=2 BI abelian action (14). For example, as was demonstrated in ref. , electric-magnetic self-duality of an N=2 action $`S(𝒲,\overline{𝒲})`$ amounts to the following N=2 supersymmetric extension of the N=1 non-local constraint (13): $$d^4xd^4\theta (𝒲^2+^2)=d^4xd^4\overline{\theta }(\overline{𝒲}^2+\overline{}^2),\frac{i}{4}=\frac{\delta S}{\delta 𝒲},$$ $`(20)`$ while it appears to be satisfied in the case of $`S_{2\mathrm{B}\mathrm{I}}`$ defined by eqs. (18) and (19). Unlike its N=1 BI counterpart, the N=2 BI action (14) or (18) does not give the full N=2 Goldstone-Maxwell action, but rather represents its low-energy part, i.e. $`S_{GM:N=4/N=2}=S_{2\mathrm{B}\mathrm{I}}+𝒪(W,\overline{W})`$. The infinitesimal parameters of the spontaneously broken (non-linearly realised) rigid symmetries can be naturally unified into a single (spacetime-independent) superfield $`\mathrm{\Lambda }=\lambda +\theta _i^\alpha \lambda _\alpha ^i+\theta _{ij}\lambda ^{ij}`$, where $`\lambda `$ is the complex parameter of the Peccei-Quinn-type symmetry associated with two spontaneously broken translations (from the the viewpoint of a D3-brane propagating in six dimensions), $`\lambda _\alpha ^i`$ are the spinor parameters of two spontaneously broken supersymmetries, whereas $`\lambda ^{ij}`$ are the parameters of spontaneously broken R-symmetry $`SU(2)`$. The natural Ansatz for the transformation laws of the non-linearly realised symmetries is given by an N=2 analogue of eq.(12) as follows : $$\delta _2𝒳=2\mathrm{\Lambda }𝒲,\delta _2𝒲=\mathrm{\Lambda }\left(1\frac{1}{4}\overline{D}^4\overline{𝒳}\right)+\mathrm{},$$ $`(21)`$ where $`𝒳`$ is the perturbative solution to the non-linear constraint (19) by iterations, to all orders in $`𝒲`$ and $`\overline{𝒲}`$, while the dots stand for some $`W`$-dependent terms needed for consistency with the second Bianchi identity (17). A variation of the N=2 BI action (18) under the non-linear transformations (21) does not vanish, but it appears to be only dependent upon higher spacetime derivatives of $`W`$ and $`\overline{W}`$. This may not be surprising since the BI action and its supersymmetric extensions were defined modulo such terms. However, it also means that the N=2 BI action has to be modified, order by order in $`W`$ and $`\overline{W}`$, in order to get the full N=2 Goldstone-Maxwell action. A derivation of the derivative-dependent terms is beyond the scope of this Letter, and we do not need them for our main purpose formulated in the title. A manifestly N=4 supersymmetric abelian BI action is not known (see, however, ref. and references therein for some partial results). ## 3 N=1 and N=2 supersymmetric NBI actions Having understood the fact that the simple non-linear constraints (11) and (19) fully determine the structure of the highly complicated abelian BI actions (6) and (14), respectively, it is natural to define the N=1 and N=2 supersymmetric NBI actions by non-abelian generalizations of eq. (11) and (19). The non-abelian (Yang-Mills) N=1 chiral superfield strength is given by the well-known formula (see, e.g., ref. for a review or an introduction) <sup>5</sup><sup>5</sup>5All superfields are now Lie algebra-valued. $$W_\alpha =\frac{1}{8}\overline{D}^2\left(e^{2V}D_\alpha e^{2V}\right),\overline{D}_\stackrel{_{{}_{}{}^{}}}{\alpha }W_\alpha =0,$$ $`(22)`$ where the real scalar gauge superfield potential $`V`$ transforms under gauge transformations with the chiral paramater $`\mathrm{\Lambda }(x,\theta ,\overline{\theta })`$ in the standard way : $$e^{2V}e^{2i\overline{\mathrm{\Lambda }}}e^{2V}e^{2i\mathrm{\Lambda }},\overline{D}_\stackrel{_{{}_{}{}^{}}}{\alpha }\mathrm{\Lambda }=0,$$ $`(23)`$ so that $`W_\alpha `$ and $`\overline{W}_\stackrel{_{{}_{}{}^{}}}{\alpha }`$ transform covariantly, viz. $$W_\alpha e^{2i\mathrm{\Lambda }}W_\alpha e^{2i\mathrm{\Lambda }},\overline{W}_\stackrel{_{{}_{}{}^{}}}{\alpha }e^{2i\overline{\mathrm{\Lambda }}}\overline{W}_\stackrel{_{{}_{}{}^{}}}{\alpha }e^{2i\overline{\mathrm{\Lambda }}}.$$ $`(24)`$ The non-abelian gauge-covariant generalization of eq. (11) is given by $$\mathrm{\Phi }=\frac{1}{2}\mathrm{\Phi }\overline{D}^2\left(e^{2V}\overline{\mathrm{\Phi }}e^{2V}\right)+\frac{1}{2}W^2,$$ $`(25)`$ where $`\mathrm{\Phi }`$ is the N=1 chiral superfield Lagrangian that transforms like $`W_\alpha `$ under the gauge transformations. The invariant action reads $$S_{1\mathrm{N}\mathrm{B}\mathrm{I}}=d^4xd^2\theta \mathrm{tr}\mathrm{\Phi }+\mathrm{h}.\mathrm{c}.,$$ $`(26a)`$ or $$S_{1\mathrm{N}\mathrm{B}\mathrm{I}}^{(s)}=d^4xd^2\theta \mathrm{Str}\mathrm{\Phi }+\mathrm{h}.\mathrm{c}.$$ $`(26b)`$ The NBI actions (26a) and (26b) are supersymmetric and gauge-invariant, while they both have the single overall trace. The symmetrized trace in eq. (26b) is supposed to be applied to the gauge-covariant operators only, by definition. It is instructive to take a look at the structure of the quartic $`(F^4)`$ terms in the actions (26), which arise from the standard ‘adjoint chiral matter’ term, $$d^4xd^4\theta (\mathrm{S})\mathrm{tr}\varphi e^{2V}\overline{\varphi }e^{2V},\varphi =W^2.$$ $`(27)`$ It is straightforward to verify that taking the trace as in eq. (26a) results in the non-abelian generalization of the Euler-Heisenberg Lagrangian, in the bosonic sector of the component expansion of the action (26a), $$\frac{1}{4}\left[\mathrm{tr}(F^2)^2+\mathrm{tr}(F\stackrel{~}{F})^2\right].$$ $`(28)`$ In contrast, taking the symmetrized trace, as in eq. (26b), exactly yields the $`F^4`$-terms appearing in the expansion of the bosonic NBI Lagrangian (3b) . Hence, if one insists on the choice (3b) of the bosonic NBI action, its supersymmetric extension in compact form is provided by eq. (26b) — cf. ref. . Supersymmetry alone does not provide a resolution between the two different actions (26a) and (26b), so that more physical input is apparently needed. One natural option is to restrict the gauge group $`G`$ to its (abelian) Cartan subgroup and then impose the condition of the generalized (non-abelian) self-duality for the resulting supersymmetric field theory of rank $`G`$ (abelian) gauge fields along the lines of ref. . As was argued in ref. , this may ultimately support the symmetrized trace. Still, in the absence of more physical reasons, the ordinary trace in eq. (26a) is much simpler, being dependent of only two matrix building blocks, $`W^2`$ and $`\overline{W}^2`$ (or $`F^2`$ and $`F\stackrel{~}{F}`$), and their covariant derivatives (see below). The action (3a) does not seem to have a nice supersymmetric generalization. It is possible to rewrite the action (26) into the manifestly gauge-invariant and N=1 supersymmetric form, by using the N=1 supersymmetric gauge-covariant derivatives in superspace, which satisfy the standard N=1 super-Yang-Mills constraints $$\{_\alpha ,_\beta \}=\{\overline{}_\stackrel{_{{}_{}{}^{}}}{\alpha },\overline{}_\stackrel{_{{}_{}{}^{}}}{\beta }\}=0,\{_\alpha ,\overline{}_\stackrel{_{{}_{}{}^{}}}{\beta }\}=2i_{\alpha \stackrel{_{{}_{}{}^{}}}{\beta }},$$ $$_\alpha ,_{\beta \stackrel{_{{}_{}{}^{}}}{\beta }}=2i\epsilon _{\alpha \beta }\widehat{\overline{W}}_\stackrel{_{{}_{}{}^{}}}{\beta },\overline{}_\stackrel{_{{}_{}{}^{}}}{\alpha },_{\beta \stackrel{_{{}_{}{}^{}}}{\beta }}=2i\epsilon _{\stackrel{_{{}_{}{}^{}}}{\alpha }\stackrel{_{{}_{}{}^{}}}{\beta }}\widehat{W}_\beta ,$$ $`(29)`$ where $`\widehat{W}_\alpha `$ is the N=1 covariantly-chiral gauge superfield strength, $$\overline{}_\stackrel{_{{}_{}{}^{}}}{\alpha }\widehat{W}_\alpha =0,\overline{}_\stackrel{_{{}_{}{}^{}}}{\alpha }\widehat{\overline{W}}{}_{}{}^{\stackrel{_{{}_{}{}^{}}}{\alpha }}=^\alpha \widehat{W}_\alpha .$$ $`(30)`$ Equation (26) then takes the form $$S_{1\mathrm{N}\mathrm{B}\mathrm{I}}=d^4xd^2\theta (\mathrm{S})\mathrm{tr}\widehat{\mathrm{\Phi }}+\mathrm{h}.\mathrm{c}.,$$ $`(31)`$ where the N=1 covariantly-chiral Lagrangian $`\widehat{\mathrm{\Phi }}`$ is the perturbative (iterative) solution to the manifestly gauge-covariant and supersymmetric nonlinear superfield constraint $$\widehat{\mathrm{\Phi }}=\frac{1}{2}\widehat{\mathrm{\Phi }}\overline{}^2\widehat{\overline{\mathrm{\Phi }}}+\frac{1}{2}\widehat{W}^2.$$ $`(32)`$ It is not difficult to generalize eqs. (31) and (32) further to the case of N=2 supersymmetry, by doing a similar construction in N=2 superspace. The standard N=2 superspace constraints, defining the off-shell N=2 supersymmetric Yang-Mills theory, are given by \| $`\{\overline{}_{\stackrel{_{{}_{}{}^{}}}{\alpha }i},\overline{}_{\stackrel{_{{}_{}{}^{}}}{\beta }j}\}=2\epsilon _{\stackrel{_{{}_{}{}^{}}}{\alpha }\stackrel{_{{}_{}{}^{}}}{\beta }}\epsilon _{ij}\widehat{𝒲},`$ \| $`\{_\alpha ^i,_\beta ^j\}=2\epsilon _{\alpha \beta }\epsilon ^{ij}\widehat{\overline{𝒲}},`$ \| \| --- \| --- \| \| $`_{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }},\overline{}_{\stackrel{_{{}_{}{}^{}}}{\beta }i}=i\epsilon _{\stackrel{_{{}_{}{}^{}}}{\alpha }\stackrel{_{{}_{}{}^{}}}{\beta }}_{\alpha i}\widehat{𝒲},`$ \| $`_{\alpha \stackrel{_{{}_{}{}^{}}}{\alpha }},_{}^{i}{}_{\beta }{}^{}=i\epsilon _{\alpha \beta }\overline{}_\stackrel{_{{}_{}{}^{}}}{\alpha }^i\widehat{\overline{𝒲}},`$ \| $$\{_{}^{i}{}_{\alpha }{}^{},\overline{}_{\stackrel{_{{}_{}{}^{}}}{\beta }j}\}=2i\delta ^i{}_{j}{}^{}_{\alpha \stackrel{_{{}_{}{}^{}}}{\beta }}^{},$$ $`(33)`$ where the non-abelian N=2 gauge superfield strength $`\widehat{𝒲}`$ obeys the off-shell constraints (N=2 Bianchi identities) $$\overline{}_{\stackrel{_{{}_{}{}^{}}}{\alpha }i}\widehat{𝒲}=0\mathrm{and}_{ij}\widehat{𝒲}=\overline{}_{ij}\widehat{\overline{𝒲}}.$$ $`(34)`$ I use the following book-keeping notation: \| $`_{ij}=_{(i}^\alpha _{j)\alpha }=_{ji},`$ \| $`\overline{}_{ij}=\overline{}_{\stackrel{_{{}_{}{}^{}}}{\alpha }(i}\overline{}_{}^{\stackrel{_{{}_{}{}^{}}}{\alpha }}{}_{j)}{}^{}=\overline{}_{ji},`$ \| \| --- \| --- \| \| $`_{\alpha \beta }=_{i(\alpha }_{\beta )}^i=_{\beta \alpha },`$ \| $`\overline{}_{\stackrel{_{{}_{}{}^{}}}{\alpha }\stackrel{_{{}_{}{}^{}}}{\beta }}=\overline{}_{i(\stackrel{_{{}_{}{}^{}}}{\alpha }}\overline{}_{\stackrel{_{{}_{}{}^{}}}{\beta })}^i=\overline{}_{\stackrel{_{{}_{}{}^{}}}{\beta }\stackrel{_{{}_{}{}^{}}}{\alpha }},`$ \| $`(35)`$ where all symmetrizations have unit weight. The non-abelian generalization of the $`\overline{D}^4`$ operator, which converts the (covariantly) anti-chiral N=2 superfields into (covariantly) chiral N=2 superfields, is most easily (and unambiguously) identified in the $`SL(4,𝐂)`$ notation of ref. , by combining fundamental $`SL(2,𝐂)`$ and $`SU(2)`$ indices into a single (fundamental) $`SL(4,𝐂)`$ index $`\underset{¯}{a}=(\stackrel{_{{}_{}{}^{}}}{\alpha },i)=1,2,3,4`$. The $`\overline{}`$-algebra of eq. (33) in the $`SL(4,𝐂)`$ notation takes the familiar Dirac-type-form $$\{\overline{}_{\underset{¯}{a}},\overline{}_{\underset{¯}{b}}\}=2C_{\underset{¯}{a}\underset{¯}{b}}\widehat{𝒲},\overline{}_{\underset{¯}{a}}\widehat{𝒲}=0,$$ $`(36)`$ with the constant metric $`C`$, $`C^2=1`$ and $`C^T=C`$. The desired gauge-covariant operator is just given by the ‘$`\gamma _5`$-type’ top product $$\overline{}^4=\frac{1}{4!}\epsilon ^{\underset{¯}{a}\underset{¯}{b}\underset{¯}{c}\underset{¯}{d}}\overline{}_{\underset{¯}{a}}\overline{}_{\underset{¯}{b}}\overline{}_{\underset{¯}{c}}\overline{}_{\underset{¯}{d}}.$$ $`(37)`$ In the notation (35) it reads $$\overline{}^4=\frac{1}{24}\left(\overline{}_{ij}\overline{}^{ij}\overline{}_{\stackrel{_{{}_{}{}^{}}}{\alpha }\stackrel{_{{}_{}{}^{}}}{\beta }}\overline{}^{\stackrel{_{{}_{}{}^{}}}{\alpha }\stackrel{_{{}_{}{}^{}}}{\beta }}\right)\frac{2}{3}\widehat{𝒲}^2.$$ $`(38)`$ The N=2 supersymmetric non-abelian Born-Infeld action is given by $$S_{2\mathrm{N}\mathrm{B}\mathrm{I}}=\frac{1}{4}d^4xd^4\theta (\mathrm{S})\mathrm{tr}\widehat{𝒳}+\mathrm{h}.\mathrm{c}.,$$ $`(39)`$ whose N=2 covariantly chiral Lagrangian $`\widehat{𝒳}`$ is the perturbative (iterative) solution to the N=2 superfield constraint $$\widehat{𝒳}=\frac{1}{4}\widehat{𝒳}\overline{}^4\widehat{\overline{𝒳}}+\widehat{𝒲}^2.$$ $`(40)`$ ## 4 Conclusion The proposed N=1 and N=2 supersymmetric NBI actions in components contain only even powers of $`F`$, while they reduce to the known super-Born-Infeld actions in the abelian case. Both actions enjoy ‘auxiliary freedom’ by keeping the auxiliary fields $`\stackrel{}{D}`$ (in the Wess-Zumino gauge) away from propagation, with $`\stackrel{}{D}=0`$ being a solution to their equations of motion. Unlike the supersymmetric abelian BI actions (sect. 2), their supersymmetric non-abelian counterparts (sect. 3) are dependent of the gauge superfields not only via their gauge superfield strengths but also directly (via the gauge-covariant derivatives). This does not allow us to extend the notion of abelian electric-magnetic duality to the supersymmetric NBI actions. Similarly, it is unclear to us whether our supersymmetric non-abelian BI actions admit any Goldstone-Yang-Mills interpretation. It would be also interesting to investigate the structure of BPS solutions to the new supersymmetric NBI actions and find a precise relation between these actions and the non-abelian Dirac-Born-Infeld actions describing clusters of D3-branes with ‘deformed’ (non-linear) supersymmetry. A connection to noncommutative geometry seems to exist along the lines of ref. too. ## Acknowledgement I thank S. J. Gates Jr. and E. A. Ivanov for discussions. I also acknowledge kind hospitality extended to me at the University of Maryland in College Park during preparation of this paper.
warning/0005/astro-ph0005224.html	ar5iv	text	# limit-from𝛽-model and cooling flows in X–ray clusters of galaxies ## 1 INTRODUCTION To constrain the physical parameters of extended X-ray sources (e.g. groups and clusters of galaxies), the observed surface brightness can be either geometrically deprojected or, more simply, fitted with a model obtained from an assumed distribution of the gas density. Given the hydrostatic equilibrium within the cluster, the gravitational potential supports both the gas and the galaxies distribution. If the latter is approximated via the King approximation (1962) to the inner portions of an isothermal sphere (Lane-Emden equation in Binney & Tremaine 1987), the gas density is then written as: $$\rho _{\mathrm{gas}}=\rho _0(1+x^2)^{3\beta /2},$$ (1) where $`x=r/r_\mathrm{c}`$ and $`r_\mathrm{c}`$ is the core radius of the distribution. The surface brightness profile observed at the projected radius $`b`$, $`S(b)`$, is the projection on the sky of the plasma emissivity, $`ϵ(r)`$: $$S(b)=_{b^2}^{\mathrm{}}\frac{ϵdr^2}{\sqrt{r^2b^2}}.$$ (2) The emissivity is equal to $$ϵ(r)=\mathrm{\Lambda }(T_{\mathrm{gas}})n_\mathrm{p}^2\mathrm{erg}\mathrm{s}^1\mathrm{cm}^3,$$ (3) where $`n_\mathrm{p}=\rho _{\mathrm{gas}}/(2.21\mu m_\mathrm{p})`$ is the proton density and the cooling function, $`\mathrm{\Lambda }(T_{\mathrm{gas}})`$, depends upon the mechanism of the emission (mainly due to bremsstrahlung at $`T_{\mathrm{gas}}>2.5`$ keV). Assuming isothermality and a $`\beta `$-model for the gas density (eq. 1), the surface brightness profile has an analytic solution (eq. 3.196.2 in Gradshteyn and Ryzhik 1965): $`S(b)`$ $`=`$ $`n_0^2r_\mathrm{c}\mathrm{\Lambda }(T_{\mathrm{gas}})B(3\beta 0.5,0.5)\left[1+\left({\displaystyle \frac{b}{r_\mathrm{c}}}\right)^2\right]^{0.53\beta }`$ (4) $`=`$ $`S_0(1+x^2)^{0.53\beta },`$ where the validity of the beta function $`B(a,b)`$ puts the strict constraint $`3\beta >0.5`$ and the cooling function $`\mathrm{\Lambda }(T_{\mathrm{gas}})`$ does not change radially. The $`\beta `$model (Cavaliere & Fusco-Femiano 1976, 1978) provides a good representation of the observed surface brightness and has the advantage to easily constrain the gas density distribution. Elsewhere (Ettori 2000), I consider the effect of the presence of a temperature gradient in the estimate of the $`\beta `$model parameters. In this paper, I will focus on the deficiency of the $`\beta `$model in modeling in a satisfactory way the central emission from cooling flow clusters of galaxies. The cooling flows (e.g. Sarazin 1988, Fabian 1994) result in an enhancement of the gas density in the central region due to the high cooling efficiency in the cluster core. Recently, there have been attempts to model this excess in emission with a generic double $`\beta `$model, i.e. the sum of two components of surface brightness (Ikebe et al. 1996, Xu et al. 1998, Mohr et al. 1999, Reiprich & Böhringer 1999). The correlation between the presence of a cooling flow and the necessity for a second $`\beta `$model is well indicated from this figure: Peres et al. (1998) find that 40 per cent of a flux limited sample of 55 clusters of galaxies has a deposition rate of more than 100 $`M_{}`$ yr<sup>-1</sup>; this is the same percentage of the clusters in the Mohr et al. sample that are better modeled with a double $`\beta `$model instead of a single one (18 out of 45 clusters). The double $`\beta `$model as sum of two $`S(b)`$ in eq. 4, however, is not physically meaningful. In fact, a single isothermal temperature is usually assumed for both different density components that, therefore, are not in equilibrium. Moreover, data with high spatial resolution do not show evidence of a second inner core radius. I present in this work a simple geometrical and physical model of the emission from cooling flow clusters of galaxies. This model relies on recent spectral evidence that the cluster plasma can be described as a gas with two phases, one related to the cooling gas and the other to the ambient medium. Assuming that the extended intracluster gas density, $`n_{\mathrm{gas}}`$, is well described by a $`\beta `$model, I show in the following section that an analytic expression for $`S(b)`$ can be obtained to describe the surface brightness from cooling flow clusters of galaxies. In Section 3, I apply this model to real data of clusters with or without cooling flows. This model allows to handle the emissivity due to each component. I discuss the physical implications of this in Section 4. In Section 5, I present some concluding remarks. ## 2 The two–phases emission model Recent spectral analyses of cooling flow clusters of galaxies (Allen et al. 2000, White 2000) have shown how the spectral capabilities of the present instruments are unable to resolve all the fine structures of a multiphase gas, allowing just a modeling with a two-phase component, one that describes the emission from the central cooling gas and the other that takes into account the extended emission from the ambient medium. These observational results provide us with a simple and natural model for the total cluster emission: an inner cold phase confined within $`r=r_{\mathrm{cool}}`$ and overlapping the diffuse, ambient gas (see Fig. 1). We assume that the two components coexist within $`r_{\mathrm{cool}}`$, whereas only the ambient plasma fills the cluster volume shell at radius above $`r_{\mathrm{cool}}`$. The total cluster emissivity is then $`ϵ(r,T)=ϵ_{\mathrm{cool}}+ϵ_{\mathrm{amb}}`$, where $$ϵ=\{\begin{array}{c}\mathrm{\Lambda }(T_{\mathrm{cool}})n_{\mathrm{p},\mathrm{cool}}(r)^2+\mathrm{\Lambda }(T_{\mathrm{amb}})n_{\mathrm{p},\mathrm{amb}}(r)^2,r<r_{\mathrm{cool}}\hfill \\ \mathrm{\Lambda }(T_{\mathrm{amb}})n_{\mathrm{p},\mathrm{amb}}(r)^2,r>r_{\mathrm{cool}}\hfill \end{array}$$ (5) from the definition in eq. 3. This simple model provides an analytic expression for the surface brightness profile defined in eq. 2: $`S(b)`$ $`=`$ $`{\displaystyle _{b^2}^{\mathrm{}}}{\displaystyle \frac{ϵdr^2}{\sqrt{r^2b^2}}}`$ (6) $`=`$ $`{\displaystyle _{b^2}^{r_{\mathrm{cool}}^2}}{\displaystyle \frac{ϵ_{\mathrm{cool}}dr^2}{\sqrt{r^2b^2}}}+{\displaystyle _{b^2}^{\mathrm{}}}{\displaystyle \frac{ϵ_{\mathrm{amb}}dr^2}{\sqrt{r^2b^2}}}`$ $`=`$ $`S_{\mathrm{cool}}(b)+S_{\mathrm{amb}}(b)`$ where the integration limits in $`S_{\mathrm{cool}}(b)`$ contains the boundary of the inner region at $`r=r_{\mathrm{cool}}`$. Now, I integrate the emissivity along the line of sight. $`S_{\mathrm{amb}}(b)`$ is still eq. 4. To integrate $`S_{\mathrm{cool}}(b)`$, one needs the assumption that the only scale parameter of the gas density is the dimension of the cooling region, $`r_{\mathrm{cool}}`$. Considering that we are in the regime $`(r/r_{\mathrm{cool}})<1`$, I can move the sign ’–’ from the exponent to the radix and derive a $`\beta `$model in the from of $$n_{\mathrm{p},\mathrm{cool}}=n_{0,\mathrm{cool}}\left[1\left(\frac{r}{r_{\mathrm{cool}}}\right)^2\right]^{1.5\beta _{\mathrm{cool}}}.$$ (7) The behavior of this profile ensures that the gas density within the cooling region has no other parameter scale than the dimension of the region itself and $`ϵ_{\mathrm{cool}}`$ goes to zero when $`rr_{\mathrm{cool}}`$. Then, I can integrate analytically $`S_{\mathrm{cool}}(b)`$ (eq. 3.196.3 in Gradshteyn and Ryzhik 1965): $`S_{\mathrm{cool}}(b)`$ $`=`$ $`S_{\mathrm{cool}}(0)\left[1\left({\displaystyle \frac{b}{r_{\mathrm{cool}}}}\right)^2\right]^{0.5+3\beta _{\mathrm{cool}}},`$ $`S_{\mathrm{cool}}(0)`$ $`=`$ $`n_{0,\mathrm{cool}}^2r_{\mathrm{cool}}\mathrm{\Lambda }(T_{\mathrm{cool}})B(0.5,3\beta _{\mathrm{cool}}+1),`$ (8) In the equations above, $`T_{\mathrm{cool}}`$ and $`T_{\mathrm{amb}}`$ represent the two gas temperatures corresponding to the cooling region and to the ambient of the cluster, respectively. ## 3 COMPARISON WITH THE DATA I have applied this model to observations of clusters of galaxies that can map the emission in regions well beyond the cluster core to disentangle the effect of the two components. Moreover, I have considered clusters with evidence of a large cooling flow and the Coma cluster (ROR: rp800005n00, exposure time: 20.0 ksec, $`z=0.0232`$), a well-known example of a no-cooling flow cluster. In particular, I have analyzed, as described in Ettori and Fabian (1999), the ROSAT PSPC images of A1795 (ROR: rp800105n00, exposure time: 33.3 ksec, $`z=0.062`$) and A2199 (ROR: rp800644n00, exposure time: 33.9 ksec, $`z=0.030`$), that have a deposition rate larger than 100 $`M_{}`$ yr<sup>-1</sup> (Allen et al. 2000) and present in the literature detailed spectral analyses of the ASCA dataset to be used as reference. In Table 1, I quote the results obtained by fitting the azimuthally averaged profiles of the surface brightness between 0 and $`R_{\mathrm{out}}`$, where the brightness value is larger than 3 times the uncertainty in that radial bin. I perform the following fits (see Fig. 2): (i) single $`\beta `$model, (ii) the double $`\beta `$model presented here with 6 parameters $`[S_{\mathrm{cool}}(0),r_{\mathrm{cool}},\beta _{\mathrm{cool}},S_{\mathrm{amb}}(0),r_\mathrm{c},\beta _{\mathrm{amb}}]`$, (iii) the double $`\beta `$model with 5 parameters, i.e. fixing $`r_\mathrm{c}=r_{\mathrm{cool}}`$. The decrease of the $`\chi ^2`$ is significant where a cooling flow is present. Where this decrease is not meaningful (or not present), like in Coma cluster, the single $`\beta `$model still represents a good description of the data. Where a two-phase model provides a significantly better fit, I find that the F–test shows no statistical improvement with a 6 parameters fit (cf. Table 1). Therefore, I use in the following considerations the best fit results obtained with a 5 parameters fit, i.e. $`r_\mathrm{c}=r_{\mathrm{cool}}`$. This is not in contradiction with the present observational results. Allen (2000) quotes the cooling radii obtained from deprojection analysis of 30 cooling flow clusters images. Ettori & Fabian (1999) estimate the core radii for 23 of these clusters using a single $`\beta `$model over the radial range \[0.1, 1\] $`r_{500}`$, the radius at which the mean cluster density is 500 the background value. The distribution of the ratio, $`r_\mathrm{c}/r_{\mathrm{cool}}`$, has a median value of 1.33, an average of 1.61 and a standard deviation of 1.19, and can be considered consistent with $`1`$. For the clusters in exam here, I measure a $`r_\mathrm{c}/r_{\mathrm{cool}}`$ ratio of 1.32$`(\pm 0.24`$, 90 per cent confidence level) and 0.84 $`(\pm 0.15)`$ for A1795 and A2199, respectively. I remind, however, that I am using a different definition of $`r_{\mathrm{cool}}`$ than the one adopted in the standard spatial analysis: in the latter, $`r_{\mathrm{cool}}`$ is the radius where the cluster cooling time first exceeds the Hubble time, whereas in this work $`r_{\mathrm{cool}}`$ defines the boundary of the central cool phase of the gas. It is worth to note that another version of a ‘5 parameters’ fit, in which $`\beta _{\mathrm{cool}}`$ is fixed equal to $`\beta _{\mathrm{amb}}`$ and $`r_\mathrm{c},r_{\mathrm{cool}}`$ are left free to vary, provides a significantly worse $`\chi ^2`$ than the one obtained by using the ‘5 parameters’ fit adopted here. Finally, the $`\chi ^2`$ obtained with the models above does not vary significantly from the $`\chi ^2`$ measured after the fit with other models which are not strictly based upon a physical framework, like the sum of two standard $`\beta `$models with 6 free parameters (e.g. Reiprich & Böhringer 1999) or, fixed the slope to a common value, with 5 free parameters (Mohr et al. 1999). ## 4 DISCUSSION The use of this physically meaningful model allows us to directly handle each of the two gas distributions, one that describes the profile of the gas related to the cooling flow and the other that model the ambient gas. As shown above, there is no statistical justification for using the 6 parameter fit. Therefore, I consider hereafter the case $`r_\mathrm{c}=r_{\mathrm{cool}}`$. We use in the following discussion the definition of the central gas density for a cluster at redshift $`z`$, that is $$n_\mathrm{p}(0)=\left[\frac{4\pi (1+z)^4S(0)}{1.21r_\mathrm{c}\mathrm{\Lambda }(T)B}\right]^{0.5}$$ (9) where $`n_\mathrm{p}(0)`$ is in cm<sup>-3</sup>, $`r_\mathrm{c}`$ is the core or cooling radius in cm, $`S(0)`$ is the central surface brightness in cts s<sup>-1</sup> ster<sup>-1</sup>, $`B`$ is the proper beta function and the cooling function $`\mathrm{\Lambda }(T)`$ is in unit of cts s<sup>-1</sup> cm<sup>5</sup>. Using the condition that the two phases have to be in pressure equilibrium, I can now put constraints on their temperatures. To do this, I consider the mean properties of each phase to handle integrated values instead of differential ones, because of the simple assumption that each phase is represented by a single temperature that does not depend on the radius. Therefore, each phase density, $`\overline{n}(<r)`$, is the integral of the radial density, $`n(r)`$, over the volume occupied from that phase (between 0 and $`r_{\mathrm{cool}}`$ for the inner phase; between 0 and $`r=Xr_{\mathrm{cool}}`$ for the outer component) divided by the integrated volume. Then, $$T_{\mathrm{cool}}=T_{\mathrm{amb}}\left(\frac{n_{\mathrm{amb}}}{n_{\mathrm{cool}}}\right)=T_{\mathrm{amb}}g^{1/(2\alpha )}I^{2/(2\alpha )}=T_{\mathrm{amb}}f$$ (10) where I have made use of the relation in eq. 9 $`\mathrm{\Lambda }(T)n_0^2S(0)/B(a,b)`$, I have assumed $`\mathrm{\Lambda }(T)T^\alpha `$ ($`\alpha 0.5`$ for only bremsstrahlung emission observed by broad–band instruments), and I have defined $`g`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Lambda }(T_{\mathrm{amb}})n_{\mathrm{amb}}(0)^2}{\mathrm{\Lambda }(T_{\mathrm{cool}})n_{\mathrm{cool}}(0)^2}}={\displaystyle \frac{S_{\mathrm{amb}}(0)}{S_{\mathrm{cool}}(0)}}{\displaystyle \frac{B_{\mathrm{cool}}}{B_{\mathrm{amb}}}}`$ $`I`$ $`=`$ $`{\displaystyle \frac{\overline{n_{\mathrm{amb}}}(<r)}{n_{\mathrm{amb}}(0)}}{\displaystyle \frac{n_{\mathrm{cool}}(0)}{\overline{n_{\mathrm{cool}}}(<r_{\mathrm{cool}})}}`$ $`=`$ $`{\displaystyle \frac{1}{X^3}}{\displaystyle \frac{_0^X(1+x^2)^{1.5\beta _{\mathrm{amb}}}x^2𝑑x}{_0^1(1x^2)^{1.5\beta _{\mathrm{cool}}}x^2𝑑x}}`$ $`=`$ $`{\displaystyle \frac{I_X[1.5]/X^3}{0.5B(1.5,1.5\beta _{\mathrm{cool}}+1)}}`$ $`f`$ $`=`$ $`g^{1/(2\alpha )}I^{2/(2\alpha )},`$ (11) with $`_0^1(1x^2)^{1.5\beta _{\mathrm{cool}}}x^2𝑑x=B(1.5,1.5\beta _{\mathrm{cool}}+1)/2`$ and $`I_X[a]=_0^X(1+x^2)^{a\beta _{\mathrm{amb}}}x^2𝑑x`$. However, one generally measures a single emission-weighted temperature, $`\overline{T}`$. Given the considerations above, I can now disentangle the two components (if any) using eq. 10 in the following relation: $`\overline{T}`$ $`=`$ $`{\displaystyle \frac{Tϵ𝑑V}{ϵ𝑑V}}={\displaystyle \frac{_{i=1,2}T_iϵ_i𝑑V_i}{_{i=1,2}ϵ_i𝑑V_i}}`$ (12) $`=`$ $`T_{\mathrm{amb}}{\displaystyle \frac{1+fB(1.5,3\beta _{\mathrm{cool}}+1)(2gI_X[3])^1}{1+B(1.5,3\beta _{\mathrm{cool}}+1)(2gI_X[3])^1}}`$ $`=`$ $`T_{\mathrm{amb}}F,`$ where I still use the relation $`\mathrm{\Lambda }(T)n_0^2S(0)/B(a,b)`$, calculate $`_0^1(1x^2)^{3\beta _{\mathrm{cool}}}x^2𝑑x=B(1.5,3\beta _{\mathrm{cool}}+1)/2`$ and adopt the symbols $`f,g,I_X`$ defined in eq. 11. In the equations above, both the function $`f`$ and $`F`$ have to be smaller than 1 by definition. Their behaviour, however, depends strongly upon $`X`$, the radius in unit of $`r_{\mathrm{cool}}`$ up to where the outer phase extends and can be represented with a single temperature. Figure 3 shows how the function $`f`$ and $`F`$ depend upon $`X`$: $`f`$ diminishes significantly due to the presence of $`X^3`$ in $`I`$ (eq. 11), whereas $`F`$ converges quite rapidly (at $`X`$ 4), providing a robust estimate on the $`\overline{T}/T_{\mathrm{amb}}`$ ratio. Therefore, even if we are not able to constrain the ratio between the temperature of the two phases due to the uncertainty of the extension of the outer component, we can assess the ambient temperature, $`T_{\mathrm{amb}}`$, in a cooling flow cluster with an emission weighted temperature, $`\overline{T}`$, just using the azimuthally averaged surface brightness profile. I show in Table 2 the constraints on $`F=\overline{T}/T_{\mathrm{amb}}`$ obtained from the spatial fit of the surface brightness profiles of the clusters in exam and compare these values to the two-temperature spectral results in Allen et al. (2000), Markevitch et al. (1998, 1999) and White (2000). The agreement is remarkably good with the results of Markevitch and collaborators and White, which assume a two-phases gas for their spectral model in a way similar to the one I have adopted for the physical framework described above. (Note that Markevitch et al. measure the isothermal ambient temperature excising the cooling region, whereas White adds a cooling flow component in the spectral fit). On the other hand, the disagreement with the results of Allen and collaborators can be explained with the more complex model that they adopt, where an absorption intrinsic to the cluster is combined with the cooling flow only. Several aspects of the cluster physical characteristics are affected from the inclusion of a cooling flow in the modeling of the surface brightness with a $`\beta `$model. With respect to the single $`\beta `$model, one expects (i) excess in the gas density in the cooling region, (ii) change in the $`\beta _{\mathrm{amb}}`$ value, (iii) variations in the gas ambient temperature. I present in Table 3 some of the more interesting physical quantities that can be evaluated with the equations above and given an emission-weighted temperature, $`\overline{T}`$ (from Allen et al. 2000: $`k\overline{T}=5.40\pm 0.05`$ and $`4.16\pm 0.03`$ keV for A1795 and A2199, respectively). For example, if one identifies the inner component with the cooling flow, a proper description of its gas distribution is now available through eq. 7. The luminosity of the intracluster gas can then be estimated without the contribution of the emission from the cool phase: $$L_{\mathrm{amb}}=1.21n_\mathrm{p}^2(r)\mathrm{\Lambda }(T_{\mathrm{amb}})4\pi r^2𝑑r$$ (13) where the integral is computed upon the cluster volume and $`\mathrm{\Lambda }(T)`$ is here in erg s<sup>-1</sup> cm<sup>3</sup>. Using only the cluster surface brightness profile and a broad–band emission-weighted gas temperature, and applying eq. 12 and 13, I will investigate in a forthcoming paper the effects on the clusters luminosity–temperature relation of the presence of significant cooling flows (see, e.g., the results from spectral analyses in Allen & Fabian 1998, Markevitch 1998). Appreciable corrections can also affect $`M_{\mathrm{gas}}`$, $`M_{\mathrm{tot}}`$ and the terms of the so-called $`\beta `$problem (Mushotzky 1984, Edge & Stewart 1991) due to the variation of the $`\beta `$ value (for A1795 and A2199, $`\beta _{\mathrm{amb}}`$ increases by 20 and 10 per cent when compared to the $`1\beta `$ model fit results, respectively). In the two-phases model described here, it is simple to calculate the gas and the total mass: $`M_{\mathrm{gas}}`$ is the integral of $`\rho _{\mathrm{gas}}=\rho _{\mathrm{cool}}+\rho _{\mathrm{amb}}=2.21\mu m_\mathrm{p}(n_{\mathrm{p},\mathrm{cool}}+n_{\mathrm{p},\mathrm{amb}})`$ upon the cluster volume and, in particular, $$M_{\mathrm{gas}}(r>r_{\mathrm{cool}})=M_{\mathrm{gas},\mathrm{cool}}+_0^r\rho _{\mathrm{amb}}4\pi r^2𝑑r,$$ (14) where $`M_{\mathrm{gas},\mathrm{cool}}=2\pi \rho _{0,\mathrm{cool}}r_{\mathrm{cool}}^3B(1.5,1.5\beta _{\mathrm{cool}}+1)27(3)\times 10^{11}M_{}`$ for A1795 (A2199) for an assumed $`T_{\mathrm{cool}}=(f/F)\overline{T}=0.2\overline{T}`$; the total gravitating mass is given by the application of the hydrostatic equilibrium, $$M_{\mathrm{tot}}=\frac{r^2}{\mu m_\mathrm{p}G\rho _{\mathrm{gas}}}\frac{d(T_{\mathrm{cool}}\rho _{\mathrm{cool}}+T_{\mathrm{amb}}\rho _{\mathrm{amb}})}{dr},$$ (15) and is proportional to $`r_{\mathrm{cool}}\beta _{\mathrm{amb}}T_{\mathrm{amb}}x^3(1+x^2)^1`$ at $`r>r_{\mathrm{cool}}`$. ## 5 CONCLUSIONS I have presented a new analytic formula to model the total surface brightness profile of clusters of galaxies where a two-phases intracluster gas can be assumed. This scenario is consistent with the present results of spectral analyses of the central regions of clusters that harbour a cooling flow. The use of this formula allows to properly disentangle the contribution of the cooling flow to the cluster emissivity using only the spatial distribution of the X-ray photons. After removing the contamination from the cooling flow, I show how some relevant physical parameters are affected, like, for example, the ambient gas temperature (see Table 2). In a forthcoming paper, I will investigate the systematic changes in the temperature, luminosity and mass (cf. Table 3) of a sample of clusters of galaxies and how these variations affect the relations among these quantities. ## ACKNOWLEDGEMENTS I thank Anna, Sara and Carlo. I acknowledge the support of the Royal Society. Andy Fabian and David White are thanked for an useful reading of the manuscript, and the anonymous referee for comments which improved this work.
warning/0005/hep-th0005262.html	ar5iv	text	# 1 Introduction ## 1 Introduction The aim of this paper is to get a relation of the partition function of topological twisted $`N=4`$ gauge theory with Seiberg-Witten invariants in four manifolds. The partition function is given by Euler number of instanton moduli space in some conditions. We will show that the Euler number labeled by instanton number $`k`$ is expressed by Seiberg-Witten invariants when the sum of Euler number and signature of the base four manifolds vanishes. This result gives us the formulas to get the partition function of the twisted $`N=4`$ gauge theory by Seiberg-Witten invariants. The partition functions of the $`N=4`$ Yang-Mills theories on some four manifolds are calculated by Vafa-Witten with topological field theory . It is an $`SL(2,Z)`$ modular form. $`SL(2,Z)`$ transformation is understood as an extension of Montonen-Olive duality . So the duality relation is apparent in that partition function. This duality is deeply connected with the Hilbert scheme picture of instanton moduli space . But, in general, instanton moduli space has variety compactification and the sum of Euler number of any compactified moduli space is not necessarily a modular form. Actually, in our calculus, the partition function is not modular form with no contrivance. On the other side, $`N=4`$ gauge theory is given by the toroidal compactification of 10-dim $`N=1`$ gauge theory on a 4-dim manifold. (Note that ”compactification” is used two ways.) So the theory is interpreted as a low energy theory of the Heterotic or TypeI string theory. Recent developments of string theory show us many evidences of duality relation in field theory. In our case, Vafa shows us one method to link the compactified instanton moduli space with the Hilbert scheme . This fact implies that a choice of compactification is understood in string theory better than field theory. We discuss the problem of compactification and duality later. For our purpose we use a similar tool to . They used the non-abelian monopole theory and related the Donaldson invariants to Seiberg-Witten invariants without using duality . We also calculate the partition function in low energy limit of cohomological field theory and there is no request of S-duality. This is the most different point from Dijkgraaf-Park-Schroers . They have determined the partition function of $`N=4`$ supersymmetric Yang-Mills theory on a Kähler surface, using S-duality. Their result is given by Seiberg-Witten invariants, too. So it is interesting to compare our results with theirs. What we do first is to extend the instanton moduli space to non-abelian monopole moduli . In usual cohomological field theory, it was done in . Vafa-Witten theory is constructed as a balanced topological field theory (we denote it as BTFT in the following) . BTFT has no ghost number anomaly, and its partition function is a sum of Euler number of given zero-section space under vanishing theorem. In the 2nd section, we will construct the non-abelian monopole theory as BTFT and investigate some character of the theory. The vanishing theorem is an obstruction to construct the partition function as the sum of Euler number of the monopole moduli, and to get a relation with Vafa-Witten theory. We do not study this case closer in this paper. In the 3rd section, we get the formulas between the partition function of a twisted $`N=4`$ Yang-Mills theory and Seiberg-Witten invariants. To get them, we break the balance of topological charge. The tools in this paper were used in getting a relation of Donaldson invariants and Seiberg-Witten invariants . We use a model which has a gauge multiplet that is balanced and a hypermultiplet that is not balanced. We call the model unbalanced topological QCD. Vacuum expectation value of an observable is calculated and the relation between Euler number of instanton moduli space and Seiberg-Witten invariants is obtained if vanishing theorem is applicable and the sum of Euler number and signature of the four manifolds vanishes. The comparison with is also made in this section. At the last section, we discuss some remaining problems and the possibility of extension. ## 2 Balanced Topological QCD In this section, we construct a Balanced Topological QCD (BTQCD), which is a twisted $`N=4`$ Yang-Mills theory coupled with massive hypermultiplets in the fundamental representation . ### 2.1 Balanced Topological QCD Let $`X`$ be a compact Riemannian four manifold and $`E`$ be an $`SU(2)`$-bundle over $`X`$. The bundle $`E`$ is classified by the instanton number $$k=\frac{1}{8\pi ^2}_XTrFF,$$ (2.1) where $`Tr`$ is the trace in the fundamental representation of $`SU(2)`$ and $`F\mathrm{\Omega }_X^2(𝒢_E)`$ is the adjoint valued curvature 2-form on $`X`$. We denote the group of gauge transformations by $`𝒢`$, i.e. elements of $`𝒢`$ are sections of $`P`$, where $`P`$ is the associated principal $`SU(2)`$-bundle over $`X`$. We pick a $`spin^c`$ structure $`c`$ on $`X`$ and consider the associated $`spin^c`$ bundle $`W_c^\pm `$. Let $`𝒜`$ be the space of all connections on $`P`$ and $`\mathrm{\Gamma }(W_c^+E)(\mathrm{\Gamma }(W_c^{}E))`$ the space of the sections of the $`spin^c`$ bundle twisted by the vector bundle $`E`$. After twisting, the complex boson in the hypermultiplet becomes a section of $`\mathrm{\Gamma }(W_c^+E)(\mathrm{\Gamma }(W_c^{}E))`$; $$q\mathrm{\Gamma }(W_c^+E),q^{}\mathrm{\Gamma }(\overline{W}_c^+\stackrel{~}{E}),$$ $$B\mathrm{\Gamma }(W_c^{}E),B^{}\mathrm{\Gamma }(\overline{W}_c^{}\stackrel{~}{E}),$$ (2.2) where $`\stackrel{~}{E}`$ denotes the vector bundle conjugate to $`E`$. The $`spin^c`$ Dirac operator $$\sigma ^\mu D_\mu :\mathrm{\Gamma }(W_c^+E)\mathrm{\Gamma }(W_c^{}E),$$ (2.3) is the Dirac operator for the $`spin^c`$ bundle twisted by $`E`$. We will sometimes denote $`\sigma ^\mu D_\mu `$ by $`\overline{)}D`$ or $`\overline{)}D_c^E`$. Throughout this paper, we restrict our attention to the case that the gauge group is $`SU(2)`$ and the theory is coupled with hypermultiplets in the fundamental representation. #### algebra of BTQCD In this paragraph, the algebra of BTQCD is given. We introduce two global supercharges $`Q_\pm `$ carrying an additive quantum number (ghost number) $`U=\pm 1`$. When they act on fields in the adjoint representation, they satisfy the following commutation relations: $$Q_+^2=\delta _\theta ^g,\{Q_+,Q_{}\}=\delta _c^g,Q_{}^2=\delta _{\overline{\theta }}^g,$$ (2.4) where $`\delta _\theta ^g`$ denotes the gauge transformation generated by adjoint scalar field $`\theta \mathrm{\Omega }_X^0(𝒢_E)`$ and we adopt $`\delta _\theta ^gA_\mu =D_\mu \theta ,\delta _\theta ^gB_{+\mu \nu }=i[B_{+\mu \nu },\theta ],\delta _\theta ^gc=i[c,\theta ]`$. When they act on fields in the fundamental representation, they satisfy the following commutation relations: $$Q_+^2=\delta _\theta ^g,\{Q_+,Q_{}\}=\delta _c^g,Q_{}^2=\delta _{\overline{\theta }}^g,$$ (2.5) where we also introduce $`U(1)`$ global transformation generated by $`miR`$ and we adopt $`\delta _\theta ^gq=(i\theta +m)q,\delta _\theta ^gq^{}=q^{}(i\theta m),\delta _\theta ^gB=(i\theta +m)B,\delta _\theta ^gB^{}=B^{}(i\theta m)`$. The relative sign difference between (2.4) and (2.5) is simply the difference of representations. A simple explanation is the following. One can construct a field $`J^a`$ in the adjoint representation with a pair of fields $`q,q^{}`$ in the fundamental representation, $$J^aq^{}T^aq.$$ (2.6) Using above transformations, one can check (2.4) follows from (2.5), $`Q_+^2J^a`$ $`=`$ $`Q_+^2(q^{}T^aq)`$ (2.7) $`=`$ $`(\delta _\theta ^gq^{})T^aq+q^{}T^a(\delta _\theta ^gq)`$ $`=`$ $`i[q^{}Tq,\theta ]^a`$ $`=`$ $`\delta _\theta ^gJ^a.`$ Note that the relative sign difference between (2.4) and (2.5) is consistent with this derivation. The recipe for giving mass to fields in the fundamental representation by global symmetry is considered by Hyun-Park-Park(H-P-P). We define $`\delta _\pm `$ transformations $`\delta _\pm [Q_\pm ,\}`$. $`\delta _\pm `$ transformations are given in Appendix A. See also the references . #### action of BTQCD Using above fields and transformations, we define the action of BTQCD as $$h^2S=\sqrt{g},$$ (2.8) where $$=\delta _+\delta _{}.$$ (2.9) $``$ is described with fields in the previous paragraph and has ghost number $`0`$. The general recipe for constructing a balanced topological field theory is given by Moore et.al. $``$ is explicitly given by, $``$ $`=`$ $`(B_+^{\mu \nu a}s_{+\mu \nu }^a)(\chi _+^{I\mu \nu a}\psi _{B\mu \nu }^a)(\chi _{B\mu }^{IIa}\psi ^{\mu a})+(i{\displaystyle \frac{1}{3}}B_{+}^{\mu \nu }{}_{}{}^{a}[B_{+\mu \rho },B_{+\nu \sigma }]^ag^{\rho \sigma })`$ (2.10) $`+(B^\alpha s_\alpha )(\chi _q^{I\alpha }\psi _{B\alpha })(\chi _{B\dot{\alpha }}^{II}\psi _q^{\dot{\alpha }})`$ $`+(s^\alpha B_\alpha )+(\psi _B^\alpha \chi _{q\alpha }^I)+(\psi _{q\dot{\alpha }}^{}\chi _B^{II\dot{\alpha }})`$ $`+(\xi ^a\eta ^a),`$ where $$s_+^{\mu \nu }=F_+^{\mu \nu }+q^{}\overline{\sigma }^{\mu \nu }q$$ (2.11) $$s_\alpha =(\overline{)}Dq)_\alpha .$$ (2.12) Finally, full lagrangian is given by $`^{full}=\delta _+\delta _{}.`$ (2.13) Explicit expression of this lagrangian is given in appendix A. This lagrangian (LABEL:full) is different from in matter fields ($`q,B`$ etc.) and also different from H-P-P in dual fields ($`B_+^{\mu \nu },c,B`$ etc.). But due to its construction, it is balanced. ### 2.2 Fixed Point In this subsection, we study the nature of the action given in subsection 2.1. Here in particular we investigate the fixed points and vanishing theorem. #### Fixed Point To check the nature of lagrangian, we decompose the bosonic part of lagrangian (LABEL:full) $$_{boson}^{full}=_{boson}^{eq}+_{boson}^{pro},$$ (2.14) where $`_{boson}^{eq}`$ $`=`$ $`H_+^{I\mu \nu a}\{H_{+\mu \nu }^{Ia}(s_{+\mu \nu }^ai[B_{+\mu \rho },B_{+\nu \sigma }]^ag^{\rho \sigma }i[B_{+\mu \nu },c]^a)\}`$ $`H_B^{II\rho a}\{H_{B\rho }^{IIa}(2D^\mu B_{+\mu \rho }^a+iB^{}\sigma _\rho T^aqiq^{}\overline{\sigma }_\rho T^aBD_\rho c^a)\}`$ $`H_q^{I\alpha }\{H_{q\alpha }^I(s_\alpha +icB_\alpha +m_cB_\alpha )\}+(h.c.)`$ $`H_{B\dot{\alpha }}^{II}\{H_B^{II\dot{\alpha }}((\overline{)}DB)^{\dot{\alpha }}+(\overline{\sigma }^{\mu \nu }B_{+\mu \nu }q)^{\dot{\alpha }}+icq^{\dot{\alpha }}+m_cq^{\dot{\alpha }})\}+(h.c.)`$ and $`_{boson}^{pro}`$ $`=`$ $`\{[\theta ,\overline{\theta }]^a[\overline{\theta },\theta ]^a[c,\theta ]^a[c,\overline{\theta }]^a+[B_+^{\mu \nu },\overline{\theta }]^a[B_{+\mu \nu },\theta ]^a\}+D_\mu \overline{\theta }^aD^\mu \theta ^a`$ (2.16) $`+(iq^{}\overline{\theta }q^{}\overline{m})(i\theta q+mq)+(iq^{}\theta q^{}m)(i\overline{\theta }q+\overline{m}q)`$ $`+(iB^{}\overline{\theta }B^{}\overline{m})(i\theta B+mB)+(iB^{}\theta B^{}m)(i\overline{\theta }B+\overline{m}B).`$ $`_{boson}^{eq}`$ is defining the moduli space that we want to consider and $`_{boson}^{pro}`$ is induced for the projection to gauge normal direction. (LABEL:2.18) lagrangian is rewritten as $`_{boson}^{eq}`$ $`=`$ $`Hsquareterms`$ (2.17) $`+{\displaystyle \frac{1}{4}}(s_{+\mu \nu }^ai[B_{+\mu \rho },B_{+\nu \sigma }]^ag^{\rho \sigma })^2{\displaystyle \frac{1}{4}}([B_{+\mu \nu },c]^a)^2`$ $`+{\displaystyle \frac{1}{4}}(2D^\mu B_{+\mu \rho }^a+iB^{}\sigma _\rho T^aqiq^{}\overline{\sigma }_\rho T^aB)^2+{\displaystyle \frac{1}{4}}(D_\rho c^a)^2`$ $`+{\displaystyle \frac{1}{2}}\|s\|^2+{\displaystyle \frac{1}{2}}\|icB+m_cB\|^2`$ $`+{\displaystyle \frac{1}{2}}\|(\overline{)}D^{}B)^{\dot{\alpha }}+(\overline{\sigma }^{\mu \nu }B_{+\mu \nu }q)^{\dot{\alpha }}\|^2+{\displaystyle \frac{1}{2}}\|icq+m_cq\|^2.`$ Thus we have the following fixed point equations $$F_{+\mu \nu }+q^{}\overline{\sigma }_{\mu \nu }qi[B_{+\mu \rho },B_{+\nu \sigma }]g^{\rho \sigma }=0$$ $$2D_\mu B_+^{\mu \nu }+iB^{}\sigma ^\nu qiq^{}\overline{\sigma }^\nu B=0$$ $$s=\overline{)}Dq=0$$ $$\overline{)}D^{}B+\overline{\sigma }_{\mu \nu }B_+^{\mu \nu }q=0$$ $$D_\nu \theta =D_\nu c=D_\nu \overline{\theta }=0$$ $$[\theta ,\overline{\theta }]=[c,\theta ]=[c,\overline{\theta }]=[B_+^{\mu \nu },\theta ]=[B_+^{\mu \nu },\overline{\theta }]=[B_+^{\mu \nu },c]=0$$ $$(i\theta +m)q=(i\overline{\theta }+\overline{m})q=(ic+m_c)q=0$$ $$q^{}(i\theta m)=q^{}(i\overline{\theta }\overline{m})=q^{}(icm_c)=0$$ $$(i\theta +m)B=(i\overline{\theta }+\overline{m})B=(ic+m_c)B=0$$ $$B^{}(i\theta m)=B^{}(i\overline{\theta }\overline{m})=B^{}(icm_c)=0.$$ (2.18) If hypermultiplet fields are set to zero ($`q=q^{}=B=B^{}=0`$), then above equations are Vafa-Witten equations . Thus we call above equations Extended Vafa-Witten equations. #### problem In the previous paragraph, we have obtained fixed point equations of BTQCD. The equations for fermionic zero-modes are just the linearization of the fixed point equation and the condition that they are orthogonal to gauge orbits. Due to the balanced structure each fermionic zero-mode has a partner with the opposite $`U`$-number. Thus there is no ghost-number anomaly and the partition function is well defined, i.e.there is no need to insert observables. We want to compute the partition function of BTQCD. According to Vafa-Witten if an appropriate vanishing theorem holds, the partition function becomes the sum of Euler number of moduli space which we want to calculate. Roughly speaking, vanishing theorem is understood as the condition that dual fields ($`B_{+\mu \nu },c,B,B^{}`$ etc.) are to be zero and the dimensions of their moduli space become zero, when we choose an appropriate metric. But we could not verify that vanishing theorem holds in this model. To compare the result of this section to that of the next section, we give the only result to compute the partition function of BTQCD on the condition that vanishing theorem holds. ### 2.3 result In this subsection, we give the result of computing the path integral of BTQCD. We define partition function of BTQCD as $$Z=\frac{1}{Vol𝒢(2\pi )^\mathrm{\Omega }}𝒟W𝒟\psi _W𝒟Q^{}𝒟\psi _Q^{}𝒟Q𝒟\psi _Qe^S,$$ (2.19) where $$W=A_\mu ,B_+^{\mu \nu },H_B^\mu ,H_+^{\mu \nu },\theta ,c,\overline{\theta }$$ $$\psi _W=\psi _\mu ,\psi _B^{\mu \nu },\chi _B^{II\mu },\chi _+^{I\mu \nu },\xi ,\eta $$ $$Q=q,B,H_q^I,H_B^{II}$$ $$\psi _Q=\psi _q,\psi _B,\chi _q^I,\chi _B^{II}$$ $$\mathrm{\Omega }=dimofH^{}s.$$ (2.20) Here we denote auxiliary fields as $`H_B^\mu ,H_+^{\mu \nu },H_q^I,H_B^{II}`$, and we call auxiliary fields for $`Y`$ as $`H^{}sofY`$ in the following. $`dimofH^{}s`$ is a number of the auxiliary fields. After path-integrations of transverse part we get the partition function as the sum of two branches, according to the methods of the next section, $$Z=Z^{VW}+Z^{BU(1)SW}.$$ (2.21) $`Z^{VW}`$ is a contribution from branch 1) (gauge symmetry is unbroken), and corresponds to Vafa-Witten partition function. $`Z^{BU(1)SW}`$ is a contribution from branch 2) (gauge symmetry is broken to $`U(1)`$), and corresponds to balanced $`U(1)`$ monopole theory. The fixed point equations of the balanced $`U(1)`$ monopole theory are $$F_{+\mu \nu }^3+\frac{1}{2}q_{}^{}{}_{1}{}^{}\overline{\sigma }_{\mu \nu }q_1=0$$ $$2_\mu B_{+}^{\mu \nu }{}_{}{}^{3}+i\frac{1}{2}B_{}^{}{}_{1}{}^{}\sigma ^\nu q_1i\frac{1}{2}q_{}^{}{}_{1}{}^{}\overline{\sigma }^\nu B_1=0$$ $$\overline{)}D^3q_1=0$$ $$\overline{)}D^3B_1+\frac{1}{2}\overline{\sigma }_{\mu \nu }B_{+}^{\mu \nu }{}_{}{}^{3}q_1=0.$$ (2.22) Where $`F_{+\mu \nu }^3`$ is a curvature of $`U(1)`$ left symmetry after breaking $`SU(2)`$ and the labels of $`q_1`$ are $`B_1`$ are the ones of color. Since we do not know the vanishing theorem for dual fields ($`B_{+\mu \nu }^{}{}_{}{}^{3},B_1,B_{}^{}{}_{1}{}^{}`$) from (2.22), we stop to investigate this model further more in this paper. ## 3 Unbalanced Topological QCD In this section, we compute a correlation function of an appropriate BRS exact operator in Unbalanced Topological QCD. As a result, we can describe Euler number of instanton moduli space with Seiberg-Witten invariants. We have a similar but not the same expression to Dijkgraaf et.al, because we treat the different theory from theirs. We discuss this point at the end of this section. ### 3.1 Unbalanced Topological QCD Here we construct Unbalanced Topological QCD, which is a twisted $`N=4`$ Yang-Mills theory coupled with only one massive hypermultiplet in the fundamental representation (we denote it as UBTQCD in the following). Alternatively one get a UBTQCD, when one set one massive hypermultiplet $`(B,\psi _B,\chi _B^{II},H_B^{II})`$ of BTQCD in the previous section to zero (we call this process breaking balanced structure). #### algebra of UBTQCD The algebra of UBTQCD is given as a part of the BTQCD algebra. Contrary to the previous section, we only consider the global supercharge $`Q_+`$. When it acts on adjoint (fundamental) fields, it satisfies the following commutation relation: $$Q_+^2=\delta _\theta ^g(\delta _\theta ^g).$$ (3.1) We adopt the same $`\delta _+`$ transformations as the previous section and appendix A of (a.2.1)$``$(a.3.2). #### action of UBTQCD We define the action of UBTQCD as $$h^2S=d^4x\sqrt{g}$$ (3.2) where $$=\delta _+\mathrm{\Psi }.$$ (3.3) We explicitly give $`\mathrm{\Psi }`$ as; $`\mathrm{\Psi }`$ $`=`$ $`\chi _+^{I\mu \nu a}\{H_{+\mu \nu }^{Ia}(s_{+\mu \nu }^ai[B_{+\mu \rho },B_{+\nu \sigma }]^ag^{\rho \sigma }i[B_{+\mu \nu },c]^a)\}`$ (3.4) $`\chi ^{II\rho a}\{H_{B\rho }^{IIa}(2D^\mu B_{+\mu \rho }^aD_\rho c^a)\}`$ $`\chi _q^{I\alpha }\{H_{q\alpha }^Is_\alpha \}`$ $`\{H_q^{I\alpha }s^\alpha \}\chi _{q\alpha }^I`$ $`+\{i[\theta ,\overline{\theta }]^a\eta ^ai\xi ^a[c,\overline{\theta }]^a\}+i[B_+^{\mu \nu },\overline{\theta }]^a\psi _{B\mu \nu }^a+D_\mu \overline{\theta }^a\psi ^{\mu a}`$ $`(iq_{\dot{\alpha }}^{}\overline{\theta }\overline{m}q_{\dot{\alpha }}^{})\psi _q^{\dot{\alpha }}\psi _{q\dot{\alpha }}^{}(i\overline{\theta }q^{\dot{\alpha }}+\overline{m}q^{\dot{\alpha }}),`$ where $$s_+^{\mu \nu a}=F_{+}^{\mu \nu }{}_{}{}^{a}+q^{}\overline{\sigma }^{\mu \nu }T^aq$$ (3.5) $$s^\alpha =(\overline{)}Dq)^\alpha .$$ (3.6) Finally the full lagrangian is given by $`^{full}`$ $`=`$ $`\delta _+\mathrm{\Psi }`$ (3.8) $`=`$ $`H_+^{I\mu \nu a}\{H_{+\mu \nu }^{Ia}(s_{+\mu \nu }^ai[B_{+\mu \rho },B_{+\nu \sigma }]^ag^{\rho \sigma }i[B_{+\mu \nu },c]^a)\}`$ $`\chi _+^{I\mu \nu a}\{i[\chi _{+\mu \nu }^I,\theta ]^a+2D_\mu \psi _\nu ^a+\psi _q^{}\overline{\sigma }_{\mu \nu }T^aq+q^{}\overline{\sigma }_{\mu \nu }T^a\psi _q2i[B_{+\mu \rho },\psi _{B\nu \sigma }]^ag^{\rho \sigma }`$ $`i[\psi _{B\mu \nu },c]^ai[B_{+\mu \nu },\xi ]^a\}`$ $`H_B^{II\rho a}\{H_{B\rho }^{IIa}(2D^\mu B_{+\mu \rho }^aD_\rho c^a)\}`$ $`\chi _B^{II\rho a}\{i[\chi _{B\rho }^{II},\theta ]^a2D^\mu \psi _{B\mu \rho }^a2i[\psi ^\mu ,B_{+\mu \rho }]^aD_\rho \xi ^ai[\psi _\rho ,c]^a\}`$ $`H_q^{I\alpha }\{H_{q\alpha }^Is_\alpha \}`$ $`\chi _q^I\{\overline{)}D\psi _q+\sigma _\rho i\psi ^\rho q\}`$ $`+(h.c.abovetwolines)`$ $`\{[\theta ,\overline{\theta }]^a[\overline{\theta },\theta ]^a[c,\theta ]^a[c,\overline{\theta }]^a+[B_+^{\mu \nu },\overline{\theta }]^a[B_{+\mu \nu },\theta ]^a\}+D_\mu \overline{\theta }^aD^\mu \theta ^a`$ $`i[\theta ,\eta ]^a\eta ^a+i\xi ^a[\xi ,\overline{\theta }]^a+i\xi ^a[c,\eta ]^a+i[\psi _B^{\mu \nu },\overline{\theta }]^a\psi _{B\mu \nu }^a+i[B_+^{\mu \nu },\eta ]^a\psi _{B\mu \nu }^a+D_\mu \eta ^a\psi ^{\mu a}+i[\psi _\mu ,\overline{\theta }]^a\psi ^{\mu a}`$ $`+(iq^{}\overline{\theta }q^{}\overline{m})(i\theta q+mq)+(iq^{}\theta q^{}m)(i\overline{\theta }q+\overline{m}q)`$ $`+2\psi _q^{}(i\overline{\theta }+\overline{m})\psi _q2\chi _q^I(i\theta +m)\chi _q^I(iq^{}\eta q^{}\eta _m)\psi _q+\psi _q^{}(i\eta q+\eta _mq).`$ Notice that lagrangian (3.8) is given by lagrangian (LABEL:full) of previous section if $`(B,\psi _B,H_B^{II},\chi _B^{II})`$ is set to zero. ### 3.2 Fixed Point In this subsection, we study the nature of the action given in subsection 3.1. Here in particular we investigate the fixed points and some observable to insert. #### Fixed Point To check the nature of lagrangian, we decompose the bosonic part of lagrangian (3.8) $$_{boson}^{full}=_{boson}^{eq}+_{boson}^{pro},$$ (3.9) where $`_{boson}^{eq}`$ $`=`$ $`H_+^{I\mu \nu a}\{H_{+\mu \nu }^{Ia}(s_{+\mu \nu }^ai[B_{+\mu \rho },B_{+\nu \sigma }]^ag^{\rho \sigma }i[B_{+\mu \nu },c]^a)\}`$ (3.10) $`H_B^{II\rho a}\{H_{B\rho }^{IIa}(2D^\mu B_{+\mu \rho }^aD_\rho c^a)\}`$ $`H_q^{I\alpha }\{H_{q\alpha }^Is_\alpha \}+(h.c.)`$ and $`_{boson}^{pro}`$ $`=`$ $`\{[\theta ,\overline{\theta }]^a[\overline{\theta },\theta ]^a[c,\theta ]^a[c,\overline{\theta }]^a+[B_+^{\mu \nu },\overline{\theta }]^a[B_{+\mu \nu },\theta ]^a\}+D_\mu \overline{\theta }^aD^\mu \theta ^a`$ (3.11) $`+(iq^{}\overline{\theta }q^{}\overline{m})(i\theta q+mq)+(iq^{}\theta q^{}m)(i\overline{\theta }q+\overline{m}q).`$ $`_{boson}^{eq}`$ is defining the moduli space that we want to consider here and $`_{boson}^{pro}`$ is induced for the projection to gauge normal direction. (3.10) lagrangian is transformed into $`_{boson}^{eq}`$ $`=`$ $`\{H_{+\mu \nu }^{Ia}{\displaystyle \frac{1}{2}}(s_{+\mu \nu }^ai[B_{+\mu \rho },B_{+\nu \sigma }]^ag^{\rho \sigma }i[B_{+\mu \nu },c]^a)\}^2`$ (3.12) $`\{H_{B\rho }^{IIa}{\displaystyle \frac{1}{2}}(2D^\mu B_{+\mu \rho }^aD_\rho c^a)\}^2`$ $`2\|H_{q\alpha }^I{\displaystyle \frac{1}{2}}s_\alpha \|^2`$ $`+{\displaystyle \frac{1}{4}}(s_{+\mu \nu }^ai[B_{+\mu \rho },B_{+\nu \sigma }]^ag^{\rho \sigma }i[B_{+\mu \nu },c]^a)^2`$ $`+{\displaystyle \frac{1}{4}}(2D^\mu B_{+\mu \rho }^aD_\rho c^a)^2`$ $`+{\displaystyle \frac{1}{2}}\|s_\alpha \|^2.`$ Thus we have the following fixed point equations $$F_{+\mu \nu }+q^{}\overline{\sigma }_{\mu \nu }qi[B_{+\mu \rho },B_{+\nu \sigma }]g^{\rho \sigma }i[B_{+\mu \nu },c]=0$$ $$2D_\mu B_+^{\mu \nu }D^\nu c=0$$ $$s=\overline{)}Dq=0$$ $$D_\nu \theta =D_\nu \overline{\theta }=0$$ $$[\theta ,\overline{\theta }]=[c,\theta ]=[c,\overline{\theta }]=[B_+^{\mu \nu },\theta ]=[B_+^{\mu \nu },\overline{\theta }]=0$$ $$(i\theta +m)q=(i\overline{\theta }+\overline{m})q=0$$ $$q^{}(i\theta m)=q^{}(i\overline{\theta }\overline{m})=0.$$ (3.13) #### problem In the previous paragraph, we have obtained the fixed point equations of UBTQCD. In the same way as the previous section, the equations for fermionic zero-modes are just the linearization of the fixed point equations and the conditions that they are orthogonal to gauge orbits. Compared with the previous section, UBTQCD does not have balanced structure. In particular the hypermultiplet does not have balanced structure, while adjoint representation fields still have balanced structure. The partition function of unbalanced theory becomes zero due to its ghost number anomaly when the moduli space dimension of mater field is non-zero. Thus to get an well-defined path integral, we have to insert some observable. One can think an observable $$I=d^4x(q^{}(i\theta +m)q+\psi _q^{}\psi _q).$$ (3.14) Note that this observable itself is BRS exact, i.e. $$I=\delta _+\frac{1}{2}d^4x(\psi _q^{}q+q^{}\psi _q).$$ (3.15) Thus the expectation value of $`I`$ is zero according to Ward-Takahashi identity, and the expectation value of $`e^I`$ becomes zero when this theory has ghost number anomaly. However as we will see, we obtain non-trivial results. ### 3.3 branch In this subsection, we will show that the fixed point equations are decomposed to two branches. We take a similar treatment for (3.13) to . Equations $$D_\nu \theta =D_\nu \overline{\theta }=0,[\theta ,\overline{\theta }]=0$$ (3.16) imply that $`\theta ,\overline{\theta }`$ can be diagonalized in the fixed points. If connection $`A_\mu `$ are irreducible, $`\theta ,\overline{\theta }`$ should be zero (the gauge symmetry is unbroken ). If connection $`A_\mu `$ are reducible, $`\theta ,\overline{\theta }`$ can be non-zero (the gauge symmetry is broken down to $`U(1)`$). When these solutions are applied to $$(i\theta +m)q=(i\overline{\theta }+\overline{m})q=0$$ $$q^{}(i\theta m)=q^{}(i\overline{\theta }\overline{m})=0,$$ (3.17) we have two branches; branch 1) $`\theta =\overline{\theta }=0`$ and $`q=q^{}=0`$ or branch 2) $`\theta =\theta ^3T^30,\overline{\theta }=\overline{\theta }^3T^30`$ and $`q0,q^{}0`$. Note that in the branch 2) we choose unbroken $`U(1)`$ as $`T^3`$ direction without a loss of generality. #### branch 1) $`\theta =\overline{\theta }=0`$ and $`q=q^{}=0`$, i.e. the gauge symmetry is unbroken. Remaining fixed point equations are $$F_+^{\mu \nu }i[B_{+\mu \rho },B_{+\nu \sigma }]g^{\rho \sigma }=0,2D_\mu B_+^{\mu \nu }=0,D_\mu c=0$$ $$[B_+^{\mu \nu },c]=0.$$ (3.18) Here one may apply the same condition as Vafa-Witten to induce the vanishing theorem, and get the moduli space of $$F_+^{\mu \nu }=0.$$ (3.19) #### branch 2) $`\theta =\theta ^3T^30,\overline{\theta }=\overline{\theta }^3T^30`$ and $`q0,q^{}0`$, i.e.the gauge symmetry is broken to $`U(1)`$. Thus the bundle $`E`$ splits into line bundles, $`E=LL^1`$ with $`LL=k`$. Then equations (3.17) are $$(i\theta ^3T^3+m)q=\left(\begin{array}{cc}\frac{i}{2}\theta ^3+m& 0\\ 0& \frac{i}{2}\theta ^3+m\end{array}\right)\left(\begin{array}{c}q_1\\ q_2\end{array}\right)=0$$ $$(i\overline{\theta }^3T^3+\overline{m})q=\left(\begin{array}{cc}\frac{i}{2}\overline{\theta }^3+\overline{m}& 0\\ 0& \frac{i}{2}\overline{\theta }^3+\overline{m}\end{array}\right)\left(\begin{array}{c}q_1\\ q_2\end{array}\right)=0$$ $$q^{}(i\theta ^3T^3m)=\left(\begin{array}{cc}q_1^{}& q_2^{}\end{array}\right)\left(\begin{array}{cc}\frac{i}{2}\theta ^3m& 0\\ 0& \frac{i}{2}\theta ^3+m\end{array}\right)=0$$ $$q^{}(i\overline{\theta }^3T^3\overline{m})=\left(\begin{array}{cc}q_1^{}& q_2^{}\end{array}\right)\left(\begin{array}{cc}\frac{i}{2}\overline{\theta }^3\overline{m}& 0\\ 0& \frac{i}{2}\overline{\theta }^3+\overline{m}\end{array}\right)=0.$$ (3.20) Thus the only non-trivial solutions for $`q`$ are either $$q=\left(\begin{array}{c}q_1\\ 0\end{array}\right),q^{}=\left(\begin{array}{cc}q_1^{}& 0\end{array}\right)and\frac{i}{2}\theta ^3+m=\frac{i}{2}\overline{\theta }^3+\overline{m}=0$$ (3.21) or $$q=\left(\begin{array}{c}0\\ q_2\end{array}\right),q^{}=\left(\begin{array}{cc}0& q_2^{}\end{array}\right)and\frac{i}{2}\theta ^3+m=\frac{i}{2}\overline{\theta }^3+\overline{m}=0.$$ (3.22) Throughout this paper we pick the non-trivial solutions for $`q`$ as $`q_10and\theta ^3=2im`$. In this branch the equations $$[c,\theta ]=[c,\overline{\theta }]=[B_+^{\mu \nu },\theta ]=[B_+^{\mu \nu },\overline{\theta }]=0$$ (3.23) imply that non-zero solutions of $`B_+^{\mu \nu },c`$ have the same direction $`T^3`$ as $`\theta `$. Finally we get remaining equations $$F_{+\mu \nu }^3+\frac{1}{2}q_1^{}\overline{\sigma }_{\mu \nu }q_1=0$$ $$2^\mu B_{+\mu \nu }^3=^\mu c^3=0$$ $$\sigma ^\mu 𝒟_\mu q_1=0,$$ (3.24) where $`^\mu `$ is the covariant derivative in respect of Levi-Civita connection of background metric $`g^{\mu \nu }`$. Here we reinterpret $`U(1)U(1)`$ (gauge $`U(1)`$ and $`spin^cU(1)`$) as a new $`U(1)`$ ($`spin^c^{}U(1)`$), or alternately we redefine $`W_c^+\zeta =W_c^{}^+`$ as a different $`spin^c`$ structure $`c^{}=c+2\zeta `$, i.e., $`det(W_c^+\zeta )=L_c\zeta ^2`$. As a result, (3.24) can be interpreted as a perturbed Seiberg-witten monopole equation for the $`spin^c`$ structure $`c^{}`$ as well as H-P-P and $`B_+,c`$ equations for $`T^3`$ direction. ### 3.4 Gaussian integral In this subsection we compute the path integral of UBTQCD. According to Appendix, we could evaluate the exact path integral of this theory. In this subsection, we only denote the diagonal part of the big matrix (see Appendix) to read the right contribution easily. As we have already mentioned in subsection 3.2, we have to insert some observable of fundamental fields to get an well-defined path integral. Thus we define expectation value of $`e^I`$ as $$<e^I>_{m,c,k}=\frac{1}{Vol𝒢(2\pi )^\mathrm{\Omega }}𝒟W𝒟\psi _W𝒟Q^{}𝒟\psi _Q^{}𝒟Q𝒟\psi _Qe^{S+I},$$ (3.25) where $$W=A_\mu ,B_+^{\mu \nu },H_B^\mu ,H_+^{\mu \nu },\theta ,c,\overline{\theta }$$ $$\psi _W=\psi _\mu ,\psi _B^{\mu \nu },\chi _B^{II\mu },\chi _+^{I\mu \nu },\xi ,\eta $$ $$Q=q,H_q^I$$ $$\psi _Q=\psi _q,\chi _q^I$$ $$I=d^4x(q^{}(i\theta +m)q+\psi _q^{}\psi _q)$$ $$\mathrm{\Omega }=dimofH^{}s.$$ (3.26) In a general computation of the path integral of TFT, it is sufficient to keep only quadratic terms for the transverse degrees and compute the one-loop approximations which give a result exactly . Now let us see that in each branch, what is transverse degrees of freedom. Picking a Riemannian metric $`g`$ , we rescale $`gtg`$ and take $`t\mathrm{}`$ limit. In branch 1), the gauge symmetry is unbroken and the matter fields decouple as the transverse degrees of freedom. In branch 2), the gauge symmetry is broken down to $`U(1)`$ and the hypermultiplet reduces to one of its color. The suppressed color degrees of freedom for hypermultiplet and the components of the $`N=4`$ vector multiplet which do not belong to the Cartan subalgebra part become the transverse degrees of freedom. On the other hand, the path integrals for the non-transverse degrees should be computed exactly. These path integrals correspond to the path integral of Vafa-Witten theory in branch 1) and the path integral of $`U(1)`$ monopole theory and $`U(1)B_+,c`$ theory in branch 2). We will use the notation $`<O>_{m,c,k}`$ for the VEV evaluated in the massive UBTQCD for a given $`spin^c`$ and instanton number $`k`$. #### result of branch 1) In this branch, the degrees of freedom for the hypermultiplet become the transverse degrees of freedom. One can decompose the lagrangian (3.8) into two parts $$^{VW}(1)+^t(1),$$ (3.27) where the Vafa-Witten part $`^{VW}(1)`$ $`=`$ $`H_+^{I\mu \nu a}\{H_{+\mu \nu }^{Ia}(F_{+\mu \nu }^ai[B_{+\mu \rho },B_{+\nu \sigma }]^ag^{\rho \sigma }i[B_{+\mu \nu },c]^a)\}`$ $`\chi _+^{I\mu \nu a}\{i[\chi _{+\mu \nu }^I,\theta ]^a+2D_\mu \psi _\nu ^a2i[B_{+\mu \rho },\psi _{B\nu \sigma }]^ag^{\rho \sigma }i[\psi _{B\mu \nu },c]^ai[B_{+\mu \nu },\xi ]^a\}`$ $`H_B^{II\rho a}\{H_{B\rho }^{IIa}(2D^\mu B_{+\mu \rho }^aD_\rho c^a)\}`$ $`\chi _B^{II\rho a}\{i[\chi _{B\rho }^{II},\theta ]^a2D^\mu \psi _{B\mu \rho }^a2i[\psi ^\mu ,B_{+\mu \rho }]^aD_\rho \xi ^ai[\psi _\rho ,c]^a\}`$ $`\{[\theta ,\overline{\theta }]^a[\overline{\theta },\theta ]^a[c,\theta ]^a[c,\overline{\theta }]^a+[B_+^{\mu \nu },\overline{\theta }]^a[B_{+\mu \nu },\theta ]^a\}+D_\mu \overline{\theta }^aD^\mu \theta ^a`$ $`+i[\theta ,\eta ]^a\eta ^a+i\xi ^a[\xi ,\overline{\theta }]^a+i\xi ^a[c,\eta ]^a+i[\psi _B^{\mu \nu },\overline{\theta }]^a\psi _{B\mu \nu }^a`$ $`+i[B_+^{\mu \nu },\eta ]^a\psi _{B\mu \nu }^a+D_\mu \eta ^a\psi ^{\mu a}+i[\psi _\mu ,\overline{\theta }]^a\psi ^{\mu a}`$ and a quadratic lagrangian due to the transverse degrees $`^t(1)`$ $`=`$ $`H_q^{I\alpha }\{H_{q\alpha }^Is_\alpha \}`$ (3.29) $`\chi _q^I\overline{)}D\psi _q`$ $`+(h.c.abovetwolines)`$ $`2q^{}\overline{m}mq+2\psi _q^{}\overline{m}\psi _q2\chi _q^Im\chi _q^I`$ $`=`$ $`2\|H_q^I+\mathrm{}\|^22m\|\chi _q^I+\mathrm{}\|^2`$ $`{\displaystyle \frac{1}{2}}q^{}(\overline{)}D^{}\overline{)}D+4m\overline{m})q+{\displaystyle \frac{1}{2m}}\psi _q^{}(\overline{)}D^{}\overline{)}D+4m\overline{m})\psi _q`$ One can rewrite the path integral (3.25) in this branch as $`<e^I>_{m,c,k}(1)`$ $`=`$ $`\underset{Z_{m,c,k}^{VW}(1)}{\underset{}{{\displaystyle \frac{1}{Vol𝒢(2\pi )^\mathrm{\Omega }^{}}}{\displaystyle 𝒟W𝒟\psi _We^{S^{VW}(1)}}}}\underset{Z_{m,c,k}^t(1)}{\underset{}{{\displaystyle \frac{1}{(2\pi )^{\mathrm{\Omega }^{\prime \prime }}}}{\displaystyle 𝒟Q^{}𝒟\psi _Q^{}𝒟Q𝒟\psi _Qe^{S^t(1)+I(1)}}}},`$ where $$h^2S^{VW}(1)=d^4x\sqrt{g}^{VW}(1)$$ $$h^2S^t(1)=d^4x\sqrt{g}^t(1)$$ $$I(1)=d^4x\sqrt{g}(q^{}mq+\psi _q^{}\psi _q)$$ $$\mathrm{\Omega }^{}=dimofadjointH^{}s$$ $$\mathrm{\Omega }^{\prime \prime }=dimoffundamentalH^{}s.$$ (3.31) For the Vafa-Witten part $`Z_{m,c,k}^{VW}(1)`$, we completely follow Vafa-Witten. Thus we have $$Z_{m,c,k}^{VW}(1)\chi _k,$$ (3.32) where $`\chi _k`$ stands for the Euler number of instanton moduli space with instanton number $`k`$ and $``$ means equality under keeping the vanishing theorem as shown in Vafa-Witten. Note that the existence of the vanishing theorem in the previous section is unknown, but, in this case, we have some examples to which we apply the vanishing theorem . When the vanishing theorem is not applicable, we denote this part as $`Z_{m,c,k}^{VW}(1)`$ itself. We discuss the problem of compactification of moduli space later. For the transverse part $`Z_{m,c,k}^t(1)`$, we first perform $`H_q^I,\chi _q^I`$ integral and get $$\frac{1}{(2\pi )^{\mathrm{\Omega }^{\prime \prime }}}\frac{[det(2m)]_{(\chi _q^I,\chi _q^I)}}{[det(\frac{1}{\pi })]_{(H_q^I,H_q^I)}}=[det(\frac{m}{2\pi })]_\mathrm{\Gamma }^{}=(\frac{m}{2\pi })^{dim(\mathrm{\Gamma }_{\lambda >0}^{}Ker(\overline{)}D^{}))}.$$ (3.33) Second we perform $`q,\psi _q`$ integral for zero and non-zero mode respectively and get $$\frac{[det(\frac{\overline{)}D^2}{2m})]_{(\psi _q^{},\psi _q)_{non0}}}{[det(\frac{\overline{)}D^2}{4\pi })]_{(q^{},q)_{non0}}}\frac{[det(1)]_{(\psi _q^{},\psi _q)_0}}{[det(\frac{m}{2\pi })]_{(q^{},q)_0}}=(\frac{2\pi }{m})^{dim(\mathrm{\Gamma }_{\lambda >0}^+Ker(\overline{)}D)}.$$ (3.34) Note that this expression is not exact, but is sufficient to get the right contribution (see Appendix). Collecting (3.33) and (3.34), one can get $$Z_{m,c,k}^t(1)=(\frac{m}{2\pi })^{dim(\mathrm{\Gamma }_{\lambda >0}^{}Ker(\overline{)}D^{}))}(\frac{2\pi }{m})^{dim(\mathrm{\Gamma }_{\lambda >0}^+Ker(\overline{)}D)}=(\frac{2\pi }{m})^{index(\overline{)}D_c^E)}.$$ (3.35) Finally for $`<e^{\widehat{v}}>_{m,c,k}(1)`$, one can get $$<e^{\widehat{v}}>_{m,c,k}(1)=Z_{m,c,k}^{VW}(1)Z_{m,c,k}^t(1)=Z_{m,c,k}^{VW}(1)(\frac{2\pi }{m})^{index(\overline{)}D_c^E)}\chi _k(\frac{2\pi }{m})^{index(\overline{)}D_c^E)},$$ (3.36) where $``$ stands for results in the vanishing theorem case. #### result of branch 2) In this branch, the gauge symmetry is broken down to $`U(1)`$. The components of any field which do not belong to the Cartan subalgebra part become the transverse variables. That is the $`\pm `$ components of adjoint fields , i.e. $`T_\pm =T_1\pm iT_2`$. And the components of the hypermultiplet with the suppressed color index become the transverse variable. One can decompose the lagrangian (3.8) into two parts $$^{U(1)}(2)+^t(2),$$ (3.37) where $`^{U(1)}(2)`$ is the lagrangian of $`U(1)`$ UBTQCD, and $`^t(2)`$ is the quadratic lagrangian due to the transverse degrees. $`U(1)`$ part $`^{U(1)}(2)`$ can be further decomposed into two parts $$^{U(1)}(2)=_{mono}^{U(1)}(2)+_{B_+,c}^{U(1)}(2),$$ (3.38) $`_{mono}^{U(1)}(2)`$ $`=`$ $`H_+^{I3\mu \nu }\{H_{+\mu \nu }^{I3}(F_{+\mu \nu }^3+{\displaystyle \frac{1}{2}}q_1^{}\overline{\sigma }_{\mu \nu }q_1)\}`$ (3.39) $`\chi _+^{I3\mu \nu }\{2_\mu \psi _\nu ^3+{\displaystyle \frac{1}{2}}\psi _{q}^{}{}_{1}{}^{}\overline{\sigma }_{\mu \nu }q_1+{\displaystyle \frac{1}{2}}q_{}^{}{}_{1}{}^{}\overline{\sigma }_{\mu \nu }\psi _{q}^{}{}_{1}{}^{}\}`$ $`H_{q}^{I}{}_{1}{}^{}\{H_{q}^{I}{}_{1}{}^{}\overline{)}Dq_1\}`$ $`\chi _{q}^{I}{}_{1}{}^{}\overline{)}D\psi _{q}^{}{}_{1}{}^{}`$ $`+(h.c.abovetwolines)`$ $`+_\mu \overline{\theta }^3^\mu \theta ^3+_\mu \eta ^3\psi ^{3\mu }`$ $`+2(i{\displaystyle \frac{1}{2}}q_{}^{}{}_{1}{}^{}\overline{\theta }^3q_{}^{}{}_{1}{}^{}\overline{m})(i{\displaystyle \frac{1}{2}}\theta ^3q_1+mq_1)`$ $`+2\psi _{q}^{}{}_{1}{}^{}(i{\displaystyle \frac{1}{2}}\overline{\theta }^3+\overline{m})\psi _{q}^{}{}_{1}{}^{}2\chi _{q}^{I}{}_{1}{}^{}(i{\displaystyle \frac{1}{2}}\theta ^3+m)\chi _{q}^{I}{}_{1}{}^{}`$ $`(i{\displaystyle \frac{1}{2}}q_{}^{}{}_{1}{}^{}\eta ^3q_{}^{}{}_{1}{}^{}\eta _m)\psi _{q}^{}{}_{1}{}^{}+\psi _{q}^{}{}_{1}{}^{}(i{\displaystyle \frac{1}{2}}\eta ^3q_1+\eta _mq_1),`$ and $$_{B_+,c}^{U(1)}(2)=H_B^{II3\rho }\{H_{B\rho }^{II3}(2^\mu B_{+\mu \rho }^3_\rho c^3)\}\chi _B^{II3\rho }\{2^\mu \psi _{B\mu \rho }^3_\rho \xi ^3\},$$ (3.40) where the first part $`_{mono}^{U(1)}(2)`$ is $`U(1)`$ monopole theory, and the second part $`_{B_+,c}^{U(1)}(2)`$ is $`U(1)B_+,c`$ theory. The quadratic lagrangian due to the transverse degrees $`^t(2)`$ is $`^t(2)`$ $`=`$ $`4\|H_{+\mu \nu }^{I+}+\mathrm{}\|^28m\|\chi _{+\mu \nu }^{I+}+\mathrm{}\|^2+16m^2\|\overline{\theta }^++\mathrm{}\|^28m\|\eta ^++\mathrm{}\|^2`$ (3.41) $`A_\mu ^+\{(D^{3+}D^{3+})^{\mu \nu }+(D^3D^3)^{\mu \nu }\stackrel{~}{B}_+^{3\mu \rho }\stackrel{~}{B}_{+\rho }^{3\nu }{\displaystyle \frac{1}{2}}\stackrel{~}{q}_1^{}\overline{\sigma }^\mu \sigma ^\nu \stackrel{~}{q}_1+((\stackrel{~}{c}^3)^2+16m\overline{m})g^{\mu \nu }\}A_\nu ^{}`$ $`+{\displaystyle \frac{1}{2m}}\psi _\mu ^+\{(D^{3+}D^{3+})^{\mu \nu }+(D^3D^3)^{\mu \nu }\stackrel{~}{B}_+^{3\mu \rho }\stackrel{~}{B}_{+\rho }^{3\nu }{\displaystyle \frac{1}{2}}\stackrel{~}{q}_1^{}\overline{\sigma }^\mu \sigma ^\nu \stackrel{~}{q}_1+((\stackrel{~}{c}^3)^2+16m\overline{m})g^{\mu \nu }\}\psi _\nu ^{}`$ $`4\|H_B^{II+}+\mathrm{}\|^28m\|\chi _B^{II+}+\mathrm{}\|^2`$ $`B_{+\mu \nu }\{(D^{3+}D^{3+})^{\mu \rho }g^{\nu \sigma }4\stackrel{~}{B}_+^{3\mu \nu }\stackrel{~}{B}_+^{3\rho \sigma }+((\stackrel{~}{c}^3)^2+16m\overline{m})g^{\mu \rho }g^{\nu \sigma }\}B_{+\rho \sigma }`$ $`+{\displaystyle \frac{1}{2m}}\psi _{B\mu \nu }\{(D^{3+}D^{3+})^{\mu \rho }g^{\nu \sigma }4\stackrel{~}{B}_+^{3\mu \nu }\stackrel{~}{B}_+^{3\rho \sigma }+((\stackrel{~}{c}^3)^2+16m\overline{m})g^{\mu \rho }g^{\nu \sigma }\}\psi _{B\rho \sigma }`$ $`c^+\{D^3D^3\stackrel{~}{B}_+^{3\mu \nu }\stackrel{~}{B}_{+\mu \nu }^3(\stackrel{~}{c}^3)^216m\overline{m}\}c^{}`$ $`+{\displaystyle \frac{1}{2m}}\xi ^+\{D^3D^3\stackrel{~}{B}_+^{3\mu \nu }\stackrel{~}{B}_{+\mu \nu }^3(\stackrel{~}{c}^3)^216m\overline{m}\}\xi ^{}`$ $`2\|H_{q}^{I}{}_{2}{}^{}+\mathrm{}\|^24m\|\chi _{q}^{I}{}_{2}{}^{}+\mathrm{}\|^2`$ $`{\displaystyle \frac{1}{2}}q_2^{}\{\overline{)}D^3\overline{)}D^32\overline{\sigma }_{\mu \nu }\stackrel{~}{q}_1\stackrel{~}{q}_1^{}\overline{\sigma }^{\mu \nu }+16m\overline{m}\}q_2`$ $`+{\displaystyle \frac{1}{4m}}\psi _{q}^{}{}_{2}{}^{}\{\overline{)}D^3\overline{)}D^3+2\overline{\sigma }_{\mu \nu }\stackrel{~}{q}_1\stackrel{~}{q}_1^{}\overline{\sigma }^{\mu \nu }+16m\overline{m}\}\psi _{q}^{}{}_{2}{}^{}`$ $`+(crossterms).`$ One can rewrite the path integral (3.25) in this branch as $`<e^I>_{m,c,k}(2)`$ $`=`$ $`\underset{Z_{m,c,k}^{U(1)}(2)}{\underset{}{{\displaystyle \frac{1}{Vol𝒢^3(2\pi )^\mathrm{\Omega }^{}}}{\displaystyle 𝒟W^3𝒟\psi _W^3𝒟Q_1^{}𝒟\psi _{Q1}^{}𝒟Q_1𝒟\psi _{Q1}e^{S^{U(1)}(2)}}}}`$ $`\underset{Z_{m,c,k}^t(2)}{\underset{}{{\displaystyle \frac{1}{Vol𝒢^\pm (2\pi )^{\mathrm{\Omega }^{\prime \prime }}}}{\displaystyle 𝒟W^\pm 𝒟\psi _W^\pm 𝒟Q_2^{}𝒟\psi _{Q2}^{}𝒟Q_2𝒟\psi _{Q2}e^{S^t(2)+I(2)}}}}`$ where $$h^2S^{U(1)}(2)=d^4x\sqrt{g}^{U(1)}(2)$$ $$h^2S^t(2)=d^4x\sqrt{g}^t(2)$$ $$I(2)=d^4x\sqrt{g}(q_2^{}2mq_2+\psi _{q}^{}{}_{2}{}^{}\psi _{q}^{}{}_{2}{}^{})$$ $$\mathrm{\Omega }^{}=dimofH^{}sofnontransversedegrees$$ $$\mathrm{\Omega }^{\prime \prime }=dimofH^{}softransversedegrees.$$ (3.43) For $`U(1)`$ monopole part, we have $$Z_{mono}^{U(1)}=\frac{1}{Vol𝒢^3(2\pi )^{\mathrm{\Omega }^{\prime \prime \prime }}}𝒟W_A^3𝒟\psi _{W_A^3}𝒟Q_1^{}𝒟\psi _{Q1}^{}𝒟Q_1𝒟\psi _{Q1}e^{S_{mono}^{U(1)}(2)},$$ (3.44) where $$W_A^3=A_{}^{\mu }{}_{}{}^{3},H_{+\mu \nu }^3,\theta ^3,\overline{\theta }^3,q_1,H_{q}^{I}{}_{1}{}^{}$$ $$\psi _{W_A}^3=\psi _A^{\mu 3},\chi _{+\mu \nu }^3,\eta ^3,\psi _{q}^{}{}_{1}{}^{},\chi _{q}^{I}{}_{1}{}^{}$$ $$h^2S_{mono}^{U(1)}(2)=d^4x\sqrt{g}_{mono}^{U(1)}(2)$$ $$\mathrm{\Omega }^{\prime \prime \prime }=dimofH^{}sofU(1)SWpart.$$ (3.45) For this part we follow H-P-P . In a simple type manifold we only need to consider the zero-dimensional moduli space of the Seiberg-Witten monopoles (we call $`(x)`$). Here we denote $`spin^c`$ structure $`c^{}`$ that we have already mentioned in subsection 3.3 by $`2x`$ if $`c^{}`$ satisfies the condition of the zero-dimensional moduli space ($`dim(c^{})=\frac{c^{}c^{}}{4}\frac{2\chi +3\sigma }{4}=0`$), and we call this $`spin^c`$ structure $`x`$ Seiberg-Witten basic class. The moduli space $`(x)`$ consists of a finite set of points. First for the contributions of the zero-dimensional moduli space $`(x)`$, we have $$𝒩n_x,$$ (3.46) where $`𝒩`$ is the standard renormalization due to the local operators constructed from metric and depends only on $`\chi `$ and $`\sigma `$ . $`n_x`$ is the sum of the number of points counted with a sign and is called the Seiberg-Witten invariants. For the total contribution to $`U(1)`$ monopole part (3.44), we have to sum (3.46) with all basic classes $`x`$ and get $$Z_{mono}^{U(1)}=𝒩\underset{x}{}n_x,$$ (3.47) For $`U(1)B_+,c`$ part we have $$Z_{B_+}^{U(1)}=\frac{1}{(2\pi )^{\mathrm{\Omega }^{\prime \prime \prime \prime }}}𝒟W_{B_+,c}^3𝒟\psi _{W_{B_+,c}}^3e^{S_{B_+,c}^1(2)},$$ (3.48) where $$W_{B_+,c}^3=B_{+}^{\mu \nu }{}_{}{}^{3},H_B^{\mu 3},c^3$$ $$\psi _{W_{B_+,c}}^3=\psi _B^{\mu \nu 3},\chi _B^{II\mu 3},\xi ^3$$ $$h^2S_{B_+,c}^{U(1)}(2)=d^4x\sqrt{g}^{U(1)B_+,c}(2)$$ (3.49) $$\mathrm{\Omega }^{\prime \prime \prime \prime }=dimofH^{}sofU(1)B_+,cpart.$$ $`Z_{B_+}^{U(1)}`$ is the partition function of the cohomological field theory with the fixed point $$^\mu B_{+\mu \nu }^{}{}_{}{}^{3}=0,_\nu c^3=0.$$ (3.50) This partition function is sum of the $`\pm 1`$ when there are only isolated solutions as usual. The condition that the $`Z_{B_+}^{U(1)}`$ is non-zero is that the dimensions of the moduli space of the 0 section defined by (3.50) becomes zero. In fact the virtual dimension of this moduli space is calculated to be $$\mathrm{\Delta }=index(d^++d)=\frac{1}{2}(\chi +\sigma ),$$ (3.51) where $`\chi `$ and $`\sigma `$ are Euler number and signature of $`X`$ respectively. Thus $`\mathrm{\Delta }=0`$ is a condition that we get non-trivial results. We discuss this point later. Finally we get $$Z_{m,c,k}^{U(1)}(2)=𝒩Z_{B_+}^{U(1)}\underset{x}{}n_x.$$ (3.52) Now we evaluate the transverse integral $`Z_{m,c,k}^t(2)`$. Following H-P-P, we choose a unitary gauge in which $$\theta _\pm =0,$$ (3.53) where $$\theta =\theta ^3T^3+\theta ^+T^++\theta ^{}T^{}.$$ (3.54) In this gauge $`\theta `$ has values on the maximal torus (Cartan sub-algebra). By following the standard Faddev-Povov gauge fixing procedure, we introduce a new nilpotent BRST operator $`\delta `$ with the algebra $$\delta \theta _\pm =\pm iC_\pm \theta _3,\delta C_\pm =0,\delta \theta _3=0,\delta \overline{C}_\pm =b_\pm ,\delta b_\pm =0,$$ (3.55) where $`C_\pm `$ and $`\overline{C}_\pm `$ are anti-commuting ghosts and anti-ghosts, respectively, and $`b_\pm `$ are commuting auxiliary fields. The action for gauge fixing terms reads $`S_{m,gauge}(2)`$ $`=`$ $`\delta {\displaystyle \frac{1}{mh^2}}{\displaystyle d^4x\sqrt{g}(\theta _+\overline{C}_{}+\overline{C}_+\theta _{})}`$ (3.56) $`=`$ $`{\displaystyle \frac{1}{mh^2}}{\displaystyle d^4x\sqrt{g}\{\theta _+b_{}+b_+\theta _{}+iC_+\theta _3\overline{C}_{}+i\overline{C}_+\theta _3C_{}\}}`$ $`=`$ $`{\displaystyle \frac{1}{mh^2}}{\displaystyle d^4x\sqrt{g}\{\theta _+b_{}+b_+\theta _{}C_+2m\overline{C}_{}\overline{C}_+2mC_{}\}}.`$ From the second line to the third line, we take weak coupling limit and replace $`\theta ^3`$ with $`2im`$. Note that this action has ghost number $`0`$. Now consider the transverse part involving adjoint fields. We perform $`b_\pm ,C_\pm ,\overline{C}_\pm ,\overline{\theta }^\pm ,\eta ^\pm `$ integral and get $$[det(im)]_{\mathrm{\Omega }^0}^{\frac{1}{2}2}[det(2)]_{\mathrm{\Omega }^0}^{\frac{1}{2}2}[det(\frac{16m^2}{\pi })]_{\mathrm{\Omega }^0}^{\frac{1}{2}}[det(8m)]_{\mathrm{\Omega }^0}^{\frac{1}{2}}=[det(2\pi m)]_{\mathrm{\Omega }^0}^{\frac{1}{2}}=(2\pi m)^{\frac{1}{2}dim(\mathrm{\Omega }_{\lambda >0}^0Ker(D^3))}.$$ (3.57) * $`H_+^\pm ,\chi _+^\pm `$ integral $$[det(2\pi m)]_{\mathrm{\Omega }^{2+}}^{\frac{1}{2}}=(2\pi m)^{\frac{1}{2}dim(\mathrm{\Omega }_{\lambda >0}^{2+}Ker(D^{3+}))}.$$ (3.58) * $`H_B^\pm ,\chi _B^\pm `$ integral $$[det(2\pi m)]_{\mathrm{\Omega }^1}^{\frac{1}{2}}=(2\pi m)^{\frac{1}{2}dim(\mathrm{\Omega }_{\lambda >0}^1Ker(D^{3+}+D^3))}.$$ (3.59) * $`A^\pm ,\psi ^\pm `$ integral for non-zero mode $$\frac{[det(\frac{D^3D^3+D^{3+}D^{3+}}{m})]_{\psi _{non0}^\pm }^{\frac{1}{2}}}{[det(\frac{D^3D^3+D^{3+}D^{3+}}{2\pi })]_{A_{non0}^\pm }^{\frac{1}{2}}}=[det(\frac{2\pi }{m})]_{\mathrm{\Omega }_{non0}^1}^{\frac{1}{2}}=(\frac{2\pi }{m})^{\frac{1}{2}dim(\mathrm{\Omega }_{\lambda >0}^1)}.$$ (3.60) * $`B_+^\pm ,\psi _B^\pm `$ integral for non-zero mode $$\frac{[det(\frac{D^{3+}D^{3+}}{m})]_{\chi _{B}^{\pm }{}_{non0}{}^{}}^{\frac{1}{2}}}{[det(\frac{D^{3+}D^{3+}}{2\pi })]_{B_{+}^{\pm }{}_{non0}{}^{}}^{\frac{1}{2}}}=[det(\frac{2\pi }{m})]_{\mathrm{\Omega }_{non0}^{2+}}^{\frac{1}{2}}=(\frac{2\pi }{m})^{\frac{1}{2}dim(\mathrm{\Omega }_{\lambda >0}^{2+})}$$ (3.61) * $`c^\pm ,\xi ^\pm `$ integral for non-zero mode $$\frac{[det(\frac{D^3D^3}{m})]_{\xi _{non0}^\pm }^{\frac{1}{2}}}{[det(\frac{D^3D^3}{2\pi })]_{c_{non0}^\pm }^{\frac{1}{2}}}=[det(\frac{2\pi }{m})]_{\mathrm{\Omega }_{non0}^0}^{\frac{1}{2}}=(\frac{2\pi }{m})^{\frac{1}{2}dim(\mathrm{\Omega }_{\lambda >0}^0)}.$$ (3.62) Now we collect all the contributions of the adjoint transverse part and get $`{\displaystyle \frac{1}{(2\pi )^{\mathrm{\Omega }_{ado}^{\prime \prime }}}}(2\pi m)^{\frac{1}{2}dim(\mathrm{\Omega }_{\lambda >0}^0Ker(D^3))}(2\pi m)^{\frac{1}{2}dim(\mathrm{\Omega }_{\lambda >0}^{2+}Ker(D^{3+}))}(2\pi m)^{\frac{1}{2}dim(\mathrm{\Omega }_{\lambda >0}^1Ker(D^{3+}+D^3))}`$ $`({\displaystyle \frac{2\pi }{m}})^{\frac{1}{2}dim(\mathrm{\Omega }_{\lambda >0}^1)}({\displaystyle \frac{2\pi }{m}})^{\frac{1}{2}dim(\mathrm{\Omega }_{\lambda >0}^{2+})}({\displaystyle \frac{2\pi }{m}})^{\frac{1}{2}dim(\mathrm{\Omega }_{\lambda >0}^0)}`$ $`=`$ $`1`$ Remaining transverse integral is fundamental part. First we perform $`H_{q}^{I}{}_{2}{}^{},\chi _{q}^{I}{}_{2}{}^{}`$ integral and get $$\frac{1}{(2\pi )^{\mathrm{\Omega }_{fun}^{\prime \prime }}}\frac{[det(4m)]_{(\chi _{q2}^{},\chi _{q2})}}{[det(\frac{1}{\pi })]_{(H_{q2}^{},H_q^2)}}=[det(\frac{m}{\pi })]_\mathrm{\Gamma }^{}=(\frac{m}{\pi })^{dim(\mathrm{\Gamma }_{\lambda >0}^{}Ker((\overline{)}D^3)^{}))}.$$ (3.64) Next we perform $`q_2,\psi _{q2}`$ integral for non-zero and zero mode respectively and get $$\frac{[det(\frac{\overline{)}D^{}\overline{)}D}{4m})]_{(\psi _{q2}^{},\psi _{q}^{}{}_{2}{}^{})_{non0}}}{[det(\frac{\overline{)}D^{}\overline{)}D}{4\pi })]_{(q_2^{},q_2)_{non0}}}\frac{[det(1)]_{(\psi _{q2}^{},\psi _{q}^{}{}_{2}{}^{})_0}}{[det(\frac{m}{\pi })]_{(q_2^{},q_2)_0}}=[det(\frac{\pi }{m})]_{\mathrm{\Gamma }^+}=(\frac{\pi }{m})^{dim(\mathrm{\Gamma }_{\lambda >0}^+Ker((\overline{)}D^3)))}.$$ (3.65) Collecting (3.64) and (3.65), one can get $$(\frac{m}{\pi })^{dim(\mathrm{\Gamma }_{\lambda >0}^{}Ker((\overline{)}D^3)^{}))}(\frac{\pi }{m})^{dim(\mathrm{\Gamma }_{\lambda >0}^+Ker((\overline{)}D^3)))}=(\frac{\pi }{m})^{index(\overline{)}D^3)}$$ (3.66) From (LABEL:3.59) and (3.66) we get $$Z_{m,c,k}^t(2)=(\frac{\pi }{m})^{index(\overline{)}D^3)}.$$ (3.67) Finally for $`<e^I>_{m,c,k}(2)`$ we get $$<e^I>_{m,c,k}(2)=Z_{m,c,k}^{U(1)}(2)Z_{m,c,k}^2(2)=𝒩(\frac{\pi }{m})^{index(\overline{)}D^3)}Z_{B_+}^{U(1)}\underset{x}{}n_x.$$ (3.68) #### synthesis As we have already mentioned , $`<e^I>_{m,c,k}`$ itself is zero. However from above two paragraphs each branch has non-trivial contributions. Thus we have finally $`0`$ $`=`$ $`Z_{m,c,k}^{VW}\left({\displaystyle \frac{2\pi }{m}}\right)^{index\overline{)}D_c^E}+𝒩Z_{B_+}{\displaystyle \underset{x}{}}n_x\left({\displaystyle \frac{\pi }{m}}\right)^{index\overline{)}D^3}`$ (3.69) $``$ $`\chi _k\left({\displaystyle \frac{2\pi }{m}}\right)^{index\overline{)}D_c^E}+𝒩Z_{B_+}{\displaystyle \underset{x}{}}n_x\left({\displaystyle \frac{\pi }{m}}\right)^{index\overline{)}D^3},`$ where the last expression is valid in the vanishing theorem case. In general $`index\overline{)}D_c^E`$ is calculated to be $$index\overline{)}D_c^E=k+\frac{rank(E)}{8}(cc\sigma ).$$ (3.70) In this case, $$c\zeta =index\overline{)}D_c^E+2\mathrm{\Delta }.$$ (3.71) The Dirac operator $`\overline{)}D^3`$ which operate on $`q_2,\psi _{q_2}`$, and so on is necessary to be understood as the Dirac operator with the connection given by $`c2\zeta `$. Then $`index\overline{)}D_c^3`$ $`=`$ $`0+{\displaystyle \frac{1}{8}}((c2\zeta )(c2\zeta )\sigma )`$ (3.72) $`=`$ $`{\displaystyle \frac{1}{8}}(cc4k+4index\overline{)}D_c^E8\mathrm{\Delta }\sigma ).`$ (3.73) Thus we get a relation $$index\overline{)}D^3=index\overline{)}D_c^E\mathrm{\Delta }$$ (3.74) Inserting (3.74) into (3.69), since $`m`$ is a free parameter, we get non-trivial result only in the case $`\mathrm{\Delta }=0`$. Remember that $`\mathrm{\Delta }`$ is also the dimension of the moduli space of $`U(1)B_+,c`$ theory. Thus the condition $`\mathrm{\Delta }=0`$ is consistent with defining $`Z_{B_+}^{U(1)}`$ (3.48). $`\mathrm{\Delta }=0`$ is also consistent with geographic condition, for example simple type condition($`b_2^+>1`$), Furuta theory($`b_2\frac{5}{4}\|\sigma \|+2`$) and $`\frac{11}{8}`$ conjecture($`b_2\frac{11}{8}\|\sigma \|`$) . Finally under the condition $`\mathrm{\Delta }=0`$, from (3.69) we have $$\chi _kZ_{m,c,k}^{VW}=𝒩Z_{B_+}\underset{x}{}n_x\left(\frac{1}{2}\right)^{index\overline{)}D_c^E}.$$ (3.75) Note that above $`x`$ satisfies that $`xx=\frac{2\chi +3\sigma }{4}`$ and $`x=\frac{c+2\zeta }{2}`$. We think the Vafa-Witten partition as the sum of (3.75) with weight $`e^{\tau k}`$, where $`\tau `$ is a parameter. But the sum of this partition function don not clarify modular invariance since $`\mathrm{\Delta }=0`$ is special case which do not depend on the coupling $`\tau `$ in topological twisted model . Additionally we do not assume duality, then there is no guarantee that our partition function has modular invariance and is same as Vafa-Witten’s. We suppose that the difference come from compactification of the moduli space. We do not use the duality relation and our model is not asymptotic-free theory. So, there is possibility that compactification in our theory is not the same as the one in Hilbert scheme. Thus we can describe the twisted $`N=4`$ Yang-Mills partition function that may not be the same as Vafa-Witten’s partition function with Seiberg-Witten invariants. Our expression is similar to Dijkgraaf. The most significant difference is $`\tau `$ dependence. Dijkgraaf’s is $`\tau `$ dependent, while ours is $`\tau `$ independent. The reason why Dijkgraaf’s partition function depends on $`\tau `$ is that they treat the physical $`N=4`$ Yang-Mills theory itself. According to Labastida, the $`N=4`$ Yang-Mills theory depends on $`\tau `$. On the other hand we treat UBTQCD, which is the twisted $`N=4`$ Yang-Mills theory coupled with a fundamental hypermultiplet. As we mention above, this difference may cause breaking the modular invariance. In other words, our theory is not conformal invariant, and $`\tau `$ is not possible to be a good parameter. But our computation is done without assumption like duality relation. If there is difference we have to interpret the origin of the difference occurred from compactification . ## 4 Conclusion We have studied the balanced topological QCD and its broken balance theory and got relations of the partition function of twisted $`N=4`$ SU(2) Yang-Mills theory with the partition function of twisted abelian QCD. This relation is understood in several ways. For example, the sum of Euler number of instanton moduli space , which is invariant under $`SL(2,Z)`$ transformation, is described by Seiberg-Witten invariants when $`\mathrm{\Delta }=0`$ and the vanishing theorem is valid. In other cases there is no vanishing theorem like §5.4 in , we got a similar but not the same formulas under the condition of $`\mathrm{\Delta }=0`$. There is no other reasons to understand the difference from the result of Vafa-Witten and Dijkgraaf et al. than the difference of compactification. Some problems are left for our future work. When $`\mathrm{\Delta }0`$, can we obtain any similar non-trivial results without assumption of duality relation? We may obtain them by simple reformation. But it is difficult to expect that the partition function have the nature of modular invariance in naive reformation. We are interested in a connection with the duality and a compactification. How can we obtain the modular invariant partition function with no assumption of duality? We have some hints of this question but no answer. As we saw in section 2, vanishing theorem of BTQCD is not studied in this paper. If the theorem exists, we get the sum of Euler number of non-abelian monopole moduli space as the partition function of the BTQCD. It is an interesting work to investigate the nature of the partition function because the theory has the branches that contain both Vafa-Witten theory and Seiberg-Witten theory. Acknowledgment We are grateful to H.Kanno for helpful suggestions and observations and a critical reading of the manuscript. We also would like to thank M.Furuta and K.Ono for valuable discussion. A.S. is supported by JSPS Research Fellowships for Young Scientists. ## Appendix A the BRS algebra and the BTQCD action We give the BRS algebra and the lagrangian of BTQCD explicitly in this appendix. ### A.1 Algebra $`\delta _\pm `$ transformations are given as follows. $$\{\begin{array}{c}\delta _{}A_\mu =\chi _\mu ^{II}\\ \delta _{}\chi _\mu ^{II}=\delta _g^{\overline{\theta }}A_\mu =D_\mu \overline{\theta }\\ \delta _{}\psi _\mu =\delta _g^cA_\mu H_{B\mu }^{II}=D_\mu cH_{B\mu }^{II}\\ \delta _{}H_{B\mu }^{II}=\delta _g^c\chi _{B\mu }^{II}+\delta _+\delta _g^{\overline{\theta }}A_\mu =i[\chi _{B\mu }^{II},c]+\delta _+\delta _g^{\overline{\theta }}A_\mu \end{array}$$ (a.1.1) $$\{\begin{array}{c}\delta _{}B_+^{\mu \nu }=\chi _+^{I\mu \nu }\\ \delta _{}\chi _+^{I\mu \nu }=\delta _g^{\overline{\theta }}B_+^{\mu \nu }=i[B_+^{\mu \nu },\overline{\theta }]\\ \delta _{}\psi _B^{\mu \nu }=\delta _g^cB_+^{\mu \nu }H_+^{I\mu \nu }=i[B_+^{\mu \nu },c]H_+^{I\mu \nu }\\ \delta _{}H^{I\mu \nu }=\delta _g^c\chi _+^{I\mu \nu }+\delta _+\delta _g^{\overline{\theta }}B_+^{\mu \nu }=i[\chi _+^{I\mu \nu },c]+\delta _+\delta _g^{\overline{\theta }}B_+^{\mu \nu }\end{array}$$ (a.1.2) $$\{\begin{array}{c}\delta _{}q^{\dot{\alpha }}=\chi _B^{II\dot{\alpha }}\\ \delta _{}\chi _B^{II\dot{\alpha }}=\delta _g^{\overline{\theta }}q^{\dot{\alpha }}=i\overline{\theta }q^{\dot{\alpha }}+\overline{m}q^{\dot{\alpha }}\\ \delta _{}\psi _q^{\dot{\alpha }}=\delta _g^cq^{\dot{\alpha }}H_B^{II\dot{\alpha }}=icq^{\dot{\alpha }}+m_cq^{\dot{\alpha }}H_B^{II\dot{\alpha }}\\ \delta _{}H_B^{II\dot{\alpha }}=\delta _g^c\chi _B^{II\dot{\alpha }}\delta _+\delta _g^{\overline{\theta }}q^{\dot{\alpha }}=ic\chi _B^{II\dot{\alpha }}+m_c\delta _+\delta _g^{\overline{\theta }}q^{\dot{\alpha }}\end{array}$$ (a.1.3) $$\{\begin{array}{c}\delta _{}B_\alpha =\chi _{q\alpha }^I\\ \delta _{}\chi _{q\alpha }^I=\delta _g^{\overline{\theta }}B_\alpha =i\overline{\theta }B_\alpha +\overline{m}B_\alpha \\ \delta _{}\psi _{B\alpha }=\delta _g^cB_\alpha H_{q\alpha }^I=icB_\alpha +m_cB_\alpha H_{q\alpha }^I\\ \delta _{}H_{q\alpha }^I=\delta _g^c\chi _{q\alpha }^I\delta _+\delta _g^{\overline{\theta }}B_\alpha =ic\chi _{q\alpha }^I+m_c\chi _{q\alpha }^I\delta _+\delta _g^{\overline{\theta }}B_\alpha \end{array}$$ (a.1.4) $$\{\begin{array}{c}\delta _{}q_{\dot{\alpha }}^{}=\chi _{B\dot{\alpha }}^{II}\\ \delta _{}\chi _{B\dot{\alpha }}^{II}=\delta _g^{\overline{\theta }}q_{\dot{\alpha }}^{}=iq_{\dot{\alpha }}^{}\overline{\theta }\overline{m}q_{\dot{\alpha }}^{}\\ \delta _{}\psi _{q\dot{\alpha }}^{}=\delta _g^cq_{\dot{\alpha }}^{}H_{B\dot{\alpha }}^{II}=iq_{\dot{\alpha }}^{}cm_cq_{\dot{\alpha }}^{}H_{B\dot{\alpha }}^{II}\\ \delta _{}H_{B\dot{\alpha }}^{II}=\delta _g^c\chi _{B\dot{\alpha }}^{II}\delta _+\delta _g^{\overline{\theta }}q_{\dot{\alpha }}^{}=i\chi _{B\dot{\alpha }}^Icm_c\chi _{B\dot{\alpha }}^I\delta _+\delta _g^{\overline{\theta }}q_{\dot{\alpha }}^{}\end{array}$$ (a.1.5) $$\{\begin{array}{c}\delta _{}B^\alpha =\chi _q^{I\alpha }\\ \delta _{}\chi _q^{I\alpha }=\delta _g^{\overline{\theta }}B^\alpha =iB^\alpha \overline{\theta }\overline{m}B^\alpha \\ \delta _{}\psi _B^\alpha =\delta _g^cB^\alpha H_q^{I\alpha }=iB^\alpha cm_cB^\alpha H_q^{I\alpha }\\ \delta _{}H_q^{I\alpha }=\delta _g^c\chi _q^{I\alpha }\delta _+\delta _g^{\overline{\theta }}B^\alpha =i\chi _q^{I\alpha }cm_c\chi _q^{I\alpha }\delta _+\delta _g^{\overline{\theta }}B^\alpha .\end{array}$$ (a.1.6) $`\delta _+`$ transformations are given by $$\{\begin{array}{c}\delta _+A_\mu =\psi _\mu \\ \delta _+\psi _\mu =\delta _g^\theta A_\mu =D_\mu \theta \\ \delta _+\chi _{B\mu }^{II}=H_{B\mu }^{II}\\ \delta _+H_{B\mu }^{II}=\delta _g^\theta \chi _{B\mu }^{II}=i[\chi _{B\mu }^{II},\theta ]\end{array}$$ (a.2.1) $$\{\begin{array}{c}\delta _+B_+^{\mu \nu }=\psi _B^{\mu \nu }\\ \delta _+\psi _B^{\mu \nu }=\delta _g^\theta B_+^{\mu \nu }=i[B_+^{\mu \nu },\theta ]\\ \delta _+\chi _+^{I\mu \nu }=H_+^{I\mu \nu }\\ \delta _+H_+^{I\mu \nu }=\delta _g^\theta \chi _+^{I\mu \nu }=i[\chi _+^{I\mu \nu },\theta ]\end{array}$$ (a.2.2) $$\{\begin{array}{c}\delta _+q^{\dot{\alpha }}=\psi _q^{\dot{\alpha }}\\ \delta _+\psi _q^{\dot{\alpha }}=\delta _g^\theta q_q^{\dot{\alpha }}=(i\theta q_q^{\dot{\alpha }}+mq_q^{\dot{\alpha }})\\ \delta _+\chi _B^{II\dot{\alpha }}=H_B^{II\dot{\alpha }}\\ \delta _+H_B^{II\dot{\alpha }}=\delta _g^\theta \chi _B^{II\dot{\alpha }}=(i\theta \chi _B^{II\dot{\alpha }}+m\chi _B^{II\dot{\alpha }})\end{array}$$ (a.2.3) $$\{\begin{array}{c}\delta _+B_\alpha =\psi _{B\alpha }\\ \delta _+\psi _{B\alpha }=\delta _g^\theta B_\alpha =(i\theta B_\alpha +mB_\alpha )\\ \delta _+\chi _{q\alpha }^I=H_{q\alpha }^I\\ \delta _+H_{q\alpha }^I=\delta _g^\theta \chi _{q\alpha }^I=(i\theta \chi _{q\alpha }^I+m\chi _{q\alpha }^I)\end{array}$$ (a.2.4) $$\{\begin{array}{c}\delta _+q_{\dot{\alpha }}^{}=\psi _{q\dot{\alpha }}^{}\\ \delta _+\psi _{q\dot{\alpha }}^{}=\delta _g^\theta q_{\dot{\alpha }}^{}=(iq_{\dot{\alpha }}^{}\theta mq_{\dot{\alpha }}^{})\\ \delta _+\chi _{B\dot{\alpha }}^{II}=H_{B\dot{\alpha }}^{II}\\ \delta _+H_{B\dot{\alpha }}^{II}=\delta _g^\theta \chi _{B\dot{\alpha }}^{II}=(i\chi _{B\dot{\alpha }}^{II}\theta m\chi _{B\dot{\alpha }}^{II})\end{array}$$ (a.2.5) $$\{\begin{array}{c}\delta _+B^\alpha =\psi _B^\alpha \\ \delta _+\psi _B^\alpha =\delta _g^\theta B^\alpha =(iB^\alpha \theta mB^\alpha )\\ \delta _+\chi _q^{I\alpha }=H_q^{I\alpha }\\ \delta _+H_q^{I\alpha }=\delta _g^\theta \chi _q^{I\alpha }=(i\chi _q^{I\alpha }\theta m\chi _q^{I\alpha }).\end{array}$$ (a.2.6) Transformations for $`c,\theta ,\overline{\theta },m,m_c,\overline{m}`$ are given by $$\{\begin{array}{cc}\delta _+\theta =0& \\ \delta _{}\theta =\xi & ,\delta _+\xi =\delta _g^\theta c=i[c,\theta ]\\ \delta _+c=\xi & ,\delta _{}\xi =\delta _g^{\overline{\theta }}\theta =i[\theta ,\overline{\theta }]\\ \delta _{}c=\eta & ,\delta _+\eta =\delta _g^\theta \overline{\theta }=i[\overline{\theta },\theta ]\\ \delta _+\overline{\theta }=\eta & ,\delta _{}\eta =\delta _g^{\overline{\theta }}c=i[c,\overline{\theta }]\\ \delta _{}\overline{\theta }=0& \end{array}$$ (a.3.1) $$\{\begin{array}{cc}\delta _+m=0& \\ \delta _{}m=\xi _m& ,\delta _+\xi _m=0\\ \delta _+m_c=\xi _m& ,\delta _{}\xi _m=0\\ \delta _{}m_c=\eta _m& ,\delta _+\eta _m=0\\ \delta _+\overline{m}=\eta _m& ,\delta _{}\eta _m=0\\ \delta _{}\overline{m}=0.& \end{array}$$ (a.3.2) ### A.2 action of BTQCD We write down the lagrangian of BTQCD explicitly in this paragraph. $`\delta _{}`$ is given as $`\delta _{}`$ $`=`$ $`\delta _{}(B_+^{\mu \nu a}s_{+\mu \nu }^a)\delta _{}(\chi _+^{I\mu \nu a}\psi _{B\mu \nu }^a)\delta _{}(\chi _{B\mu }^{IIa}\psi ^{\mu a})+\delta _{}(i{\displaystyle \frac{1}{3}}B_{+}^{\mu \nu }{}_{}{}^{a}[B_{+\mu \rho },B_{+\nu \sigma }]^ag^{\rho \sigma })`$ (a.5) $`+\delta _{}(B^\alpha s_\alpha )\delta _{}(\chi _q^{I\alpha }\psi _{B\alpha })\delta _{}(\chi _{B\dot{\alpha }}^{II}\psi _q^{\dot{\alpha }})`$ $`+\delta _{}(s^\alpha B_\alpha )+\delta _{}(\psi _B^\alpha \chi _{q\alpha }^I)+\delta _{}(\psi _{q\dot{\alpha }}^{}\chi _B^{II\dot{\alpha }})`$ $`+\delta _{}(\xi ^a\eta ^a)`$ $`=`$ $`\chi _+^{I\mu \nu a}\{H_{+\mu \nu }^{Ia}(s_{+\mu \nu }^ai[B_{+\mu \rho },B_{+\nu \sigma }]^ag^{\rho \sigma }i[B_{+\mu \nu },c]^a)\}`$ $`\chi ^{II\rho a}\{H_{B\rho }^{IIa}(2D^\mu B_{+\mu \rho }^a+iB^{}\sigma _\rho T^aqiq^{}\overline{\sigma }_\rho T^aBD_\rho c^a)\}`$ $`\chi _q^{I\alpha }\{H_{q\alpha }^I(s_\alpha +icB_\alpha +m_cB_\alpha )\}`$ $`\chi _{B\dot{\alpha }}^{II}\{H_B^{II\dot{\alpha }}((\overline{)}D^{}B)^{\dot{\alpha }}+(\overline{\sigma }^{\mu \nu }B_{+\mu \nu }q)^{\dot{\alpha }}+icq^{\dot{\alpha }}+m_cq^{\dot{\alpha }})\}`$ $`\{H_q^{I\alpha }(s^\alpha iB^\alpha cm_cB^\alpha )\}\chi _{q\alpha }^I`$ $`\{H_{B\dot{\alpha }}^{II}((\overline{)}D^{}B)_{\dot{\alpha }}^{}+(q^{}B_{+\mu \nu }\overline{\sigma }^{\mu \nu })_{\dot{\alpha }}iq_{\dot{\alpha }}^{}cm_cq_{\dot{\alpha }}^{})\}\chi _B^{II\dot{\alpha }}`$ $`+\{i[\theta ,\overline{\theta }]^a\eta ^ai\xi ^a[c,\overline{\theta }]^a\}+i[B_+^{\mu \nu },\overline{\theta }]^a\psi _{B\mu \nu }^a+D_\mu \overline{\theta }^a\psi ^{\mu a}`$ $`(iB^\alpha \overline{\theta }\overline{m}B^\alpha )\psi _{B\alpha }(iq_{\dot{\alpha }}^{}\overline{\theta }\overline{m}q_{\dot{\alpha }}^{})\psi _q^{\dot{\alpha }}`$ $`\psi _B^\alpha (i\overline{\theta }B_\alpha +\overline{m}B_\alpha )\psi _{q\dot{\alpha }}^{}(i\overline{\theta }q^{\dot{\alpha }}+\overline{m}q^{\dot{\alpha }}),`$ The full lagrangian is given as $`^{full}`$ $`=`$ $`\delta _+\delta _{}`$ $`=`$ $`H_+^{I\mu \nu a}\{H_{+\mu \nu }^{Ia}(s_{+\mu \nu }^ai[B_{+\mu \rho },B_{+\nu \sigma }]^ag^{\rho \sigma }i[B_{+\mu \nu },c]^a)\}`$ $`\chi _+^{I\mu \nu a}\{i[\chi _{+\mu \nu }^I,\theta ]^a+2D_\mu \psi _\nu ^a+\psi _q^{}\overline{\sigma }_{\mu \nu }T^aq+q^{}\overline{\sigma }_{\mu \nu }T^a\psi _q2i[B_{+\mu \rho },\psi _{B\nu \sigma }]^ag^{\rho \sigma }`$ $`i[\psi _{B\mu \nu },c]^ai[B_{+\mu \nu },\xi ]^a\}`$ $`H_B^{II\rho a}\{H_{B\rho }^{IIa}(2D^\mu B_{+\mu \rho }^a+iB^{}\sigma _\rho T^aqiq^{}\overline{\sigma }_\rho T^aBD_\rho c^a)\}`$ $`\chi _B^{II\rho a}\{i[\chi _{B\rho }^{II},\theta ]^a2D^\mu \psi _{B\mu \rho }^a2i[\psi ^\mu ,B_{+\mu \rho }]^a`$ $`+i\psi _B^{}\sigma _\rho T^aq+iB^{}\sigma _\rho T^a\psi _qi\psi _q^{}\overline{\sigma }_\rho T^aBiq^{}\overline{\sigma }_\rho T^a\psi _BD_\rho \xi ^ai[\psi _\rho ,c]^a\}`$ $`H_q^{I\alpha }\{H_{q\alpha }^I(s_\alpha +icB_\alpha +m_cB_\alpha )\}`$ $`\chi _q^I\{\overline{)}D\psi _q+\sigma _\rho i\psi ^\rho q+i\xi B+ic\psi _B+\xi _mB+m_c\psi _B\}`$ $`+(h.c.abovetwolines)`$ $`H_{B\dot{\alpha }}^{II}\{H_B^{II\dot{\alpha }}((\overline{)}D^{}B)^{\dot{\alpha }}+(\overline{\sigma }^{\mu \nu }B_{+\mu \nu }q)^{\dot{\alpha }}+icq^{\dot{\alpha }}+m_cq^{\dot{\alpha }})\}`$ $`\chi _B^{II}\{\overline{)}D^{}\psi _B\overline{\sigma }_\rho i\psi ^\rho B+(\overline{\sigma }^{\mu \nu }\psi _{B\mu \nu }q)+(\overline{\sigma }^{\mu \nu }B_{+\mu \nu }\psi _q)+i\xi q+ic\psi _q+\xi _mq+m_c\psi _q\}`$ $`+(h.c.abovetwolines)`$ $`\{[\theta ,\overline{\theta }]^a[\overline{\theta },\theta ]^a[c,\theta ]^a[c,\overline{\theta }]^a+[B_+^{\mu \nu },\overline{\theta }]^a[B_{+\mu \nu },\theta ]^a\}+D_\mu \overline{\theta }^aD^\mu \theta ^a`$ $`+i[\theta ,\eta ]^a\eta ^a+i\xi ^a[\xi ,\overline{\theta }]^a+i\xi ^a[c,\eta ]^a+i[\psi _B^{\mu \nu },\overline{\theta }]^a\psi _{B\mu \nu }^a`$ $`+i[B_+^{\mu \nu },\eta ]^a\psi _{B\mu \nu }^a+D_\mu \eta ^a\psi ^{\mu a}+i[\psi _\mu ,\overline{\theta }]^a\psi ^{\mu a}`$ $`+(iq^{}\overline{\theta }q^{}\overline{m})(i\theta q+mq)+(iq^{}\theta q^{}m)(i\overline{\theta }q+\overline{m}q)`$ $`+2\psi _q^{}(i\overline{\theta }+\overline{m})\psi _q2\chi _q^I(i\theta +m)\chi _q^I(iq^{}\eta q^{}\eta _m)\psi _q+\psi _q^{}(i\eta q+\eta _mq)`$ $`+(iB^{}\overline{\theta }B^{}\overline{m})(i\theta B+mB)+(iB^{}\theta B^{}m)(i\overline{\theta }B+\overline{m}B)`$ $`+2\psi _B^{}(i\overline{\theta }+\overline{m})\psi _B2\chi _B^{II}(i\theta +m)\chi _B^{II}(iB^{}\eta B^{}\eta _m)\psi _B+\psi _B^{}(i\eta B+\eta _mB).`$ ## Appendix B the path integral of the transverse part As we have mentioned in the first part of section3.4, the path integral in 3.4 is not exact, but it amounts to the right result that we will derive in this section. In computation, we take the weak coupling limit. When we replace the non-transverse fields with the fixed point values, we denote $`Y_{nontrans}`$ by $`\stackrel{~}{Y}_{nontrans}`$. In particular the fixed points of $`\theta ,\overline{\theta }`$ are given as $`\theta =\overline{\theta }=0`$ in branch 1) and $`\theta ^3=2im,\overline{\theta }=2i\overline{m}`$ in branch 2). We also discuss the different treatment from the path integral in 3.4 at the end of this section. See , too. ### B.1 branch 1) and its big matrix In branch 1), the path integral of the transverse part is $`Z_{m,c,k}^t(1)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^{\mathrm{\Omega }^{\prime \prime }}}}{\displaystyle 𝒟Q^{}𝒟\psi _Q^{}𝒟Q𝒟\psi _Qe^{S^t(1)+I(1)}}`$ (b.1) where $$h^2S^t(1)=d^4x\sqrt{g}^t(1)$$ $$I(1)=d^4x\sqrt{g}(q^{}mq+\psi _q^{}\psi _q)$$ $$\mathrm{\Omega }^{\prime \prime }=dimoffundamentalH^{}s.$$ (b.2) For $`^t(1)`$(3.29), we denote $`^t(1)`$ $`=`$ $`2\|H_q^I+\mathrm{}\|^22m\|\chi _q^I+\mathrm{}\|^2`$ (b.3) $`{\displaystyle \frac{1}{2}}q^{}M^b(1)q+{\displaystyle \frac{1}{2m}}\psi _q^{}M^f(1)\psi _q,`$ where $$M^b(1)=M^f(1)=\overline{)}D^{}\overline{)}D+4m\overline{m}.$$ (b.4) In general $`M^b(M^f)`$ is matrix and is not necessarily diagonalized. $`M^b`$ and $`M^f`$ may not be the same as we will see soon. We call $`M^b(1)`$ big matrix of branch 1). Before computing (b.1), we briefly review the notion of $`index\overline{)}D`$. One can decompose $`q\mathrm{\Gamma }^+,H_q^I\mathrm{\Gamma }^{}`$ into $$\overline{)}D^{}\overline{)}Dq^\lambda =\lambda q^\lambda $$ $$\overline{)}D\overline{)}D^{}H_{q}^{I}{}_{}{}^{\lambda }=\lambda H_{q}^{I}{}_{}{}^{\lambda }.$$ (b.5) These decomposition is called spectra decomposition. Note that if $`\lambda >0`$ then $`q^\lambda `$ and $`H_{q}^{I}{}_{}{}^{\lambda }`$ are isomorphic. However if $`\lambda =0`$ then $`q^\lambda `$ and $`H_{q}^{I}{}_{}{}^{\lambda }`$ are not isomorphic. $`index\overline{)}D`$ measures the difference between $`\mathrm{\Gamma }_{\lambda =0}^+`$ and $`\mathrm{\Gamma }_{\lambda =0}^{}`$, and is defined as $`index\overline{)}D`$ $`=`$ $`dim\mathrm{\Gamma }_{\lambda =0}^+dim\mathrm{\Gamma }_{\lambda =0}^{}`$ (b.6) $`=`$ $`dimKer\overline{)}DdimKer\overline{)}D^{},`$ where we denote $`\mathrm{\Gamma }_{\lambda =0}^+=Ker\overline{)}D,\mathrm{\Gamma }_{\lambda =0}^+=Ker\overline{)}D^{}`$. In computing (b.1), (b.6) emerges when non-kinetic part and off-diagonal part of $`M`$ are able to be ignored (in this branch simply $`m\overline{m}`$ terms in (b.4)). This process is achieved by diagonalization and field redefinition. Then we get the expression (b.1) as $`index\overline{)}D`$. Conversely it is enough to get this expression that we consider only kinetic diagonal part of $`M`$ in the path integral. Now we perform the path integral of the transverse part of branch 1) explicitly. First for $`H_q^I,\chi _q^I`$ integral, $$\frac{1}{(2\pi )^{\mathrm{\Omega }^{\prime \prime }}}\frac{[det(2m)]_{(\chi _q^I,\chi _q^I)}}{[det(\frac{1}{\pi })]_{(H_q^I,H_q^I)}}=[det(\frac{m}{2\pi })]_\mathrm{\Gamma }^{}=(\frac{m}{2\pi })^{dim(\mathrm{\Gamma }_{\lambda >0}^{}Ker(\overline{)}D^{}))}.$$ (b.7) Note that the transformation at the second equality is necessary to derive $`index\overline{)}D`$. For $`q,\psi _q`$ integral for non-zero mode, $$\frac{[det(\frac{\overline{)}D^2}{2m})]_{(\psi _q^{},\psi _q)_{non0}}}{[det(\frac{\overline{)}D^2}{4\pi })]_{(q^{},q)_{non0}}}=(\frac{2\pi }{m})^{dim(\mathrm{\Gamma }_{\lambda >0}^+)}.$$ (b.8) For $`q,\psi _q`$ integral for zero mode, we consider that the integrant of this path integral comes only from observable $`I(1)`$ and we get $$\frac{[det(1)]_{(\psi _q^{},\psi _q)_0}}{[det(\frac{m}{2\pi })]_{(q^{},q)_0}}=(\frac{2\pi }{m})^{dimKer(\overline{)}D)}.$$ (b.9) From (b.8) and (b.9) $$(\frac{2\pi }{m})^{dim(\mathrm{\Gamma }_{\lambda >0}^+Ker(\overline{)}D))}.$$ (b.10) Collecting (b.7) and (b.10), $$Z_{m,c,k}^t=(\frac{m}{2\pi })^{dim(\mathrm{\Gamma }_{\lambda >0}^{}Ker(\overline{)}D^{}))}(\frac{2\pi }{m})^{dim(\mathrm{\Gamma }_{\lambda >0}^+Ker(\overline{)}D))}=(\frac{2\pi }{m})^{index\overline{)}D}.$$ (b.11) Note that $`dim(\mathrm{\Gamma }_{\lambda >0}^+)`$ and $`dim(\mathrm{\Gamma }_{\lambda >0}^{})`$ cancel each other. ### B.2 branch 2) and its big matrix In branch 2), the path integral of the transverse part is $`Z_{m,c,k}^t(2)`$ $`=`$ $`{\displaystyle \frac{1}{Vol𝒢^\pm (2\pi )^{\mathrm{\Omega }^{\prime \prime }}}}{\displaystyle 𝒟W^\pm 𝒟\psi _W^\pm 𝒟Q_2^{}𝒟\psi _{Q2}^{}𝒟Q_2𝒟\psi _{Q2}e^{S^t(2)+I(2)}},`$ where $$h^2S^t(2)=d^4x\sqrt{g}^t(2)$$ $$I(2)=d^4x\sqrt{g}(q_2^{}2mq_2+\psi _{q}^{}{}_{2}{}^{}\psi _{q}^{}{}_{2}{}^{})$$ $$\mathrm{\Omega }^{\prime \prime }=dimofH^{}softransversedegrees.$$ (b.13) For $`^t(2)`$ (3.41), we denote $`^t(2)`$ $`=`$ $`4\|H_{+\mu \nu }^{I+}+\mathrm{}\|^28m\|\chi _{+\mu \nu }^{I+}+\mathrm{}\|^2+16m^2\|\overline{\theta }^++\mathrm{}\|^28m\|\eta ^++\mathrm{}\|^2`$ (b.14) $`4\|H_B^{II+}+\mathrm{}\|^28m\|\chi _B^{II+}+\mathrm{}\|^2`$ $`2\|H_{q}^{I}{}_{2}{}^{}+\mathrm{}\|^24m\|\chi _{q}^{I}{}_{2}{}^{}+\mathrm{}\|^2`$ $`Y^TM^b(2)Y+{\displaystyle \frac{1}{2m}}\mathrm{\Psi }_Y^TM^f(2)\mathrm{\Psi }_Y,`$ where $`Y^T,\mathrm{\Psi }_Y^T`$ are raw vectors, $$Y^T=(A_\mu ^+,B_{+\nu \rho }^{}{}_{}{}^{+},c^+,q_{}^{}{}_{2}{}^{})$$ (b.15) $$\mathrm{\Psi }_Y^T=(\psi _\mu ^+,\psi _{B\nu \rho }^{}{}_{}{}^{+},\xi ^+,\psi _{q}^{}{}_{2}{}^{}),$$ (b.16) and $`Y,\mathrm{\Psi }_Y`$ are column vectors, $$Y=(A_\sigma ^{},B_{+\gamma \delta }^{}{}_{}{}^{},c^{},q_2)$$ (b.17) $$\mathrm{\Psi }_Y=(\psi _\sigma ^{},\psi _{B\gamma \delta }^{}{}_{}{}^{},\xi ^{},\psi _{q}^{}{}_{2}{}^{}).$$ (b.18) To derive the result (3.67) from (LABEL:a.7) (b.14), we can neglect the non-kinetic terms and off-diagonal part of $`M(2)`$ (we will give explicitly later). There is the contribution from Faddeev-Popov determinant of $`\theta ^\pm =0`$ gauge and it is possible to discard the path integral of $`Y^\pm `$ for zero mode according to the balanced structure of adjoint fields. In this remaining subsection, we concentrate on giving $`M(2)`$ explicitly. $`M^b(2)(M^f(2))`$ can be decomposed into $$M^b(2)=\left(\begin{array}{cc}M_{AA}^b& M_{Aq}^b\\ M_{q^{}A}^b& M_{q^{}q}^b\end{array}\right).$$ (b.19) We denote matrix element of $`M^b(2)`$ (or $`M_{AA}^b,M_{Aq}^b,M_{q^{}A}^b,M_{q^{}q}^b`$) by $`\{M^b(2)\}^{A^{+\mu }B_+^{\gamma \delta }},\{M_{Aq}^b\}^{A^{+\mu }q}`$ etc. #### diagonal part of $`M^b(2)(M^f(2))`$ $`\{M_{AA}^b\}^{A^{+\mu }A^\sigma }`$ $`=`$ $`\{M_{AA}^f\}^{A^{+\mu }A^\sigma }`$ $`=`$ $`(D^{3+}D^{3+})^{\mu \sigma }+(D^3D^3)^{\mu \sigma }\stackrel{~}{B}_+^{3\mu \rho }\stackrel{~}{B}_{+\rho }^{3\sigma }{\displaystyle \frac{1}{2}}\stackrel{~}{q}_1^{}\overline{\sigma }^\mu \sigma ^\sigma \stackrel{~}{q}_1+((\stackrel{~}{c}^3)^2+16m\overline{m})g^{\mu \sigma }`$ $`\{M_{AA}^b\}^{B_+^{+\nu \rho }B_+^{\gamma \delta }}`$ $`=`$ $`\{M_{AA}^f\}^{B_+^{+\nu \rho }B_+^{\gamma \delta }}`$ $`=`$ $`(D^{3+}D^{3+})^{\nu \rho }g^{\gamma \delta }4\stackrel{~}{B}_+^{3\nu \rho }\stackrel{~}{B}_+^{3\gamma \delta }+((\stackrel{~}{c}^3)^2+16m\overline{m})g^{\nu \rho }g^{\gamma \delta }`$ $`\{M_{AA}^b\}^{c^+c^{}}`$ $`=`$ $`\{M_{AA}^f\}^{c^+c^{}}`$ (b.22) $`=`$ $`D^3D^3\stackrel{~}{B}_+^{3\mu \nu }\stackrel{~}{B}_{+\mu \nu }^3(\stackrel{~}{c}^3)^216m\overline{m}`$ $`\{M_{q^{}q}^b\}^{q^{}q}`$ $`=`$ $`\overline{)}D^3\overline{)}D^32\overline{\sigma }_{\mu \nu }\stackrel{~}{q}_1\stackrel{~}{q}_1^{}\overline{\sigma }^{\mu \nu }+16m\overline{m}`$ (b.23) $`\{M_{q^{}q}^f\}^{q^{}q}`$ $`=`$ $`\overline{)}D^3\overline{)}D^3+2\overline{\sigma }_{\mu \nu }\stackrel{~}{q}_1\stackrel{~}{q}_1^{}\overline{\sigma }^{\mu \nu }+16m\overline{m}`$ (b.24) Note that $`\{M_{q^{}q}^b\}^{q^{}q}`$ and $`\{M_{q^{}q}^f\}^{q^{}q}`$ are different. #### off-diagonal part of $`M_{AA}^b(M_{AA}^f)`$ $`\{M_{AA}^b\}^{A^{+\mu }B_+^{\rho \sigma }}`$ $`=`$ $`\{M_{AA}^b\}^{A^{+\mu }B_+^{\rho \sigma }}`$ (b.25) $`=`$ $`i(\stackrel{}{D^{3+}})_\nu \stackrel{~}{c}^3g^{\mu \rho }g^{\nu \sigma }2i\stackrel{~}{B}_+^{3\mu \rho }(D^{3+})^\rho i(\stackrel{}{D^3})^\mu \stackrel{~}{B}_+^{3\rho \sigma }`$ $`2i(\stackrel{}{D^{3+}})_\nu \stackrel{~}{B}_+^{3\mu \rho }g^{\nu \sigma }`$ $`\{M_{AA}^b\}^{A^{+\mu }c^{}}`$ $`=`$ $`\{M_{AA}^b\}^{A^{+\mu }c^{}}`$ (b.26) $`=`$ $`i(\stackrel{}{D^{3+}})_\nu \stackrel{~}{B}_+^{3\mu \nu }2i\stackrel{~}{B}_+^{3\mu \rho }(D^3)_\rho +i(\stackrel{}{D^3})^\mu \stackrel{~}{c}^3`$ $`\{M_{AA}^b\}^{B_+^{+\nu \rho }c^{}}`$ $`=`$ $`\{M_{AA}^b\}^{B_+^{+\nu \rho }c^{}}`$ (b.27) $`=`$ $`(\stackrel{}{D^{3+}})^\nu (D^3)^\rho 2\stackrel{~}{c}^3\stackrel{~}{B}_+^{3\nu \rho }`$ #### off-diagonal part of $`M_{Aq}^b(M_{Aq}^f)`$ Here using $$Y^\pm =\frac{1}{2}(Y^1iY^2),$$ (b.28) we denote $`\{M_{Aq}^b\}^{Y^1q},\{M_{Aq}^b\}^{Y^2q}`$ and $`\{M_{Aq}^f\}^{Y^1q},\{M_{Aq}^f\}^{Y^2q}`$ instead of $`\{M_{Aq}^b\}^{Y^+q},\{M_{Aq}^f\}^{Y^+q}`$. The reason why we cannot denote $`\{M_{Aq}^b\}^{Y^+q},\{M_{Aq}^f\}^{Y^+q}`$ is that there are terms $`(D_\mu ^3A_\nu ^+D_\nu ^3A_\mu ^+)_+q_{}^{}{}_{2}{}^{}\overline{\sigma }^{\mu \nu }\stackrel{~}{q}_1`$ and $`i(\overline{)}D^3q_2)^{}\overline{)}A^{}\stackrel{~}{q}_1`$ exist simultaneously in $`^t(2)`$ (b.14) , for example. $$\{M_{Aq}^b\}^{A^{1\mu }q}=\frac{1}{2}(\stackrel{}{D^{3+}})^\nu \stackrel{~}{q}_1^{}\overline{\sigma }_{\mu \nu }i\frac{1}{4}\stackrel{~}{q}_1^{}\overline{\sigma }_\mu \overline{)}D^3i\frac{1}{4}(\stackrel{}{D^{3+}})^\mu \stackrel{~}{q}_1^{}$$ (b.29) $$\{M_{Aq}^b\}^{A^{2\mu }q}=i\frac{1}{2}(\stackrel{}{D^{3+}})^\nu \stackrel{~}{q}_1^{}\overline{\sigma }_{\mu \nu }\frac{1}{4}\stackrel{~}{q}_1^{}\overline{\sigma }_\mu \overline{)}D^3\frac{1}{4}(\stackrel{}{D^{3+}})^\mu \stackrel{~}{q}_1^{}$$ (b.30) $$\{M_{Aq}^f\}^{A^{1\mu }q}=\frac{1}{2}(\stackrel{}{D^{3+}})^\nu \stackrel{~}{q}_1^{}\overline{\sigma }_{\mu \nu }i\frac{1}{4}\stackrel{~}{q}_1^{}\overline{\sigma }_\mu \overline{)}D^3i\frac{1}{4}(\stackrel{}{D^{3+}})^\mu \stackrel{~}{q}_1^{}$$ (b.31) $$\{M_{Aq}^b\}^{A^{2\mu }q}=i\frac{1}{2}(\stackrel{}{D^{3+}})^\nu \stackrel{~}{q}_1^{}\overline{\sigma }_{\mu \nu }\frac{1}{4}\stackrel{~}{q}_1^{}\overline{\sigma }_\mu \overline{)}D^3\frac{1}{4}(\stackrel{}{D^{3+}})^\mu \stackrel{~}{q}_1^{}$$ (b.32) $$\{M_{Aq}^b\}^{B^{1\nu \rho }q}=i\frac{1}{2}\stackrel{~}{c}^3\stackrel{~}{q}_1^{}\overline{\sigma }^{\nu \rho }+\frac{1}{4}\stackrel{~}{B}^{3\nu \rho }\stackrel{~}{q}_1^{}+i\stackrel{~}{B}_{+\mu }^{3\rho }\overline{\sigma }^{\mu \nu }\stackrel{~}{q}_1^{}$$ (b.33) $$\{M_{Aq}^b\}^{B^{2\nu \rho }q}=\frac{1}{2}\stackrel{~}{c}^3\stackrel{~}{q}_1^{}\overline{\sigma }^{\nu \rho }i\frac{1}{4}\stackrel{~}{B}^{3\nu \rho }\stackrel{~}{q}_1^{}\stackrel{~}{B}_{+\mu }^{3\rho }\overline{\sigma }^{\mu \nu }\stackrel{~}{q}_1^{}$$ (b.34) $$\{M_{Aq}^f\}^{B^{1\nu \rho }q}=i\frac{1}{2}\stackrel{~}{c}^3\stackrel{~}{q}_1^{}\overline{\sigma }^{\nu \rho }+\frac{1}{4}\stackrel{~}{B}^{3\nu \rho }\stackrel{~}{q}_1^{}i\stackrel{~}{B}_{+\mu }^{3\rho }\overline{\sigma }^{\mu \nu }\stackrel{~}{q}_1^{}$$ (b.35) $$\{M_{Aq}^f\}^{B^{2\nu \rho }q}=\frac{1}{2}\stackrel{~}{c}^3\stackrel{~}{q}_1^{}\overline{\sigma }^{\nu \rho }i\frac{1}{4}\stackrel{~}{B}^{3\nu \rho }\stackrel{~}{q}_1^{}+\stackrel{~}{B}_{+\mu }^{3\rho }\overline{\sigma }^{\mu \nu }\stackrel{~}{q}_1^{}$$ (b.36) $$\{M_{Aq}^b\}^{c^1q}=i\frac{1}{2}\stackrel{~}{B}_+^{3\mu \nu }\overline{\sigma }_{\mu \nu }\stackrel{~}{q}_1^{}\frac{1}{4}\stackrel{~}{c}^3\stackrel{~}{q}_1^{}$$ (b.37) $$\{M_{Aq}^b\}^{c^2q}=\frac{1}{2}\stackrel{~}{B}_+^{3\mu \nu }\overline{\sigma }_{\mu \nu }\stackrel{~}{q}_1^{}+i\frac{1}{4}\stackrel{~}{c}^3\stackrel{~}{q}_1^{}$$ (b.38) $$\{M_{Aq}^f\}^{c^1q}=i\frac{1}{2}\stackrel{~}{B}_+^{3\mu \nu }\overline{\sigma }_{\mu \nu }\stackrel{~}{q}_1^{}\frac{1}{4}\stackrel{~}{c}^3\stackrel{~}{q}_1^{}$$ (b.39) $$\{M_{Aq}^f\}^{c^2q}=\frac{1}{2}\stackrel{~}{B}_+^{3\mu \nu }\overline{\sigma }_{\mu \nu }\stackrel{~}{q}_1^{}+i\frac{1}{4}\stackrel{~}{c}^3\stackrel{~}{q}_1^{}$$ (b.40) For (b.29)$``$(b.40), one can fine the relation $$\left(\begin{array}{c}\{M_{Aq}^f\}^{Y^1q}\\ \{M_{Aq}^f\}^{Y^2q}\end{array}\right)=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right)\left(\begin{array}{c}\{M_{Aq}^b\}^{Y^1q}\\ \{M_{Aq}^b\}^{Y^2q}\end{array}\right).$$ (b.41) Using above explicit matrix elements (LABEL:a.20)$``$(b.40), we can perform the path integral (LABEL:a.7) directly, instead of neglecting non-kinetic off-diagonal part of $`M^b(2)(M^f(2))`$. Then we have a crucial obstacle from the difference between (b.23) and (b.24), while the obstacle from (b.29)$``$(b.40) is resolved by the relation (b.41). This obstacle tells us that the contributions from (b.23) and (b.24) is not $`1`$ and that the result (3.67) is effective up to order of square of $`\stackrel{~}{q}_1`$. (In fact this problem does not appear when we treat adjoint matter instead of fundamental matter. Thus we think that this problem comes from the choice of the representation of matter fields.) However the contributions from (b.23) and (b.24) becomes $`1`$ in $`\stackrel{~}{q}_10`$ limit after path integration. Thus we estimate that the contributions from (b.23) and (b.24) to be $`1`$ in the case that the result $`Z_{m,c,k}^t(2)`$ (3.67) is topological. This is why it is enough to estimate the path integral with the indexes that the only kinetic terms in diagonal block are counted from the big matrices in section 3.
warning/0005/astro-ph0005601.html	ar5iv	text	# Proton and electron acceleration through magnetic turbulence in relativistic outflows ## 1 Introduction Measurements of radio and optical polarization of blazar jet sources indicate that nonthermal lepton synchrotron emission is the dominant radiation mechanism at these frequencies (Begelman et al. begelman (1984)). The radio through X-ray afterglow and the prompt soft-gamma ray emission observed (Costa et al. costa (1997), van Paradijs et al. vanparadijs (1997), Djorgovski et al. djorgovski (1997), Frail et al. frail (1998)) from gamma-ray bursts (GRBs) is attributed to the same process (Tavani tavani (1996), Paczyński & Rhoads paczynski (1993), Mészáros & Rees mr93 (1993)). In the standard blast-wave physics that has been developed (Mészáros & Rees meszaros97 (1997); Vietri vietri97 (1997), Waxman waxman97 (1997)) to explain the long wavelength afterglow emission from GRBs, particles captured from the external medium energize the blast wave and cause blast wave deceleration. When a blast wave sweeps up material from the surrounding medium, the free energy of the captured particles initially resides in the more massive protons, nucleons, or charged grains and dust, if the latter can survive capture by a relativistic wind, as assumed here. The less massive nonthermal electrons and positrons emit, however, most of the radiant energy. In order that high radiative efficiencies are possible, it is therefore required that there is efficient transfer of energy from the swept-up massive particles to electrons, or to very high-energy protons emitting synchrotron radiation or photomeson secondaries. Major uncertainties in modeling relativistic outflows involve magnetic field generation and the mechanism for dissipating and transferring the free energy of protons, ions and dust particles to lighter particles (Mészáros, Rees & Papathanassiou meszaros93 (1993), Chiang & Dermer chiang (1999)). Nonthermal electrons and protons may both undergo first-order Fermi shock acceleration in relativistic blast waves to form power-law distributions reaching to very high energies, and the origin of ultra-high energy cosmic rays (UHECRs) has been attributed to shock acceleration of protons in the GRB blast waves (Vietri vietri (1995), Waxman waxman (1995)). This mechanism has recently been called into question (Gallant & Achterberg gallant (1999)). Following the first shock crossing by a particle, subsequent cycles do not permit large gains of energy because the particle is captured by the shock before it has been scattered through a large angle. A stochastic mechanism for particle energization in relativistic shocks avoids these difficulties, but requires strong magnetic fields to retain UHECRs within the shell (Waxman waxman (1995)). However, the detailed physics of stochastic gyroresonant acceleration was not considered. In this paper, we show that the magnetic turbulence generated by the capture of heavy charged particles can accelerate lighter charged particles to high energies through stochastic gyroresonant acceleration. This could account for the hard spectra observed from GRBs and blazars during their flaring states. The acceleration of particles to high energies could produce UHECRs through gyroresonant acceleration in relativistic outflows. The question of the survivability of dust against sublimation by the intense radiation fields of the GRB is addressed in the Appendix. ## 2 Particle energization in relativistic outflows ### 2.1 The relativistic pick-up model A two-stream instability is formed when a blast wave with an entrained magnetic field encounters a medium of density $`n_i^{}`$. According to the recent work of Pohl & Schlickeiser (pohl (2000), hereafter referred to as PSpohl (2000)), this causes the captured protons and electrons to rapidly isotropize in the blast wave plasma. The isotropization process operates on the comoving time scale (eq. (74) of PSpohl (2000)) $`t_{\mathrm{iso}}(\mathrm{s})0.7n_{b,8}^{1/2}/\mathrm{\Gamma }_{300}n_i^{}`$, where $`\mathrm{\Gamma }=300\mathrm{\Gamma }_{300}`$ is the bulk Lorentz factor of the blast wave and $`n_{b,8}`$ denotes the comoving blastwave density in units of $`10^8`$ protons cm<sup>-3</sup>. The blast wave particle density in the comoving frame is given by $`n_b(\mathrm{cm}^3)`$ $`E/(\mathrm{\Gamma }_0m_pc^24\pi r^2\mathrm{\Delta }r)`$$`10^8E_{54}/r_{16}^3`$, and is implied by depositing an initial energy $`10^{54}E_{54}`$ ergs/(4$`\pi `$ sr) in a blastwave shell with comoving radius $`\mathrm{\Delta }rr/\mathrm{\Gamma }_0`$, where $`\mathrm{\Gamma }_0`$ is the initial blastwave Lorentz factor. Alfvenic magnetic turbulence is generated by the incoming protons and electrons with an energy density (eq. (53) of PSpohl (2000)) $`\mathrm{\Delta }U_A\beta _Am_pc^2n_i^{}\mathrm{\Gamma }\sqrt{\mathrm{\Gamma }^21}`$, where $`\beta _Ac`$ is the Alfven speed. The total energy density of the swept-up material in the comoving frame is $`\mathrm{\Delta }U_i=m_pc^2n_i^{}\mathrm{\Gamma }\sqrt{\mathrm{\Gamma }^21}`$ (Blandford & McKee bm76 (1976)). Thus the fraction of incoming energy converted to transverse plasma wave turbulence in the process of isotropizing the particles is $`\beta _A`$. PSpohl (2000) calculated the emitted radiation in the thick-target limit when the captured, isotropized protons undergo nuclear inelastic collisions with thermal particles in the blastwave. Although this process involves the fraction $`(1\beta _A)`$ of the incoming energy retained in the protons, it operates on the comoving pion production time scale $`t_{pp}(\mathrm{s})1.4\times 10^7/n_{b,8}`$, where $`n_{b,8}`$ denotes the comoving blastwave density in units of $`10^8`$ protons cm<sup>-3</sup>. If the acceleration rate is sufficiently rapid, pion production could produce low-level prompt and extended emission in GRBs or blazars if $`n_b`$ remains at the estimated level. Moreover, luminous radiation from pion-decay electrons and positrons is only produced if $`t_{pp}`$ is much shorter than the diffusive proton escape time scale from the blast wave, given by $`t_{\mathrm{esc}}(\mathrm{\Delta }r)^2/\kappa =r^2/(\mathrm{\Gamma }_0ct_{\mathrm{iso}})^2`$, noting that the spatial diffusion coefficient $`\kappa c^2t_{\mathrm{iso}}`$. This requires that $`n_{b,8}^{1/2}n_i^{}r_{16}^2>>15\mathrm{\Gamma }_{300}`$, which holds well in jets but is only marginally satisfied in GRB blast waves. ### 2.2 The role of captured charged dust Here we extend the work of PSpohl (2000) by considering the process of channeling the plasma turbulence energy into the energy of the swept-up primary electrons through gyroresonant interactions. As demonstrated in that paper, incoming charged particles with mass $`m_j`$, charge $`Q_je`$, and Lorentz factor $`\mathrm{\Gamma }`$ generate Alfvenic turbulence with parallel wavenumbers $`\|k\|>R_j^1`$, where the Larmor radius $`R_j=m_jc^2\sqrt{\mathrm{\Gamma }^21}/(Q_j\|e\|B_0)`$ and $`B_0`$ is the entrained magnetic field in the frame of the blastwave. For simplicity, we assume that the magnetic field is parallel to the blast wave velocity, and generalize the calculation of PSpohl (2000) by including the effects of negatively charged dust. According to eqs. (48)-(50) of PSpohl (2000), we obtain the power spectrum of backward ($``$) and forward (+) propagating Alfven waves given by $`I_{}(k)\|Z(k)\|`$ and $`I_+(k)I_0^2(k)/\|Z(k)\|`$, where $`I_0(k)`$ is the initial turbulence spectrum before isotropization, and $$Z(k)=\frac{1}{2}\frac{\omega _{p,i}}{c}\frac{n_p}{n_b}\frac{B^2}{k^2}$$ $$\times \underset{j=e,p,d}{}\frac{Q_jn_j}{n_p}[1(R_jk)^1]H[\|k\|R_j^1].$$ (1) The term $`\omega _{p,i}(\mathrm{s}^1)=1.3\times 10^3\sqrt{n_b}`$ is the proton plasma frequency and $`H[x]`$ denotes the Heaviside step function. This expression assumes $`I_0(k)\|Z(k)\|`$ and makes use of the charge-neutrality condition $`Q_dn_d+n_e=n_p=\mathrm{\Gamma }n_i^{}`$, where we assume that all dust particles have the same charge. A number of processes influence the charging of dust grains in a proton-electron-plasma. In the absence of any other process due to the higher electron mobility an initially uncharged dust particle will absorb more electrons than protons and thus attain a negative potential of order $`2.5k_BT_e/e`$ (Spitzer 1968). However, a number of competing processes like neutral-atom collisions and ion impacts on the grain surface lead to electron loss from the dust grains. Also, a significant ultraviolet photon flux can charge the grain by the photoelectric effect to a positive potential. All these processes, together with poorly known grain properties, make an exact determination of the grain charge impossible (see e.g. McKee et al. (1977) for the variety of grain properties resulting from model calculations under various assumptions). Here we follow the argumentation of Ellison, Drury & Meyer (1997). If the grain potential is $`\varphi `$, then the charge on a spherical (of radius $`a`$) grain is of order $`q=eQ_d4\pi ϵ_0a\varphi `$, so that $`Q_d700(a/(10^7\mathrm{m}))(\varphi /10\mathrm{V})`$. The number of atoms in the grain will be of order $`[a/(10^{10}\mathrm{m})]^3`$, so that with the mean atomic weight $`\mu `$ of the grain atoms the entire grain atomic weight is $`A_G=\mu [a/(10^{10}\mathrm{m})]^3`$. This implies a dust mass of $`m_d=\mu m_p[a/(10^{10}\mathrm{m})]^3=210^{10}m_pa_7^3`$ if we adopt $`\mu =20`$ for silicate grains. Combined with the charge estimate this yields for the charge-weighted dust/proton mass ratio $`Y_d=m_d/(Q_dm_p)=310^7a_7^2`$. It is important to note that the densities of the incoming protons, electrons and dust particles are much smaller than the density of the blast wave plasma $`n_b^{}`$, so that the properties of the low-frequency Alfven waves, carried by the blast wave plasma, are not modified by the captured particles. Mathematically, this corresponds to the statement that the real part of the Alfvenic dispersion relation is solely determined by the blast wave electrons and protons. However, the free energy of beams of incoming protons, electrons and dust particles gives rise to a non-zero imaginary part of the Alfvenic dispersion relation leading to a positive growth rate of backward moving (in the blast wave plasma) Alfven waves. Moreover, for relativistic beams (Miller miller (1982)) the growth rate of longitudinal electrostatic waves is much smaller than the growth rate of transverse Alfvenic waves, so that the free energy is dissipated by isotropising the incoming beam (i.e. capture) rather than heating the blast wave plasma by plateauing. Two comments on the equilibrium wave spectrum (1) are appropriate: (1) According to the quasilinear relaxation theory the equilibrium Alfven wave spectrum given in Eq. (1) is achieved formally after infinitely long time $`t_{\mathrm{ql}}\mathrm{}`$ by transferring the free energy in the initial particle-beam-distributions into plasma waves. Numerical simulations of the electrostatic beam instability (see e.g. Grognard 1975) indicate that this asymptotic equilibrium distribution is established after $`t_{\mathrm{ql}}100t_{\mathrm{iso}}`$. (2) The wave spectrum in Eq. (1) is calculated from quasilinear wave kinetic equations that only take into account the wave growth from the unstable particle beams but neglect wave-wave interactions. Since the latter scale with the total magnetic field fluctuation energy density $`(\delta B)^2`$, this is justified as long as the total wave energy density is small to the initial energy density of the beam particles. Because according to Eqs. (53) and (64) of PSpohl (2000) $`(\delta B)^2/(4\pi w_0)=\beta _A<<1`$ this is indeed fulfilled. The turbulence spectrum (1) leads to the quasilinear isotropization time scale $`t_{\mathrm{iso}}`$ given above. This time scale and $`t_{\mathrm{ql}}=100t_{\mathrm{iso}}`$ are much shorter than the blastwave crossing time $`t_{lc}(\mathrm{s})\mathrm{\Delta }r/c1000r_{16}/\mathrm{\Gamma }_{300}`$. Fig. 1 shows the resulting turbulence spectrum of backward propagating, left-handed circularly polarized Alfven waves for a charge-weighted dust/proton density ratio $`ϵ=Q_dn_d/n_p=0.01`$, and a charge-weighted dust/proton mass ratio $`Y_d`$ $`=m_d/(Q_dm_p)`$ $`=10^6`$. The relative contributions of dust to the turbulence spectrum (1) is determined primarily by the parameters $`ϵ`$ and $`Y_d`$. Most of the magnetic turbulence energy density resides at low wave numbers and is provided by the charged dust for this value of $`ϵ`$. The total dust mass $`M_d=m_dn_d=Y_dϵM_p`$, where the total proton mass $`M_p=m_pn_p`$. Thus the total mass in dust for the calculation in Fig. 1 is $`10^4`$ times the mass in the ionized gas, as might be found in the weakly ionized, low temperature star-forming environments that harbor GRB sources. By contrast, the relative mass in dust to that in protons would be much less in the enviroments of active galactic nuclei, where the intense UV and X-ray radiation field would photo-disintegrate the dust. ### 2.3 Stochastic gyroresonant acceleration of electrons and protons The primary electrons and protons may increase their energy by stochastic gyroresonant acceleration with the turbulence spectrum given by Eq.(1). Since a fraction $`\beta _A`$ of the incoming free energy is converted into transverse plasma wave turbulence, these particles can tap at most a fraction $`\beta _A`$ of the swept-up energy via this process. Efficient energy transfer therefore requires that $`\beta _A0.01`$. The e-folding time scale for stochastic gyroresonant acceleration is $`t_{\mathrm{acc}}\beta _A^2t_{\mathrm{iso}}`$, whereas the diffusive escape time $`t_{\mathrm{esc}}R^2/\kappa =t_{lc}^2/t_{iso}`$ (Barbosa barbosa (1979), Schlickeiser schlickeiser89 (1989)). Because the turbulence spectrum in Eq.(1) is $`k^2`$, the acceleration, escape, and isotropization time scales are independent of particle energy. The comoving deceleration time scale $`t_{\mathrm{dec}}`$(s)$`=1.3\times 10^4(E_{54}/n_i^{}\mathrm{\Gamma }_{300}^5)^{1/3}`$ (Mészáros & Rees mr93 (1993)). Provided that $`\beta _A0.03n_{b,8}^{1/4}(\mathrm{\Gamma }_{300}^2/E_{54}n_i^2)^{1/3}`$, the acceleration time scale is much shorter than the time scale for evolutionary changes of the blast wave due to bulk deceleration. In this case, a steady-state calculation of the equilibrium particle energy spectrum is justified, which results from the balance of gyroresonant acceleration, diffusive escape, deceleration and radiative losses. The steady-state kinetic equation for the phase-space density $`f_e`$ of relativistic electrons is (Schlickeiser schlickeiser84 (1984)) $$\frac{1}{\gamma ^2}\frac{d}{d\gamma }[\frac{\gamma ^4}{t_\mathrm{a}}\frac{df_e}{d\gamma }+(\frac{\gamma ^3}{t_{\mathrm{dec}}}+\frac{\gamma ^4}{t_{\mathrm{syn}}})f_e]\frac{f_e}{t_\mathrm{E}}=S_e\delta (\gamma \mathrm{\Gamma })$$ (2) where $`S_e=\sqrt{\mathrm{\Gamma }^21}n_e^{}c/(4\pi \mathrm{\Gamma }r)`$ and $`A`$ is the area of the blast wave that is effective in sweeping up material. If no dust is present, then the electrons can only be accelerated to $`\gamma <\mathrm{\Gamma }R_p/R_e=\mathrm{\Gamma }m_p/m_e`$. In this range, $`t_a=t_{\mathrm{acc}}`$ and $`t_E=t_{\mathrm{esc}}`$. If dust is present, then electrons can be accelerated throughout the range $`\mathrm{\Gamma }m_p/m_e\gamma <\gamma _{\mathrm{e},\mathrm{max}}=\mathrm{\Gamma }m_d/(m_eQ_d)`$, but the acceleration and escape times are modified according to the relations $`t_a=t_{\mathrm{acc}}/ϵ`$ and $`t_E=ϵt_{\mathrm{esc}}`$. The radiative loss time scale is $`t_{\mathrm{syn}}(\mathrm{s})=6\pi m_ec/(\sigma _TB_0^2)=7.7\times 10^8B_0^2`$. Likewise, the steady-state kinetic equation for the phase-space density $`f_p`$ of relativistic protons is $$\frac{1}{\gamma ^2}\frac{d}{d\gamma }[\frac{ϵ\gamma ^4}{t_{\mathrm{acc}}}\frac{df_p}{d\gamma }+\gamma ^3(\frac{1}{t_{pp}}+\frac{1}{t_{\mathrm{dec}}})f_p]$$ $$\frac{f_p}{ϵt_{\mathrm{esc}}}=S_p\delta (\gamma \mathrm{\Gamma }),$$ (3) where $`S_p=\sqrt{\mathrm{\Gamma }^21}n_i^{}c/(4\pi \mathrm{\Gamma }r)`$. Eq.(3) holds in the presence of dust when $`\gamma <\gamma _{\mathrm{p},\mathrm{max}}`$ $`=\mathrm{\Gamma }R_d/R_p=`$ $`\mathrm{\Gamma }m_d/(m_pQ_d)`$ $`=\mathrm{\Gamma }Y_d`$, and assumes that effects of photopion production do not limit the acceleration of protons to the highest energies. Evidently, the parameter $`Y_d`$ and the bulk Lorentz factor $`\mathrm{\Gamma }`$ determine the maximum proton Lorentz factor $`\gamma _{\mathrm{p},\mathrm{max}}`$. If no dust is present, then protons are not significantly accelerated by gyroresonant processes unless long wavelength MHD turbulence is generated through processes not treated here. The exact solution to Eq.(2) for electrons accelerated by turbulence with a spectrum given by Eq.(1) is $$f_e(\mathrm{\Gamma }<\gamma <\gamma _{\mathrm{e},\mathrm{max}})=t_aS_e\frac{\stackrel{~}{\mathrm{\Gamma }}[\mu (3a)/2]}{\stackrel{~}{\mathrm{\Gamma }}[1+2\mu ]}$$ $$\frac{\mathrm{\Gamma }^{\mu +(a+1)/2}\gamma ^{\mu (a+3)/2}}{\gamma _s^{2\mu }}\mathrm{exp}(\gamma /\gamma _s)\times $$ $$M(\mu \frac{3a}{2},1+2\mu ,\frac{\mathrm{\Gamma }}{\gamma _s})U(\mu \frac{3a}{2},1+2\mu ,\frac{\gamma }{\gamma _s})$$ (4) (Schlickeiser schlickeiser84 (1984)), where $`\mu =\sqrt{\lambda +(3a)^2/4}`$, $`\lambda =t_a/t_E`$, $`a=t_a/t_{\mathrm{dec}}`$, $`\stackrel{~}{\mathrm{\Gamma }}(a_0)`$ denotes the gamma function, and $`M(a_1,a_2,a_3)`$ and $`U(a_1,a_2,a_3)`$ are the Kummer functions of the first and second kind, respectively. The value $`\gamma _s=t_{\mathrm{syn}}/t_{\mathrm{acc}}370\mathrm{\Gamma }_{300}n_i^{}/n_{b,8}^{3/2}`$ is the electron Lorentz factor where the electron synchrotron loss time scale equals the acceleration time scale. In the regime $`\gamma \mathrm{\Gamma }`$, an accurate approximation is obtained by taking the asymptotic expansion of the function $`U`$ for small arguments, resulting in $$f_e(\mathrm{\Gamma }<\gamma <\gamma _{\mathrm{e},\mathrm{max}})$$ $$\frac{t_aS_e}{2\mu }\mathrm{\Gamma }^{\mu +(1+a)/2}\gamma ^{(3+a)/2\mu }\mathrm{exp}(\gamma /\gamma _s).$$ (5) The exact solution to Eq.(3) for protons accelerated by turbulence with a spectrum given by Eq.(1) is $$f_p(\mathrm{\Gamma }<\gamma <\gamma _{\mathrm{p},\mathrm{max}})=$$ $$\frac{t_{\mathrm{acc}}S_p}{2ϵ\mu _p}\mathrm{\Gamma }^{\mu _p+(a_p+1)/2}\gamma ^{(a_p+3)/2\mu _p},$$ (6) where $`\mu _p=\sqrt{\lambda _p+(3a_p)^2/4}`$, $`\lambda _p=t_{\mathrm{acc}}/ϵ^2t_{\mathrm{esc}}`$, and $`a_p=t_{\mathrm{acc}}[t_{\mathrm{pp}}^1+t_{\mathrm{dec}}^1]/ϵ`$. In Fig. 2, we show the resulting proton and electron spectra, choosing $`\mathrm{\Gamma }_{300}=1`$, $`\beta _A=0.03`$, $`r_{16}=0.1`$, $`n_{b,8}=0.1`$, $`ϵ=0.1`$, $`Y_d=10^6`$, and $`n_p^{}=1`$ cm<sup>-3</sup>. As can be seen, the electron spectrum $`dN_e/d\gamma \gamma ^1\mathrm{exp}(\gamma /11700)`$ is very flat below the synchrotron cutoff, whereas the proton spectrum displays a spectrum $`dN_p/d\gamma \gamma ^{1.36}`$ below the maximum Lorentz factor $`\gamma _{\mathrm{p},\mathrm{max}}=3\times 10^8`$. For these parameters, we find comoving acceleration times of $`330(\beta _A/0.03)^2`$ s for electrons, and $`10`$ times larger for the protons. The choosen value for $`\beta _A`$ implies a magnetic field strength of 14 Gauss. Larger values of $`\beta _A`$ will shorten the acceleration time scales and improve the accelerated particle efficiency. Requiring that the gyroradii of the swept-up dust particles $`R_d=Y_dR_p`$ in a magnetic field of strength $`14B_{14}`$ Gauss is smaller than the comoving shell radius $`10^{14}r_{16}\mathrm{\Gamma }_{100}^1`$ cm limits the charge-weighted dust/proton mass ratio to values smaller than $`Y_d<510^6r_{16}B_{14}\mathrm{\Gamma }_{100}^2`$, and consequently the maximum proton Lorentz factor in the comoving blast wave frame is $`\gamma _{p,\mathrm{max}}=\mathrm{\Gamma }Y_d<510^8r_{16}B_{14}\mathrm{\Gamma }_{100}^1`$. Fig. 3 shows the resulting synchrotron emission spectrum produced by the electron distribution function (5). For the parameters used here, the synchrotron spectrum is very hard. Particle spectral indices $`dN/d\gamma \gamma ^p`$ with $`p1`$ can be produced through this process, leading to synchrotron spectra $`F_\nu \nu ^\alpha `$ with $`\alpha =1/3`$, corresponding to the index of the elementary synchrotron emissivity formula for monoenergetic electrons. In standard shock acceleration theory (Blandford & Eichler 1987), particles are accelerated with $`p2`$, leading to values of $`\alpha 1/2`$. Harder synchrotron spectra can however be produced through nonlinear effects in first-order Fermi acceleration, or by introducing a low-energy cutoff in the electron distribution function. Synchrotron spectra from flaring blazar sources (Catanese et al. catanese (1997), Cohen et al. cohen (1997), Pian et al. pian (1998)) and GRBs display very hard spectra requiring nonthermal lepton spectra with $`p<2`$, consistent with the stochastic Fermi acceleration mechanism proposed here. The hardening above the ankle of the cosmic ray spectrum and beyond the ZGK cutoff could be related to a hard UHECR component formed through stochastic acceleration in GRB blast waves (Waxman waxman (1995), Takeda et al. takeda (1998)). Because turbulence is generated in the process of capturing and isotropizing particles from the external environoment, a wave spectrum that can yield very high energy electrons and protons through gyroresonant acceleration must be present in the relativistic blast waves. If dust survives capture by the magnetized plasma wind, then long wavelength turbulence can be generated that could accelerate protons to ultra-high energies on time scales shorter than the duration over which the blast wave decelerates to nonrelativistic speeds. Gyroresonant acceleration of particles by MHD turbulence provides a mechanism to transfer energy from heavier to lighter species, and overcomes difficulties associated with accelerating UHECRs in relativistic blast waves through the first-order Fermi mechanism. ## 3 Summary and conclusions We have shown that efficient and rapid stochastic gyroresonant acceleration of electrons and protons occurs when interstellar protons and dust are captured by a relativistic outflow. The MHD plasma wave turbulence, created within the blast wave plasma by the penetrating protons and charged dust particles, is transferred from the heavier to the lighter particles through the gyroresonant acceleration process. This mechanism is particularly attractive for the energization of nonthermal electrons in GRB sources, as radiation modelling often requires approximate energy equipartition between electrons and protons (e.g., Katz katz (1994), Beloborodov & Demianski beloborodov (1995), Smolsky & Usov smolsky (1996), Chiang & Dermer chiang (1999)). Capture by and isotropization within a blast wave, as discussed by Pohl & Schlickeiser (pohl (2000)), provides primary electrons with an energy 1836 times smaller than that of the protons. As shown here, the subsequent acceleration of electrons in the proton-induced turbulence leads to approximate energy equipartition of these nonthermal electrons. Moreover, if in addition charged dust particles penetrate the blast wave plasma, we find that further acceleration of the particles is possible to an energy determined by the radiative loss processes of the electrons and protons. For some parameter regimes, this could accelerate particles to $`10^{20}`$ eV and account for UHECR production. By solving the appropriate kinetic equations, we calculated the resulting particle energy spectra of the accelerated protons and electrons, and the observable synchrotron radiation spectrum of the accelerated electrons. The latter peaks at hard X-ray energies and is very flat, in agreement with observed photon spectra of GRBs and some blazars during their flaring states. A crucial issue, particularly for the stochastic acceleration of protons, is the penetration of dust. While large amounts of dust are present in the environments of star-forming regions in galaxies, it is also necessary that the dust not be destroyed during capture by the blast wave or by the intense radiation associated with captured nonthermal particles in the blast wave. We address some aspects of this issue in Appendix A. ###### Acknowledgements. RS gratefully acknowledges partial support of his work by the Deutsche Forschungsgemeinschaft through Sonderforschungsbereich 191 and by the Verbundforschung, grant DESY-05AG9PCA. The work of CD is supported by the Office of Naval Research and NASA grant DPR S-13756G. ## 4 Appendix A: Survivability of dust in GRB radiation field A major uncertainty in this study is whether the dust can survive the intense GRB radiation field by sublimation before the blast wave arrives and sweeps up the dust particle. In the inhomogeneous environment surrounding a massive star/GRB progenitor (Dermer & Böttcher db00 (2000)), it is probable that the blast wave will impact a dense dusty cloud after passing through a nearly evacuated region. In that case, there will be little radiation in advance of the blast wave, because the blast wave will only start to be energized and radiate after it encounters the cloud. Under such circumstances, the survivability of a large fraction of the dust to sublimation is assured. In the ideal case where a GRB blast wave passes through a uniform density medium before intercepting dust particles, we outline a calculation of dust sublimation and identify parameter sets that are compatible with dust survivability. We follow the treatment of Waxman & Draine (wd99 (1999)) for dust sublimation. Although reverse shock emission could be important for dust sublimation, we neglect it here. It has only been detected from GRB 990123 and can be strongly suppressed if, for example, the blast wave shell is thin and the reverse shock traverses the shell before becoming relativistic (Sari et al. snp96 (1996)). In this case the optical-UV emission irradiating the dust, which is the photon energy range most important for dust sublimation (Waxman & Draine wd99 (1999)), can be described by the $`L_\nu `$ = $`L_{\nu ,\mathrm{max}}(\nu /\nu _0)^{1/3}`$ portion of the synchrotron emissivity spectrum where the quantity $`L_{\nu ,\mathrm{max}}`$ = $`4\pi n_i^{}x^3m_ec^2\sigma _\mathrm{T}\mathrm{\Gamma }B_0/9e`$ (Sari et al. spn98 (1998)), where $`e`$ is the proton charge. Using standard blast-wave physics, $`B_0`$(G) $`=0.388\mathrm{\Gamma }\sqrt{n_i^{}e_B}`$, where $`e_B<1`$ is the magnetic-field parameter. To give good fits to the prompt phase of GRBs, it is necessary that the electrons are in the weakly cooled regime (Chiang & Dermer chiang (1999), Dermer et al. dbc00 (2000)). For a weakly cooled spectrum, $`\nu _0=\nu _m=\mathrm{\Gamma }\gamma _m^2eB/2\pi m_ec`$, where $`\gamma _m=e_e\mathrm{\Gamma }(p_e2)m_p/[m_e(p_e1)]`$ and $`e_e>0.1`$. Taking the electron injection index $`p_e2.2`$ \- 3.0 = 2.5, as implied by burst spectroscopy in the prompt and afterglow phase, we find $`\nu _m`$(Hz) =$`4.1\times 10^{11}e_e^2\mathrm{\Gamma }^4\sqrt{n_i^{}e_B}`$ (see Sari et al. spn98 (1998), Dermer et al. dbc00 (2000) for more details). The luminosity (ergs s<sup>-1</sup>) between the 1 eV $`=2.41\times 10^{14}`$ Hz and 7.5 eV band is $$L_{17.5}1.28\times 10^5n_i^{4/3}(e_B/e_e^2)^{1/3}x^3\mathrm{\Gamma }^{2/3}$$ (7) $$2.0\times 10^{46}(\frac{n_i^{}}{100})^{4/3}(\frac{e_B}{10^4})^{1/3}(\frac{e_e}{0.5})^{2/3}(\frac{x}{10^{16}\mathrm{cm}})^3(\frac{\mathrm{\Gamma }}{300})^{2/3},$$ which also defines our standard parameter set. The relationship of $`\mathrm{\Gamma }`$ and $`x`$ to observer time $`t`$ for an adiabatic blast wave decelerating in a uniform surrounding medium is $$\frac{\mathrm{\Gamma }}{\mathrm{\Gamma }_0}=\frac{1}{\sqrt{1+(4\tau )^{3/4}}},\mathrm{and}\frac{x}{x_d}\frac{\tau }{1+4^{1/4}\tau ^{3/4}},$$ (8) where the deceleration radius $`x_d=(3E_0/8\pi \mathrm{\Gamma }_0^2m_pc^2n_i^{})^{1/3}`$ $``$ $`2.1\times 10^{16}(E_{54}/\mathrm{\Gamma }_{300}^2n_2)^{1/3}`$ cm, the dimensionless time $`\tau =t/t_d=ct\mathrm{\Gamma }_0^2c/x_d`$, and the apparent isotropic explosion energy is $`10^{54}E_{54}`$ ergs. For the sublimation of dust, we consider the total energy change due to (1) heating by the GRB in the 1 - 7.5 eV band, (2) thermal reradiation by dust, and (3) the energy carried away by the sublimed particles. In this approximation, the heating rate $`dT/dt`$ due to dust energization is given through $$[4\pi a^3k_\mathrm{B}(\frac{\rho }{m})]^1\frac{dT}{dt}=\frac{dE}{dt}=\frac{L_{17.5}(t)a^2(t)}{4\pi r^2}$$ $$4\pi a^2(t)\sigma _{\mathrm{SB}}T^412\pi k_\mathrm{B}T(\frac{\rho }{m})a^2(t)\frac{da}{dt}.$$ (9) The rate of shrinkage of dust grain radius $`a(t)`$ with time $`t`$, is given by $$\frac{da}{dt}=(\frac{m}{\rho })^{1/3}\nu _0\mathrm{exp}(b/k_\mathrm{B}T),$$ (10) (Waxman & Draine wd99 (1999)), where $`\rho /m=10^{23}`$ cm<sup>-3</sup>, $`\nu _0=1\times 10^{15}`$ s<sup>-1</sup>, and $`b/k_\mathrm{B}=7\times 10^4`$ K for refractory grains. We look for solutions of Eqs. (7)- (10) where the dust grain is closer than the blast wave to the GRB source at the moment the dust becomes completely sublimed. There are no interesting solutions for our standard parameters; in this case, dust at $`10^{16}`$ cm becomes sublimed when the blast wave has only reached $`x=1.6\times 10^{15}`$ cm. It is easy, however, to find solutions for slightly less energetic GRBs in more tenuous environments. For example, the blast wave will reach dust particles before they sublime if the dust resides within $`x9\times 10^{15}`$ cm when $`E_0=10^{52}`$ ergs and $`n_i^{}=1`$ cm<sup>-3</sup>. These values of the energy and density are not too different than the values found when modeling the afterglow of GRB 980425 (Wijers & Galama wg99 (1999)). Thus we find that dust will survive destruction by sublimation in some, though not all, GRB environments if the surrounding medium is uniform. A clumpy surrounding medium improves the chances of dust survival.
warning/0005/quant-ph0005045.html	ar5iv	text	# Supersymmetric and Shape-Invariant Generalization for Non-resonant Jaynes-Cummings Systems ## I Introduction The integrability condition called shape-invariance originates in supersymmetric quantum mechanics . The separable positive-definite Hamiltonian $`\widehat{H}_1=\widehat{A}^{}\widehat{A}`$ is called shape-invariant if the condition $$\widehat{A}(a_1)\widehat{A}^{}(a_1)=\widehat{A}^{}(a_2)\widehat{A}(a_2)+R(a_1),$$ (1) is satisfied . In this equation $`a_1`$ and $`a_2`$ represent parameters of the Hamiltonian. The parameter $`a_2`$ is a function of $`a_1`$ and the remainder $`R(a_1)`$ is independent of the dynamical variables such as position and momentum. Even though not all exactly-solvable problems are shape-invariant , shape invariance, especially in its algebraic formulation , has proven to be a powerful technique to study exactly-solvable systems. In a previous paper we used shape-invariance to calculate the energy eigenvalues and eigenfunctions for the Hamiltonian $$\widehat{𝐇}=\widehat{A}^{}\widehat{A}+\frac{1}{2}[\widehat{A},\widehat{A}^{}]\left(\widehat{\sigma }_3+1\right)+\sqrt{\mathrm{}\mathrm{\Omega }}\left(\widehat{\sigma }_+\widehat{A}+\widehat{\sigma }_{}\widehat{A}^{}\right),$$ (2) where $$\widehat{\sigma }_\pm =\frac{1}{2}\left(\widehat{\sigma }_1\pm i\widehat{\sigma }_2\right),$$ (3) and $`\widehat{\sigma }_i`$, with $`i=1,\mathrm{\hspace{0.17em}2},\mathrm{and}\mathrm{\hspace{0.17em}\hspace{0.17em}3}`$, are the Pauli matrices. This is a generalization of the Jaynes-Cummings Hamiltonian . A different, but related problem was considered in Ref. . Our goal in this paper is to study a further generalization of the Jaynes-Cummings Hamiltonian by introducing a term proportional to $`\sigma _3`$ with an arbitrary coefficient (the so-called non-resonant limit). In addition to the energy levels we study the time-evolution and the population inversion factor. Introducing the similarity transformation that replaces $`a_1`$ with $`a_2`$ in a given operator $$\widehat{T}(a_1)\widehat{O}(a_1)\widehat{T}^{}(a_1)=\widehat{O}(a_2)$$ (4) and the operators $$\widehat{B}_+=\widehat{A}^{}(a_1)\widehat{T}(a_1)$$ (5) $$\widehat{B}_{}=\widehat{B}_+^{}=\widehat{T}^{}(a_1)\widehat{A}(a_1),$$ (6) the condition of Eq. (1) can be written as a commutator $$[\widehat{B}_{},\widehat{B}_+]=\widehat{T}^{}(a_1)R(a_1)\widehat{T}(a_1)R(a_0),$$ (7) where we used the identity $$R(a_n)=\widehat{T}(a_1)R(a_{n1})\widehat{T}^{}(a_1),$$ (8) valid for any $`n`$. The ground state of the Hamiltonian $`\widehat{H}_1=\widehat{A}^{}\widehat{A}=\widehat{B}_+\widehat{B}_{}`$ satisfies the condition $$\widehat{A}\psi _0=0=\widehat{B}_{}\psi _0;$$ (9) and the unnormalized $`n`$-th excited state is given by $$\psi _n\left(\widehat{B}_+\right)^n\psi _0$$ (10) with the eigenvalue $$_n=\underset{k=1}{\overset{n}{}}R(a_k).$$ (11) We note that the Hamiltonian of Eq. (2) can also be written as $$\widehat{𝐇}=\left[\begin{array}{cc}\widehat{T}& 0\\ 0& \pm 1\end{array}\right]\widehat{𝐡}_\pm \left[\begin{array}{cc}\widehat{T}^{}& 0\\ 0& \pm 1\end{array}\right],$$ (12) where $$\widehat{𝐡}_\pm =\widehat{B}_+\widehat{B}_{}+\frac{1}{2}R(a_0)\left(\widehat{\sigma }_3+1\right)\pm \sqrt{\mathrm{}\mathrm{\Omega }}\left(\widehat{\sigma }_+\widehat{B}_{}+\widehat{\sigma }_{}\widehat{B}_+\right).$$ (13) ## II The Generalized Non-resonant Jaynes-Cummings Hamiltonian The standard Jaynes-Cummings model, normally used in quantum optics, idealizes the interaction of matter with electromagnetic radiation by a simple Hamiltonian of a two-level atom coupled to a single bosonic mode. This Hamiltonian has a fundamental importance to the field of quantum optics and it is a central ingredient in the quantized description of any optical system involving the interaction between light and atoms. The Jaynes-Cummings Hamiltonian defines a molecule, a composite system formed from the coupling of a two-state system and a quantized harmonic oscillator. In this case, its non-resonant expression can be written as $$\widehat{𝐇}=\widehat{A}^{}\widehat{A}+\frac{1}{2}[\widehat{A},\widehat{A}^{}]\left(\widehat{\sigma }_3+1\right)+\sqrt{\mathrm{}\mathrm{\Omega }}\left(\widehat{\sigma }_+\widehat{A}+\widehat{\sigma }_{}\widehat{A}^{}\right)+\mathrm{}\mathrm{\Delta }\widehat{\sigma }_3,$$ (14) where $`\mathrm{\Omega }`$ is a constant related with the coupling strength and $`\mathrm{\Delta }`$ is a constant related with the detuning of the system. Following Ref. we introduce the operator $$\widehat{𝐒}=\widehat{\sigma }_+\widehat{A}+\widehat{\sigma }_{}\widehat{A}^{},$$ (15) where the operators $`\widehat{A}`$ and $`\widehat{A}^{}`$ satisfy the shape invariance condition of Eq. (1). Using this definition we can decompose the non-resonant Jaynes-Cummings Hamiltonian in the form $$\widehat{𝐇}=\widehat{𝐇}_o+\widehat{𝐇}_{int},$$ (16) where $`\widehat{𝐇}_o=\widehat{𝐒}^2,`$ (18) $`\widehat{𝐇}_{int}=\sqrt{\mathrm{}\mathrm{\Omega }}\widehat{𝐒}+\mathrm{}\mathrm{\Delta }\widehat{\sigma }_3.`$ (19) First, we search for the eigenstates of $`\widehat{𝐒}^2`$. In this case it is more convenient to work with its $`B`$-operator expression, which can be written as $$\widehat{𝐒}^2=\left[\begin{array}{cc}\widehat{T}& 0\\ 0& \pm 1\end{array}\right]\left[\begin{array}{cc}\widehat{B}_{}\widehat{B}_+& 0\\ 0& \widehat{B}_+\widehat{B}_{}\end{array}\right]\left[\begin{array}{cc}\widehat{T}^{}& 0\\ 0& \pm 1\end{array}\right]\left[\begin{array}{cc}\widehat{H}_2& 0\\ 0& \widehat{H}_1\end{array}\right],$$ (20) where $`\widehat{H}_2=\widehat{T}\widehat{B}_{}\widehat{B}_+\widehat{T}^{}`$. Note the freedom of sign choice in Eq. (20), which results in two possible decompositions of $`\widehat{𝐒}^2`$. Next, we introduce the states $$\mathrm{\Psi }^{(\pm )}=\left[\begin{array}{cc}\widehat{T}& 0\\ 0& \pm 1\end{array}\right]\left[\begin{array}{c}C_m^{(\pm )}m\\ C_n^{(\pm )}n\end{array}\right],$$ (21) where $`C_{m,n}^{(\pm )}C_{m,n}^{(\pm )}[R(a_1),R(a_2),R(a_3),\mathrm{}]`$ are auxiliary coefficients and, $`m`$ and $`n`$ are the abbreviated notation for the states $`\psi _m`$ and $`\psi _n`$ of Eq. (10). Using Eqs. (7), (20) and (21), the commutation between $`\widehat{H}_1`$ and a function of $`R(a_k)`$, and the $`\widehat{T}`$-operator unitary condition, one gets $`\widehat{𝐒}^2\mathrm{\Psi }^{(\pm )}`$ $`=`$ $`\left[\begin{array}{cc}\widehat{T}& 0\\ 0& \pm 1\end{array}\right]\left[\begin{array}{cc}\widehat{B}_+\widehat{B}_{}+R(a_0)& 0\\ 0& \widehat{B}_+\widehat{B}_{}\end{array}\right]\left[\begin{array}{c}C_m^{(\pm )}m\\ C_n^{(\pm )}n\end{array}\right]`$ (22) $`=`$ $`\left[\begin{array}{cc}\widehat{T}& 0\\ 0& \pm 1\end{array}\right]\left[\begin{array}{cc}_m+R(a_0)& 0\\ 0& _n\end{array}\right]\left[\begin{array}{c}C_m^{(\pm )}m\\ C_n^{(\pm )}n\end{array}\right].`$ (23) However, using Eqs. (8) and (11) one can write $`\widehat{T}\left[_m+R(a_0)\right]\widehat{T}^{}`$ $`=`$ $`\widehat{T}\left[R(a_1)+R(a_2)+\mathrm{}+R(a_m)+R(a_0)\right]\widehat{T}^{}`$ (24) $`=`$ $`R(a_2)+R(a_3)+\mathrm{}+R(a_{m+1})+R(a_1)=_{m+1}.`$ (25) Hence the states $$\mathrm{\Psi }_m^{(\pm )}=\left[\begin{array}{cc}\widehat{T}& 0\\ 0& \pm 1\end{array}\right]\left[\begin{array}{c}C_m^{(\pm )}m\\ C_{m+1}^{(\pm )}m+1\end{array}\right],m=0,1,2,\mathrm{}$$ (26) are the normalized eigenstates of the operator $`\widehat{𝐒}^2`$ $$\widehat{𝐒}^2\mathrm{\Psi }_m^{(\pm )}=_{m+1}\mathrm{\Psi }_m^{(\pm )}.$$ (27) We observe that the orthonormality of the wavefunctions imply in the following relations among the $`C`$’s: $`\mathrm{\Psi }_m^{(\pm )}\mathrm{\Psi }_m^{(\pm )}`$ $`=`$ $`\left[C_m^{(\pm )}\right]^2+\left[C_{m+1}^{(\pm )}\right]^2=1`$ (29) $`\mathrm{\Psi }_m^{()}\mathrm{\Psi }_m^{(\pm )}`$ $`=`$ $`C_m^{(\pm )}C_m^{()}C_{m+1}^{(\pm )}C_{m+1}^{()}=0.`$ (30) Since $`\widehat{𝐒}^2`$ and $`\widehat{𝐇}_{int}`$ commute then it is possible to find a common set of eigenstates. We can use this fact to determine the eigenvalues of $`\widehat{𝐇}_{int}`$ and the relations among the $`C`$’s coefficients. For that we need to calculate $$\widehat{𝐇}_{int}\mathrm{\Psi }_m^{(\pm )}=\lambda _m^{(\pm )}\mathrm{\Psi }_m^{(\pm )},$$ (31) where $`\lambda _m^{(\pm )}`$ are the eigenvalues to be determined. Using Eqs. (15), (16) and (26), the last eigenvalue equation can be rewritten in a matrix form as $$\alpha \left[\begin{array}{cc}\beta & \widehat{T}\widehat{B}_{}\\ \widehat{B}_+\widehat{T}^{}& \beta \end{array}\right]\left[\begin{array}{cc}\widehat{T}& 0\\ 0& \pm 1\end{array}\right]\left[\begin{array}{c}C_m^{(\pm )}m\\ C_{m+1}^{(\pm )}m+1\end{array}\right]=\lambda _m^{(\pm )}\left[\begin{array}{c}C_m^{(\pm )}m\\ C_{m+1}^{(\pm )}m+1\end{array}\right],$$ (32) where $`\alpha =\sqrt{\mathrm{}\mathrm{\Omega }}`$ and $`\beta =\mathrm{}\mathrm{\Delta }/\alpha `$. Since the $`C`$’s coefficients commute with the $`\widehat{A}`$ or $`\widehat{A}^{}`$ operators, then the last matrix equation permits to obtain the following equations $`\left[\alpha \beta \lambda _m^{(\pm )}\right]\left(\widehat{T}C_m^{(\pm )}\widehat{T}^{}\right)\widehat{T}m\pm \alpha C_{m+1}^{(\pm )}\widehat{T}\widehat{B}_{}m+1=0`$ (34) $`\alpha \left(\widehat{T}C_m^{(\pm )}\widehat{T}^{}\right)\widehat{B}_+m\left[\alpha \beta +\lambda _m^{(\pm )}\right]C_{m+1}^{(\pm )}m+1=0.`$ (35) Introducing the operator $$\widehat{Q}^{}=\left(\widehat{B}_+\widehat{B}_{}\right)^{1/2}\widehat{B}_+$$ (36) one can write the normalized eigenstate of $`\widehat{H}_1`$ as $$m=\left(\widehat{Q}^{}\right)^m0,$$ (37) and, with Eqs. (36) and (37) we can show that $`\widehat{B}_+m=\sqrt{_{m+1}}m+1,`$ (39) $`\widehat{T}\widehat{B}_{}m+1=\sqrt{_{m+1}}\widehat{T}m.`$ (40) Using Eqs. (37), then Eqs. (32) take the form $`\left\{\left[\alpha \beta \lambda _m^{(\pm )}\right]\left(\widehat{T}C_m^{(\pm )}\widehat{T}^{}\right)\pm \alpha \sqrt{_{m+1}}C_{m+1}^{(\pm )}\right\}\widehat{T}m=0`$ (42) $`\left\{\alpha \sqrt{_{m+1}}\left(\widehat{T}C_m^{(\pm )}\widehat{T}^{}\right)\left[\alpha \beta +\lambda _m^{(\pm )}\right]C_{m+1}^{(\pm )}\right\}m+1=0.`$ (43) From Eqs. (II) it follows that $$\lambda _m^{(\pm )}=\pm \alpha \sqrt{_{m+1}+\beta ^2},$$ (44) and $$C_{m+1}^{(\pm )}=\left(\frac{\sqrt{_{m+1}+\beta ^2}\beta }{\sqrt{_{m+1}}}\right)\left(\widehat{T}C_m^{(\pm )}\widehat{T}^{}\right).$$ (45) Eqs. (27) and (45) imply that $$C_{m+1}^{(\pm )}=C_m^{()},$$ (46) and the eigenstates and eigenvalues of the generalized non-resonant Jaynes-Cummings Hamiltonians can be written as $$E_m^{(\pm )}=_{m+1}\pm \sqrt{\mathrm{}\mathrm{\Omega }_{m+1}+\mathrm{}^2\mathrm{\Delta }^2},$$ (47) and $$\mathrm{\Psi }_m^{(\pm )}=\left[\begin{array}{cc}\widehat{T}& 0\\ 0& \pm 1\end{array}\right]\left[\begin{array}{c}C_m^{(\pm )}m\\ C_m^{()}m+1\end{array}\right],m=0,1,2,\mathrm{}$$ (48) a) The Resonant Limit From these general results we can verify two important and simple limiting cases. The first one corresponds to the resonant situation, for which $`\mathrm{\Delta }=0`$ $`(\beta =0)`$. Using these conditions in Eqs. (45) and (47) and Eqs. (27) we get $$E_m^{(\pm )}=_{m+1}\pm \sqrt{\mathrm{}\mathrm{\Omega }_{m+1}},$$ (49) and $$C_{m+1}^{(\pm )}=\widehat{T}C_m^{(\pm )}\widehat{T}^{}=C_m^{(\pm )}=\frac{1}{\sqrt{2}}.$$ (50) Therefore the Jaynes-Cummings resonant eigenstate is given by $$\mathrm{\Psi }_m^{(\pm )}=\frac{1}{\sqrt{2}}\left[\begin{array}{cc}\widehat{T}& 0\\ 0& \pm 1\end{array}\right]\left[\begin{array}{c}m\\ m+1\end{array}\right],m=0,1,2,\mathrm{}$$ (51) These particular results are shown in the Ref. . b) The Standard Jaynes-Cummings Limit The second important limit corresponds to the standard Jaynes-Cummings Hamiltonian, related with the harmonic oscillator system. In this limit we have that $`\widehat{T}=\widehat{T}^{}1`$, $`\widehat{B}_{}\widehat{a}`$, $`\widehat{B}_+\widehat{a}^{}`$, $`\mathrm{\Delta }=\omega \omega _o`$ and $`_{m+1}=(m+1)\mathrm{}\omega `$. Using these conditions in the Eqs. (45), (47) and Eqs. (27) we conclude that $$E_m^{(\pm )}=(m+1)\mathrm{}\omega \pm \mathrm{}\sqrt{\mathrm{\Omega }\omega (m+1)+(\omega \omega _o)^2},$$ (52) and $$C_{m+1}^{(\pm )}=\gamma _m^{(\pm )}C_m^{(\pm )}=C_m^{()}=\frac{1}{\sqrt{1+\left(\gamma _m^{()}\right)^2}},$$ (53) where $`\gamma _m^{(\pm )}=\sqrt{1+\delta _m^2}\delta _m,`$ (55) $`\delta _m={\displaystyle \frac{\omega \omega _o}{\sqrt{(m+1)\mathrm{\Omega }\omega }}}.`$ (56) Therefore the standard Jaynes-Cummings eigenstate, written in a matrix form, is given by $$\mathrm{\Psi }_m^{(\pm )}=\frac{1}{\sqrt{1+\left(\gamma _m^{(\pm )}\right)^2}}\left[\begin{array}{cc}1& 0\\ 0& \pm \gamma _m^{(\pm )}\end{array}\right]\left[\begin{array}{c}m\\ m+1\end{array}\right],m=0,1,2,\mathrm{}$$ (57) These results are shown in many papers, in particular, in the Ref. . ## III The Time Evolution of the System To study the time-dependent Schrödinger equation for a Jaynes-Cummings system in non-resonant situation $$i\mathrm{}\frac{}{t}\mathrm{\Psi }(t)=\left(\widehat{𝐇}_o+\widehat{𝐇}_{int}\right)\mathrm{\Psi }(t)$$ (58) we can write the wavefunction as $$\mathrm{\Psi }(t)=\mathrm{exp}\left(i\widehat{𝐇}_ot/\mathrm{}\right)\mathrm{\Psi }_i(t),$$ (59) and, by substituting this into Schrödinger equation and taking into account the commutation property between $`\widehat{𝐇}_o`$ and $`\widehat{𝐇}_{int}`$, we obtain $$i\mathrm{}\frac{}{t}\mathrm{\Psi }_i(t)=\widehat{𝐇}_{int}\mathrm{\Psi }_i(t).$$ (60) We introduce the evolution matrix $`\widehat{𝐔}_i(t,0)`$: $$\mathrm{\Psi }_i(t)=\widehat{𝐔}_i(t,0)\mathrm{\Psi }_i(0).$$ (61) which satisfies the equation $$i\mathrm{}\frac{}{t}\widehat{𝐔}_i(t,0)=\widehat{𝐇}_{int}\widehat{𝐔}_i(t,0),$$ (62) that is, in matrix form, written as $$i\mathrm{}\left[\begin{array}{cc}\widehat{U}_{11}^{}& \widehat{U}_{12}^{}\\ \widehat{U}_{21}^{}& \widehat{U}_{22}^{}\end{array}\right]=\alpha \left[\begin{array}{cc}\beta & \widehat{T}\widehat{B}_{}\\ \widehat{B}_+\widehat{T}^{}& \beta \end{array}\right]\left[\begin{array}{cc}\widehat{U}_{11}& \widehat{U}_{12}\\ \widehat{U}_{21}& \widehat{U}_{22}\end{array}\right],$$ (63) where the primes denote the time derivative. One fast way to diagonalize the evolution matrix differential equation is by differentiating Eq. (62) with respect to time. We find $$i\mathrm{}\frac{^2}{t^2}\widehat{𝐔}_i(t,0)=\widehat{𝐇}_{int}\frac{}{t}\widehat{𝐔}_i(t,0)=\frac{1}{i\mathrm{}}\widehat{𝐇}_{int}^2\widehat{𝐔}_i(t,0),$$ (64) which can be written as $$\left[\begin{array}{cc}\widehat{U}_{11}^{\prime \prime }& \widehat{U}_{12}^{\prime \prime }\\ \widehat{U}_{21}^{\prime \prime }& \widehat{U}_{22}^{\prime \prime }\end{array}\right]=\left[\begin{array}{cc}\widehat{\omega }_1& 0\\ 0& \widehat{\omega }_2\end{array}\right]\left[\begin{array}{cc}\widehat{U}_{11}& \widehat{U}_{12}\\ \widehat{U}_{21}& \widehat{U}_{22}\end{array}\right],$$ (65) where $`\mathrm{}\widehat{\omega }_1=\alpha \sqrt{\widehat{T}\widehat{B}_{}\widehat{B}_+\widehat{T}^{}+\beta ^2}=\sqrt{\mathrm{}\mathrm{\Omega }\widehat{H}_2+(\mathrm{}\mathrm{\Delta })^2},`$ (67) $`\mathrm{}\widehat{\omega }_2=\alpha \sqrt{\widehat{B}_+\widehat{B}_{}+\beta ^2}=\sqrt{\mathrm{}\mathrm{\Omega }\widehat{H}_1+(\mathrm{}\mathrm{\Delta })^2}.`$ (68) Since by initial conditions $`\widehat{𝐔}_i(0,0)=\widehat{𝐈}`$, then we can write the solution of the evolution matrix differential equation (64) as $$\widehat{𝐔}_i(t,0)=\left[\begin{array}{cc}\mathrm{cos}(\widehat{\omega }_1t)& \mathrm{sin}(\widehat{\omega }_1t)\widehat{C}\\ \mathrm{sin}(\widehat{\omega }_2t)\widehat{D}& \mathrm{cos}(\widehat{\omega }_2t)\end{array}\right],$$ (69) and the $`\widehat{C}`$ and $`\widehat{D}`$ operators can be determined by the unitarity conditions $$\widehat{𝐔}_i^{}(t,0)\widehat{𝐔}_i(t,0)=\widehat{𝐔}_i(t,0)\widehat{𝐔}_i^{}(t,0)=\widehat{𝐈}.$$ (70) In the appendix A we show that the unitarity conditions (70) imply $`\widehat{C}=\widehat{D}^{}={\displaystyle \frac{i}{(\widehat{H}_2)^{1/4}}}\sqrt{\widehat{T}\widehat{B}_{}}`$ (72) $`\widehat{D}=\widehat{C}^{}.`$ (73) Therefore, we can write the final expression of the time evolution matrix of the system as $$\widehat{𝐔}_i(t,0)=\left[\begin{array}{cc}\mathrm{cos}(\widehat{\omega }_1t)& \mathrm{sin}(\widehat{\omega }_1t)\widehat{C}\\ \mathrm{sin}(\widehat{\omega }_2t)\widehat{C}^{}& \mathrm{cos}(\widehat{\omega }_2t)\end{array}\right].$$ (74) For Jaynes-Cummings systems an important physical quantity to see how the system under consideration evolves in time is the population inversion factor , defined by $$\widehat{𝐖}(t)\widehat{\sigma }_+(t)\widehat{\sigma }_{}(t)\widehat{\sigma }_{}(t)\widehat{\sigma }_+(t)=\widehat{\sigma }_3(t),$$ (75) where the time dependence of the operators is related with the Heisenberg picture. In this case, the time evolution of the population inversion factor will be given by $$\frac{d\widehat{\sigma }_3(t)}{dt}=\frac{1}{i\mathrm{}}\widehat{𝐔}_i^{}(t,0)[\widehat{\sigma }_3,\widehat{𝐇}]\widehat{𝐔}_i(t,0),$$ (76) and since we have $$[\widehat{\sigma }_3,\widehat{𝐇}]=\alpha [\widehat{\sigma }_3,\widehat{𝐒}]=2\alpha \widehat{𝐒}\widehat{\sigma }_3,$$ (77) then Eq. (76) can be written as $$\frac{d\widehat{\sigma }_3(t)}{dt}=\frac{2i\alpha }{\mathrm{}}\widehat{𝐒}(t)\widehat{\sigma }_3(t).$$ (78) We can obtain a differential equation with constant coefficients for $`\widehat{\sigma }_3(t)`$ by taking the time derivative of Eq. (78) $$\frac{d^2\widehat{\sigma }_3(t)}{dt^2}=\frac{2i\alpha }{\mathrm{}}\left\{\frac{d\widehat{𝐒}(t)}{dt}\widehat{\sigma }_3(t)+\widehat{𝐒}(t)\frac{d\widehat{\sigma }_3(t)}{dt}\right\}.$$ (79) Having in mind that $$\frac{d\widehat{𝐒}(t)}{dt}=\frac{1}{i\mathrm{}}\widehat{𝐔}_i^{}(t,0)[\widehat{𝐒},\widehat{𝐇}]\widehat{𝐔}_i(t,0),$$ (80) and that $$[\widehat{𝐒},\widehat{𝐇}]=\alpha \beta [\widehat{𝐒},\widehat{\sigma }_3]=2\alpha \beta \widehat{𝐒}\widehat{\sigma }_3,$$ (81) we conclude that $$\frac{d\widehat{𝐒}(t)}{dt}=\frac{2i\alpha \beta }{\mathrm{}}\widehat{𝐒}(t)\widehat{\sigma }_3(t).$$ (82) Using Eqs. (78) and (82) into Eq. (79) we obtain $$\frac{d^2\widehat{\sigma }_3(t)}{dt^2}+\widehat{𝚯}^2\widehat{\sigma }_3(t)=\widehat{𝐅}(t)$$ (83) where $`\widehat{𝚯}^2={\displaystyle \frac{4\alpha ^2}{\mathrm{}^2}}\widehat{𝐒}^2`$ (85) $`\widehat{𝐅}(t)={\displaystyle \frac{4\alpha ^2\beta }{\mathrm{}^2}}\widehat{𝐔}_i^{}(t,0)\widehat{𝐒}\widehat{𝐔}_i(t,0).`$ (86) Eq. (83) corresponds to a non-homogeneous linear differential equation for $`\widehat{\sigma }_3(t)`$ with constant coefficients since $`\widehat{𝐒}^2`$ and $`\widehat{𝐇}`$ commute and, therefore, $`\widehat{𝚯}`$ is a constant of the motion. The general solution of this differential equation can be written as $$\widehat{\sigma }_3(t)=\widehat{\sigma }^H(t)+\widehat{\sigma }^P(t),$$ (87) and each matrix element of the homogeneous solution, satisfies the differential equation $$\frac{d^2\widehat{\sigma }_{jk}^H(t)}{dt^2}+\widehat{\nu }_j^2\widehat{\sigma }_{jk}^H(t)=0,j,k=1,\mathrm{or}\mathrm{\hspace{0.33em}2},$$ (88) with $`\mathrm{}\widehat{\nu }_1=2\alpha \sqrt{\widehat{T}\widehat{B}_{}\widehat{B}_+\widehat{T}^{}}=2\sqrt{\mathrm{}\mathrm{\Omega }\widehat{H}_2},`$ (90) $`\mathrm{}\widehat{\nu }_2=2\alpha \sqrt{\widehat{B}_+\widehat{B}_{}}=2\sqrt{\mathrm{}\mathrm{\Omega }\widehat{H}_1}.`$ (91) The solution of Eq. (88) is given by $$\widehat{\sigma }_{jk}^H(t)=\widehat{y}_j(t)\widehat{c}_{jk}+\widehat{z}_j(t)\widehat{d}_{jk},$$ (92) where $`\widehat{y}_j(t)=\mathrm{cos}(\widehat{\nu }_jt)`$ (94) $`\widehat{z}_j(t)=\mathrm{sin}(\widehat{\nu }_jt),`$ (95) and the coefficients $`\widehat{c}_{jk}`$ and $`\widehat{d}_{jk}`$ can be determined by the initial conditions. The matrix elements of the particular solution of the $`\widehat{\sigma }_3(t)`$ differential equation need to satisfy $$\frac{d^2\widehat{\sigma }_{jk}^P(t)}{dt^2}+\widehat{\nu }_j^2\widehat{\sigma }_{jk}^P(t)=\widehat{F}_{jk}(t),j,k=1,\mathrm{or}\mathrm{\hspace{0.33em}2},$$ (96) and can be obtained by the variation of parameter or by Green function methods, giving $$\widehat{\sigma }_{jk}^P(t)=\widehat{\nu }_j^1\left\{\widehat{z}_j(t)_0^t𝑑\xi \widehat{y}_j(\xi )\widehat{F}_{jk}(\xi )\widehat{y}_j(t)_0^t𝑑\xi \widehat{z}_j(\xi )\widehat{F}_{jk}(\xi )\right\},$$ (97) where we used that the Wronskian of the system of solutions $`\widehat{y}_j(t)`$ and $`\widehat{z}_j(t)`$ is given by $`\widehat{\nu }_j`$. After we determine the elements of the $`\widehat{𝐅}(t)`$-matrix, it is necessary to resolve the integrals in Eq. (97) to obtain the explicit expression of the particular solution. In the appendix B we show that, using Eqs. (15), (74), and (83), it is possible to conclude that these matrix elements can be written as $`\widehat{\sigma }_{11}^P(t)`$ $`=`$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_1^1\sqrt{\widehat{T}\widehat{B}_{}}\left\{\widehat{z}_2(t)𝒢_{CS}^{(+)}(t;\widehat{\nu }_2,\widehat{\omega }_2,\widehat{\omega }_1)\widehat{y}_2(t)𝒢_{SS}^{(+)}(t;\widehat{\nu }_2,\widehat{\omega }_2,\widehat{\omega }_1)\right\}\widehat{H}_2^{1/4}`$ (99) $`+`$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_1^1\widehat{H}_2^{1/4}\left\{\widehat{z}_1(t)𝒢_{SC}^{()}(t;\widehat{\nu }_1,\widehat{\omega }_1,\widehat{\omega }_2)\widehat{y}_1(t)𝒢_{CC}^{()}(t;\widehat{\nu }_1,\widehat{\omega }_1,\widehat{\omega }_2)\right\}\sqrt{\widehat{B}_+\widehat{T}^{}},`$ (100) $`\widehat{\sigma }_{12}^P(t)`$ $`=`$ $`{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_1^1\sqrt{\widehat{T}\widehat{B}_{}}\left\{\widehat{z}_2(t)𝒢_{CC}^{(+)}(t;\widehat{\nu }_2,\widehat{\omega }_2,\widehat{\omega }_1)\widehat{y}_2(t)𝒢_{SC}^{(+)}(t;\widehat{\nu }_2,\widehat{\omega }_2,\widehat{\omega }_1)\right\}\sqrt{\widehat{T}\widehat{B}_{}}`$ (101) $`+`$ $`{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_1^1\widehat{H}_2^{1/4}\left\{\widehat{z}_1(t)𝒢_{SS}^{()}(t;\widehat{\nu }_1,\widehat{\omega }_1,\widehat{\omega }_2)+\widehat{y}_1(t)𝒢_{CS}^{()}(t;\widehat{\nu }_1,\widehat{\omega }_1,\widehat{\omega }_2)\right\}\widehat{H}_1^{1/4},`$ (102) $`\widehat{\sigma }_{21}^P(t)`$ $`=`$ $`{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_2^1\sqrt{\widehat{B}_+\widehat{T}^{}}\left\{\widehat{z}_1(t)𝒢_{CC}^{(+)}(t;\widehat{\nu }_1,\widehat{\omega }_1,\widehat{\omega }_2)\widehat{y}_1(t)𝒢_{SC}^{(+)}(t;\widehat{\nu }_1,\widehat{\omega }_1,\widehat{\omega }_2)\right\}\sqrt{\widehat{B}_+\widehat{T}^{}}`$ (103) $`+`$ $`{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_2^1\widehat{H}_1^{1/4}\left\{\widehat{z}_2(t)𝒢_{SS}^{()}(t;\widehat{\nu }_2,\widehat{\omega }_2,\widehat{\omega }_1)\widehat{y}_2(t)𝒢_{CS}^{()}(t;\widehat{\nu }_2,\widehat{\omega }_2,\widehat{\omega }_1)\right\}\widehat{H}_2^{1/4},`$ (104) $`\widehat{\sigma }_{22}^P(t)`$ $`=`$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_2^1\sqrt{\widehat{B}_+\widehat{T}^{}}\left\{\widehat{z}_1(t)𝒢_{CS}^{(+)}(t;\widehat{\nu }_1,\widehat{\omega }_1,\widehat{\omega }_2)\widehat{y}_1(t)𝒢_{SS}^{(+)}(t;\widehat{\nu }_1,\widehat{\omega }_1,\widehat{\omega }_2)\right\}\widehat{H}_1^{1/4}`$ (105) $`+`$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_2^1\widehat{H}_1^{1/4}\left\{\widehat{z}_2(t)𝒢_{SC}^{()}(t;\widehat{\nu }_2,\widehat{\omega }_2,\widehat{\omega }_1)+\widehat{y}_2(t)𝒢_{CC}^{()}(t;\widehat{\nu }_2,\widehat{\omega }_2,\widehat{\omega }_1)\right\}\sqrt{\widehat{T}\widehat{B}_{}},`$ (106) where $`\gamma =4\alpha ^2\beta /\mathrm{}^2`$, and the auxiliary functions are given by $$𝒢_{XY}^{(\pm )}(t;\widehat{p},\widehat{q},\widehat{r})=_{XY}(t;\widehat{p}\widehat{q},\widehat{r})\pm _{XY}(t;\widehat{p}+\widehat{q},\widehat{r}),X,Y=C\mathrm{or}S,$$ (107) with $`_{CC}(t;\widehat{x},\widehat{w})`$ $``$ $`{\displaystyle _0^t}𝑑\xi \mathrm{cos}(\widehat{x}\xi )\mathrm{cos}(\widehat{w}\xi )`$ (109) $`=`$ $`{\displaystyle \underset{m,n=0}{\overset{\mathrm{}}{}}}(1)^{m+n}{\displaystyle \frac{\widehat{x}^{2m}\widehat{w}^{2n}}{(2m)!(2n)!}}{\displaystyle \frac{t^{2m+2n+1}}{(2m+2n+1)}}`$ (110) $`_{CS}(t;\widehat{x},\widehat{w})`$ $``$ $`{\displaystyle _0^t}𝑑\xi \mathrm{cos}(\widehat{x}\xi )\mathrm{sin}(\widehat{w}\xi )`$ (111) $`=`$ $`{\displaystyle \underset{m,n=0}{\overset{\mathrm{}}{}}}(1)^{m+n}{\displaystyle \frac{\widehat{x}^{2m}\widehat{w}^{2n+1}}{(2m)!(2n+1)!}}{\displaystyle \frac{t^{2m+2n+2}}{(2m+2n+2)}}`$ (112) $`_{SC}(t;\widehat{x},\widehat{w})`$ $``$ $`{\displaystyle _0^t}𝑑\xi \mathrm{sin}(\widehat{x}\xi )\mathrm{cos}(\widehat{w}\xi )`$ (113) $`=`$ $`{\displaystyle \underset{m,n=0}{\overset{\mathrm{}}{}}}(1)^{m+n}{\displaystyle \frac{\widehat{x}^{2m+1}\widehat{w}^{2n}}{(2m+1)!(2n)!}}{\displaystyle \frac{t^{2m+2n+2}}{(2m+2n+2)}}`$ (114) $`_{SS}(t;\widehat{x},\widehat{w})`$ $``$ $`{\displaystyle _0^t}𝑑\xi \mathrm{sin}(\widehat{x}\xi )\mathrm{sin}(\widehat{w}\xi )`$ (115) $`=`$ $`{\displaystyle \underset{m,n=0}{\overset{\mathrm{}}{}}}(1)^{m+n}{\displaystyle \frac{\widehat{x}^{2m+1}\widehat{w}^{2n+1}}{(2m+1)!(2n+1)!}}{\displaystyle \frac{t^{2m+2n+3}}{(2m+2n+3)}}.`$ (116) With these results for the particular solution we can conclude that $$\widehat{\sigma }_{ij}^P(0)=0=\frac{d\widehat{\sigma }_{ij}^P(0)}{dt}.$$ (117) Now, using Eqs. (78), (87), (92), (117) and the initial conditions, we have $`\left[\widehat{\sigma }_3(0)\right]_{ij}=\widehat{c}_{ij}`$ (119) $`\left[{\displaystyle \frac{d\widehat{\sigma }_3(0)}{dt}}\right]_{ij}={\displaystyle \frac{2i\alpha }{\mathrm{}}}\left[\widehat{𝐒}(0)\widehat{\sigma }_3(0)\right]_{ij}=\widehat{\nu }_i\widehat{d}_{ij}.`$ (120) Therefore, the final expression for the elements of the population inversion matrix of the system can be written as $$[\widehat{\sigma }_3(t)]_{ij}=\mathrm{cos}(\widehat{\nu }_it)\left[\widehat{\sigma }_3(0)\right]_{ij}+\frac{2i\alpha }{\mathrm{}}\mathrm{sin}(\widehat{\nu }_it)\widehat{\nu }_i^1\left[\widehat{𝐒}(0)\widehat{\sigma }_3(0)\right]_{ij}+\widehat{\sigma }_{ij}^P(t).$$ (121) Again, using these final results we can verify two important and simple limiting cases. a) The Resonant Limit The first one corresponds to the resonant situation $`(\mathrm{\Delta }=0)`$. Eqs. (65), (74), (88) and (III) allow us to conclude that, in this case, the evolution matrix of the system is given by $$\widehat{𝐔}_i(t,0)=\left[\begin{array}{cc}\mathrm{cos}\left(\frac{1}{2}\widehat{\nu }_1t\right)& \mathrm{sin}\left(\frac{1}{2}\widehat{\nu }_1t\right)\widehat{C}\\ \mathrm{sin}\left(\frac{1}{2}\widehat{\nu }_2t\right)\widehat{C}^{}& \mathrm{cos}\left(\frac{1}{2}\widehat{\nu }_2t\right)\end{array}\right].$$ (122) and the elements of the population inversion of the system are $$[\widehat{\sigma }_3(t)]_{ij}=\mathrm{cos}(\widehat{\nu }_it)\left[\widehat{\sigma }_3(0)\right]_{ij}+\frac{2i\alpha }{\mathrm{}}\mathrm{sin}(\widehat{\nu }_it)\widehat{\nu }_i^1\left[\widehat{𝐒}(0)\widehat{\sigma }_3(0)\right]_{ij}.$$ (123) b) The Standard Jaynes-Cummings Limit This second important limit corresponds to the case of the harmonic oscillator system, and in this limit we have that $`\widehat{T}=\widehat{T}^{}1`$, $`\widehat{B}_{}\widehat{a}`$, $`\widehat{B}_+\widehat{a}^{}`$ and $`[\widehat{a},\widehat{a}^{}]=\mathrm{}\omega `$. With these conditions the operators $`\widehat{\omega }_1`$ and $`\widehat{\omega }_2`$ commute, and this fact permits to evaluate the integrals related with the particular solution of the population inversion elements using trigonometric product relations. Using that and the expressions obtained in the appendix B, after a considerable amount of algebra and trigonometric product relations we can show that is possible to write the expressions for the $`\widehat{\sigma }_{ij}^P(t)`$-matrix elements as $`\widehat{\sigma }_{11}^P(t)`$ $`=`$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_1^1\sqrt{\widehat{a}}\left\{𝒦_S(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_2)𝒦_S(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_2)\right\}\left(\widehat{a}\widehat{a}^{}\right)^{1/4}`$ (125) $``$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_1^1\left(\widehat{a}\widehat{a}^{}\right)^{1/4}\left\{𝒦_S(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_1)𝒦_S(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_1)\right\}\sqrt{\widehat{a}^{}}`$ (126) $`\widehat{\sigma }_{12}^P(t)`$ $`=`$ $`{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_1^1\sqrt{\widehat{a}}\left\{𝒦_C(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_2)𝒦_C(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_2)\right\}\sqrt{\widehat{a}}`$ (127) $``$ $`{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_1^1\left(\widehat{a}\widehat{a}^{}\right)^{1/4}\left\{𝒦_C(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_1)𝒦_C(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_1)\right\}\left(\widehat{a}^{}\widehat{a}\right)^{1/4}`$ (128) $`\widehat{\sigma }_{21}^P(t)`$ $`=`$ $`{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_2^1\sqrt{\widehat{a}^{}}\left\{𝒦_C(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_1)+𝒦_C(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_1)\right\}\sqrt{\widehat{a}^{}}`$ (129) $``$ $`{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_2^1\left(\widehat{a}^{}\widehat{a}\right)^{1/4}\left\{𝒦_C(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_2)𝒦_C(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_2)\right\}\left(\widehat{a}\widehat{a}^{}\right)^{1/4}`$ (130) $`\widehat{\sigma }_{22}^P(t)`$ $`=`$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_2^1\sqrt{\widehat{a}^{}}\left\{𝒦_S(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_1)+𝒦_S(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_1)\right\}\left(\widehat{a}^{}\widehat{a}\right)^{1/4}`$ (131) $``$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{\nu }_2^1\left(\widehat{a}^{}\widehat{a}\right)^{1/4}\left\{𝒦_S(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_2)+𝒦_S(t;\widehat{\omega }_2,\widehat{\omega }_1,\widehat{\nu }_2)\right\}\sqrt{\widehat{a}},`$ (132) where, now, the auxiliary functions are given by $`𝒦_S(t;\widehat{p},\widehat{q},\widehat{r})`$ $`=`$ $`{\displaystyle \frac{\widehat{r}\mathrm{sin}\left[\left(\widehat{p}+\widehat{q}\right)t\right]\left(\widehat{p}+\widehat{q}\right)\mathrm{sin}\left(\widehat{r}t\right)}{\widehat{r}^2\left(\widehat{p}+\widehat{q}\right)^2}}`$ (134) $`𝒦_C(t;\widehat{p},\widehat{q},\widehat{r})`$ $`=`$ $`{\displaystyle \frac{\widehat{r}\mathrm{cos}\left[\left(\widehat{p}+\widehat{q}\right)t\right]\widehat{r}\mathrm{cos}\left(\widehat{r}t\right)}{\widehat{r}^2\left(\widehat{p}+\widehat{q}\right)^2}}.`$ (135) Considering the expressions above we may easily verify that the particular solution for the population inversion factor must still satisfy the initial conditions (117). Therefore, in this case the final expression for the population inversion factor has the same form given by Eq. (121), with $`\mathrm{}\widehat{\nu }_1=2\alpha \sqrt{\widehat{a}\widehat{a}^{}},\mathrm{}\widehat{\nu }_2=2\alpha \sqrt{\widehat{a}^{}\widehat{a}},`$ (137) $`\mathrm{}\widehat{\omega }_1=\alpha \sqrt{\widehat{a}\widehat{a}^{}+\beta ^2},\mathrm{}\widehat{\omega }_2=\alpha \sqrt{\widehat{a}^{}\widehat{a}+\beta ^2}.`$ (138) ## IV Conclusions In this article we extended our earlier work on bound-state problems which represent two-level systems. The corresponding coupled-channel Hamiltonians generalize the Jaynes-Cummings non-resonant Hamiltonian. If we take the starting Hamiltonian to be the simplest shape-invariant system, namely the harmonic oscillator, our results reduce to those of the standard non-resonant Jaynes-Cummings approach, which has been extensively used to model a two-level atom-single field mode interaction whose detuning it is not null. Another possible extension of our model is to consider intensity-dependent couplings. This will be taken up in the following paper . ## ACKNOWLEDGMENTS This work was supported in part by the U.S. National Science Foundation Grants No. PHY-9605140 and PHY-0070161 at the University of Wisconsin, and in part by the University of Wisconsin Research Committee with funds granted by the Wisconsin Alumni Research Foundation. A.B.B. acknowledges the support of the Alexander von Humboldt-Stiftung. M.A.C.R. acknowledges the support of Fundação de Amparo à Pesquisa do Estado de São Paulo (Contract No. 98/13722-2). A.N.F.A. acknowledges the support of Fundação Coordenação de Aperfeiçoamento de Pessoal de Nível Superior (Contract No. BEX0610/96-8). A.B.B. is grateful to the Max-Planck-Institut für Kernphysik and H.A. Weidenmüller for the very kind hospitality. Appendix A Here we give the steps used to obtain the specific form of the operators $`\widehat{C}`$ and $`\widehat{D}`$. Using Eq. (69) into the unitary condition equation (70) actually we can show that the $`\widehat{C}`$ and $`\widehat{D}`$ operators need to satisfy the following six conditions $`\widehat{C}\widehat{C}^{}=\widehat{C}^{}\widehat{C}=1`$ (140) $`\widehat{D}\widehat{D}^{}=\widehat{D}^{}\widehat{D}=1`$ (141) $`\widehat{D}^{}\mathrm{sin}(\widehat{\omega }_2t)=\mathrm{sin}(\widehat{\omega }_1t)\widehat{C}`$ (142) $`\widehat{D}\mathrm{cos}(\widehat{\omega }_1t)=\mathrm{cos}(\widehat{\omega }_2t)\widehat{C}^{}.`$ (143) At this point we can use the following property $`\sqrt{\widehat{T}\widehat{B}_{}}\widehat{\omega }_2`$ $`=`$ $`\sqrt{\widehat{T}\widehat{B}_{}}\sqrt{\alpha ^2\widehat{B}_+\widehat{B}_{}+\beta ^2}/\mathrm{}`$ (144) $`=`$ $`\sqrt{\alpha ^2\widehat{T}\widehat{B}_{}\widehat{B}_+\widehat{B}_{}+\widehat{T}\widehat{B}_{}\beta ^2}/\mathrm{}`$ (145) $`=`$ $`\sqrt{\alpha ^2\widehat{T}\widehat{B}_{}\widehat{B}_+\widehat{T}^{}\widehat{T}\widehat{B}_{}+\beta ^2\widehat{T}\widehat{B}_{}}/\mathrm{}`$ (146) $`=`$ $`\sqrt{\alpha ^2\widehat{T}\widehat{B}_{}\widehat{B}_+\widehat{T}^{}+\beta ^2}/\mathrm{}\sqrt{\widehat{T}\widehat{B}_{}}`$ (147) $`=`$ $`\widehat{\omega }_1\sqrt{\widehat{T}\widehat{B}_{}}.`$ (148) Then, with this result we have $$\sqrt{\widehat{T}\widehat{B}_{}}\widehat{\omega }_2^2=\sqrt{\widehat{T}\widehat{B}_{}}\widehat{\omega }_2\widehat{\omega }_2=\widehat{\omega }_1\sqrt{\widehat{T}\widehat{B}_{}}\widehat{\omega }_2=\widehat{\omega }_1^2\sqrt{\widehat{T}\widehat{B}_{}},$$ (149) and finally, by induction, we conclude that $$\sqrt{\widehat{T}\widehat{B}_{}}\widehat{\omega }_2^n=\widehat{\omega }_1^n\sqrt{\widehat{T}\widehat{B}_{}}.$$ (150) In the same way, $`\sqrt{\widehat{B}_+\widehat{T}^{}}\widehat{\omega }_1`$ $`=`$ $`\sqrt{\widehat{B}_+\widehat{T}^{}}\sqrt{\alpha ^2\widehat{T}\widehat{B}_{}\widehat{B}_+\widehat{T}^{}+\beta ^2}/\mathrm{}`$ (151) $`=`$ $`\sqrt{\alpha ^2\widehat{B}_+\widehat{T}^{}\widehat{T}\widehat{B}_{}\widehat{B}_+\widehat{T}^{}+\widehat{B}_+\widehat{T}^{}\beta ^2}/\mathrm{}`$ (152) $`=`$ $`\sqrt{\alpha ^2\widehat{B}_+\widehat{B}_{}\widehat{B}_+\widehat{T}^{}+\beta ^2\widehat{B}_+\widehat{T}^{}}/\mathrm{}`$ (153) $`=`$ $`\sqrt{\alpha ^2\widehat{B}_+\widehat{B}_{}+\beta ^2}/\mathrm{}\sqrt{\widehat{B}_+\widehat{T}^{}}`$ (154) $`=`$ $`\widehat{\omega }_2\sqrt{\widehat{B}_+\widehat{T}^{}}.`$ (155) Then, with this result we have $$\sqrt{\widehat{B}_+\widehat{T}^{}}\widehat{\omega }_1^2=\sqrt{\widehat{B}_+\widehat{T}^{}}\widehat{\omega }_1\widehat{\omega }_1=\widehat{\omega }_2\sqrt{\widehat{B}_+\widehat{T}^{}}\widehat{\omega }_1=\widehat{\omega }_2^2\sqrt{\widehat{B}_+\widehat{T}^{}},$$ (156) and finally, again by induction, we get $$\sqrt{\widehat{B}_+\widehat{T}^{}}\widehat{\omega }_1^n=\widehat{\omega }_2^n\sqrt{\widehat{B}_+\widehat{T}^{}}.$$ (157) Using the properties given by Eqs. (150) and (157) and the forms of $`\widehat{C}`$, $`\widehat{D}`$ operators, defined by Eqs. (70), we can verify that $$\widehat{C}\widehat{C}^{}=\widehat{D}^{}\widehat{D}=\frac{i}{\widehat{H}_2^{1/4}}\sqrt{\widehat{T}\widehat{B}_{}}\sqrt{\widehat{B}_+\widehat{T}^{}}\frac{(i)}{\widehat{H}_2^{1/4}}=\frac{1}{\widehat{H}_2^{1/4}}\sqrt{\widehat{H}_2}\frac{1}{\widehat{H}_2^{1/4}}=1,$$ (158) and $$\widehat{C}^{}\widehat{C}=\widehat{D}\widehat{D}^{}=\sqrt{\widehat{B}_+\widehat{T}^{}}\frac{(i)}{\widehat{H}_2^{1/4}}\frac{i}{\widehat{H}_2^{1/4}}\sqrt{\widehat{T}\widehat{B}_{}}=\sqrt{\widehat{B}_+\widehat{T}^{}}\frac{1}{\sqrt{\widehat{H}_2}}\sqrt{\widehat{H}_2}\frac{1}{\sqrt{\widehat{B}_+\widehat{T}^{}}}=1.$$ (159) Also using the series expansion of the trigonometric functions, we can show that $`\widehat{D}^{}\mathrm{sin}(\widehat{\omega }_2t)`$ $`=`$ $`{\displaystyle \frac{i}{\widehat{H}_2^{1/4}}}\sqrt{\widehat{T}\widehat{B}_{}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}(1)^n{\displaystyle \frac{\left(\widehat{\omega }_2t\right)^{2n+1}}{(2n+1)!}}`$ (160) $`=`$ $`{\displaystyle \frac{i}{\widehat{H}_2^{1/4}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}(1)^n\sqrt{\widehat{T}\widehat{B}_{}}{\displaystyle \frac{\left(\widehat{\omega }_2t\right)^{2n+1}}{(2n+1)!}}`$ (161) $`=`$ $`{\displaystyle \frac{i}{\widehat{H}_2^{1/4}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}(1)^n{\displaystyle \frac{\left(\widehat{\omega }_1t\right)^{2n+1}}{(2n+1)!}}\sqrt{\widehat{T}\widehat{B}_{}}`$ (162) $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}(1)^n{\displaystyle \frac{\left(\widehat{\omega }_1t\right)^{2n+1}}{(2n+1)!}}{\displaystyle \frac{i}{\widehat{H}_2^{1/4}}}\sqrt{\widehat{T}\widehat{B}_{}}`$ (163) $`=`$ $`\mathrm{sin}(\widehat{\omega }_1t)\widehat{C},`$ (164) where we used the commutation between $`\widehat{H}_2`$ and $`\widehat{\omega }_1`$ (see Appendix B). In the same way we can prove that $`\widehat{D}\mathrm{cos}(\widehat{\omega }_1t)`$ $`=`$ $`\sqrt{\widehat{B}_+\widehat{T}^{}}{\displaystyle \frac{i}{\widehat{H}_2^{1/4}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}(1)^n{\displaystyle \frac{\left(\widehat{\omega }_1t\right)^{2n}}{(2n)!}}`$ (165) $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}(1)^n\sqrt{\widehat{B}_+\widehat{T}^{}}{\displaystyle \frac{\left(\widehat{\omega }_1t\right)^{2n}}{(2n)!}}{\displaystyle \frac{i}{\widehat{H}_2^{1/4}}}`$ (166) $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}(1)^n{\displaystyle \frac{\left(\widehat{\omega }_2t\right)^{2n}}{(2n)!}}\sqrt{\widehat{B}_+\widehat{T}^{}}{\displaystyle \frac{i}{\widehat{H}_2^{1/4}}}`$ (167) $`=`$ $`\mathrm{cos}(\widehat{\omega }_2t)\widehat{C}^{}.`$ (168) Again, we used the commutation between $`\widehat{H}_2`$ and $`\widehat{\omega }_1`$. Appendix B In this appendix we show the necessary steps to obtain the explicit expressions of the particular solution elements of the population inversion factor. To resolve the integrals in Eq. (97), first we need to determine the elements of the $`\widehat{𝐅}(t)`$-matrix. To do that we can use Eqs. (15), (83), and (74) to write down $`\widehat{F}_{11}(t)`$ $`=`$ $`\gamma \left\{\mathrm{cos}(\widehat{\omega }_1t)\widehat{T}\widehat{B}_{}\mathrm{sin}(\widehat{\omega }_2t)\widehat{C}^{}+\widehat{C}\mathrm{sin}(\widehat{\omega }_2t)\widehat{B}_+\widehat{T}^{}\mathrm{cos}(\widehat{\omega }_1t)\right\}`$ (170) $`=`$ $`i\gamma \left\{\sqrt{\widehat{T}\widehat{B}_{}}\mathrm{cos}(\widehat{\omega }_2t)\mathrm{sin}(\widehat{\omega }_1t)\widehat{H}_2^{1/4}\widehat{H}_2^{1/4}\mathrm{sin}(\widehat{\omega }_1t)\mathrm{cos}(\widehat{\omega }_2t)\sqrt{\widehat{B}_+\widehat{T}^{}}\right\}`$ (171) $`\widehat{F}_{12}(t)`$ $`=`$ $`\gamma \left\{\mathrm{cos}(\widehat{\omega }_1t)\widehat{T}\widehat{B}_{}\mathrm{cos}(\widehat{\omega }_2t)\widehat{C}\mathrm{sin}(\widehat{\omega }_2t)\widehat{B}_+\widehat{T}^{}\mathrm{sin}(\widehat{\omega }_1t)\widehat{C}\right\}`$ (172) $`=`$ $`\gamma \left\{\sqrt{\widehat{T}\widehat{B}_{}}\mathrm{cos}(\widehat{\omega }_2t)\mathrm{cos}(\widehat{\omega }_1t)\sqrt{\widehat{T}\widehat{B}_{}}+\widehat{H}_2^{1/4}\mathrm{sin}(\widehat{\omega }_1t)\mathrm{sin}(\widehat{\omega }_2t)\widehat{H}_1^{1/4}\right\}`$ (173) $`\widehat{F}_{21}(t)`$ $`=`$ $`\gamma \left\{\mathrm{cos}(\widehat{\omega }_2t)\widehat{B}_+\widehat{T}^{}\mathrm{cos}(\widehat{\omega }_1t)\widehat{C}^{}\mathrm{sin}(\widehat{\omega }_1t)\widehat{T}\widehat{B}_{}\mathrm{sin}(\widehat{\omega }_2t)\widehat{C}^{}\right\}`$ (174) $`=`$ $`\gamma \left\{\sqrt{\widehat{B}_+\widehat{T}^{}}\mathrm{cos}(\widehat{\omega }_1t)\mathrm{cos}(\widehat{\omega }_2t)\sqrt{\widehat{B}_+\widehat{T}^{}}+\widehat{H}_1^{1/4}\mathrm{sin}(\widehat{\omega }_2t)\mathrm{sin}(\widehat{\omega }_1t)\widehat{H}_2^{1/4}\right\}`$ (175) $`\widehat{F}_{22}(t)`$ $`=`$ $`\gamma \left\{\widehat{C}^{}\mathrm{sin}(\widehat{\omega }_1t)\widehat{T}\widehat{B}_{}\mathrm{cos}(\widehat{\omega }_2t)+\mathrm{cos}(\widehat{\omega }_2t)\widehat{B}_+\widehat{T}^{}\mathrm{sin}(\widehat{\omega }_1t)\widehat{C}\right\}`$ (176) $`=`$ $`i\gamma \left\{\sqrt{\widehat{B}_+\widehat{T}^{}}\mathrm{cos}(\widehat{\omega }_1t)\mathrm{sin}(\widehat{\omega }_2t)\widehat{H}_1^{1/4}\widehat{H}_1^{1/4}\mathrm{sin}(\widehat{\omega }_2t)\mathrm{cos}(\widehat{\omega }_1t)\sqrt{\widehat{T}\widehat{B}_{}}\right\},`$ (177) where $`\gamma =4\alpha ^2\beta /\mathrm{}^2`$. Here we used the properties (ACKNOWLEDGMENTS), (150) and (157), together with the following operators relations $`\widehat{C}\sqrt{\widehat{B}_+\widehat{T}^{}}=\sqrt{\widehat{T}\widehat{B}_{}}\widehat{C}^{}=i\widehat{H}_2^{1/4}`$ (179) $`\sqrt{\widehat{B}_+\widehat{T}^{}}\widehat{C}=\widehat{C}^{}\sqrt{\widehat{T}\widehat{B}_{}}=i\widehat{H}_1^{1/4}.`$ (180) Now, keeping in mind that $`[\widehat{\nu }_j,\widehat{\omega }_j]=0`$, $`(j=1,\mathrm{or}\mathrm{\hspace{0.33em}2})`$, so we may use the trigonometric relationships involving products of trigonometric functions with arguments $`\widehat{\nu }_jt`$ and $`\widehat{\omega }_jt`$ (since we have $`\mathrm{exp}(\widehat{\nu }_jt)\mathrm{exp}(\pm \widehat{\omega }_jt)=\mathrm{exp}[(\widehat{\nu }_j\pm \widehat{\omega }_j)t]`$). Then, using those relationships, the following commutators $$[\widehat{\nu }_1,\widehat{H}_2]=[\widehat{\omega }_1,\widehat{H}_2]=[\widehat{\nu }_2,\widehat{H}_1]=[\widehat{\omega }_2,\widehat{H}_1]=0,$$ (181) and the same properties (ACKNOWLEDGMENTS), (150) and (157), we can show that $`\widehat{y}_1(t)\widehat{F}_{11}(t)`$ $`=`$ $`i{\displaystyle \frac{\gamma }{2}}\sqrt{\widehat{T}\widehat{B}_{}}\left\{\mathrm{cos}\left[(\widehat{\nu }_2\widehat{\omega }_2)t\right]\mathrm{sin}(\widehat{\omega }_1t)+\mathrm{cos}\left[(\widehat{\nu }_2+\widehat{\omega }_2)t\right]\mathrm{sin}(\widehat{\omega }_1t)\right\}\widehat{H}_2^{1/4}`$ (183) $`+`$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{H}_2^{1/4}\left\{\mathrm{sin}\left[(\widehat{\nu }_1\widehat{\omega }_1)t\right]\mathrm{cos}(\widehat{\omega }_2t)\mathrm{sin}\left[(\widehat{\nu }_1+\widehat{\omega }_1)t\right]\mathrm{cos}(\widehat{\omega }_2t)\right\}\sqrt{\widehat{B}_+\widehat{T}^{}}`$ (184) $`\widehat{y}_1(t)\widehat{F}_{12}(t)`$ $`=`$ $`{\displaystyle \frac{\gamma }{2}}\sqrt{\widehat{T}\widehat{B}_{}}\left\{\mathrm{cos}\left[(\widehat{\nu }_2\widehat{\omega }_2)t\right]\mathrm{cos}(\widehat{\omega }_1t)+\mathrm{cos}\left[(\widehat{\nu }_2+\widehat{\omega }_2)t\right]\mathrm{cos}(\widehat{\omega }_1t)\right\}\sqrt{\widehat{T}\widehat{B}_{}}`$ (185) $`+`$ $`{\displaystyle \frac{\gamma }{2}}\widehat{H}_2^{1/4}\left\{\mathrm{sin}\left[(\widehat{\nu }_1+\widehat{\omega }_1)t\right]\mathrm{sin}(\widehat{\omega }_2t)\mathrm{sin}\left[(\widehat{\nu }_1\widehat{\omega }_1)t\right]\mathrm{sin}(\widehat{\omega }_2t)\right\}\widehat{H}_1^{1/4}`$ (186) $`\widehat{y}_2(t)\widehat{F}_{21}(t)`$ $`=`$ $`{\displaystyle \frac{\gamma }{2}}\sqrt{\widehat{B}_+\widehat{T}^{}}\left\{\mathrm{cos}\left[(\widehat{\nu }_1\widehat{\omega }_1)t\right]\mathrm{cos}(\widehat{\omega }_2t)+\mathrm{cos}\left[(\widehat{\nu }_1+\widehat{\omega }_1)t\right]\mathrm{cos}(\widehat{\omega }_2t)\right\}\sqrt{\widehat{B}_+\widehat{T}^{}}`$ (187) $`+`$ $`{\displaystyle \frac{\gamma }{2}}\widehat{H}_1^{1/4}\left\{\mathrm{sin}\left[(\widehat{\nu }_2+\widehat{\omega }_2)t\right]\mathrm{sin}(\widehat{\omega }_1t)\mathrm{sin}\left[(\widehat{\nu }_2\widehat{\omega }_2)t\right]\mathrm{sin}(\widehat{\omega }_1t)\right\}\widehat{H}_2^{1/4}`$ (188) $`\widehat{y}_2(t)\widehat{F}_{22}(t)`$ $`=`$ $`i{\displaystyle \frac{\gamma }{2}}\sqrt{\widehat{B}_+\widehat{T}^{}}\left\{\mathrm{cos}\left[(\widehat{\nu }_1\widehat{\omega }_1)t\right]\mathrm{sin}(\widehat{\omega }_2t)+\mathrm{cos}\left[(\widehat{\nu }_1+\widehat{\omega }_1)t\right]\mathrm{sin}(\widehat{\omega }_2t)\right\}\widehat{H}_1^{1/4}`$ (189) $`+`$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{H}_1^{1/4}\left\{\mathrm{sin}\left[(\widehat{\nu }_2\widehat{\omega }_2)t\right]\mathrm{cos}(\widehat{\omega }_1t)\mathrm{sin}\left[(\widehat{\nu }_2+\widehat{\omega }_2)t\right]\mathrm{cos}(\widehat{\omega }_1t)\right\}\sqrt{\widehat{T}\widehat{B}_{}}.`$ (190) In a similar way, we can show that $`\widehat{z}_1(t)\widehat{F}_{11}(t)`$ $`=`$ $`i{\displaystyle \frac{\gamma }{2}}\sqrt{\widehat{T}\widehat{B}_{}}\left\{\mathrm{sin}\left[(\widehat{\nu }_2\widehat{\omega }_2)t\right]\mathrm{sin}(\widehat{\omega }_1t)+\mathrm{sin}\left[(\widehat{\nu }_2+\widehat{\omega }_2)t\right]\mathrm{sin}(\widehat{\omega }_1t)\right\}\widehat{H}_2^{1/4}`$ (192) $``$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{H}_2^{1/4}\left\{\mathrm{cos}\left[(\widehat{\nu }_1\widehat{\omega }_1)t\right]\mathrm{cos}(\widehat{\omega }_2t)\mathrm{cos}\left[(\widehat{\nu }_1+\widehat{\omega }_1)t\right]\mathrm{cos}(\widehat{\omega }_2t)\right\}\sqrt{\widehat{B}_+\widehat{T}^{}}`$ (193) $`\widehat{z}_1(t)\widehat{F}_{12}(t)`$ $`=`$ $`{\displaystyle \frac{\gamma }{2}}\sqrt{\widehat{T}\widehat{B}_{}}\left\{\mathrm{sin}\left[(\widehat{\nu }_2\widehat{\omega }_2)t\right]\mathrm{cos}(\widehat{\omega }_1t)+\mathrm{sin}\left[(\widehat{\nu }_2+\widehat{\omega }_2)t\right]\mathrm{cos}(\widehat{\omega }_1t)\right\}\sqrt{\widehat{T}\widehat{B}_{}}`$ (194) $``$ $`{\displaystyle \frac{\gamma }{2}}\widehat{H}_2^{1/4}\left\{\mathrm{cos}\left[(\widehat{\nu }_1+\widehat{\omega }_1)t\right]\mathrm{sin}(\widehat{\omega }_2t)\mathrm{cos}\left[(\widehat{\nu }_1\widehat{\omega }_1)t\right]\mathrm{sin}(\widehat{\omega }_2t)\right\}\widehat{H}_1^{1/4}`$ (195) $`\widehat{z}_2(t)\widehat{F}_{21}(t)`$ $`=`$ $`{\displaystyle \frac{\gamma }{2}}\sqrt{\widehat{B}_+\widehat{T}^{}}\left\{\mathrm{sin}\left[(\widehat{\nu }_1\widehat{\omega }_1)t\right]\mathrm{cos}(\widehat{\omega }_2t)+\mathrm{sin}\left[(\widehat{\nu }_1+\widehat{\omega }_1)t\right]\mathrm{cos}(\widehat{\omega }_2t)\right\}\sqrt{\widehat{B}_+\widehat{T}^{}}`$ (196) $``$ $`{\displaystyle \frac{\gamma }{2}}\widehat{H}_1^{1/4}\left\{\mathrm{cos}\left[(\widehat{\nu }_2+\widehat{\omega }_2)t\right]\mathrm{sin}(\widehat{\omega }_1t)\mathrm{cos}\left[(\widehat{\nu }_2\widehat{\omega }_2)t\right]\mathrm{sin}(\widehat{\omega }_1t)\right\}\widehat{H}_2^{1/4}`$ (197) $`\widehat{z}_2(t)\widehat{F}_{22}(t)`$ $`=`$ $`i{\displaystyle \frac{\gamma }{2}}\sqrt{\widehat{B}_+\widehat{T}^{}}\left\{\mathrm{sin}\left[(\widehat{\nu }_1\widehat{\omega }_1)t\right]\mathrm{sin}(\widehat{\omega }_2t)+\mathrm{sin}\left[(\widehat{\nu }_1+\widehat{\omega }_1)t\right]\mathrm{sin}(\widehat{\omega }_2t)\right\}\widehat{H}_1^{1/4}`$ (198) $``$ $`i{\displaystyle \frac{\gamma }{2}}\widehat{H}_1^{1/4}\left\{\mathrm{cos}\left[(\widehat{\nu }_2\widehat{\omega }_2)t\right]\mathrm{cos}(\widehat{\omega }_1t)\mathrm{cos}\left[(\widehat{\nu }_2+\widehat{\omega }_2)t\right]\mathrm{cos}(\widehat{\omega }_1t)\right\}\sqrt{\widehat{T}\widehat{B}_{}}.`$ (199) The non-commutativity between the operators $`\widehat{\omega }_1`$ and $`\widehat{\omega }_2`$ imply that to calculate the integrals involving the terms given by Eqs. (181) and (ACKNOWLEDGMENTS) we need to use the series expansion of the trigonometric functions. In this case the integrals can be easily done because the time variable can be considered as a parameter factor. Finally, inserting these results into Eq. (97) it is trivial to find the expression (III) for the matrix elements of the particular solution.
warning/0005/cond-mat0005159.html	ar5iv	text	# Analytical results for a trapped, weakly-interacting Bose-Einstein condensate under rotation \[ ## Abstract We examine the problem of a repulsive, weakly-interacting and harmonically trapped Bose-Einstein condensate under rotation. We derive a simple analytic expression for the energy incorporating the interactions when the angular momentum per particle is between zero and one and find that the interaction energy decreases linearly as a function of the angular momentum in agreement with previous numerical and limiting analytical studies. PACS numbers: 03.75.Fi, 05.30.Jp, 67.40.Db, 67.40.Vs \] One of the questions that followed the experimental realization of Bose-Einstein condensation in vapours of alkali atoms is how they behave under rotation. The experimental observation of vortex states in a two-component system has been reported by Matthews et al. , and the existence of vortex states in a stirred one-component Bose-Einstein condensate has been demonstrated by Madison et al. . Theoretical studies have also addressed this problem both in the Thomas-Fermi limit of strong interactions between the atoms and in the limit of weak interactions . In the limit of weak interactions and for a given angular momentum $`L`$, there is a large degeneracy corresponding to the ways of distributing $`L`$ units of angular momentum among $`N`$ harmonically confined atoms. Diagonalizing the interaction within this space of degenerate states, Bertsch and Papenbrock determined the ground state energy numerically for a limited number of particles and an effective repulsive interaction between them; the same method can evidently provide the complete excitation spectrum of the system. One remarkable observation is that, for $`2LN`$, the ground-state energy appears to have a particularly simple form with an interaction energy which decreases linearly with $`L`$. This result has also been obtained analytically in Ref. in the region $`LN`$ with use of a Bogoliubov transformation, and it is also in agreement with the study of this problem using the mean-field approximation . Here, we present an analytic derivation of this linear dependence of the ground state energy on $`L`$. Our starting point is the hamiltonian $`\widehat{H}=\widehat{H}_0+\widehat{V}`$, where $`\widehat{H}_0={\displaystyle \underset{i}{}}{\displaystyle \frac{\mathrm{}^2}{2M}}\mathbf{}_i^2+{\displaystyle \underset{i}{}}{\displaystyle \frac{1}{2}}M\omega ^2r_i^2`$ (1) includes the kinetic energy of the particles and the potential energy due to the trapping potential, and $`\widehat{V}={\displaystyle \frac{1}{2}}U_0{\displaystyle \underset{ij}{}}\delta (𝐫_i𝐫_j)`$ (2) is the interaction energy, which we assume to be a contact interaction. Here, $`M`$ is the mass of the atoms, $`\omega `$ is the frequency of the trapping potential, which is assumed to be isotropic, and $`U_0=4\pi \mathrm{}^2a/M`$ is the strength of the effective two-body interaction, with $`a`$ being the scattering length for atom-atom collisions. We assume that $`a>0`$, i.e., we treat only the case of repulsive interactions between the atoms. Further, we restrict our attention to the domain $`nU_0\mathrm{}\omega ,`$ (3) where $`n`$ is the density of the condensed atoms. If the system has angular momentum $`L`$, this condition allows us to work within the subspace of states which are degenerate in the absence of interactions. All other states differ by an energy of order $`\mathrm{}\omega `$, which is much larger than $`nU_0`$ by assumption. It is convenient to work in the basis $$\|N_0,N_1,N_2,\mathrm{}\|0^{N_0},1^{N_1},2^{N_2},\mathrm{},$$ (4) where $`N_k`$ is the number of particles with angular momentum $`\mathrm{}k`$. These occupation numbers must be chosen to respect the obvious constraints on particle number and angular momentum: $`{\displaystyle \underset{k}{}}N_k=N;{\displaystyle \underset{k}{}}kN_k=L.`$ (5) We introduce annihilation and creation operators $`a_k`$ and $`a_k^{}`$, which destroy or create one particle with angular momentum $`\mathrm{}k`$. The number operator $`\widehat{N}`$ and the angular-momentum operator $`\widehat{L}`$ can then be expressed as $`\widehat{N}={\displaystyle \underset{k}{}}a_k^{}a_k;\widehat{L}={\displaystyle \underset{k}{}}\mathrm{}ka_k^{}a_k.`$ (6) We also introduce the raising and lowering operators, $`\widehat{L}_\pm `$, as $`\widehat{L}_+={\displaystyle \underset{k}{}}\sqrt{k+1}a_{k+1}^{}a_k;\widehat{L}_{}={\displaystyle \underset{k}{}}\sqrt{k+1}a_k^{}a_{k+1}.`$ (7) Finally, we write the interaction $`\widehat{V}`$, Eq. (2), as $`\widehat{V}={\displaystyle \frac{U_0}{2}}{\displaystyle \underset{K}{}}\left({\displaystyle \underset{k^{}=0}{\overset{K}{}}}{\displaystyle \frac{a_k^{}^{}}{\sqrt{I_k^{}}}}{\displaystyle \frac{a_{Kk^{}}^{}}{\sqrt{I_{Kk^{}}}}}\right)J_K`$ (8) $`\times \left({\displaystyle \underset{k=0}{\overset{K}{}}}{\displaystyle \frac{a_k}{\sqrt{I_k}}}{\displaystyle \frac{a_{Kk}}{\sqrt{I_{Kk}}}}\right),`$ (9) where $`I_k=k!`$ and $`J_K=K!/2^K`$. It is readily seen that $`[\widehat{L},\widehat{L}_\pm ]=\widehat{L}_\pm ,`$ (10) which shows that the operators $`\widehat{L}_\pm `$ raise or lower the angular momentum by one unit when applied to the states of Eq. (4). Further, both the interaction, $`\widehat{V}`$, and the hamiltonian commute with $`\widehat{L}_\pm `$, $`[\widehat{V},\widehat{L}_\pm ]=[\widehat{H},\widehat{L}_\pm ]=0.`$ (11) This shows that, if $`\|L,N`$ is an eigenstate of $`\widehat{V}`$ with eigenvalue $`_{L,N}`$, the state $`\widehat{L}_+\|L,N`$ will also be an eigenstate of $`\widehat{V}`$ with the same eigenenergy but with $`L+1`$ units of angular momentum. We now proceed to the demonstration that the spectrum is indeed linear if $`LN`$. For any given eigenstate with a given $`N`$ and $`L`$, the repeated application of $`\widehat{L}_+`$ will generate infinitely many new eigenstates with angular momenta $`L+1,L+2,\mathrm{}`$. These states are obtained by adding angular momentum to the center of mass coordinate while leaving the wavefunction unchanged in the relative coordinates of the particles. Since the interaction energy depends only on relative coordinates and since the center of mass variables separate for a harmonic potential, the interaction energy is the same for each of these states. For any given $`N`$ and $`L`$, it is useful to pay special attention to those “intrinsic” eigenstates, $`\|L,N_{\mathrm{int}}`$, which do not have any center of mass excitation. The observation that the interaction energy in the ground state decreases with increasing $`L`$ indicates that the ground state wavefunction is an intrinsic state. When $`\widehat{L}_+`$ acts on any eigenstate $`\|L1,N`$ (i.e., intrinsic or not), it creates a state $`\|L,N`$ with center of mass excitation. Such states are necessarily orthogonal to intrinsic states with angular momentum $`L`$, and thus $`{}_{\mathrm{int}}{}^{}L,N\|\widehat{L}_+\|L1,N=L1,N\|\widehat{L}_{}\|L,N_{\mathrm{int}}=0.`$ (12) Since the states $`\|L1,N`$ form a complete set with angular moment $`L1`$, we conclude from Eq. (12) that $`\widehat{L}_{}\|L,N_{\mathrm{int}}=0.`$ (13) Let us now consider two particular basis states, $`\|A`$ and $`\|B`$, for any $`L<N`$ in the occupation number representation of Eq. (4), $`\|A=\|NL+1,L2,1;\|B=\|NL,L,0.`$ (14) States $`\|A`$ and $`\|B`$ are not eigenstates of $`\widehat{V}`$ but will in general be components of these eigenstates. Similarly, consider the state $`\|C`$ $`\|C=\|NL+1,L1,0.`$ (15) Note that these three states have $`N_k=0`$ for all $`k3`$. Under the action of $`\widehat{L}_{}`$, $`\|A`$ and $`\|B`$ are the only states which can lead to state $`\|C`$. Therefore, if an intrinsic eigenstate $`\|L,N_{\mathrm{int}}`$ includes these basis states in the combination $`\alpha \|A+\beta \|B`$, we see from Eq. (13) that $`\widehat{L}_{}(\alpha \|A+\beta \|B)\|C=0.`$ (16) We can satisfy Eq. (16) either by choosing $`\alpha =\beta =0`$ or by insisting that $`{\displaystyle \frac{\beta }{\alpha }}=\left({\displaystyle \frac{2(L1)}{L(NL+1)}}\right)^{1/2}.`$ (17) The second option is interesting since it uniquely determines the ratio of two components in $`\|L,N_{\mathrm{int}}`$. Since the state $`\|L,N_{\mathrm{int}}`$ is an eigenstate of $`\widehat{V}`$, it follows that $`B\|\widehat{V}\|L,N_{\mathrm{int}}=_{L,N}B\|L,N_{\mathrm{int}}.`$ (18) However, the only component of the state $`\|L,N_{\mathrm{int}}`$ which can be connected to state $`\|B`$ by the interaction $`\widehat{V}`$ is precisely the linear superposition $`\alpha \|A+\beta \|B`$. Since the ratio $`\beta /\alpha `$ is known, this fact can be used to determine the eigenvalue $`_{L,N}`$. Since $`\|A`$ and $`\|B`$ involve particles with $`k=0,1`$, and 2 only, the relevant part of the interaction $`\widehat{V}`$ is thus $`\widehat{V}={\displaystyle \frac{1}{2}}U_0(a_0^{}a_0^{}a_0a_0+2a_1^{}a_0^{}a_1a_0+{\displaystyle \frac{1}{2}}a_1^{}a_1^{}a_1a_1`$ (19) $`+{\displaystyle \frac{1}{\sqrt{2}}}a_1^{}a_1^{}a_2a_0+{\displaystyle \frac{1}{\sqrt{2}}}a_2^{}a_0^{}a_1a_1).`$ (20) The required matrix elements are $`{\displaystyle \frac{\beta }{2}}B\|a_0^{}a_0^{}a_0a_0+2a_1^{}a_0^{}a_1a_0+{\displaystyle \frac{1}{2}}a_1^{}a_1^{}a_1a_1\|B=`$ (21) $`{\displaystyle \frac{\beta }{2}}[(NL)(NL1)+2L(NL)+{\displaystyle \frac{1}{2}}L(L1)],`$ (22) and $`{\displaystyle \frac{\alpha }{2\sqrt{2}}}B\|a_1^{}a_1^{}a_2a_0\|A={\displaystyle \frac{\alpha }{2\sqrt{2}}}\sqrt{(NL1)(L1)L}.`$ (23) It now follows directly from Eqs. (17), (18), (22), and (23) that $`_{L,N}={\displaystyle \frac{N(2NL2)}{4}}v_0,`$ (24) in agreement with the numerical results of Refs. . Here $`v_0=U_0\|\mathrm{\Phi }_0(𝐫)\|^4𝑑𝐫`$, and $`\mathrm{\Phi }_0`$ is the ground state wavefunction of the harmonic oscillator. To be precise, we have demonstrated that an intrinsic eigenstate either has the eigenvalue Eq. (24) or does not contain the components $`\|A`$ and $`\|B`$. Completeness indicates that at least one intrinsic eigenstate with eigenvalue given by Eq. (24) must exist. As we mentioned earlier, the interaction energy of Eq. (24) varies linearly with $`L`$ for the contact two-body potential of Eq. (2). This feature is more general. As long as one deals with a two-body potential that conserves the angular momentum and commutes with $`\widehat{L}_\pm `$, this linearity persists, with the slope $`_{L,N}/L`$ being dependent on the specific form of the interaction. We have exploited the special features of harmonically trapped, weakly interacting bosons to determine a single eigenvalue of $`\widehat{V}`$ from the properties of two (in general exponentially small) components of the wavefunction. The present elementary argument does not enable us to construct the full eigenfunction or to demonstrate that this eigenstate describes the ground state of the system. While the present method can be applied to other choices of $`\widehat{V}`$, there appears to be no simple generalization to the interesting case $`L>N`$. We thank B. Mottelson, C. J. Pethick and S. M. Reimann for helpful discussions. G.M.K. was supported by the European Commission, TMR program, contract No. ERBFMBICT 983142. G.M.K. would like to thank the Foundation of Research and Technology, Hellas (FORTH) for its hospitality.
warning/0005/cond-mat0005166.html	ar5iv	text	# Theory of the Fano Resonance in the STM Tunneling Density of States due to a Single Kondo Impurity ## Abstract The conduction electron density of states nearby single magnetic impurities, as measured recently by scanning tunneling microscopy (STM), is calculated, taking into account tunneling into conduction electron states only. The Kondo effect induces a narrow Fano resonance in the conduction electron density of states, while scattering off the d-level generates a weakly energy dependent Friedel oscillation. The line shape varies with the distance between STM tip and impurity, in qualitative agreement with experiments, but is very sensitive to details of the band structure. For a Co impurity the experimentally observed width and shift of the Kondo resonance are in accordance with those obtained from a combination of band structure and strongly correlated calculations. Recently, several groups have demonstrated using scanning tunneling microscopy that a magnetic Kondo impurity adsorbed on the surface of a normal metal causes a narrow, resonance-like structure in the electronic surface density of states (DOS), whose asymmetric line shape resembles that of a Fano resonance . The experiments were performed with single Ce atoms on Ag as well as with single Co atoms on Au and Cu surfaces by measuring the I-V characteristics of the tunneling current through the tip of a scanning tunneling microscope (STM) placed close to the surface and at a small distance $`R`$ from the magnetic atom (see Fig. 1(inset)). Although the Kondo resonance, formed in the local conduction electron density of states (LDOS) at the Fermi energy $`\epsilon _F`$ due to resonant spin flip scattering, has been known for a long time , the precise line shape was not studied earlier because of the limited spatial resolution available in experiments. Wide tunnel junctions, which were proposed as measurement devices , probe the averaged DOS rather than the LDOS. Such experiments exhibited a giant zero bias resistance peak or a weak conductance peak, the latter induced by assisted tunneling through impurities in the tunnel barrier . In the present Letter we present a detailed microscopic study of the Kondo line shape as measured by STM in the vicinity of a single magnetic ion. We assume that the STM current is predominantly due to tunneling into the conduction LDOS, i.e. we neglect direct tunneling into the d- or f-level of a Co or Ce ion. This assumption is justified, because the d- or f-level is localized deeply in the atomic core, and is sufficient to explain the experimental findings as seen below. When a discrete (single-particle) level is coupled to the conduction electron sea (bare DOS $`\rho _0`$) via a hybridization $`V`$, there is a twofold effect: (1) the discrete level is broadened and (2) the continuous conduction LDOS becomes modified. The resulting line shape in the LDOS is called Fano resonance, in reminiscence of the first study of this problem in the context of atomic physics . Here we generalize this problem to the interacting case, i.e. when the discrete level arises from a many-body effect, like the Kondo resonance. It is shown that the many-body correlation effects and the consecutive Fano line shape in the conduction LDOS may be understood in separate steps, thus greatly simplifying the theoretical treatment as compared to other studies . For concreteness we focus on a Co atom on Au and use the Anderson model with a fivefold orbital degeneracy of the d-level $`\epsilon _d<0`$, $`m=1\mathrm{}5`$, $`H`$ $`=`$ $`H_o+\epsilon _d{\displaystyle \underset{m,\sigma }{}}d_{m,\sigma }^{}d_{m,\sigma }^{}+V{\displaystyle \underset{m\sigma 𝐤}{}}\left(d_{m,\sigma }^{}a_{𝐤,\sigma }^{}+\text{h.c.}\right)`$ (1) $`+`$ $`U{\displaystyle \underset{(m,\sigma )(m^{},\sigma ^{})}{}}d_{m,\sigma }^{}d_{m,\sigma }^{}d_{m^{},\sigma ^{}}^{}d_{m^{},\sigma ^{}}^{}`$ (2) where $`H_o=_𝐤\epsilon _𝐤a_{𝐤,\sigma }^{}a_{𝐤,\sigma }`$ is the kinetic energy of the conduction electrons, and $`a_{𝐤,\sigma }`$, $`d_{m,\sigma }`$ are the electron operators in the conduction band and in the d-level, respectively. $`U`$ is the Coulomb repulsion between two electrons in any of the local levels. There are two types of resonances in the DOS of the impurity d-level: ($`i`$) the Co d-levels with a broadening $`\mathrm{\Delta }=\pi \|V\|^2\rho _0`$ and ($`ii`$) the Kondo resonance whose width is given by the Kondo temperature $`T_K`$. In order to make contact with experiment, one must obtain realistic estimates for the parameters of the model. Therefore, we have applied the semi-relativistic, screened Korringa-Kohn-Rostoker method in combination with the local spin-density approximation (LSDA) for calculating the self-consistent electronic structure of a Co impurity placed onto a Au(111) surface. The LSDA assumes a spontaneous magnetization and consecutive Stoner-like splitting of the local level, which is fictitious for a single impurity. Although this does not account for properties of dynamical origin like the splitting into lower and upper Hubbard states or the Kondo effect, it is an accurate method to determine static quantities like the average impurity occupation number $`n_d=_{m\sigma }d_{m,\sigma }^{}d_{m,\sigma }^{}`$. This is because these are determined on time scales much shorter than the spin flip time $`\tau =h/k_BT_K`$, so that the local moment is effectively static for this purpose. Our LSDA results are summarized in Fig. 1: The orbital degeneracy is lifted due to crystal field splitting, and $`\mathrm{\Delta }0.2`$eV. The on-site Coulomb repulsion $`U`$ is proportional to the LSDA Stoner splitting and may be estimated as $`U=2.8`$eV . Mainly due to sp-d hybridization, the Co d-levels are shifted downward compared to bulk Co, so that $`n_d8.8`$ instead of $`n_d=7`$ expected from the nuclear charge of Co. The excess charge is compensated by a positive conduction electron depletion cloud around the impurity. It follows that of the 11 possible charge states of the Co d-level, $`z=0,1,\mathrm{},10`$, the system fluctuates only between $`z=8,9,10`$. Using Hund’s rule, these may be identified with the empty ($`z=8`$, $`\sigma =1`$, fixed), singly ($`z=9`$, $`\sigma =\pm 1/2`$) and doubly ($`z=10`$, $`\sigma =0`$) occupied states of an effective spin-1/2 Anderson impurity model without orbital degeneracy. The level energy $`\overline{\epsilon }_d`$ and Coulomb repulsion $`\overline{U}`$ of the effective model are different from those of the original model, Eq. (2). They may be extracted as follows : According to Eq. (2) and neglecting fluctuations, the charge states $`z=0,1,\mathrm{},10`$ of the d-level have energies $`E(z)=z\epsilon _d+Uz(z1)/2`$ (Fig. 2(a)). Since $`E(z)`$ must have its minimum (ground state) at $`z=n_d8.8`$, we have $`\epsilon _d/U=(n_d1/2)`$. The energies of the singly and doubly occupied orbital, measured relative to the empty state, $`E(8)`$, are given by $`\overline{\epsilon }_d=E(9)E(8)`$, $`2\overline{\epsilon }_d+\overline{U}=E(10)E(8)`$ and may be estimated as $`\overline{\epsilon }_d=0.84`$eV, $`\overline{U}=2.84`$eV. The corresponding Co d-orbital spectral function $`A_d(\omega )`$, calculated using the non-crossing approximation (NCA) , is shown in Fig. 2 (b). In addition to the broad peaks located near $`\overline{\epsilon }_d`$ and $`\overline{\epsilon }_d+\overline{U}`$ the narrow Kondo resonance is clearly seen. The Kondo temperature, $`T_K=D\sqrt{2\mathrm{\Delta }/(\pi D)}\mathrm{exp}\{1/(2\rho _0J)\}`$, with the spin exchange coupling $`J=\mathrm{\Delta }/(\pi \rho _0)[1/\|\overline{\epsilon }_d\|+1/\|\overline{\epsilon }_d+\overline{U}\|]`$, may be estimated as $`T_K52K`$. Here $`D`$ is a band cutoff provided by the bulk and (111) surface bands of Au, respectively; we assume $`D\epsilon _F^{Au}=5.5`$eV. It follows from the parabolic form of $`E(z)`$ (Fig. 2 (a)) that $`\|\overline{\epsilon }_d\|<(>)\overline{\epsilon }_d+\overline{U}_d`$, if $`n_d<(>)\mathrm{\hspace{0.33em}9}`$. As a consequence of level repulsion, the Kondo resonance is somewhat shifted upward (present case, $`n_d=8.8`$; shown in Fig. 2 (c)) or downward from the Fermi level. According to local Fermi liquid theory , for $`T<T_K`$ the Kondo resonance is a pure potential scatterer. For the analytical treatment below we may, therefore, model $`A_d(\omega )=\frac{1}{\pi }\text{Im}G_d(\omega i\delta )`$ crudely as a sum of three Lorentzians. This corresponds to $`G_d(\omega i\delta )`$ $`=`$ $`{\displaystyle \frac{Z_d}{\omega \overline{\epsilon }_di\mathrm{\Delta }}}+{\displaystyle \frac{Z_U}{\omega \overline{\epsilon }_d\overline{U}i\mathrm{\Delta }}}`$ (3) $`+`$ $`{\displaystyle \frac{Z_K}{\omega \epsilon _KiT_K}},`$ (4) where $`\epsilon _K`$ is the position of the Kondo resonance. $`Z_d`$, $`Z_U`$ and $`Z_K`$ are the appropriate strengths of the poles. In the Kondo regime, at low temperatures ($`TT_K`$), unitary scattering implies $`Z_K\stackrel{(<)}{}\pi T_K/\mathrm{\Delta }`$. We now turn to the conduction electron LDOS as measured by the STM tip at a distance $`R`$ from the impurity. It is related to the electronic field operator smeared around the tip position R as $`\mathrm{\Psi }_𝐑=\mathrm{\Psi }(𝐱)U_𝐑(𝐱)d^2x`$, where $`U_𝐑(𝐱)`$ is a form factor centered at $`𝐱=𝐑`$, and $`\mathrm{\Psi }(𝐱)`$ is the two-dimensional Fourier transform of $`a_{𝐤,\sigma }`$. According to the Anderson model, the exact conduction electron $`t`$-matrix is given in terms of the d-electron Green’s function as $`t_\sigma (i\omega _n)=\frac{\mathrm{\Delta }}{\pi \rho _0}G_{d,\sigma }(i\omega _n)`$. Then the correction to the conduction electron Green’s function due to the presence of the impurity, $`\delta 𝒢_{R,\sigma }(i\omega _n)=𝒢_{R,\sigma }(i\omega _n)𝒢_{R,\sigma }^{(0)}(i\omega _n)`$, is $$\delta 𝒢_{R,\sigma }(i\omega _n)=𝒢_{R,\sigma }^{(0)}(i\omega _n)t_\sigma (i\omega _n)𝒢_{R,\sigma }^{(0)}(i\omega _n),$$ (5) where $`𝒢_{R,\sigma }^{(0)}(i\omega _n)=T_\tau \left(\mathrm{\Psi }_\sigma (0)\mathrm{\Psi }_{R,\sigma }^{}\right)_{\omega _n}`$. The measured perturbation in the tunneling LDOS at distance $`R`$ is $$\delta \rho _R(\omega )=\frac{1}{\pi }\text{Im}\delta 𝒢_{R,\sigma }(\omega i\delta ).$$ (6) Using Eq. (5), it can be expressed as $`\delta \rho _{R,\sigma }(\omega )`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}\left[\text{Im}𝒢_{R,\sigma }^{(0)}(\omega i\delta )\right]^2\times `$ (7) $`\left[\left(q_{R,\sigma }^21\right)\text{Im}t_\sigma (\omega i\delta )+2q_{R,\sigma }\text{Re}t_\sigma (\omega i\delta )\right],`$ (8) where we have defined $$q_{R,\sigma }=\frac{\text{Re}𝒢_{R,\sigma }^{(0)}(\omega i\delta )}{\text{Im}𝒢_{R,\sigma }^{(0)}(\omega i\delta )}.$$ (9) In the following we drop the spin index $`\sigma `$. $`\text{Im}𝒢_R^{(0)}(\omega i\delta )`$ and $`q_R`$ depend on $`R`$, but on the scale $`T_K`$ very weakly on the energy; thus from here on $`\omega =0`$ is taken in these quantities. At $`R=0`$, $`q_{R=0}`$ is identical to the asymmetry parameter $`q`$ of the Fano theory , and if $`t(\omega i\delta )`$ is a simple (single-particle) level, $`t(\omega i\delta )=(\mathrm{\Delta }/\pi \rho _0)/(\omega \epsilon _di\mathrm{\Delta })`$, the line shape of Eq. (7) reduces to Fano’s well-known expression $`\rho (\omega )=\rho _0+\delta \rho =\rho _0(x+q)^2/(x^2+1)`$, with $`x=(\omega \epsilon _d)/\mathrm{\Delta }`$ . However, Eq. (7) is not limited to the non-interacting case. E.g. in the Kondo problem, correlations are contained in $`t(i\omega _n)`$ and may be treated separately (see above), while the Fano line shape is due to mixing of real and imaginary parts in Eq. (7). The line shape is sensitive both to the scattering phase shift contained in $`t(\omega i\delta )`$ and to the space dependent phase of the free conduction electron wave function exhibited by $`q_R`$. In particular, it depends crucially on the local charge density and on the details of the band structure . In the Kondo regime ($`\|\overline{\epsilon }_d\|,\overline{U}\mathrm{\Delta }`$, $`TT_K`$), for $`\omega \stackrel{<}{}T_K`$, $`\|\epsilon _K\|<T_K\mathrm{\Delta }`$, and using the (approximate) unitarity relation $`Z_K\stackrel{(<)}{}\pi T_K/\mathrm{\Delta }`$ (see remark after Eq. (3)) the $`t`$-matrix may be cast in the form $$\text{Re}t(\omega i\delta )=\frac{1}{\rho _0}\left\{\frac{\epsilon }{1+\epsilon ^2}+\beta +𝒪(\frac{T_K}{\mathrm{\Delta }})\right\}$$ (11) $$\text{Im}t(\omega i\delta )=\frac{1}{\rho _0}\left\{\frac{1}{1+\epsilon ^2}\beta \frac{\mathrm{\Delta }}{\overline{\epsilon }_d}+𝒪(\frac{T_K}{\mathrm{\Delta }})\right\},$$ (12) where $`\beta =(Z_d/\pi )\overline{\epsilon }_d\mathrm{\Delta }/(\overline{\epsilon }_d^2+\mathrm{\Delta }^2)`$, and $`\epsilon =(\omega \epsilon _K)/T_K`$. In $`\beta `$ we have taken into account only the 9-fold occupied d-state, as this is closest to $`\epsilon _F`$ and the other charge states give only a small contribution to the frequency independent part of $`t(\omega )`$. Thus the final expression for the LDOS correction is $$\delta \rho _R(\omega )=\frac{[\text{Im}𝒢_{R,\sigma }^{(0)}(\omega i\delta )]^2}{\pi \rho _0}\left\{\frac{2q_R\epsilon +q_R^21}{\epsilon ^2+1}+C_R\right\},$$ (13) where $`C_R=\beta [2q_R(q_R^21)\mathrm{\Delta }/\overline{\epsilon }_d]`$ arises from potential scattering by the d-level and corresponds to a weakly energy dependent Friedel oscillation. The first part coming from the scattering by the Kondo resonance gives a Fano line shape in the tunneling LDOS, controlled by the parameter $`q_R`$. Eq. (13) can be fitted to the experimental data for a Co atom on a Au (111) surface with excellent agreement (Fig. 3). From the fit parameters (see the figure caption of Fig. 3) we can conclude that $`\epsilon _K>0`$. The value of the Kondo temperature $`T_K50`$K obtained from the fit is substantially smaller than the bulk value ($`T_K>300`$ K for Co ) as the coupling of the impurity to the conduction electrons is weaker on the surface. Both $`T_K`$ and the shift $`\epsilon _K`$ of the Kondo resonance are consistent with the predictions of the NCA calculation in combination with the LSDA for Co on a Au(111) surface. The remaining task is to calculate $`q_R`$ and $`C_R`$ as a function of $`R`$. We have considered tunneling of electrons from the tip (1) into the 3-dimensional Au bulk states as well as (2) into the 2-dimensional Au(111) surface band . In both cases a free electron-like band structure was assumed, neglecting, e.g., in case (2) corrections due to the herringbone surface reconstruction . In case (1) the bulk Au Fermi wave number was taken as $`k_F^{(b)}=1.21\AA ^1`$, while in case (2) it was extracted from the surface band structure mapped out in Ref. by scanning tunneling spectroscopy, $`k_F^{(s)}=0.189\AA ^1`$. A window function was adopted for the form factor $`U_𝐑(𝐱)`$, and the scattering by the Co d-orbital was taken to be isotropic, as the scattering phase shift does not strongly depend on angular momentum. From Eq. (13), the Lorentzian line shape with a weight $`A=\text{Re}[𝒢_R^{(0)}]^2/(\pi \rho _0)`$ is formed at those distances $`R`$ where $`\text{Re}𝒢_R^{(0)}=0`$ ($`A<0`$) or $`\text{Im}𝒢_R^{(0)}=0`$ ($`A>0`$). For case (1), the R dependence of the LDOS is demonstrated in Fig. 4. The line shape changes periodically between asymmetric Fano and Lorentzian line shape, the period given by the Friedel wave length, $`\lambda _F^{(b)}=\pi /k_F^{(b)}2.6\AA `$. The overall amplitude decreases with increasing distance. For case (2), the results are similar, however with an oscillation period of $`\lambda _F^{(s)}=\pi /k_F^{(s)}16.6\AA `$. The latter is in good quantitative agreement with the measured period of approximately $`16\AA `$, indicating that the two-dimensional Au(111) surface band gives an important contribution to the total LDOS measured by STM. However, the precise dependence of the line shape on $`R`$ is not reproduced by our simplifying assumption of a free electron band structure. In particular, the model calculation predicts an initially increasing $`q_R`$ as $`R`$ grows from $`R=0`$, while the fit of Eq. (13) to the experiments implies a decreasing $`q_R`$. The precise $`R`$-dependence of $`q_R`$ will require taking into account the detailed band structure as well as the additional scattering phase shift induced by the charge of the Co ion and its spatially extended screening cloud. This is beyond the scope of the present Letter. Summarizing, we have shown that the Kondo resonance in the magnetic impurity d-level density of states ($`T<T_K`$) causes a Fano line shape in the density of states measured by STM even if we take into account only the tunneling into the conduction electron states. As expected, the Fano shape arises independently of whether the local level is a single-particle orbital or a many-body resonance like in the Kondo effect. Superimposed on the Fano line is a weakly energy dependent Friedel oscillation induced by potential scattering off the broad d-level. From the fit to the experimental data for Co on Au(111) surfaces we conclude that the Kondo resonance is shifted upward from the Fermi level, and $`T_K`$ for Co on the surface is substantially reduced compared to the bulk value. By combining electronic structure calculations (LSDA) with methods for strong correlations (NCA) these findings are reproduced semi-quantitatively. In particular, the shift of the Kondo resonance is due to level repulsion between the d-level and the Kondo resonance. We demonstrated the dependence of the line shape on the distance of the tip from the impurity by using bulk as well as surface state Green’s functions. We emphasize, however, that details of the band structure have to be taken into account in order to reproduce this distance dependence quantitatively. We acknowledge useful discussions with R. Berndt, M. F. Crommie and W.-D. Schneider. J.K. is grateful for the hospitality of the Condensed Matter Physics group, Technical University of Budapest, where part of this work was performed. This work was supported by the Humboldt foundation (A.Z.), by the OTKA Postdoctoral Fellowship D32819 (O.Ú.), by Hungarian grants OTKA T024005, T029813, T030240 and by DFG through SFB 195 (J.K.).
warning/0005/hep-lat0005020.html	ar5iv	text	# Quantum Spin Formulation of the Principal Chiral Model ## 1 INTRODUCTION In the construction of QCD as a quantum link model was proposed. In this paper we follow this new approach to quantising field theories, known as D-theory to reformulate the principal chiral model. In this approach, classical fields are replaced by quantum operators. In order to compensate for the loss in the number of degrees of freedom brought about by passing from continuous field variables to discrete eigenspectra of the quantum operators, the theory is formulated with an additional dimension, which later disappears; classical fields emerge as low-energy excitations of the discrete variables, provided the $`(d+1)`$-dimensional theory has massless excitations: When the extent of the extra dimension becomes small in units of the correlation length $`\xi `$, the $`d`$-dimensional field theory emerges via dimensional reduction. The guiding principles in formulating such a field theory are symmetry considerations. In particular, the quantum operators are Lie group generators chosen so that the Hamiltonian has the same symmetries on a quantum level as the original action has on a classical level. The virtue of the D-theory formulation of field theories is that the quantum partition function $`Z=\text{Tr}\left(\mathrm{exp}\left(\beta 𝐇\right)\right)`$ is a trace over the discrete eigenvalues of the quantum operators that make up the Hamiltonian. This promises for a much more efficient treatment by numerical techniques, and it is hoped that cluster algorithms similar to those used for Heisenberg model calculations can be developed. In the remainder of this article, we formulate the principal chiral model as a quantum spin model. Using coherent state path integral techniques we find an effective action for the low-energy excitations of the model and finally we show how the two-dimensional principal chiral model of classical fields emerges via dimensional reduction. ## 2 PRINCIPAL CHIRAL MODEL, $`d=2`$ ### 2.1 D-Theory Formulation The action of the two-dimensional continuum target theory is the following: $$S[U]=\frac{1}{2g^2}d^2x\text{Tr}\left(_\mu U^{}(x)_\mu U(x)\right),$$ (1) where the $`U(x)`$ are unitary $`N\times N`$ matrices. On the lattice, derivatives are replaced by finite differences, and we are free to add and subtract constant terms in the action to obtain a lattice regularised action of the form $$S_{\text{lat}}[U]=\frac{1}{g^2}\underset{<xy>}{}\text{Tr}\left(U_x^{}U_y+V(U_xU_x^{})\right).$$ (2) Note that $`U_xU_x^{}=I`$, so that the potential term is only a constant. Here it has no influence on the physics, but it motivates a similar term in the D-theory Hamiltonian, which will then have some effect on the low energy effective theory of the model. The target theory has a global $`U(N)_LU(N)_R`$ symmetry of the form $`U_xU_x^{}=LU_xR^{}`$, where $`L`$ and $`R`$ are unitary matrices. It is known that this symmetry breaks to a $`U(N)_V`$ symmetry in the ground state. Let us now replace the classical fields $`U_x^{ij}`$ by quantum operators $`𝐔_x^{ij}`$ and write down a D-theory Hamiltonian, which evolves the two-dimensional system in an additional Euclidean time direction: $`𝐇={\displaystyle \frac{J}{2N}}{\displaystyle \underset{x,\mu }{}}[𝐔_x^{ij}\left(𝐔_{x+\mu }^{ij}\right)^{}+𝐔_{x+\mu }^{ij}\left(𝐔_x^{ij}\right)^{}`$ (3) $`+V\left(𝐔_x^{ij}\left(𝐔_x^{kj}\right)^{}\right)+\text{h.c.}].`$ We would like this Hamiltonian to have an $`U(N)_LU(N)_R`$ symmetry, i.e. $`[𝐆_L^a,𝐇]=[𝐆_R^a,𝐇]=0`$, where $`𝐆_L^a`$ and $`𝐆_R^b`$ are mutually commuting sets of $`U(N)`$ generators. This can be realised by embedding $`u(N)_L`$ and $`u(N)_R`$ diagonally in the algebra of $`U(2N)`$ as follows: Let $`\{\lambda ^a\}`$ be matrices of the defining representation of $`su(N)`$, with commutation relations given by $`[\lambda ^a,\lambda ^b]=2if_c^{ab}\lambda ^c`$. Then, $$\begin{array}{cc}\hfill [𝐆_L^a,𝐆_L^b]=2if_c^{ab}𝐆_L^c,& [𝐆_R^a,𝐆_R^b]=2if_c^{ab}𝐆_R^c,\hfill \\ \hfill [𝐆_R^a,𝐔^{ij}]=𝐔^{ik}\lambda _{kj}^a,& [𝐆_L^a,𝐔^{ij}]=\lambda _{ik}^a𝐔^{kj},\hfill \end{array}$$ $$\begin{array}{cc}\hfill [𝐆_R^a,𝐆_L^b]& =[𝐓,𝐆_L^a]=[𝐓,𝐆_R^a]=0,\hfill \\ \hfill [𝐓,𝐔^{ij}]& =2𝐔^{ij},\hfill \\ \hfill [\mathrm{}𝐔^{ij},\mathrm{}𝐔^k\mathrm{}]& =[\mathrm{}𝐔^{ij},\mathrm{}𝐔^k\mathrm{}]\hfill \\ & =i\left(\delta _{ik}\mathrm{}\lambda _j\mathrm{}^a𝐆_R^a+\delta _j\mathrm{}\mathrm{}\lambda _{ik}^a𝐆_L^a\right),\hfill \\ \hfill [\mathrm{}𝐔^{ij},\mathrm{}𝐔^k\mathrm{}]& =\hfill \\ & i\left(\delta _{ik}\mathrm{}\lambda _j\mathrm{}^a𝐆_R^a\delta _j\mathrm{}\mathrm{}\lambda _{ik}^a𝐆_L^a+\frac{2}{N}\delta _{ik}\delta _j\mathrm{}𝐓\right).\hfill \end{array}$$ If we restrict ourselves to representations of $`u(2N)`$ which correspond to rectangular Young tableaux with $`N`$ rows and $`n_c`$ columns, we can use a fermionic basis of rishons for our representation : $`𝐒_i^j`$ $`=𝐆_L^{ij}={\displaystyle \underset{x}{}}\left(\mathrm{}_x^{i\alpha }\mathrm{}_x^{j\alpha }{\displaystyle \frac{n_c}{2}}\delta ^{ij}\right),`$ $`𝐒_{N+i}^{N+j}`$ $`=𝐆_R^{ij}={\displaystyle \underset{x}{}}\left(r_x^{i\alpha }r_x^{j\alpha }{\displaystyle \frac{n_c}{2}}\delta ^{ij}\right),`$ $`𝐓`$ $`={\displaystyle \underset{x}{}}\left(r_x^{i\alpha }r_x^{i\alpha }\mathrm{}_x^{i\alpha }\mathrm{}_x^{i\alpha }\right),`$ $`𝐒_i^{N+j}`$ $`=𝐔_x^{ij}=\mathrm{}_x^{i\alpha }r_x^{j\alpha },`$ $`𝐒_{N+i}^j`$ $`=\left(𝐔_x^{}\right)^{ij}=\left(𝐔_x^{ji}\right)^{}=r_x^{i\alpha }\mathrm{}_x^{j\alpha },`$ (4) $`{\displaystyle \underset{i}{}}(\mathrm{}^{i\alpha }\mathrm{}^{i\beta }`$ $`+r^{i\alpha }r^{i\beta })=\delta ^{\alpha \beta }N,`$ (5) where $`\alpha =1,\mathrm{},n_c`$ is an additional colour index and $`i,j=1,\mathrm{},N`$. For convenience, we have chosen these generators not to be traceless. We then have $`𝐆_{L(R)}^a=\lambda _{ij}^a𝐆_{L(R)}^{ij}`$. The constraint (5) at each lattice point is needed to obtain the correct representation. For $`J>0`$, this model is antiferromagnetic and we consider a bipartite lattice, made up of sublattices A and B. For sites on sublattice A we choose one representation, and for sites on sublattice B we choose its conjugate representation. In our case, both representations share the same Young tableau, so they are unitarily equivalent. But they need not be the same representation, as we shall see in the next section. Note that the properties of $`𝐇`$ are completely determined, once the representation of $`u(2N)`$ has been specified. In particular, the physics will be independent of whether the generators are represented by fermionic or bosonic operators. ### 2.2 Low Energy Effective Theory Following reference we can set up a coherent state path integral of the form $`Z=𝒟Qe^S`$, where $`S=`$ $$\begin{array}{c}\frac{n_c}{4}_0^\beta 𝑑\tau _0^1𝑑u\left[\text{Tr}\left(Q(\tau ,u)\frac{Q(\tau ,u)}{u}\frac{Q(\tau ,u)}{\tau }\right)\right]\hfill \\ \hfill _0^\beta 𝑑\tau H(Q(\tau )).\end{array}$$ (6) Without going into the details, in this derivation it is important to note that coherent states $`\|q`$ are labeled by $`GL(N,)`$ matrices $`q`$ and that $`q\|𝐒_\alpha ^\beta (x)\|q=\eta _x(n_c/2)Q_\alpha ^\beta (x)`$. \[$`\eta _x=+1(1)`$ for $`x`$ in sublattice $`A`$($`B`$).\] The Q-field is of the form: $$\begin{array}{c}Q=\mathrm{exp}\left[\left(\begin{array}{cc}0& q\\ q^{}& 0\end{array}\right)\right]\hfill \\ \hfill \times \left(\begin{array}{cc}I_N& 0\\ 0& I_N\end{array}\right)\mathrm{exp}\left[\left(\begin{array}{cc}0& q\\ q^{}& 0\end{array}\right)\right],\end{array}$$ (7) where $`qGL(N,)`$. The following boundary conditions hold: $$Q(\tau ,0)=Q(\tau ^{},0),\text{for all}\tau ,\tau ^{};$$ $$Q(\tau ,1)=Q(\tau );Q(0,u)=Q(\beta ,u).$$ (8) We now decompose the matrix field $`q`$ into a unitary matrix field $`U`$ and a hermitian matrix field $`C`$: $`q=CU`$. Under $`U(N)_LU(N)_R`$ transformations, these fields transform as follows: $`qq^{}=LqR^{}`$, $`UU^{}=LUR^{}`$ and $`CC^{}=LCL^{}`$. Substituting this coset decomposition into (7) we find, $$Q=\left(\begin{array}{cc}\mathrm{cos}(2C)& \mathrm{sin}(2C)U\\ U^{}\mathrm{sin}(2C)& U^{}\mathrm{cos}(2C)U\end{array}\right).$$ (9) Now define: $`B\mathrm{sin}(2C)`$ and choose $`V[(\mathrm{sin}(2C))UU^{}(\mathrm{sin}(2C))]=V(B^2)`$ such that its coefficients are of order one and the minimum of $`3B^2+V(B^2)`$ is attained for $`B`$ of the form $`B=bI,\mathrm{\hspace{0.33em}0}<b<1`$. While all other terms in the Hamiltonian are proportional to $`a^2`$ ($`a`$ is the lattice spacing), the term $`3B^2+V(B^2)`$ is not, so that when taking the continuum limit $`a0`$, any fluctuations of this term around its minimum value are suppressed. The field $`B`$ is thus frozen in the value $`bI`$. Next, we decompose the field $`U(x)`$ into staggered and uniform components: $$U(x)\mathrm{\Omega }(x)\sqrt{1a^2L^{}(x)L(x)}+\eta _xaL(x),$$ $$\mathrm{\Omega }(x)L(x)+L^{}(x)\mathrm{\Omega }(x)=0,$$ (10) and $`\mathrm{\Omega }(x)`$ is unitary. Then integrate out the $`L`$ field to find an effective action for the long-wavelength uniform fluctuations of the form $$\begin{array}{c}S=\frac{n_c}{2}\eta _x\underset{x}{}_0^\beta 𝑑\tau \text{Tr}\left(\sqrt{1b^2}\mathrm{\Omega }_\tau \mathrm{\Omega }^{}\right)\hfill \\ \hfill +_0^\beta 𝑑\tau d^2x\frac{\rho _s}{2}\text{Tr}\left(_\mu \mathrm{\Omega }_\mu \mathrm{\Omega }^{}+\frac{1}{c^2}_\tau \mathrm{\Omega }_\tau \mathrm{\Omega }^{}\right).\end{array}$$ (11) Here, $`\rho _s=Jb^2n_c^2/2N`$ is the spin stiffness and $`c=Jb^2n_ca/(N\sqrt{1b^2})`$ is the spin wave velocity. The first term is a Berry Phase term. We therefore get a low energy effective action for the Goldstone modes associated with the symmetry breaking pattern $`U(N)_LU(N)_RU(N)`$. Finally, we assume that the correlation length is much larger than the extent of the additional time dimension, $`\xi \beta c`$, so that the system is dimensionally reduced to a two-dimensional model. We can then integrate over $`\tau `$. From , we have that for the two-dimensional system, $`\xi \mathrm{exp}(2\pi /(g^2N))=\mathrm{exp}(2\pi \beta \rho _s/N)`$, which is indeed consistent with $`\xi \beta c`$ in the zero-temperature $`(\beta \mathrm{})`$ limit. ## 3 SUMMARY AND CONCLUSIONS We have reformulated the principal chiral model as a quantum spin model, in which the fields take on discrete values (the eigenvalues of the operators). To compensate for this restriction, we introduced an extra time direction. Having chosen a particular representation for the operators in the Hamiltonian, we were then able to show that the dimensionally reduced theory emerges from the low-energy theory of the collective excitations of the discrete variables. We note that the representation chosen for the operators influences the symmetry breaking pattern. The extension of these results to quantum link QCD is currently under investigation. The virtue of the D-theory formulation of field theories is that numerical simulations should be much easier than in conventional field theory formulations, due to the discrete nature of the field variables.
warning/0005/cond-mat0005214.html	ar5iv	text	# Simple fluids with complex phase behavior 0pt0.4pt ## Abstract We find that a system of particles interacting through a simple isotropic potential with a softened core is able to exhibit a rich phase behavior including: a liquid-liquid phase transition in the supercooled phase, as has been suggested for water; a gas-liquid-liquid triple point; a freezing line with anomalous reentrant behavior. The essential ingredient leading to these features resides in that the potential investigated gives origin to two effective core radii. PACS numbers: 61.20.Gy, 64.70.-p Fluid-fluid or liquid-liquid transitions have been proposed for many different substances such as supercooled water,C,S,Ga,Se,Te,$`I_2`$,Cs,$`SiO_2`$ etc.. The transition may occur between stable phases as in carbon or between metastable phases as has been suggested in supercooled water. In general, these substances also exhibit different amorphous states as well as multiple crystalline structures. The complex intermolecular interactions characterizing the above materials are usually modelled through anisotropic potentials depending on the molecular orientation and thus the possibility that simple fluids interacting through isotropic potentials may exhibit similar exotic phase behaviors is a challenging field of investigation. As argued by Mishima and Stanley the competition between expanded and compact structures in fluid-fluid (liquid-liquid) transitions suggests that the potential should possess two equilibrium positions. The most obvious form with such feature is one with two wells. Such potentials were shown to give rise to water-like thermodynamic anomalies, though the presence of a new critical point could not be directly observed. Another form of interparticle interaction which could produce different equilibrium positions is that in which there is a region of negative curvature in the repulsive core: these so-called softened-core potentials were proposed by Stell and Hemmer who argued that they might produce a second transition if a first already exists. Recently, through a mixed numerical-mean field type calculation, it was found that a potential consisting of a softened-core plus an infinite range van der Waals attractive term may give rise to a second critical point. Very recently, molecular dynamics simulation showed for a softened-core potential with an attractive square well evidences of a transition between two fluid phases in the supercooled region. The purpose of this Letter is to report the findings of a study of the phase behavior of a system of particles interacting through a potential consisting of a hard core with a finite repulsive shoulder and an attractive square well. Our analysis, based mainly on thermodynamically self-consistent ($`TSC`$) integral equations for fluids and partly on Monte-Carlo ($`MC`$) simulations, shows the existence of a liquid-gas critical point in the stable fluid phase and of a liquid-liquid critical point in the supercooled region. The liquid-gas and liquid-liquid coexistence lines meet in a gas-liquid-liquid triple point. Moreover, the behavior of the freezing line, estimated through one-phase criteria, such as the Hansen-Verlet ($`HV`$) rule and the entropic criterion based on the analysis of residual multiparticle entropy, is consistent with the existence of multiple crystalline structures in the solid phase. This is the first time that a microscopic theory linking directly the behavior of the system to the form of the interparticle pair potential, predicts for a simple fluid the existence of a liquid-liquid critical point and of a gas-liquid-liquid triple point. The chosen potential has a repulsive part $`V_{rep}(r)`$ consisting of a hard core of radius $`r_0=\sigma `$ and a repulsive square shoulder of height $`ϵ`$ and radius $`r_1=2.5\sigma `$, plus an attractive component $`V_{attr}(r)`$ having the form of a square well of depth $`1.25ϵ`$ extending from $`r_1=2.5\sigma `$ to $`r_2=3\sigma `$. To reach a thorough comprehension of the role played by the different components of the potential we first study its purely repulsive part, and then consider the effect of adding the attractive component. As a preliminary investigation we study through $`MC`$ simulation a system of particles interacting through the potential $`V_{rep}(r)`$ in $`2D`$, since in this case the spatial arrangement of the particles can be easily visualized. Snapshots representing the configuration of the system at a low temperature and different densities are shown in Fig.1, together with the corresponding structure factors $`S(k)`$. As the density increases, the system rapidly undergoes a very strong ordering (Fig.1(b)) indicated by the huge first peak of $`S(k)`$ (the splitted second peak of $`S(k)`$ can be related to the presence of lattice dislocations evident in the figure). Upon further increasing the density, particles begin to enter the soft core and many hard cores get in contact: the first peak of $`S(k)`$ decreases while the second peak of $`S(k)`$ is considerably enhanced (Fig.1(d)). Finally at very high densities (not shown in the figure) the third peak of $`S(k)`$ becomes the highest one, while the first two peaks are greatly depressed. The progressive “rising and falling” of the peaks of the structure factors reflects the turning on/off of different effective length scales. When the temperature $`T`$ and the density $`\rho `$ are sufficiently small, the soft core is practically impenetrable and the particles behave as hard spheres of radius $`r_1`$. As $`T`$ and $`\rho `$ increase, more and more particles penetrate the soft core until this becomes scarcely influent and the system is essentially equivalent to an assembly of hard spheres of radius $`r_0`$. In general, the system can be considered a “mixture” of two populations of hard spheres, of radius $`r_0`$ and $`r_1`$ respectively. The relative concentration of the two species is fixed by the values of $`T`$ and $`\rho `$. Thus, in contrast to standard simple fluids, the system has three possible length scales: $`r_1`$, $`r_{10}=(r_1+r_0)/2`$ and $`r_0`$, and as many indicators of structural ordering, namely the peaks of the structure factor corresponding to the wavevectors $`k_1`$, $`k_{10}`$, $`k_0`$, associated to these lengths. We stress the crucial role of the finite repulsive shoulder; due to its presence particles can be in one of two “states”: this is the essential feature which opens the possibility of liquid-liquid-immiscibility in a pure substance. Let us now consider a system of particles interacting through the potential $`V_{rep}(r)`$ in $`3D`$. We study its structural and thermodynamical properties using the $`TSC`$ Roger-Young ($`RY`$) integral equation. According to $`HV`$ rule a fluid is expected to undergo crystallization when the first (main) peak of $`S(k)`$ attains the value $`2.85`$. This statement refers to simple fluids with a single length scale. In our case, different length scales come into play and one must consider all the associated indicators of structural ordering. As shown in Fig.2, their behavior, upon increasing the density at constant $`T`$, appears analogous to that observed in $`2D`$, with each peak in turn rising and then “falling” down (except that associated to the hard core radius which becomes higher and higher as $`\rho `$ increases). The overall trend is confirmed by $`MC`$ calculation, though the theory tends to overestimate the heights of $`S(k_{10})`$ and $`S(k_0)`$ at high densities. In Fig.3 we show the locii of the points of the plane $`T`$,$`\rho `$ for which $`S(k_1)`$, $`S(k_{10})`$ and $`S(k_0)`$ are equal to $`2.85`$. The freezing line predicted through a straightforward application of $`HV`$ rule to the “anomalous” simple fluid investigated, should coincide with the line which bounds the region where at least one of the peaks of the structure factor is greater than $`2.85`$. This line shows a reentrant behavior which suggests that the system may undergo crystallization upon the density decreasing. The behavior can be related to the presence of the finite repulsive shoulder. In fact, the fraction of particles penetrating the soft core increases with the density (at constant $`T`$) thus “generating” additional space for the system. This effect is particularly important where one length scale begins to become less effective in favor of the smaller one. In these regions the phenomenon may overcompensate the general decrease of the space available to the system upon the density increasing thus causing a tendency of the system to become less ordered: accordingly, the freezing line may have a negative derivative. We note, however, that this behavior might be overestimated in Fig.3. In fact, since different length scales are contemporarily effective and contribute to the global structural ordering of the system, we expect that a more realistic freezing line has a smoother behavior which “blends” together the upper portions of the different curves of Fig.3. The inset in Fig.3 shows a magnification of the low temperature-low density region. Here, since only one length scale ($`r_1`$) is effective, the prediction of $`HV`$ is expected to be reliable. The freezing line starts nearly vertical at a density $`\rho ^{}0.06`$ which corresponds to the freezing density $`\rho r_1^30.945`$ of a fluid of hard spheres of radius $`r_1`$. As $`\rho `$ increases, the freezing line bends and shows a reentrant behavior: the phenomenon is associated with the very onset of the soft core penetration which, for the reasons above discussed, has a disordering effect on the system. This is reflected in the anomalous behavior shown by the isothermal compressibility $`\chi _T`$, which suddenly increases with the density (see same figure). Correspondingly, the first peak of $`S(k)`$ abruptly decreases: this occurs also at lower temperatures (where the value of $`S(k_1)`$ remains however greater than $`2.85`$) along a line coinciding approximately with the extension of the portion of the freezing line with negative derivative. This line meets the $`T=0`$ axis at a density $`\rho ^{}0.09`$ which corresponds to the closest packing of hard spheres of radius $`r_1`$ (occurring at $`\rho r_1^3=\sqrt{2}`$). The above results lead to conclude that the region shadowed in the inset of Fig.3 corresponds to an expanded solid phase of the system. We now investigate the phase behavior of a system of particles interacting through the full potential $`V_{rep}(r)`$ plus $`V_{attr}(r)`$. Calculations are performed making use of the $`HMSA`$ $`TSC`$ equation, which is better suited than $`RY`$ equation for interparticle interactions including an attractive component and reduces to this for purely repulsive potentials. The phase diagram of the system is shown in Fig.4. Two coexistence curves occur, each terminating at a critical point, denoted $`C1`$ and $`C2`$. The critical densities and temperatures are respectively $`\rho _{C1}^{}=0.06`$, $`T_{C1}^{}=1.3`$, and $`\rho _{C2}^{}=0.77`$, $`T_{C2}^{}=0.55`$. These values were estimated using the rectilinear diameter rule and the scaling relationship for the width of the coexistence curve with the non-classical exponent $`\beta 0.325`$. Below $`T_{C1}`$ the system separates into a gas and a liquid phase. The liquid phase is not unique since, below $`T_{C2}`$ ($`T_{C2}<T_{C1})`$, separates into distinct low-density ($`LD`$) and high-density ($`HD`$) phases. Since the critical point $`C2`$ is well below the freezing line, the liquid-liquid transition occurs between metastable phases in the supercooled region of the system. This feature recalls the scenario proposed for water, but in that case the liquid-liquid coexistence line is expected to start from $`C2`$ running at higher pressures as $`T`$ decreases. In the system investigated, the contrary is observed, this line running at lower pressures as $`T`$ decreases (see inset of Fig.4). This makes possible a new feature: the simultaneous coexistence of three fluid phases. In fact, the gas-liquid and the liquid-liquid coexistence lines meet in a gas-liquid-liquid triple ($`GLL`$) point which lays in the supercooled phase ($`\rho _{GLL}^{}0.57`$, $`T_{GLL}^{}0.53`$). The addition of the attractive well causes a shift towards higher temperatures of the locii of the points where $`S(k_1)`$, $`S(k_{10})`$ and $`S(k_0)`$ are equal to $`2.85`$, while their location in density remains essentially unaltered. Only a small portion of the line $`S(k_1)=2.85`$ is visible since it lays almost entirely in a region, correspondingly approximately to the gas-liquid spinodal decomposition, where the theory is unstable. The freezing line corresponding to $`HV`$ rule and that estimated through the entropic criterion are in substantial agreement at low and high densities (pressures), yielding close estimates of the gas-liquid-solid ($`GLS`$) triple point ($`\rho _{GLS}^{}0.29`$, $`T_{GLS}^{}0.97`$). On the contrary, these lines differ markedely in the intermediate region, where the first one shows a very evident reentrant behavior, while in the second approach this can be appreciated only numerically. For the reasons exposed in the purely repulsive case, the second result is probably more reliable. Thus, as shown in the inset in Fig.4, the freezing temperature increases initially with pressure, then remains essentially constant (according to the entropic criterion) in the pressure range $`1P^{}2`$, ($`P=P\sigma ^3/ϵ`$) and eventually increases again with pressure. Though our estimate of the freezing line is based solely on one-phase criteria, its shape, with branches having distinctly different slopes, is consistent with the possibility that structural changes occur in the solid state of the system. Consequently, transitions may be possible between solid phases of the system investigated. The results presented in this Letter show that a pure model system, with a softened-core isotropic potential, may have a rich phase behavior with features typical of substances characterized by much more complex anisotropic interactions. The possibility that simple substances related to the model investigated may exist (or may be “realized”) in nature well deserves to be investigated. We wish to thank F. Sciortino, P.Ballone, P.V. Giaquinta, C.Caccamo, S.Dugdale, for helpful suggestions and for interesting and stimulating discussions. One of us (G.M.) wishes to thank H.E.Stanley for stimulating his interest on the subject.
warning/0005/hep-th0005273.html	ar5iv	text	# 1 Introduction ## 1 Introduction There has been considerable work aimed at formulating models of quantum field theory on non–commutative spaces. The motivation is to obtain new insights into the UV–divergences and the problem of renormalization. On some simple noncommutative spaces, it is now possible to formulate quantum field theories. In some case, the UV divergences are completely regularized , while in others they persist . Moreover, it was realized that such noncommutative spaces are in fact induced by certain sectors of string theory, particularly open strings ending on $`D`$–branes with a background $`B`$ field . This is both a valuable source of physical insights, as well as a vindication of a more “puristic” approach of studying such spaces per se. In particular, spaces with quantum group symmetries have also been studied from a more formal approach. While quantum groups appear naturally in the context of 2–dimensional conformal field theories , a formulation of a quantum field theory based on such spaces has proved to be difficult. Recently, Alekseev, Recknagel and Schomerus have found that spherical $`D`$–branes in the $`SU(2)`$ WZW model are seen by open strings ending on them (in an appropriate background) as certain quasi–associative algebras, which are closely related to $`q`$–deformations of fuzzy spheres. Here $`q`$ is related to the level $`k`$ of the WZW model by the formula $$q=\mathrm{exp}(\frac{i\pi }{k+2}).$$ (1.1) We shall take this as sufficient motivation to study in detail the $`q`$–deformed fuzzy spheres, and to formulate field theories on them. The algebra found in is (weakly) non–associative, and covariant under $`SU(2)`$. Using a so–called Drinfeld twist, it can be transformed into an associative algebra which we call $`𝒮_{q,N}^2`$. It is covariant under the “quantum group” $`U_q(su(2))`$, which is the quantized universal enveloping algebra of Drinfeld and Jimbo . Here $`N`$ is an integer related to a particular boundary condition on the $`D`$–brane in $`SU(2)`$ WZW model. After reviewing the undeformed fuzzy sphere, we define $`𝒮_{q,N}^2`$ in Section 2 for both $`q`$ and $`\|q\|=1`$. As an algebra, it is simply a finite–dimensional matrix algebra, equipped with additional structure such as an action of $`U_q(su(2))`$, a covariant differential calculus, a star structure, and an integral. For $`q`$, this is precisely the “discrete” series of Podleś spheres . The case $`\|q\|=1`$, which is most relevant to string theory, has apparently not been studied in detail in the literature. In Section 3, we develop the non–commutative differential geometry on $`𝒮_{q,N}^2`$, using an approach which is suitable for both $`q`$ and $`\|q\|=1`$. The differential calculus turns out to be rather elaborate, but quite satisfactory. We are able to show, in particular, that in both cases there exists a 3–dimensional exterior differential calculus with real structure and a Hodge star, and we develop a frame formalism . This allows us to write Lagrangians for field theories on $`𝒮_{q,N}^2`$. In particular, the fact that the tangential space is 3–dimensional unlike in the classical case turns out to be very interesting physically, and is related to recent results on Chern–Simons actions on the $`D`$–branes. Using these tools, we study in Section 4 actions for scalar fields and abelian gauge fields on $`𝒮_{q,N}^2`$. The latter case is particularly interesting, since it turns out that certain actions for gauge theories arise in a very natural way in terms of polynomials of one–forms. In particular, the kinetic terms arise automatically due to the noncommutativity of the space. Moreover, because the calculus is 3–dimensional, the gauge field consists of a usual (abelian) gauge field plus a (pseudo) scalar in the classical limit. This is similar to a Kaluza–Klein reduction. One naturally obtains analogs of Yang–Mills and Chern–Simons actions, again because the calculus is 3–dimensional. In a certain limit where $`q=1`$, such actions were shown to arise from open strings ending on $`D`$–branes in the $`SU(2)`$ WZW model . The gauge theory actions for $`q1`$ suggest a new version of gauge invariance, where the gauge “group” is a quotient $`U_q(su(2))/I`$, which can be identified with the space of functions on the deformed fuzzy sphere. This is discussed in Section 4.2. Finally in Section 5, we give the precise relation of $`𝒮_{q,N}^2`$ to the quasi–associative algebra of functions on $`D`$–branes found in , using a Drinfeld–twist. In this paper, we shall only consider the first–quantized situation; the second quantization is postponed to a forthcoming paper . The latter turns out to be necessary for implementing the symmetry $`U_q(su(2))`$ on the space of fields in a fully satisfactory way. ## 2 The $`q`$–deformed fuzzy sphere ### 2.1 Review of the undeformed case We briefly recall the definition of the “standard” fuzzy sphere . Much information about the standard unit sphere $`S^2`$ in $`𝐑^\mathrm{𝟑}`$ is encoded in the infinite dimensional algebra of polynomials generated by $`\stackrel{~}{x}=(\stackrel{~}{x}_1,\stackrel{~}{x}_2,\stackrel{~}{x}_3)𝐑^\mathrm{𝟑}`$ with the defining relations $$[\stackrel{~}{x}_i,\stackrel{~}{x}_j]=0,\underset{i=1}{\overset{3}{}}\stackrel{~}{x}_i^2=r^2$$ (2.1) The algebra of functions on the fuzzy sphere is defined as the finite algebra $`𝒮_N^2`$ generated by $`\widehat{x}=(\widehat{x}_1,\widehat{x}_2,\widehat{x}_3)`$, with relations $$[\widehat{x}_i,\widehat{x}_j]=i\lambda _N\epsilon _{ijk}\widehat{x}_k,\underset{i=1}{\overset{3}{}}\widehat{x}_i^2=r^2.$$ (2.2) The real parameter $`\lambda _N>0`$ characterizes the non-commutativity. These relations are realized in a suitable finite–dimensional irreducible unitary representations of the $`SU(2)`$ group. This is most conveniently done using the Wigner-Jordan realization of the generators $`\widehat{x}_i,i=1,2,3`$, in terms of two pairs of annihilation and creation operators $`A_\alpha ,A_{}^{+}{}_{}{}^{\alpha },\alpha =\pm \frac{1}{2}`$, which satisfy $$[A_\alpha ,A_\beta ]=[A_{}^{+}{}_{}{}^{\alpha },A_{}^{+}{}_{}{}^{\beta }]=0,[A_\alpha ,A_{}^{+}{}_{}{}^{\beta }]=\delta _\alpha ^\beta ,$$ (2.3) and act on the Fock space $``$ spanned by the vectors $$\|n_1,n_2=\frac{1}{\sqrt{n_1!n_2!}}(A_{}^{+}{}_{}{}^{\frac{1}{2}})^{n_1}(A_{}^{+}{}_{}{}^{\frac{1}{2}})^{n_2}\|0.$$ (2.4) Here $`\|0`$ is the vacuum defined by $`A_i\|0=0`$. The operators $`\widehat{x}_i`$ take the form $$\widehat{x}_i=\frac{\lambda _N}{\sqrt{2}}A_{}^{+}{}_{}{}^{\alpha ^{}}\epsilon _{\alpha ^{}\alpha }\sigma _i^{\alpha \beta }A_\beta .$$ (2.5) Here $`\epsilon _{\alpha \alpha ^{}}`$ is the antisymmetric tensor (spinor metric), and $`\sigma _i^{\alpha \beta }`$ are the Clebsch-Gordan coefficients, that is rescaled Pauli–matrices. The number operator is given by $`\widehat{N}=_\alpha A_{}^{+}{}_{}{}^{\alpha }A_\alpha `$. When restricted to the $`(N+1)`$-dimensional subspace $$_N=\{A_{}^{+}{}_{}{}^{\alpha _1}\mathrm{}A_{}^{+}{}_{}{}^{\alpha _N}\|0(N\text{creation operators})\}.$$ (2.6) it yields for any given $`N=0,1,2,\mathrm{}`$ the irreducible unitary representation in which the parameters $`\lambda _N`$ and $`r`$ are related as $$\frac{r}{\lambda _N}=\sqrt{\frac{N}{2}\left(\frac{N}{2}+1\right)}.$$ (2.7) The algebra $`𝒮_N^2`$ generated by the $`\widehat{x}_i`$ is clearly the simple matrix algebra $`Mat(N+1)`$. Under the adjoint action of $`SU(2)`$, it decomposes into the direct sum $`(1)(3)(5)\mathrm{}(2N+1)`$ of irreducible representations of $`SO(3)`$ . ### 2.2 The $`q`$–deformed fuzzy sphere The fuzzy sphere $`𝒮_N^2`$ is invariant under the action of $`SO(3)`$, or equivalently under the action of $`U(so(3))`$. We shall define finite algebras $`𝒮_{q,N}^2`$ generated by $`x_i`$ for $`i=1,0,1`$, which have completely analogous properties to those of $`𝒮_N^2`$, but which are covariant under the quantized universal enveloping algebra $`U_q(su(2))`$. This will be done for both $`q`$ and $`q`$ a phase, including the appropriate reality structure. In the first case, the $`𝒮_{q,N}^2`$ will turn out to be the “discrete series” of Podleś’ quantum spheres . Here we will study them more closely from the above point of view. However, we also allow $`q`$ to be a root of unity, with certain restrictions. In a twisted form, this case does appear naturally on $`D`$–branes in the $`SU(2)`$ WZW model, as was shown in . In order to make the analogy to the undeformed case obvious, we shall perform a $`q`$–deformed Jordan–Wigner construction, which is covariant under $`U_q(su(2))`$. To fix the notation, we recall the basic relations of $`U_q(su(2))`$ $`[H,X^\pm ]`$ $`=`$ $`\pm 2X^\pm ,`$ $`[X^+,X^{}]`$ $`=`$ $`{\displaystyle \frac{q^Hq^H}{qq^1}}=[H]_q`$ (2.8) where the $`q`$–numbers are defined as $`[n]_q=\frac{q^nq^n}{qq^1}`$. The action of $`U_q(su(2))`$ on a tensor product of representations is encoded in the coproduct<sup>4</sup><sup>4</sup>4We use the opposite coproduct than in the standard conventions, but nevertheless the invariant tensors and $`\widehat{R}`$ –matrices will be the standard ones. The reason for this is explained in Appendix A. $`\mathrm{\Delta }(H)`$ $`=`$ $`H1+1H`$ $`\mathrm{\Delta }(X^\pm )`$ $`=`$ $`X^\pm q^{H/2}+q^{H/2}X^\pm .`$ (2.9) The antipode and the counit are given by $`S(H)`$ $`=`$ $`H,S(X^+)=q^1X^+,S(X^{})=qX^{},`$ $`\epsilon (H)`$ $`=`$ $`\epsilon (X^\pm )=0.`$ (2.10) The star structure is related to the Cartan–Weyl involution $`\theta (X^\pm )=X^{},\theta (H)=H`$, and will be discussed below. All symbols will now be understood to carry a label “$`q`$”, which we shall omit. An algebra $`𝒜`$ is called an $`U_q(su(2))`$–module algebra if there exists an action $`U_q(su(2))\times 𝒜`$ $``$ $`𝒜,`$ $`(u,a)`$ $``$ $`ua`$ (2.11) which satisfies $`u(ab)=(u_{(1)}a)(u_{(2)}b)`$ for $`a,b𝒜`$. Here $`\mathrm{\Delta }(u)=u_{(1)}u_{(2)}`$ is the Sweedler notation for the coproduct. Consider $`q`$–deformed creation and anihilation operators $`A_\alpha ,A_{}^{+}{}_{}{}^{\alpha }`$ for $`\alpha =\pm \frac{1}{2}`$, which satisfy the relations (cp. ) $`A_{}^{+}{}_{}{}^{\alpha }A_\beta `$ $`=`$ $`\delta _\beta ^\alpha +q\widehat{R}_{\beta \delta }^{\alpha \gamma }A_\gamma A_{}^{+}{}_{}{}^{\delta }`$ $`(P^{})_{\gamma \delta }^{\alpha \beta }A_\alpha A_\beta `$ $`=`$ $`0`$ $`(P^{})_{\gamma \delta }^{\alpha \beta }A_{}^{+}{}_{}{}^{\delta }A_{}^{+}{}_{}{}^{\gamma }`$ $`=`$ $`0`$ (2.12) where $`\widehat{R}_{\beta \delta }^{\alpha \gamma }=q(P^+)_{\beta \delta }^{\alpha \gamma }q^1(P^{})_{\beta \delta }^{\alpha \gamma }`$ is the decomposition of the $`\widehat{R}`$ –matrix of $`U_q(su(2))`$ into the projection operators on the symmetric and antisymmetric part. They can be written as $`(P^{})_{\gamma \delta }^{\alpha \beta }`$ $`=`$ $`{\displaystyle \frac{1}{[2]_q}}\epsilon ^{\alpha \beta }\epsilon _{\gamma \delta },`$ $`(P^+)_{\gamma \delta }^{\alpha \beta }`$ $`=`$ $`\sigma _i^{\alpha \beta }\sigma _{\gamma \delta }^i`$ (2.13) Here $`\epsilon _{\alpha \beta }`$ is the $`q`$–deformed invariant antisymmetric tensor, and $`\sigma _{\alpha \beta }^i`$ are the $`q`$–deformed Clebsch–Gordan coefficients; they are given explicitly in the Appendix. The factor $`[2]_q^1`$ arises from the relation $`\epsilon ^{\alpha \beta }\epsilon _{\alpha \beta }=[2]_q`$. The above relations are covariant under $`U_q(su(2))`$, and define a left $`U_q(su(2))`$–module algebra. We shall denote the action on the generators with lower indices by $$uA_\alpha =A_\beta \pi _\alpha ^\beta (u),$$ (2.14) so that $`\pi _\beta ^\alpha (uv)=\pi _\gamma ^\alpha (u)\pi _\beta ^\gamma (v)`$ for $`u,vU_q(su(2))`$. The generators with upper indices transform in the contragredient representation, which means that $$A_\alpha ^+:=\epsilon _{\alpha \beta }A_{}^{+}{}_{}{}^{\beta }$$ (2.15) transforms in the same way under $`U_q(su(2))`$ as $`A_\alpha `$. We consider again the corresponding Fock space $``$ generated by the $`A_{}^{+}{}_{}{}^{\alpha }`$ acting on the vacuum $`\|0`$, and its sectors $$_N=\{A_{}^{+}{}_{}{}^{\alpha _1}\mathrm{}A_{}^{+}{}_{}{}^{\alpha _N}\|0(N\text{creation operators})\}.$$ (2.16) It is well–known that these subspaces $`_N`$ are $`N+1`$–dimensional, as they are when $`q=1`$, and it follows that they form irreducible representations of $`U_q(su(2))`$ (at root of unity, this will be true due to the restriction (2.36) we shall impose). This will be indicated by writing $`_N=(N+1)`$, and the decomposition of $``$ into irreducible representations is $$=_0_1_2\mathrm{}=(1)(2)(3)\mathrm{}$$ (2.17) Now we define $$\widehat{Z}_i=A_{}^{+}{}_{}{}^{\alpha ^{}}\epsilon _{\alpha \alpha ^{}}\sigma _i^{\alpha \beta }A_\beta $$ (2.18) and $$\widehat{N}=\underset{\alpha }{}A_{}^{+}{}_{}{}^{\alpha ^{}}\epsilon _{\alpha \alpha ^{}}\epsilon ^{\alpha \beta }A_\beta .$$ (2.19) After some calculations, these operators can be shown to satisfy the relations $`\epsilon _k^{ij}\widehat{Z}_i\widehat{Z}_j={\displaystyle \frac{q^1}{\sqrt{[2]_q}}}(q^1[2]_q\lambda \widehat{N})\widehat{Z}_k`$ (2.20) $`\widehat{Z}^2:=g^{ij}\widehat{Z}_i\widehat{Z}_j=q^2{\displaystyle \frac{[2]_q+\widehat{N}}{[2]_q}}\widehat{N}`$ (2.21) Here $`\lambda =(qq^1)`$, $`g^{ij}`$ is the $`q`$–deformed invariant tensor for spin 1 representations, and $`\epsilon _k^{ij}`$ is the corresponding $`q`$–deformed Clebsch–Gordan coefficient; they are given in the Appendix. Moreover, one can verify that $`\widehat{N}A_{}^{+}{}_{}{}^{\alpha }`$ $`=`$ $`q^3A_{}^{+}{}_{}{}^{\alpha }+q^2A_{}^{+}{}_{}{}^{\alpha }\widehat{N},`$ $`\widehat{N}A_\alpha `$ $`=`$ $`q^1A_\alpha +q^2A_\alpha \widehat{N},`$ (2.22) which implies that $$[\widehat{N},\widehat{Z}_i]=0.$$ On the subspace $`_N`$, the “number” operator $`\widehat{N}`$ takes the value $$\widehat{N}_N=q^{N2}[N]_q_N$$ (2.23) It is convenient to introduce also an undeformed number operator $`\widehat{n}`$ which has eigenvalues $$\widehat{n}_N=N_N,$$ in particular $`\widehat{n}A_\alpha =A_\alpha (\widehat{n}1)`$. On the subspaces $`_N`$, the relations (2.20) become $`\epsilon _k^{ij}x_ix_j`$ $`=`$ $`\mathrm{\Lambda }_Nx_k,`$ (2.24) $`xx`$ $`:=`$ $`g^{ij}x_ix_j=r^2.`$ (2.25) Here the variables have been rescaled to $`x_i`$ with $$x_i=r\frac{q^{\widehat{n}+2}}{\sqrt{[2]_qC_N}}\widehat{Z}_i.$$ The $`r`$ is a real number, and we have defined $`C_N`$ $`=`$ $`{\displaystyle \frac{[N]_q[N+2]_q}{[2]_q^2}},`$ $`\mathrm{\Lambda }_N`$ $`=`$ $`r{\displaystyle \frac{[2]_{q^{N+1}}}{\sqrt{[N]_q[N+2]_q}}}.`$ (2.26) Using a completeness relation (see Appendix A), (2.24) can equivalently be written as $$(P^{})_{kl}^{ij}x_ix_j=\frac{1}{[2]_{q^2}}\mathrm{\Lambda }_N\epsilon _{kl}^nx_n.$$ (2.27) There is no $`i`$ in the commutation relations, because we use a weight basis instead of Cartesian coordinates. One can check that these relations precisely reproduce the “discrete” series of Podleś’ quantum spheres (after another rescaling), see , Proposition 4.II. Hence we define $`𝒮_{q,N}^2`$ to be the algebra generated by the variables $`x_i`$ acting on $`_N`$. Equipped with a suitable star structure and a differential structure, this will be the $`q`$–deformed fuzzy sphere. It is easy to see that the algebra $`𝒮_{q,N}^2`$ is simply the full matrix algebra $`Mat(N+1)`$, i.e. it is the same algebra as $`𝒮_N^2`$ for $`q=1`$. This is because $`_N`$ is an irreducible representation of $`U_q(su(2))`$. To see it, we use complete reducibility of the space of polynomials in $`x_i`$ of degree $`k`$ to conclude that it decomposes into the direct sum of irreducible representations $`(1)(3)(5)\mathrm{}(2k+1)`$. Counting dimensions and noting that $`x_1^N0(2N+1)`$, it follows that $`dim(𝒮_{q,N}^2)=(N+1)^2=dimMat(N+1)`$, and hence $$𝒮_{q,N}^2=(1)(3)(5)\mathrm{}(2N+1).$$ (2.28) This is true even if $`q`$ is a root of unity provided (2.36) below holds, a relation which will be necessary for other reasons as well. This is the decomposition of the functions on the $`q`$–deformed fuzzy sphere into $`q`$–spherical harmonics, and it is automatically truncated. Note however that not all information about a (quantum) space is encoded in its algebra of functions; in addition, one must specify for example a differential calculus and symmetries. For example, the action of $`U_q(su(2))`$ on $`𝒮_{q,N}^2`$ is different from the action of $`U(su(2))`$ on $`𝒮_N^2`$. The covariance of $`𝒮_{q,N}^2`$ under $`U_q(su(2))`$ can also be stated in terms of the quantum adjoint action. It is convenient to consider the cross–product algebra $`U_q(su(2))\text{ }\times 𝒮_{q,N}^2`$, which as a vector space is equal to $`U_q(su(2))𝒮_{q,N}^2`$, equipped with an algebra structure defined by $$ux=(u_{(1)}x)u_{(2)}.$$ (2.29) Here the $``$ denotes the action of $`uU_q(su(2))`$ on $`x𝒮_{q,N}^2`$. Conversely, the action of $`U_q(su(2))`$ on $`𝒮_{q,N}^2`$ can be written as $`ux=u_{(1)}xSu_{(2)}`$. The relations (2.29) of $`U_q(su(2))\text{ }\times 𝒮_{q,N}^2`$ are automatically realized on the representation $`_N`$. Since both algebras $`𝒮_{q,N}^2`$ and $`U_q(su(2))`$ act on $`_N`$ and generate the full matrix algebra $`Mat(N+1)`$, it must be possible to express the generators of $`U_q(su(2))`$ in terms of the $`\widehat{Z}_i`$. The explicit relation can be obtained by comparing the relations (2.24) with (2.29). One finds $`X^+q^{H/2}`$ $`=`$ $`q^{N+3}\widehat{Z}_1,`$ $`X^{}q^{H/2}`$ $`=`$ $`q^{N+1}\widehat{Z}_1,`$ $`q^H`$ $`=`$ $`{\displaystyle \frac{[2]_{q^{N+1}}}{[2]_q}}+{\displaystyle \frac{q^{N+2}(qq^1)}{\sqrt{[2]_q}}}\widehat{Z}_0,`$ (2.30) if acting on $`_N`$. In fact, this defines an algebra map $$j:U_q(su(2))𝒮_{q,N}^2$$ (2.31) which satisfies $$j(u_{(1)})xj(Su_{(2)})=ux$$ (2.32) for $`x𝒮_{q,N}^2`$ and $`uU_q(su(2))`$. This is analogous to results in . We shall often omit $`j`$ from now on. In particular, $`𝒮_{q,N}^2`$ is the quotient of the algebra $`U_q(su(2))`$ by the relation (2.25)<sup>5</sup><sup>5</sup>5more precisely, its finite–dimensional representation.. The relations (2.32) and those of $`U_q(su(2))`$ can be verified explicitly using (2.20). Moreover, one can verify that it is represented correctly on $`_N`$ by observing that $`X^+(A_{1/2}^+)^N\|0=0`$, which means that $`(A_{1/2}^+)^N\|0`$ is the highest–weight vector of $`_N`$. ### 2.3 Reality structure for $`q`$ In order to define a real quantum space, we must also construct a star structure, which is an involutive anti–linear anti–algebra map. For real $`q`$, the algebra (2.12) is consistent with the following star structure $`(A_\alpha )^{}`$ $`=`$ $`A_{}^{+}{}_{}{}^{\alpha }`$ $`(A_{}^{+}{}_{}{}^{\alpha })^{}`$ $`=`$ $`A_\alpha `$ (2.33) This can be verified using the standard compatibility relations of the $`\widehat{R}`$ –matrix with the invariant tensor . On the generators $`x_i`$, it implies the relation $$x_i^{}=g^{ij}x_j,$$ (2.34) as well as the equality $$\widehat{N}^{}=\widehat{N}.$$ The algebras $`𝒮_{q,N}^2`$ are now precisely Podleś’ “discrete” $`C^{}`$ algebras $`\stackrel{~}{S}_{q,c(N+1)}^2`$. Using (2.30), this is equivalent to $$H^{}=H,(X^\pm )^{}=X^{},$$ (2.35) which is the star–structure for the compact form $`U_q(su(2))`$. It is well–known that there is a unique Hilbert space structure on the subspaces $`_N`$ such that they are unitary irreducible representations of $`U_q(su(2))`$. Then the above star is simply the operator adjoint. ### 2.4 Reality structure for $`q`$ a phase When $`q`$ is a phase, finding the correct star structure is not quite so easy. The difference with the case $`q`$ is that $`\mathrm{\Delta }(u^{})=()\mathrm{\Delta }^{}(u)`$ for $`\|q\|=1`$ and $`uU_q(su(2))`$, where $`\mathrm{\Delta }^{}`$ denotes the flipped coproduct. We shall define a star only on the algebra $`𝒮_{q,N}^2`$ generated by the $`x_i`$, and not on the full algebra generated by $`A_\alpha `$ and $`A_\alpha ^+`$. There appears to be an obvious choice at first sight, namely $`x_i^{}=x_i`$, which is indeed consistent with (2.24). However, it is the wrong choice for our purpose, because it induces the noncompact star structure $`U_q(sl(2,))`$. Instead, we define a star–structure on $`𝒮_{q,N}^2`$ as follows. The algebra $`U_q(su(2))`$ acts on the space $`𝒮_{q,N}^2`$, which generically decomposes as $`(1)(3)\mathrm{}(2N+1)`$. This decomposition should be a direct sum of unitary representations of the compact form of $`U_q(su(2))`$, which means that the star structure on $`U_q(su(2))`$ should be (2.35), as it is for real $`q`$. There is a slight complication, because not all finite–dimensional irreducible representations are unitary if $`q`$ is a phase . However, all representations with dimension $`2N+1`$ are unitary provided $`q`$ has the form $$q=e^{i\pi \phi },\text{with}\phi <\frac{1}{2N}.$$ (2.36) This will be assumed from now on. As was pointed out before, we can consider the algebra $`𝒮_{q,N}^2`$ as a quotient of $`U_q(su(2))`$ via (2.30). It acts on $`_N`$, which is an irreducible representation of $`U_q(su(2))`$, and hence has a natural Hilbert space structure. We define the star on the operator algebra $`𝒮_{q,N}^2`$ by the adjoint (that is by the matrix adjoint in an orthonormal basis), hence by the star (2.35) using the identification (2.30). There is a very convenient way to write down this star structure on the generators $`x_i`$, similar as in . It involves an element $`\omega `$ of an extension of $`U_q(su(2))`$ introduced by and , which implements the Weyl reflection on irreducible representations. The essential properties are $`\mathrm{\Delta }(\omega )`$ $`=`$ $`^1\omega \omega ,`$ (2.37) $`\omega u\omega ^1`$ $`=`$ $`\theta S^1(u),`$ (2.38) $`\omega ^2`$ $`=`$ $`vϵ,`$ (2.39) where $`v`$ and $`ϵ`$ are central elements in $`U_q(su(2))`$ which take the values $`q^{N(N+2)/2}`$ resp. $`(1)^N`$ on $`_N`$. Here $`=_1_2U_q(su(2))U_q(su(2))`$ is the universal $`R`$ element. In a suitable (weight) basis of a unitary representation of $`U_q(su(2))`$, the matrix representing $`\omega `$ is given the invariant tensor, $`\pi _j^i(\omega )=q^{N(N+2)/4}g^{ij}`$, and $`\omega ^{}=\omega ^1`$. This is discussed in detail in . From now on, we denote with $`\omega `$ the element in $`𝒮_{q,N}^2`$ which represents this element on $`_N`$. We claim that the star structure on $`𝒮_{q,N}^2`$ as explained above is given by the following formula: $$x_i^{}=\omega x_i\omega ^1=x_jL_{}^{}{}_{k}{}^{j}q^2g^{ki},$$ (2.40) where $$L_{}^{}{}_{j}{}^{i}=\pi _j^i(_1^1)_2^1$$ (2.41) as usual ; a priori, $`L_{}^{}{}_{i}{}^{i}U_q(su(2))`$, but it is understood here as an element of $`𝒮_{q,N}^2`$ via (2.30). One can easily verify using $`(\epsilon _k^{ij})^{}=\epsilon _k^{ji}`$ (for $`\|q\|=1`$) that (2.40) is consistent with the relations (2.24) and (2.25). In the limit $`q1`$, $`L_{}^{}{}_{j}{}^{i}\delta _j^i`$, therefore (2.40) agrees with (2.34) in the classical limit. Hence we define the $`q`$–deformed fuzzy sphere for $`q`$ a phase to be the algebra $`𝒮_{q,N}^2`$ equipped with the star–structure (2.40). To show that (2.40) is correct in the sense explained above, it is enough to verify that it induces the star structure (2.35) on $`U_q(su(2))`$, since both $`U_q(su(2))`$ and $`𝒮_{q,N}^2`$ generate the same algebra $`Mat(N+1)`$. This can easily be seen using (2.38) and (2.30). A somewhat related conjugation has been proposed in using the universal element $``$. ### 2.5 Invariant integral The integral on $`𝒮_{q,N}^2`$ is defined to be the unique functional on $`𝒮_{q,N}^2`$ which is invariant under the (quantum adjoint) action of $`U_q(su(2))`$. It is given by the projection on the trivial sector in the decomposition (2.28). We claim that it can be written explicitly using the quantum trace: $$\underset{𝒮_{q,N}^2}{}f(x_i):=4\pi r^2\frac{1}{[N+1]_q}\text{Tr}_q(f(x_i))=4\pi r^2\frac{1}{[N+1]_q}\text{Tr}(f(x_i)q^H)$$ (2.42) for $`f(x_i)𝒮_{q,N}^2`$, where the trace is taken on $`_N`$. Using $`S^2(u)=q^Huq^H`$ for $`uU_q(su(2))`$, it follows that $$\underset{𝒮_{q,N}^2}{}fg=\underset{𝒮_{q,N}^2}{}S^2(g)f.$$ (2.43) This means that it is indeed invariant under the quantum adjoint action, $`{\displaystyle \underset{𝒮_{q,N}^2}{}}uf(x_i)`$ $`=`$ $`{\displaystyle \underset{𝒮_{q,N}^2}{}}u_1f(x_i)S(u_2)`$ (2.44) $`=`$ $`{\displaystyle \underset{𝒮_{q,N}^2}{}}S^1(u_2)u_1f(x_i)=\epsilon (u){\displaystyle \underset{𝒮_{q,N}^2}{}}f(x_i),`$ using the identification (2.30). The normalization constant is obtained from $$\text{Tr}_q(1)=\text{Tr}(q^H)=q^N+q^{N2}+\mathrm{}+q^N=[N+1]_q$$ on $`_N`$, so that $`\underset{𝒮_{q,N}^2}{}1=4\pi r^2`$. ###### Lemma 2.1 Let $`f𝒮_{q,N}^2`$. Then $$\left(\underset{𝒮_{q,N}^2}{}f\right)^{}=\underset{𝒮_{q,N}^2}{}f^{}$$ (2.45) for real $`q`$, and $$\left(\underset{𝒮_{q,N}^2}{}f\right)^{}=\underset{𝒮_{q,N}^2}{}f^{}q^{2H}$$ (2.46) for $`q`$ a phase, with the appropriate star structure (2.34) respectively (2.40). In (2.46), we use (2.30). Proof Assume first that $`q`$ is real, and consider the functional $$I_{q,N}(f):=\text{Tr}(f^{}q^H)^{}$$ for $`f𝒮_{q,N}^2`$. Then $`I_{q,N}(uf)`$ $`=`$ $`\text{Tr}((u_1fS(u_2))^{}q^H)^{}`$ (2.47) $`=`$ $`\text{Tr}(S^1((u^{})_2)f^{}(u^{})_1q^H)^{}=\text{Tr}(f^{}(u^{})_1S((u^{})_2)q^H)^{}`$ $`=`$ $`\epsilon (u)I_{q,N}(f),`$ where $`(S(u))^{}=S^1(u^{})`$ and $`()\mathrm{\Delta }(u)=\mathrm{\Delta }(u^{})`$ was used. Hence $`I_{q,N}(f)`$ is invariant as well, and (2.45) follows using uniqueness of the integral (up to normalization). For $`\|q\|=1`$, we define $$\stackrel{~}{I}_{q,N}(f):=\text{Tr}(f^{}q^H)^{}$$ with the star structure (2.40). Using $`(S(u))^{}=S(u^{})`$ and $`()\mathrm{\Delta }(u)=\mathrm{\Delta }^{}(u^{})`$, an analogous calculation shows that $`\stackrel{~}{I}_{q,N}`$ is invariant under the action of $`U_q(su(2))`$, which again implies (2.46). $``$$``$ For $`\|q\|=1`$, the integral is neither real nor positive, hence it cannot be used for a GNS construction. Nevertheless, it is clearly the appropriate functional to define an action for field theory, since it is invariant under $`U_q(su(2))`$. To find a way out, we introduce an auxiliary antilinear algebra–map on $`𝒮_{q,N}^2`$ by $$\overline{f}=S^1(f^{})$$ (2.48) where $`S`$ is the antipode on $`U_q(su(2))`$, using (2.30). Note that $`S`$ preserves the relation (2.25), hence it is well–defined on $`𝒮_{q,N}^2`$. This is not a star structure, since $$\overline{\overline{f}}=S^2f$$ for $`\|q\|=1`$. Using (2.30), one finds in particular $$\overline{x_i}=g^{ij}x_j.$$ (2.49) This is clearly consistent with the relations (2.24) and (2.25). We claim that (2.46) can now be stated as $$\left(\underset{𝒮_{q,N}^2}{}f\right)^{}=\underset{𝒮_{q,N}^2}{}\overline{f}\text{for }\|q\|=1.$$ (2.50) To see this, observe first that $$\text{Tr}(S(f))=\text{Tr}(f),$$ (2.51) which follows either from the fact that $`\widehat{I}_{q,N}(f):=\text{Tr}(S^1(f)q^H)=\text{Tr}(S^1(q^Hf))`$ is yet another invariant functional, or using $`\omega f\omega ^1=\theta S^1(f)`$ together with the observation that the matrix representations of $`X^\pm `$ in a suitable basis are real. This implies $$\text{Tr}_q(f^{}q^{2H})=\text{Tr}(f^{}q^H)=\text{Tr}(S(q^HS^1(f^{})))=\text{Tr}(q^HS^1(f^{}))=\text{Tr}_q(\overline{f}),$$ (2.52) and (2.50) follows. Now we can write down a positive inner product on $`𝒮_{q,N}^2`$: ###### Lemma 2.2 The sesquilinear forms $$(f,g):=\underset{𝒮_{q,N}^2}{}f^{}g\text{for}q$$ (2.53) and $$(f,g):=\underset{𝒮_{q,N}^2}{}\overline{f}g,\text{for}\|q\|=1$$ (2.54) are hermitian, that is $`(f,g)^{}=(g,f)`$, and satisfy $$(f,ug)=(u^{}f,g)$$ (2.55) for both $`q`$ and $`\|q\|=1`$. They are positive definite provided (2.36) holds for $`\|q\|=1`$, and define a Hilbert space structure on $`𝒮_{q,N}^2`$. Proof For $`q`$, we have $`(f,ug)`$ $`=`$ $`{\displaystyle \underset{𝒮_{q,N}^2}{}}f^{}u_1gSu_2={\displaystyle \underset{𝒮_{q,N}^2}{}}S^1(u_2)f^{}u_1g={\displaystyle \underset{𝒮_{q,N}^2}{}}(S((u^{})_2)^{}f^{}((u^{})_1)^{}g`$ (2.56) $`=`$ $`{\displaystyle \underset{𝒮_{q,N}^2}{}}\left((u^{})_1fS(u^{})_2\right)^{}g=(u^{}f,g),`$ and hermiticity is immediate. For $`\|q\|=1`$, consider $`(f,ug)`$ $`=`$ $`{\displaystyle \underset{𝒮_{q,N}^2}{}}\overline{f}u_1gSu_2={\displaystyle \underset{𝒮_{q,N}^2}{}}S^1(u_2)S^1(f^{})u_1g={\displaystyle \underset{𝒮_{q,N}^2}{}}S^1\left((u^{})_1fS(u^{})_2\right)^{}g`$ (2.57) $`=`$ $`(u^{}f,g).`$ Hermiticity follows using (2.50): $$(f,g)^{}=\underset{𝒮_{q,N}^2}{}\overline{\overline{f}}\overline{g}=\underset{𝒮_{q,N}^2}{}S^2(f)\overline{g}=\underset{𝒮_{q,N}^2}{}\overline{g}f=(g,f).$$ Using the assumption (2.36) for $`\|q\|=1`$, it is not difficult to see that they are also positive–definite. $``$$``$ ## 3 Differential Calculus In order to write Lagrangians, it is convenient to use the notion of an (exterior) differential calculus . A covariant differential calculus over $`𝒮_{q,N}^2`$ is a graded bimodule $`\mathrm{\Omega }_{q,N}^{}=_n\mathrm{\Omega }_{q,N}^n`$ over $`𝒮_{q,N}^2`$ which is a $`U_q(su(2))`$–module algebra, together with an exterior derivative $`d`$ which satisfies $`d^2=0`$ and the graded Leibnitz rule. We define the dimension of a calculus to be the rank of $`\mathrm{\Omega }_{q,N}^1`$ as a free right $`𝒮_{q,N}^2`$–module. ### 3.1 First–order differential forms Differential calculi for the Podleś sphere have been studied before . It turns out that 2–dimensional calculi do not exist for the cases we are interested in; however there exists a unique 3–dimensional module of 1–forms. As opposed to the classical case, it contains an additional “radial” one–form. This will lead to an additional scalar field, which will be discussed later. By definition, it must be possible to write any term $`x_idx_j`$ in the form $`_kdx_kf_k(x)`$. Unfortunately the structure of the module of 1–forms turn out to be not quadratic, rather the $`f_k(x)`$ are polynomials of order up to 3. In order to make it more easily tractable and to find suitable reality structures, we will construct this calculus using a different basis. First, we will define the bimodule of 1–forms $`\mathrm{\Omega }_{q,N}^1`$ over $`𝒮_{q,N}^2`$ which is covariant under $`U_q(su(2))`$, such that $`\{dx_i\}_i`$ is a free right $`𝒮_{q,N}^2`$–module basis, together with a map $`d:𝒮_{q,N}^2\mathrm{\Omega }_{q,N}^1`$ which satisfies the Leibnitz rule. Higher–order differential forms will be discussed below. Consider a basis of one–forms $`\xi _i`$ for $`i=1,0,1`$ with the covariant commutation relations<sup>6</sup><sup>6</sup>6they are not equivalent to $`u\xi _i=u_{(1)}\xi _iu_{(2)}`$. $$x_i\xi _j=\widehat{R}_{ij}^{kl}\xi _kx_l,$$ (3.1) using the $`(3)(3)`$ $`\widehat{R}`$–matrix of $`U_q(su(2))`$. It has the projector decomposition $$\widehat{R}_{ij}^{kl}=q^2(P^+)_{ij}^{kl}q^2(P^{})_{ij}^{kl}+q^4(P^0)_{ij}^{kl},$$ (3.2) where $`(P^0)_{ij}^{kl}=\frac{1}{[3]_q}g^{kl}g_{ij}`$, and $`(P^{})_{ij}^{kl}=_n\frac{1}{[2]_{q^2}}\epsilon _n^{kl}\epsilon _{ij}^n`$. The relations (3.1) are consistent with (2.25) and (2.24), using the braiding relations $`\widehat{R}_{ij}^{kl}\widehat{R}_{lu}^{rs}\epsilon _n^{ju}`$ $`=`$ $`\epsilon _t^{kr}\widehat{R}_{in}^{ts},`$ (3.3) $`\widehat{R}_{ij}^{kl}\widehat{R}_{lu}^{rs}g^{ju}`$ $`=`$ $`g^{kr}\delta _i^s`$ (3.4) and the quantum Yang–Baxter equation $`\widehat{R}_{12}\widehat{R}_{23}\widehat{R}_{12}=\widehat{R}_{23}\widehat{R}_{12}\widehat{R}_{23}`$, in shorthand–notation . We define $`\mathrm{\Omega }_{q,N}^1`$ to be the free right module over $`𝒮_{q,N}^2`$ generated by the $`\xi _i`$. It is clearly a bimodule over $`𝒮_{q,N}^2`$. To define the exterior derivative, consider $$\mathrm{\Theta }:=x\xi =x_i\xi _jg^{ij},$$ (3.5) which is a singlet under $`U_q(su(2))`$. It turns out (see Appendix B) that $`[\mathrm{\Theta },x_i]0\mathrm{\Omega }_{q,N}^1`$. Hence $$df:=[\mathrm{\Theta },f(x)]$$ (3.6) defines a nontrivial derivation $`d:𝒮_{q,N}^2\mathrm{\Omega }_{q,N}^1`$, which completes the definition of the calculus up to first order. In particular, it is shown in Appendix B that $$dx_i=\mathrm{\Lambda }_N\epsilon _i^{nk}x_n\xi _k+(qq^1)(qx_i\mathrm{\Theta }r^2q^1\xi _i).$$ (3.7) Since all terms are linearly independent, this is a 3–dimensional first–order differential calculus, and by the uniqueness it agrees with the 3–dimensional calculus in . In view of (3.7), it is not surprising that the commutation relations between the generators $`x_i`$ and $`dx_i`$ are very complicated ; will not write them down here. The meaning of the $`\xi `$–forms will become more clear in Section 3.4. Using (7.5) and the relation $`\xi x=q^4x\xi `$, one finds that $$xdx=\left(\mathrm{\Lambda }_N^2+([2]_{q^2}2)r^2\right)\mathrm{\Theta }.$$ On the other hand, this must be equal to $`x_i\mathrm{\Theta }x_jg^{ij}r^2\mathrm{\Theta }`$, which implies that $$x_i\mathrm{\Theta }x_jg^{ij}=\alpha r^2\mathrm{\Theta }$$ with $$\alpha =[2]_{q^2}1\frac{\mathrm{\Lambda }_N^2}{r^2}=1\frac{1}{C_N}.$$ (3.8) Combining this, it follows that $$dxx=r^2\frac{1}{C_N}\mathrm{\Theta }=xdx.$$ (3.9) Moreover, using the identity (7.7) one finds $$\epsilon _i^{jk}x_jdx_k=(\alpha q^2)r^2\epsilon _i^{jk}x_j\xi _k\mathrm{\Lambda }_Nr^2\xi _i+q^2\mathrm{\Lambda }_Nx_i\mathrm{\Theta },$$ (3.10) which together with (3.7) yields $$\xi _i=\frac{q^2}{r^2}\mathrm{\Theta }x_i+\frac{q^2C_N\mathrm{\Lambda }_N}{r^4}\epsilon _i^{jk}x_jdx_kq^2(1q^2)\frac{C_N}{r^2}dx_i.$$ (3.11) ### 3.2 Higher–order differential forms Podleś has constructed an extension of the above 3–dimensional calculus including higher–order forms for a large class of quantum spheres. This class does not include ours, however, hence we will give a different construction based on $`\xi `$–variables, which will be suitable for $`q`$ a phase as well. Consider the algebra $$\xi _i\xi _j=q^2\widehat{R}_{ij}^{kl}\xi _k\xi _l$$ (3.12) which is equivalent to $`(P^+)_{kl}^{ij}\xi _i\xi _j=0`$, $`(P^0)_{kl}^{ij}\xi _i\xi _j=0`$ where $`P^+`$ and $`P^0`$ are the projectors on the symmetric components of $`(3)(3)`$ as above; hence the product is totally ($`q`$–) antisymmetric. It is not hard to see (and well–known) that the dimension of the space of polynomials of order $`n`$ in the $`\xi `$ is $`(3,3,1)`$ for $`n=(1,2,3)`$, and zero for $`n>3`$, as classically. We define $`\mathrm{\Omega }_{q,N}^n`$ to be the free right $`𝒮_{q,N}^2`$–module with the polynomials of order $`n`$ in $`\xi `$ as basis; this is covariant under $`U_q(su(2))`$. Then $`\mathrm{\Omega }_{q,N}^n`$ is in fact a (covariant) $`𝒮_{q,N}^2`$–bimodule, since the commutation relations (3.1) between $`x`$ and $`\xi `$ are consistent with (3.12), which follows from the quantum Yang–Baxter equation. There remains to construct the exterior derivative. To find it, we first note that (perhaps surprisingly) $`\mathrm{\Theta }^20`$, rather $$\mathrm{\Theta }^2=\frac{q^2\mathrm{\Lambda }_N}{[2]_{q^2}}\epsilon ^{ijk}x_i\xi _j\xi _k.$$ (3.13) The $`\epsilon ^{ijk}`$ is defined in (7.8). By a straightforward but lengthy calculation which is sketched in Appendix B, one can show that $$dx_idx_jg^{ij}+\frac{r^2}{C_N}\mathrm{\Theta }^2=0.$$ We will show below that an extension of the calculus to higher–order forms exists; then this can be rewritten as $$d\mathrm{\Theta }\mathrm{\Theta }^2=0.$$ (3.14) The fact that $`\mathrm{\Theta }^20`$ makes the construction of the extension more complicated, since now $`\alpha ^{(n)}[\mathrm{\Theta },\alpha ^{(n)}]_\pm `$ does not define an exterior derivative. To remedy this, the following observation is useful: the map $`_H:\mathrm{\Omega }_{q,N}^1`$ $``$ $`\mathrm{\Omega }_{q,N}^2,`$ $`\xi _i`$ $``$ $`{\displaystyle \frac{q^2\mathrm{\Lambda }_N}{[2]_{q^2}}}\epsilon _i^{jk}\xi _j\xi _k`$ (3.15) defines a left–and right $`𝒮_{q,N}^2`$–module map; in other words, the commutation relations between $`\xi _i`$ and $`x_j`$ are the same as between $`_H(\xi _i)`$ and $`x_j`$. This follows from the braiding relation (3.3). This is in fact the natural analogue of the Hodge–star on 1–forms in our context, and will be discussed further below. Here we note the important identity $$\alpha (_H\beta )=(_H\alpha )\beta $$ (3.16) for any $`\alpha ,\beta \mathrm{\Omega }_{q,N}^1`$, which is proved in Appendix B. Now (3.13) can be stated as $$_H(\mathrm{\Theta })=\mathrm{\Theta }^2,$$ (3.17) and applying $`_H`$ to $`df=[\mathrm{\Theta },f(Y)]`$ one obtains $$[\mathrm{\Theta }^2,f(x)]=_Hdf(x).$$ (3.18) Now we define the map $`d:\mathrm{\Omega }_{q,N}^1`$ $``$ $`\mathrm{\Omega }_{q,N}^2,`$ $`\alpha `$ $``$ $`[\mathrm{\Theta },\alpha ]_+_H(\alpha ).`$ (3.19) It is easy to see that this defines a graded derivation from $`\mathrm{\Omega }_{q,N}^1`$ to $`\mathrm{\Omega }_{q,N}^2`$, and the previous equation implies immediately that $$(dd)f=0.$$ In particular, $$d\xi _i=(1q^2)\xi \mathrm{\Theta }+\frac{q^2\mathrm{\Lambda }_N}{[2]_{q^2}}\epsilon _i^{jk}\xi _j\xi _k.$$ (3.20) To complete the differential calculus, we extend it to $`\mathrm{\Omega }_{q,N}^3`$ by $`d:\mathrm{\Omega }_{q,N}^2`$ $``$ $`\mathrm{\Omega }_{q,N}^3,`$ $`\alpha ^{(2)}`$ $``$ $`[\mathrm{\Theta },\alpha ^{(2)}].`$ (3.21) As is shown in Appendix B, this satisfies indeed $$(dd)\alpha =0\text{for any}\alpha \mathrm{\Omega }_{q,N}^1.$$ It is easy to see that the map (3.21) is non–trivial. Moreover there is precisely one monomial of order 3 in the $`\xi `$ variables, given by $$\mathrm{\Theta }^3=\frac{q^6\mathrm{\Lambda }_Nr^2}{[2]_{q^2}[3]_q}\epsilon ^{ijk}\xi _i\xi _j\xi _k,$$ (3.22) which commutes with all functions on the sphere, $$[\mathrm{\Theta }^3,f]=0$$ (3.23) for all $`f𝒮_{q,N}^2`$. Finally, we complete the definition of the Hodge star operator by $$_H(1)=\mathrm{\Theta }^3,$$ (3.24) and by requiring that $`(_H)^2=id`$. ### 3.3 Star structure A $``$–calculus (or a real form of $`\mathrm{\Omega }_{q,N}^{}`$) is a differential calculus which is a graded $``$–algebra such that the star preserves the grade, and satisfies $`(\alpha ^{(n)}\alpha ^{(m)})^{}`$ $`=`$ $`(1)^{nm}(\alpha ^{(m)})^{}(\alpha ^{(n)})^{},`$ $`(d\alpha ^{(n)})^{}`$ $`=`$ $`d(\alpha ^{(n)})^{}`$ (3.25) for $`\alpha ^{(n)}\mathrm{\Omega }_{q,N}^n`$; moreover, the action of $`U_q(su(2))`$ must be compatible with the star on $`U_q(su(2))`$. Again, we have to distinguish the cases $`q`$ and $`\|q\|=1`$. #### 1) $`q`$. In this case, the star structure must satisfy $$(dx_i)^{}=g^{ij}dx_j,x_i^{}=g^{ij}x_j,$$ (3.26) which by (2.25) implies $$\mathrm{\Theta }^{}=\mathrm{\Theta }.$$ (3.27) Using (3.11), it follows that $`\xi _i^{}`$ $`=`$ $`g^{ij}\xi _j+q^2(qq^1){\displaystyle \frac{[2]_qC_N}{r^2}}g^{ij}dx_j`$ (3.28) $`=`$ $`g^{ij}\xi _jq^2(qq^1){\displaystyle \frac{[2]_qC_N}{r^2}}g^{ij}\left(\mathrm{\Lambda }_N\epsilon _j^{kl}x_k\xi _l(qq^1)(qx_j\mathrm{\Theta }q^1r^2\xi _j)\right).`$ To show that this is indeed compatible with (3.1), one needs the following identity $$q^2(qq^1)\frac{[2]_qC_N}{r^2}(dx_ix_j\widehat{R}_{ij}^{kl}x_kdx_l)=(\mathrm{𝟏}(\widehat{R}^2)_{ij}^{kl})\xi _kx_l$$ (3.29) which can be verified with some effort, see Appendix B. In particular, this shows that if one imposed $`x_i\xi _j=(\widehat{R}^1)_{ij}^{kl}\xi _kx_l`$ instead of (3.1), one would obtain an equivalent calculus. This is unlike in the flat case, where one has two inequivalent calculi . Moreover, one can show that this real form is consistent with (3.12). #### 2) $`\|q\|=1`$. In view of (2.40), it is easy to see that the star structure in this case is $$(\xi _i)^{}=q^4\omega \xi _i\omega ^1,x_i^{}=\omega x_i\omega ^1.$$ (3.30) Recall that $`\omega `$ is a particular unitary element of $`𝒮_{q,N}^2`$ introduced in Section 2.4. It is obvious using $`(\widehat{R}_{ij}^{kl})^{}=(\widehat{R}^1)_{ji}^{lk}`$ that this is an involution which is consistent with (3.1), and one can verify that $$\mathrm{\Theta }^{}=\mathrm{\Theta }.$$ (3.31) This also implies $$[\omega ,\mathrm{\Theta }]=0,$$ hence $$(dx_i)^{}=\omega dx_i\omega ^1.$$ (3.32) Finally, $`_H`$ is also compatible with the star structure: $$(_H(\alpha ))^{}=_H(\alpha ^{})$$ (3.33) where $`\alpha \mathrm{\Omega }_{q,N}^1`$, for both $`q`$ and $`\|q\|=1`$. This is easy to see for $`\alpha =\xi _i`$ in the latter case, and for $`\alpha =dx_i`$ in the case $`q`$. This implies that indeed $`(d\alpha ^{(n)})^{}=d(\alpha ^{(n)})^{}`$ for all $`n`$. We summarize the above results: ###### Theorem 3.1 The definitions (3.19), (3.21) define a covariant differential calculus on $`\mathrm{\Omega }_{q,N}^{}=_{n=0}^3\mathrm{\Omega }_{q,N}^n`$ over $`𝒮_{q,N}^2`$ with $`dim(\mathrm{\Omega }_{q,N}^n)=(1,3,3,1)`$ for $`n=(0,1,2,3)`$. Moreover, this is a $``$–calculus with the star structures (3.26) and (3.30) for $`q`$ and $`\|q\|=1`$, respectively. ### 3.4 Frame formalism On many noncommutative spaces , it is possible to find a particularly convenient set of one–forms (a “frame”) $`\theta _a\mathrm{\Omega }^1`$, which commute with all elements in the function space $`\mathrm{\Omega }^0`$. Such a frame exists here as well, and in terms of the $`\xi _i`$ variables, it takes a similar form to that of . Consider the elements $`\theta ^a`$ $`=`$ $`\mathrm{\Lambda }_NS(L_{}^{+}{}_{j}{}^{a})g^{jk}\xi _k\mathrm{\Omega }_{q,N}^1,`$ (3.34) $`\lambda _a`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Lambda }_N}}x_iL_{}^{+}{}_{a}{}^{i}𝒮_{q,N}^2.`$ (3.35) where as usual $`L_{}^{+}{}_{j}{}^{i}`$ $`=`$ $`_1\pi _j^i(_2),`$ (3.36) $`S(L_{}^{+}{}_{j}{}^{i})`$ $`=`$ $`_1^1\pi _j^i(_2^1)`$ (3.37) are elements of $`U_q(su(2))`$, which we consider here as elements in $`𝒮_{q,N}^2`$ via (2.31). Then the following holds: ###### Lemma 3.2 $`[\theta ^a,f]`$ $`=`$ $`0,`$ (3.38) $`df`$ $`=`$ $`[\lambda _a,f]\theta ^a,`$ (3.39) $`\mathrm{\Theta }=x_i\xi _jg^{ij}`$ $`=`$ $`\lambda _a\theta ^a.`$ (3.40) for any $`f𝒮_{q,N}^2`$. In this sense, the $`\lambda _a`$ are dual to the frame $`\theta ^b`$. They satisfy the relations $`\lambda _a\lambda _bg^{ba}`$ $`=`$ $`{\displaystyle \frac{1}{q^4\mathrm{\Lambda }_N^2}}r^2,`$ $`\lambda _a\lambda _b\epsilon _c^{ba}`$ $`=`$ $`{\displaystyle \frac{1}{q^2}}\lambda _c,`$ $`\theta ^a\theta ^b`$ $`=`$ $`q^2\widehat{R}_{cd}^{ba}\theta ^d\theta ^c`$ (3.41) $`d\theta ^a`$ $`=`$ $`\lambda _b[\theta ^a,\theta ^b]_++{\displaystyle \frac{1}{q^2[2]_{q^2}}}\epsilon _{bc}^a\theta ^c\theta ^b`$ $`_H\theta ^a`$ $`=`$ $`{\displaystyle \frac{1}{q^2[2]_{q^2}}}\epsilon _{bc}^a\theta ^c\theta ^b`$ $`\theta ^a\theta ^b\theta ^c`$ $`=`$ $`\mathrm{\Lambda }_N^2{\displaystyle \frac{q^6}{r^2}}\epsilon ^{cba}\mathrm{\Theta }^3`$ (3.42) In particular in the limit $`q=1`$, this becomes $`\lambda _a=\frac{1}{\mathrm{\Lambda }_N}x_a`$, and $`dx_a=\epsilon _{ab}^cx_c\theta ^b`$, using (3.7). Proof Using $$S(L_{}^{+}{}_{j}{}^{i})x_k=x_l(\widehat{R}^1)_{jk}^{ln}S(L_{}^{+}{}_{n}{}^{i})$$ (which follows from (2.29)) and $`\mathrm{\Delta }(S(L_{}^{+}{}_{j}{}^{i}))=S(L_{}^{+}{}_{j}{}^{n})S(L_{}^{+}{}_{n}{}^{i})`$, it is easy to check that $`[\theta ^a,x_i]=0`$ for all $`i,a`$, and (3.38) follows. (3.40) follows immediately from $`L_{}^{+}{}_{a}{}^{i}S(L_{}^{+}{}_{j}{}^{a})=\delta _j^i`$, and To see (3.42), one needs the well–known relation $`L_{}^{+}{}_{r}{}^{l}L_{}^{+}{}_{s}{}^{k}g^{sr}=g^{kl}`$, as well as $`L_{}^{+}{}_{r}{}^{l}L_{}^{+}{}_{s}{}^{k}\epsilon _n^{sr}=\epsilon _m^{kl}L_{}^{+}{}_{n}{}^{m}`$; the latter follows from the quasitriangularity of $`U_q(su(2))`$. The commutation relations among the $`\theta `$ are obtained as in by observing $`\theta ^a\theta ^b`$ $`=`$ $`\mathrm{\Lambda }_N\theta ^aS(L_{}^{+}{}_{n}{}^{b})g^{nl}\xi _l`$ (3.43) $`=`$ $`\mathrm{\Lambda }_NS(L_{}^{+}{}_{n}{}^{b})\theta ^ag^{nl}\xi _l`$ $`=`$ $`\mathrm{\Lambda }_N^2S(L_{}^{+}{}_{n}{}^{b})S(L_{}^{+}{}_{j}{}^{a})g^{jk}g^{nl}\xi _k\xi _l,`$ using the commutation relations $`\widehat{R}_{ij}^{kl}SL_{}^{+}{}_{n}{}^{i}SL_{}^{+}{}_{m}{}^{j}=SL_{}^{+}{}_{i}{}^{k}SL_{}^{+}{}_{j}{}^{l}\widehat{R}_{nm}^{ij}`$, as well as (3.4). The remaining relations can be checked similarly. $``$$``$ ### 3.5 Integration of forms As classically, it is natural to define the integral over the forms of the highest degree, which is $`3`$ here. Since any $`\alpha ^{(3)}\mathrm{\Omega }_{q,N}^3`$ can be written in the from $`\alpha ^{(3)}=f\mathrm{\Theta }^3`$, we define $$\alpha ^{(3)}=f\mathrm{\Theta }^3:=\underset{𝒮_{q,N}^2}{}f$$ (3.44) by (2.42), so that $`\mathrm{\Theta }^3`$ is the volume form. This definition is natural, since $`[\mathrm{\Theta }^3,f]=0`$. Integrals of forms with degree $`3`$ will be set to zero. This integral satisfies an important cyclic property, as did the quantum trace (2.43). To formulate it, we extend the map $`S^2`$ from $`𝒮_{q,N}^2`$ to $`\mathrm{\Omega }_{q,N}^{}`$ by $$S^2(\xi _i)=q^H\xi _i,$$ extended as an algebra map. Then the following holds (see Appendix B): $$\alpha \beta =S^2(\beta )\alpha $$ (3.45) for any $`\alpha ,\beta \mathrm{\Omega }_{q,N}^{}`$ with $`deg(\alpha )+\mathrm{deg}(\beta )=3`$. Now Stokes theorem follows immediately: $$𝑑\alpha ^{(2)}=[\mathrm{\Theta },\alpha ^{(2)}]=0$$ (3.46) for any $`\alpha ^{(2)}\mathrm{\Omega }_{q,N}^2`$, because $`S^2\mathrm{\Theta }=\mathrm{\Theta }`$. This purely algebraic derivation is also valid on some other spaces . Finally we establish the compatibility of the integral with the star structure. From $`\mathrm{\Theta }^{}=\mathrm{\Theta }`$ and (2.45), we obtain $$(\alpha ^{(3)})^{}=(\alpha ^{(3)})^{}\text{for}q.$$ (3.47) For $`\|q\|=1`$, we have to extend the algebra map $`\overline{f}`$ (2.48) to $`\mathrm{\Omega }_{q,N}^{}`$. It turns out that the correct definition is $$\overline{\xi _i}=q^4g^{ij}\xi _j+q^2(qq^1)\frac{[2]_qC_N}{r^2}g^{ij}dx_j,$$ (3.48) extended as an antilinear algebra map; compare (3.28) for $`q`$. To verify that this is compatible with (3.1) and (3.12) requires the same calculations as to verify the star structure (3.28) for $`q`$. Moreover one can check using (3.10) that $$\overline{dx_i}=g^{ij}dx_j,$$ (3.49) which implies that $`\overline{\mathrm{\Theta }}=\mathrm{\Theta }`$, and $`\overline{_H(\alpha )}`$ $`=`$ $`_H(\overline{\alpha }),`$ $`\overline{d\alpha }`$ $`=`$ $`d\overline{\alpha },`$ $`\overline{\overline{\alpha }}`$ $`=`$ $`S^2\alpha `$ (3.50) for any $`\alpha \mathrm{\Omega }_{q,N}^{}`$. Hence we have $$(\alpha ^{(3)})^{}=\overline{\alpha ^{(3)}}\text{for}\|q\|=1.$$ (3.51) ## 4 Actions and fields ### 4.1 Scalar fields With the tools provided in the previous sections, it is possible to construct actions for 2–dimensional euclidean field theories on the $`q`$–deformed fuzzy sphere. We start with scalar fields, which are simply elements $`\psi 𝒮_{q,N}^2`$. The obvious choice for the kinetic term is $`S_{kin}[\psi ]`$ $`=`$ $`i{\displaystyle \frac{r^2}{\mathrm{\Lambda }_N^2}}{\displaystyle (d\psi )^{}_H𝑑\psi }\text{for}q,`$ $`S_{kin}[\psi ]`$ $`=`$ $`{\displaystyle \frac{r^2}{\mathrm{\Lambda }_N^2}}{\displaystyle \overline{d\psi }_H𝑑\psi }\text{for}\|q\|=1,`$ (4.1) which, using Stokes theorem, can equivalently be written in the form $`S_{kin}[\psi ]`$ $`=`$ $`i{\displaystyle \frac{r^2}{\mathrm{\Lambda }_N^2}}{\displaystyle }\psi ^{}(d_Hd)\psi ={\displaystyle \frac{r^2}{\mathrm{\Lambda }_N^2}}i{\displaystyle \underset{𝒮_{q,N}^2}{}}\psi ^{}(_Hd_Hd)\psi \text{for}q,`$ $`S_{kin}[\psi ]`$ $`=`$ $`{\displaystyle \frac{r^2}{\mathrm{\Lambda }_N^2}}{\displaystyle }\overline{\psi }(d_Hd)\psi ={\displaystyle \frac{r^2}{\mathrm{\Lambda }_N^2}}{\displaystyle \underset{𝒮_{q,N}^2}{}}\overline{\psi }(_Hd_Hd)\psi \text{for}\|q\|=1.`$ They are real $$S_{kin}[\psi ]^{}=S_{kin}[\psi ]$$ (4.3) for both $`q`$ and $`\|q\|=1`$, using the reality properties established in the previous sections. The fields can be expanded in terms of the irreducible representations $$\psi (x)=\underset{K,n}{}a^{K,n}\psi _{K,n}(x)$$ (4.4) according to (2.28), with coefficients $`a^{K,n}`$; this corresponds to the first–quantized case. However, in order to ensure invariance of the actions under $`U_q(su(2))`$ (or a suitable subset thereof), we must assume that $`U_q(su(2))`$ acts on products of fields via the $`q`$–deformed coproduct. This can be implemented consistently only after a “second quantization”, such that the coefficients in (4.4) generate a $`U_q(su(2))`$–module algebra. This will be presented in a forthcoming paper . One can also consider real fields, which have the form $`\psi (x)^{}`$ $`=`$ $`\psi (x)\text{for}q,`$ $`\overline{\psi (x)}`$ $`=`$ $`q^{H/2}\psi (x)q^{H/2}\text{for}\|q\|=1.`$ (4.5) This is preserved under the action of a certain real sector $`𝒢U_q(su(2))`$ (4.30); the discussion is completely parallel to the one below (4.32) in the next section, hence we will not give it here. Clearly $`_Hd_Hd`$ is the analog of the Laplace operator for functions, which can also be written in the usual form $`d\delta +\delta d`$, with $`\delta =_Hd_H`$. It is hermitian by construction. We wish to evaluate it on the irreducible representations $`\psi _K(2K+1)`$, that is, on spin–$`K`$ representations. The result is the following: ###### Lemma 4.1 If $`\psi _K𝒮_{q,N}^2`$ is a spin $`K`$ representation, then $$_Hd_Hd\psi _K=\frac{2}{[2]_qC_N}[K]_q[K+1]_q\psi _K.$$ (4.6) The proof is in Appendix B. It is useful to write down explicitly the hermitian forms associated to the above kinetic action. Consider $`S_{kin}[\psi ,\psi ^{}]`$ $`=`$ $`i{\displaystyle \frac{r^2}{\mathrm{\Lambda }_N^2}}{\displaystyle (d\psi )^{}_H𝑑\psi ^{}}\text{for}q,`$ $`S_{kin}[\psi ,\psi ^{}]`$ $`=`$ $`{\displaystyle \frac{r^2}{\mathrm{\Lambda }_N^2}}{\displaystyle \overline{d\psi }_H𝑑\psi ^{}}\text{for}\|q\|=1.`$ (4.7) Using Lemma 2.2, it follows immediately that they satisfy $`S_{kin}[\psi ,\psi ^{}]^{}`$ $`=`$ $`S_{kin}[\psi ^{},\psi ],`$ $`S_{kin}[\psi ,u\psi ^{}]`$ $`=`$ $`S_{kin}[u^{}\psi ,\psi ^{}]`$ (4.8) for both $`q`$ and $`\|q\|=1`$. To be explicit, let $`\psi _{K,n}`$ be an orthonormal basis of (2K+1). We can be assume that it is a weight basis, so that $`n`$ labels the weights from $`K`$ to $`K`$. Then it follows that $$S_{kin}[\psi _{K,n},\psi _{K^{},m}]=c_K\delta _{K,K^{}}\delta _{n,m}$$ (4.9) for some $`c_K`$. Clearly one can also consider interaction terms, which could be of the form $$S_{int}[\psi ]=\underset{𝒮_{q,N}^2}{}\psi \psi \psi ,$$ (4.10) or similarly with higher degree. ### 4.2 Gauge fields Gauge theories arise in a very natural way on $`𝒮_{q,N}^2`$. For simplicity, we consider only the analog of the abelian gauge fields here. They are simply one–forms $$B=B_a\theta ^ar\mathrm{\Omega }_{q,N}^1,$$ (4.11) which we expand in terms of the frames $`\theta ^a`$ introduced in Section 3.4. Notice that they have 3 independent components, which reflects the fact that calculus is 3–dimensional. Loosely speaking, the fuzzy sphere does see a shadow of the 3–dimensional embedding space. One of the components is essentially radial and should be considered as a scalar field, however it is naturally tied up with the other 2 components of $`B`$. We will impose the reality condition $`B^{}`$ $`=`$ $`B\text{for}q,`$ $`\overline{B}`$ $`=`$ $`q^{H/2}Bq^{H/2}\text{for}\|q\|=1.`$ (4.12) Since only 3–forms can be integrated, the most simple candidates for Langrangians that can be written down have the form $$S_3=\frac{1}{r^2\mathrm{\Lambda }_N^2}B^3,S_2=\frac{1}{r^2\mathrm{\Lambda }_N^2}B_HB,S_4=\frac{1}{r^2\mathrm{\Lambda }_N^2}B^2_HB^2.$$ (4.13) They are clearly real, with the reality condition (4.12); the factor $`i`$ for real $`q`$ is omitted here. We also define $$F:=B^2_HB,$$ (4.14) for reasons which will become clear below. The meaning of the field $`B`$ becomes obvious if one writes it in the form $$B=\mathrm{\Theta }+A,B_a=\frac{1}{r}\lambda _a+A_a$$ (4.15) While $`B`$ and $`\mathrm{\Theta }`$ become singular in the limit $`N\mathrm{}`$, $`A`$ remains well–defined. Using $`F`$ $`=`$ $`dA+A^2,`$ $`{\displaystyle A\mathrm{\Theta }^2}`$ $`=`$ $`{\displaystyle 𝑑A\mathrm{\Theta }}={\displaystyle }_HA\mathrm{\Theta },`$ $`{\displaystyle A^2\mathrm{\Theta }}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle (AdA+A_HA)}`$ (4.16) which follow from (3.19), one finds $`S_2`$ $`=`$ $`{\displaystyle \frac{1}{r^2\mathrm{\Lambda }_N^2}}{\displaystyle A_HA}+2A\mathrm{\Theta }^2`$ $`S_3`$ $`=`$ $`{\displaystyle \frac{1}{r^2\mathrm{\Lambda }_N^2}}{\displaystyle A^3}+{\displaystyle \frac{3}{2}}(AdA+A_HA)+3A\mathrm{\Theta }^2+\mathrm{\Theta }^3`$ (4.17) and $$S_{YM}:=\frac{1}{r^2\mathrm{\Lambda }_N^2}F_HF=\frac{1}{r^2\mathrm{\Lambda }_N^2}(dA+A^2)_H(dA+A^2).$$ (4.18) The latter action (which is a linear combination of $`S_2,S_3`$, and $`S_4`$) is clearly the analog of the Yang–Mills action, which in the classical limit contains a gauge field and a scalar, as we will see below. In the limit $`q1`$, it reduces to the action considered in . The actions $`S_3`$ and $`S_2`$ alone contain terms which are linear in $`A`$, which would indicate that the definition of $`A`$ (4.15) is not appropriate. However, the linear terms cancel in the following linear combination $$S_{CS}:=\frac{1}{3}S_3\frac{1}{2}S_2=\frac{2\pi }{3\mathrm{\Lambda }_N^2}+\frac{1}{2}\frac{1}{r^2\mathrm{\Lambda }_N^2}A𝑑A+\frac{2}{3}A^3.$$ (4.19) Notice that the “mass term” $`A_HA`$ has also disappeared. This form is clearly the analog of the Chern–Simons action. It is very remarkable that it exists on $`𝒮_{q,N}^2`$, which is related to the fact that the calculus is 3–dimensional. In the case $`q=1`$, this is precisely what has been found recently in the context of 2–branes on the $`SU(2)`$ WZW model . In terms of the components (4.11), $`B^2=B_aB_b\theta ^a\theta ^br^2`$, and $`_HB=\frac{r}{q^2[2]_{q^2}}B_a\epsilon _{bc}^a\theta ^c\theta ^b`$. Moreover, it is easy to check that $`_H(\theta ^b\theta ^c)`$ $`=`$ $`q^2\epsilon _a^{cb}\theta ^a,`$ $`\theta ^a_H\theta ^b`$ $`=`$ $`\mathrm{\Lambda }_N^2{\displaystyle \frac{q^4}{r^2}}g^{ba}\mathrm{\Theta }^3,`$ $`\theta ^a\theta ^b_H\theta ^c\theta ^d`$ $`=`$ $`[2]_{q^2}\mathrm{\Lambda }_N^2{\displaystyle \frac{q^8}{r^2}}(P^{})_{a^{}b^{}}^{dc}g^{b^{}b}g^{a^{}a}\mathrm{\Theta }^3=\mathrm{\Lambda }_N^2{\displaystyle \frac{q^8}{r^2}}\epsilon _n^{dc}\epsilon _m^{ba}g^{nm}\mathrm{\Theta }^3.`$ (4.20) Hence $`F`$ $`=`$ $`(B_aB_b+{\displaystyle \frac{1}{q^2r[2]_{q^2}}}B_c\epsilon _{ba}^c)\theta ^a\theta ^br^2=({\displaystyle \frac{\lambda _a}{r}}A_b+A_a{\displaystyle \frac{\lambda _b}{r}}+A_aA_b+{\displaystyle \frac{1}{q^2r[2]_{q^2}}}A_c\epsilon _{ba}^c)\theta ^a\theta ^br^2`$ (4.21) $`=`$ $`F_{ab}\theta ^a\theta ^br^2,`$ where we define $`F_{ab}`$ to be totally antisymmetric, i.e. $`F_{ab}=(P^{})_{ba}^{b^{}a^{}}F_{a^{}b^{}}`$ using (3.41). This yields $`S_{YM}`$ $`=`$ $`q^8[2]_{q^2}{\displaystyle \underset{𝒮_{q,N}^2}{}}F_{ab}F_{cd}(P^{})_{a^{}b^{}}^{dc}g^{b^{}b}g^{a^{}a},`$ To understand these actions better, we write the gauge fields in terms of “radial” and “tangential” components, $$A_a=\frac{x_a}{r}\varphi +A_a^t$$ (4.23) where $`\varphi `$ is defined such that $$x_aA_b^tg^{ab}=0;$$ (4.24) this is always possible. However to get a better insight, we consider the case $`q=1`$, and take the classical limit $`N\mathrm{}`$ in the following sense: for a given (smooth) field configuration in $`𝒮_N^2`$, we use the sequence of embeddings of $`𝒮_{q,N}^2`$ to approximate it for $`N\mathrm{}`$. Then terms of the form $`[A_a^t,A_b^t]`$ vanish in the limit (since the fields are smooth in the limit). The curvature then splits into a tangential and radial part, $`F_{ab}=F_{ab}^t+F_{ab}^\varphi `$, where<sup>7</sup><sup>7</sup>7the pull–back of $`F`$ to the 2–sphere in the classical case is unaffected by this split $`F_{ab}^t`$ $`=`$ $`{\displaystyle \frac{1}{2r}}\left([\lambda _a,A_b^t][\lambda _b,A_a^t]+A_c^t\epsilon _{ba}^c\right),`$ $`F_{ab}^\varphi `$ $`=`$ $`{\displaystyle \frac{1}{2r^2}}\left(\epsilon _{ab}^cx_c\varphi +[\lambda _a,\varphi ]x_b[\lambda _b,\varphi ]x_a\right).`$ (4.25) Moreover, $`x^aF_{ab}^t`$ $``$ $`{\displaystyle \frac{1}{4r}}[x^a\lambda _a,A_b^t]{\displaystyle \frac{1}{2r}}[\lambda _b,x^aA_a^t]=0,`$ $`x^a[\lambda _a,\varphi ]`$ $``$ $`{\displaystyle \frac{1}{2}}[x^a\lambda _a,\varphi ]=0`$ (4.26) in the classical limit, which implies that $`{\displaystyle \underset{S^2}{}}F_{ab}^tF_{}^{\varphi }{}_{}{}^{ab}`$ $`=`$ $`{\displaystyle \underset{S^2}{}}{\displaystyle \frac{1}{2r^2}}\epsilon _{ab}^nx_n\varphi F_{}^{t}{}_{}{}^{ab},`$ $`{\displaystyle \underset{S^2}{}}F_{ab}^\varphi F_{}^{\varphi }{}_{}{}^{ab}`$ $`=`$ $`{\displaystyle \underset{S^2}{}}{\displaystyle \frac{1}{2r^2}}\left(\varphi ^2+[\lambda _a,\varphi ][\lambda ^a,\varphi ]\right)`$ in the limit. Therefore we find $$S_{YM}=\underset{S^2}{}\left(2F_{ab}^tF_{}^{t}{}_{ab}{}^{}+\frac{2}{r^2}\epsilon _{ab}^nx_n\varphi F_{}^{t}{}_{}{}^{ab}+\frac{1}{r^2}(\varphi ^2+[\lambda _a,\varphi ][\lambda ^a,\varphi ])\right)$$ (4.27) in the limit, as in . Similarly, the Chern–Simons action (4.19 ) becomes $`S_{CS}`$ $``$ $`{\displaystyle \frac{2\pi }{3\mathrm{\Lambda }_N^2}}+{\displaystyle \frac{1}{2r^2\mathrm{\Lambda }_N^2}}{\displaystyle 𝑑A^t(A^t+2\mathrm{\Lambda }_N\mathrm{\Theta }\varphi )}\mathrm{\Lambda }_N^2\varphi ^2\mathrm{\Theta }^3`$ (4.28) $`=`$ $`{\displaystyle \frac{2\pi }{3\mathrm{\Lambda }_N^2}}+{\displaystyle \frac{1}{2r}}{\displaystyle \underset{S^2}{}}F_{ab}^t(A_c^t+2{\displaystyle \frac{x_c}{r}}\varphi )\epsilon ^{abc}{\displaystyle \frac{1}{2r^2}}{\displaystyle \underset{S^2}{}}\varphi ^2`$ for $`N\mathrm{}`$. In the flat limit $`r\mathrm{}`$, the term $`F_{ab}^tA_c^t\epsilon ^{abc}`$ vanishes because of (4.26), leaving the $`F\varphi `$ coupling term (after multiplying with $`r`$). Back to finite $`N`$ and $`q1`$. To further justify the above definition of curvature (4.14), we consider the zero curvature condition, $`F=0`$. In terms of the $`B`$ fields, this is equivalent to $$\epsilon _c^{ba}B_aB_b+\frac{1}{q^2r}B_c=0$$ (4.29) which is (up to rescaling) the same as equation (2.24) with opposite multiplication<sup>8</sup><sup>8</sup>8this can be implemented e.g. using the antipode of $`U_q(su(2))`$; in particular, the solutions $`B_a𝒮_{q,N}^2`$ are precisely all possible representations of $`U_q^{op}(su(2))`$ in the space of $`N+1`$–dimensional matrices. They are of course classified by the number of partitions of $`Mat(N+1)`$ into blocks with sizes $`n_1,\mathrm{},n_k`$ such that $`n_i=N+1`$, as in the case $`q=1`$. #### Gauge invariance. We have seen that actions which describe gauge theories in the limit $`q=1`$ arise very naturally on $`𝒮_{q,N}^2`$ (as on certain other higher–dimensional $`q`$–deformed spaces ). However, it is less obvious in which sense they are actually gauge–invariant for $`q1`$. For $`q=1`$, the appropriate gauge transformation is $`BUBU^1`$, for a unitary element $`U𝒮_N^2`$. This transformation does not work for $`q1`$, because of (3.45). Instead, we propose the following: let $``$ $`=`$ $`\{\gamma U_q(su(2)):\epsilon (\gamma )=0,\gamma ^{}=S\gamma \},`$ $`𝒢`$ $`=`$ $`\{\gamma U_q(su(2)):\epsilon (\gamma )=1,\gamma ^{}=S\gamma \}=e^{}`$ (4.30) for $`q`$, and $``$ $`=`$ $`\{\gamma U_q(su(2)):\epsilon (\gamma )=0,\gamma ^{}=S_0\gamma \},`$ $`𝒢`$ $`=`$ $`\{\gamma U_q(su(2)):\epsilon (\gamma )=1,\gamma ^{}=S_0\gamma \}=e^{}`$ (4.31) for $`\|q\|=1`$, where $`S_0(u)=q^{H/2}S(u)q^{H/2}`$. Clearly $``$ is a subalgebra (without 1) of $`U_q(su(2))`$, and $`𝒢`$ is closed under multiplication. Using the algebra map $`j`$ (2.31), $`𝒢`$ can be mapped to some real sector of the space of functions on the fuzzy sphere. Now consider the following “gauge” transformations: $$Bj(\gamma _{(1)})Bj(S\gamma _{(2)})\text{for}\gamma 𝒢.$$ (4.32) It can be checked easily that these transformations preserve the reality conditions (4.12) for both real $`q`$ and $`\|q\|=1`$. In terms of components $`B=B_a\theta ^ar`$, this transformation is simply (suppressing $`j`$) $$B_a\gamma _{(1)}B_aS\gamma _{(2)}=\gamma B_a,$$ (4.33) which is the rotation of the fields $`B_a𝒮_{q,N}^2`$ considered as scalar fields<sup>9</sup><sup>9</sup>9notice that this is not the rotation of the one–form $`B`$, because $`\gamma _{(1)}\xi _iS\gamma _{(2)}\gamma \xi _i`$, i.e. the rotation $`\gamma U_q(su(2))`$ does not affect the index $`a`$ because of (3.38). In terms of the $`A_a`$ variables, this becomes $$A_a\theta ^a\gamma _{(1)}A_aS\gamma _{(2)}\theta ^a+\gamma _{(1)}d(S\gamma _{(2)})=(\gamma A_a)\theta ^a+\gamma _{(1)}d(S\gamma _{(2)}),$$ (4.34) using (3.6) and (3.38). Hence these transformations are a mixture of rotations of the components (first term) and “pure gauge transformations” (second term). Moreover, the radial and tangential components get mixed. To understand these transformations better, consider $`q=1`$. Then we have two transformations of a given gauge field $`B_a`$, the first by conjugation with an unitary element $`U𝒮_{q,N}^2`$, and the second by (4.32). We claim that the respective spaces of inequivalent gauge fields are in fact equivalent. Indeed, choose e.g. $`a=1`$; then there exists a unitary $`U𝒮_{q,N}^2`$ such that $`U^1B_aU`$ is a diagonal matrix with real entries. On the other hand, using a suitable $`\gamma 𝒢`$ and recalling (2.28), one can transform $`B_a`$ into the form $`B_a=_ib_i(x_0)^i`$ with real $`b_i`$, which is again represented by a diagonal matrix in a suitable basis. Hence at least generically, the spaces of inequivalent gauge fields are equivalent. One can also see more intuitively that (4.34) corresponds to an abelian gauge transformation in the classical limit. Consider again $`\gamma (x)=e^{ih(x)}`$ with $`h(x)^{}=Sh(x)`$, approximating a smooth function in the limit $`N\mathrm{}`$. Using properly rescaled variables $`x_i`$, one can see using (2.30) that if viewed as an element in $`U(su(2))`$, $`\gamma `$ approaches the identity, that is $`\gamma A_a(x)A_a(x)`$ in the classical limit. Now write the functions on $`𝒮_N^2`$ in terms of the variables $`x_1`$ and $`x_1`$, for example. Then (2.30) yields $$(\mathrm{𝟏}S)\mathrm{\Delta }(x_i)=x_i11x_i,$$ (4.35) for $`i=\pm 1`$, and one can see that $$\gamma _{(1)}[\lambda _i,S\gamma _{(2)}]_ih(x_i)$$ (4.36) in the (flat) classical limit. Hence (4.34) indeed becomes a gauge transformation in the classical limit. To summarize, we found that the set of gauge transformations in the noncommutative case is a (real sector of a) quotient of $`U_q(su(2))`$, and can be identified with the space of (real) functions on $`𝒮_{q,N}^2`$ using the map $`j`$. However, the transformation of products of fields is different from the classical case. Classically, the gauge group acts on products “componentwise”, which means that the coproduct is trivial. Here, we must assume that $`U_q(su(2))`$ acts on products of fields via the $`q`$–deformed coproduct, so that the above actions are invariant under gauge transformations, by (2.44). In particular, the “gauge group” has become a real sector of a Hopf algebra. Of course, this can be properly implemented on the fields only after a “second quantization”, as in the case of rotation invariance (see Section 4.1). This will be presented in a forthcoming paper . This picture is also quite consistent with observations of a BRST–like structure in $`U_q(so(2,3))`$ at roots of unity, see . Finally, we point out that the above actions are invariant under a global $`U_q(su(2))`$ symmetry, by rotating the frame $`\theta ^a`$. ## 5 Drinfel’d twists and the relation with $`D`$–branes Finally we relate our $`q`$–deformed fuzzy sphere to the effective algebra of functions on spherical $`D`$–branes in the $`SU(2)`$ WZW model at level $`k`$, as determined by Alekseev, Recknagel and Schomerus . Their result is as follows. The $`D`$–branes (more precisely their boundary conditions) are classified by an integer $`N`$ which satisfies<sup>10</sup><sup>10</sup>10$`N`$ is denoted by $`2\alpha `$ in $`0Nk`$. The Hilbert space of the associated boundary conformal field theory decomposes into irreducible representations of the affine Lie algebra $`\widehat{su(2)}_k`$. One can assign abstract elements $`\{Y_i^I\}_{I,i}`$ to the boundary vertex operators (primary fields), where $`I`$ ranges from $`0`$ to $`\mathrm{min}(N,kN)`$, and $`IiI`$. The $`\{Y_i^I\}_i`$ form irreducible spin $`I`$ representations of the horizontal algebra $`su(2)`$, and are interpreted as the analog of spherical harmonics on the $`D`$–brane; in particular, there exist only finitely many of them. The algebra induced by the OPE of the corresponding boundary vertex operators is given by $$Y_i^IY_j^J=\underset{K,k}{}\left[\begin{array}{ccc}I& J& K\\ i& j& k\end{array}\right]\left\{\begin{array}{ccc}I& J& K\\ N/2& N/2& N/2\end{array}\right\}_qY_k^K\text{with}q=e^{\frac{i\pi }{k+2}},$$ (5.1) where the sum goes from $`K=0`$ to $`\mathrm{min}(I+J,kIJ,N,kN)`$. This is a finite, noncommutative, quasiassociative algebra $`𝒜`$. Here the first bracket denotes the Clebsch–Gordon coefficients of $`SU(2)`$, and the curly brackets denote the $`q`$–deformed $`6J`$–symbols of $`U_q(su(2))`$. The latter arise from the fusion matrices of the underlying conformal field theory, which have been known to be related to quantum groups for a long time . In the present paper, we only consider roots of unity $`q`$ which satisfy (2.36). This means that $`Nk/2`$ in the above situation, so that we can only consider a certain subset of the allowed boundary conditions here. There will be some qualitative changes in the remaining cases, which we do not consider in the present paper. The reason for the non–associativity of the algebra $`𝒜`$ is a mixing of $`q`$–deformed and undeformed group theory objects. However as was already indicated in , one can sometimes “twist” this algebra using a Drinfeld–twist into an associative one. In particular this can be done if $`q`$ satisfies (2.36), in the following way: On the same vector space $`𝒜`$, we define a new multiplication by $$a\stackrel{~}{}b:=(^{1(1)}a)(^{1(2)}b)=(^1(ab)).$$ (5.2) Here $`a,b𝒜`$, and $$=^{(1)}^{(2)}U(g)U(g)$$ (5.3) is the Drinfeld twist in Sweedler–notation. We can ignore some fine points here since we only consider certain representations of $``$. The twist relates the undeformed Clebsch–Gordan coefficients to the deformed ones as follows: $$\left[\begin{array}{ccc}I& J& K\\ i& j& k^{}\end{array}\right](g^{(K)})^{k^{}k}=\left[\begin{array}{ccc}I& J& K\\ i^{}& j^{}& k^{}\end{array}\right]_q(g_q^{(K)})^{k^{}k}\pi _i^i^{}(^{(1)})\pi _j^j^{}(^{(2)})$$ (5.4) Here we have raised indices using $`(g_q^{(K)})^{k^{}k}`$, which is the $`q`$–deformed invariant tensor, and we will assume that $`(g^{(K)})^{k^{}k}=\delta ^{k^{}k}`$ (in an orthonormal basis). It should be noted that even though the abstract element $``$ exists only for generic (more precisely formal) $`q`$, the representations of $``$ which are needed above do exist at roots of unity, assuming the restrictions (2.36) on $`q`$; this is because the Clebsch–Gordan decomposition is then still analytic in $`q`$. Hence the twisted multiplication rule for the generators $`Y_i^I`$ is $$Y_i^I\stackrel{~}{}Y_j^J=\underset{K,k}{}\left[\begin{array}{ccc}I& J& K\\ i& j& k^{}\end{array}\right]_q(g_q^{(K)})^{k^{}k}\left\{\begin{array}{ccc}I& J& K\\ N/2& N/2& N/2\end{array}\right\}_qY_k^K.$$ (5.5) As was already pointed out in , this defines an associative algebra. We claim that this is precisely the algebra $`𝒮_{q,N}^2`$, which in turn is the matrix algebra $`Mat(N+1)`$. To see this, we reconsider the algebra $`𝒮_{q,N}^2`$ from a group–theoretic point of view: Let now $`Y_i^I𝒮_{q,N}^2`$ be an irreducible spin $`I`$ representation of $`U_q(su(2))`$, according to the decomposition (2.28). In acts on the Fock space $`_N`$ (2.16), which in turn is a spin $`N/2`$ representation of $`U_q(su(2))`$, with a basis of the form $`(A^+\mathrm{}A^+)_r\|0`$. Hence if we denote with $`\pi (Y_i^I)_s^r`$ the matrix which represents the operator $`Y_i^I`$ on $`_N`$, we can conclude that it is proportional to the Clebsch–Gordan coefficient of the decomposition $`(2I+1)(N+1)(N+1)`$; this is the Wigner–Eckart theorem. Hence in a suitable normalization of the basis $`Y_i^I`$, we can write $$\pi (Y_i^I)_s^r=(g_q^{(N/2)})^{rr^{}}\left[\begin{array}{ccc}N/2& I& N/2\\ r^{}& i& s\end{array}\right]_q=\left[\begin{array}{ccc}I& N/2& N/2\\ i& s& r^{}\end{array}\right]_q(g_q^{(N/2)})^{r^{}r}.$$ (5.6) Therefore the matrix representing the operator $`Y_i^IY_j^J`$ is given by $`\pi (Y_i^I)_s^r`$ $`\pi (Y_j^J)_t^s=(g_q^{(N/2)})^{rr^{}}\left[\begin{array}{ccc}N/2& I& N/2\\ r^{}& i& s\end{array}\right]_q(g_q^{(N/2)})^{ss^{}}\left[\begin{array}{ccc}N/2& J& N/2\\ s^{}& j& t\end{array}\right]_q`$ (5.11) $`=`$ $`{\displaystyle \underset{K}{}}\left\{\begin{array}{ccc}N/2& J& N/2\\ I& N/2& K\end{array}\right\}_q\left[\begin{array}{ccc}I& J& K\\ i& j& k^{}\end{array}\right]_q(g_q^{(K)})^{k^{}k}\left[\begin{array}{ccc}K& N/2& N/2\\ k& t& r^{}\end{array}\right]_q(g_q^{(N/2)})^{r^{}r}`$ (5.18) $`=`$ $`{\displaystyle \underset{K}{}}\left\{\begin{array}{ccc}I& J& K\\ N/2& N/2& N/2\end{array}\right\}_q\left[\begin{array}{ccc}I& J& K\\ i& j& k^{}\end{array}\right]_q(g_q^{(K)})^{k^{}k}\pi (Y_k^K)_t^r.`$ (5.23) Here we used the identity $$\left\{\begin{array}{ccc}N/2& J& N/2\\ I& N/2& K\end{array}\right\}_q=\left\{\begin{array}{ccc}I& J& K\\ N/2& N/2& N/2\end{array}\right\}_q,$$ (5.25) which is proved in . This calculation is represented graphically in Figure 1, which shows that it essentially reduces to the definition of the $`6j`$–symbols. Therefore the algebra of $`𝒮_{q,N}^2`$ is precisely (5.5), which is a twist of the algebra (5.1) found in . In a sense, this twisting is similar to deformation quantization; however, $`𝒮_{q,N}^2`$ is a $`U_q(su(2))`$–module algebra, while (5.1) is a $`U(su(2))`$–module algebra. ## 6 Summary and outlook In this paper, we have studied the $`q`$–deformed fuzzy sphere $`𝒮_{q,N}^2`$, which is an associative algebra which is covariant under $`U_q(su(2))`$, for real $`q`$ and $`q`$ a phase. In the first case, this is the same as the “discrete series” of Podleś quantum spheres. We develop the formalism of differential forms and frames, as well as integration. We then briefly consider scalar and gauge field theory on this space. It appears that $`𝒮_{q,N}^2`$ is a nice and perhaps the simplest example of quantum spaces which are covariant under a quantum group. This makes it particularly well suited for studying field theory, an endeavour which has proved to be rather difficult on other $`q`$–deformed spaces. We are able to write hermitian actions for scalar and gauge fields, including analogs of Yang–Mills and Chern–Simons actions. In particular, the form of the actions for gauge theories suggests a new type of gauge symmetry, where the role of the gauge group is played by $`U_q(su(2))`$, which can be mapped onto the space of functions on $`𝒮_{q,N}^2`$. This suggests that formulating field theory on quantized spaces which are less trivial than the ones corresponding to a Moyal product on flat spaces requires new approaches, and may lead to interesting new insights. The main motivation for doing this is the discovery that a quasi–associative twist of $`𝒮_{q,N}^2`$ arises on spherical $`D`$–branes in the $`SU(2)_k`$ WZW model, for $`q`$ a root of unity. In view of this result, we hope that the present formalism may be useful to formulate a low–energy effective field theory induced by open strings ending on the $`D`$–branes. This in turn inspires one to consider some kind of second quantization of field theories on $`𝒮_{q,N}^2`$, corresponding to a loop expansion and many–particle states. It is quite interesting that also from a more formal point of view, such a second quantization turns out to be necessary for a satisfactory definition of symmetries in such a field theory. This will be presented in a future publication . Moreover, while the question of using either the quasi–associative algebra (5.1) or the associative $`𝒮_{q,N}^2`$ may ultimately be a matter of taste, the latter does suggest certain forms for Lagrangians, induced by the differential calculus. It would be very interesting to compare this with a low–energy effective action induced from string theory. In this paper we have only considered spaces which correspond to a subset of the allowed boundary conditions discussed in . The remaining cases will show some qualitatively new features, and are postponed for future work. #### Acknowledgements: We would like to thank A. Alekseev, C.-S. Chu, B. Cerchiai, G. Fiore, J. Pawelczyk, E. Scheidegger and J. Wess for useful and stimulating discussions. H. S. and J. M. are grateful for hospitality at the Erwin–Schrödinger Institut in Vienna, where this work was initiated. H. G. and J. M. would also like to thank J. Wess and the Max–Planck Institut für Physik in München for hospitality and financial support. ## 7 Appendix A: Invariant tensors Before giving the explicit forms of the invariant tensors used in this paper, we briefly explain our conventions and the relation to the literature. The quantum spaces in and in much of the standard literature are defined as left $`Fun_q(G)`$–comodule algebras. This is equivalent to right $`U_q(g)`$–module algebras. However it is more intuitive to work with left module algebras. This can be achieved using the antipode, $`uf(x):=f(x)S(u)`$; moreover if $`\pi _j^i(u)`$ is the fundamental representation, then $`ux_i=x_j\pi _i^j(u)`$ where $`x_i=g_{ij}x^j`$. However the coproduct then becomes reversed, $`ufg=(u_2f)(u_1g)`$. We have incorporated this by defining $`U_q(g)`$ with the reversed coproduct (2.9) and antipode. This means that our $`U_q^{}U_q^+`$ (where $`U_q^\pm `$ denotes the Borel subalgebras) is obtained from the usual one by flipping the tensor components. For example, our $`\widehat{R}_+^+=\pi _+^{}(_1)\pi _{}^+(_2)=(qq^1)`$ in the fundamental representation of $`U_q(su(2))`$, where $`\pm `$ labels the weights. Then the characteristic equation and all compatibility relations with the invariant tensors have the same form as usual, and are obtained from the standard ones by flipping all horizontal indices. The $`q`$–deformed epsilon–symbol (“spinor metric”) for spin 1/2 representations is given by $$\epsilon ^+=q^{\frac{1}{2}},\epsilon ^+=q^{\frac{1}{2}},$$ (7.1) all other components are zero. The corresponding tensor with lowered indices is $`\epsilon _{\alpha \beta }=\epsilon ^{\alpha \beta }`$ and satisfies $`\epsilon ^{\alpha \beta }\epsilon _{\beta \gamma }=\delta _\gamma ^\alpha `$. In particular, $`\epsilon ^{\alpha \beta }\epsilon _{\alpha \beta }=(q+q^1)=[2]_q`$. The $`q`$–deformed sigma–matrices, i.e. the Clebsch–Gordon coefficients for $`(3)(2)(2)`$, are given by $`\sigma _1^{++}`$ $`=`$ $`1=\sigma _1^{},`$ $`\sigma _0^+`$ $`=`$ $`{\displaystyle \frac{q^{\frac{1}{2}}}{\sqrt{[2]_q}}},\sigma _0^+={\displaystyle \frac{q^{\frac{1}{2}}}{\sqrt{[2]_q}}}`$ (7.2) in an orthonormal basis, and $`\sigma _i^{\alpha \beta }=\sigma _{\alpha \beta }^i`$. They are normalized such that $`_{\alpha \beta }\sigma _{\alpha \beta }^i\sigma _j^{\alpha \beta }=\delta _j^i`$. That is, they define a unitary map (at least for $`q`$). The $`q`$–deformed invariant tensor for spin 1 representations is given by $$g^{11}=q^1,g^{00}=1,g^{11}=q,$$ (7.3) all other components are zero. Then $`g_{\alpha \beta }=g^{\alpha \beta }`$ satisfies $`g^{\alpha \beta }g_{\beta \gamma }=\delta _\gamma ^\alpha `$, and $`g^{\alpha \beta }g_{\alpha \beta }=q^2+1+q^2=[3]_q`$. The Clebsch–Gordon coefficients for $`(3)(3)(3)`$, i.e. the $`q`$–deformed structure constants, are given by $$\begin{array}{cc}\epsilon _1^{10}=q^1,\hfill & \epsilon _1^{01}=q,\hfill \\ \epsilon _0^{00}=(qq^1),\hfill & \epsilon _0^{11}=1=\epsilon _0^{11},\hfill \\ \epsilon _1^{01}=q^1,\hfill & \epsilon _1^{10}=q\hfill \end{array}$$ (7.4) in an orthonormal basis, and $`\epsilon _{ij}^k=\epsilon _k^{ij}`$. They are normalized such that $`_{ij}\epsilon _{ij}^n\epsilon _m^{ij}=[2]_{q^2}\delta _m^n`$. Moreover, the following identities hold: $`\epsilon _{ij}^ng^{jk}`$ $`=`$ $`\epsilon _i^{nk}`$ (7.5) $`g_{ij}\epsilon _{kl}^j`$ $`=`$ $`\epsilon _{ik}^jg_{jl}`$ (7.6) $`\epsilon _i^{nk}\epsilon _k^{lm}\epsilon _i^{km}\epsilon _k^{nl}`$ $`=`$ $`g^{nl}\delta _i^m\delta _i^ng^{lm}`$ (7.7) which can be checked explicitly. In view of (7.6), the $`q`$–deformed totally ($`q`$–)antisymmetric tensor is defined as follws: $$\epsilon ^{ijk}=g^{in}\epsilon _n^{jk}=\epsilon _n^{ij}g^{nk}.$$ (7.8) It is invariant under the action of $`U_q(su(2))`$. ## 8 Appendix B: some proofs #### Proof of (3.7): Using the identity $$\mathrm{𝟏}=q^2\widehat{R}+(1+q^4)P^{}+(1q^6)P^0,$$ (8.1) (7.5), (3.4), and the braiding relation (3.1) we can calculate the commutation relation of $`\mathrm{\Theta }`$ with the generators $`x_i`$: $`x_i\mathrm{\Theta }`$ $`=`$ $`x_i(x_j\xi _tg^{jt})`$ $`=`$ $`q^2\widehat{R}_{ij}^{kl}x_kx_l\xi _tg^{jt}+q^2\mathrm{\Lambda }_N\epsilon _{ij}^nx_n\xi _tg^{jt}+{\displaystyle \frac{r^2}{[3]}}(1q^6)g_{ij}\xi _tg^{jt}`$ $`=`$ $`q^2\mathrm{\Theta }x_i+q^2\mathrm{\Lambda }_N\epsilon _{ij}^nx_n\xi _tg^{jt}+r^2{\displaystyle \frac{(1q^6)}{[3]_q}}\xi _i`$ $`=`$ $`q^2\mathrm{\Theta }x_i+q^2\mathrm{\Lambda }_N\epsilon _i^{nk}x_n\xi _k+r^2q^3(qq^1)\xi _i,`$ which yields (3.7). #### Proof of (3.13) and (3.14): Using (3.12), one has $$\mathrm{\Theta }\xi _i=q^2\xi _i\mathrm{\Theta },$$ (8.2) which implies $`\mathrm{\Theta }\mathrm{\Theta }=\mathrm{\Theta }(x\xi )=dx\xi q^2\mathrm{\Theta }\mathrm{\Theta }`$, hence $$(1+q^2)\mathrm{\Theta }^2=dx\xi .$$ On the other hand, (3.7) yields $$dx\xi =\mathrm{\Lambda }_Nx_i\epsilon _j^{kl}\xi _k\xi _lg^{ij}q^3(qq^1)\mathrm{\Theta }^2,$$ and combining this it follows that $$\mathrm{\Theta }^2=\frac{q^2\mathrm{\Lambda }_N}{[2]_{q^2}}x_i\xi _k\xi _l\epsilon ^{ikl}.$$ We wish to relate this to $`dx_idx_jg^{ij}`$, which is proportional to $`d\mathrm{\Theta }`$. Using the relations $`\epsilon _i^{nk}x_n\xi _k=q^2\epsilon _i^{nk}\xi _nx_k`$, $`\mathrm{\Theta }=q^4\xi x`$, (7.6) and (7.7), one can show that $$\epsilon _i^{nk}x_n\xi _k\epsilon _j^{ml}x_m\xi _lg^{ij}=\mathrm{\Lambda }_Nx_i\xi _k\xi _l\epsilon ^{ikl}+q^2\mathrm{\Theta }^2$$ which using (3.7) implies $$dx_idx_jg^{ij}=\frac{1}{C_N}r^2\mathrm{\Theta }^2.$$ #### Proof of $`dd=0`$ on $`\mathrm{\Omega }_{q,N}^1`$. First, we calculate $`[\mathrm{\Theta },d\xi _i]`$ $`=`$ $`(q^21)(q^2+1)\left(\xi _i\mathrm{\Theta }^2+{\displaystyle \frac{q^2\mathrm{\Lambda }_N}{[2]_{q^2}}}\epsilon _i^{kl}\xi _k\xi _l\mathrm{\Theta }\right)`$ $`=`$ $`(q^21)(q^2+1)(\xi _i(_H\mathrm{\Theta })(_H\xi _i)\mathrm{\Theta })=0`$ using (8.2), (3.17), and (3.16). This implies that $`d(d(f\xi _i))`$ $`=`$ $`[\mathrm{\Theta },df\xi _i+fd\xi _i]`$ $`=`$ $`[\mathrm{\Theta },df]_+\xi _idf[\mathrm{\Theta },\xi _i]_++dfd\xi _i+f[\mathrm{\Theta },d\xi _i]`$ $`=`$ $`ddf+_H(df)\xi _idf(d\xi _i+_H(\xi _i))+dfd\xi _i`$ $`=`$ $`df_H\xi _i+(_Hdf)\xi _i=0`$ by (3.16) for any $`f𝒮_{q,N}^2`$. This proves $`dd=0`$ on $`\mathrm{\Omega }_{q,N}^1`$. #### Proof of (3.16). First, we show that $$(_H\xi _i)\xi _j=\xi _i(_H\xi _j),$$ (8.3) which is equivalent to $$\epsilon _i^{nk}\xi _n\xi _k\xi _j=\xi _i\epsilon _j^{nk}\xi _n\xi _k.$$ Now $`\mathrm{\Omega }_{q,N}^3`$ is one–dimensional as module over $`𝒮_{q,N}^2`$, generated by $`\mathrm{\Theta }^3`$ (3.22), which in particular is a singlet under $`U_q(su(2))`$. This implies that $`\epsilon _i^{nk}\xi _n\xi _k\xi _j`$ $`=`$ $`(P^0)_{ij}^{rs}\epsilon _r^{nk}\xi _n\xi _k\xi _s`$ (8.4) $`=`$ $`{\displaystyle \frac{q^6[2]_{q^2}}{\mathrm{\Lambda }_Nr^2}}g_{ij}\mathrm{\Theta }^3`$ $`=`$ $`(P^0)_{ij}^{rs}\xi _r\epsilon _s^{nk}\xi _n\xi _k=\xi _i\epsilon _j^{nk}\xi _n\xi _k,`$ as claimed. Now (3.16) follows immediately using the fact that $`_H`$ is a left–and right $`𝒮_{q,N}^2`$–module map. #### Reality structure for $`q`$: These are the most difficult calculations, and they are needed to verify (3.48) as well. First, we have to show that (3.1) is compatible with the star structure (3.28). By a straightforward calculation, one can reduce the problem to proving (3.29). We verify this by projecting this quadratic equation to its spin 0, spin 1, and spin 2 part. The first two are easy to check, using (3.10) in the spin 1 case. To show the spin 2 sector, it is enough to consider (3.29) for $`i=j=1`$, by covariance. This can be seen e.g. using $`[x_1,\epsilon _1^{ij}x_i\xi _j]=q^2\mathrm{\Lambda }_Nx_1\xi _1`$, which in turn can be checked using (8.1), (3.3) and (7.4). Next, we show that (3.12) is compatible with the star structure (3.28). This can be reduced to $$(q^2\widehat{R}q^2\widehat{R}^1)_{ij}^{kl}dx_k\xi _l=q^2(qq^1)\frac{[2]_qC_N}{r^2}(\mathrm{𝟏}+q^2\widehat{R})_{ij}^{kl})dx_kdx_l$$ The spin 0 part is again easy to verify, and the spin 1 part vanishes identically (since then $`\widehat{R}`$ has eigenvalue $`q^2`$). For the spin 2 part, one can again choose $`i=j=1`$, and verify it e.g. by comparing with the differential of equation (3.29). #### Proof of (3.45): Since $`\mathrm{\Omega }_{q,N}^{}`$ is finitely generated and because of (2.43) and $`[\mathrm{\Theta }^3,f]=0`$, it is enough to consider $`\beta =\xi _k`$. In this case, the claim reduces to $$\xi _i\xi _j\xi _k=S^2(\xi _k)\xi _i\xi _j.$$ Now $`S^2(\xi _k)=D_k^l\xi _l`$, where $`D_k^l=\delta _k^lq^{2r_l}`$ with $`r_l=(2,0,2)`$ for $`l=(1,0,1)`$, respectively. Since $`\xi _i\xi _j=\frac{1}{[2]_{q^2}}\epsilon _{ij}^n(\epsilon _n^{rs}\xi _r\xi _s)`$, there remains to show that $`(\epsilon _n^{rs}\xi _r\xi _s)\xi _k=S^2(\xi _k)(\epsilon _n^{rs}\xi _r\xi _s)`$. By (8.4), this is equivalent to $$g_{nk}\mathrm{\Theta }^3=D_k^lg_{ln}\mathrm{\Theta }^3,$$ which follows from the definition of $`D_k^l`$. #### Proof of Lemma 4.6: Using (3.18), (3.23) and $`d\mathrm{\Theta }^2=0`$, we have $`d_Hd\psi `$ $`=`$ $`d(\psi \mathrm{\Theta }^2\mathrm{\Theta }^2\psi )=(d\psi )\mathrm{\Theta }^2\mathrm{\Theta }^2d\psi `$ (8.5) $`=`$ $`(d\psi )\mathrm{\Theta }^2+[\mathrm{\Theta }^2,\psi ]\mathrm{\Theta }`$ $`=`$ $`(d\psi )\mathrm{\Theta }^2+(_Hd\psi )\mathrm{\Theta }.`$ To proceed, we need to evaluate $`d\psi _K`$. Because it is an irreducible representation, it is enough to consider $`\psi _K=(x_1)^K`$. From (3.29) and using $`\xi _1x_1=q^2x_1\xi _1`$, it follows that $$dx_1x_1=q^2x_1dx_1\frac{q^2}{C_N}r^2x_1\xi _1,$$ since $`\widehat{R}`$ can be replaced by $`q^2`$ here. By induction, one finds $$dx_1x_1^k=x_1^k\left(q^{2k}dx_1[k]_{q^2}\frac{q^2}{C_N}r^2\xi _1\right),$$ (8.6) and by an elementary calculation it follows that $$d(x_1^{k+1})=[k+1]_qx_1^k\left(q^kdx_1\frac{q^2}{[2]_qC_N}[k]_qr^2\xi _1\right).$$ (8.7) Moreover, we note that using (3.16) $$(\xi _i\mathrm{\Theta }+_H\xi _i)\mathrm{\Theta }=\xi _i(_H\mathrm{\Theta })+(_H\xi _i)\mathrm{\Theta }=2(_H\xi _i)\mathrm{\Theta }=\frac{2}{r^2}x_i\mathrm{\Theta }^3.$$ (8.8) The last equality follows easily from (8.4) and (8.2). Similarly $$(dx_i\mathrm{\Theta }+_Hdx_i)\mathrm{\Theta }=2_Hdx_i\mathrm{\Theta }=2dx_i\mathrm{\Theta }^2.$$ (8.9) Now we can continue (8.5) as $`d_Hdx_1^K`$ $`=`$ $`(dx_1^{K1}\mathrm{\Theta }+_Hdx_1^{K1})\mathrm{\Theta }`$ (8.10) $`=`$ $`[K]_qx_1^{K1}\left(2q^{K1}dx_1\mathrm{\Theta }^22{\displaystyle \frac{q^2}{[2]_qC_N}}[K1]_qx_1\mathrm{\Theta }^3\right).`$ Finally it is easy to check that $$dx_i\mathrm{\Theta }^2=\frac{1}{C_N}x_i\mathrm{\Theta }^3,$$ (8.11) and after a short calculation one finds (4.6).
warning/0005/hep-th0005028.html	ar5iv	text	# Hagedorn inflation of D-branes ## 1 Introduction Inflation is a beautiful solution to several difficult problems in cosmology; the horizon problem, the flatness problem, and the monopole problem. In field theory however, the standard way to obtain inflation is to add a positive cosmological constant (which has a negative pressure $`p=\rho )`$. This ingredient is, without doubt, the least attractive feature of standard inflation and it is generally extremely difficult to control its adverse effects (e.g. the graceful exit and moduli problems). It is worth asking therefore if there are other forms of matter that can have a negative pressure and hence give an accelerating scale factor. In this paper we introduce a candidate that actually has this property – open strings on D-branes at temperatures close to the string scale. At sufficiently high temperatures and densities fundamental strings enter a curious ‘long string’ Hagedorn phase . To date applications of this phase have been quite limited in string cosmology because the thermodynamics is governed by the finite temperature partition function. A rigorous analysis therefore requires nothing less than solving the string system in a cosmological setting, a difficult problem that might at best be tractable only in a few special cases. Moreover, in order to understand the effect of macroscopic phenomena such as winding modes we need the microcanonical ensemble (as we shall discuss) – an ensemble that does not particularly lend itself to cosmological applications. In this paper we show that, for certain systems, it is possible to bypass these technical difficulties by using a classical random walk picture to model the behaviour of the strings in cosmological backgrounds. The particular systems we will focus on are D-branes in the weak coupling limit . In this limit one can separate the energy momentum tensor into two components; a localized component corresponding to the D-brane tension, and a diffuse component that spreads into the bulk corresponding to open string excitations of the brane. (At the risk of causing confusion we will often refer to the latter as a ‘bulk’ component.) We will in addition allow a bulk cosmological constant, although our focus in this paper will be on the cosmology when the combined effect of the brane and bulk cosmological constants is subdominant. The random walk picture allows us a first glimpse at the cosmological effects of a primordial Hagedorn phase of open strings on branes and we find two interesting types of behaviour. * The first we call Hagedorn inflation. We will show that the transverse ‘bulk’ components of a D-brane’s energy-momentum tensor can be negative. If all of the transverse dimensions have winding modes, this negative ‘pressure’ causes the brane to power law inflate along its length with a scale factor that varies as $`at^{4/3}`$ even in the absence of a nett cosmological constant (as shown in eq.(72)). If there are transverse dimensions that are large (in the sense that the string modes are not space-filling in these directions), then we can find exponential inflation (as shown in eq.(73)). * If there is a small but negative cosmological constant, the Universe can enter a stable but oscillating phase; i.e. a ‘bouncing’ universe. The nett effect of the Hagedorn phase is to soften the singular behaviour associated with the collapse. Such singularity smoothing is a familiar aspect of strings, but the nice feature here is that we find it in a purely perturbative regime. We should at this point also emphasize a general observation that we make, namely that the diffuse stringy component can have a dominant effect on the cosmology even in the weak coupling limit. At first sight this may be somewhat surprising given that the intrinsic tension energy of a D-brane goes like $`\rho _{br}1/g_s1/\widehat{\kappa }`$ where $`g_s`$ is the string coupling and $`\widehat{\kappa }`$ is the effective gravitational coupling. However, we will see that the cosmological effects of the two components are proportional to $`\widehat{\kappa }^4\rho _{br}^2`$ and $`\widehat{\kappa }^2\rho `$ for the brane and diffuse ‘bulk’ components respectively. Then the contribution of the brane component is proportional to $`\widehat{\kappa }^4\rho _{br}^2\widehat{\kappa }^2`$. Since $`\frac{1}{g_s}\stackrel{>}{}\rho \stackrel{>}{}\mathrm{\hspace{0.17em}1}`$ to be in the Hagedorn phase, if e.g. the transverse volumes are of order unity in string units, then the contribution of the bulk component is $`\widehat{\kappa }\widehat{\kappa }^2\rho \widehat{\kappa }^2`$. Hence the cosmological effect of the diffuse bulk component can be dominant when $`g_s`$ becomes small. We begin in sections 2 and 3 by deriving the energy-momentum tensor, the principal ingredient of Einstein’s equations. Since the results can be understood rather intuitively, this part of the discussion is organized so that cosmologists (and indeed anybody else) can skip the bulk of sections 2 and 3 concerning string thermodynamics and proceed directly to the energy-momentum tensor which is summarized at the end of section 3. The thermodynamic discussion of section 2 gives a detailed introduction to the behaviour of both type I and type IIA/B open strings on D-branes as calculated from the microcanonical ensemble in a flat background. In particular we discuss the importance of macroscopic modes such as winding modes. Much of this section is a collation of results from ref.. We then reintroduce the classical random walk picture paying special attention to the meaning of quantities such as average string length. In section 3 we use the thermodynamic results to calculate the energy momentum tensor $`T_{\mu \nu }`$ of the Hagedorn phase. $`T_{\mu \nu }`$ enters into the higher dimensional Einstein’s equations and determines the cosmology, and in particular we show that open string winding modes gives negative transverse components. For convenience the results for $`T_{\mu \nu }`$ are summarized at the end of section 3 where we also discuss the heuristic interpretation of this negative pressure. Armed with the energy-momentum tensor, we examine the resultant cosmology. In sections 4 and 5, we solve the equations of motion with various ansätze, and find the advertised inflationary behaviour as well as bouncing solutions with singularity smoothing behaviour. We will, purely for definiteness, consider adiabatic systems in solving the evolution equations of the universe. Under the assumption of adiabaticity, the inflationary growth period drives down the temperature of the system; eventually the temperature drop causes the universe to leave the Hagedorn regime, and consequently inflation ends automatically. However adiabatic systems are probably unable to provide a realistic scenario with sufficient inflation. In section 6 we therefore discuss how, in non-adiabatic systems, inflation can be sustained. We conclude in section 7. ## 2 The Hagedorn phase and random walks This section presents some background thermodynamics needed to get $`T_{\mu \nu }`$ (which enters into Einstein’s equations) in the Hagedorn regime. In section 2.1 we review the thermodynamic properties of D-branes in toroidal compactifications. These compactifications allow us to use the microcanonical ensemble, which is defined in terms of global parameters such as total energy and volume. The importance of the microcanonical ensemble is that it allows a rigorous understanding of the effect of winding modes. This understanding enables us in sections 2.2 and 2.3 to extend the analysis to more general universes of any shape; in particular, we can determine those cases in which thermodynamics and hence cosmology can be studied. Essentially, we will argue that thermodynamics only makes sense in a cosmological setting for those systems in which the canonical and microcanonical ensembles agree. Where this is not the case the systems are dominated by large scale fluctuations. For the systems in which the two ensembles agree, we will derive a expression for the partition function based on a heuristic random walk argument. In particular this expression gives a geometric understanding of the partition function that allows us to discuss its validity in limits of high energy density, small volume, etc. In section 3, we then use this partition function to find the energy momentum tensor, $`T_{\mu \nu }`$. Readers whose principal interest is cosmology may wish to read subsection 2.1 and then proceed directly to the summary of the results for $`T_{\mu \nu }`$ in section 3.4. ### 2.1 String thermodynamics and the Hagedorn phase The Hagedorn phase arises in theories containing fundamental strings because they have a large number of internal degrees of freedom. Indeed, because of the existence of many oscillator modes, the density of states grows exponentially with energy $`\epsilon `$, $`\omega (\epsilon )\epsilon ^be^{\beta _H\epsilon }`$, where the inverse Hagedorn temperature $`\beta _H`$ (where $`\beta =1/T`$) and the exponent $`b`$ depend on the particular theory in question (for example heterotic or type II) . For type I,IIA,IIB strings the numerical value of the inverse Hagedorn temperature is $`\beta _H=2\sqrt{2}\pi `$ in string units. It is easy to see that thermodynamic quantities, such as the entropy, are liable to diverge at the Hagedorn temperature; obtaining the partition function $`Z`$ with the canonical ensemble and multiplying by the usual Boltzmann factor $`e^{\beta \epsilon }`$, one finds an integral for the partition function (for large $`\epsilon `$) $`Z𝑑\epsilon \epsilon ^be^{(\beta _H\beta )\epsilon }`$ which diverges at $`\beta =\beta _H`$ for $`b2`$. If an infinite amount of energy is required to reach $`T_H`$, then we say that the system is limiting. If not then the system is said to be non-limiting. Already the simple canonical ensemble above indicates that the Hagedorn temperature is a limiting temperature for the $`b2`$ cases . The remaining systems seem to be non-limiting, and until recently it was thought that this might imply some kind of Hagedorn phase transition (drawing strongly on the analogy with the quark-hadron phase transition) to more fundamental degrees of freedom. However, in ref. it was shown that all string systems are in fact limiting, including arbitrary D$`p`$-branes (i.e. open strings attached to $`p`$ dimensional defects). The limiting/non-limiting question is a rather subtle one, and in order to resolve it one must first allow for large scale fluctuations by using the microcanonical ensemble (as pointed out in the early papers of Carlitz and Frautschi ) and second, retain full volume dependence (as pointed out for heterotic strings by Deo et al and for open strings in ref.). Once both of these factors are included, it becomes clear that in all cases the Hagedorn temperature is truly a limiting temperature rather than an indication of some sort of phase transition. Perhaps the best evidence for this is that the Hagedorn phase completes (by entropy matching) a phase diagram which includes other non-perturbative phases such as black holes. (There is a sense in which the entire Hagedorn phase can be thought of as a first-order phase transition from Yang-Mills/supergravity degrees of freedom to black-branes/black-holes. We return to this point later.) Let us summarize the rigorous results from ref. for open and closed strings in a toroidally compactified space. The most direct route to the thermodynamics is to evaluate the one loop partition function with the Euclidean time coordinate, $`\tau `$, compactified with radius $`\beta `$. However for the random walk discussion later, it is useful to begin with the density of states $`\omega (\epsilon )`$ of a single string of energy $`\epsilon `$ from which the same results are obtained; $$\omega (\epsilon )=\{\begin{array}{cc}\beta _H\frac{V_{}}{V_{}}V_o\frac{e^{\beta _H\epsilon }}{(\beta _H\epsilon )^{\gamma _o+1}}\hfill & \text{ open}\hfill \\ \beta _HV_c\frac{e^{\beta _H\epsilon }}{(\beta _H\epsilon )^{\gamma _c+1}}\hfill & \text{ closed.}\hfill \end{array}$$ In the above $`V_{}`$ and $`V_{}`$ are the volumes transverse and perpendicular to the D-brane. For open strings, $$\gamma _o=\frac{d_o}{2}1,$$ (1) where $`d_o`$ is the number of dimensions transverse to the brane in which there are no windings and $`V_o`$ is the volume of this space (if there are windings in all dimensions, then $`V_o=𝒪(1)`$ in string units). Similarly, for closed strings, $$\gamma _c=\frac{d_c}{2}$$ (2) where $`d_c`$ is the number of dimensions in which closed strings have no windings and, again, $`V_c`$ is the volume of this space. Note that $`\gamma _o`$ and $`\gamma _c`$ are $`\epsilon `$-dependent critical exponents, because winding modes are quenched or activated depending on the string energy. Below we shall use $`\gamma `$ to stand for either $`\gamma _c`$ or $`\gamma _o`$ as appropriate. We now collect the results obtained in ref. in the microcanonical ensemble working in an approximation to the thermodynamic limit. The two main types of behaviour are the single-long-string or non-limiting $`\mathrm{𝐍𝐋}`$ behaviour (see discussion below where we discuss the consequences of placing ‘non-limiting’ systems in thermal contact with limiting systems), with entropy density $$\sigma S/V_{}$$ (3) of the form $$\sigma _{\mathrm{𝐍𝐋}[\gamma ]}\beta _H\rho \frac{1+\gamma }{V_{}}\mathrm{log}(\rho ),$$ (4) and the various types of ‘limiting’ behaviour $`𝐋[\gamma ]`$ $$\sigma _{𝐋[\gamma ]}\{\begin{array}{cc}\beta _H\rho +2\sqrt{\rho /V_{}}\hfill & \mathrm{if}\gamma =1\hfill \\ & \\ \beta _H\rho +\frac{\gamma 1}{\gamma }\rho ^{\frac{\gamma }{\gamma 1}}\hfill & \mathrm{if}\gamma =\frac{1}{2},\frac{1}{2}\hfill \\ & \\ \beta _H\rho +\mathrm{log}(\rho )\hfill & \mathrm{if}\gamma =0\hfill \\ & \\ \beta _H\rho e^\rho \hfill & \mathrm{if}\gamma =1,\hfill \end{array}$$ up to positive constants of $`𝒪`$ (1) in string units. Here $$\rho E/V_{},$$ (5) where $`E`$ is the total energy of the system. Note that the microcanonical ensemble results are expressed in terms of global parameters such as total energy which are valid in a toroidal compactification. We shall find that it is only possible to generalize the discussion for those systems in which the canonical and microcanonical results agree. The L\[-1\] system is our ‘standard’ high energy regime. From eq.(2), we see that it corresponds to $`d_o=0`$, so that all dimensions have windings and it is the system which is always reached provided that the volumes are finite and the energy density is high enough. These systems are equivalent to $`D1`$ branes where $`D`$ is the total number of space-time dimensions (in other words freely moving open strings). (Once there are many windings, we can T dualize the Dirichlet directions so that they become Neumann directions much smaller than the string scale – the winding modes become a spectrum of Kaluza-Klein modes indicating open strings which are energetic enough to probe all of the $`D1`$ Neumann dimensions – even the small ones.) There are two other sorts of behaviour, marginal limiting $`\mathrm{𝐌𝐋}`$ and weak limiting $`\mathrm{𝐖𝐋}`$ with entropy densities $`\sigma _{\mathrm{𝐌𝐋}}`$ $``$ $`\beta _H\rho \rho ^{2D3}(V_{D1})^{2D4}e^{\rho R^{D3}}`$ $`\sigma _{\mathrm{𝐖𝐋}}`$ $``$ $`\beta _c\rho +{\displaystyle \frac{\beta _cf}{2}}\rho {\displaystyle \frac{\beta _c^2V_{}f^2}{24}}\rho ^2`$ (6) respectively, where $`f=V_{}/V_{}`$. Whether a system is limiting or not is a function of the global parameters of the system such as total energy, $`E`$, and volumes $`V_{}`$$`V_{}`$. For example, imagine increasing the energy of an L\[-1/2\] system. These systems are characteristic of an intermediate energy phase of open string systems possessing one large transverse dimension without windings. As the energy is increased windings will eventually be excited in this direction. If the transverse dimension is of order $`R_{}`$, this happens at an energy threshold $`E=R_{}^2`$, and for higher energies the thermodynamic behaviour changes to that of an L\[-1\] system, the universal high-energy regime<sup>4</sup><sup>4</sup>4 Note that the decoupling of winding modes implied by these equations of state is universal and in particular independent of the zero-temperature vacuum contribution to the partition function which may or may not be finite as the transverse radius becomes infinite . This is because in the finite temperature piece of the partition function the winding modes are always accompanied by a Boltzmann-like suppression factor. A related point is that as well as the above contributions there is a suppressed contribution to the entropy density of order $`\sigma _{YM}T^{D2\gamma _o3}`$ corresponding to a Yang-Mills gas in the $`D2\gamma _o3`$ dimensions with windings and KK modes. In the Euclidean approach this contribution comes from a UV cut-off in the Schwinger integral and represents a field-theoretic contribution which is subdominant for temperatures close to $`1/\beta _H`$ but again is universal.. The NL behaviour is so called because the temperature is higher than $`T_H`$ as can easily be seen when we calculate the temperature of the subsystems from $`S_i(E_i)`$ $`=`$ $`\mathrm{log}\mathrm{\Omega }_i(E_i)`$ $`T_i^1`$ $`=`$ $`\beta _i=S_i/E_i,`$ (7) where $`\mathrm{\Omega }_i`$ is the microcanonical density of states in the subsystem $`i`$. The non-limiting systems formally obey $$T_𝐋<T_H<T_{\mathrm{𝐍𝐋}}.$$ (8) However it is important to realise that this simply means that the $`\mathrm{𝐍𝐋}`$ regimes are transients for a finite ten-dimensional volume. Equilibrium can never be achieved when they are in thermal contact with the surrounding (colder) limiting system of closed strings. We can compare the microcanonical ensemble results reviewed in the previous two pages to the canonical ensemble by taking the thermodynamic limit (i.e. letting $`V_{\|\|}`$ become infinite whilst keeping the density on the brane $`\rho `$ constant). If we consider strings attached to a single brane embedded in $`D=10`$ space time dimensions, then there is universal agreement between the canonical and microcanonical ensembles when there are $`0<d_{}<4`$ transverse dimensions. In this case heading towards the thermodynamic limit quenches winding modes – the transverse dimensions can be thought of as effectively infinite in this limit. However when $`d_{}>4`$ (D$`p`$-branes in ten dimensions with $`p<5`$), modes as large as the transverse dimensions can be quenched if $`V_{}R_{}^2`$, giving $`\mathrm{𝐍𝐋}`$ behaviour, or activated into a $`𝐋[1]`$ system for $`V_{}\stackrel{>}{}\sqrt{V_{}}`$ (when the saddle point approximation is valid). Taking $`V_{}\mathrm{}`$ with $`\rho `$ fixed and $`R_{}`$ constant leads to the onset of L$`[1]`$ behaviour in the thermodynamic limit. $`\mathrm{𝐖𝐋}`$ behaviour occurs in an intermediate region of parameter space where windings are being quenched in an L\[-1\] system and the usual saddle point approximation is invalid. Finally closed strings in the thermodynamic limit have a critical dimension $`D=3`$, where $`D`$ is the number of spacetime dimensions. For $`D3`$, we get standard canonical behaviour, $`𝐋[(D1)/2]`$. On the other hand, for $`D>3`$ winding modes generate ‘marginally limiting’ behaviour $`\mathrm{𝐌𝐋}`$. Some examples of different situations are shown (somewhat impressionistically) in figure 1 where the dimension dependence is evident. The entropic preference of strings for branes with higher dimensionality has the obvious interpretation that, when the strings are volume filling, the larger dimension branes ‘cut across’ more strings. Note that in the microcanonical ensemble we always use the definition of temperature as derived from the microcanonical definition of entropy, $`\beta =_ES(E,V,N)`$ where $`E`$ is the energy of the total system. The discussion above illustrates the main reason for this rather convoluted set of definitions ($`Z(\beta )\mathrm{\Omega }S(E)\beta `$); the effect of winding modes cannot fully be explored using the canonical ensemble. For example the canonical ensemble simply doesn’t know about the ML and WL systems. The microcanonical ensemble also allows us to look at cases where parts of the system never come into equilibrium. This is true for instance for a D$`3`$-brane placed in a reservoir of closed strings. As we have seen, for large transverse volumes, a D$`3`$-brane formally has a temperature that is higher than $`T_H`$ and so any thermal energy that it may have quickly evaporates into the surrounding bath of closed strings. Only the microcanonical ensemble allows one to compare the temperatures of different subsystems. Indeed the failure to do so leads to the false conclusion that the system as a whole might be non-limiting. Moreover, strictly speaking the thermodynamic limit (i.e. the assumption on which the canonical ensemble rests) does not actually exist for open strings. The phase diagram derived in ref. tells us that when trying to take the thermodynamic limit one inevitably encounters a black hole phase. Thus it makes more sense to define the system with global quantities such as total energy, and work in finite volumes taking an approximation to the naive thermodynamic limit if possible. ### 2.2 The random walk interpretation Despite its many advantages the microcanonical ensemble has one distinct disadvantage; in many cosmological backgrounds global properties such as total energy cannot be defined. We therefore need to be able to model these aspects of the thermodynamics with a more flexible approach. The random walk interpretation provides just such flexibility. For example, we can easily understand the NL/L cross over if we interpret a string of energy $`\epsilon `$ as a random walk of length $`\epsilon `$ in string units. The size of such a random walk is $`\sqrt{\epsilon }`$ so that, roughly speaking, to have modes of size $`R_{}`$ we require that the total energy $`E=\rho V_{}V_{}R_{}^2`$. In this classical picture the NL/L cross over occurs precisely when there are modes as large as the transverse volume<sup>5</sup><sup>5</sup>5We should caution against taking this picture too literally although, as we are about to see, it does better than might be expected in explaining the thermodynamics. In particular the way to measure the ‘size’ of the string is to put it in a box with a certain energy and see at what radius it begins to excite winding modes. This should be borne in mind in order to avoid getting into circular arguments.. One helpful feature of the random walk interpretation is that we no longer have to bother about the precise topological properties of the space, merely its size in string units. On the other hand what we have gained from the rigorous analysis of the microcanonical ensemble is confidence in our understanding of which systems are stable, when one can use naive canonical ensemble results (i.e. as an approximation to the thermodynamic limit when the systems are limiting) and when the random walk interpretation is valid. An analysis of random walks was the approach taken by Lowe and Thorlacius for closed strings and Lee and Thorlacius for open strings attached to D-branes . These authors studied the Boltzmann equations for the average number of interacting strings sections of different lengths. The equilibrium solutions for closed string loops of length $`\mathrm{}`$ for example show that the (single string) density of states must be of the form $`\omega (\epsilon )\frac{1}{\epsilon }e^{\beta _H\epsilon }`$, where $`\epsilon =\sigma \mathrm{}`$ and the temperature is related to the average total length $`\overline{L}`$ of string in the ensemble by $$\beta _H=\beta \frac{1}{\overline{L}}.$$ (9) This supports the classical interpretation of the string as a random walk with, $`\epsilon =\sigma \mathrm{}`$, where $`\sigma `$ is the string tension encoding information about, for example, the step-size. In ref. it was shown that such a random walk interpretation also accounts for the volume dependence as follows. First consider the distribution function $`\omega (\epsilon )`$ for closed strings in $`D`$ large space-time dimensions. The energy $`\epsilon `$ of the string is proportional to the length of the random walk. The number of walks with a fixed starting point and a given length $`\epsilon `$ grows exponentially as $`\mathrm{exp}(\beta _H\epsilon )`$. Since the walk must be closed, this overcounts by a factor of the volume of the walk, which we shall denote by $$V(\mathrm{walk})=W.$$ (10) Finally, there is a factor of $`V_{D1}`$ from the translational zero mode, and a factor of $`1/\epsilon `$ because any point in the closed string can be a starting point. The end result is $$\omega (\epsilon )_{\mathrm{closed}}V_{D1}\frac{1}{\epsilon }\frac{e^{\beta _H\epsilon }}{W}.$$ (11) Now, the volume of the walk is proportional to $`\epsilon ^{(D1)/2}`$ if it is well-contained in the volume $`(R\sqrt{\epsilon })`$, or roughly $`V_{D1}`$ if it is space-filling $`(R\sqrt{\epsilon })`$. From here we get the standard result . We have $$\omega (\epsilon )_{\mathrm{closed}}V_{D1}\frac{e^{\beta _H\epsilon }}{\epsilon ^{(D+1)/2}}$$ (12) in $`D`$ effectively non-compact space-time dimensions, and $$\omega (\epsilon )_{\mathrm{closed}}\frac{e^{\beta _H\epsilon }}{\epsilon }$$ (13) in an effectively compact space. Note that these results agree with those of eq.(1); for example to obtain the same result as eq.(12), we take $`d_c=D1`$, the number of space dimensions. We can customize this analysis for open strings attached to a brane by a slight modification of the combinatorics. (The more general case of intersecting $`p,q`$ branes was discussed in ref.). The leading exponential degeneracy of a random walk of length $`\epsilon `$ with a fixed starting point in say the D$`p`$-brane is the same as for closed strings: $`\mathrm{exp}(\beta _H\epsilon )`$. Fixing also the end-point at a particular point of the D$`p`$-brane requires the factor $`1/W`$ to cancel the overcounting, just as in the closed string case. Now, both end-points move freely in the part of the brane occupied by the walk. This gives a further degeneracy factor $$W_{}^2$$ (14) from the positions of the end-points. Finally, the overall translation of the walk in the excluded transverse volume gives a factor $`V_{}/W_{}`$. The final result is: $$\omega (\epsilon )_{\mathrm{open}}\frac{V_{}}{W_{}}\mathrm{exp}(\beta _H\epsilon ).$$ (15) Thus we see the sensitivity of the density of states to the effective volume of the random walk in the transverse directions. If the walk is well-contained in these directions, $`(R_{}\sqrt{\epsilon })`$, then we find $`W_{}\epsilon ^{d_{}/2}`$ and $$\omega (\epsilon )_{\mathrm{open}}\frac{V_{}}{\epsilon ^{d_{}/2}}\mathrm{exp}(\beta _H\epsilon )R_{}\sqrt{\epsilon }.$$ (16) On the other hand, if it is space-filling in the transverse directions $`(R_{}\sqrt{\epsilon })`$, the transverse volume of the walk is just $`W_{}V_{}`$ and we find $$\omega (\epsilon )_{\mathrm{open}}\frac{V_{}}{V_{}}\mathrm{exp}(\beta _H\epsilon )R_{}\sqrt{\epsilon },$$ (17) in agreement with eq.(2.1). ### 2.3 Random walks in a cosmological background The above thermodynamics involved a Euclidean metric. How can we adapt the results for a cosmological background? First we assumed that the density of states increased as $`\mathrm{exp}(\beta _H\epsilon )`$. In a sense the parameter $`\epsilon `$ is now no longer the energy but becomes merely something to be integrated over. In the case of a non-trivial metric the most natural interpretation of $`\epsilon `$ is that it is the proper length of a string in the bulk and certainly we can always go to the local inertial frame in which a small portion of the string has the usual Euclidean energy $``$ length equivalence. In addition it is clear that the average number of strings of a given proper length is well defined and so the arguments of Lee, Lowe and Thorlacius then give us the correct distribution of proper string lengths in equilibrium. Furthermore eq.(9) then gives us a working definition of temperature in terms of average proper string lengths. To put these arguments on a firmer footing we now derive one aspect of the strings’ behaviour in a cosmological background that will be useful in the discussions that follow; namely that in a slowly varying background a classical string of proper length $`\mathrm{}\epsilon `$ occupies a volume of proper-size $`\sqrt{\epsilon }`$. Consider a classical string of proper length $`\mathrm{}=\epsilon /\sigma `$ (where $`\sigma 1`$ in string units) with one end point fixed in a $`D1`$ dimensional space. From our observations above we would expect such a string to have a density of states $`\omega (\epsilon )e^{\beta _H\epsilon }`$. The crucial point is that we can arbitrarily divide this string into $`N`$ small strings each of which has one end free and one end that is fixed to the end of the previous string. Consequently the density of states of the large single string is the same as that of $`N`$ small strings each of which has one end fixed and with energies obeying $`_i\epsilon _i=\epsilon `$ where $`\epsilon _i`$ is the energy of each string in its local inertial frame. By choosing a sufficiently large $`N`$ we may always use the flat space approximation to evaluate the density of states of each small string rigorously (with one end fixed and one end free), $$\omega _i(\epsilon _i)=e^{\beta _H\epsilon _i}.$$ (18) In particular we find the same $`\beta _H`$ for all the strings. The total density of states is then given by $$\omega (\epsilon )=\underset{i}{}\left(\omega _id\epsilon _i\right)\delta (\epsilon \underset{i}{}\epsilon _i)=e^{\beta _H\epsilon },$$ (19) as required. It is now clear how to find the region occupied by the string. We measure this by determining the gyration, defined as the average size of the fluctuations of the free end from the fixed end, measured along null geodesics passing through the latter. Thermodynamics in the local inertial frames indicates that each small string has a spherically symmetric gyration with a radius $`\sqrt{\epsilon _i}`$. To find the gyration of the large string we must (since they are average fluctuations) add those of the subsystems in quadrature. Hence, regardless of which geodesic we choose, the combined system has a gyration $`\sqrt{_i\epsilon _i}=\sqrt{\epsilon }`$. Proceeding now to the volume dependence of $`\omega (\epsilon )`$, we first make the usual quasi-equilibrium approximation that equilibrium is established much more quickly than any change in the metric so that the metric may be taken to be approximately constant in time when evaluating properties such as density. In order to simplify matters, we also make the (by now) familiar assumption that the metric is factorizable into parallel dimensions $`x`$ and transverse ones $`y`$ $$ds^2=n^2dt^2+g_{\|\|ij}dx^idx^j+g_{nm}dy^ndy^m,$$ (20) with the brane lying at $`y^n=0`$. We define parallel and total volumes as $`V_{\|\|}(y^{})`$ $`=`$ $`{\displaystyle 𝑑x𝑑y\delta (yy^{})\sqrt{g_{\|\|}}}`$ $`W_{\|\|}(y^{})`$ $`=`$ $`{\displaystyle 𝑑x𝑑y\delta (yy^{})\sqrt{g_{\|\|}}\eta (y,x)}`$ $`V`$ $`=`$ $`{\displaystyle 𝑑x𝑑y\sqrt{g_{\|\|}}\sqrt{g_{}}}`$ $`W`$ $`=`$ $`{\displaystyle 𝑑x𝑑y\sqrt{g_{\|\|}}\sqrt{g_{}}\eta (y,x)}`$ (21) where $`\eta `$ is a function that is one in the region of the random walk and zero elsewhere, and where $`g_{\|\|}`$ and $`g_{}`$ are the determinants of the metric in the brane dimensions and transverse dimensions respectively. Note that we are using $`x`$ and $`y`$ as shorthand for all the parallel and transverse coordinates. We also define an averaging over the extra dimension with an overbar, $$\overline{O}=\frac{𝑑y\sqrt{g_{}}O(y)}{𝑑y^{}\sqrt{g_{}}}.$$ (22) For notational convenience we treat $`\sqrt{g_{}(y)}`$ as a ‘$`y`$-dependent transverse volume’, $$V_{}(y)=\sqrt{g_{}(y)}𝑑y^{},$$ (23) so that the actual transverse volume is written $$\overline{V_{}}=𝑑y\sqrt{g_{}(y)}.$$ (24) Note that $`V=\overline{V_{}}\overline{V_{}}`$. From the discussion above, we have $`W_{\|\|}=\epsilon ^{d_{\|\|}/2}`$ reflecting the fact that the string has typical size $`\sqrt{\epsilon }`$. For example if we take a slice at $`y=y^{}`$ this should be $$W_{\|\|}=\epsilon ^{d_{\|\|}/2}=𝑑x𝑑y\delta (yy^{})\sqrt{g_{\|\|}}\eta =𝑑x\sqrt{g_{\|\|}}(y^{})\eta (y^{})$$ (25) so that $`\eta `$ must compensate for $`\sqrt{g_{\|\|}}`$ in order to make this volume independent of $`y^{}`$. In order to examine the thermodynamics we now need to decide when we can use the microcanonical results. We first divide the parallel dimensions into small locally flat patches with $`\sqrt{g_{}}`$ approximately independent of $`x`$. In each patch it is consistent to use the microcanonical results, provided there are no long range fluctuations. Thus we are generally prohibited from examining the NL systems since adjacent patches will not be in equilibrium. <sup>6</sup><sup>6</sup>6Also note that by examining the thermodynamics in a small region we do not artificially quench Kaluza-Klein modes as long as the local patches are much larger than the string scale. This is not the case in the perpendicular directions however because of the presence of macroscopic winding modes. Hence in a local patch with average volume $`\overline{V_{}}`$ we may define an energy density by for example $`\rho =E/\overline{V_{}}`$ where $`E`$ is the total energy in the patch. However, for the thermodynamics, we must always include the whole transverse volume $`\overline{V_{}}`$ to avoid artifically unquenching winding modes. The random walk argument proceeds as before with one change. The factor that corrects for the translation of the walk in the excluded volume is not $`V_{}/W_{}`$ but rather $`V/W`$ for the L\[-1\] system. In a toroidal compactification these are of course the same for the L\[-1\] case. Here however, we must take account of the fact that the walk is squeezed by the metric so that it may have more ‘room’ at one side of the compactified dimension than the other. This factor can also be written as $$\frac{V}{W}=\frac{\overline{V_{}}}{W_{}},$$ (26) and we now find, for example, that the density of states in the L\[-1\] system is $$\omega (\epsilon )\frac{\overline{V_{}}}{\overline{V_{}}}e^{\beta _H\epsilon },$$ (27) i.e. the same expression as in the flat space case but with all volumes averaged over transverse dimensions as in eqs.(22) and (23). It is now possible to find $`\mathrm{log}Z(\beta ,\overline{V_{}},\overline{V_{}})`$ from $`\omega (\epsilon )`$ by integrating over $`\epsilon `$ with a Boltzmann weighting, which we do for the various sytems in the following section. Note that the $`\beta `$ appearing in the resulting partition function is the random walk definition of ‘inverse temperature’ whose precise physical interpretation is given by eq.(9). ## 3 Stress-energy tensor $`T_{\mu \nu }`$ in a bulk Hagedorn phase We now use these thermodynamic results to find the bulk energy momentum tensor during the Hagedorn regime. We may find the energy momentum tensor from $$T_\nu ^\mu =2\frac{g^{\mu \rho }}{\sqrt{g}}\frac{\delta \mathrm{log}Z(\beta ,\overline{V_{}},\overline{V_{}})}{\delta g^{\rho \nu }},$$ (28) where $`\beta ,\overline{V_{}},\overline{V_{}}`$ are given by eqs.(9),(2.3),(22),(23). We will treat the functional derivative with respect to $`g_{\mu \nu }`$ in the following way. We assume that small changes in the metric correspond to making small changes in the volumes in $`\mathrm{log}Z`$, e.g. for a single extra dimension $$\frac{\delta Z}{\delta g_{55}}=𝑑x^{}\frac{\delta Z}{\delta \overline{V}_{}(x^{})}\frac{\delta \overline{V}_{}(x^{})}{\delta g_{55}}.$$ (29) Then, in the case of only one extra dimension, we can write $$\overline{V_{}}=𝑑y\sqrt{g_{55}}=\frac{d^5x\sqrt{g_{55}}}{d^4x}=\frac{1}{v\beta }d^5x\sqrt{g_{55}},$$ (30) and $$\frac{\delta \overline{V_{}}}{\delta g_{55}}=\frac{1}{2\sqrt{g_{55}}v\beta }.$$ (31) Then from eqs.(2931) we can determine the functional derivative in eq.(28). Our ansatz automatically means that $`T_{05}=0`$ and hence $`G_{05}=0`$; in other words we are not considering energy exchange between the brane and the bulk. (In general there might be energy flux between the two.) We begin by evaluating the energy momentum tensor for our ‘standard’ case of $`𝐋[1]`$ (windings in all transverse dimensions) appropriate to high energies and densities; the resulting $`T_{\mu \nu }`$ is presented in eq.(3.1). We then proceed to $`𝐋[\gamma 1]`$ cases (in which windings are quenched in some of the transverse dimensions) and present the resulting $`T_{\mu \nu }`$ in eq.(56). ### 3.1 $`T_{\mu \nu }`$ at high energies and densities; L\[-1\] We now apply eq.(28) to the very high energy regime in which there are modes stretching across the whole space in all transverse dimensions. In the nomenclature of ref. these are the limiting L\[-1\] systems since $`d_o=0`$. We first need the principal cosmological ingredient, the partition function. The expression for the density of states in the L\[-1\] system is $$\omega (\epsilon )\frac{\overline{V_{}}}{\overline{V_{}}}e^{\beta _H\epsilon },$$ (32) which, when integrated with a Boltzmann weighting, gives $$\mathrm{log}Z(\beta ,\overline{V_{}},\overline{V_{}})=2\frac{\overline{V_{}}^2\beta _H^2}{\overline{V_{}}(\beta ^2\beta _H^2)}+a_c\overline{V_{}}\overline{V_{}}\rho _c\frac{(\beta ^2\beta _H^2)}{2\beta _H},$$ (33) where $`a_c`$ and $`\rho _c`$ are a critical pressure and energy density (defined with reference to the brane dimensions, i.e. with dimensions $`E_c/\overline{V_{}}`$) which are of order unity in string units . Here, subscript-$`c`$ refers to critical quantities to remain in the Hagedorn phase. These are the successive terms in a saddle point approximation. After doing the functional derivative we can replace $`\beta `$ by using the saddle point result for the L\[-1\] system, $$2\frac{\overline{V_{}}\beta _H^2}{\overline{V_{}}(\beta ^2\beta _H^2)}\sqrt{\frac{\overline{V_{}}}{\overline{V_{}}}\beta _H(\rho \rho _c)},$$ (34) where as before, $`\rho =E/\overline{V_{}}`$. For the saddle point approximation to be valid we require $$(\rho \rho _c)R_{}^4.$$ (35) This is also the condition for the leading (first) term in $`\mathrm{log}Z`$ to dominate over the $`a_c`$ term, and for the $`a_c`$ term to dominate over the $`\rho _c`$ term. We neglect higher order terms. In order to write the resulting energy-momentum tensor, we define the parameters $$\alpha _{\|\|}=\frac{\overline{V_{}}}{V_{}}\text{ ; }\alpha _{}=\frac{\overline{V_{}}}{V_{}}$$ (36) and a pressure $$\widehat{p}_1=\frac{\sqrt{\rho }}{\overline{V_{}}^{3/2}}.$$ (37) In the following section we shall argue that $`\alpha _{\|\|}1`$ wherever the Hagedorn phase is the dominant phase. If the off-diagonal components of the metric are small, then the variation with respect to $`g^{\mu \nu }`$ gives the bulk $`\widehat{T}_{\mu \nu }`$ terms; $`\widehat{T}_0^0`$ $`=`$ $`\widehat{\rho }={\displaystyle \frac{\rho \alpha _{\|\|}\alpha _{}}{\overline{V_{}}}}`$ $`\widehat{T}_i^i`$ $`=`$ $`\widehat{p}_1`$ $`\widehat{T}_m^m`$ $`=`$ $`\widehat{p}_1(2\alpha _{\|\|}1)\widehat{p}_1,`$ (38) where $`i`$ labels the parallel $`x`$ directions and $`m`$ labels the transverse $`y`$ directions. We have dropped the $``$ notation for the expectation value of the energy momentum tensor, but shall continue to use hats to signify bulk properties such as the bulk density above, $`\widehat{\rho }`$. It is instructive to consider the conservation equations derived from the Bianchi identity $`T_{;\mu }^{\mu \nu }=0`$. For example, restricting ourselves to the often discussed scenario with one extra dimension, $`y`$, and metric of the form $$ds^2=n^2dt^2+a^2dx^2+b^2dy^2,$$ (39) the E-M conservation equations may then be written (we are using the same notation as in ref.) $`{\displaystyle \frac{d\widehat{T}_0^0}{dt}}+(\widehat{T}_0^0\widehat{T}_i^i)3{\displaystyle \frac{\dot{a}}{a}}+(\widehat{T}_0^0\widehat{T}_5^5){\displaystyle \frac{\dot{b}}{b}}`$ $`=`$ $`0`$ $`{\displaystyle \frac{d\widehat{T}_5^5}{dy}}+\widehat{T}_5^5\left({\displaystyle \frac{n^{}}{n}}+{\displaystyle \frac{a^{}}{a}}\right)+{\displaystyle \frac{n^{}}{n}}\widehat{T}_0^03{\displaystyle \frac{a^{}}{a}}\widehat{T}_i^i`$ $`=`$ $`0.`$ (40) We shall see shortly that $`\alpha _{\|\|},\alpha _{}1`$ so the first equation is a conservation law for entropy; $$\frac{d(\sigma \overline{V_{}})}{dt}\left(1+\sqrt{\overline{V_{}}/\rho }\right)2\sqrt{\overline{V_{}}/\rho }.\frac{d}{dt}\left(\overline{V_{}}\sqrt{\rho /\overline{V_{}}}\right)$$ (41) where $$\sigma =\rho +2\sqrt{\rho /\overline{V_{}}}.$$ (42) Recall from ref. that the saddle point dominance condition is $$\sqrt{\overline{V_{}}/\rho }\frac{1}{R_{}^2}1.$$ (43) Hence to order $`1/R_{}^2`$ we simply recover the expected entropy conservation law, $`\frac{d(\sigma \overline{V_{}})}{dt}=0`$. Note that the entropy is evenly distributed because of the dominant first term; indeed to first order $`\widehat{T}_0^0=\widehat{\sigma }(y)`$ where $`\widehat{\sigma }`$ is the local entropy density in the bulk. ### 3.2 $`T_{\mu \nu }`$ for the L\[$`\gamma 1`$\] systems. As the energy density of an L\[-1\] system falls below certain energy thresholds, the modes that are sensitive to the size of the transverse dimensions (e.g. winding modes in toroidal compactifictions) become quenched. Once this happens, the compactness of the quenched directions is of no further consequence for the thermodynamics, which aquires a different critical exponent, $`\gamma `$, and can even temporarily become NL (although the name ‘non-limiting’ is probably misleading for the reasons discussed in the Introduction and in ref.). If there are many transverse dimensions, and they are anisotropic, there may be a few energy thresholds and $`\gamma `$ may assume several (increasing) values before the density drops below the Hagedorn density $`\rho _c`$ and the system finally leaves the Hagedorn phase and enters the Yang-Mills phase. In order to discuss the thermodynamics of these intermediate systems, we first introduce the function $`\eta ^{}(y)`$ which gives support in the regions occupied by the random walking strings in the various transverse dimensions. (Note that $`\eta (x,y)`$ described the shape of a single string whereas $`\eta ^{}(y)`$ is for all strings, so $`\eta ^{}(y)=0`$ indicates that strings do not extend this far into the bulk.) We define an averaging over the region given support by $`\eta ^{}`$; $$\overline{O}_\eta ^{}=\frac{𝑑y\eta ^{}\sqrt{g_{}}O(y)}{𝑑y^{}\eta ^{}\sqrt{g_{}}}.$$ (44) The single string density of states density of states now involves the ratio of the volumes $`V_\eta ^{}`$ $`=`$ $`{\displaystyle 𝑑x𝑑y\eta ^{}\sqrt{g_{}}\sqrt{g_{\|\|}}}`$ $`W`$ $`=`$ $`{\displaystyle 𝑑x𝑑y\eta ^{}\eta \sqrt{g_{}}\sqrt{g_{\|\|}}}`$ (45) (where $`\eta `$ is defined below eq. (2.3)) so that $$\frac{V_\eta ^{}}{W}=\frac{\overline{V_{\|\|}}_\eta ^{}}{W_{\|\|}}.$$ (46) The expression for the single string density of states becomes $$\omega (\epsilon )\frac{\overline{V_{\|\|}}_\eta ^{}}{\overline{V_{}}_\eta ^{}}\frac{e^{\beta _H\epsilon }}{\epsilon ^\gamma }$$ (47) where in accord with our previous definition of $`\eta `$, we have assumed that the strings occupy a proper volume $`\epsilon ^{d_0/2}`$ in the $`d_0`$ dimensions where the modes do not fill the whole of space. As before we have $`\gamma _0=d_0/21`$. Defining $$\overline{f}=\frac{\overline{V_{\|\|}}_\eta ^{}}{\overline{V_{}}_\eta ^{}}$$ (48) we may now use the expressions for the singular part of $`\mathrm{log}Z`$ found in ref.; $$\mathrm{log}Z_{sing}\{\begin{array}{cc}\mathrm{\Gamma }(\gamma )\overline{f}(\beta \beta _c)^\gamma ,\hfill & \gamma \text{Z}^+\{0\}\hfill \\ & \\ \frac{(1)^{\gamma +1}}{\mathrm{\Gamma }(\gamma +1)}\overline{f}(\beta \beta _c)^\gamma \mathrm{log}(\beta \beta _c)\hfill & \gamma \text{Z}^+\{0\}.\hfill \end{array}$$ In the microcanonical ensemble, these expressions gives us the same relations as in eq.(2.1) with $`\rho `$ replaced by $$\overline{\rho }=\frac{E}{\overline{V_{}}_\eta ^{}},$$ (49) whence, substituting the appropriate expression for $`\overline{\rho }`$ from eq.(2.1) for each $`\gamma `$, we find the following for the L\[$`\gamma 1`$\] systems; $`\widehat{T}_0^0`$ $`=`$ $`\eta ^{}\widehat{\rho }`$ $`\widehat{T}_i^i`$ $`=`$ $`\eta ^{}\widehat{p}_\gamma `$ $`\widehat{T}_m^m`$ $`=`$ $`\{\begin{array}{cc}\eta ^{}\widehat{p}_\gamma \left(2\alpha _{\|\|}1\right)\eta ^{}\widehat{p}_\gamma \hfill & \text{windings}\hfill \\ 0\hfill & \text{no windings,}\hfill \end{array}`$ (51) where $$\widehat{p}_\gamma =\{\begin{array}{cc}\frac{1}{\overline{V_{}}_\eta ^{}^{3/2}}\overline{\rho }^{\frac{\gamma }{\gamma 1}}\hfill & \mathrm{if}\gamma =\frac{1}{2},\frac{1}{2}\hfill \\ & \\ \frac{1}{\overline{V_{}}_\eta ^{}^{3/2}}\mathrm{log}\overline{\rho }\hfill & \mathrm{if}\gamma =0\hfill \\ & \\ \frac{1}{\overline{V_{}}_\eta ^{}^{3/2}}e^{\overline{\rho }}\hfill & \mathrm{if}\gamma =1.\hfill \end{array}$$ Because, as we see later, $`\alpha _{\|\|,}1`$, we shall henceforth drop the overbar notation. ### 3.3 NL\[$`\gamma `$\], WL and ML systems. For the NL systems, $`\gamma >1`$ and the microcanonical calculation tells us that $$E=\frac{(1+\gamma )}{\beta \beta _c}.$$ (52) In these cases we have already argued that they cannot be in equilibrium with the remaining systems since the temperature that we formally derive from the microcanonical ensemble has $`T>T_H`$. Indeed eq.(9) also indicates the average total length of strings in the ensemble cannot be positive. The fact that there is no equilibrium solution to the Boltzmann equations is merely a different way of seeing that these systems are transient. Little more information can be gained from the random walk picture in this case. However it seems likely that, as in the flat space microcanonical calculation, these systems tend to lose energy to closed strings in a cosmological setting. The WL systems correspond to open strings that do not satisfy the saddle point condition. This is a cross over region where winding modes are just beginning to be excited, and the corresponding specific heat is therefore small (hence ‘weak limiting’). Unfortunately this case is also difficult to analyse because the microcanonical ensemble does not agree with the canonical result (indeed there is no canonical equivalent of these systems). Consequently it is difficult to find anything that we may interpret as $`T_\mu ^\nu `$ for these cases. The same is true for the ML systems which correspond to closed strings in the Hagedorn regime. Fortunately all of these systems have lower entropy than the limiting systems and can be neglected in the cosmology. ### 3.4 Summary and discussion of $`T_{\mu \nu }`$ We now summarize the results for the energy-momentum tensor. We first drop the overline notation of the previous section and simply redefine $`V_{}`$ and $`V_{}`$ to be the transverse and parallel volumes covariantly averaged over the region of the transverse dimensions covered by the strings. (See eqs.(2.3),(23),(44).) We define an energy density $`\rho `$ of strings $$\rho =\frac{E}{V_{}}.$$ (53) For the limiting systems with critical exponent $`\gamma `$, L\[$`\gamma `$\], we find that the ‘bulk’ components of the energy momentum tensor are given by $`\widehat{T}_0^0`$ $`=`$ $`\widehat{\rho }`$ $`\widehat{T}_i^i`$ $`=`$ $`\widehat{p}_\gamma `$ $`\widehat{T}_m^m`$ $``$ $`\{\begin{array}{cc}\widehat{p}_\gamma \hfill & \text{transverse with windings}\hfill \\ 0\hfill & \text{transverse without windings,}\hfill \end{array}`$ (56) where $`\widehat{\rho }=\rho /V_{}`$ and $$\widehat{p}_\gamma =\{\begin{array}{cc}\frac{1}{V_{}^{3/2}}\rho ^{\frac{\gamma }{\gamma 1}}\hfill & \mathrm{if}\gamma =1,\frac{1}{2},\frac{1}{2}\hfill \\ & \\ \frac{1}{V_{}^{3/2}}\mathrm{log}\rho \hfill & \mathrm{if}\gamma =0\hfill \\ & \\ \frac{1}{V_{}^{3/2}}e^\rho \hfill & \mathrm{if}\gamma =1,\hfill \end{array}$$ wherever there are strings present, and zero otherwise. The approximation in the transverse components is valid when the parallel volume changes only by a small fraction over the transverse directions. For the remaining NL, WL, ML systems we do not know how to consider the cosmology or indeed whether they have any meaning at all in a cosmological setting. In the NL case large scale fluctations make it impossible to have a stable macroscopic system. For the WL and ML systems, one cannot properly take into account the effect of winding modes. The fact that the canonical and microcanonical ensembles disagree in a toroidally compactified space tells us that they play a dominant role in the thermodynamics, but it is not possible to formulate the required microcanonical treatment in cosmology since well defined global parameters such as total energy are generally lacking. As expected $`\widehat{T}_0^0`$ resembles the local energy density of strings. The $`\widehat{T}_i^i`$ represents a relatively small pressure coming from Kaluza-Klein modes in the Neumann directions and $`\widehat{T}_m^m`$ is a negative pressure coming from winding modes in the Dirichlet directions. If we T-dualize the Dirichlet directions these ‘winding modes’ also become Kaluza-Klein modes in Neumann-directions and $`T_m^m`$ becomes positive. Thus negative $`T_m^m`$ reflects the fact that we have T-dualized a dimension much smaller than the string scale thereby reversing the pressure. For this reason negative $`T_\mu ^\mu `$ is expected to be a general feature of space-filling excitations in transverse dimensions. Those who are familiar with T-duality may suspect an apparent conflict with this reversal of pressure when we T-dualize the extra dimensions. However we stress that T-duality is maintained. Pressure is merely defined as the change in free energy with increase in volume. Thus since free energy is invariant under T-duality the pressure must reverse sign. On the other hand the cosmological consequences in the T-dual system must of course be the same. In particular gravitons do not propagate in the extra dimensions of the T-dualized system (i.e. their Kaluza-Klein modes are heavy) since the extra dimensions are much smaller than the string scale. However gravitational degrees of freedom come from the closed string sector in the bulk which has both Kaluza-Klein modes and winding modes in all directions. Hence in the T-dual system the light gravitational degrees of freedom in the extra dimensions are winding modes. The original transverse components of Einstein’s equations (which are valid only in the approximation that the Dirichlet directions are much larger than the string scale) must in the T-dual system be replaced by the equations of motion derived from the action now expressed as a sum over gravitational winding sectors. The upshot is that T-duality dictates that the resulting equations have the same consequences for the brane cosmology as the original system. ## 4 Cosmological Equations in $`D=5`$ In the previous sections we obtained the bulk stress-energy tensor at temperatures close to the Hagedorn temperature for the most general limiting cases in arbitrary numbers of dimensions. We now proceed to a discussion of the cosmology for which we will for the most part consider a ‘toy-model’ system with the following restrictions. The first restriction is that we implicitely be considering a D-brane configuration that has 3 large parallel dimensions (i.e. the ‘observable universe’) and only one transverse dimension $`y`$ that supports winding modes. Such a configuration could possibly arise in a 5 dimensional intersection of D-branes. A more realistic case would probably include all $`D=10`$ dimensions but the Einstein equations are significantly more complicated in this case and we leave their full examination to a later paper. However in section 5.6 we will discuss why we expect to find the same behaviour in higher dimensions. The second restriction we make is simply for definiteness; we assume adiabaticity when solving Einstein equations. We shall find in this section and in section 5 both power law and exponential inflation. In principal, as we will see in section 6, it is therefore possible to solve the horizon problem even with adiabaticity. In practice however it is not possible to introduce a priori the entire entropy of the observable universe (i.e. $`S10^{88}`$) onto a primordial brane of order the string scale without either destroying the brane or having an unreasonably small string coupling (i.e. $`g_s10^{88}`$!). In section 6 therefore we shall discuss possible deviations from adiabaticity. We use the metric of eq.(39), $$ds^2=n^2dt^2+a^2dx^2+b^2dy^2,$$ (57) which foliates the space into flat, homogeneous, and isotropic spatial 3-planes. Here $`𝐱=x_1,x_2,x_3`$ are the coordinates on the spatial 3-planes while $`y`$ is the coordinate of the extra dimension. For simplicity, we make a further restriction by imposing $`Z_2`$ symmetry under $`yy`$. Without any loss of generality we choose the 3-brane of the ‘observable universe’ to be fixed at $`y=0`$. (Even if the brane is moving we can always reparameterize the theory so that the 3-brane is fixed in coordinate space, although in doing so we lose further freedom of gauge choice .) With our assumption of homegeneity we can associate a scale factor with the parallel dimensions $`a(t,y)`$ and one for the ‘extra’ transverse dimensions $`b(t,y)`$. We define $`a_0(t(\tau ))=a(t,0)`$ as the scale factor describing the expansion of the 3-brane where $`t(\tau )𝑑\tau n(\tau ,y=0)`$ is the proper time of a comoving observer. Several authors \- have presented the bulk Einstein equations but for completeness we briefly restate the results; $`\widehat{G}_{00}`$ $`=`$ $`3\left\{{\displaystyle \frac{\dot{a}}{a}}\left({\displaystyle \frac{\dot{a}}{a}}+{\displaystyle \frac{\dot{b}}{b}}\right){\displaystyle \frac{n^2}{b^2}}\left[{\displaystyle \frac{a^{\prime \prime }}{a}}+{\displaystyle \frac{a^{}}{a}}\left({\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{b^{}}{b}}\right)\right]\right\}=\widehat{\kappa }^2\widehat{T}_{00},`$ (58) $`\widehat{G}_{ii}`$ $`=`$ $`{\displaystyle \frac{a^2}{b^2}}\left\{{\displaystyle \frac{a^{}}{a}}\left({\displaystyle \frac{a^{}}{a}}+2{\displaystyle \frac{n^{}}{n}}\right){\displaystyle \frac{b^{}}{b}}\left({\displaystyle \frac{n^{}}{n}}+2{\displaystyle \frac{a^{}}{a}}\right)+2{\displaystyle \frac{a^{\prime \prime }}{a}}+{\displaystyle \frac{n^{\prime \prime }}{n}}\right\}`$ (59) $`+`$ $`{\displaystyle \frac{a^2}{n^2}}\left\{{\displaystyle \frac{\dot{a}}{a}}\left({\displaystyle \frac{\dot{a}}{a}}+2{\displaystyle \frac{\dot{n}}{n}}\right)2{\displaystyle \frac{\ddot{a}}{a}}+{\displaystyle \frac{\dot{b}}{b}}\left(2{\displaystyle \frac{\dot{a}}{a}}+{\displaystyle \frac{\dot{n}}{n}}\right){\displaystyle \frac{\ddot{b}}{b}}\right\}=\widehat{\kappa }^2\widehat{T}_{ii},`$ $`\widehat{G}_{05}`$ $`=`$ $`3\left({\displaystyle \frac{n^{}}{n}}{\displaystyle \frac{\dot{a}}{a}}+{\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{\dot{b}}{b}}{\displaystyle \frac{\dot{a}^{}}{a}}\right)=\widehat{\kappa }^2\widehat{T}_{05},`$ (60) $`\widehat{G}_{55}`$ $`=`$ $`3\left\{{\displaystyle \frac{a^{}}{a}}\left({\displaystyle \frac{a^{}}{a}}+{\displaystyle \frac{n^{}}{n}}\right){\displaystyle \frac{b^2}{n^2}}\left[{\displaystyle \frac{\dot{a}}{a}}\left({\displaystyle \frac{\dot{a}}{a}}{\displaystyle \frac{\dot{n}}{n}}\right)+{\displaystyle \frac{\ddot{a}}{a}}\right]\right\}=\widehat{\kappa }^2\widehat{T}_{55},`$ (61) where we use the notation of ref. in which $$\widehat{\kappa }^2=8\pi \widehat{G}=8\pi /M_5^3,$$ (62) where $`M_5`$ is the five-dimensional Planck mass and the dots and primes denote differentiation with respect to $`t`$ and $`y`$, respectively. Note that, as stated earlier, our ansatz implies that $`T_{05}=0`$ in the bulk. The bulk equations only apply in an open region that does not include the boundary. It is the Israel conditions which connect the boundary and the bulk -; $$\underset{\pm \text{faces}}{}\left(K_{\mu \nu }Kh_{\mu \nu }\right)=t_{\mu \nu },$$ (63) where $`K_{\mu \nu }`$ is the extrinsic curvature, the sum over faces is for each side of the boundary surface and we have defined $$t_{\mu \nu }=\frac{2}{\sqrt{h}}\frac{\delta S_{boundary}}{\delta h^{\mu \nu }}$$ (64) as the energy momentum tensor on the boundary. We will assume that this energy momentum tensor on the boundary can be written in a perfect fluid form. For example, $$t_0^0=\rho _{br}$$ (65) and $$t_1^1=p_{br},$$ (66) where $`\rho _{br}`$ and $`p_{br}`$ are the energy density and pressure, respectively, measured by a comoving observer. We can now find the Israel conditions specific to the metric in eq.(39). Because we are considering a brane at a $`Z_2`$ symmetry fixed plane, from the sum over faces we get two identical terms, merely resulting in a factor of two. In a more general spacetime, the Israel conditions would still hold, but one would have to explicitly add the contributions from the two sides of the boundary surfaces since these contributions would no longer be equal. With the $`Z_2`$ symmetry, the Israel condition becomes, $`3[a^{}/a]_0`$ $`=`$ $`\widehat{\kappa }^2b_0\rho _{br}`$ $`3[n^{}/n]_0`$ $`=`$ $`\widehat{\kappa }^2b_0(2\rho _{br}+3p_{br})`$ (67) The same result can be found by adding $`\delta (y)t_\nu ^\mu `$ to $`T_\nu ^\mu `$ and equating the coefficients of $`\delta (y)`$ in Einstein’s equations. It is interesting at this point to consider the relative contributions of the stringy excitations and of the D-brane tension to the cosmology. In particular we, somewhat counterintuitively, find that the diffuse stringy component can have a dominant effect on the cosmology even in the weak coupling limit. To see this, first note that the effective gravitational coupling $`\widehat{\kappa }`$ is given in terms of the string coupling $`g_s`$ by $$\frac{1}{\widehat{\kappa }^2}\frac{m_s^{D2}Vol_{D5}}{g_s^2}M_5^3.$$ (68) where $`Vol_{D5}`$ is the volume of the compactification from $`D=10`$ to $`5`$ dimensions, $`m_s`$ is the string scale and $`M_5`$ is the 5-dimensional Planck mass. The intrinsic tension of a D-brane is given by $$T_{Dp}=\frac{1}{(2\pi )^pg_s}\frac{1}{\widehat{\kappa }\sqrt{Vol_{D5}}}$$ (69) For convenience we will henceforth assume that $`Vol_{D5}m_s^5`$ so that $`g_s`$ and $`\widehat{\kappa }`$ are of the same order of magnitude in string units (where $`m_s1`$). (This need not be the case if some of the other space dimensions are compactified with a radius much larger or smaller than the string scale.) Thus the intrinsic tension of the D-brane satisfies $`\rho _{br}1/\widehat{\kappa }`$, and in the small coupling limit one might expect it to dominate the cosmology. However, this is not the case. If we substitute the Israel matching conditions into the cosmological equations (58-61), we see that the contributions of the D-brane tension and of the stringy components in the cosmological equations are of order $$\widehat{\kappa }^4\rho _{br}^2\text{ and }\widehat{\kappa }^2\frac{\rho }{V_{}}\text{ , }\widehat{\kappa }^2\widehat{p}_\gamma $$ respectively. When the density of the string gas is close to the critical Hagedorn density, $`\rho _c1`$ (the effective lower bound), all the terms are of order $`\widehat{\kappa }^2`$ and are therefore comparable. Since the density of the string gas can be significantly larger than the lower bound of $`\rho _c1`$, clearly the bulk stringy contribution can be significantly larger than the brane contribution. We can further see the dominance of the bulk stringy component in the weak coupling limit. All our work has assumed that $$\rho \stackrel{<}{}T_{Dp}.$$ (70) If this condition is violated, the thermal energy of the brane is larger than its rest mass and one would expect our perturbative treatment of the D-brane to break down (see ref. for possible outcomes.) (Note that this bound does not apply directly to orientifolds which are stuck at fixed points.) Since $`T_{Dp}\frac{1}{\widehat{\kappa }}`$, we can see that the amount of thermal energy that can be loaded onto the brane increases as the coupling gets weaker. When this bound is saturated, the stringy contribution to the cosmological equations can be $`𝒪(\widehat{\kappa })`$ while that from the intrinsic tension is $`𝒪(\widehat{\kappa }^2)`$; thus weak coupling is advantageous for string dominance in the cosmology. This rather surprising fact will become clear in the examples of the next section. ## 5 Results: Behaviour of scale factors $`a(t)`$ and $`b(t)`$ In this section we present results for the behaviour of the scale factors $`a`$ and $`b`$. As discussed above, we will for simplicity and definiteness mainly consider 5 dimensional, adiabatic systems. In sections 5.6 and 6 we will consider how it may be possible to relax these restrictions. Our results are obtained by solving Einstein’s equations together with the constraints provided by the Israel conditions. We use the energy momentum tensor derived above appropriate to a primordial Hagedorn epoch. First we will consider the regime L\[$`\gamma `$\] under the assumption that the brane is energy/pressureless and that $`\dot{b}=0`$. We also discuss some constraints imposed by the stability of the extra dimension. We then catalogue more general types of behaviour that can occur for $`\dot{b}0`$ and discuss a semi-realistic example in which $`b`$ is constrained. ### 5.1 Hagedorn inflation; $`\mathrm{\Lambda }_{br}=\mathrm{\Lambda }_{bulk}=\dot{b}=0`$ Let us first impose $`\dot{b}=0`$ simply as an external condition; i.e. the extra dimension is fixed in time. We also set the cosmological constants both on the brane and in the bulk to zero, $`\mathrm{\Lambda }_{br}=\mathrm{\Lambda }_{bulk}=0`$. All other contributions to the energy momentum localized on the brane (for example massless Yang-Mills degrees of freedom or possibly an NL subsystem) have an entropy that is subdominant to the limiting bulk degrees of freedom as long as $`\rho >\rho _c`$ (where $`\rho _c`$ is a critical density of order 1). For the moment we will therefore neglect them and set $`a^{}(t,0)=n^{}(t,0)=0`$. In this discussion the brane at $`y=0`$ is playing no role in determining the evolution of the cosmology; the scale factor $`a_0`$ changes purely as a result of the bulk equations. (We shall see shortly that consistency of the full solutions requires a second brane at $`y=l`$). A. Solutions to the $`T_{55}`$ equation It is simple to solve the 55 equation for the scale factor $`a(t,y)`$ at the origin, $$a(t,0)=a_0(t),$$ (71) for all the L$`[\gamma ]`$ systems. A1. L\[-1\]: In the L\[-1\] systems (the high energy case with windings in all transverse dimensions) we see that $`T_5^5\sqrt{\rho }a^{3/2}`$ (where, again, $`\rho `$ is the local energy density measured with respect to the volume $`V_{}`$) and hence, $$a_0(t)t^{4/3},$$ (72) i.e. power law inflation. A2. L: We find that $`\gamma =0`$ gives us a period of exponential inflation. To find this solution we write $`a_0^2=\mathrm{exp}F`$ and the $`T_5^5`$ equation becomes $$\ddot{F}+\dot{F}^2=\frac{2\widehat{\kappa }^2}{3}\left(\mathrm{log}\rho (0)\frac{3}{2}F\right)$$ which may easily be solved to give $$\frac{a_0(t)}{a_0(0)}=\mathrm{exp}\left(\frac{t^2\widehat{\kappa }^2}{8V_{}^{\frac{3}{2}}}+t\sqrt{\frac{\widehat{\kappa }^2}{6V_{}^{\frac{3}{2}}}\mathrm{log}(\rho (0)e^{3/4})}\right)$$ (73) where $`\rho (0)`$ is the initial density at $`t=0`$. Initially the second term in the exponent dominates and there is exponential inflation. This solution has an automatic end to inflation when $`da_0(t)/dt0`$, i.e. when the two terms in the exponent are roughly of the same magnitude. Hagedorn inflation ends at time $`t\widehat{\kappa }^1\sqrt{V_{}^{\frac{3}{2}}\mathrm{log}\rho (0)}`$ which corresponds to $`\mathrm{log}\frac{a_0^3(t)}{a_0^3(0)}\mathrm{log}\rho (0)`$. During a period of adiabaticity, the entropy and energy density dilute as $`1/a_0^3`$ because the Hagedorn phase is always to a first approximation like pressureless matter ($`\widehat{p}_1\widehat{\rho }`$). Hence we find that the above condition for inflation to end happens when $`\rho (t)1`$ (in string units), just as the system is dropping out of the Hagedorn phase. A3. Other systems: We summarize the energy momentum tensors and cosmological behavior of $`a_0(t)`$ for the remaining systems in table 1, where $`p=3`$ in our ‘toy-model’. In all these cases it should be noted that during a period of adiabaticity, the entropy density dilutes as $`1/a_0^3`$. The reason for the various different types of behaviour is of course due to $`T_5^5`$ which may drop off more slowly than $`1/a_0^3`$ and in the L case drops only logarithmically with the expansion. Note inflationary (superluminal) expansion for $`𝐋`$\[$`\gamma 0]`$. In adiabatic systems the Hagedorn regime and hence the inflationary behaviour we have found eventually come to an end. For a system to be in the Hagedorn regime requires an entropy density higher than the critical Hagedorn density (of order 1 in string units). Below this density the energy momentum tensor and hence the cosmology is governed by the massless relativistic Yang-Mills gas (the gas is present on the brane even in the Hagedorn regime but is subdominant). Thus there is no problem exiting from the inflationary behaviour. In fact the main issue is how long inflation can last. We return to this question later when we discuss how inflation can be sustained. B. Solutions to the $`T_{00}`$ and $`T_{ii}`$ equations We can get an approximate solution for the remaining equations which is valid in all the regions of interest as follows. For small $`y`$ we write $`a(t,y)`$ $`=`$ $`a_0(t)+{\displaystyle \frac{y^2}{2}}a_2(t)+\mathrm{}`$ $`n(t,y)`$ $`=`$ $`1+{\displaystyle \frac{y^2}{2}}n_2(t)+\mathrm{}`$ (74) where henceforth we normalize $`n(t,0)=1`$. First, given our constraints above, the $`05`$ equation is trivially satisfied. As before we solve the 55 equation with $`\widehat{p}_\gamma =V_{}^{3/2}\rho ^{2/3q}`$ and find $$a_0(t)At^q,$$ (75) where $$b_0^{3/2}A^{2/q}=\frac{\widehat{\kappa }^2}{3q(2q1)}.$$ (76) Solving the 00 and $`ii`$ equations with $`T_0^0=\widehat{\kappa }^2\rho /V_{}=\widehat{\kappa }^2\widehat{\rho }`$ we find $`a_2(t)`$ $`=`$ $`b_0^2a_0\left({\displaystyle \frac{q^2}{t^2}}{\displaystyle \frac{\kappa ^2\widehat{\rho }}{3}}\right)`$ $`n_2(t)`$ $`=`$ $`b_0^2\left({\displaystyle \frac{q^22q}{t^2}}+{\displaystyle \frac{\kappa ^2(2\widehat{\rho }+3\widehat{p}_\gamma )}{3}}\right).`$ (77) Phase diagram constraints on the transverse dimension and $`\alpha _{\|\|}1`$ One important aspect of the approximate solution is that it is valid in all Hagedorn dominated regions. Obviously the approximation requires that the $`y^2`$ terms are small. Looking at $`a(t,y)`$ we see that the condition is $`a_2y^2<a_0`$ or $$\widehat{\kappa }^2\widehat{\rho }<\frac{1}{R_{}^2}.$$ (78) Recall from ref. that the Hagedorn regime occupies a small-coupling region between the Yang Mills phase and the black brane phase. Thus if the energy density is too high, thermodynamics is dominated by black branes and our entire perturbative derivation ceases to be valid. In addition there is a ‘holographic’ constraint which is saturated when the entropy is equal to that of a black D$`p`$-brane filling the entire transverse volume (for a review see ref.); $$\widehat{\kappa }^2\widehat{\rho }V_{}\widehat{\kappa }^2\rho V_{}^{\frac{D3p}{D1p}},$$ (79) where as before (see eq.(69)) we assume for convenience that the volume of the compactification to $`D`$ dimensions is $`1`$ so that $`\widehat{\kappa }g_s`$. In the case that the $`9p`$ transverse dimensions are isotropic, this is the same as eq.(78). The present case is (for convenience) taken to be one in which there is only one transverse dimension of size $`R_{}`$ which is much larger than the string scale. We then set $`D=5`$ and $`p=3`$ and again find eq.(78). Thus our approximate solution for the particular case above breaks down precisely where the holographic bound is saturated - when the event horizon of a black-brane fills the transverse dimensions. Note that, depending on the dimensionalities, the entire calculation (i.e. based on Hagedorn regime thermodynamics) can already break down at much lower densities corresponding merely to black brane dominance. Considering for definiteness black-hole type manifolds with planar asympotics, $`S_\beta ^1\times T_L^p`$, it was seen in ref. that when $$\widehat{\kappa }^2\rho 1$$ (80) the dominant component in the thermodynamics (i.e. that phase with the largest entropy) is the horizon of a Schwarzschild-like black brane. Since $`R_{}`$ is larger than the string scale, eq.(80) can be satisfied at a density much lower than the density of the holographic bound above. Adding in the constraint that the Yang-Mills entropy be less than the Hagedorn, we find that the phase diagram restrictions can be summarized as $$\text{min}\{R_{}^{D3p},1\}\widehat{\kappa }^2\rho \widehat{\kappa }^2N^2$$ (81) where $`N`$ is the number of branes situated at $`y=0`$. (The additional requirement of $`\rho \rho _c1`$ explains why we need to be at a small string coupling, $`\widehat{\kappa }^21`$.) One consequence of the phase diagram constraints is that in the full solution $`a`$ is not allowed to vary significantly in the range $`y[l,l]`$, and hence, as promised, $`\alpha _{\|\|}(0)1`$ in all the regions of validity. Note that $`a^{\prime \prime }`$ and $`n^{\prime \prime }`$ must still compensate for $`T_0^0`$ being greater than $`T_5^5`$ (the saddle-point dominance condition) and these two requirements are compatible only at weak coupling. Stabilization and the second brane A nice check of our approximate solution is that it should obey the stabilization condition discussed in for the case where $`T_{05}=0`$ (as we are assuming throughout), $$𝑑y\sqrt{g_{55}}(T_0^0+T_i^i2T_5^5)=0.$$ (82) To see that it indeed does without any further assumptions we need to consider what happens to the solution at the boundaries of the compact space. Assume that the large transverse dimension is compactified on a range $`[l,l]`$ and recall that for convenience we are assuming $`Z_2`$ symmetry under $`yy`$. Thus at $`y=l`$ the solutions have a discontinuity in the derivatives of $`a`$ and $`n`$ which is effectively a second brane. Using the discontinuity condition at $`y=l`$, $`a_l^{}a_{+l}^{}`$ $`=`$ $`{\displaystyle \frac{\widehat{\kappa }^2}{3}}\rho _{br}a(l)b_0`$ $`n_l^{}n_{+l}^{}`$ $`=`$ $`+{\displaystyle \frac{\widehat{\kappa }^2}{3}}(3p_{br}+2\rho _{br})n(l)b_0,`$ (83) we find that this corresponds to additional terms in the energy-momentum tensor given support at $`y=l`$ by a delta function; $`\mathrm{\Delta }T_0^0`$ $`=`$ $`\delta (yl)\rho _{br}=\delta (yl)R_{}\left(\widehat{\rho }3{\displaystyle \frac{q^2}{\widehat{\kappa }^2t^2}}\right)`$ $`\mathrm{\Delta }T_i^i`$ $`=`$ $`\delta (yl)p_{br}=\delta (yl)R_{}\left(\widehat{p}_\gamma {\displaystyle \frac{3q^22q}{\widehat{\kappa }^2t^2}}\right),`$ (84) where $`R_{}=b_0l`$ and where we have neglected terms of order $`\widehat{\kappa }^4\widehat{\rho }^2R_{}^3`$ in accord with our previous discussion. Note that the signs of the brane contributions are the opposite of those coming from the bulk, so that eq.(82) now reduces to $$𝑑y\left(\frac{2q^2q}{t^2}+T_5^5\right)=0.$$ (85) This equation is automatically satisfied since it is the 55 component of Einstein’s equations. In the stabilized system (i.e. with the $`\dot{b}=0`$ ansatz) the brane we find at $`y=l`$ inevitably has a peculiar equation of state that does not resemble a Yang-Mills or cosmological constant on the brane. Nevertheless the stabilization condition should always be satisfied so that, in a full model which includes a stabilization mechanism (given our assumptions for $`a^{}=b^{}=n^{}=0`$ at $`y=0`$ and $`Z_2`$ symmetry), the gravitational sector should behave as if there were a brane with these characteristics located at $`y=l`$. ### 5.2 L\[$`\gamma `$\] with $`\frac{\widehat{\kappa }^2}{12}\mathrm{\Lambda }_{br}^2\mathrm{\Lambda }_{bulk}=0`$ and $`\dot{b}0`$ We now study the case where the extra dimension is not stabilized but there is still no nett cosmological constant. By ‘nett’ we mean that the contribution from $`a^{}`$ and $`n^{}`$ in $`G_5^5`$ cancels the bulk contribution on the RHS of the $`T_5^5`$ equation. The condition for this is $$\mathrm{\Lambda }_{nett}=\frac{\widehat{\kappa }^2}{12}\mathrm{\Lambda }_{br}^2+\mathrm{\Lambda }_{bulk}=0.$$ (86) In this case we could again, by studying the 00 and ii equations, try to find an approximate solution which is valid throughout the Hagedorn regime as we did in the previous subsection. However we would not learn anything further by doing this since we could always independently adjust $`a_2(t)`$ and $`n_2(t)`$ to satisfy the equations. (There is nothing other than the phase diagram constraints to predetermine these parameters.) In addition eq.(82) need not be satisfied (and in fact the integral is $`\dot{b}`$), but one should keep in mind the fact that the full solution again has discontinuities in $`a^{}`$ and $`n^{}`$ at $`y=l`$ so that in general there should be a brane (or something that can compensate in a similar way) situated at $`y=l`$ even when $`\dot{b}0`$. The next subsection shows that by eliminating the second brane at $`y=l`$ one gains an extra constraint on $`a^{\prime \prime }`$ or $`n^{\prime \prime }`$ and hence useful additional information from the $`ii`$ and 00 equations. Therefore we will first focus on the 55 equation. Since $`\alpha 1`$, for $`\gamma =1,\frac{1}{2},\frac{1}{2}`$, the pressure is $$\widehat{p}_\gamma a^{\frac{3\gamma }{\gamma 1}}b^{\frac{3}{2}}.$$ (87) We will now catalogue some possible types of behaviour for $`a`$ with different ansätze for $`b`$. Taking $`\mathrm{\Lambda }_{nett}=0`$, we find a family of power law solutions for $`\gamma =1,\frac{1}{2},\frac{1}{2}`$ of the form $`a_0(t)`$ $``$ $`At^q`$ $`b_0(t)`$ $``$ $`Bt^r`$ (88) where subscript-0 indicates values at $`y=0`$, $`q`$ $`=`$ $`{\displaystyle \frac{\gamma 1}{2\gamma }}({\displaystyle \frac{4}{3}}r),`$ $`A^{\frac{3\gamma }{\gamma 1}}B^{\frac{3}{2}}`$ $`=`$ $`{\displaystyle \frac{\widehat{\kappa }^2}{3q(2q1)}},`$ (89) and where $`r`$ is arbitrary. Note that if we choose $`r=4/3`$ then $`bt^{4/3}`$ and $`a`$ is stationary whatever the value of $`\gamma `$ (the complementary situation to $`b`$ stationary and $`a`$ undergoing power law inflation in eq.(72)). We also find a family of hyperbolic solutions of the form $`{\displaystyle \frac{a_0(t)}{a_0(0)}}`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{sinh}2C(t+t_1)}{\mathrm{sinh}2Ct_1}}\right)^{\frac{1}{2}}`$ $`{\displaystyle \frac{b_0(t)}{b_0(0)}}`$ $`=`$ $`\{\begin{array}{cc}\left(\frac{\mathrm{sinh}2C(t+t_1)}{\mathrm{sinh}2Ct_1}\right)^{\frac{\gamma }{1\gamma }}\hfill & \mathrm{if}\gamma =1,\frac{1}{2},\frac{1}{2}\hfill \\ & \\ \mathrm{log}\left(\frac{\mathrm{sinh}2C(t+t_1)}{\mathrm{sinh}2Ct_1}\right)\hfill & \mathrm{if}\gamma =0\hfill \\ & \\ \mathrm{exp}\left(\frac{\mathrm{sinh}2C(t+t_1)}{\mathrm{sinh}2Ct_1}\right)\hfill & \mathrm{if}\gamma =1.\hfill \end{array}`$ (91) where $`C`$ and $`t_1`$ are defined by $`C^2`$ $`=`$ $`{\displaystyle \frac{\widehat{\kappa }^2}{6}}\widehat{p}_\gamma (0)`$ $`{\displaystyle \frac{\dot{a}_0(0)}{a_0(0)}}`$ $`=`$ $`C\mathrm{coth}2Ct_1,`$ (92) $`\widehat{p}_\gamma `$ is given by eq.(3.4), and where e.g. $`\widehat{p}_\gamma (0),\rho (0)\widehat{p}_\gamma (t=0),\rho (t=0)`$. Note that the solutions with $`t_1\pm \mathrm{}`$ describe the scale factor exponentially increasing or decreasing; e.g. for $`\gamma =1,\frac{1}{2},\frac{1}{2}`$ in this limit we have $$\begin{array}{cccc}\frac{a_0(t)}{a_0(0)}\hfill & \hfill & \mathrm{exp}\left(\pm Ct\right)\hfill & \text{ , }t_1\pm \mathrm{}\hfill \\ \frac{b_0(t)}{b_0(0)}\hfill & \hfill & \mathrm{exp}\left(\pm \frac{2\gamma }{1\gamma }Ct\right)\hfill & \text{ , }t_1\pm \mathrm{}.\hfill \end{array}$$ In the hyperbolic solutions $`T_5^5`$ is constant in time (see eq. (87). Both expansion and collapse of $`a`$ are possible, since $`b`$ compensates appropriately (in such a way as to keep $`T_5^5`$ constant). It seems rather paradoxical that the compensation can go either way depending on the sign of $`\gamma `$. For example, in the $`\gamma =1/2`$ solution of eq.(5.2), if $`a(t)`$ is expanding then, in order to keep $`T_5^5`$ constant, $`b`$ has to expand exponentially as well. The solution in the previous subsection for the $`\dot{b}=0`$, $`\gamma =0`$ system (section 5.1.A2, eq.(73)) is similar, although there is a slight difference because there $`b`$ is constant, resulting in the $`e^{t^2}`$ ‘exit’ term. We should also remember (bearing in mind our previous discussion) that if adiabaticity is assumed then the density $`T_0^0`$ is always diluted exponentially fast as $`T_0^0\mathrm{exp}\left(Ct\frac{3\gamma }{1\gamma }\right)`$. Finally we should add that which, if any, of these solutions is appropriate depends on the initial value we choose for $`\dot{b}_0`$. ### 5.3 L\[$`\gamma `$\] with $`\frac{\widehat{\kappa }^2}{12}\mathrm{\Lambda }_{br}^2\mathrm{\Lambda }_{bulk}0`$ and $`\dot{b}0`$ Next we consider the case of additional cosmological constants in the brane and bulk. We transfer the $`a^{}`$ and $`n^{}`$ terms to the RHS of the 55 equation and write $$T_{5eff}^5=\left(\widehat{p}_\gamma +\mathrm{\Lambda }_{nett}\right),$$ (93) where $$\mathrm{\Lambda }_{nett}=\mathrm{\Lambda }_{bulk}\frac{\widehat{\kappa }^2}{12}\mathrm{\Lambda }_{br}^2.$$ (94) First we can see (rather trivially) that the hyperbolic solutions of the previous subsection in eq.(5.2) still exist if we modify $`C`$ to $`C^{}`$, where $$C^2=\frac{\widehat{\kappa }^2}{6}\left(\mathrm{\Lambda }_{nett}+\widehat{p}_\gamma (0)\right),$$ (95) and provided $`C^2>0`$. Not surprisingly we recover the usual cosmological constant driven inflation when we set $`\rho =0`$. When $`C^2<0`$ (which is the case for $`\mathrm{\Lambda }_{nett}<0`$ and $`\|\mathrm{\Lambda }_{nett}\|>\widehat{p}_\gamma (0)`$), we find a singular collapsing solution for $`a`$; $`{\displaystyle \frac{a_0(t)}{a_0(0)}}`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{sin}2\|C^{}\|(t+t_1)}{\mathrm{sin}2\|C^{}\|t_1}}\right)^{\frac{1}{2}}`$ $`{\displaystyle \frac{b_0(t)}{b_0(0)}}`$ $`=`$ $`\{\begin{array}{cc}\left(\frac{\mathrm{sin}2\|C^{}\|(t+t_1)}{\mathrm{sin}2\|C^{}\|t_1}\right)^{\frac{\gamma }{1\gamma }}\hfill & \mathrm{if}\gamma =1,\frac{1}{2},\frac{1}{2}\hfill \\ & \\ \mathrm{log}\left(\frac{\mathrm{sin}2\|C^{}\|(t+t_1)}{\mathrm{sin}2\|C^{}\|t_1}\right)\hfill & \mathrm{if}\gamma =0\hfill \\ & \\ \mathrm{exp}\left(\frac{\mathrm{sin}2\|C^{}\|(t+t_1)}{\mathrm{sin}2\|C^{}\|t_1}\right)\hfill & \mathrm{if}\gamma =1.\hfill \end{array}`$ (97) where $`t_1`$ is defined by $$\frac{\dot{a}_0(0)}{a_0(0)}=\|C^{}\|\mathrm{cot}2\|C^{}\|t_1.$$ (98) For the power law solutions eq.(5.2), we note that the Hagedorn contribution to $`T_5^5`$ varies as $`1/t^2`$ and therefore decreases in time compared to the constant $`\mathrm{\Lambda }_{nett}`$ contributions. Thus we can consider the 55 equation in the two limits that $`T_{5eff}^5`$ in eq.(93) is dominated either by the cosmological constant term or by the Hagedorn term. For example, if $`\gamma =1,\frac{1}{2},\frac{1}{2}`$, we have $`\begin{array}{ccc}a_0(t)\hfill & \hfill & A(tt_0)^q\hfill \\ b_0(t)\hfill & \hfill & B(tt_0)^r\hfill \end{array}\}`$ $`T_{5eff}^5\widehat{p}_\gamma `$ (101) $`\begin{array}{ccc}\frac{a_0(t)}{a_0(0)}\hfill & =\hfill & \left(\frac{\mathrm{sinh}2C^{}(t+t_1)}{\mathrm{sinh}2C^{}t_1}\right)^{\frac{1}{2}}\hfill \\ \frac{b_0(t)}{b_0(0)}\hfill & =\hfill & \left(\frac{\mathrm{sinh}2C^{}(t+t_1)}{\mathrm{sinh}2C^{}t_1}\right)^{\frac{\gamma }{1\gamma }}\hfill \end{array}\}`$ $`\|T_{5eff}^5\|\|\mathrm{\Lambda }_{nett}\|\|\widehat{p}_\gamma \|,C^2>0`$ (104) $`\begin{array}{ccc}\frac{a_0(t)}{a_0(0)}\hfill & =\hfill & \left(\frac{\mathrm{sin}2\|C^{}\|(t+t_1)}{\mathrm{sin}2\|C^{}\|t_1}\right)^{\frac{1}{2}}\hfill \\ \frac{b_0(t)}{b_0(0)}\hfill & =\hfill & \left(\frac{\mathrm{sin}2\|C^{}\|(t+t_1)}{\mathrm{sin}2\|C^{}\|t_1}\right)^{\frac{\gamma }{1\gamma }}\hfill \end{array}\}`$ $`\|T_{5eff}^5\|\|\mathrm{\Lambda }_{nett}\|\|\widehat{p}_\gamma \|,C^2<0`$ (107) where $`t_0`$ is the time at which $`a`$ is released from its small initial value. ### 5.4 Bouncing universe with softened singularity If $`\mathrm{\Lambda }_{nett}<0`$ the above solutions admit an oscillating universe. For example if $`r<4/3`$, we find a universe in which a negative cosmological constant causes $`a`$ to follow the sin curve of eq.(107) which reaches a maximum at $`\|C^{}\|(t+t_1)=\pi /4`$ before heading towards zero. However, unlike the case of pure negative cosmological constant in ordinary FRW, we do not hit a singularity because, when the scale factor $`a`$ becomes small, $`\mathrm{\Lambda }_{nett}+\widehat{p}_\gamma `$ in eq. (93) becomes positive and $`T_5^5`$ is dominated by $`\widehat{p}_\gamma `$. Instead of hitting a singularity, $`a`$ reaches some minimum value and rebounds upwards on the power law solution of eq.(101) (with $`t_0`$ now marking the time of the rebound, $`\dot{a}(t_0,0)=0`$). The scale factor $`a`$ follows the power law curve until the Hagedorn contribution to $`T_5^5`$ (which varies as $`1/(tt_0)^2`$) drops below the cosmological constant piece, at which point we pick up the sin curve of eq.(107) and have completed a single oscillation. The oscillation is around the completely static solution (with $`\dot{a}=\ddot{a}=\dot{b}=0`$) where the pressure $`\widehat{p}_\gamma `$ is exactly compensated by a negative cosmological constant, $$\widehat{p}_\gamma =\mathrm{\Lambda }_{nett}.$$ (108) The behaviour of $`b`$ depends on the value of $`\gamma `$. For $`\gamma =\frac{1}{2}`$ we find that $`b`$ contracts and expands with $`a`$. In this case, at small $`a`$ we must have a power law solution with $`r>4/3`$ so that the contracting/expanding power law solution matches onto a contracting/expanding $`\mathrm{sin}`$ solution. On the other hand for $`\gamma =1`$ or $`\frac{1}{2}`$, $`b`$ contracts and expands out of phase with $`a`$, and we require $`0<r<\frac{4}{3}`$ at small $`a`$. Clearly if, at small $`a`$, $`r0`$ then from eq.(101) $`b`$ is almost static. We stress that this singularity smoothing behaviour may or may not occur depending on the initial conditions (e.g. $`\dot{b}`$) and additional constraints on the system. For example, it does occur in numerical solutions with $`\dot{b}=0`$. It is encouraging to see that one contribution of stringy physics to the cosmology can be to soften the singularity that would otherwise appear at $`a=0`$. This is a familiar aspect of string theory, but an appealing feature of the Hagedorn phase is that we can find it purely perturbatively. ### 5.5 Full solution with an example of a ‘physical’ brane In the previous subsections we discussed the cosmological behaviour for various ansätze. However, apart from the $`\dot{b}=0`$ constraint (which can be motivated by some unknown stabilization mechanism) the other ansätze are unmotivated, and we have still to show that the behaviour can arise in ‘realistic’ physical systems. The reason for this unwanted freedom is that the system of equations is underconstrained – there are four parameters ($`a_0,b_0,a_2,n_2`$) but only three independent equation (00, $`ii`$$`55`$). Therefore, we now consider an example in which we impose an additional constraint coming from a particular choice of equation of state on the brane. This allows us to find a full solution, in which the behaviour observed in the previous subsection arises in certain limits. As we saw when we discussed stabilization, there is a discontinuity in $`a^{}`$ and $`n^{}`$ which corresponds to a brane located at $`y=l`$. Generally, the equation of state for this brane will look rather unrealistic (as indeed it does for the static $`\dot{b}=0`$ solution) and we are forced to account for it by invoking unknown contributions from the gravitational sector. Our first assumption is therefore that there is no second brane (or rather, in view of the need to conserve Ramond-Ramond charge, that only the brane at $`y=0`$ has any cosmological effect). With this assumption, continuity at $`y=l`$ requires $`a_l^{}=n_l^{}=0`$ and we must drop the approximation that the single brane at $`y=0`$ (our would-be universe) is ‘empty’ (in the sense that $`a_0^{}=n_0^{}=0`$) and consider its energy and momentum as well. The Israel conditions tell us that this energy density is given by $`3[a^{}/a]_0`$ $`=`$ $`\widehat{\kappa }^2b_0\rho _{br}`$ $`3[n^{}/n]_0`$ $`=`$ $`\widehat{\kappa }^2b_0(2\rho _{br}+3p_{br}).`$ (109) We now make an additional assumption about the equation of state of the brane at $`y=0`$; we assume it is also Hagedorn-like, i.e. as well as the bulk L\[$`\gamma `$\] degrees of freedom there is an additional localized Hagedorn component with excitations near the brane. Even though we do not understand the thermodynamic behaviour of this subsystem, we can be fairly confident that its energy momentum tensor is like pressureless matter (as all the Hagedorn systems are to a first approximation – they all have $`\|\rho _{br}\|\|p_{br}\|`$). Note that, to simplify this example, we will neglect the brane tension contribution which acts like $`\mathrm{\Lambda }_{br}`$. For $`p_{br}0`$ it is therefore sufficient to impose the relation $$n^{}/n=2a^{}/a$$ (110) for all $`y`$. Integrated, this gives $$n=\frac{a_0(t)^2}{a(t)^2}.$$ (111) Together with $`a_l^{}=n_l^{}=0`$ this means that the approximate solutions for $`a(t,y)`$ and $`n(t,y)`$ are going to be of the form $`a(t,y)`$ $`=`$ $`a_0(t)+a_1(t)\left(\|y\|{\displaystyle \frac{y^2}{2l}}\right)`$ $`n(t,y)`$ $`=`$ $`12{\displaystyle \frac{a_1(t)}{a_0(t)}}\left(\|y\|{\displaystyle \frac{y^2}{2l}}\right).`$ (112) The 05 equation may easily be solved and we find $$b\lambda (y)=a^{}a^2,$$ (113) where $`\lambda (y)`$ is a constant of integration. Note that this and eq.(5.5) imply that $$\frac{\widehat{\kappa }^2}{3}\rho _{br}=\frac{\lambda _0}{a_0(t)^3}$$ (114) so that the brane density is indeed diluted by the expansion of the scale factor $`a`$ like pressureless matter. We now apply this to the 00, $`ii`$ and 55 Einstein equations which now are more constrained by the continuity at $`y=l`$ and by our choice of equation of state. Somewhat fortuitously the condition in eq.(110) gives a cancellation of the $`a^{\prime \prime }`$ and $`n^{\prime \prime }`$ terms and the equations at $`y=0`$ reduce to $`{\displaystyle \frac{\dot{a}_0}{a_0}}\left({\displaystyle \frac{\dot{a}_0}{a_0}}+{\displaystyle \frac{\dot{b}_0}{b_0}}\right){\displaystyle \frac{\lambda _0^2}{a_0^6}}{\displaystyle \frac{\lambda _0^{}}{a_0^3b_0}}`$ $`=`$ $`{\displaystyle \frac{\widehat{\kappa }^2}{3}}\widehat{\rho }`$ $`3{\displaystyle \frac{\lambda _0^2}{a_0^6}}+{\displaystyle \frac{\dot{a}_0}{a_0}}\left({\displaystyle \frac{\dot{a}_0}{a_0}}+2{\displaystyle \frac{\dot{b}_0}{b_0}}\right)+2{\displaystyle \frac{\ddot{a}_0}{a_0}}+{\displaystyle \frac{\ddot{b}_0}{b_0}}`$ $`=`$ $`\widehat{\kappa }^2\widehat{p}_\gamma `$ $`{\displaystyle \frac{\lambda _0^2}{a_0^6}}+\left({\displaystyle \frac{\dot{a}_0}{a_0}}\right)^2+{\displaystyle \frac{\ddot{a}_0}{a_0}}`$ $`=`$ $`+{\displaystyle \frac{\widehat{\kappa }^2}{3}}\widehat{p}_\gamma .`$ (115) Dominance of the $`\widehat{p}`$ pressure term in the 55 equation is only possible if $$\lambda _0^2/a_0^6(\dot{a}/a)^2\|\widehat{\kappa ^2}\widehat{p}_\gamma \|\widehat{\kappa ^2}\widehat{\rho }$$ (116) and so the 00 equation must be dominated by the $`\lambda ^{}`$ term. Solving gives $$\frac{\lambda (y)}{\lambda _0}1\frac{\rho }{\rho _{br}}\frac{y}{l}.$$ (117) Consistency of our solutions therefore requires that $`\rho _{br}\stackrel{>}{}\rho `$ and hence by eq.(116) $$\widehat{\kappa }^2\rho \frac{\widehat{p}_\gamma }{\rho }\frac{1}{V_{}}.$$ (118) As for the previous solutions, this corresponds to the holographic constraint in the phase diagram, and implies small coupling. Thus we again observe that, because only $`\widehat{\kappa }^4\rho _{br}^2`$ appears in the Einstein equations, the cosmology can be dominated by the ‘bulk’ energy-density even when $`\rho <\rho _{br}`$, provided that the coupling is small. For the L\[-1\] systems the $`ii`$ and $`55`$ equations may be solved by substituting $`d=\frac{a_0(t)b_0(t)}{a_0(0)b_0(0)}`$ and $`c=\frac{a_0(t)^2}{a_0(0)^2}`$ so that $$\frac{\ddot{c}}{c}=\frac{2\widehat{\kappa }^2}{3}\widehat{p}_1;\frac{\ddot{d}}{d}=\frac{4\widehat{\kappa }^2}{3}\widehat{p}_1=\frac{4\widehat{\kappa }^2}{3}\widehat{p}_1(0)d^{3/2}.$$ (119) Defining $$w=\dot{d}_0(0)^2\frac{16\widehat{\kappa }^2\widehat{p}_1(0)}{3}\sqrt{d_0(0)}$$ we find power law and hyperbolic behaviour emerging in different limits. When $$tt_0\frac{w^{\frac{3}{2}}}{\widehat{\kappa }^4\widehat{p}_1(0)^2}$$ (120) we have power law behaviour (eqn.5.2) with $$\left(\frac{4}{3}r\right)\left(\frac{5}{3}2r\right)=\frac{1}{\widehat{\kappa }^2\widehat{p}_1(0)}.$$ (121) So for large $`\widehat{p}_1`$, $`r4/3`$, $`5/6`$. On the other hand when $$tt_0\frac{w^{\frac{3}{2}}}{\widehat{\kappa }^4\widehat{p}_1(0)^2}$$ (122) we have $`\sqrt{d}w/\widehat{\kappa }^4\widehat{p}_1(0)^2\sqrt{d(0)}`$ and hence we find the hyperbolic solution of eq.(104) with expanding $`a`$ and collapsing $`b`$. ### 5.6 More dimensions Finally we briefly remark on the extension to higher numbers of dimensions. In type I,IIA/B models the total number of dimensions is $`D=10`$ and the energy momentum tensor is therefore of the form $`T_\mu ^\nu =diag(T_0^0,\widehat{p}_\gamma ,\widehat{p}_\gamma \mathrm{},\widehat{p}_\gamma ,\widehat{p}_\gamma \mathrm{},0\mathrm{},0)`$. A D$`p`$-brane has $`9p`$ large Dirichlet directions giving $$1\gamma \text{min}\{1,\frac{7p}{2}\}.$$ (123) As we discussed in the introduction, it is always true that $`\gamma 1`$ because for $`p<5`$, the only possible limiting system has excited winding modes and is the $`\gamma =1`$ system. A diagonal metric with universal scale factors is consistent with the above form of $`T_\mu ^\nu `$. We choose the metric to be a function of a single transverse dimension which we label $`y`$, and add $`D2p`$ additional transverse dimensions labelled $`y_m`$, ($`m=p+2,..D`$); $$ds^2=n^2dt^2+a(t,y)^2dx_i^2+b(t,y)^2dy^2+d(t,y)^2dy_mdy_m.$$ (124) This metric gives 4 independent equations (00,$`ii`$,$`55`$,$`mm`$) for the higher dimensional systems. Most importantly, $`a^{\prime \prime }`$, $`n^{\prime \prime }`$ and $`d^{\prime \prime }`$ do not appear in the 55 components of $`G_\mu ^\nu `$ . Consequently the 55 equation may again be treated separately from the $`00`$, $`ii`$, $`mm`$ equations. The latter 3 equations may always be solved independently (as in the $`D=5`$ case studied in detail in this paper) using the three unknowns $`a^{\prime \prime }`$, $`n^{\prime \prime }`$ and $`d^{\prime \prime }`$. Thus, although we will not discuss the higher dimensional case in detail here, we can immediately see on dimensional grounds that the solutions to the $`55`$ equations with $`\dot{d}=\dot{b}=0`$ must be as shown in table 1 (with the appropriate value of $`p`$). An amusing aspect of this is that power law inflation in L\[-1\] systems (i.e. the universal high energy system) requires $`p3`$ regardless of the total number of dimensions. ## 6 Sustaining inflation and solving the Horizon Problem One of the most striking features of the cosmological solutions we have found is that they automatically predict a superluminal growth of the scale factor which is usually an ingredient of standard inflation. The remaining point that needs to be addressed is that our analysis assumed adiabaticity, and in order to realistically solve any cosmological problems, this assumption (as we will see) almost certainly requires modification. In this section we discuss solving the horizon problem beginning with the causality condition and then addressing the issue of adiabaticity. We then discuss how nonabiabatic effects may contribute to a sustained period of inflation. ### 6.1 Causality Condition: An estimate of the size of the observable universe today is given by the distance light could travel between photon decoupling and now, $$d_{obs}a(t_o)_{t_{dec}}^{t_o}𝑑t/a(t).$$ (125) Note that $`d_{obs}=𝒪(1)\times (t_ot_{dec})`$ for $`at^p`$ and $`p=𝒪(1)`$ between $`t_{dec}`$ and $`t_o`$. Here $`a`$ is the scale factor. We can compare the comoving size of the observable universe to the comoving size of a causally connected region at some earlier time $`t_p`$: $$d_{hor}(t_p)/a(t_p)=_0^{t_p}\frac{dt}{a(t)}.$$ (126) The observable universe today fits inside a causally connected region at $`t_p`$ if $$\frac{d_{hor}(t_p)}{a(t_p)}\frac{d_{obs}}{a(t_o)}.$$ (127) Here, subscript-$`o`$ refers to today and subscript-$`p`$ indicates some primordial time before the inflationary period of interest. If condition eq.(127) is met, then the horizon size at $`t_p`$ (before nucleosynthesis) is large enough to allow for a causal explanation of the smoothness of the universe today. Note that more creative explanations of large scale smoothness may not involve comparing these two patches. For example, in the context of the brane scenarios, one might imagine that two regions of our observable universe which seem to be causally disconnected might in fact have talked to each other because of a geodesic between them that went off our brane, into the bulk, and then back onto our brane at some distant point, as demonstrated by Chung and Freese . In the remainder of our discussion here we restrict ourselves to the case where eq.(127) is relevant; this is certainly the case for all the brane and boundary inflation models considered to date. For power law expansion of the scale factor both before $`t_p`$ and after $`t_{dec}`$ (which may or may not be the case), we can take $`tH^1`$ during these periods. The causality condition eq.(127) then becomes $$\frac{1}{a_pH_p}\frac{1}{H_oa_o},$$ (128) or, equivalently, $$\frac{t_p}{a_p}\frac{t_{today}}{a_{today}}.$$ (129) We take the Hubble constant today to be $$H_o=\alpha _o^{1/2}T_o^2/m_{pl}(t_o),$$ (130) where $`\alpha _o=(8\pi /3)(\pi ^2/30)g_{}(t_o)\eta _o`$, $`g_{}`$ is the number of relativistic degrees of freedom and $`\eta (t_o)10^410^5`$ is the ratio today of the energy density in matter to that in radiation. From eq.(128) we can then see that accelerated expansion of the scale factor with $`\ddot{a}>0`$ is required to solve the horizon problem. As we have shown throughout the paper , such accelerated expansion can easily occur during the Hagedorn regime. We must here assume that, within an intitial horizon patch of size $`a_p`$, there has been enough smoothing that it is sensible to talk about a single value of the temperature and entropy density. This same assumption must be made in every inflation model. ### 6.2 The issue of non-adiabaticity Adiabatic Hagedorn inflating systems can in principle already solve the horizon problem (as opposed to standard inflation in which nonadiabaticity is required). These systems could start with a large initial entropy due to the proximity of the initial temperature to the Hagedorn temperature. Subsequently, the entropy remains constant as the temperature drops and the scale factor $`a`$ grows by a large amount. However in practice it is difficult to solve the horizon problem assuming adiabaticity: it is not possible to introduce a priori the entire entropy of the observable universe (i.e. $`S10^{88}`$) onto a primordial brane of order the string scale without either destroying the brane or having an unreasonably small string coupling (i.e. $`g_s10^{88}`$!). Consider for example the $`𝐋[1]`$ system with $`at^{4/3}`$. Then from the first entry in Table I, we see that the entropy on a world-volume $`a_p^3`$ is $$S_p\beta _Ha_p^3\rho \beta _Ha_p^3\frac{1}{(\beta \beta _H)^2}\frac{1}{V_{}}.$$ (131) We will assume that the entropy on the 3-brane today is $`S(a_{today}T_{today})^3`$. We obtain the temperature/time relation during the Hagedorn phase from $`\rho a^3t^4`$ by equating this expression for $`\rho `$ with that from the first entry in Table I. We find $`t_p=V_{}^{1/4}\sqrt{\beta _p\beta _H}`$. Eq.(127) can therefore be satisfied if $$\frac{S_0}{S_p}>\frac{m_s}{T_{today}}\left(V_{}(\beta _p\beta _H)^2\right)^{1/4}=\frac{m_s}{T_{today}}\rho _p^{1/4}.$$ (132) Hence even an adiabatic system here can in principle solve the horizon problem: provided that the primordial density is large enough (or the primordial temperature sufficiently close to $`T_H`$) the horizon problem is solved even if $`S_0=S_p`$. Indeed, we can see from eqn.(131) that the initial entropy density can be as large as we like provided that we are prepared to accept a $`\beta `$ that is arbitrarily close to $`\beta _H`$. The serious practical difficulty arises however when we consider the stability bound in eq.(70) since it clearly implies that non-perturbative effects will be important unless the string coupling is fantastically weak. Indeed, generally if one assumes adiabaticity in attempting to solve the horizon problem, one requires the entire entropy of the observable universe to be present in string excitations at $`t=t_p`$. Hence, if the volumes are initially of order the string scale, then the energy density must be enormous, $`\rho (t_p)10^{88}`$ (e.g. in eqn.132). After 60 $`e`$-folds of inflation $`\rho (t)`$ falls below $`\rho _c`$, the critical density which is needed to be in the Hagedorn phase, and the Yang-Mills phase takes over. However, for such a large initial value of the energy density, brane stability at $`t=t_p`$ (i.e. the constraint in eq.(70)) requires an initial value of $`g_s<10^{88}`$. Because of this it is worth looking more critically at the supposed adiabaticity. This is a crucial assumption that almost certainly requires modification in a de Sitter background since the latter possesses a horizon with its own associated entropy. Hence, even if adiabaticity were the correct criterion (which it is not), it would only be meaningful for the coupled string/de Sitter system since there is a backreaction of the horizon on the string gas. The upshot for strings in a de Sitter background is that they are unstable to fluctuations and that this instability can sustain a period of de Sitter inflation . This phenomenon is well known for strings once they are in the de Sitter-like phase<sup>7</sup><sup>7</sup>7We qualify ‘de Sitter’ because the exponential solutions we have do not possess the full O(1,4) de Sitter symmetry.. However the missing ingredient that the present study adds is an explanation for how the universe enters a de Sitter like phase in the first place. Our mechanism gets the universe into the locally de Sitter phase, whereupon the mechanism of refs. keeps it there without having to assume adiabaticity. To achieve a sustained period of inflation therefore, we could just invoke the findings of refs.. However for the remainder of this section we briefly elaborate on this property of strings, using purely heuristic arguments, in order to indicate a possible direction for future study. Let us return to the density of states and consider instead a truly static metric in which the relevant Killing field is the time-translation, $$K^\mu \frac{}{t}.$$ (133) De Sitter space is an example of a static metric; in an appropriate coordinate system the metric is $$ds^2=dt^2(1H^2r^2)+\frac{1}{1H^2r^2}dr^2+r^2d\mathrm{\Omega }^2$$ (134) and the Killing field has components $`(1,0,0,0)`$. We can therefore discuss thermodynamics by compactifying on an imaginary time coordinate with $`itit+\beta `$. A suitable modification to discuss the present case is to add an additional $`y`$ coordinate with a metric $`g_{55}=b(y)^2`$ which is time independent in this coordinate system and also independent of $`r`$; $$ds^2=(1H^2r^2)dt^2+\frac{1}{1H^2r^2}dr^2+r^2d\mathrm{\Omega }^2+b(y)^2dy^2.$$ (135) Under this assumption each 4 dimensional leaf of the foliated space is a de Sitter space and the geometry unambiguously fixes the Hawking temperature of the Horizon, $`2\pi /H=\beta `$. It is well established (see for example ref. and references therein) that the combined (geometric plus thermal) system obeys a generalized second law which is that the total entropy, $$S=S_{matter}+S_{horizon}=S_{matter}+\frac{1}{4\widehat{\kappa }^2}A_{horizon},$$ (136) obeys $`\delta S>0`$, where $`A_{horizon}`$ is the area of the horizon. The constant-time surface $`\mathrm{\Sigma }`$ extends to the horizon $`r1/H=\beta /2\pi `$ so that $$V_{\|\|}=Vol(S_{d_{\|\|}})r^{d_{\|\|}}\beta ^{d_{\|\|}};A_{horizon}=Vol(S_{d_{\|\|}1})r^{d_{\|\|}1}\beta ^{d_{\|\|}1}$$ (137) where $`Vol(S_{d_{\|\|}})`$ is the volume of the unit sphere in $`d_{\|\|}`$ dimensions. Thus, assuming equilibrium at the temperature $`\beta ^1`$, the total entropy of an L\[-1\] string gas in a de Sitter background is of the form $$S\frac{\beta ^{d_{\|\|}}}{V_{}(\beta \beta _H)^2}+\frac{\beta ^{d_{\|\|}1}}{4\pi \widehat{\kappa }^2}.$$ (138) This function has a minimum at $$T_{crit}=T_H\left(1𝒪\left((\widehat{\kappa }^2/V_{})^{1/3}\right)\right).$$ (139) Below $`T_{crit}`$ fluctuations tend to drive the universe towards low temperature, i.e. out of the de Sitter phase<sup>8</sup><sup>8</sup>8We recognize that, in making this heuristic argument, we are on rather thin ice since we are simultaneously assuming a string coupling which is small enough for our approximations to be correct, but large enough to maintain equilibrium between the different degrees of freedom.. The horizon has a negative specific heat which becomes relatively larger as the temperature drops so that equilibrium can be maintained (the condition being $`\|C_V\|>C_{V+}`$). Above $`T_{crit}`$ in the presence of a thermal gas of strings the flow is in the opposite direction towards an asymptotically limiting Hubble constant, $`H_{asymp}=2\pi T_H`$. Thus the generalized second law indicates that fluctuations in the geometry tend to drive the universe even further into the de Sitter phase provided that the Hubble constant is greater than the critical value. The entropy that the horizon loses when the Hubble constant is increased is more than offset by the huge increase in string entropy. Moreover, since the specific heat of the string gas becomes larger when the temperature increases, inevitably equilibrium is lost and energy flow to the strings becomes apparently limitless. There is a slight difference between this picture and the results of ref. which we should comment on. In ref., the density never goes above a critical density, $`\rho _c^{}`$, and indeed asymptotes to it from below as one goes back in time. In our case, on the other hand, we must have $`\rho >\rho _c1`$ for the calculation to be valid (i.e. to be in the Hagedorn regime). Indeed in the present paper the energy density diverges as we approach the Hagedorn temperature. The difference is because ref. was concerned with stretched strings and introduced a cut-off in the momentum integral in order to find the fractal dimension of the string. This ‘coarse graining’ put an artificial upper limit on the amount of string that can be packed into the volume. In fact they found the maximum fractal dimension (i.e. 2) only at the critical density $`\rho _c^{}`$. A fractal dimension of 2 is typical of the random walk behaviour which exists once Hagedorn behaviour sets in. So the coarse graining in ref. effectively removed the Hagedorn behaviour and consequently ref. found that $`\rho =\rho _c^{}`$ gave $`\beta =\beta _H`$. Conversely, the present paper begins in the regime $`\rho >\rho _c`$ where the calculations of ref. end. ## 7 Conclusion and discussion In this paper we have studied the possible cosmological implications of the Hagedorn regime of open strings on D-branes in the weak coupling limit. Our main result in sections 2-5 is that, due to the non-extensive dependence of the free energy on the volumes, a gas of open strings can exhibit negative pressure leading naturally to a period of power law or even exponential inflation – Hagedorn inflation. We also find that the open string gas can dominate the cosmological evolution at weak coupling even though the D-brane tension becomes large in this limit. Hagedorn inflation also has a natural exit since any significant cooling can cause a change in the thermodynamics if winding modes become quenched or if the density drops below the critical density, $`\rho _c1`$, needed for the entropy of the Hagedorn phase to be dominant. Such a cooling can be caused by a sudden adiabatic increase in the transverse radius or by the inflation itself. We find this ‘easy-exit’ feature of open-string Hagedorn inflation to be of its most appealing features. In addition, we found that a small but negative cosmological constant, can cause the universe to enter a stable but oscillating phase. The effect of the Hagedorn phase is to soften the singular behaviour associated with the bounce. The most striking aspect of our discussion is probably the existence of negative pressure. One might therefore ask how general a feature this is expected to be. By T-dualizing we argued that we can put the negative pressure down to the fact that in any particular direction, the gravitational degrees of freedom have both Kaluza-Klein modes and winding modes whereas the open strings (which are dominant in the entropy) have only one or the other. Hence we expect negative pressure to be possible whenever there are large space-filling modes that dominate the entropy. In section 6 we speculated on how the inflation might be sustained through the well known phenomenon of string instability in a de Sitter background. Hagedorn inflation may be thought of as a first example in the search for alternatives to the cosmological constant within the framework of string/brane systems. One possible direction for further study in this area is connected with the fact that we have throughout been taking the string coupling to be weak enough so that the brane tension does not play an important role in the cosmology. It is therefore interesting to ask if new cosmological effects might arise from large scale fluctuations in the branes themselves . The arguments of section 6 indicate that if it does then the brane driven inflation may be qualitatively different to the string driven inflation discussed here. This is because open strings are one dimensional objects with $`SE+const\sqrt{E}`$ whereas fluctuating $`p`$-branes have an entropy $$SE^{\frac{2p}{p+1}}.$$ (140) On calculating $`\beta =S/E`$ we see that $`p1`$ branes do not have the divergent behaviour which is the defining feature of strings. The thermodynamics of these objects has been the subject of much study and there is probably more to be learnt here. Finally, an interesting connected issue which we did not discuss is related to the effect of brane melting discussed in ref., in which the non-perturbative aspect of D-brane thermal production must be taken into account. At the present time it is hard to make any quantitative estimate of these effects on the cosmological solutions we have been discussing here, but we hope to be able to address this question in future work. ### Acknowledgements We thank Dan Chung, Cedric Deffayet, Emilian Dudas, Keith Olive, Geraldine Servant, Carlos Savoy and Richard Woodard for discussions. S.A.A. and I.I.K. thank Jose Barbón and Eliezer Rabinovici for a previous collaboration and discussions concerning this work. S.A.A. thanks the C.E.A. Saclay for support. K.F. acknowledges support from the Department of Energy through a grant to the University of Michigan. K.F. thanks CERN in Geneva, Switzerland and the Max Planck Institut fuer Physik in Munich, Germany for hospitality during her stay. I.I.K. is supported in part by PPARC rolling grant PPA/G/O/1998/00567, the EC TMR grant FMRX-CT-96-0090 and by the INTAS grant RFBR - 950567.
warning/0005/astro-ph0005446.html	ar5iv	text	# Neutrino Degeneracy and Decoupling: New Limits from Primordial Nucleosynthesis and the Cosmic Microwave Background ## 1 INTRODUCTION The universe appears to be charge neutral to very high precision (Lyttleton & Bondi (1959)). Hence, any universal net lepton number beyond the net baryon number must reside entirely within the neutrino sector. Since the present relic neutrino number asymmetry is not directly observable there is no firm experimental basis for postulating that the lepton number for each species is zero. It is therefore possible for the individual lepton numbers $`L_l`$ of the universe to be large compared to the baryon number of the universe, $`B`$, while the net total lepton number is small $`(LB)`$. Furthermore, it has been suggested that even if the baryon number asymmetry is small, the total lepton number could be large in the context of $`SU(5)`$ and $`SO(10)`$ Grand Unified Theories (GUT’s) (Harvey & Kolb (1981); Fry & Hogan (1982); Dolgov & Kirilova (1991); Dolgov (1992)). It has also been proposed recently (Casas, Cheng, & Gelmini (1999)) that models based upon the Affleck-Dine scenario of baryogenesis (Affleck & Dine (1985)) might generate naturally lepton number asymmetry which is seven to ten orders of magnitude larger than the baryon number asymmetry. Neutrinos with large lepton asymmetry and masses $`0.07\text{ eV}`$ might even explain the existence of cosmic rays with energies in excess of the Greisen-Zatsepin-Kuzmin cutoff (Greisen (1966); Gelmini & Kusenko (1999)). Degenerate, massive (2.4 eV) neutrinos may even be required (Larsen & Madsen (1995)) to provide a good fit to the power spectrum of large scale structure in mixed dark matter models. It is, therefore, equally important for both particle physics and cosmology to carefully scrutinize the limits which cosmology places on the allowed range of both the lepton and baryon asymmetries. The consequences of a large lepton asymmetry and associated neutrino degeneracy for big bang nucleosynthesis (BBN) have been considered in many papers. Models with degenerate $`\nu _e`$ (Wagoner, Fowler, & Hoyle (1967); Terasawa & Sato (1985); Scherrer (1983); terasawa88 (1988); Kajino & Orito 1998a ), both $`\nu _e`$ and $`\nu _\mu `$ (Yahil & Beaudet (1976); Beaudet & Goret (1976); Beaudet & Yahil (1977)), and for three degenerate neutrino species (David & Reeves (1980); Olive et al. (1991); Kang & Steigman (1992); Starkman (1992); Kajino & Orito 1998b ; Kim & Lee (1995); Kim, Kim & Lee (1998)) have been analyzed. The effects of the degeneracy of electron type neutrinos on inhomogeneous BBN models were also considered in Kajino & Orito (1998a, 1998b). Constraints on lepton asymmetry also arise from the requirement that sufficient structure develop by the present time (Kang & Steigman (1992)) and from the power spectrum of fluctuations in galaxies (Larsen & Madsen (1995)) and the cosmic microwave background temperature (Kinney & Riotto (1999); Lesgourgues & Pastor (1999); Hannestad (2000)). The present work differs from those listed above primarily in that we carefully reexamine the neutrino decoupling temperature. Although neutrino degeneracy has been treated in considerable detail previously (cf. Kang & Steigman (1992)) we reexamine several issues which were not treated in that paper or elsewhere. We show that there are corrections to the neutrino reaction rates which cause the neutrino decoupling temperature to be significantly higher than previous estimates. These corrections include a proper accounting of effects of degeneracy on the abundances of the weakly-reacting species present. We also consider corrections in the the electron and photon effective masses and the relative velocities of interacting species. In addition we make a more careful continuous treatment of the various particle annihilation epochs. These corrections are important since it becomes easier for neutrinos to decouple before the annihilation of various particles and even before the QCD transition. The fact that neutrinos may decouple when there were many particle degrees of freedom causes the relic neutrino temperature to be much lower than that of the standard nondegenerate big bang. This allows for new solutions for BBN in which significant baryon density is possible without violating the various abundance constraints. We find that for decoupling temperatures just above the QCD epoch it is also possible to find models in which the structure constraint and even the CMB power spectrum constraint are marginally satisfied. In this paper we explore a physically plausible range of baryon and lepton asymmetric cosmological models in which the electron neutrino is only slightly degenerate while the muon and tau neutrinos have nearly equal (but opposite-sign) significant chemical potentials. We use these models to deduce new cosmological constraints on the baryon and lepton content of the universe. We emphasize that previous studies of BBN and the CMB temperature fluctuations with highly degenerate neutrinos have not exhaustively scrutinized the possibility of important effects from the lower implied neutrino temperature when neutrinos decouple before various particle annihilations and/or the QCD epoch. In what follows we first discuss here how the observable BBN yields in a neutrino-degenerate universe impose bounds on the baryon and lepton asymmetries which still allow a large neutrino degeneracy and baryon density (even $`\Omega _\mathrm{b}=1`$). We show that there exists a one parameter family of allowed neutrino-degenerate models which satisfy all of the constraints from BBN. We next examine other cosmological non-nucleosynthesis constraints, i.e. the cosmic microwave background (CMB) fluctuations and the time scale for the development of structure. We show that these constraints can be marginally satisfied for a limited range of highly degenerate models from the fact the relic neutrino temperature is much lower than that of the standard nondegenerate big bang. We also discuss how a determination of the neutrino degeneracy parameters could constrain the neutrino mass spectrum from the implied neutrino contribution $`\mathrm{\Omega }_\nu `$ to the closure density. ## 2 LEPTON ASYMMETRY & NEUTRINO DECOUPLING In this section, we review the basic relations which define the magnitude of neutrino degeneracy and summarize the cosmological implications. Radiation and relativistic particles dominate the evolution of the early hot big bang. In particular, relativistic neutrinos with masses less than the neutrino decoupling temperature, mν <𝒪(TD) MeV <subscript𝑚𝜈𝒪subscript𝑇𝐷similar-to MeVm_{\nu}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}{\cal O}(T_{D})\sim\mbox{~{}MeV}, contributed an energy density greater than that due to photons and charged leptons. Therefore, a small modification of neutrino properties can significantly change the expansion rate of the universe. The energy density of massive degenerate neutrinos and antineutrinos for each species is described by the usual Fermi-Dirac distribution functions $`f_\nu `$ and $`f_{\overline{\nu }}`$, $`\rho _\nu +\rho _{\overline{\nu }}=`$ (1) $`{\displaystyle \frac{1}{2\pi ^2}}{\displaystyle _0^{\mathrm{}}}𝑑pp^2\sqrt{p^2+m_\nu ^2}(f_\nu (p)+f_{\overline{\nu }}(p)).`$ where, $`p`$ denotes the magnitude of the 3-momentum, and $`m_\nu `$ is the neutrino mass. Here and throughout the paper we use natural units( $`\mathrm{}=c=k_B=1`$). The distribution functions are $$\begin{array}{ccc}\hfill f_\nu (p)=\frac{1}{\mathrm{exp}\left(p/T_\nu \xi _\nu \right)+1},& & \\ \hfill f_{\overline{\nu }}(p)=\frac{1}{\mathrm{exp}\left(p/T_\nu +\xi _\nu \right)+1},& & \end{array}$$ (2) where the degeneracy parameter $`\xi _\nu `$ is defined in term of the neutrino chemical potential, $`\mu _\nu `$, as $`\xi _\nu \mu _\nu /T_\nu `$. It will have a nonzero value if a lepton asymmetry exists. In the limit of massless neutrinos, the energy density in neutrinos becomes $$\rho _\nu +\rho _{\overline{\nu }}=\frac{7}{8}\frac{\pi ^2}{15}\underset{i}{}T_{\nu _i}^4\left[1+\frac{15}{7}\left(\frac{\xi _{\nu _i}}{\pi }\right)^4+\frac{30}{7}\left(\frac{\xi _{\nu _i}}{\pi }\right)^2\right],$$ (3) from which it is clear that a neutrino degeneracy in any species tends to increase the energy density. The net lepton asymmetry $`L`$ of the universe can be expressed to high accuracy as $$\begin{array}{ccc}\hfill L=\underset{l=e,\mu ,\tau }{}L_l,& & \\ \hfill L_l=\frac{n_{\nu _l}n_{\overline{\nu }_l}}{n_\gamma }.& & \end{array}$$ (4) This is analogous to the baryon-to-photon ratio $`\eta (n_Bn_{\overline{B}})/n_\gamma `$. Here, $`n_{\nu _l}`$ $`(n_{\overline{\nu }_l})`$ are the number densities for each neutrino (anti-neutrino) species, $`n_\gamma `$ is the photon number density, and $`n_B`$ ($`n_{\overline{B}}`$) is the (anti) baryon number density. Once the temperature drops sufficiently below the muon rest mass, say T <10 MeV <𝑇10 MeVT\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}10\mbox{~{}MeV}, all charged leptons except for electrons and positrons will have decayed away. Overall charge neutrality then demands that the difference between the number densities of electrons and positrons equal the proton number density. Hence, any electron degeneracy is negligibly small and any significant lepton asymmetry must reside in the neutrino sector. After the epoch of $`e^\pm `$ annihilation, the magnitudes of the lepton and baryon asymmetries are conserved. They are equal to the present value in the absence of any subsequent baryon and/or lepton number-violating processes. Elastic scattering reactions, such as $`\nu _l(\overline{\nu }_l)+l^\pm \nu _l(\overline{\nu }_l)+l^\pm `$, keep the neutrinos in kinetic equilibrium. Annihilation and creation processes which can change their number density, like $`\nu _l+\overline{\nu }_ll+\overline{l}`$, $`\nu _l+l^{}\nu _l^{}+l`$, etc, maintain the neutrinos in chemical equilibrium. When the rates for these weak interactions become slower than the Hubble expansion rate, neutrinos decouple and begin a ”free expansion”. Their number densities continue to diminish as $`1/R^3`$ and their momenta red-shift by a factor $`1/R`$, where $`R`$ is the cosmic scale factor. However, this decoupling has no effect on the shape of the distribution functions. Relativistic neutrinos and antineutrinos continue to be described by the Fermi-Dirac distributions of Eq. (2). Since the individual lepton number is believed to be conserved, the degeneracy parameters $`\xi _{\nu _l}`$ remain constant after decoupling. ### 2.1 Neutrino Decoupling Temperature When one estimates the present density of relic neutrinos one must consider the effect of the changing number of degrees of freedom for the remaining interacting particles. For example, once the neutrinos are totally decoupled, they are not heated by subsequent pair annihilations. Hence, their temperature $`T_\nu `$ is lower than the temperature $`T_\gamma `$ of photons (and any other electromagnetically interacting particles) by a factor $`y_\nu =T_\nu /T_\gamma `$. In the standard non-degenerate cosmology, with three flavors of massless, non-degenerate neutrinos which decouple just before the $`e^\pm `$ pair annihilation epoch, the present ratio of the neutrino to photon temperatures is given by $`y_\nu =(4/11)^{1/3}`$. Neutrinos and anti-neutrinos drop out of thermal equilibrium with the background thermal plasma when the weak reaction rate per particle, $`\mathrm{\Gamma }`$, becomes slower than the expansion rate, $`H`$. Since the decoupling occurs quickly, it is a good approximation to estimate a characteristic decoupling temperature $`T_D`$, at which the slowest weak-equilibrating reaction per particle $`\mathrm{\Gamma }`$ becomes slower than the expansion rate. The expansion rate is simply given by the Friedmann equation $$H=\sqrt{(8/3)\pi G\underset{i}{}\rho _i},$$ (5) where the sum is over all species contributing mass energy (mostly relativistic particles). Obviously, the universal expansion in neutrino-degenerate models is more rapid because of the higher neutrino mass-energy density (e.g. Equation 3). This causes the decoupling to occur sooner and at a higher temperature than in the non-degenerate case. Our estimates of the weak decoupling temperature as a function of neutrino degeneracy are shown in Figure 1 and compared with the previous estimates Kang & Steigman (1992). Our decoupling temperatures are significantly higher than those estimated in Kang & Steigman (1992) for reasons which we now explain. Significant neutrino degeneracy will cause the weak reaction rate per particle, $`\mathrm{\Gamma }`$, to be slower than the nondegenerate case because the neutrino final states are occupied (Kang & Steigman (1992); Freese, Kolb, & Turner (1983)). More importantly, the existence of a lepton asymmetry represented by a finite positive neutrino chemical potential will cause the number density of antineutrinos to be significantly less than that of neutrinos. Hence, the annihilation rate per neutrino $`\mathrm{\Gamma }_\nu =n_{\overline{\nu }}\sigma _{\nu \overline{\nu }}`$ can be significantly less than the reaction rate per antineutrino, electron or positron. In Kang & Steigman (1992) the weak reaction rate was estimated from the reaction rate per electron annihilation, $`e^++e^{}\nu _i+\overline{\nu _i}`$. We will call this $`\mathrm{\Gamma }_e`$. One could however also consider the reaction rate per positron, $`\mathrm{\Gamma }_{e^+}`$, the reaction rate per neutrino species $`i`$, $`\mathrm{\Gamma }_{\nu _i}`$, or the reaction rate per antineutrino, $`\mathrm{\Gamma }_{\overline{\nu }_i}`$. In the limit of high temperature and no neutrino degeneracy, all of these reaction rates are identical. However, when for example the neutrinos have a large positive chemical potential, the neutrino annihilation rate $`\mathrm{\Gamma }_{\nu _i}=n_{\overline{\nu }_i}\sigma v`$ becomes the rate determining reaction. This is because of the scarcity of anti-neutrinos with which to interact. It is then this reaction rate and not $`\mathrm{\Gamma }_e`$ which should be used to determine when the weak reactions are no longer able to maintain chemical equilibrium. In Kang & Steigman (1992) expressions were derived for the phase space integrals for $`\sigma v`$ which reduce to $$\sigma v=\frac{8\left(a^2+b^2\right)}{27\pi \zeta ^2(3)}I(\xi )G_FT^2,$$ (6) where $`I(\xi )`$ is a phase space integration factor given in that paper, $`\zeta `$ is the Riemann zeta function, and $`G_F`$ is the Fermi coupling constant. The quantities $`a`$ and $`b`$ relate to the Weinberg angle, $`a=1/2+2\mathrm{sin}^2\theta _W`$ for electron neutrinos, and $`a=1/22\mathrm{sin}^2\theta _W`$ for muon and tau neutrinos, while $`b=1/2`$. Using this, we derive the following expression for the decoupling temperature corresponding to each neutrino species, $`i`$, $$T_D(\xi _i)\left[\frac{4.38}{a^2+b^2}\right]^{1/3}\left[\frac{\sqrt{g_{eff}^{}(\xi )}}{I(\xi )}\right]^{1/3}\times \left(\frac{n_{\overline{i}}}{n_e}\right)^{1/3}\mathrm{MeV},$$ (7) where $`\overline{i}`$ denotes the matter (or antimatter) counter part for the species under consideration. The formula A.8 of Kang & Steigman (1992) omits the latter density ratio factor which is only unity in the limit of no degeneracy. Thus, their $`\mathrm{\Gamma }`$ used to deduce the decoupling temperature was too large. We have checked the above result against a full eight-dimensional integration of the Boltzmann equation and find our results to be correct. There are also two other differences between our results and those of Kang & Steigman (1992): The first is that we have included finite temperature corrections to the mass of the electron and photon (Fornengo, Kim, & Song (1997)). The second is that we have calculated the average neutrino annihilation rate in the cosmic comoving frame. In this frame the Møller velocity is used instead of the relative velocity for the integration of the collision term in the Boltzmann equations (Gondolo, & Gelmini (1991); Enqvist et al. (1992)). These corrections also slightly increase the neutrino decoupling temperatures. Therefore even in the nondegenerate case we find slightly higher decoupling temperatures than those of Kang & Steigman (1992): $`T_D(\xi _\nu =0)2.33\text{ MeV}\mathrm{for}\nu _\mathrm{e},`$ $`T_D(\xi _\nu =0)3.90\text{ MeV}\mathrm{for}\nu _{\mu ,\tau }.`$ (8) These values are in good agreement with those of Enqvist et al. (1992). If we remove the above corrections, we recover the decoupling temperatures of (Kang & Steigman 1992), i.e. $`T_D(\xi _\nu =0)1.98`$ MeV for $`\nu _e`$, and $`T_D(\xi _\nu =0)3.30`$ MeV for $`\nu _{\mu ,\tau }`$. ### 2.2 Relic Neutrino Temperature One does not need to increase the decoupling temperature by much before photon heating by annihilations becomes important for determining the relic neutrino temperature. For illustration, Figure 2 shows the ratio of muon to photon energy densities as a function of temperature in units of the muon rest mass $`m_\mu =105`$ MeV. A similar curve could be drawn for any massive species. One can see that even at a temperature of only 20% of the muon rest mass, muons still contribute about 10% of the mass energy density, and hence, can affect the ratio of the photon to neutrino temperatures as these remaining muons annihilate. Combining Figures 1 and 2, one can see that even for a degeneracy parameter of $`\xi _\nu 6`$, the decoupling temperature is at 20% of $`m_\mu `$. For the case of highly degenerate neutrinos (ξνμ >10.8 >subscript𝜉subscript𝜈𝜇10.8\hbox{$\xi_{\nu_{\mu}}$}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}}10.8 and ξντ >9.8 >subscript𝜉subscript𝜈𝜏9.8\hbox{$\xi_{\nu_{\tau}}$}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}}9.8), $`T_D(\xi _{\nu _i})`$ can exceed the muon rest energy and even the QCD phase transition temperature. If the neutrinos decouple early, they are not heated as the number of particle degrees of freedom change. Hence, the ratio of the neutrino to photon temperatures, $`T_\nu /T_\gamma `$, is reduced. The computation of the ratio of the final present neutrino temperature to the photon temperature is straightforward. Basically, since the universe is a closed system, the relativistic entropy is conserved within a comoving volume. That is; $$R^3s=\mathrm{Constant},$$ (9) where the entropy density $`s`$ is defined, $$s\frac{(\rho _i+p_i\mu _in_i)}{T}=\frac{2\pi ^2}{45}g_{eff}T^3,$$ (10) and the sum is over all species present. Since the mass-energy is dominated by relativistic particles $`s`$ can be written in terms of the effective number of particle degrees of freedom, $$g_{eff}=\underset{i=bosons}{}g_i\left(\frac{T_i}{T}\right)^3+\frac{7}{8}\underset{i=fermions}{}g_i\left(\frac{T_i}{T}\right)^3.$$ (11) As each species annihilates, $`sR^3`$ remains constant. Therefore, the temperature of the remaining species increases by a factor of $`(g_{eff}^{before}/g_{eff}^{after})^{1/3}`$. This accounts for the usual heating of photons relative to neutrinos due to $`e^\pm `$ pair annihilations by a factor of $`(11/4)^{1/3}`$. Note, that in the computation of $`g_{eff}`$ from equations (10) and (11) it is important to evaluate the energy densities continuously (cf. Figure 2) and not simply assume abrupt annihilation as the temperature approaches the rest energy of each particle as was done in Kang & Steigman (1992). Figure 3 illustrates $`g_{eff}`$ as a function of temperature from 1 MeV to 1 TeV. In each calculation $`g_{eff}`$ depends upon the relic neutrino temperature and therefore when the neutrinos decoupled. For purposes of illustration only, we drawn this figure for three neutrino types which do not decouple until $``$ MeV. To construct this figure we have included: mesons ($`\pi `$, $`\rho `$, $`\varphi `$, $`\omega `$, $`\eta ,\eta ^{}`$); leptons ($`e`$, $`\mu `$, $`\tau `$, $`\nu _e`$, $`\nu _\mu `$, $`\nu _\tau `$); a QCD phase transition at 150 MeV; light quarks and gluons ($`u,d,gluon`$); s-quarks (with $`m_s=120`$ MeV); c-quarks (with $`m_c=1200`$ MeV); b-quarks (with $`m_b=`$ 4250 MeV); W, Z-bosons (with $`m_Z`$ = 80 GeV); and t-quarks (with $`m_t=`$ 173 GeV). One can see from this figure that if neutrinos decouple before the QCD phase transition, there is substantial heating of photons due to the large change in degrees of freedom at this epoch. Figure 4 shows the final ratio of neutrino temperature today to that of the standard non-degenerate big bang for three neutrino flavors. The biggest drop in temperature for all three neutrino flavors occurs for $`\xi _\nu 10`$. This corresponds to a decoupling temperature above the cosmic QCD phase transition. The low temperature is the result of the decrease in particle degrees of freedom during this phase transition. This discontinuity will have important consequences in the subsequent discussions. For comparison, the same computation from Figure 1 of Kang & Steigman (1992) is shown. One can see that there is a substantial difference between the two results. These differences stem from two effects. On the one hand the decoupling temperature is higher in the present work for a given $`\xi _{\nu _i}`$ because of the various corrections as explained above. This is the reason that the discontinuity in relic neutrino temperature occurs at $`\xi _{\nu _{\mu ,\tau }}11`$ instead of 15 as in Kang & Steigman (1992). The second difference is that here we have calculated the proper thermodynamic energy density for each fermion and boson species continuously (cf. Eq. 1), whereas in Kang & Steigman (1992) an approximation was made that each particle species instantaneously annihilates as the temperature drops below its rest energy. This is the reason for the step function appearance of the curve in Kang & Steigman (1992). Some discussion of book keeping here is in order as it explains why we observe the jump in relic temperature when neutrinos decouple above the QCD transition. Let us consider a neutrino which decouples just above the QCD transition. Then the relevant fermion degrees of freedom consist mainly of 3 $`q\overline{q}`$ colors (treating the $`s`$-quark as relativistic) $`\times `$ 3 flavors contributing $`g_f=31.5`$, plus $`e^+e^{}`$ and $`\mu ^+\mu ^{}`$ plus 2 remaining nondecoupled flavors of $`\nu \overline{\nu }`$ lepton pairs. These contributing another $`g_f=10.5`$. For the bosons there are 8 gluons and 1 photon contributing a total $`g_b=18`$. While the species are in equilibrium the total $`g_{eff}=60`$. If no more species decouple before the electron pair annihilation epoch, the fermion degrees of freedom before electron-positron annihilation (at $`T1`$ MeV) include only 2 flavors of $`\nu \overline{\nu }`$ plus $`e^+e^{}`$ so that $`g_f=7`$, while for bosons only photons remain so that $`g_b=2`$ and $`g_{eff}=9`$. The relic neutrino temperature is therefore reduced by an additional factor due to these differences in particle degrees of freedom $`(T_\nu /T_\gamma )/[4/11]^{1/3}=(9/60)^{1/3}0.53`$ (cf. Figure 4). For all three neutrino flavors the temperature begins to decrease relative to the standard value for a degeneracy parameter as small as $`\xi 5`$. This is because some relic $`\mu \overline{\mu }`$ pairs are still present even at temperatures well below the muon rest energy (cf. Figure 2). The first neutrino species to be affected as the degeneracy increases are the $`\nu _\mu `$ and $`\nu _\tau `$. They decouple at a higher temperature than $`\nu _e`$ even in the standard nondegenerate big bang because the electrons continue to experience charged-current interactions to lower temperature. The muon neutrinos exhibit a slightly different behavior than $`\nu _\tau `$ for degeneracy parameters $`\xi _{\nu _\mu }>5`$ because the $`\mu \overline{\mu }`$ density is large enough at the decoupling temperature for charged-current interactions to be significant. This maintains equilibrium for $`\nu _\mu `$ to somewhat lower temperatures even for degenerate neutrinos. This causes the $`\nu _\mu `$ decoupling temperature to be lower (cf. Figure 1) and the relic temperature to be slightly higher than the $`\nu _\tau `$ temperature for degeneracy parameters between 5 and 9. ## 3 PRIMORDIAL NUCLEOSYNTHESIS ### 3.1 Current Status Although the homogeneous BBN model has provided strong support for the standard hot big-bang cosmology, possible conflicts have emerged between the predictions of BBN abundances as the astronomical data have become more precise in recent years. One difficulty has been imposed by recent detections of a low deuterium abundance (Burles & Tytler 1998a ; Burles & Tytler 1998b , see also Levshakov, Tytler, & Burles (1999)) in Lyman-$`\alpha `$ absorption systems along the line of sight to high red-shift quasars. The low D/H favors a high baryon content universe and a high primordial <sup>4</sup>He abundance, Yp >0.244 >subscript𝑌𝑝0.244Y_{p}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}}0.244. This is inconsistent with at least some of the reported constraints from measurements of a low primordial abundance of <sup>4</sup>He, $`Y_p0.235\pm 0.003`$, in low-metallicity extragalactic H II regions (Olive & Steigman (1995); Steigman (1996); Hata et al. (1996); Olive, Steigman, & Skillman (1997); Kajino & Orito 1998a ; Piembert & Piembert (2000)). This situation is exacerbated by recent detailed analyses (Esposito et al. (1999); Lopez & Turner 1999) of the theoretical uncertainties in the weak interactions affecting the neutron to proton ratio at the onset of primordial nucleosynthesis. These results require a positive net correction to the theoretically determined <sup>4</sup>He mass fraction $`Y_p`$ of $`+0.004`$ to $`+0.005`$ or $``$2%. We also note that the low deuterium abundance is marginally inconsistent with the <sup>7</sup>Li abundance inferred by measurements of lithium in Population II halo stars (Ryan et al. (1999); Kajino et al. (2000)). Significant depletion of lithium from these stars, or a lower reaction rate for primordial lithium production may be required. Another potential difficulty has been imposed by recent X-ray observations of rich clusters (White et al. (1993); White & Fabian (1995); David, Jones, & Forman (1995); Bahcall, Lubin, & Dorman (1995)). The implied baryonic contribution to the closure density is $`0.08\Omega _\mathrm{b}h_{50}^{3/2}/\mathrm{\Omega }_M0.22`$ (Tytler et al. (2000)), where $`\mathrm{\Omega }_M`$ is the total matter (dark plus visible) contribution, and $`h_{50}`$ is the Hubble constant $`H_0`$ in units of 50 km s<sup>-1</sup> Mpc<sup>-1</sup>. Consistency with the limits ($`0.03\Omega _\mathrm{b}h_{50}^20.06`$) from homogeneous BBN (Walker et al. (1991); Smith et al. (1993); Copi, Schramm & Turner (1995); Schramm & Mathews (1995); Olive, Steigman, & Walker (1999)) then requires that $`0.14\mathrm{\Omega }_Mh_{50}^{1/2}0.75`$. Hence, matter dominated cosmological models (for example with $`H_0=75`$ km s<sup>-1</sup> Mpc<sup>-1</sup> and $`\mathrm{\Omega }_M0.61`$) can be in conflict with BBN. Another possible conflict has emerged in recent measurements of the power spectrum of the cosmic microwave background (CMB) temperature fluctuations. A suppression of the second acoustic peak in the power spectrum has been recently reported in the balloon based CMB sky maps from the BOOMERANG (Bernardls, P. et al. (2000)) and MAXIMA-1 (Hanany et al. (2000)) collaborations. Indeed, the derivation of cosmological parameters from these new data sets (Lange et al. (2000); Balbi et al. (2000)) indicate a preference for a large baryon density. For example, the optimum cosmological models consistent with the BOOMERANG data (Lange et al. (2000)) indicate $`\mathrm{\Omega }_bh^2=0.032\pm 0.004`$ $`(1\sigma )`$ for fits in which there is no prior restriction on the range of $`\mathrm{\Omega }_bh^2`$. Similarly, the likelihood analysis based upon the MAXIMA-1 data (Balbi et al. (2000)) indicates $`\mathrm{\Omega }_bh^2=0.025\pm 0.005`$. These results are to be compared with the best current $`(1\sigma )`$ limit from standard primordial nucleosynthesis without neutrino degeneracy, $`\mathrm{\Omega }_bh^2=0.019\pm 0.0024`$ (Olive, Steigman, & Walker (1999);Tytler et al. (2000)) which come from the deuterium abundance in Lyman-$`\alpha `$ clouds observed along the line of sight to background quasars. The CMB data indicate a baryon content which is at least 1-2 $`\sigma `$ above the optimum BBN value. Indeed, $`\mathrm{\Omega }_bh^20.20`$ would require a primordial helium abundance of $`Y_p0.25`$ which is beyond even the most generous adopted limits from observation (Olive, Steigman, & Walker (1999)). Thus, it seems that that the CMB data hint at (though not necessarily prove) the need for a modification of BBN which allows higher values of $`\mathrm{\Omega }_bh^2`$ while still satisfying the constraints from light element abundances. Such conditions are easily satisfied by neutrino-degenerate models. ### 3.2 Neutrino-Degenerate BBN Most previous works have only considered the effects of neutrino degeneracy on the light elements <sup>4</sup>He, D, and <sup>7</sup>Li. Recently, the predicted abundance of other elements such as <sup>6</sup>Li, <sup>9</sup>Be, and <sup>11</sup>B in a neutrino-degenerate universe were also studied (Kim & Lee (1995); Kim, Kim & Lee (1998)). Here we investigate the effects of lepton asymmetry on the predicted abundances of heavier elements ($`12A18`$) as well as these light elements. Non-zero lepton numbers primarily affect nucleosynthesis in two ways (Wagoner, Fowler, & Hoyle (1967); Terasawa & Sato (1985); terasawa88 (1988); Scherrer (1983); Yahil & Beaudet (1976); Beaudet & Goret (1976); Beaudet & Yahil (1977); David & Reeves (1980); Olive et al. (1991); Kang & Steigman (1992); Starkman (1992); Kajino & Orito 1998a ; Kim & Lee (1995); Kim, Kim & Lee (1998)). First, degeneracy in any neutrino species leads to an increased universal expansion rate independently of the sign of $`\xi _{\nu _e}`$ (cf. Eqs. 3 and 5). As a result, the weak interactions that maintain neutrons and protons in statistical equilibrium decouple earlier. This effect alone would lead to an enhanced neutron-to-proton ratio at the onset of the nucleosynthesis epoch and increased <sup>4</sup>He production. Secondly, a non-zero electron neutrino degeneracy can directly affect the equilibrium n/p ratio at weak-reaction freeze out. The equilibrium n/p ratio is related to the electron neutrino degeneracy by $`\mathrm{n}/\mathrm{p}=\mathrm{exp}\{(\mathrm{\Delta }M/T_{np})\xi _{\nu _e}\}`$, where $`\mathrm{\Delta }M`$ is the neutron-proton mass difference and $`T_{np}`$ is the freeze-out temperature for the weak reactions converting protons to neutrons and vice versa. This effect leads to either increased or decreased <sup>4</sup>He production, depending upon the sign of $`L_e`$ or $`\xi _{\nu _e}`$. There is also a third effect which we emphasize in this paper. As discussed in the previous section, $`T_\nu /T_\gamma `$ can be reduced if the neutrinos decouple at an earlier epoch. This lower temperature reduces the energy density of the highly degenerate neutrinos during the BBN era, and hence, slows down the expansion of the universe. This leads to decreased <sup>4</sup>He production. We show in the next section that the allowed values for $`\xi _{\nu _e},\xi _{\nu _\mu },\xi _{\nu _\tau }`$ and $`\Omega _\mathrm{b}`$ which satisfy the light-element abundance constraints are significantly changed for large degeneracy (ξνμ,ξντ >10 >subscript𝜉subscript𝜈𝜇subscript𝜉subscript𝜈𝜏10\hbox{$\xi_{\nu_{\mu}}$},~{}\hbox{$\xi_{\nu_{\tau}}$}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}}10) compared to the results of previous studies. ### 3.3 Summary of Light-Element Constraints The primordial light element abundances deduced from observations have been reviewed in a number of recent papers (Olive, Steigman, & Walker (1999): Nolett & Burles (2000); Steigman (2000); Tytler et al. (2000)). There are several outstanding uncertainties. For primordial helium there is an uncertainty due to the fact that deduced abundances tend to fall in one of two possible values, one high and the other low. Hence, for <sup>4</sup>He we adopt a wide range: $$0.226Y_p0.247.$$ For deuterium there is a similar possibility for either a high or low value. Here, however, we adopt the generally accepted low values of Tytler et al. (2000), $$2.9\times 10^5\text{D/H}4.0\times 10^5.$$ For primordial lithium there is some uncertainty from the possibility that old halo stars may have gradually depleted their primordial lithium. Because of this possibility we adopt the somewhat conservative constraint: $$1.26\times 10^{10}\text{7}\text{Li}/\mathrm{H}3.5\times 10^{10}$$ ### 3.4 Nucleosynthesis Model As we shall see, shifts in the relic neutrino temperature can dramatically affect the abundance yields (Kajino & Orito 1998b ). We now explore the parameter space of neutrino degeneracy and baryonic content to reinvestigate the range of models compatible with the constraints from light element abundances. For the present work we have applied a standard big bang code with all reactions updated up to A=18. \[However, in the present discussion only reactions involving nuclei up to A=15 are significant.\] In this way possible effects of lepton asymmetry on heavier element abundances could be analyzed along with the light elements. In this context a recent compilation of the nuclear reaction rates relevant to the production of <sup>11</sup>B (Orito, Kajino, & Oberhummer (1998)) was useful because several important LiBeB(a,x) and (n,$`\gamma `$) reaction rates in the literature sometimes differ from one another by 2-3 orders of magnitude. The calculated abundances of heavier elements based upon these rates can also differ from one another by 1-2 orders of magnitude. We carry out BBN calculations which include all of the recent compilations of reaction rates relevant to the production of isotopes including those that are heavier than <sup>7</sup>Li up to <sup>18</sup>O (i.e. Orito, Kajino, & Oberhummer (1998); Mohr, Herndl, & Oberhummer (1999); Angulo et al. (1999); Herndl et al. (1999); Wagemans et al. (1999); and any other previous estimates are considered). ### 3.5 Neutrino Degeneracy Parameters We have explored a broad range of the parameter space of neutrino-degenerate models. The main effects of the inclusion of either $`\nu _\mu `$ or $`\nu _\tau `$ degeneracy on BBN is an enhancement of energy density of the universe. The values for $`\xi _{\nu _\mu }`$ and $`\xi _{\nu _\tau }`$ primarily affect the expansion rate. This means that $`\xi _{\nu _\mu }`$ and $`\xi _{\nu _\tau }`$ are roughly interchangeable as far as their effects on nucleosynthesis are concerned. Furthermore, we expect that the net total lepton number is small though the lepton number for individual species could be large. Hence, it is perhaps most plausible to assume that the absolute values of $`\|\xi _{\nu _\mu }\|`$ and $`\|\xi _{\nu _\tau }\|`$ are nearly equal. Therefore, in what follows, we describe results for $`\|\xi _{\nu _\mu }\|=\|\xi _{\nu _\tau }\|\xi _{\nu _{\mu ,\tau }}`$. This reduces the parameter space to three quantities: $`\Omega _\mathrm{b}`$, $`\xi _{\nu _e}`$, and $`\|\xi _{\nu _\mu }\|=\|\xi _{\nu _\tau }\|`$. The calculations were conducted by first choosing a value for $`\mathrm{\Omega }_bh_{50}^2`$. It was then necessary to find a value of $`\xi _{\nu _e}`$ for which the light element constraints could be satisfied for some value of $`\xi _{\nu _{\mu ,\tau }}`$. ### 3.6 Results As an illustration, Figure 5 shows a calculations of the helium abundance as a function of $`\xi _{\nu _{\mu ,\tau }}`$ in a model with $`\mathrm{\Omega }_bh_{50}^2=0.3`$. It was found in this model that the helium constraint could be satisfied for $`\xi _{\nu _e}`$ in the range of 0.79 to 0.94. This figure is for $`\xi _{\nu _e}=0.9`$. Our results for $`\xi _{\nu _{\mu ,\tau }}<10`$ are consistent with those of Kang & Steigman (1992) if we run for a similar baryon-to-photon ratio $`\eta _{10}=5`$ or $`100`$ and utilize the neutron half life adopted in that paper. In our calculations for helium (Figure 5) and other light nuclides (Figure 6) we see two different jumps in the nucleosynthesis yields. There is a small one for $`\xi _{\nu _{\mu ,\tau }}10`$ and a large one for $`\xi _{\nu _{\mu ,\tau }}11.4`$. These two jumps are due to the fact that we consider two degenerate neutrino species which decouple at slightly different temperatures (due to the presence of muons to interact with $`\nu _\mu `$). Therefore, one species ($`\nu _\tau `$) decouples above the QCD transition for a smaller value for $`\xi _{\nu _{\mu ,\tau }}`$. However, this jump has a smaller effect on nucleosynthesis because so many degrees of freedom are still carried by the other degenerate neutrino species ($`\nu _\mu `$) which decouples later. The main difference between our results and Kang & Steigman (1992), however, is that the helium abundance is significantly changed at high values of $`\xi _{\nu _{\mu ,\tau }}`$. This results from our higher decoupling temperatures and therefore the lower relic neutrino temperatures during primordial nucleosynthesis (See Figure 4). This effect of varying $`T_\nu /T_\gamma `$ for large $`\xi _\nu `$ on BBN has not been adequately explored in the previous studies. Figure 6 illustrates the effects on the other light-element abundances for this particular parameter set. This figure shows that for moderate values of $`\xi _{\nu _{\mu ,\tau }}`$ the main effect is that weak reaction freeze-out occurs at a higher temperature. The resultant enhanced n/p ratio increases the abundances of the neutron-rich light elements, D, $`{}_{}{}^{3}\mathrm{H}`$ and <sup>7</sup>Li, while the <sup>7</sup>Be abundance decreases. Regarding <sup>7</sup>Li and <sup>7</sup>Be, the enhanced expansion rate from neutrino degeneracy affects the yield of $`A=7`$ elements in two different ways. These elements are produced mainly by the nuclear electromagnetic capture reactions: $`t(\alpha ,\gamma )\text{7}\text{Li}`$ and $`\text{3}\text{He}(\alpha ,\gamma )\text{7}\text{Be}`$. Hence, the production of these elements begins at a later time and lower temperature than the other light elements because they require time for a significant build up of the reacting A = 3 and 4 elements. In a neutrino-degenerate universe, however, the increased expansion rate, hastens the freeze-out of these reactions from nuclear statistical equilibrium (NSE) resulting in reduced $`A=7`$ yields relative to the nondegenerate case. However, the enhanced n/p ratio also increases the tritium abundance in NSE. This effect tends to offset the effect of rapid expansion on the production of <sup>7</sup>Li. The net result is more <sup>7</sup>Li production. As in the case of primordial helium, there is a rapid change in the nucleosynthesis yields on Figure 6 once the $`\nu _\mu `$ and $`\nu _\tau `$ decoupling temperatures separately exceed the epoch of the QCD phase transition. The ensuing lower neutrino temperature during primordial nucleosynthesis then resets the abundances to those of a lower effective degeneracy. The two discontinuities in Figure 6 at $`\xi _{\nu _{\mu ,\tau }}=10`$ and 11.4 correspond to the points at which respectively the tau or muon neutrino decoupling temperatures separately exceed the QCD phase transition temperature (cf. Figure 4). Figure 7(a) summarizes the allowed regions of the $`\xi _{\nu _e}`$ vs. $`\xi _{\nu _{\mu ,\tau }}`$ plane based upon the various indicated light element constraints in a universe with $`\Omega _\mathrm{b}h_{50}^2=0.1`$. The usually identified allowed region (cf. Kohri, Kawasaki, & Sato (1997)) for small $`\xi _{\nu _e}0.2`$ and $`\xi _{\nu _{\mu ,\tau }}2`$ is apparent. Indeed, it has been argued (cf. Kohri, Kawasaki, & Sato (1997)) that such degeneracy may be essential to explain the differences in the constraints from primordial helium and deuterium. Figures 7(b) and 7(c) show the same plots, but for $`\Omega _\mathrm{b}h_{50}^2=0.2`$ and $`\Omega _\mathrm{b}h_{50}^2=0.3`$, respectively. The decline in the primordial deuterium abundance for models in which $`\xi _{\nu _{\mu ,\tau }}>10`$ allows for new regions of the parameter space in which the light element constraints can be accommodated. This new allowed region for large degeneracy persists as the baryon density is increased. Figure 8 highlights the basic result of this study. It shows that there exists a single parameter (either $`\xi _{\nu _{\mu ,\tau }}`$ or $`\Omega _\mathrm{b}`$$`h_{50}^2`$) family of neutrino degenerate models allowed by BBN. For low $`\Omega _\mathrm{b}`$$`h_{50}^2`$ models, only the usual low values for $`\xi _{\nu _e}`$ and $`\xi _{\nu _{\mu ,\tau }}`$ are allowed. Between $`\Omega _\mathrm{b}h_{50}^2`$ 0.187 and 0.3, however, more than one allowed region emerges. For Ωbh502 >0.4 >Ωbh5020.4\hbox{${\sl\Omega}_{\rm b}$}\hbox{$h^{2}_{50}$}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}}0.4 only the large degeneracy solution is allowed. Neutrino degeneracy can even allow baryonic densities up to $`\Omega _\mathrm{b}h_{50}^2=1`$. This result has been noted previously (cf. Kang & Steigman (1992): Starkman (1992)). What is different here is that the high $`\Omega _\mathrm{b}`$$`h_{50}^2`$ models are made possible for smaller values of $`\xi _{\nu _e}`$ due to the higher neutrino decoupling temperatures deduced in the present work. As we shall see, this allows the circumvention of other cosmological constraints as well suggesting that baryons and degenerate neutrinos might provide a larger contribution to the universal closure density than has previously been derived. Figures 9-11 illustrate the elemental abundances obtained in the family of allowed models as a function of $`\Omega _\mathrm{b}`$$`h_{50}^2`$. Figure 11 in particular allows us to consider whether there exists an abundance signature in other elements which might distinguish this new degenerate neutrino solution from standard BBN. For the most part the yields of the light and heavy species are similar to those of the standard non-degenerate big bang. However, the boron abundance exhibits an increase with baryon density due to alpha captures on <sup>7</sup>Li. Thus, in principle, an anomalously high boron abundance together with beryllium and <sup>6</sup>Li similar to that expected from standard BBN might be a signature of neutrino-degeneracy. ## 4 OTHER COSMOLOGICAL CONSTRAINTS We have seen that a new parameter space in the constraints from light elements on BBN emerges in neutrino-degenerate models just from the fact that the relic neutrino temperature is substantially diminished when degeneracy pushes neutrino decoupling to an earlier epoch. The viability of this solution requires large neutrino degeneracy. Hence, it becomes necessary to reexamine constraints posed from other cosmological considerations. ### 4.1 Structure Formation It has been argued (Kang & Steigman (1992)) that large neutrino degeneracy is ruled out from the implied delay in galaxy formation in such a hot dark matter universe. This argument is summarized (Kang & Steigman (1992)) as follows: Neutrino degeneracy speeds up the expansion rate by a factor $$S_0^2(\xi _\nu )=1+0.135(F(\xi _\nu )3)0.135F(\xi _\nu ),$$ (12) where $`F(\xi _\nu )`$ is an effective energy density factor (Kang & Steigman (1992)) for neutrinos. For massless neutrinos we have from Equations 3 and 5, $$F(\xi _\nu )=\underset{i}{}F(\xi _{\nu _i})=\underset{i}{}\left(\frac{T_{\nu _i}}{T_\gamma }\right)_{Nuc}^4\left[1+\frac{15}{7}\left(\frac{\xi _{\nu _i}}{\pi }\right)^4+\frac{30}{7}\left(\frac{\xi _{\nu _i}}{\pi }\right)^2\right],$$ (13) where $`(T_{\nu _i}/T_\gamma )_{Nuc}`$ is the neutrino-to-photon temperature ratio before $`e^\pm `$ pair annihilation. This ratio is unity in the standard model with little or no degeneracy. With the increased mass-energy in degenerate neutrinos, the time of matter-radiation equality occurs at smaller redshift and there is less time for the subsequent growth of fluctuations. The redshift $`z_{eq}(\xi _\nu )`$ for matter-radiation equality for a neutrino-degenerate universe can be written in terms of the redshift for a nondegenerate universe and the speed-up factor, $$1+z_{eq}(\xi _\nu )=S_0^2(1+z_{eq}(\xi _\nu =0)),$$ (14) For the present matter-dominated universe without neutrino degeneracy and $`\mathrm{\Omega }_M1`$ we have 1+zeq(ξν=0) <1.06×104 <1subscript𝑧𝑒𝑞subscript𝜉𝜈01.06superscript1041+z_{eq}(\xi_{\nu}=0)\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}1.06\times 10^{4} (Kang & Steigman (1992)). Furthermore, one demands that the fluctuation amplitude $`A(z)`$ grows at least linearly with redshift, i.e. $`A(z)>(1+z_{eq})A(z_{eq})`$. One also requires that the amplitude at least reaches unity by the present time, 1A(zeq) <(1+zeq) <A(z=0)A(zeq). <1𝐴subscript𝑧𝑒𝑞1subscript𝑧𝑒𝑞 <𝐴𝑧0𝐴subscript𝑧𝑒𝑞{1\over A(z_{eq})}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}(1+z_{eq})\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}{A(z=0)\over A(z_{eq})}~{}. (15) Thus, the requirement of sufficient growth in initial perturbations places a bound on the speed-up. Namely, S02 <1.06×1041+zeq(ξν) <1.06×104A(zeq). <superscriptsubscript𝑆021.06superscript1041subscript𝑧𝑒𝑞subscript𝜉𝜈 <1.06superscript104𝐴subscript𝑧𝑒𝑞S_{0}^{2}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}{1.06\times 10^{4}\over 1+z_{eq}(\xi_{\nu})}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}1.06\times 10^{4}A(z_{eq})~{}. (16) Then requiring that A(zeq) <103 <𝐴subscript𝑧𝑒𝑞superscript103A(z_{eq})\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}10^{-3} (Steigman & Turner (1985)) leads to the constraint that F(ξν) <10.6/.13579 <𝐹subscript𝜉𝜈10.6.13579F(\xi_{\nu})\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}10.6/.135\approx 79. Figure 13 shows the $`F(\xi _\nu )`$ and $`F(\xi _{\nu _i})`$ calculated as a function of $`\xi _{\nu _{\mu ,\tau }}`$ for the allowed models of Figure 8. Since our interesting parameter regions in Figure 8 satisfy $`\xi _{\nu _e}\xi _{\nu _{\mu ,\tau }}`$, only $`\nu _{\mu ,\tau }`$ contribute significantly to the total $`F(\xi _\nu )`$. Also shown are values of $`F(\xi _{\nu _i})`$ if only one neutrino species was degenerate. The dashed line gives the constraint F(ξν) <79 <𝐹subscript𝜉𝜈79F(\xi_{\nu})\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}79 from Kang & Steigman (1992). For comparison, the two-dot dashed line also shows the $`F(\xi _\nu )`$ for a single species from Kang & Steigman (1992). For low values of ξνμ,τ <4 <subscript𝜉subscript𝜈𝜇𝜏4\hbox{$\xi_{\nu_{\mu,\tau}}$}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}4, our $`F(\xi _{\nu _i})`$ values for a single degenerate species are nearly identical to those of Kang & Steigman (1992). However, as $`\xi _{\nu _i}`$ increases, our curves are lower due to the fact that we treat the annihilation epochs continuously (cf. Figure 4) rather than discretely, except for the QCD transition. By chance, our limit in the low degeneracy range for two degenerate neutrinos, (ξνμ,τ <6.2 <subscript𝜉subscript𝜈𝜇𝜏6.2\hbox{$\xi_{\nu_{\mu,\tau}}$}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}6.2) is close to the single species limit (ξiKS <6.9 <superscriptsubscript𝜉𝑖𝐾𝑆6.9\xi_{i}^{KS}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}6.9) of Kang & Steigman (1992). In the present work, however, the single species limit for moderate degeneracy increases to ξi <8.2 <subscript𝜉𝑖8.2\xi_{i}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}8.2. Moreover, in the present work we now find that there also exists a new allowed region with 11.4 <ξνμ,τ <13 <11.4subscript𝜉subscript𝜈𝜇𝜏 <1311.4\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}\hbox{$\xi_{\nu_{\mu,\tau}}$}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}13 for which this growth constraint is satisfied even for two degenerate species. For only one degenerate species the new allowed range expands to 11.4 <ξνμ <16 <11.4subscript𝜉subscript𝜈𝜇 <1611.4\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}\hbox{$\xi_{\nu_{\mu}}$}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}16. An upper limit of ξνμ,τ <13 <subscript𝜉subscript𝜈𝜇𝜏13\hbox{$\xi_{\nu_{\mu,\tau}}$}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}13 for two degenerate neutrino species corresponds to allowed BBN models with a baryon fraction as large as Ωbh502 <0.25 <subscriptΩ𝑏superscriptsubscript5020.25\Omega_{b}h_{50}^{2}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}0.25 (cf. Figure 8). ### 4.2 Neutrino Mass Constraint Figure 12 summarizes the neutrino contribution to the closure density $`\mathrm{\Omega }_\nu `$ as a function of neutrino degeneracy and mass. This figure assumes the plausible model of nearly degenerate $`\nu _\mu `$ and $`\nu _\tau `$ masses and negligible $`\nu _e`$ mass. For this figure $`\mathrm{\Omega }_\nu `$ refers to the combined contributions from both $`\nu _\mu `$ and $`\nu _\tau `$. The contribution changes for large degeneracy due to the lower present-day neutrino temperature. This figure can be used to constrain the masses of the $`\nu `$ and $`\tau `$ neutrino types in different cosmological models. For example, if we assume a model with $`\mathrm{\Omega }_bh_{50}^2=0.1`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.6`$, and $`\mathrm{\Omega }_\nu =0.3`$, then we would find that the masses of both the $`\nu _\mu `$ and $`\nu _\tau `$ must be < <\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}2 eV, if these two species are to provide the neutrino contribution to the closure density. Studies of large scale structure also constrain neutrino masses and degeneracy in hot-dark-matter (HDM) and mixed-dark-matter (MDM) models. Indeed, at least some neutrino mass may presently be required to account for the observed power spectrum of galactic and microwave background structure. It has been argued (Primack et al. (1995)) from considerations of structure formation in the early universe that two neutrino flavors ($`\nu _\mu ,\nu _\tau `$) may have a rest mass of $`2.4`$ eV, compatible with all neutrino oscillation experiments. This postulate solves the main problem of cold dark matter (CDM) models, i.e. production of too much structure on small scales. Furthermore, Larsen & Madsen (1995) argue that neutrino degeneracy is required in a MDM model with 2.4 eV neutrinos to obtain an optimum fit to the power spectrum.. If we take 2.4 eV as the given mass of $`\nu _\mu `$ and $`\nu _\tau `$, then for a Ων <0.9 <subscriptΩ𝜈0.9\Omega_{\nu}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}0.9 $`\mathrm{\Omega }_bh_{50}^2=0.1`$ MDM cosmology we would deduce from Figure 12 that the maximum degeneracy for two species with this mass would correspond to ξνμ,τ <2.5 <subscript𝜉subscript𝜈𝜇𝜏2.5\hbox{$\xi_{\nu_{\mu,\tau}}$}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}2.5 similar to the values used in Larsen & Madsen (1995). We suggest that in light of the present results a similar study of the galactic power spectrum for $`\xi _{\nu _{\mu ,\tau }}11`$ and $`m_\nu 0.1`$ eV should be undertaken as well. ### 4.3 Cosmic Microwave Background Constraint Perhaps, the most stringent remaining constraint on neutrino degeneracy comes from its effect on the cosmic microwave background. Several recent works (Kinney & Riotto (1999); Lesgourgues & Pastor (1999); Hannestad (2000)) have shown that neutrino degeneracy can dramatically alter the power spectrum of the CMB. The essence of this constraint is that degenerate neutrinos increase the energy density in radiation at the time of photon decoupling in addition to delaying the time of matter-radiation energy-density equality as discussed above. One effect of this is to increase the amplitude of the first acoustic peak in the CMB power spectrum at $`l200`$. For example, based upon a $`\chi ^2`$ analysis (Lineweaver & Barbosa (1998)) of 19 experimental points and window functions, Lesgourgues and Pastor (1999) concluded that $`\xi _\nu 6`$ for a single degenerate neutrino species. However, in the existing CMB constraint calculations (Kinney and Riotto 1999: Lesgourgues and Pastor 1999; Hannestad 2000) only small degeneracy parameters with the standard relic neutrino temperatures were utilized in their derived constraint. Hence, the possible effects of a diminished relic neutrino temperature for large neutrino degeneracy need to be reconsidered. To investigate this we have done calculations of the CMB power spectrum, $`\mathrm{\Delta }T^2=l(l+1)C_l/2\pi `$ based upon the CMBFAST code of Seljak & Zaldarriago (1996). For the optimum neutrino-degenerate models ($`\xi _{\nu _{\mu ,\tau }}11`$) and a neutrino contribution $`\mathrm{\Omega }_\nu 0.25`$ we deduce from the solid curve on Figure 9 that the neutrino mass is $`m_{\nu _{\mu ,\tau }}0.1`$ eV and therefore unimportant during the photon decoupling epoch. Therefore, we only consider massless neutrinos here. For massless neutrinos it can be proven (Lesgourgues & Pastor (1999)) that the only effect of neutrino degeneracy on the CMB is to increase the background pressure and energy density of relativistic particles (cf. Eqs. 1-3). We have in this way explicitly modified CMBFAST code to account for the contribution of massless degenerate neutrinos with a relic temperature ratio $`y_\nu =T_\gamma /T_\nu `$ as given in Figure 4. We have evaluated $`\chi ^2`$ for fits to the CMB power spectrum, based upon the ”radical compression” technique as described in Bond, Jaffe & Knox (2000). We have used the latest 69 observational points and window functions available from the web page given in that paper. The advantage of this approach is that the non-Gaussian experimental uncertainties in the power spectrum are correctly weighted in the evaluation of the goodness of fit. For the purposes of the present study, we take as a benchmark the ”All” case best fit $`\mathrm{\Omega }=1`$ model of Dodelson & Knox (2000) who derived cosmological parameters based upon this same data set and compression technique. Although there is some degeneracy in the cosmological parameters they deduced an optimum fit to the power spectrum for $`\mathrm{\Omega }_bh^2=0.019`$, $`H_0=65`$ km sec<sup>-1</sup> Mpc<sup>-1</sup>, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.69`$, $`\mathrm{\Omega }_M=0.31`$, $`\tau =0.17`$, and $`n=1.12`$, where $`\mathrm{\Omega }_M`$ is the total matter contribution, $`\tau `$ is the reionization parameter, and $`n`$ is the ”tilt” of the power spectrum. This benchmark is plotted as the dashed curve in Figure 14. For this case we find $`\chi ^2=101`$. \[Note that our $`\chi ^2`$ is slightly different from that quoted in Dodelson & Knox (2000) because we use different binning of the power spectrum\]. For comparison the ”radical compression” of the CMB data into 14 bins used in this work is also shown (Bond et al. (2000)). Rather than to do an exhaustive parameter search we have taken an approach similar to Lesgourgues & Pastor (1999). That is, we fix several representative cosmological models and then study their goodness of fit to the CMB data. The best case for large neutrino degeneracy will be for a value of the degeneracy parameter $`\xi _{\nu _{\mu ,\tau }}`$ such that neutrino decoupling occurs just before the QCD phase transition. This is the value for which the relic neutrino energy density is a local minimum (cf. Figure 13). For the present work this corresponds to $`\xi _{\nu _{\mu ,\tau }}=11.4`$, $`\xi _{\nu _e}=0.73`$, $`\Omega _\mathrm{b}h_{50}^2=0.187`$ models. In what follows we fix $`\xi _{\nu _{\mu ,\tau }}`$, $`\xi _{\nu _e}`$, and $`\Omega _\mathrm{b}`$$`h_{50}^2`$ at these values and refer to this as the large degeneracy model. We have found \[as did Lesgourgues & Pastor (1999)\] that the currently favored $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ models give a poor fit to the data even with no degeneracy. Adding neutrino degeneracy to an $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ model only makes the fit worse. The main problem is that the first acoustic peak increases in amplitude and moves to larger $`l`$. Hence, even though a local minimum develops for large degeneracy, the $`\chi ^2`$ is substantially increased and large neutrino degeneracy is probably ruled out for $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$ models. For smaller $`\mathrm{\Omega }_\mathrm{\Lambda }`$ a local minimum develops in the $`\chi ^2`$ for both small values of degeneracy $`\xi _{\nu _{\mu ,\tau }}1`$ and large degeneracy $`\xi _{\nu _{\mu ,\tau }}=11.4`$. As pointed out in Lesgourgues & Pastor (1999), the best case for neutrino degeneracy is with $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ models. However, those models are probably ruled out by observations of type Ia supernovae at high redshift (Garnavich, et al. (1998); Perlmutter et al. 1998a ; Perlmutter et al. 1998b ; Riess et al. (1998)). At the $`3\sigma `$ confidence level for $`\mathrm{\Omega }=1`$ models, the type Ia Supernova data are consistent with $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7\pm 0.3`$ Hence, we take $`\mathrm{\Omega }_\mathrm{\Lambda }=0.4`$ as a plausible cosmological model which is marginally consistent with the type Ia results. Nevertheless, for purposes of illustration, we have also made a search for optimum parameters for matter dominated $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$, $`\mathrm{\Omega }=1`$ models. The reason low $`\mathrm{\Omega }_\mathrm{\Lambda }`$ models are favored is that they shift the first acoustic peak back to lower $`l`$. Larger values of $`H_0`$ also slightly improve the fit by shifting the first acoustic peak to lower $`l`$ and decreasing the baryon density for fixed $`\Omega _\mathrm{b}`$$`h_{50}^2`$ which lowers the peak amplitude. We take $`H_0=65\pm 10`$ ($`h_{50}=1.3\pm 0.2`$) as a reasonable range (Dodelson & Knox (2000)), and therefore utilize $`h_{50}=1.5`$ as the optimum Hubble parameter for the neutrino-degenerate models. This implies that $`\mathrm{\Omega }_b=0.084`$ for the large degeneracy models. The ionization parameter does not particularly help the fits as it mainly serves to decrease the amplitude of both the first and second peaks in the power spectrum. We therefore set $`\tau =0`$ for the large degeneracy models. The only remaining adjustable parameter of the fits is then the tilt parameter $`n`$. Values of $`n`$ slightly below unity also help with the amplitude and location of the first acoustic peak. The solid line on Figure 14 shows a $`\mathrm{\Omega }_\mathrm{\Lambda }=0.4`$ model for which $`n=0.78`$. For this fit $`\mathrm{\Delta }\chi ^2=27`$ which makes this large degeneracy model marginally consistent with the data at a level of $`5.2\sigma `$. The dotted line in Figure 14 shows the matter dominated $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ best fit model with $`n=0.83`$. For this fit $`\mathrm{\Delta }\chi ^2=9`$ which makes this large degeneracy model consistent with the data at the level of $`3\sigma `$. As can be seen from Figure 14, a model with large neutrino degeneracy seems marginally acceptable based upon the presently uncertain power spectrum. The main differences in the fits between the large degeneracy models and our adopted benchmark model are that the first peak is shifted to slightly higher $`l`$ values and the second peak is somewhat suppressed. It thus becomes important to quantify the amplitude of the second peak in order to constrain the large degeneracy models proposed herein. Indeed, after the present fits were completed a suppression of the second acoustic peak in the power spectrum was reported in the high-resolution BOOMERANG (Bernardls, P. et al. (2000); Lange et al. (2000)) and MAXIMA-1 (Hanany et al. (2000): Balbi et al. (2000)) results. We have not yet analyzed the goodness of fit to these data as the experimental window functions are not yet available. In a subsequent paper we will examine the implications of those data in detail. For purposes of illustration, however, we compare the fit models of Figure 14 with the published BOOMERANG and MAXIMA-1 power spectra in Figure 15. Here one can clearly see that the suppression of the second acoustic peak in the observed power spectrum is consistent with our derived neutrino-degenerate models. In particular, the MAXIMA-1 results are in very good agreement with the predictions of the neutrino-degenerate cosmological models described herein. There is, however, a calibration uncertainty between these two sets (Hanany et al. (2000)). If one only considers the BOOMERANG results alone, the diminished amplitude of the first acoustic peak probably tightens the constraint for low neutrino degeneracy models (cf. Hannestad (2000)) although even for this set alone, a high degeneracy model is probably still acceptable (Lesgourgues & Peloso (2000)). It is clear, that these new data sets will substantially improve the goodness of fit for the neutrino-degenerate models (Lesgourgues & Peloso (2000)). Moreover, both data sets seem to require an increase in the baryonic contribution to the closure density as allowed in our neutrino-degenerate models. ## 5 CONCLUSIONS We have discussed how the relic neutrino temperature is substantially diminished in cosmological models with a large neutrino degeneracy. We have shown that all of the BBN light-element abundance constraints (assuming some destruction of <sup>7</sup>Li) can be satisfied for a single-parameter family of cosmological models in which significant neutrino degeneracies and large values of $`\mathrm{\Omega }_bh^2`$ can exist. The requirement that large scale structure become nonlinear in sufficient time can also be satisfied for models with either moderate degeneracy (ξνμ,τ <6.2 <subscript𝜉subscript𝜈𝜇𝜏6.2\hbox{$\xi_{\nu_{\mu,\tau}}$}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}6.2 and Ωbh502 <0.22 <Ωbh5020.22\hbox{${\sl\Omega}_{\rm b}$}\hbox{$h^{2}_{50}$}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}0.22) or large neutrino degeneracy (11.4 <ξνμ,τ <13 <11.4subscript𝜉subscript𝜈𝜇𝜏 <1311.4\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}\hbox{$\xi_{\nu_{\mu,\tau}}$}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}13 and 0.187 <Ωbh502 <0.25 <0.187Ωbh502 <0.250.187\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}\hbox{${\sl\Omega}_{\rm b}$}\hbox{$h^{2}_{50}$}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}0.25). We have also shown that even the constraint from neutrino-degeneracy effects on fluctuations of the cosmic microwave background temperature may be marginally avoided for models with ΩΛ <0.4 <subscriptΩΛ0.4\Omega_{\Lambda}\mathrel{\mathchoice{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\displaystyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\textstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}{\lower 3.6pt\vbox{\halign{$\mathsurround=0pt\scriptscriptstyle\hfil#\hfil$\cr<\crcr\sim\crcr}}}}0.4, $`\xi _{\nu _{\mu ,\tau }}11`$, and $`\Omega _\mathrm{b}h_{50}^20.2`$. At present, the power spectrum of the CMB is the most stringent constraint. Nevertheless, neutrino-degenerate models can be found which are marginally consistent at the 3-5$`\sigma `$ level. This tight constraint is due, at least in part, to a suppression of the second acoustic peak in the spectrum. It is therefore encouraging that the recent BOOMERANG and MAXIMA-1 results suggest that such a suppression in the second acoustic peak may indeed occur in agreement with the expectations of the large neutrino degeneracy, high $`\Omega _\mathrm{b}`$ models proposed here. Thus, high resolution microwave background observations become even more important as a means to quantify the limits to (or existence of) possible cosmological neutrino degeneracy. Based upon the current analysis, we conclude that all of the cosmological constraints on large neutrino degeneracy are marginally satisfied when a careful accounting of the neutrino decoupling and relic neutrino temperature is made. It will, therefore, be most interesting to see what further constraints can be placed on this possibility from the soon to be available space-based high resolution CMB observations such as the NASA MAP and ESA Planck missions. One of the authors (GJM) wishes to acknowledge the hospitality of the National Astronomical Observatory of Japan where much of this work was done. This work has been supported in part by the Grant-in-Aid for Scientific Research (10640236, 10044103, 11127220, 12047233) of the Ministry of Education, Science, Sports, and Culture of Japan, and also in part by NSF Grant (PHY-9901241 at OSU) along with DoE Nuclear Theory Grant (DE-FG02-95-ER40394 at UND).
warning/0005/cs0005027.html	ar5iv	text	# 1 Overview ## 1 Overview The paper begins by reviewing traditional approaches to surface representation and inference. Then the new field representation and inference paradigm is introduced within the context of maximally informative (MI) inference , early ideas appearing in . The knowledge representation distribution is introduced and discussed in the context of MI inference. Then, using the MI inference approach, the here-named Generalized Kalman Filter (GKF) equations are derived for a specific example instance of inferring a surface height field. The GKF equations motivate a location-dependent adaptive scale or multigrid approach to the MI inference of continuous-basis fields. ## 2 Introduction: Surface representation ### 2.1 Traditional methods Many methods for representing surfaces have been utilized previously, however these methods involve representing the surface by a discrete basis field, perhaps with a deterministic interpolation defined (bi-linear, tensor B-splines, etc.) to provide a definition for the surface at points intermediate to the discrete field. Probability distributions or densities of these discrete fields then often take the form of normalized exponentials of sums of clique energy functions, and produce a construct commonly known as a Markov Random Field. (See Geman , for an often cited example.) There are several immediate observations on these approaches: * The surface remains unspecified at points intermediate to the discrete field, except by the often undefined notion of interpolation. * When interpolation is not defined, the discrete field probability distribution says nothing about the probability distribution of surface at points intermediate to the discrete field points. * When interpolation is defined then, given a value of the discrete field, there is no uncertainty in the surface intermediate to the discrete field points. There is a deterministic mapping from any given discrete field to the corresponding continuous surface. In particular, when the discrete field basis covers a fixed grid on the $`(x,y)`$ plane with $`z`$ heights at each grid point, known here as a height field, all sampling of the surface intermediate to the fixed grid is determined at the scale of the fixed grid. This is generally not physical, see next. * The surface distribution is not an intrinsic property of any physical surface, rather a post-hoc imposition of the analyst attempting a useful regularization. For instance, necessary scaling properties are ignored: Moving a camera closer to the surface, for example, so that the density of sample points on the physical surface increases, is not properly represented in the fixed basis of the discrete field distribution; there is no consistency imposed that requires a subsampled set of points to have the same probability density that one would find by marginalizing the surface distribution over the sample points not in the subsampling. ### 2.2 Scaling consistency The consistency condition mentioned in the last section, which must be imposed on probability distributions for continuous fields is: > Scaling of sample points consistency: For $`SA`$ indices of discrete field variables, > > $$P(X_S)=P(X_A)𝑑X_{AS}$$ > (1) Note that equation 1 is a condition which must be imposed on the distributions which any modelling system learns where it is sensible to supersample or subsample the field arbitrarily, as in the continuous field basis case. ### 2.3 Elements of the paradigm The rest of this paper discusses an approach to continuous field inference which corrects the deficiencies, including the intermediate value and scaling problems, of traditional discrete-basis approaches to the inference of discrete height fields, for example. The new approach is here named the Generalized Kalman Filter. There are four central objects of importance within the inference approach described in this paper, one of which is a new object to Bayesian inference: * The prior distribution for field. The prior holds all information about fields before any data is observed. * The likelihood distribution. The likelihood is predictive for data, given the field. It incorporates all of the physics of the measurement process. * The posterior distribution. The posterior distribution summarizes everything knowable about the field given assumptions of likelihood form, the prior knowledge, and all data. * The knowledge-representation (KR) distribution. Within the usual Bayesian point of view, the KR distribution is the new mathematical object. In the paradigm described in this paper the KR distribution is the object updated when new data arrives. The KR distribution is parameterized by maximally informative statistics (see ) for the learned field knowledge. Note that because the KR distribution has a finite number-of-values limitation, the KR distribution is not necessarily able to represent what could have been learned from data about the (continuous) field. Generally, the prior distribution and the KR distribution determine an approximation (possibly exact) to the field posterior distribution. It should be noted that modern computer architecture (memory and space-time) constraints appear to be the fundamental physical drivers for the utilization of the KR distribution, simply because storing the exact posterior generally requires an infinite amount of memory. In the height field inference application discussed later the KR distribution is parameterized by heights at a set of discrete basis points, but holds knowledge about a continuous basis height field. However, generally, the KR distribution may use an arbitrary set of basis functions. One advance of the GKF is that the KR distribution is naturally adaptive in both dimension and scale, allowing the learning of continuous-basis field information at the appropriate scale, where appropriate. Benefits of the approach described in this paper are that it has these information theoretically optimal features: 1. A location-dependent adaptive and scalable multigrid-like algorithm, so that only the bytes necessary to represent the learned information are stored, leading to a style of maximally sparse representation of surface knowledge; 2. A recursive updating algorithm. It will become clear that the Bayesian GKF field inference paradigm also has these properties: * It is the information learned about the field, (the KR distribution), which takes the form of a distribution over discrete values. In the surface inference example these discrete values are heights at discrete basis points. * The prior distribution for fields, in conjunction with the learned knowledge of the field held within the KR distribution determine a well-defined posterior distribution over continuous fields. * The field posterior distribution is always a well defined quantity everywhere. In the surface inference example discussed later, this continuity is at points intermediate to the discrete height field basis points of the KR distribution. * The scaling condition equation 1 is automatically imposed because the posterior distribution is a distribution over fields. As an example consider the inference of continuous surfaces: While it may seem obvious, in the case of continuous surface inference, that what one is actually representing with a discrete set of values in memory is only a part of the information which helps to determine the surface posterior distribution, it is unusual to not be discussing the height field as the primary representation of surface. It is the inherently discrete nature of the storage of information in machines which forces us into this stance - generally it is impossible to represent an arbitrary continuous field with a finite set of discrete values - one must also have another object from which to compute the intermediate values of the field. (Another way to look at the disparity between the current proposal for field inference and traditional proposals is that the traditional approaches are sufficient only for band-limited fields.) In section 3 the GKF is specialized to height fields, where an example, surface representation and learning, of the GKF paradigm is described. (The approach taken in this section is to specialize to a case that is then easily seen to generalize to the general continuous basis field inference paradigm.) The next section continues with observations on the update scheme. Further sections continue with the example special case for surface distributions with particularly tractable mathematics, and final sections provide explicit forms for the general GKF equations, a discussion on their relationship to the standard Kalman filter, a discussion on the amount of information learned at each update, and a search heuristic. Extensive appendices provide supporting mathematics for the derivations. ## 3 Surface representation and inference In this section the main ideas of the Bayesian surface representation and inference paradigm presented in this paper are given. The technique is general, though: section 4 discusses the extension to an arbitrary-basis, arbitrary-dimension field. ### 3.1 Surface distributions The surface and height field distributions (the prior, likelihood, and posterior surface and height field distributions) are discussed in this section. #### 3.1.1 Surface and height field prior distributions Consider a set $`S`$ of surfaces where each element $`𝒔S`$ is a height field, i.e. such that $`𝒔=s(x,y)`$ is real function of two variables. Write the prior probability distribution for surfaces in $`S`$ given the parameters $`\theta `$ which determine the prior distribution as $$P(𝒔\theta ).$$ (2) Consider a vector $`𝒗=(v_1,\mathrm{},v_n)`$ of discrete $`(x,y)`$ points, $`v_i=(x_i,y_i)`$. For any given surface $`𝒔`$ denote the associated vector of heights by $`𝒉(𝒔,𝒗)=(h_1(𝒔,𝒗),\mathrm{},h_n(𝒔,𝒗))`$. Write the prior distribution of the surface heights at the chosen points $`𝒗`$ as $`P(𝒉_v\theta )`$. This discrete height distribution may be found as follows: $`P(𝒉_v\theta )`$ $`=`$ $`{\displaystyle P(𝒉_v𝒔,\theta )P(𝒔\theta )𝑑𝒔}`$ (3) $`=`$ $`{\displaystyle P(𝒉_v𝒔)P(𝒔\theta )𝑑𝒔}`$ (4) $`=`$ $`{\displaystyle 𝜹(𝒉_v𝒉(𝒔,𝒗))P(𝒔\theta )𝑑𝒔}`$ (5) where the vector delta-function is defined as $`𝜹(𝒉_v𝒉(𝒔,𝒗))`$ $`=`$ $`\mathrm{\Pi }_{i=1}^n\delta (h_{v,i}h_i(𝒔,𝒗))`$ (6) Now, given that what is known is the surface heights $`𝒉_v`$ at a vector $`𝒗`$ of discrete $`(x,y)`$ points, the posterior distribution of surfaces is found from Bayes’ theorem as $`P(𝒔𝒉_v,\theta )`$ $`=`$ $`{\displaystyle \frac{P(𝒉_v𝒔,\theta )P(𝒔\theta )}{P(𝒉_v\theta )}}`$ (7) $`=`$ $`{\displaystyle \frac{P(𝒉_vs)P(𝒔\theta )}{P(𝒉_v\theta )}}`$ (8) $`=`$ $`{\displaystyle \frac{𝜹(𝒉_v𝒉(𝒔,𝒗))P(𝒔\theta )}{𝜹(𝒉_v𝒉(𝒔,𝒗))P(𝒔\theta )𝑑𝒔}}`$ (9) where the denominator distribution was found in equation 5. #### 3.1.2 Measurements: The Likelihood In general, a surface $`𝒔`$ and some other parameters $`\varphi `$ not dependent upon $`𝒔`$ (i.e. camera point spread function, camera position and direction, lighting position and direction, etc.) specify the probability distribution for data (likelihood) $`P(𝒙𝒔,\varphi ,\theta )=P(𝒙𝒔,\varphi )`$ (10) where the data distribution is independent of $`\theta `$ once $`𝒔`$ is known. #### 3.1.3 Conditioning on data: Surface and height field posterior distributions Given data, the surface posterior distribution is inferred using Bayes’ theorem as $`P(𝒔𝒙,\varphi ,\theta )`$ $`=`$ $`{\displaystyle \frac{P(𝒙𝒔,\varphi ,\theta )P(𝒔\varphi ,\theta )}{P(𝒙\varphi ,\theta )}}`$ (11) $`=`$ $`{\displaystyle \frac{P(𝒙𝒔,\varphi )P(𝒔\theta )}{P(𝒙𝒔,\varphi )P(𝒔\theta )𝑑𝒔}}`$ (12) The distribution of the surface posterior marginalized to a set of discrete points may be written using equations 1112, doing steps similar to those taken in equations 35, as $`P(𝒉_v𝒙,\varphi ,\theta )`$ $`=`$ $`{\displaystyle P(𝒉_v𝒔,𝒙,\varphi ,\theta )P(𝒔𝒙,\varphi ,\theta )𝑑𝒔}`$ (13) $`=`$ $`{\displaystyle P(𝒉_v𝒔)P(𝒔𝒙,\varphi ,\theta )𝑑𝒔}`$ (14) $`=`$ $`{\displaystyle 𝜹(𝒉_v𝒉(𝒔,𝒗))P(𝒔𝒙,\varphi ,\theta )𝑑𝒔}`$ (15) In steps similar to equations 79 the surface posterior when a height field is also known is given by $`P(𝒔𝒉_v,𝒙,\varphi ,\theta )`$ $`=`$ $`{\displaystyle \frac{P(𝒉_v,𝒙𝒔,\varphi ,\theta )P(𝒔\varphi ,\theta )}{P(𝒉_v,𝒙\varphi ,\theta )}}`$ (16) $`=`$ $`{\displaystyle \frac{P(𝒉_v𝒔)P(𝒙𝒔,\varphi )P(𝒔\theta )}{P(𝒉_v,𝒙\varphi ,\theta )}}`$ (17) $`=`$ $`{\displaystyle \frac{𝜹(𝒉_v𝒉(𝒔,𝒗))P(𝒙𝒔,\varphi )P(𝒔\theta )}{𝜹(𝒉_v𝒉(𝒔,𝒗))P(𝒙𝒔,\varphi )P(𝒔\theta )𝑑𝒔}}`$ (18) where we used the facts that, given a surface, the data and the surface heights are independent, and the surface distribution is independent of the camera and lighting parameters $`\varphi `$. ### 3.2 Approximating the posterior One motivation for approximating the surface distribution is that generally a surface is an uncountably infinite, continuous entity, and therefore there is little else which can be done to represent it exactly other than to go into, literally, infinite detail (requiring an infinite supply of memory). It is therefore useful to have an approximation scheme which, although finite, captures the relevant information provided by data. Another excellent reason for developing an approximation is mathematical tractability. Having a representation scheme which allows a tractable calculation of the posterior is a huge benefit for both computation and communication. Finally, it is of great interest to not waste computational resources while representing learned surface information. The solution to the surface representation problem presented here addresses the competition for representational resources (memory) issue in a unique manner. #### 3.2.1 The knowledge representation distribution The full posterior may be written in the form $`P(𝒔𝒙,\varphi ,\theta )`$ $`=`$ $`{\displaystyle P(𝒔𝒉_v,𝒙,\varphi ,\theta )P(𝒉_v𝒙,\varphi ,\theta )𝑑𝒉_v}`$ (19) where the distributions inside the integral appear in equations 1318. The issue of generating a finite representation is not yet resolved via equation 19 however, since storing information sufficient to determine the distributions $`P(𝒔𝒙,\varphi ,\theta )`$, and $`P(𝒔𝒉_v,𝒙,\varphi ,\theta )`$ generally requires storing an infinite set of values in a finite amount of memory, or requires that all data be stored, disallowing any discarding of data and the incremental updating of the representation. Instead, consider the following approximation where the prior conditioned on a set of heights, along with a new distribution, the knowledge representation distribution $`\widehat{P}(𝒉_v𝒙,\varphi ,\theta )`$, are substituted for the distributions inside the integral of equation 19. $`\widehat{P}(𝒔\widehat{P}(𝒉_v𝒙,\varphi ,\theta ))`$ $`=`$ $`{\displaystyle P(𝒔𝒉_v,\theta )\widehat{P}(𝒉_v𝒙,\varphi ,\theta )𝑑𝒉_v}`$ (20) It is important to note at this point that any suitable surface distribution may be substituted into the right-hand side of equation 20 for $`P(𝒔𝒉_v,\theta )`$, since it is important only that the resulting integral be capable of making a good approximation to the true posterior. Further, it is not necessary to restrict the basis $`𝒗`$ to discrete height field basis points, any suitable basis may be taken, for instance Fourier components. Although all of the calculations of this paper are carried thru with the form of 20, other forms may prove more convenient, and it is not difficult to suggest others. In particular, since equation 20 will be used in an iterative update loop later, updates that take for the right-hand side prior term the last posterior term appear quite reasonable (the corresponding GKF update equations may be found immediately from those presented later). Although conditioning on the KR distribution $`\widehat{P}(𝒉_v𝒙,\varphi ,\theta )`$ may seem strange, a good way to understand the meaning is that it is the KR distribution which is being used as a statistic for the learned surface information. The key thing to notice in equation 20 is that, with reasonable regularity conditions, choosing the points of $`𝒗`$ sufficiently dense, the approximation desired to the full posterior may become arbitrarily good. The trick will be to choose $`𝒗`$ appropriately, properly weighting the competing need to approximate arbitrarily well everywhere with the limited resources that are imposed when a finite amount of storage is available, i.e. when the dimensionality of $`𝒗`$ is fixed. This will be addressed in the next section. In the case of simple imaging systems, the point spread function and pixel diameter are good indicators of the necessary sampling scale for $`𝒗`$. In the super-resolved case, the resolution expected available from the data is the appropriate scale for $`𝒗`$. The approximation to the posterior of 20 has several properties which make it valuable: * The prior distribution $`P(𝒔𝒉_v,\theta )`$ which supplies the uncertainties associated with points of the surface not in the vector $`𝒗`$ may be chosen to have a simple form (see appendix 12.1) that is easily encoded algorithmically in finite memory. * There is a clear separation between what was already known - the prior $`P(𝒔𝒉_v,\theta )`$, and what has been learned - the KR distribution $`\widehat{P}(𝒉_v𝒙,\varphi ,\theta )`$. * There is a clear description of the scale at which information has been acquired in terms of the density and uncertainties associated with the points $`(𝒗,h(𝒔,𝒗))`$ on the surface, and in terms of the uncertainties of their positions as encoded in the KR distribution. In practice, it is useful to take a multinormal distribution over the discrete-point height field as the KR distribution. Let the parameterization of the KR distribution be $`\mathrm{\Theta }_v`$. For example, if the KR is taken to be multinormal then the parameters of that distribution are $$\mathrm{\Theta }_v(𝒙)=(𝝁_v(𝒙),\mathrm{\Sigma }_v(𝒙)),$$ (21) the mean and covariance matrix of the multinormal, where the functional dependence on $`𝒙`$ indicates a data dependency through the update procedure, and the subscript $`𝒗`$ indicates that the parameters parameterize a distribution of heights at points $`𝒗`$. Because the KR distribution and its parameters are related by a one-to-one mapping, re-write equation 20 as $`\widehat{P}(𝒔\mathrm{\Theta }_v,\theta )`$ $`=`$ $`{\displaystyle P(𝒔𝒉_v,\theta )\widehat{P}(𝒉_v\mathrm{\Theta }_v)𝑑𝒉_v}.`$ (22) In summary, we have arrived at an approximation to the surface posterior distribution, via the KR distribution, parameterized by $`\mathrm{\Theta }_v`$. ### 3.3 Updating the knowledge representation Now we discuss updating $`\mathrm{\Theta }_v`$ when new data are acquired. Temporarily restrict attention to the fixed $`𝒗`$ case. During this and the next sections refer to figure 1 for a flowchart of the general GKF update process. #### 3.3.1 Bayes’ theorem Having acquired $`\mathrm{\Theta }_v^n=\mathrm{\Theta }_v(x^n)`$, from previously seen data $`𝒙^n=(𝒙_1,\mathrm{},𝒙_n)`$ and upon seeing new data $`𝒙_{n+1}`$, the goal is to find $`\mathrm{\Theta }_v^{n+1}`$ such that the surface distribution given $`\mathrm{\Theta }_v^{n+1}`$ approximates the surface distribution given $`𝒙_{n+1}`$ and $`\mathrm{\Theta }_v^n`$. Given new data $`𝒙_{n+1}`$ in the context of the previously seen data $`𝒙^n`$ summarized by $`\mathrm{\Theta }_v^n`$, our updated surface distribution is found via Bayes’ theorem $`\widehat{P}(𝒔𝒙_{n+1},\mathrm{\Theta }_v^n,\varphi ,\theta )`$ $`=`$ $`{\displaystyle \frac{P(𝒙_{n+1}𝒔,\mathrm{\Theta }_v^n,\varphi ,\theta )\widehat{P}(𝒔\mathrm{\Theta }_v^n,\varphi ,\theta )}{\widehat{P}(𝒙_{n+1}\mathrm{\Theta }_v^n,\varphi ,\theta )}}`$ (23) $`=`$ $`{\displaystyle \frac{P(𝒙_{n+1}𝒔,\varphi )\widehat{P}(𝒔\mathrm{\Theta }_v^n,\theta )}{\widehat{P}(𝒙_{n+1}\mathrm{\Theta }_v^n,\varphi ,\theta )}}`$ $`=`$ $`{\displaystyle \frac{P(𝒙_{n+1}𝒔,\varphi )\widehat{P}(𝒔\mathrm{\Theta }_v^n,\theta )}{P(𝒙_{n+1}𝒔,\varphi )\widehat{P}(𝒔\mathrm{\Theta }_v^n,\theta )𝑑𝒔}}`$ where we defined $$\widehat{P}(𝒙_{n+1}\mathrm{\Theta }_v^n,\varphi ,\theta )=P(𝒙_{n+1}𝒔,\varphi )\widehat{P}(𝒔\mathrm{\Theta }_v^n,\theta )𝑑𝒔.$$ (24) The updated posterior $`\widehat{P}(𝒔\mathrm{\Theta }_v^n,𝒙_{n+1},\varphi ,\theta )`$ will be approximated by the $`\mathrm{\Theta }_v^{n+1}`$ parameterized KR distribution of equation 22 as $`\widehat{P}(𝒔\mathrm{\Theta }_v^{n+1},\theta )`$ $`=`$ $`{\displaystyle P(𝒔𝒉_v,\theta )\widehat{P}(𝒉_v\mathrm{\Theta }_v^{n+1})𝑑𝒉_v}.`$ (25) The approximation condition for determining $`\mathrm{\Theta }_v^{n+1}`$ is then written $`\widehat{P}(𝒔\mathrm{\Theta }_v^{n+1},\theta )`$ $``$ $`\widehat{P}(𝒔𝒙_{n+1},\mathrm{\Theta }_v^n,\varphi ,\theta )`$ (26) Equation 26 suggests we try to minimize various measures of the closeness of the two distributions. For example, one measure is the average square difference of the two distributions, $$\left\|P_1(𝒔)P_2(𝒔)\right\|^2𝑑𝒔$$ (27) but there is (apparently) no good first-principles reason to use this form. In the next section we discuss the measure of distance which lead$`𝒔`$ to the maximally informative choice of $`\mathrm{\Theta }_v^{n+1}`$. #### 3.3.2 Maximally informative inference The measure of distance which leads to the $`\mathrm{\Theta }^{n+1}`$ providing the most information about the surface distribution is the maximally informative choice for the statistic $`\mathrm{\Theta }^{n+1}`$. The condition for being maximally informative, see , is that the Kullback-Leibler distance $`D(P_1(𝒔),P_2(𝒔))`$ is minimized, where $`D(P_1(𝒔),P_2(𝒔))={\displaystyle P_1(𝒔)log\left(\frac{P_1(𝒔)}{P_2(𝒔)}\right)𝑑𝒔}`$ (28) and where the $`P`$’s above are posterior distributions of field, that is $`P_1(𝒔)`$ $`=`$ $`\widehat{P}(𝒔𝒙_{n+1},\mathrm{\Theta }_v^n,\varphi ,\theta )`$ (29) $`P_2(𝒔)`$ $`=`$ $`\widehat{P}(𝒔\mathrm{\Theta }_v^{n+1},\theta ).`$ (30) That is, Find the $`\mathrm{\Theta }^{n+1}`$ such that $`_{\mathrm{\Theta }_v^{n+1}}{\displaystyle \widehat{P}(𝒔\mathrm{\Theta }_v^n,𝒙_{n+1},\varphi ,\theta )log\left(\frac{\widehat{P}(𝒔\mathrm{\Theta }_v^n,𝒙_{n+1},\varphi ,\theta )}{\widehat{P}(𝒔\mathrm{\Theta }_v^{n+1},\theta )}\right)𝑑𝒔}=\mathrm{𝟎}`$ while at the $`\mathrm{\Theta }_v^{n+1}`$ satisfying the derivative condition above $`det\left[_{\mathrm{\Theta }_v^{n+1}}^2{\displaystyle \widehat{P}(𝒔\mathrm{\Theta }_v^n,𝒙_{n+1},\varphi ,\theta )log\left(\frac{\widehat{P}(𝒔\mathrm{\Theta }_v^n,𝒙_{n+1},\varphi ,\theta )}{\widehat{P}(𝒔\mathrm{\Theta }_v^{n+1},\theta )}\right)𝑑𝒔}\right]<0`$ i.e., the hessian is negative definite and the extremum is a local maximum. If possible, choose the global maximum. Note that the Kullback-Leibler distance is asymmetric. Generally, it is highly relevant which distribution contains the prior information and which distribution is being updated. Maximum entropy techniques reverse the roles of $`P_1`$ and $`P_2`$ which appear here. For a detailed explanation see . In the following section are some observations on the approach taken to maximally informative surface inference. Section 5 then briefly makes explicit the specific distribution forms which are assumed. The Generalized Kalman Filter update equations for the surface inference example which follow from this approach are then presented in section 6, completing the derivation of the maximally informative approach. ## 4 Observations on the update scheme Note the following: * The updating scheme described here is a maximally informative update scheme and is related to the Kalman filter. The Kalman filter is a minimum variance filtering scheme applicable in the case of fixed representation dimension. The crucial step which has been taken in the current work is the step of allowing the representation scheme to be adaptable. We have adopted the label “Generalized Kalman Filter” (GKF) to describe the idea represented here. The GKF equations are presented in section 6. * To this point we have only optimized over $`\mathrm{\Theta }_v`$. It is clear that we may also vary the number of vertices $`\left\|𝒗\right\|`$ of the representation, allowing optimization over the number of vertices. Varying the number of vertices of the representation is absolutely necessary if surface knowledge at scales smaller than the current set of vertices represents is to ever accumulate. In section 6 the GKF update equations are derived assuming that the number of vertices in the representation basis vertex set is arbitrary at each update. * Beyond allowing the number of vertices to vary, the positions of the vertices may be allowed to vary. In section 6 the GKF update equations are derived assuming that the representation basis vertex set positions are arbitrary. * Detecting when and where new vertices are necessary is a matter of observing directly in equations 28 or 3.3.2 when new data produces a lower surface uncertainty over a region, and when having smaller uncertainty at neighboring vertices is not sufficient to represent this lower uncertainty over the region. * The vertex representation for the surface knowledge is convenient, but not necessary. For example it is possible to extend a height field to a height-and-reflectance field or “arbitrary dimension field”, where the reflectance lies within a many-dimensional space. Reasonable structures for the covariance matrix allow differing correlations between reflectance values and between height values. It will be seen in in section 6 that the GKF update equations are easily used in the “arbitrary dimension field” context. * In its most abstract form, instead of having a “field”, there is simply a set of objects, while for each “object” there is an associated vector of properties, where some of the components of the property vector may be considered a location in space. In this fairly abstracted setting, the collection of objects has an associated joint probability distribution which describes the probability distribution over configurations of objects. It will be seen in in section 6 that the GKF update equations are easily understood in the “object” context. * Equation 3.3.2 which defines the quantity to be minimized is where a penalty term which indicates how many bits in hardware is available in trade for each bit of information learned from data. For example, one might penalize the KL distance by $`1/10`$th the number of bytes it takes to represent the new information gained by extending the number of points represented. The exact form of the information learned about the surface distribution contained in the KR distribution is found in section 8, where the dimensionality of the representation enters directly, and where bits-used penalty-terms may be introduced. * The previous note points out how a minimum description length method fails for this problem. It is certainly the case that that our update scheme may require much more memory (in bits) to represent the information learned than the information learned (in bits). At some point, if information at small enough scales is desired, MDL would truncate and stop. Clearly, applying MDL would then be a disaster. On the other hand, what seems to work here may be called an adaptive MDL approach. * Note that a method like maximum entropy is entirely deficient for providing distributions of surfaces: given the constraints implied by the knowledge of the distribution of the heights at discrete points: maximum entropy ignores correlations between nearby surface points no matter how close, an entirely ludicrous situation. On the other hand, a method like relative maximum entropy, based on inverting the roles of the distributions in equation 28, claims to provide the least informative inference relative to the prior information, a heuristic, difficult to justify, at best. Further, such approaches are typically based on likelihood distributions, rather than the posteriors that appear in equation 28. ## 5 Surface Distribution Forms ### 5.1 Prior For simplicity of mathematical presentation only, the prior in our surface inference example is taken multinormal over continuous, smooth height fields. One particular, conveniently chosen, representation of the prior distribution is constructed in appendix 12.1. This prior may be written in the shorthand $$P(𝒔\theta )=N(𝝁_s,\mathrm{\Sigma }_s)(𝒔)$$ (33) where $`\theta =(𝝁_s,\mathrm{\Sigma }_s)`$ is the parameter vector. The density of the height field determined by the prior $`P(𝒉_v\theta )`$ $`=`$ $`{\displaystyle P(𝒉_v𝒔)P(𝒔\theta )𝑑𝒔}`$ (34) $`=`$ $`{\displaystyle 𝜹(𝒉_v𝒉(𝒔,𝒗))P(𝒔\theta )𝑑𝒔}`$ (35) $`=`$ $`N(𝝁_v,\mathrm{\Sigma }_v)(𝒉_v)`$ (36) where $`𝝁_v=A_{vs}𝝁_s`$ $`\mathrm{\Sigma }_v=A_{vs}\mathrm{\Sigma }_sA_{vs}^T`$ (37) and the projection onto the height field is given by $`A_{vs}`$. Note that equation 37 implies that the surface density covariance is represented differently than a discrete surface distribution covariance matrix. Specifically, the projection matrix $`A_{vs}`$ is a delta-function-like operator, and $`\mathrm{\Sigma }_s`$ is a continuous function of two positions. In appendix 12.1 we show that the surface density has a compact continuous power spectrum representation, and there give the explicit form of that representation. Thus the notation of equation 37 must be considered a shorthand for the underlying continuous construct. ### 5.2 Likelihood When measurement is modelled as a linear process corrupted by gaussian noise we have $`𝒙`$ $`=`$ $`M𝒔+ϵ`$ $`ϵ`$ $``$ $`N(\mathrm{𝟎},\mathrm{\Sigma }_ϵ).`$ (38) or $`P(𝒙𝒔,\varphi )=N(M𝒔,\mathrm{\Sigma }_ϵ)(𝒙)`$ (39) where $`\varphi =(M,\mathrm{\Sigma }_ϵ)`$ is the parameter vector. ## 6 The Generalized Kalman Filter equations. In this section a concise derivation of the Generalized Kalman Filter update equations specialized to the discrete basis multinormal KR distribution of equation 22 are derived. The updated KR need not have the same basis dimension nor position as the previous KR basis, solving the problem of how to allow updates from one representation to the next, same, finer or coarser, representation. Proceeding, the KR distribution in terms of the parameterized height field of equation 22 is $$\widehat{P}(𝒔\mathrm{\Theta }_v^n,\theta )=P(𝒔𝒉_v,\theta )\widehat{P}(𝒉_v\mathrm{\Theta }_v^n)𝑑𝒉_v$$ (40) The distribution of surface given the height field from equation 9 is $`P(𝒔𝒉_v\theta )`$ $`=`$ $`{\displaystyle \frac{P(𝒉_v𝒔)P(𝒔\theta )}{P(𝒉_v\theta )}}`$ (41) $`=`$ $`{\displaystyle \frac{𝜹(𝒉_v𝒉(𝒔,𝒗))P(𝒔\theta )}{P(𝒉_v\theta )}}`$ Simplify the integral of the KR distribution to find $`\widehat{P}(𝒔\mathrm{\Theta }_v^n,\theta )`$ $`=`$ $`{\displaystyle \frac{P(𝒉_v𝒔)P(𝒔\theta )}{P(𝒉_v\theta )}\widehat{P}(𝒉_v\mathrm{\Theta }_v^n)𝑑𝒉_v}`$ (42) $`=`$ $`P(𝒔\theta ){\displaystyle 𝜹(𝒉_v𝒉(𝒔,𝒗))\frac{\widehat{P}(𝒉_v\mathrm{\Theta }_v^n)}{P(𝒉_v\theta )}𝑑𝒉_v}`$ $`=`$ $`P(𝒔\theta ){\displaystyle \frac{\widehat{P}(𝒉(𝒔,𝒗)\mathrm{\Theta }_v^n)}{P(𝒉(𝒔,𝒗)\theta )}}`$ Note how the full surface distribution is simply modified by the ratio $`{\displaystyle \frac{\widehat{P}(𝒉(𝒔,𝒗)\mathrm{\Theta }_v^n)}{P(𝒉(𝒔,𝒗)\theta )}}`$ (43) From equation 23 the Bayesian update of the KR distribution is $`\widehat{P}(𝒔𝒙_{n+1},\mathrm{\Theta }_v^n,\varphi ,\theta )`$ $`=`$ $`{\displaystyle \frac{P(𝒙_{n+1}𝒔,\varphi )\widehat{P}(𝒔\mathrm{\Theta }_v^n,\theta )}{P(𝒙_{n+1}𝒔,\varphi )\widehat{P}(𝒔\mathrm{\Theta }_v^n,\theta )𝑑𝒔}}`$ (44) $`=`$ $`{\displaystyle \frac{P(𝒙_{n+1}𝒔,\varphi )\widehat{P}(𝒔\mathrm{\Theta }_v^n,\theta )}{\widehat{P}(𝒙_{n+1}\mathrm{\Theta }_v^n,\varphi ,\theta )}}`$ Rewriting the updated distribution using equation 42 yields $`\widehat{P}(𝒔𝒙_{n+1},\mathrm{\Theta }_v^n,\varphi ,\theta )`$ $``$ $`P(𝒙_{n+1}𝒔,\varphi )P(𝒔\theta )\times {\displaystyle \frac{\widehat{P}(𝒉(𝒔,𝒗)\mathrm{\Theta }_v^n)}{P(𝒉(𝒔,𝒗)\theta )}}`$ For maximally informative inference of the new KR we minimize, from equation 28, $`D(P_1(𝒔),P_2(𝒔))`$ $`=`$ $`D(\widehat{P}(𝒔𝒙_{n+1},\mathrm{\Theta }_v^n,\varphi ,\theta ),\widehat{P}(𝒔\mathrm{\Theta }_{\overline{v}}^{n+1},\theta ))`$ $`=`$ $`{\displaystyle \widehat{P}(𝒔𝒙_{n+1},\mathrm{\Theta }_v^n,\varphi ,\theta )log\left(\frac{\widehat{P}(𝒔𝒙_{n+1},\mathrm{\Theta }_v^n,\varphi ,\theta )}{\widehat{P}(𝒔\mathrm{\Theta }_{\overline{v}}^{n+1},\theta )}\right)𝑑𝒔}`$ Note that it is not assumed here that $`𝒗`$ and $`\overline{𝒗}`$ have the same dimension. Expanding the probability distributions within the logarithm appearing above yields $`D(P_1(𝒔),P_2(𝒔))`$ $`=`$ $`{\displaystyle \widehat{P}(𝒔𝒙_{n+1},\mathrm{\Theta }_v^n,\varphi ,\theta )}`$ (47) $`\times [log\left(P(𝒉(𝒔,𝒗)\theta )\right)`$ $`+log\left(P(𝒉(𝒔,\overline{𝒗})\theta )\right)`$ $`+log\left(P(𝒙_{n+1}𝒔,\varphi )\right)`$ $`log\left(\widehat{P}(𝒙_{n+1}\mathrm{\Theta }_v^n,\varphi ,\theta )\right)`$ $`+log\left(\widehat{P}(𝒉(𝒔,𝒗)\mathrm{\Theta }_v^n)\right)`$ $`log\left(\widehat{P}(𝒉(𝒔,\overline{𝒗})\mathrm{\Theta }_{\overline{v}}^{n+1})\right)]d𝒔`$ Each term has the form of an information (or uncertainty). Together the six terms paint a descriptive picture of how information is acquired by the maximally informative update when taken as three groups of two terms: Denote by “new KR” the two terms with $`\overline{𝒗}`$ and $`\mathrm{\Theta }_{\overline{v}}^{n+1}`$, by “previous KR” the two terms with $`𝒗`$ and $`\mathrm{\Theta }_v^n`$ and no data, and by “new data” the two terms with data dependency. Now, noting the signs on these quantities, because $`D`$ is positive, the whole point of choosing a good $`\mathrm{\Theta }^{n+1}`$ approximation by minimizing $`D`$ is that $`\text{Expected information in new KR}`$ $`(\text{Expected information in previous KR}`$ $`+\text{Expected information in new data})`$ (48) or in very rough terms we may see the update as capturing the sum-total of the available knowledge $`\text{Total knowledge}=\text{Prior knowledge}+\text{New knowledge from data}`$ (49) Because only terms depending upon the update parameters $`\overline{𝒗}`$ and $`\mathrm{\Theta }_{\overline{v}}^{n+1}`$ are needed to perform the minimization, we drop the other terms at this point, and after making the multinormal substitutions for the distributions in the above we have $`\overline{D}(P_1(𝒔),P_2(𝒔))`$ $`=`$ $`{\displaystyle \widehat{P}(𝒉_{\overline{v}}𝒙_{n+1},\mathrm{\Theta }_v^n,\varphi ,\theta )log\left(N(𝝁_{\overline{v}},\mathrm{\Sigma }_{\overline{v}})(𝒉_{\overline{v}})\right)𝑑𝒉_{\overline{v}}}`$ $``$ $`{\displaystyle \widehat{P}(𝒉_{\overline{v}}𝒙_{n+1},\mathrm{\Theta }_v^n,\varphi ,\theta )log\left(N(𝝁_{\overline{v}}^{n+1},\mathrm{\Sigma }_{\overline{v}}^{n+1})(𝒉_{\overline{v}})\right)𝑑𝒉_{\overline{v}}}`$ To simplify the $`\widehat{P}`$’s appearing in equation LABEL:eq:two-d-terms, the distribution of surface given old knowledge and new data, marginalized to the height field $`\overline{𝒗}`$, is useful, as is seen by observing equations 47 and LABEL:eq:two-d-terms. Thus, consider $`\widehat{P}(𝒔𝒙_{n+1},\mathrm{\Theta }_v^n,\varphi ,\theta )`$ $``$ $`N(M(\varphi )𝒔,\mathrm{\Sigma }_ϵ^{n+1})(𝒙_{n+1})N(𝝁_s,\mathrm{\Sigma }_s)(𝒔)`$ (51) $`\times {\displaystyle \frac{N(𝝁_v^n,\mathrm{\Sigma }_v^n)(𝒉(𝒔,𝒗))}{N(𝝁_v,\mathrm{\Sigma }_v)(𝒉(𝒔,𝒗))}}`$ found by making substitutions into 6 for the assumed distributions. Since it is not necessarily the case that $`v_i\{\overline{v}_j\}`$ or that $`\overline{v}_i\{v_j\}`$. proceed by marginalizing to the union of the components of $`𝒗`$ and $`\overline{𝒗}`$, which we denote $`𝒗\overline{𝒗}`$, and then to the $`\overline{𝒗}`$ components. Let $`A_{v\overline{v},s}`$ denote the projection from $`𝒗_s`$ to $`𝒗\overline{𝒗}`$, $`A_{\overline{v},v\overline{v}}`$ denote the projection from $`𝒗\overline{𝒗}`$ to $`\overline{𝒗}`$, and $`A_{\overline{v},v}`$ denote the projection from $`𝒗`$ to $`\overline{𝒗}`$. In performing the two projections (from $`𝒗_s`$ to $`𝒗\overline{𝒗}`$, and then from $`𝒗\overline{𝒗}`$ to $`\overline{𝒗}`$) in order we find (not necessarily in most simple form), using results of appendices 12.212.5, that $`{\displaystyle \widehat{P}(𝒔𝒙_{n+1},\mathrm{\Theta }_v^n,\varphi ,\theta )𝑑𝒔}\overline{𝒗}=N(𝝁_R,\mathrm{\Sigma }_R)(𝒉_{\overline{𝒗}})`$ (52) where $`𝝁_{\overline{v}}^R=\mathrm{\Sigma }_R(\mathrm{\Sigma }_Q^1𝝁_{\overline{𝒗}}^Q+(\mathrm{\Sigma }_{\overline{v}}^n)^1𝝁_{\overline{𝒗}}^𝒏\mathrm{\Sigma }_{\overline{v}}^1𝝁_{\overline{𝒗}})`$ $`\mathrm{\Sigma }_R^1=\mathrm{\Sigma }_Q^1+(\mathrm{\Sigma }_{\overline{v}}^n)^1\mathrm{\Sigma }_{\overline{v}}^1`$ (53) and where $`𝝁_{\overline{v}}^Q=A_{\overline{v},v\overline{v}}A_{v\overline{v},s}𝝁_s^P`$ $`\mathrm{\Sigma }_Q^1=A_{\overline{v},v\overline{v}}A_{v\overline{v},s}\mathrm{\Sigma }_P^1A_{v\overline{v},s}^TA_{\overline{v},v\overline{v}}^T`$ (54) $`𝝁_s^P=\mathrm{\Sigma }_P(\mathrm{\Sigma }_s^1𝝁_s+M^T\mathrm{\Sigma }_ϵ^1𝒙_{n+1})`$ $`\mathrm{\Sigma }_P^1=\mathrm{\Sigma }_s^1+M^T\mathrm{\Sigma }_ϵ^1M`$ (55) $`𝝁_{\overline{𝒗}}^𝒏=A_{\overline{v},v}𝝁_𝒗^𝒏`$ $`(\mathrm{\Sigma }_{\overline{v}}^n)^1=A_{\overline{v},v}(\mathrm{\Sigma }_v^n)^1A_{\overline{v},v}^T`$ (56) $`𝝁_{\overline{𝒗}}=A_{\overline{v},v}𝝁_𝒗`$ $`\mathrm{\Sigma }_{\overline{v}}^1=A_{\overline{v},v}\mathrm{\Sigma }_v^1A_{\overline{v},v}^T`$ (57) $`𝝁_v=A_{v,s}𝝁_𝒔`$ $`\mathrm{\Sigma }_v^1=A_{v,s}\mathrm{\Sigma }_s^1A_{v,s}^T`$ Using the results of appendix 12.6, the quantities of equation 53 above correspond to the values of the mean and standard deviation parameters of the new KR, found at the minimum Kullback Leibler distance, i.e. the minimization is immediately apparent from those results. Thus: $`\mathrm{\Theta }_{\overline{v}}^{n+1}=(𝝁_{\overline{v}}^{n+1},\mathrm{\Sigma }_{\overline{v}}^{n+1})`$ $`𝝁_{\overline{v}}^{n+1}=𝝁_{\overline{v}}^R`$ $`\mathrm{\Sigma }_{\overline{v}}^{n+1}=\mathrm{\Sigma }_{\overline{v}}^R`$ (59) Equations 53 are the Generalized Kalman Filter (GKF) update equations for the surface inference example, yet are quite a bit more general (the necessary change of variables needed when the forward projection is nonlinear appears in appendix 12.10). Having these update equations allows one to consider updating a representation of any dimension relative to the original representation. Thus. knowledge may be represented in finer detail, corresponding to the old representation being contained in the new, knowledge may be represented in the same detail, corresponding to the case when the new representation is the same as the old representation, or knowledge may be tossed, corresponding to the case when the new representation does not contain the old representation. The maximally informative inference approach and its result of the Kullback Leibler distance on conditional posteriors led directly here to deriving the GKF and the solution of the problem of storing knowledge at scales adaptive to the actual needs of the data driving the update. The standard KF is discussed in . ## 7 Specializing the GKF When the surface of interest is itself a discrete height field, and the KR representation basis never changes in dimension nor position from that height field’s basis, then all projections appearing in equations 53 and following are identities, and the update equations simplify to the standard Kalman filter equations, in effect equations 6 only, given suitable identification of the variables. ## 8 Information learned Once a new set of parameters has been chosen, and for the purpose of evaluating the new update in the context of other possible updates at different scales, using different representational bases, it is useful to have the quantity of information about the surface distribution that is contained in the KR at the maximally informative update. Using the results of appendix 12.6 in equation LABEL:eq:two-d-terms we have this information, up to a constant, is given by $`I_R`$ $`=`$ $`C(𝒙_{n+1},\mathrm{\Theta }_v^n,\varphi ,\theta )`$ (60) $`+`$ $`{\displaystyle \frac{1}{2}}\left(Tr\left[(\mathrm{\Sigma }_R+U(𝝁_R𝝁_{\overline{v}}))\mathrm{\Sigma }_{\overline{v}}^1\right]+log(\left\|\mathrm{\Sigma }_{\overline{v}}\right\|)\right)`$ $``$ $`{\displaystyle \frac{1}{2}}\left(Tr\left[\mathrm{\Sigma }_R\mathrm{\Sigma }_R^1\right]+log(\left\|\mathrm{\Sigma }_R\right\|)\right)`$ Note that the $`d`$’s (representation basis dimensions) from the $`dlog(2\pi )`$’s of equation LABEL:eq:exp-info have cancelled. However the $`d`$’s remain hidden within the terms as matrix dimensions. When considering optimizing learned inormation against storage resources, one must weigh a separate cost in bits for the memory used against the bits learned, the expression above. Note also, interestingly the expression above contains a BIC-like $`log(d)`$ dependence term. ## 9 Search for update parameters Now that we know what the update equations for the updating of the KR distribution look like, it is worthwhile considering how an updating scheme might be implemented to acquire information at the appropriate scale. First, we dismiss the notion that we will ever be using the continuous height field $`𝒗_s`$ (the support of $`𝒔`$) at any time. None of the update equations force that to happen! Second, since we have concluded that computationally $`𝒗_s`$ is a discrete set, and since there will always be pathological cases where the surface is much rougher than we care to represent, we acknowledge that fact and proceed by presenting a useful algorithm which allows the updating of the KR while maintaining the ability to explore a large range of scales. The following multigrid-style algorithm provides the general flavor: * Choose $`𝒗_s`$ denser by several orders of scale than the current representation, and using other criteria associated with the knowledge of the data acquisition system (see below). * Choose $`\overline{𝒗}`$ at regular scales intermediate between $`𝒗_s`$ and the old KR on $`𝒗`$, compute the updates on all $`\overline{𝒗}`$ chosen at these scales. * Compute the information learned at each scale. * Plot the information learned as a function of increasing density (decreasing scale). * Choose, based on exploration of the plot, and costs associated with storing the learned information, whether to explore other octaves of scale. If Choose to explore, repeat above procedure. * If choice is to pick an informationally and storage attractive KR, do this and update the representation accordingly. In the surface reconstruction problem data often comes in the form of images. The images may come from devices with vastly different resolutions, and the known parameters of pixel size, point spread function and geometry determine the appropriate reconstruction scale. Finally adapting the surface to resolve at sub-pixel scales requires a memory-aggressive approach which extends the exploration farther out on the learning curve towards smaller, denser representation scales. ## 10 Conclusion Field inference has been generalized from the typical discrete fixed-basis setting to a continuous-basis setting. The problem of surface inference was solved in the context of continuous field inference. Using the approach of acquiring the maximally informative KR distribution, the GKF equations were found. The GKF allows the updated KR parameters to be found at any scale and/or “positions” (abstractly, basis components). The approach allows the learning of information at the relevant scales desired. It provides an information-theoretic justification for location-dependent adaptive multi-grid inference. It also effectively provides similar justification for a scale-adaptive MDL method. This is apparently the first time that the maximally informative inference of continuous-basis objects and the multigrid approach have been rigorously justified. ## 11 Acknowledgements I thank the members of the Ames Data Understanding group for their interest and comments, especially the invaluable valiant contributions of Dr. Robin D. Morris, who thoughtfully, carefully, and painstakingly spent a week-in-agony checking the maths (any remaining mistakes are fully mine, however), and Drs. Vadim Smelyanskiy and David Maluf for their comments. Finally, immense thanks go to Dr. Peter Cheeseman, for comments, and support. This project was partially supported by the NASA Ames Center for Excellence in Information Technology contract NAS-214217. ## 12 Appendices ### 12.1 Construction of a 2D surface prior In this appendix we first introduce the reader to the fourier representation of a gaussian process, then using the notions developed find the representation for a 2D gaussian process over the plane, where the correlations of the process at points $`𝒙`$ and $`𝒚`$ are proportional to $`exp(k\left\|𝒙𝒚\right\|)`$, $`k>0`$, a simple translation-invariant choice for the form of the correlation structure of the probability density of surfaces having the plane as support. The utility for the GKF of having this process is that it serves as a simply computed algorithmic representation of the prior for surfaces having the plane as support. #### 12.1.1 The discrete gaussian process Consider $`f(n,𝒄)`$, $`nZ_N=\{N,\mathrm{},1,0,1,\mathrm{},N\}`$, a discrete process with expression as the fourier expansion $$f(n,𝒄)=\underset{k=N}{\overset{N}{}}c_ke^{ikn}$$ (61) where the coefficients $`𝒄=(c_k)`$ are constrained by $`f𝑹`$ so that $`c_k=c_k^{}`$, and the $`n`$ and $`k`$ range over $`Z_N`$. Let the coefficients be random variables: $`c_k=x_k+iy_k`$ with $`x_kN(0,\sigma _k)`$ and $`y_kN(0,\sigma _k)`$ both gaussian distributed random variables with mean $`0`$ and standard deviation $`\sigma _k`$. Now, dropping the $`k`$’s, the joint density of $`(x,y)`$ is given by $$P_{x,y}(x,y)=\frac{e^{x^2/2\sigma ^2}}{\sqrt{2\pi }\sigma }\frac{e^{y^2/2\sigma ^2}}{\sqrt{2\pi }\sigma }.$$ (62) From this the joint density of $`(r,\theta )`$ where $`r=\sqrt{x^2+y^2}`$ and $`\theta =\mathrm{arctan}(y/x)`$ is given by $$P_{r,\theta }(r,\theta )=\frac{re^{r^2/2\sigma ^2}}{2\pi \sigma ^2}.$$ (63) The density of $`r`$ is given directly by integrating over $`\theta `$ $$P_r(r)=\frac{re^{r^2/2\sigma ^2}}{\sigma ^2},$$ (64) while the density of $`\theta `$ is given directly by integrating over $`r`$ $$P_\theta (\theta )=\frac{1}{2\pi }.$$ (65) Making a change of variables, the density of $`cc^{}=x^2+y^2=r^2`$ is given by the exponential distribution $$P_{cc^{}}(u)=\frac{e^{u/2\sigma ^2}}{2\sigma ^2}$$ (66) The distribution of $`c_k+c_k=2Re[c_k]=2x_k`$, $`k>0`$ is of interest because the process is real. $$P_{c+c^{}}(u)=\frac{e^{u^2/2(2\sigma )^2}}{\sqrt{2\pi }2\sigma }$$ (67) which is just a gaussian with zero mean but twice the variance of the components $`x`$ and $`y`$ of $`c`$. Note that the actual coefficients in equation 61 $`c_ke^{ikn}+c_ke^{ikn}=2Re[c_ke^{ikn}]`$ also have the distribution of equation 67 since the phase of $`c_k`$ is uniformly distributed in $`[0,2\pi ]`$. Now, given a set of integers $`\zeta Z_N`$ we may ask for the density of the sampled values of the process $`f`$ at $`𝜻=(n_1,n_2,\mathrm{},n_m)`$ $$𝒇(𝜻)=(f(n_1),f(n_2),\mathrm{},f(n_m)),$$ (68) where $`m=\left\|\zeta \right\|,n_iZ_N,i=1,\mathrm{},m`$. Define $$𝒇(𝜻,𝒄)=(f(n_1,𝒄),f(n_2,𝒄),\mathrm{},f(n_m,𝒄))$$ (69) Then the probability density function which describes the sampled values is $$P(𝒇(𝜻))=\delta (𝒇(𝜻)𝒇(𝜻,𝒄))P(𝒄)𝑑𝒄$$ (70) where $$P(𝒄)=P(c_0)\underset{k=1}{\overset{N}{}}P(c_k+c_k)$$ (71) Note that that the density of $`P(𝒇(𝜻))`$ is multivariate gaussian since the representation of $`𝒇(𝜻,𝒄)`$ as a fourier series shows that it is the sum of gaussian random vectors with components $`2Re[c_ke^{ikn}]`$. The covariances of the process are found as $`\mathrm{\Sigma }_{m,n}=E[f(m)f(n)]`$ $`=`$ $`E[f(m)f^{}(n)]`$ (72) $`=`$ $`E\left[{\displaystyle \underset{k,l=N}{\overset{N}{}}}c_kc_l^{}e^{i(kmln)}\right]`$ $`=`$ $`{\displaystyle \underset{k=N}{\overset{N}{}}}E[c_kc_k^{}]e^{ik(mn)}`$ $`=`$ $`F[E[c_kc_k^{}]](mn)`$ where we used the fact that the coefficients of different frequency are uncorrelated for $`kl`$, i.e $`E[c_kc_l^{}]=0`$ for $`kl`$. Define the power spectrum $`R(k)`$ as $$R(k)=E[c_kc_k^{}]$$ (73) Then we have that the covariance is given by the fourier transform of the power spectrum, $$\mathrm{\Sigma }_{m,n}=E[f(m)f(n)]=F[R](mn)=\mathrm{\Sigma }_{mn}$$ (74) where we have acknowledged that the covariance structure is dependent only upon the difference $`mn`$. From this we see that the inverse fourier transform of the covariance is the power spectrum, $$F^1\left[\mathrm{\Sigma }_u\right](k)=R(k)$$ (75) Finally, note that the density of $`c_kc_k^{}`$ given by equation 66 allows us to infer the parameters $`\sigma _k`$ which are the standard deviations of the gaussian processes $`x_k`$ and $`y_k`$ underlying the coefficients $`c_k`$, since from equation 66 $$E[c_kc_k^{}]=u\frac{e^{u/2\sigma _k^2}}{2\sigma _k^2}𝑑u=2\sigma _k^2$$ (76) In the next section the basis for gaussian processes developed here is extended to the continuous 2D case to compute the power spectrum of a process specified by a continuous-basis covariance structure. #### 12.1.2 The continuous-basis 2D process Similar to the development in the last section, in two dimensions, given the continuous-basis covariance $`\mathrm{\Sigma }_𝒙=exp(k\left\|𝒙\right\|)`$, $`k>0`$., the power spectrum is found as the inverse fourier transform of the covariance, i.e. $`R(𝒖=(u,v))`$ $`=`$ $`F_2^1[\mathrm{\Sigma }_𝒙](u,v)`$ (77) $`=`$ $`{\displaystyle e^{k\left\|(x,y)\right\|}e^{iux}e^{ivy}𝑑x𝑑y}`$ Make the change of variables $`(x,y)(r,\theta )`$ so that $`x=rcos(\theta )`$, $`y=rsin(\theta )`$, then $$R(u,v)=_0^{\mathrm{}}_0^{2\pi }e^{kr}e^{ir(ucos(\theta )+vsin(\theta ))}r𝑑r𝑑\theta $$ (78) For simplicity, make the further change of variables $`(u,v)(s,\varphi )`$ so that $`u=scos(\varphi )`$, $`v=ssin(\varphi )`$, so that $`R(s,\varphi )`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{2\pi }}e^{kr}e^{irs(cos(\varphi )cos(\theta )+sin(\varphi )sin(\theta ))}r𝑑r𝑑\theta `$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{2\pi }}e^{kr}e^{irscos(\theta \varphi )}r𝑑r𝑑\theta `$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}re^{kr}{\displaystyle _0^{2\pi }}e^{irscos(\theta \varphi )}𝑑\theta 𝑑r`$ $`R(s)`$ $`=`$ $`2\pi {\displaystyle _0^{\mathrm{}}}re^{kr}J_0(rs)𝑑r`$ (79) Finally, $`R(𝒖)`$ $`=`$ $`{\displaystyle \frac{2\pi k}{(\left\|𝒖\right\|^2+k^2)^{3/2}}}`$ (80) Note that we have neglected the proportionality constant $`1/2\pi `$ in the fourier transform, amounting to normalizing the delta function to $`2\pi `$, and have scaled $`𝒖`$ to units of cycles per $`2\pi `$. Note also that both the covariance of the process and the power spectrum scale with the same proportionality constant. Harmonic analysis is discussed in ### 12.2 Multinormal density MGF The moment generating function for a probability distribution $`f`$ is defined as the functional $$M[f](𝝀)=E_f[e^{Tr[U(𝝀,𝒙)]}]$$ (81) where $`U(𝒚,𝒛)`$ is defined such that $`U=[U_{ij}]`$ and $`U_{ij}(𝒚,𝒛):=y_iz_j`$, from which holds the property $$\frac{^kM[f](𝝀)}{\lambda _{i_1}\mathrm{}\lambda _{i_k}}_{𝝀=0}=E_f[x_{i_1}\mathrm{}x_{i_k}]$$ (82) i.e the moments are found as derivatives of the MGF with respect to the parameter $`𝝀`$ at $`𝝀=0`$. Take the multinormal density function for $`𝒙`$ $`P(𝒙\mathrm{\Theta })`$ $`=`$ $`N(\mathrm{\Theta })(𝒙)`$ (83) $`=`$ $`N(𝝁,\mathrm{\Sigma })(𝒙)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^{d/2}\mathrm{\Sigma }^{1/2}}}exp({\displaystyle \frac{1}{2}}Tr[U(𝒙𝝁)\mathrm{\Sigma }^1])`$ where $`U(𝒚)`$ is defined such that $`U_{ij}(𝒚):=U_{ij}(𝒚,𝒚)`$ and $`d=Dim(𝒙)`$. The MGF of $`N(\mathrm{\Theta })(𝒙)`$ is then given by $`M[N(\mathrm{\Theta })(𝒙)](𝝀)=E[e^{Tr[U(𝝀,𝒙)]}\mathrm{\Theta }]`$ $`={\displaystyle \frac{1}{(2\pi )^{d/2}\mathrm{\Sigma }^{1/2}}exp(\frac{1}{2}Tr[U(𝒙𝝁)\mathrm{\Sigma }^1]+Tr[U(𝝀,𝒙)])𝑑𝒙}`$ Minus twice the exponent of the integral above may be written as $`Tr[U(𝒙𝝁)\mathrm{\Sigma }^1]2Tr[U(𝝀,𝒙)]`$ $`=`$ $`Tr[U(𝒙(𝝁𝝀\mathrm{\Sigma }))\mathrm{\Sigma }^1]`$ (85) $`+Tr[U(𝝁)\mathrm{\Sigma }^1]`$ $`Tr[U(𝝁𝝀\mathrm{\Sigma })\mathrm{\Sigma }^1]`$ $`=`$ $`Tr[U(𝒙(𝝁𝝀\mathrm{\Sigma }))\mathrm{\Sigma }^1]`$ $`Tr[U(𝝀)\mathrm{\Sigma }]`$ $`2Tr[U(𝝀,𝝁)]`$ from which the moment generating function is immediately found as $$M[N(\mathrm{\Theta })(𝒙)](𝝀)=exp(Tr[U(𝝁,𝝀)]+\frac{1}{2}Tr[U(𝝀)\mathrm{\Sigma }])$$ (86) From the above we have $`E[x_i\mathrm{\Theta }]=\mu _i`$ $`E[(x_i\mu _i)(x_j\mu _j)\mathrm{\Theta }]=\mathrm{\Sigma }_{ij}`$ (87) which agrees with the calculation of appendix 12.2. Two things to note: 1. The inverse of $`\mathrm{\Sigma }`$ is assumed to exist. 2. All moments are determined by simple products and sums of the parameters $`(𝝁,\mathrm{\Sigma })`$. ### 12.3 Multinormal linear change of variables Letting $`𝒚=A𝒙`$ be the change of variables, where $`P(𝒙\mathrm{\Theta })=N(\mathrm{\Theta })(𝒙)`$, the MGF of the density $`P(𝒚\mathrm{\Theta })`$ is found from the MGF of the density for $`P(𝒙\mathrm{\Theta })`$ in a straightforward manner as $`M[P(𝒚\mathrm{\Theta })](𝝀)`$ $`=`$ $`E[e^{Tr[U(𝝀,𝒚)]}\mathrm{\Theta }]`$ $`=`$ $`E[e^{Tr[U(𝝀,A𝒙)]}\mathrm{\Theta }]`$ $`=`$ $`E[e^{Tr[U(A^T𝝀,𝒙)]}\mathrm{\Theta }]`$ $`=`$ $`exp(Tr[U(𝝁,A^T𝝀)]+{\displaystyle \frac{1}{2}}Tr[U(A^T𝝀)\mathrm{\Sigma }])`$ $`=`$ $`exp(Tr[U(A𝝁,𝝀)]+{\displaystyle \frac{1}{2}}Tr[U(𝝀)(A\mathrm{\Sigma }A^T)])`$ Note that the dropped subscripts $`x`$ and $`x`$ of the $`\mathrm{\Theta }`$ and $`\lambda `$ are easily determined by the context, and that the density used to take the expectation naturally changed in equation 12.3 from $`P(𝒚\mathrm{\Theta })`$ to $`P(𝒙\mathrm{\Theta })`$ without confusion. With this result and referring to equation 86 and preceding we find that the density for $`𝒚`$ is multinormal with $`𝝁_𝒚=A𝝁_𝒙`$ $`\mathrm{\Sigma }_y=A\mathrm{\Sigma }_xA^T`$ (90) Note that everywhere the condition of $`A`$ was neither mentioned nor assumed, thus $`A`$ may be a rectangular matrix or otherwise not of full rank. ### 12.4 Multinormal projections Another useful operation is that of projection onto a subset of the components of the argument of the multinormal distribution. Projections may be trivially represented as a linear operation, where the “projection matrix” is typically a rectangular matrix having the form of a unique (single) element of value $`1`$ in each row and column, zeroes elsewhere. Finding the distribution of the projected variables is equivalent to the operation of marginalizing over the components not in the projection. Let $`A`$ be the projection matrix selecting a subset of the variables of $`𝒙`$ as $`𝒚=A𝒙`$. Then, using the result of section 12.3, we immediately find integrals of the form $$N(𝝁,\mathrm{\Sigma })(𝒙)𝑑𝒙𝒚=N(A𝝁,A\mathrm{\Sigma }A^T)(𝒚)$$ (91) Both vector $`A𝝁`$ and the matrix $`A\mathrm{\Sigma }A^T`$ are now just appropriately rearranged pieces of the original vector $`𝝁`$ and matrix $`\mathrm{\Sigma }`$. Specifically, if $`y_k=x_{i_k}`$ then $`[A\mathrm{\Sigma }A^T]_{pq}=\mathrm{\Sigma }_{i_pj_q}`$. ### 12.5 Multinormal multiplication One operation which frequently occurs in Bayesian inference is that of taking the product of two multinormal distributions of the same variable and normalizing that product to find a new distribution. Finding the new $`\mathrm{\Theta }=(𝝁,\mathrm{\Sigma })`$ amounts to completing the square, but it is useful to state the result, and we do this here. Let $`\mathrm{\Theta }_1=(𝝁_1,\mathrm{\Sigma }_1)`$ and $`\mathrm{\Theta }_1=(𝝁_1,\mathrm{\Sigma }_1)`$ be the parameters of the multinormal distributions in the product. Then $`𝝁=\mathrm{\Sigma }(\mathrm{\Sigma }_1^1𝝁_\mathrm{𝟏}+\mathrm{\Sigma }_2^1𝝁_\mathrm{𝟐})`$ $`\mathrm{\Sigma }=(\mathrm{\Sigma }_1^1+\mathrm{\Sigma }_1^1)^1`$ (92) ### 12.6 Expected uncertainty in multinormals It is useful to know the expected uncertainty of one gaussian distribution in the context of another. Consider the quantity $`E[log(P(\mathrm{\Theta }_2)(𝒙))\mathrm{\Theta }_1]={\displaystyle N(𝝁_1,\mathrm{\Sigma }_1)(𝒙)log\left(N(𝝁_2,\mathrm{\Sigma }_2)(𝒙)\right)𝑑𝒙}`$ (93) which occurs in similar form in the development of the Generalized Kalman Filter (section 6) and represents the expected uncertainty, or entropy, of the surface representation in the context of the updated surface distribution. The value of this integral is found straightforwardly using the results mentioned in appendix 12.2 as $`E[log(N(𝝁_2,\mathrm{\Sigma }_2)(𝒙))\mathrm{\Theta }_1]`$ $`=`$ $`{\displaystyle \frac{1}{2}}E\left[Tr[U(𝒙𝝁_2)\mathrm{\Sigma }_2^1]\right]`$ $`+{\displaystyle \frac{d}{2}}log(2\pi )+{\displaystyle \frac{1}{2}}log(\left\|\mathrm{\Sigma }_2\right\|)`$ $`=`$ $`{\displaystyle \frac{1}{2}}Tr\left[(\mathrm{\Sigma }_1+U(𝝁_1𝝁_2))\mathrm{\Sigma }_2^1\right]`$ $`+{\displaystyle \frac{d}{2}}log(2\pi )+{\displaystyle \frac{1}{2}}log(\left\|\mathrm{\Sigma }_2\right\|)`$ ### 12.7 Maximizing the expected information Varying $`\mathrm{\Sigma }_2`$, the minimum value of the uncertainty above occurs when $`\mathrm{\Theta }_2=\mathrm{\Theta }_1`$. That this is true for the $`𝝁`$ component of $`\mathrm{\Theta }_2`$ is immediate from the positive definite quadratic nature of the first term. For the $`\mathrm{\Sigma }`$ component the following fact following from the properties of determinants and matrix inverses facilitates the result: $$\frac{\left\|\mathrm{\Sigma }\right\|}{\mathrm{\Sigma }_{kl}}=(1)^{k+l}\frac{Cof_{kl}(\mathrm{\Sigma })}{\left\|\mathrm{\Sigma }\right\|}=\mathrm{\Sigma }_{kl}^1$$ (95) ### 12.8 Notes on matrix inverses and submatrices Given the invertible matrix $`V`$, composed in the following manner of submatrices $`V_{11}`$, $`V_{12}`$, $`V_{21}`$, $`V_{22}`$, $$A=\left[\begin{array}{cc}V_{11}& V_{12}\\ V_{21}& V_{22}\end{array}\right]$$ (96) and its inverse $$A^1=\left[\begin{array}{cc}\widehat{V}_{11}& \widehat{V}_{12}\\ \widehat{V}_{21}& \widehat{V}_{22}\end{array}\right]$$ (97) then it is immediate that the following relationships hold among the submatrices $`\left[\begin{array}{cc}I_{11}& N_{12}\\ N_{21}& I_{22}\end{array}\right]`$ $`=`$ $`\left[\begin{array}{cc}V_{11}\widehat{V}_{11}+V_{12}\widehat{V}_{21}& V_{11}\widehat{V}_{12}+V_{12}\widehat{V}_{22}\\ V_{21}\widehat{V}_{11}+V_{22}\widehat{V}_{21}& V_{21}\widehat{V}_{12}+V_{22}\widehat{V}_{22}\end{array}\right]`$ (102) where $`I`$ and $`N`$ represent the identity and zero matrices respectively. Any quadratic operator $`𝒙^TQ𝒙`$ may be decomposed using projection matrices $`A`$ and $`\overline{A}`$ where these are diagonal matrices with one and zero entries only, and where $$A+\overline{A}=I$$ (103) in the following manner $`𝒙^TQ𝒙`$ $`=`$ $`𝒙^T(A+\overline{A})Q(A+\overline{A})^T𝒙`$ (104) $`=`$ $`𝒙_A^TQ_{AA}𝒙_A+𝒙_A^TQ_{A\overline{A}}𝒙_{\overline{A}}+𝒙_{\overline{A}}^TQ_{\overline{A}A}𝒙_A+𝒙_{\overline{A}}^TQ_{\overline{A}\overline{A}}𝒙_{\overline{A}}`$ Now, assume $`Q`$ is symmetric and that both it and $`Q_{AA}`$ and $`Q_{\overline{A}\overline{A}}`$ are invertible, and rewrite this form as the sum of two terms as follows $`𝒙^TQ𝒙`$ $`=`$ $`(𝒙_A𝜶)^TQ_{AA}(𝒙_A𝜶)+C(𝒙_{\overline{A}})`$ $`=`$ $`𝒙_A^TQ_{AA}𝒙_A𝒙_A^TQ_{A\overline{A}}𝒙_{\overline{A}}𝒙_{\overline{A}}^TQ_{\overline{A}A}𝒙_A+𝜶^TQ_{AA}𝜶+C(𝒙_{\overline{A}})`$ where $`𝜶=(Q_{AA})^1Q_{A\overline{A}}𝒙_{\overline{A}}`$. Thus $$C(𝒙_{\overline{A}})=𝒙_{\overline{A}}^T\left(Q_{\overline{A}\overline{A}}Q_{\overline{A}A}(Q_{AA})^1Q_{A\overline{A}}\right)𝒙_{\overline{A}}$$ (106) Applying the identities of equation 102 $$Q_{AA}\widehat{Q}_{A\overline{A}}+Q_{A\overline{A}}\widehat{Q}_{\overline{A}\overline{A}}=N_{A\overline{A}}$$ followed by $$Q_{\overline{A}A}\widehat{Q}_{A\overline{A}}+Q_{\overline{A}\overline{A}}\widehat{Q}_{\overline{A}\overline{A}}=I_{\overline{A}\overline{A}}$$ find that $$Q_{\overline{A}\overline{A}}Q_{\overline{A}A}(Q_{AA})^1Q_{A\overline{A}}=(\widehat{Q}_{\overline{A}\overline{A}})^1$$ (107) so that $$C(𝒙_{\overline{A}})=𝒙_{\overline{A}}^T(\widehat{Q}_{\overline{A}\overline{A}})^1𝒙_{\overline{A}}$$ (108) which immediately provides an alternate method for marginalizing gaussian distributions. ### 12.9 Alternate inverse forms In the GKF update equations expressions for updating inverse matrices in terms of the sum of other inverse matrices occur. Because one of the summand matrices may not be well-conditioned, it is of interest to find an expression for the updated matrix in terms of the other matrices, which explicitly is not a function of the inverse matrices. Thus, let $`P`$, $`Q`$, $`R`$ be invertible matrices such that $`P^1=Q^1+R^1`$ (109) Then we find $`P`$ $`=`$ $`QQ(Q+R)^1Q`$ (110) by the following direct substitution $`PP^1`$ $`=`$ $`(QQ(Q+R)^1Q)(Q^1+R^1)`$ (111) $`=`$ $`IQ\left[(Q+R)^1(I+QR^1)R^1\right]`$ $`=`$ $`I`$ ### 12.10 Nonlinear forward projection In the nonlinear forward projection case the projection is given by $`𝒇(𝒔)`$, where $`𝒇()`$ is a nonlinear function of $`𝒔`$ rather then the linear form $`M𝒔`$. Because the derivative of the forward projection is often a straightforward object to compute, expand $`𝒇(𝒔)`$ about the mean of the old surface, $`𝝁_s`$ $$𝒙=𝒇(𝝁_s)+\frac{𝒇}{𝒔}_{𝝁_s}(𝒔𝝁_s)+ϵ$$ (112) Letting $`M=\frac{𝒇}{𝒔}_{𝝁_s}`$ we have $`P(𝒙𝒔,\varphi )`$ $`=`$ $`N((𝒇(𝝁_s)M𝝁_s)+M𝒔,\mathrm{\Sigma }_ϵ)(𝒙)`$ $`=`$ $`N(M𝒔,\mathrm{\Sigma }_ϵ)(𝒙\mathbf{}\mathbf{(}𝒇\mathbf{(}𝝁_𝒔\mathbf{)}\mathbf{}𝑴𝝁_𝒔\mathbf{)})`$ so that the appropriate changes to be made to the GKF update equations are simply $`𝒙𝒙\mathbf{}\mathbf{(}𝒇\mathbf{(}𝝁_𝒔\mathbf{)}\mathbf{}𝑴𝝁_𝒔\mathbf{)}`$ $`M{\displaystyle \frac{𝒇}{𝒔}}_{𝝁_s}`$ (114) while everything else otherwise remains the same.
warning/0005/hep-th0005163.html	ar5iv	text	# 1 Introduction ## 1 Introduction In the context of the AdS/CFT correspondence , Polchinski and Strassler used brane polarization , to find a string theory dual to a confining gauge theory in four dimensions. They perturbed the $`𝒩=4`$ SYM theory living on N D3 branes by adding fermion masses and found a non singular dual spacetime which contained dielectric branes. Giving masses to three out of the four Weyl fermions broke supersymmetry to $`𝒩=1`$, and the resulting gauge theory had a rich structure of vacua, domain walls, flux tubes, and instantons. Motivated by their approach we perturb the 2+1 dimensional SYM theory living on a set of coincident D2 branes. We add fermion mass terms and examine its supergravity/string theory dual. Unlike the theory studied in , this theory is not conformal, and there are different dual descriptions of the unperturbed theory depending on the regime we are in . In particular, very close to the branes the dilaton blows up, and IIA string theory ceases to provide a good description; consequently, we find ourselves in M-theory. Polarization of Dp branes into D(p+2) branes can be understood in two ways. First, if we turn on a field which couples to a D(p+2) brane, a configuration with the Dp branes polarized into a wrapped D(p+2) brane is energetically favored. Alternatively, examining the superpotential in the weakly coupled theory, shows that the fields describing the coordinates of the Dp branes become noncommutative. The resulting configuration couples with a $`p+3`$ form (via a Born-Infeld term discovered in ) and has a nonzero transverse size; therefore it can be interpreted as wrapped D(p+2) brane. For other types of polarizations (D3 branes into NS5 branes or M2 branes into M5 branes the matrix model description is not well understood. The only direct description is through energy considerations. Nevertheless, S-duality or the analysis of the symmetries of the M2 brane give indirect evidence for the second picture. In , adding fermion mass perturbations results in the D3 branes being polarized into D5 or NS5 branes. If three Weyl fermions have equal masses, the $`R`$ symmetry is broken to $`SO(3)`$, and the five-branes have an $`R^4\times S^2`$ geometry. If the three fermion masses are different, the shape of the five-branes degenerates from a two-sphere into a two-ellipsoid. In , turning on four equal masses for the four complexified fermions polarizes the M2 branes into M5 branes of geometry $`R^3\times S^3`$. If the four masses are not equal, the $`SO(4)`$ invariant configuration of the M5 branes degenerates into an ellipsoid. By perturbing the theory on D2 branes with fermion mass terms, we find configuration corresponding to D2 branes polarized into D4 or NS5 branes. The first polarizations is very similar to the D3 $``$ D5 polarization found in and can be obtained from it via T-duality. The second polarization is reminiscent of the M2 $``$ M5 polarization from and can be seen as coming from it under the compactification of the eleventh dimension. The D4/D2 geometry is smooth everywhere. For the NS5 brane the dilaton blows up in the throat region. However, the rate of growth of the NS5 brane dilaton in the throat region is much lower than that for the D2 brane. Therefore the NS5/D2 geometry is less singular than that of D2 branes alone. We may also study polarizations into NS5 branes with D4 brane charge, which should be similar to the oblique polarizations of , but they are more complicated and will be the object of further study . This paper is organized as follows: In chapter 2 we perturb the IIA supergravity solution created by a large number $`N`$ of D2 branes, with the field corresponding to fermion masses in the three dimensional Yang -Mills. In chapter 3 we examine polarizations into D4 branes. We first consider a test D4 brane with large D2 brane charge in the perturbed geometry, and find that it gets polarized into $`S^2`$. We then extend the analysis to the full warped geometry created by shells of polarized D2 branes and complete the analysis of the D2-D4 system by finding its near shell solution. In chapter 4 we study the polarization of D2 branes into NS5 branes. First we examine an SO(4) invariant fermion mass term, which does not preserve any supersymmetry and polarizes the D2 branes into NS5 branes wrapped on a 3-sphere. Then we take one fermion mass to 0, which restores back $`𝒩=1`$ and squashes the sphere into an SO(3) invariant ellipsoid. We extend the analysis to the full warped geometry. Finally, in chapter 5 we find string theory duals to domain walls, baryon vertices, flux tubes and glueballs. ## 2 Perturbation of the bulk According to the extension of the AdS/CFT duality to non conformal cases , the 2+1 SYM living on a set of $`N`$ D2 branes in a given regime is dual to type IIA string theory living in the near horizon geometry of these branes. The unperturbed space-time is given by: $`ds^2`$ $`=`$ $`Z^{1/2}\eta _{\mu \nu }dx^\mu dx^\nu +Z^{1/2}dx^mdx^m,`$ $`e^\varphi `$ $`=`$ $`g_sZ^{1/4},`$ $`C_3^0`$ $`=`$ $`{\displaystyle \frac{1}{g_sZ}}dx^0dx^1dx^2,F_4^0=dC_3^0,`$ (1) where $`\mu ,\nu =0,1,2`$, $`i,j=3,\mathrm{},9`$, $`g_s`$ is the string coupling and the metric is given in the string frame. When the D2 branes are coincident $$Z=\frac{R^5}{r^5},R^5=6\pi ^2Ng_s\alpha _{}^{}{}_{}{}^{5/2}.$$ (2) The supergravity description is valid when the curvature and the dilaton are small. The Yang Mills coupling constant is related to the string coupling by $`g_{YM}^2=g_s/\sqrt{\alpha ^{}}`$. We are interested in the boundary theory in the strongly coupled regime, which corresponds to small bulk curvature. As we go close to the D2 branes, the dilaton grows (the eleventh dimension opens up), and we have to use M-theory to describe the physics. The fermions on the two brane transform in the $`\mathrm{𝟖}`$ of the $`SO(7)`$ R symmetry group . The fermion bilinear and the mass matrix transform in a symmetric traceless representation 35. This corresponds to antisymmetric three-tensors on the seven dimensional transverse space. Thus, we expect a fermion mass to correspond to a nonzero background for a three form field strength (which can only be NS-NS), and by Poincareé duality to a four form field strength (which is R-R). The relevant IIA equations of motion and Bianchi identities are: $`2d(e^{2\varphi }H_3)=F_4F_4,`$ $`d(F_4+B_2F_4)=0,`$ $`dF_4=0,`$ $`dH_3=0.`$ (3) The background $`F_4^0`$ given by (1) and the first order perturbation $`F_4^1`$, which has components only in the transverse directions are two 4-form field strengths. Since both $`F_4^1`$ and $`H_3`$ have only transverse components, it is useful to express the ten dimensional Hodge $``$ in terms of a 7-dimensional one obtained with a flat metric: $`F_4={\displaystyle \frac{1}{Z}}(_7F_4)dx^0dx^1dx^2,`$ $`H_3={\displaystyle \frac{1}{Z^{1/2}}}(_7H_3)dx^0dx^1dx^2.`$ (4) By combining (3) with (1) we obtain $`d[1/Z(_7H_3+g_sF_4^1)]=0,`$ $`d[1/Z(_7g_sF_4^1+H_3)]=0.`$ (5) We can see that $`F_4^1`$ and $`H_3`$ are tensor spherical harmonics on the transverse space. We can express them in a basis for 3 and 4 tensors respectively, given in the Appendix: $`H_3=r^p(c_1T_3+c_2V_3),F_4^1=r^q(c_3T_4+c_4V_4).`$ (6) Equation (5) suggests the use of an ansatz: $$_7T_3=T_4,$$ (7) which implies $$_7V_4=T_3V_3,$$ (8) and its Hodge dual. Combining these with the relations in the Appendix, we obtain two solutions. The first one (which will interest us here) has a non normalizable mode $`H_3=g_s\alpha /r^5(3T_35V_3),`$ $`F_4^1=\alpha /r^5(4T_45V_4),`$ (9) and a normalizable mode: $`H_3=g_s\overline{\alpha }/r^7(3T_37V_3),`$ $`F_4^1=\overline{\alpha }/r^7(4T_47V_4).`$ (10) The other solution has the modes $`H_3=4g_s\alpha /T_3,`$ $`F_4^1=4\alpha /T_4,`$ (11) and $`H_3=4g_s\alpha /r^{12}(T_3+4V_3),`$ $`F_4^1=4\alpha /r^{12}(T_43V_4).`$ (12) These modes come from the M-theory anti self-dual and self-dual four form field strengths after the eleventh dimension is compactified. We now explore the relation between the $`H_3`$ and $`F_4^1`$, which we turned on, and the $`𝒩=8`$ gauge theory on the $`N`$ D2 branes. This theory has eight Majorana fermions transforming in the $`\mathrm{𝟖}`$ of the $`SO(7)`$ R symmetry group. The theory also has seven scalars in the fundamental representation of SO(7) and a vector. In $`𝒩=2`$ language, six of these scalars can be combined with six fermions into three hypermultiplets, and one scalar can be combined with the remaining two fermions and the vector into a vector multiplet. We see that giving masses to the three hypermultiplets (six fermions) preserves $`𝒩=2`$ supersymmetry. Nevertheless, if we give mass to the other two fermions we break supersymmetry completely. In order to relate fermion masses to tensors on the transverse seven dimensional space it is convenient to group the eight fermions into four complex pairs, and to group six of the seven transverse directions into three complex pairs : $`\mathrm{\Lambda }^1`$ $`=`$ $`\lambda ^1+i\lambda ^2,\mathrm{\Lambda }^2=\lambda ^3+i\lambda ^4,\mathrm{\Lambda }^3=\lambda ^5+i\lambda ^6,\mathrm{\Lambda }^4=\lambda ^7+i\lambda ^8,`$ $`z^1`$ $`=`$ $`x^3+ix^7,z^2=x^4+ix^8,z^3=x^5+ix^9.`$ (13) In M-theory $`x^6`$ can be also grouped with the eleventh dimension into a complex pair, but this is not possible here. Under a rotation $`z^ie^{i\varphi _i}z^i`$ the fermions transform as : $`\mathrm{\Lambda }^1`$ $``$ $`e^{i(\varphi _1+\varphi _2+\varphi _3)/2}\mathrm{\Lambda }^1,`$ $`\mathrm{\Lambda }^2`$ $``$ $`e^{i(\varphi _1\varphi _2\varphi _3)/2}\mathrm{\Lambda }^2,`$ $`\mathrm{\Lambda }^3`$ $``$ $`e^{i(\varphi _1+\varphi _2\varphi _3)/2}\mathrm{\Lambda }^3,`$ $`\mathrm{\Lambda }^4`$ $``$ $`e^{i(\varphi _1+\varphi _2+\varphi _3)/2}\mathrm{\Lambda }^4.`$ (14) If we perturb the Lagrangian with the traceless symmetric (in $`\lambda _i`$) combination $$\mathrm{\Delta }L=\mathrm{Re}(m_1\mathrm{\Lambda }_1^2+m_2\mathrm{\Lambda }_2^2+m_3\mathrm{\Lambda }_3^2+m_4\mathrm{\Lambda }_4^2),$$ (15) the corresponding bulk spherical harmonics transforming in the same way under SO(7) will be: $`T_4`$ $`=`$ $`\mathrm{Re}(m_1d\overline{z}^1dz^2dz^3dx^6+m_2dz^1d\overline{z}^2dz^3dx^6`$ $`+`$ $`m_3dz^1dz^2d\overline{z}^3dx^6+m_4dz^1dz^2dz^3dx^6),`$ $`T_3`$ $`=`$ $`\mathrm{Im}(m_1d\overline{z}^1dz^2dz^3+m_2dz^1d\overline{z}^2dz^3`$ (16) $`+`$ $`m_3dz^1dz^2d\overline{z}^3+m_4dz^1dz^2dz^3).`$ We chose the signs so that $`T_4=T_3`$. The coefficient $`\alpha `$ from (9) relates the fermion mass from the gauge theory to the strength of the nonnormalizable mode, and will be determined in the next chapter. Throughout this paper we are keeping three masses equal: $`m_1=m_2=m_3=m`$. If $`m_4=0`$, we have $`𝒩=2`$ supersymmetry, and SO(3) spherical symmetry. We can use this perturbation to study polarization of D2 branes into D4 branes. If $`m_4=m`$ we have SO(4) symmetry and no supersymmetry. We can study polarization of D2 branes into $`S^3`$ wrapped NS5 branes. We can also study D2/NS5 configurations for general $`m_4`$. In the limit $`m_40`$ we recover supersymmetry. ## 3 Polarization into D4 branes wrapped on $`S^2`$. As we discussed in the introduction, there are two alternative ways to study polarization of D2 branes into D4 branes: using a matrix model description, or examining the potential for D4 branes with large D2 brane charge. In order to find the self-interacting potential for $`S^2`$ wrapped D4 branes with large D2 brane charge, we first solve a simpler problem, and then, we use it to find the potential for a general configuration. In the first subsection of this chapter we explore a probe D4 brane with D2 brane charge $`n`$, in the background created by a large number $`Nn`$ of coincident D2 branes (1) perturbed with (9). We find a supersymmetric ground state at a nonzero radius. Then, we generalize to $`N`$ D2 branes distributed uniformly over a 2-sphere. We use this later generalization to find the full self-interacting potential of a spherical D4 brane with large D2 brane charge. ### 3.1 D4 brane probes In a background formed by a large number of Dp branes the action for a probe D(p+2) brane of world-volume $`R^{p+1}\times S^2`$ carrying Dp brane charge, is $`S`$ $`=`$ $`\mu _{p+2}/g_s{\displaystyle d^{p+3}\xi Z^{(p3)/4}\sqrt{detG_{}det(G_{}+_2)}}`$ (17) $``$ $`\mu _{p+2}{\displaystyle (C_{p+3}+2\pi \alpha ^{}_2C_{p+1})},`$ where $`\mu ,\nu =0\mathrm{}p`$ and $`m,n=p+1\mathrm{}9`$. $`G_{}`$ is a metric on a two sphere with $`detG_{}=Zr^4\mathrm{sin}\theta ^2`$ and $`G_{}`$ is the metric in the $`R^{p+1}`$ directions with $`detG_{}=Z^{(p+1)/2}`$. We have also defined $$2\pi \alpha ^{}_22\pi \alpha ^{}F_2B_2.$$ (18) The Dp brane charge of the D(p+2) brane is obtained by turning on a 2-form field strength along the $`S^2`$: $`F_{\theta \varphi }=(n/2)\mathrm{sin}\theta `$ so that $$_{S^2}F_2=2\pi n,$$ (19) and $$F_{ab}F^{ab}=\frac{n^2}{2Zr^4}.$$ (20) This result does not depend on the particular form of $`Z`$. Specializing to the case of interest, $`p=2`$, we can expand the integrand of the D4 action for large D2 brane charge: $$\sqrt{det(G_{}+2\pi \alpha ^{}F)}2\pi \alpha ^{}\sqrt{detF}+\frac{detG_{}}{4\pi \alpha ^{}\sqrt{detF}}.$$ (21) Using (2), we see that the expansion is valid, provided $$4\pi ^2F_{ab}F^{ab}n^2r/Ng\sqrt{\alpha ^{}}1.$$ (22) The leading term of the Born-Infeld action corresponds to the D2 brane charge, and is canceled by the second term in the Wess-Zumino action, as expected. The sub-leading term of the Born-Infeld action gives the potential per unit longitudinal volume $$\frac{S_{BI}}{V}=\mu _4/g_s_{S^2}d^2\xi Z^{1/4}\frac{\sqrt{detG_{}}detG_{}}{4\pi \alpha ^{}\sqrt{detF}}=\frac{\mu _42r^4}{g_sn\alpha ^{}},$$ (23) where factors of $`Z`$ have canceled making the answer identical to the $`D3/D5`$ case. The integral of $`C_5`$ gives a sub-leading term in the Wess-Zumino action. Since we know $`F_4`$ we can find $`C_5`$ by Poincare duality: $$F_4=dC_5+H_3C_3.$$ (24) The relative sign of $`F_4`$ and $`H_3C_3`$ is found from the equation of motion (3). The relative sign of $`C_5`$ and $`H_3C_3`$ is determined by gauge invariance under the transformations $`\delta C_3=d\chi _2,`$ $`\delta C_5=H_3.\chi _2`$ (25) The component of $`C_5`$ which couples with the wrapped D4 brane, is the Poincare dual of $`F_4^1`$. Using (4) and (9) we find that up to a gauge choice $$C_5=\frac{2\alpha }{3R^5}dx^0dx^1dx^2S_2.$$ (26) Calling $`\zeta \alpha g_s/R^5`$, we find the contribution of the Wess-Zumino term to be $$\frac{S_{WZ}}{V}=\frac{2\mu _4\zeta }{3g_s}_{S^2}S_2.$$ (27) By suitably choosing the plane of $`S_2`$ we can have a minimum of the effective potential away from $`r=0`$. The minimum of (27) is when $`S^2`$ is in the 789 plane. A general SO(3) invariant brane configuration is obtained by rotating the sphere of radius $`r`$ in the 3-7, 4-8 and 5-9 planes by the same angle $`\theta `$, and can be parametrized by $`z=re^{i\theta }`$. Therefore, $$_{S^2}S_2=4\pi \left(3m\mathrm{Im}(zz\overline{z})+m_4Im(zzz)\right)$$ (28) Using (27) we find the two main terms in the potential to be: $$\frac{S_{WZ}+S_{BI}}{V}=\frac{2\mu _4\|z\|^4}{g_sn\alpha ^{}}\frac{8\pi \mu _4m\zeta \mathrm{Im}zz\overline{z}}{g_s}.$$ (29) Since we have supersymmetry in our theory, we expect the potential to be the square of the derivative of a superpotential, so we expect to have another term in (29) proportional to $`m^2n`$. This term comes from second order perturbation of the background metric, dilaton, and $`C_{012}`$ (via the second equation of (3)). Since $`C_{012}`$ couples to $`F_2`$, this term is proportional to $`m^2n/g_s`$. Fortunately, supersymmetry allows us to read off this term by completing the square in (29). Using the convention in we express the D2 brane collective coordinate $`z`$ in terms of the gauge theory scalar $`\varphi =z/2\sqrt{2}\pi \alpha ^{}`$ and obtain: $$\frac{S}{V}=\frac{2\mu _4\|z\|^4}{g_sn\alpha ^{}}\frac{8\pi \zeta \mu _4m\mathrm{Im}zz\overline{z}}{g_s}+\frac{8\pi ^2\zeta ^2n\mu _4m^2\|z\|^2\alpha ^{}}{g_s}=\frac{8}{g_{ym}^2n}\|\varphi ^2imn\varphi \zeta /\sqrt{2}\|^2.$$ (30) We can compare this with the classical 3-dimensional $`𝒩=2`$ Kahler potential and superpotential (which is identical to the one in \- equations (73) and (74)): $`K={\displaystyle \frac{n}{g_{ym}^2}}\overline{\mathrm{\Phi }}\mathrm{\Phi },W={\displaystyle \frac{mn}{2g_{ym}^2}}\mathrm{\Phi }^2+{\displaystyle \frac{2\sqrt{2}i}{3g_{ym}^2}}\mathrm{\Phi }^3.`$ (31) They agree for $`\zeta =1/2`$, which fixes the normalization of (9). There is a supersymmetric minimum at $$z=i\pi mn\alpha ^{}.$$ (32) Thus, a probe D4 brane with D2 charge $`n`$ has a supersymmetric ground state in the 789 plane at radius $`\|z\|`$. As one may notice, the potential (30) does not depend on $`x_6`$. Therefore, the D4 sphere can be continuously moved along $`x_6`$. This modulus can in principle be lifted by quantum corrections, and deserves a more complete investigation. ### 3.2 The full problem In this section we find the full potential for a D4 brane with D2 brane charge $`N`$ in the geometry created by itself. Since this is a self interaction problem, we have to find the potential by bringing D4 shells with D2 charge from infinity in the background created by polarized D2 branes. The only difference from the previous chapter is the new $`Z`$. Calling $`r_0`$ the radius of the 2-sphere in which the branes are polarized, $`w`$ the radius in the polarization plane and $`y`$ the radius in the transverse directions we obtain: $$Z=\frac{R^5}{6r_0w}\left(\frac{1}{(y^2+(wr_0)^2)^{3/2}}\frac{1}{(y^2+(w+r_0)^2)^{3/2}}\right).$$ (33) Far away from the branes this reduces to (2). One might expect the calculation to be much harder, but using a clever trick from we do not have to compute anything. Equations (5) imply that $$d[Z^1(H_3+g_s_7F_4)]=d_7[Z^1(H_3+g_s_7F_4)]=0;$$ (34) therefore, $`Z^1(H_3+g_s_7F_4)`$ is constant and equal to its value at $`\mathrm{}`$ (which is given by the gauge theory). In particular, (24) implies that $`C_5`$ remains unchanged; thus, the Wess- Zumino term is the same. Moreover, we have seen that the Born-Infeld term (23) does not depend on $`Z`$ either. The third term of the potential is related to the first two by supersymmetry, and therefore it is also invariant. Thus, there is no change in the potential of a probe brane if the D2 branes creating the geometry are spread out. Because of the above, the total potential is the same as the probe potential. A general ground state is given by D2 branes polarized into several shells of D4 branes, with D2 brane charges $`n_i`$. As we have discussed, the potential felt by a D4 shell depends only on its radius, and not on $`x^6`$ or on the position of the other shells. Thus, the shells may be located at different positions along $`x^6`$. This moduli space can in principle be destroyed by quantum corrections, and we can also destroy it classically by giving a mass to $`x^6`$. The action is a sum of the actions (30), each with its own $`n_i`$. The results in this section are similar to the ones obtained for the D3-D5 polarization in . Indeed, by performing a T -duality along $`x^6`$ we can obtain all the D3-D5 configurations from D2-D4 configurations. ### 3.3 The near shell solution In finding the equilibrium configuration for the D2-D4 system we assumed that $`F_2`$ is dominant in the Born-Infeld term. This approximation breaks down close to the D4 brane shell, where we need to account for the deformation of the metric due to the D4 charge. Since very close to the shell we can approximate the shell to be flat, we can match the metric (1,33) with that of $`p`$ flat D4 branes with magnetic flux . This metric is: $`ds_{string}^2`$ $`=`$ $`{\displaystyle \frac{u^{3/2}}{pa^{3/2}\pi g_s\sqrt{\alpha ^{}}}}\left[\eta _{\mu \nu }dx^\mu dx^\nu +h(d\stackrel{~}{x^3}d\stackrel{~}{x^3}+d\stackrel{~}{x^4}d\stackrel{~}{x^4})\right]+`$ $`+`$ $`{\displaystyle \frac{\sqrt{\alpha ^{}}pa^{3/2}\pi g_s}{u^{3/2}}}(du^2+u^2d\mathrm{\Omega }_4),`$ $`e^{2\varphi }`$ $`=`$ $`{\displaystyle \frac{g_s^2u^{3/2}a^3R^{3/2}}{(1+a^3u^3)}},h={\displaystyle \frac{1}{1+a^3u^3}}.`$ (35) The warped metric near the shell, at a point which we can fix without loss of generality to $`(w_1,w_2,w_3)=(0,0,r_0)`$ is: $$ds_{string}^2=\frac{\sqrt{6}r_0\rho ^{3/2}}{R^{5/2}}\eta _{\mu \nu }dx^\mu dx^\nu +\frac{R^{5/2}}{\sqrt{6}r_0\rho ^{3/2}}(dwdw+dydy),$$ (36) where $`\rho ^2=(w_3r_0)^2+y^2`$ and $`R^{3/2}=\pi g_s\sqrt{\alpha ^{}}pa^{3/2}`$ , $`R^5=6\pi ^2g_sn\alpha ^{5/2}`$. For large $`au`$ the metric (35) matches (36) provided $`u=\rho \alpha ^{1/3},a^3={\displaystyle \frac{n\alpha ^{5/2}}{g_sr_0^2\pi ^2}},\stackrel{~}{x}^2={\displaystyle \frac{\pi g_s^{3/2}\sqrt{n}\alpha ^6n^2}{p^2r_0^4}}w^2.`$ (37) Let us briefly describe the ranges of validity of the supergravity solutions. There are two regions of interest. The D4 branes can be considered approximately flat when we are very close to them. This assumption is valid when $`\rho _cr_0`$. On the other hand, the supergravity solution (36) is valid when the dilaton is small, which is always the case near the crossover distance, $`au=1`$. This gives $$\rho _c(g_sp^2/n)^{1/3}\sqrt{\alpha ^{}}.$$ (38) Thus, the near shell approximation is consistent because $`\rho _c`$ is much smaller than the size of the sphere $`r_0`$. The supergravity approximation is valid if the size of the transverse sphere $`R^{3/2}/\sqrt{a}`$ is not too small. In the conformal case this quantity is independent of $`r_0`$. In our case the radius of the transverse 4-sphere is $$\left(\frac{\pi g_s^2\alpha ^2p}{nm^2}\right)^{1/6}.$$ (39) Therefore, we can trust the supergravity approximation for small $`m`$ and /or large $`p`$. ## 4 NS5 brane probes As we have seen, a background with nonzero $`H_3`$ and $`C_5`$ supports a D4 brane at nonzero radius. Nevertheless, the field equations of motion couple $`H_3`$ with $`F_4`$, so we expect the background to support an NS5 brane at nonzero radius as well. The type IIA $`NS_5`$ is quite a mysterious object with a self-dual 2-form field living on it, and its action has been constructed only recently . We first investigate an NS5 brane with the geometry $`R^3\times S^3`$, and with a large number $`n`$ of D2 branes dissolved in it, in the background created by an even larger number $`N`$ of D2 branes. Turning on a D2 brane charge $`n`$ is done by turning on a 3-form field strength along the $`S^3`$. This has to be done carefully because the field strength is self dual. As we shall see, self duality in this case is nontrivial. ### 4.1 The NS5 brane action The NS5 brane action in a background with no RR 1-form flux is very similar to the M5 brane action. There are two formulations of the M5 brane action. The first is a manifestly covariant formulation (first obtained by Pasti Sorokin and Tonin ) which involves an auxiliary fields. The second formulation (first obtained by Perry and Schwarz in , and generalized for general gravitational backgrounds in ) has no auxiliary fields but only 5-dimensional manifest covariance . This action is more amenable to explicit calculations, and can be obtained from the first one by gauge fixing. The formulation of the NS5 brane action in is very similar to the PST M5 brane action, in that it contains an auxiliary field, which imposes the self-duality of the 2-form field when integrated out. Moreover, in order to do explicit computations, the NS5 brane action has to be gauge fixed. In the absence of background RR one-forms, the gauge fixing of NS5 brane action is similar to the reduction of the PST action to the PS action in M-theory. For zero RR one-form flux and for the dilaton background (1) the NS5 brane action becomes: $`S`$ $`=`$ $`\mu _5{\displaystyle d^6\xi [L_{BI}+L_2+L_{WZ}]},`$ $`L_{BI}`$ $`=`$ $`g_s^2Z^{1/2}\sqrt{det(G_{mn}+ig_sZ^{1/4}D_{mn})},`$ $`L_2`$ $`=`$ $`\sqrt{G}{\displaystyle \frac{1}{4_ra^ra}}_ma(D)^{mnp}D_{npq}^qa,`$ $`L_{WZ}`$ $`=`$ $`[B_6{\displaystyle \frac{1}{2}}F_3C_3],`$ (40) where $`B_6`$ and $`C_3`$ are the pullbacks of the IIA forms, $`D_3F_3C_3`$, and $`D_{mn}=(D)_{mnp}\frac{^pa}{\sqrt{_ra^ra}}`$. The first and the last terms look like Born-Infeld and Wess-Zumino terms respectively. If we fix the gauge to $`a=x^2`$, $`b_{2m}=0`$, we recover an action with reduced explicit Lorentz invariance, similar to the PS action: $`S=\mu _5{\displaystyle d^6\xi g_s^2Z^{1/2}\sqrt{det(G_{mn}+ig_sZ^{1/4}D_{mn})}}+`$ $`+{\displaystyle \frac{1}{4}}\sqrt{{\displaystyle \frac{G}{G_{22}}}}D^{mn}D_{mn2}+L_{WZ},`$ (41) where $`D_2`$ has been changed to: $`D^{mn}={\displaystyle \frac{\sqrt{G_{22}}}{3!\sqrt{G}}}ϵ^{2mnpqr}D_{pqr}.`$ (42) The equations of motion are the same as in , with an extra $`g_sZ^{1/4}`$ multiplying $`1/\sqrt{G/G_{22}}`$. When only three-form fields in orthogonal directions are turned on, the equations simplify to give: $`D_{mn2}=_2B_{mn}C_{2mn}={\displaystyle \frac{\sqrt{G_{22}}D_{mn}}{\sqrt{1+g_s^2Z^{1/2}D_{mnp}D^{mnp}}}}.`$ (43) Under this assumption also $`\sqrt{det(G_{mn}+ig_sZ^{1/4}D_{mn})}=\sqrt{G}\sqrt{1+g_s^2Z^{1/2}D_{mnp}D^{mnp}}.`$ (44) ### 4.2 The NS5 brane probe We first consider a situation with four equal fermion masses, where we have SO(4) symmetry, but no supersymmetry. In the absence of supersymmetry the masses of the seven scalars are not constrained. We study a test configuration in which a large number of D2 branes is extended in the 012 directions; a probe NS5 brane has three directions wrapped on an $`S^3`$ in the $`3456`$ plane, and the other three are parallel to the D2 branes. We call $`\widehat{ϵ}_{ijkl}`$ the numerical antisymmetric tensor restricted to the $`3456`$ plane. We also give the NS5 brane a D2 brane charge $`n`$, by turning on a 3- form field strength along $`S^3`$: $$F_3=\frac{A}{r^43!}\widehat{ϵ}_{ijkl}x^idx^jdx^kdx^l=A\mathrm{sin}^2\theta \mathrm{sin}\theta _1d\theta d\theta _1d\varphi ,$$ (45) where $$A=4\pi n(\alpha ^{})^{3/2},$$ (46) and $`\theta `$ , $`\theta _1`$ and $`\varphi `$ are the angles on the three-sphere. For now we assume $`n<N`$, so that the effect of the brane probe on the background is negligible. Equation (16) also gives us $`T_{3456}=4m`$; thus, $$C_3^1=\frac{Z}{2g_s}S_3=\frac{4m}{2g_s3!}\left(\frac{R}{r}\right)^5\widehat{ϵ}_{ijkl}x^idx^jdx^kdx^l=\frac{2mR^5}{g_sr}\mathrm{sin}^2\theta \mathrm{sin}\theta _1d\theta d\theta _1d\varphi ,$$ (47) $`C_3^0`$ is given by (1). The 6 form dual of the NS-NS 2 form can be found using (3) from which we obtain $$dB_6\frac{1}{2}C_3F_4=F_7=e^{2\mathrm{\Phi }}H_3.$$ (48) Using (1), (9) and the relations in the Appendix, we obtain $$dB_6=[Z^1(\frac{_7H_3}{g_s^2}+\frac{F_4^1}{g_s})+\frac{1}{2}d(C_3^1C_3^0)]dx^0dx^1dx^2=0.$$ (49) Thus, the first term in the Wess-Zumino action gives no contribution, like in the M-theory case . This is to be contrasted with the non-vanishing $`dC_5`$ in the D4 brane case and the non-vanishing $`dC_6`$ and $`dB_6`$ in . Since we have spherical symmetry, the value of the action is the same at every point on the 3-sphere. At the point $`x^6=r`$: $$D_{}=F_{345}C_{345}=\left[\frac{A}{r^3}\frac{2mR^5}{g_sr^4}\right]=\frac{A}{r^3}C_{345}.$$ (50) We are interested in the limit when the D2 brane charge of the NS5 brane is bigger than its NS5 brane charge. This means $`n^2g_s\sqrt{\alpha ^{}}rN`$. In this limit, $`F_{345}`$ will be dominant in the Born-Infeld term. Using the equation of motion (43) we obtain: $$_2B_{01}C_{201}=\frac{\sqrt{G_{}G_{}^1}D_{}}{\sqrt{1+g_s^2Z^{1/2}G_{}^1D_{}^2}}=\frac{D_{}Z^{3/2}}{\sqrt{1+g_s^2D_{}^2Z^1}}.$$ (51) To find $`F_{012}`$ we observe that for large D2 brane charge the right hand side of (51) gives exactly the background value of $`C_{201}`$, so the contribution of $`F_{012}`$ to the Wess-Zumino term is negligible. We have discovered a very interesting fact. The equations of motion of an NS5 brane with large D2 brane charge in the geometry (1) give rise to a “dual” 3 form equal to the 3-form field of this geometry. The bulk $`C_3`$ created by the dissolved branes is the same as the one obtained via the NS5 brane equation of motion. This is an interesting connection between the IIA supergravity and the NS5 brane equations of motion which deserves further study. Note that this result is independent of $`G_{}`$, and thus, it holds for any NS5 brane shape at any warp factor. This phenomenon is very similar to the one observed in the case of the M-theory 5-brane , and is probably a general feature of actions with a self-dual field. We now proceed to the calculation of the potential per unit longitudinal volume felt by the NS5 brane in this geometry. The dominant parts of the potential are: $`{\displaystyle \frac{S_{BI}}{2\pi ^2\mu _5V}}`$ $`=`$ $`r^3{\displaystyle \frac{Z^{1/2}}{g_s^2}}\sqrt{G_{}G_{}}\sqrt{1+g_s^2D_{}^2Z^{1/2}}{\displaystyle \frac{A}{g_sZ}}+{\displaystyle \frac{r^6}{2g_s^3A}}+{\displaystyle \frac{r^3C_{345}}{g_sZ}},`$ $`{\displaystyle \frac{S_{mixed}}{2\pi ^2\mu _5V}}`$ $`=`$ $`{\displaystyle \frac{r^3}{2}}{\displaystyle \frac{D_{}^2Z^{3/2}}{\sqrt{1+g_s^2D_{}^2Z^1}}}{\displaystyle \frac{A}{2g_sZ}}+{\displaystyle \frac{r^6}{4g_s^3A}}{\displaystyle \frac{r^3C_{345}}{2g_sZ}},`$ $`{\displaystyle \frac{S_{WZ}}{2\pi ^2\mu _5V}}`$ $`=`$ $`{\displaystyle \frac{r^3}{2}}[2B_{012345}^6+C_{012}F_{345}F_{012}C_{345}]{\displaystyle \frac{A}{2g_sZ}}.`$ (52) where the approximation is valid for $`n^2g_s\sqrt{\alpha ^{}}rN`$. We have also included the $`B^6`$ component which in the probe brane case is zero, but might be non-vanishing for a general warped geometry. The first terms represent the interaction of the dissolved D2 branes. They are proportional to $`g_s^1`$, as D brane action terms are, and they cancel because parallel D2 branes do not interact. The terms proportional to $`r^6`$ are the gravitational energy of the NS5 brane, and they are attractive. The terms proportional to $`C^3`$ are repelling. The effect of those two terms is to favor a ground state at nonzero radius. Nevertheless, as in the D4 brane case, there can be another term in the potential of the same relevance as the first two. We expect a second order correction to $`C_{012}`$, via the second equation in (3) as well as second order corrections to the metric and dilaton. Using the first order fields (9) and remembering that $`C_{012}`$ couples with $`F_{345}`$, we can see that the extra term is proportional to $`m^2Ar^2/g_s`$, and is of the same order as the first two terms. The potential is therefore: $$\frac{S}{2\pi ^2\mu _5V}=\frac{3r^6}{4g_s^3A}\frac{m}{g_s^2}r^4+\frac{cAm^2r^2}{g_s},$$ (53) where $`c`$ is a constant. Determining $`c`$ is not trivial because of the second order corrections to the metric and the dilaton. In the first part of this paper, as well as in $`c`$ is computed by invoking supersymmetry. Unfortunately, with four fermion masses, there is no supersymmetry. If we give up one fermion mass, and thus give up SO(4) invariance for SO(3) invariance, we can still use supersymmetry to determine the value of $`c`$. This will be the subject of section 4.4. We can also find $`c`$ by noticing that the last term of (53) represents a mass for the scalars. This will be the subject of chapter 4.5. ### 4.3 NS5 brane warping In the general case the D2 branes are polarized into a sphere with NS5 brane charge. We try to find the potential for this configuration. Since we are doing a ‘self-force’ problem, we have to find the potential by bringing small D2 brane spheres one by one from infinity. In classical electromagnetism this gives the familiar factor of 1/2 for the self energy of a charged object. We first consider a test NS5 brane shell with D2 charge in the geometry created by a large number of polarized D2 branes. The only difference from the previous chapter is the new $`Z`$. We can use an argument similar to the one we have used for D4 branes to show that the potential does not change. The equations (5) imply that $`_7H_3+g_sF_4^1`$ is harmonic. Its behavior at infinity is given by the boundary theory, and is the same as for trivial $`Z`$. Therefore, $$_7H_3+g_sF_4^1=2ZT_4.$$ (54) Using (49) and (52) we see that the particular combination which enters the NS5 brane action, $`B_6\frac{1}{2}C_3^1C_3^0`$, is not changed. This “conspiracy” is similar to the one in . The $`C^1C^0`$ term (coming from the BI and the mixed action) appears in the effective NS5 brane action and in the first order bulk equation of motion with the same coefficient relative to $`B_6`$. Thus, in a general $`Z`$ background $`B_6`$ and $`C_3^1`$ may change, but an NS5 brane with large D2 charge does not feel that. The large terms in (52) cancel as usual, and the terms proportional to $`r^6`$ do not depend on $`Z`$. The last term may change, but in the configurations we care about it is fixed by supersymmetry, so it does not depend on $`Z`$ either. Since the potential felt by a probe brane does not depend on the distribution of branes, the total potential is the same as the probe potential. One may also worry about the effect of the NS5 branes on themselves. The NS5 brane is a source in the Bianchi identity for $`H_3`$ (3), but does not affect the equations (5). Therefore, the terms which enter the NS5 brane action are not changed, and thus, there is no NS5 self coupling. ### 4.4 Unequal masses As we have mentioned above, we would like to study a supersymmetric case, where we can compute the third term in the potential. We will examine the case when three fermion masses are equal, and the fourth is taken to zero. In this limit supersymmetry is restored, and the 3-sphere becomes a very long SO(3) invariant ellipsoid. In (16) we chose $`m_1=m_2=m_3=m`$, and $`m_4m`$. We also assume that the ellipsoid is situated in the 3456 plane. This configuration can be generalized to other $`SO(3)`$ invariant configurations by a simple phase rotation by the same angle in the 3-7,4-8 and 5-9 planes. The ellipsoid is parametrized $`x^6`$ $`=`$ $`\alpha r\mathrm{cos}\theta ,`$ $`x^3`$ $`=`$ $`r\mathrm{sin}\theta \mathrm{cos}\theta _1,`$ $`x^4`$ $`=`$ $`r\mathrm{sin}\theta \mathrm{sin}\theta _1\mathrm{sin}\varphi ,`$ $`x^5`$ $`=`$ $`r\mathrm{sin}\theta \mathrm{sin}\theta _1\mathrm{cos}\varphi ,`$ (55) where $`\alpha `$ gives the squashing of the ellipsoid, and depends on the ratio $`m_4/m`$. Since in general the induced fields are not constant on the ellipsoid, it is more useful to use $`\theta ,\theta _1`$ and $`\varphi `$ as coordinates. The field strength $`F_{\theta \theta _1\varphi }`$ (45) remains unchanged, while $`C_{\theta \theta _1\varphi }`$ (47) is multiplied by $`\alpha `$ and modified by the change in (16). Also $$G_{}=r^6Z^{3/2}\mathrm{sin}^4\theta \mathrm{sin}^2\theta _2(\alpha ^2\mathrm{sin}^2\theta +\mathrm{cos}^2\theta ).$$ (56) We have observed in chapter 4.2 that the equations of motion of an NS5 brane with large D2 charge create a dual field almost equal to the background $`C_{012}`$, regardless of $`G_{}`$, so there is no new contribution from the Wess-Zumino term. The dominant terms of the potential do not depend on $`G_{}`$, so they have the same value as before, and they cancel as in the SO(4) invariant case. The terms which before were proportional to $`r^6/A`$ are now variable on the ellipsoid, and their value is: $$V_{r^6}_{E^3}\sqrt{G_{}}\frac{G_{}}{D_{\theta \varphi \alpha }}.$$ (57) The integral gives: $$\frac{S_{r^6}}{2\pi ^2\mu _5V}=\frac{3r^6}{4A}\frac{3\alpha ^2+1}{4}.$$ (58) We can interpret $`x^6`$ as the real part of a complex variable $`w`$, whose imaginary part is forced to zero. This picture fits very well with the intuition coming from reducing the M-theory picture to IIA. Indeed, as we go from M-theory to IIA the eleventh dimension $`x^{10}`$ shrinks to zero. We can put the back phase of $`r`$ and obtain the potential: $$\frac{S}{2\pi ^2\mu _5V}=\frac{3}{16g_s^3A}\left(3\|z\|^4\|w\|^2+\|z\|^6\right)\frac{1}{4g_s^2}\mathrm{Re}(3mwzz\overline{z}+m_4z^3\overline{w}),$$ (59) where the phase of $`w`$ does not have any meaning. In the limit $`m_40`$ the theory becomes supersymmetric, and it is possible to complete the square in (59) to obtain: $`{\displaystyle \frac{S}{2\pi ^2\mu _5V}}`$ $`=`$ $`{\displaystyle \frac{3}{16g_s^3A}}\left(3\|z\|^4\|w\|^2+\|z\|^6\right){\displaystyle \frac{1}{4g_s^2}}\mathrm{Re}(3mwzz\overline{z}+m_4z^3\overline{w})`$ (60) $`+`$ $`{\displaystyle \frac{A}{12g_s}}(3m^2\|z\|^2+m_4^2\|w\|^2).`$ As we explained, the last term in this potential comes from second order corrections to the bulk $`C_{012}`$, metric and dilaton. Since these second order fields also give the last term in the D4 brane effective potential (30), we expect the last terms to be similar. Indeed, they have the same dependence on $`g_s,m^2,r^2`$, and only differ by a numerical factor (1/2). This is not unexpected in view of the different ways the dilaton and the graviton enter the D4 respectively the NS5 brane actions. The potential has supersymmetric minima at: $`z^2={\displaystyle \frac{2Ag_s}{3}}m\sqrt{{\displaystyle \frac{m_4}{m}}},`$ $`x_6^2={\displaystyle \frac{2Ag_s}{3}}m\sqrt{{\displaystyle \frac{m}{m_4}}},`$ (61) where $`A=4\pi n(\alpha ^{})^{3/2}`$. The ellipsoid has aspect ratio $`\alpha =\sqrt{\frac{m}{m_4}}`$, and degenerates into a very long line in the supersymmetric limit. In this limit the NS5 flux on the opposite sides of the ellipsoid adds to zero. Thus, we will just have D2 branes at different positions along $`x^6`$. Like in the study of D2-D4 polarizations, $`x^6`$ is a modulus. A general vacuum configuration will be given by a combination of D4 brane spheres and NS5 brane ellipsoids, with D2 brane charges $`n_i`$ and at positions given by (32) or (61). We allow $`m_4`$ to be small but nonzero, so that we can study the NS5 vacua, and in the same time use supersymmetry. This can make the supersymmetric ground state slightly metastable (with a tunneling rate which becomes infinite as $`m_40`$ ). However, the picture in which the D2 branes are polarized into an NS5 brane ellipsoid is valid. It is interesting to fully study the case with four fermion masses turned on, which depending on the magnitude of the third term may or may not produce polarization into spherical NS5 branes. In M-theory this configuration is supersymmetric, and M2 branes are polarized into M5 branes wrapped on three-spheres. As we flow into IIA we lose the supersymmetry, but intuitively the branes should still remain polarized. ### 4.5 SO(4) invariant non supersymmetric solutions The general scalar mass term in the boundary theory is: $$m^2(\mathrm{\Phi }_3^2+\mathrm{\Phi }_4^2+\mathrm{\Phi }_5^2+\mathrm{\Phi }_7^2+\mathrm{\Phi }_8^2+\mathrm{\Phi }_9^2+\mathrm{\Phi }_6^2)+\mu _{ij}\mathrm{\Phi }_i\mathrm{\Phi }_j,$$ (62) where $`\mu _{ij}`$ are general $`7\times 7`$ traceless symmetric matrices. The scalar $`\mathrm{\Phi }_6`$ is different from the other scalars because it belongs to the vector multiplet in the supersymmetric case. From (60) we see that supersymmetry constrains the last term of (62) to be $$\mu _{ij}\mathrm{\Phi }_i\mathrm{\Phi }_j=\frac{m^2}{6}(\mathrm{\Phi }_3^2+\mathrm{\Phi }_4^2+\mathrm{\Phi }_5^2+\mathrm{\Phi }_7^2+\mathrm{\Phi }_8^2+\mathrm{\Phi }_9^26\mathrm{\Phi }_6^2).$$ (63) As we have seen in the case of D4 branes and NS5 branes, this mass term is the last term in the action of a wrapped brane. It appears in the supergravity calculation through the second equation in (3), where the fields (9) give a contribution: $$C_{012}^2m^2/g_s,$$ (64) and similar ones to the dilaton and the metric. This determines in principle the last term of the action. Nevertheless, the supergravity equations of motion satisfied by $`C_{012}`$ have homogeneous solutions, and as explained in , the solution is not determined completely by the square of the first order perturbations in (9). These perturbations, which are proportional to $`T^3T^4`$, only determine the mass of the $`L=0`$ mode. There is another $`L=2`$ piece, the second term in (62), which can be specified arbitrarily in the boundary theory. This arbitrariness is not present in the supersymmetric theory (63). Since we want to investigate an SO(4) invariant situation (in the 3456 plane), we set $`m_4=m`$ in (16), and set $`\mu _{ij}=0`$. The last term of (59) is $$\frac{S_{r^2}}{2\pi ^2\mu _5V}=\frac{A}{12g_s}(3m^2)r^2=\frac{A}{12g_s}(3m^2)\frac{8}{3\pi }_0^\pi 𝑑\theta \mathrm{sin}^2\theta (x_3^2+x_4^2+x_5^2+x_7^2+x_8^2+x_9^2),$$ (65) where the normalization of the last term is obtained using (55). This is properly interpreted as a combination of an $`L=0`$ and $`L=2`$ mode: $`{\displaystyle \frac{S_{r^2}}{2\pi ^2\mu _5V}}`$ $`=`$ $`{\displaystyle \frac{A}{12g_s}}(3m^2){\displaystyle \frac{8}{3\pi }}{\displaystyle _0^\pi }d\theta \mathrm{sin}^2\theta [{\displaystyle \frac{6}{7}}(x_3^2+x_4^2+x_5^2+x_7^2+x_8^2+x_9^2+x_6^2)`$ (66) $`+`$ $`{\displaystyle \frac{1}{7}}(x_3^2+x_4^2+x_5^2+x_7^2+x_8^2+x_9^26x_6^2)].`$ The coefficient of the $`L=0`$ term is given by $`T_3T_3`$ and is proportional to $`3m^2`$. Turning on the fourth fermion mass changes this to $`4m^2`$. Thus the SO(4) invariant action is given by $$\frac{S}{2\pi ^2\mu _5V}=\frac{3r^6}{4g_s^3A}\frac{m}{g_s^2}r^4+\frac{8Am^2r^2}{21g_s}.$$ (67) This has a local minimum at $$r=\frac{2\sqrt{g_sAm}}{3}\sqrt{1+\frac{1}{\sqrt{7}}}.$$ (68) We can see that the energy of the NS5 brane at this radius is higher than the one at the origin. It is interesting to explore if other possible D2 brane polarizations have lower energies than this configuration. We can easily find the new D4 brane effective potential for the SO(4) invariant case. As we can see from (16) the second term of (29) gets modified: $$\mathrm{Im}(3z^2\overline{z})\mathrm{Im}(3z^2\overline{z}+z^3).$$ (69) Since this is the term which keeps the brane from collapsing and is the only phase dependent term, we can see that the orientation which maximizes it is $`z=re^{i\pi /6}`$. The last term also changes, in exactly the same way as in the NS5 case. The potential felt by the D4 brane in the plane most favorable to polarization is: $$\frac{S}{V}=\frac{2\mu _4}{g_sn\alpha ^{}}(r^45/3r^3r_0+32/21r^2r_0),$$ (70) where $`r_0=\pi mn\alpha ^{}`$ is the old polarization radius. This potential has only one minimum, at $`r=0`$, and no metastable minima at $`r0`$. Therefore D2 branes cannot polarize into D4 branes. This means that the SO(4) invariant theory does not have any Higgs vacua. Unfortunately, we do not know all the vacua of the theory, so we cannot claim based on this that the lowest energy vacuum is confining. ### 4.6 The near shell metric of an NS5 brane wrapped on a 3-ellipsoid Let us try to find the metric near the NS5 brane ellipsoidal shell. As in the D4 brane case, we obtain the near shell limit of (1), and we match it with the metric near a flat NS5 brane with D2 brane flux. If we denote by $`\stackrel{}{w}`$ the coordinate in the 4-plane of the ellipsoid, and by $`y`$ the radius in the 3-plane transverse to the ellipsoid, we can express the warp factor as: $$Z=c_{ellipsoid}\frac{\sigma (\stackrel{}{r})}{((\stackrel{}{r}\stackrel{}{w})^2+y^2)^{5/2}},$$ (71) where $`\stackrel{}{r}`$ is the coordinate on the ellipsoid, $`\sigma (\stackrel{}{r})`$ is the D2 charge density, and $`c`$ is a normalization constant. The D2 charge is proportional to the absolute value of the 3-form field strength $`F_{\theta \theta _1\varphi }`$ given by (45). We are interested in the limit of $`Z`$ as we approach the ellipsoid in two points, as indicated in the Figure 2. For point 1 we evaluate $`Z`$ at $`\stackrel{}{w}=(0,w,0,0)`$: $$Z=\frac{R^5}{\pi }\frac{\mathrm{sin}^2\theta \mathrm{sin}\theta _1d\theta d\theta _1}{(r_0^2\alpha ^2\mathrm{cos}^2\theta +r_0^2\mathrm{sin}^2\theta +w^2+y^22wr_0\mathrm{sin}\theta \mathrm{cos}\theta _1)^{5/2}},$$ (72) where the normalization constant is chosen to match (2) far away from the shell. As before $`\rho ^2(wr_0)^2+y^2`$. In the $`\rho 0`$ limit we obtain $$Z\frac{2R^5}{3\pi \rho ^2r_0^3\alpha }.$$ (73) For point 2 we evaluate $`Z`$ at $`\stackrel{}{w}=(w,0,0,0)`$: $$Z=\frac{2R^5}{\pi }_0^\pi \frac{d\theta \mathrm{sin}^2\theta }{(\alpha ^2r_0^2\mathrm{cos}^2\theta 2\alpha r_0\mathrm{cos}\theta w+w^2+y^2)^{5/2}},$$ (74) where we chose normalization so that $`Z`$ matches (2) for large $`w`$ and $`y`$. This integral can be conveniently rewritten as $$Z=\frac{R^5}{2(\alpha r_0)^5}_0^\pi \frac{d\theta \mathrm{sin}^2\theta }{((\mathrm{cos}\theta x1)^2+\overline{y}^2)^{5/2}},$$ (75) where $`x=(w\alpha r_0)/(\alpha r_0),`$ $`\overline{y}=y/(\alpha r_0).`$ (76) On dimensional grounds we can write (75) in the limit $`x0`$ $`\overline{y}0`$ as $$\frac{R^5}{(\alpha r_0)^{3/2}}\frac{1}{\rho ^{7/2}}.$$ (77) where we have ignored a possible homogeneous function of $`x`$ and $`\overline{y}`$. In order to find the tension of flux tubes, and to find where strings drooping from infinity attach on the ellipsoid, we have to compare (77) with (73). For large $`\alpha `$, the warp factor near 1 is far bigger than that near 2. Thus, flux tubes attached to 1 will have less energy than those attached to 2. In what follows we will derive the near shell metric in the vicinity of point 1 and in the next chapter calculate the tension of the flux tube attached there. In , the near shell metric of D3 branes polarized into IIB NS5 branes was matched with the warped metric in the $`\rho 0`$ limit. The warp factor is $$Z_c\frac{R_c^4}{4\rho ^2r_{0}^{}{}_{c}{}^{2}},,$$ (78) where we put subscript $`c`$ on all quantities from to avoid confusion. In it was explained how to obtain the metric for IIA NS5 branes with D2 brane charge by T-dualizing the corresponding IIB NS5 metric with three brane charge. Keeping all constants of the original IIB solution intact, we obtain<sup>1</sup><sup>1</sup>1Compare for example (38) and (35) of with (103) of .: $`ds_{string}^2`$ $`=`$ $`{\displaystyle \frac{2r_{0}^{}{}_{c}{}^{}(\rho ^2+\rho _c^2)^{1/2}}{R_c^2}}\eta _{\mu \nu }dx^\mu dx^\nu +{\displaystyle \frac{R_c^2}{2r_{0}^{}{}_{c}{}^{}(\rho ^2+\rho _c^2)^{1/2}}}(dw^idw^i)`$ (79) $`+`$ $`{\displaystyle \frac{R_c^2(\rho ^2+\rho _c^2)^{1/2}}{2r_{0}^{}{}_{c}{}^{}}}(dw^4dw^4+dydy),`$ $`e^{2\varphi }`$ $`=`$ $`g_s^2\alpha ^2{\displaystyle \frac{\sqrt{\rho ^2+\rho _c^2}}{\rho ^2}},r_{0}^{}{}_{c}{}^{}=m_c\pi \alpha ^{}g_sN,`$ $`\rho _c`$ $`=`$ $`{\displaystyle \frac{2r_{0}^{}{}_{c}{}^{}\alpha ^{}}{R_c^2}},R_c^4=4\pi g_sN\alpha ^2,`$ where the $`i`$’s run from 1 to 3. Let us stress that the constants in the above formula do not have any direct physical interpretation in type IIA theory. We have a metric which matches (1) with (78). What we need is a metric which matches (1) with (73). The required metric in IIA language is obtained by a formal identification of $`Z`$ and $`Z_c`$ in the large $`\rho `$ limit. Then $`\pi g_sNm_c^2=3\pi r_0^3\alpha /R^5,r_0^2={\displaystyle \frac{8\pi ^2}{3\alpha }}N\alpha ^{3/2}g_sm.`$ (80) We also obtain $$\rho _{\mathrm{crossover}}r_0\left(\frac{\alpha ^3m\sqrt{\alpha ^{}}}{N\alpha g_s}\right)^{1/4}.$$ (81) The crossover distance is small compared to $`r_0`$ which makes the near shell approximation valid. Since the transverse three sphere has size $`\alpha ^{}`$ for a single NS5 branes, string theory corrections are marginal for the near shell metric <sup>2</sup><sup>2</sup>2The action for coincident NS5 branes is not known.. This completes the derivation of near shell metric for NS5 brane ellipsoids with D2 flux. ### 4.7 The near shell metric of an NS5 brane wrapped on a 3-sphere If the NS5 brane is polarized into a 3-sphere the warp factor is $`Z={\displaystyle \frac{2R^5}{\pi }}{\displaystyle _0^\pi }{\displaystyle \frac{\mathrm{sin}^2\theta d\theta }{(r_0^2+w^22r_0w\mathrm{cos}\theta +y^2)^{5/2}}}=`$ $`{\displaystyle \frac{R^5}{(r_0^2+w^2+y^2)^{5/2}}}{}_{2}{}^{}F_{1}^{}({\displaystyle \frac{5}{4}},{\displaystyle \frac{7}{4}};2;{\displaystyle \frac{4r_0^2w^2}{(r_0^2+w^2+y^2)^2}}).`$ (82) Introducing as before the variable $`\rho =(wr_0)^2+y^2`$ and working in the $`\rho 0`$ limit we find $$Z\frac{2R^5}{3\pi \rho ^2r_0^3}.$$ (83) Let us note that we could have obtained (83) from (73) by simply making the ellipsoid a sphere ($`\alpha =1`$). However $`r_0`$ for the ellipsoid is different than $`r_0`$ for the sphere. The metric is given by (79), where instead of (80) we make the identification $$\pi g_sNm_c^2=3\pi r_0^3/R^5,r_0=\frac{4\sqrt{g_sN\pi \alpha ^{3/2}m}}{3}\sqrt{1+\frac{1}{\sqrt{7}}}.$$ (84) We will use this formula in the next section to calculate the tension of flux tubes in this vacuum. ## 5 Gauge Theory In this section we discuss the string theory duals of the gauge theory objects.First we compute the tension of flux tubes, and show that a vacuum dual to D2 branes polarized into NS5 branes is confining. Then we give string theory interpretations of domain walls, condensates, baryon vertices and glueballs. ### 5.1 Flux tubes In the framework of the AdS/CFT duality, there is a correspondence between Wilson lines in the boundary theory and fundamental string world sheets in the bulk. This correspondence was not initially extended to non conformal theories because the geometry is not valid throughout the bulk. For Dp branes with $`p<3`$ the string world-sheet would probe regions near the branes where supergravity is not valid because of the large dilaton. Brane polarization solves this problem by smoothing out the geometry. Therefore, we can do concrete computations of flux tube tensions in various vacua. In a vacuum with D2 branes polarized into D4 branes, a fundamental string lowered from the boundary attaches to the D4 brane shell. It costs no energy to move its ends apart on the boundary. Therefore, vacua with only D4 branes are screening. Another way to see this is by noticing that the parallel components of $`G_{\mu \nu }`$ in (35) go to zero as we approach the D4 branes. Let us find the tension of a hanging string in a vacuum with NS5 branes. Since $`Z`$ diverges at small distances from the D2-NS5 shell, the string can lower its tension by getting closer to the shell. However, at some point the geometry (1) ceases to be valid, and we need to use the near shell metric (79). In contrast to the D4 case, $`G_{}`$ has a minimum value of $`\pi \alpha ^{}m_c^2g_sN`$. Thus, a flux tube has a finite tension: $$\tau _e=m_c^2g_sN/2=4\pi \left(\frac{2m}{3}\right)^{3/2}\frac{\sqrt{Ng_s/\alpha }}{\alpha ^{1/4}}.$$ (85) The tension of the flux tube scales as expected with $`g_{YM}^2N`$ and vanishes in the limit of a strongly squashed ellipsoid. The tension of the fundamental string attached to a spherically wrapped NS5 brane is $$\tau _e=\frac{16m^{3/2}\sqrt{g_sN/\pi }}{27\alpha ^{1/4}}\left(1+\frac{1}{\sqrt{7}}\right)^{3/2}.$$ (86) ### 5.2 Domain walls Since our theory has many discrete vacua, we expect to have domain walls. The simplest to understand are the domain walls between vacua containing D4 branes. If the two vacua have the same number of D4 branes, they meet on a 2-sphere where they exchange the D2 brane charge. The tension of the domain wall is the bending tension of this configuration. If the vacua have different numbers of D4 branes, they meet on a 2-sphere, and by charge conservation the appropriate number of D4 branes span the 3-ball that has the 2-sphere as the boundary. Let us try to understand the domain wall between a vacuum with an $`S^2`$ brane wrapped D4 brane and the one with an NS5 brane wrapped on a 3-ellipsoid. The NS5 brane is extended in the 3456 direction, while the D4 brane is extended in the 789 directions. As they get close they bend towards each other, but since the D4 is lighter it bends more. By charge conservation at the intersection point there is an NS5 brane with D4 brane charge on it, spanning the squashed ball that has the ellipsoidal boundary. There are also other vacua in this theory, characterized by oblique configurations of NS5 branes with D4 flux. They will be studied together with their domain walls in . ### 5.3 Condensates The coefficient of the normalizable mode in the polarized brane background determines the value of gauge theory condensates. The highest normalizable mode (10) decays at infinity as $`r^7`$. Its coefficient $`\overline{\alpha }`$ gives the value of a condensate containing the fermion bilinear and its supersymmetric partners. As an example consider the normalizable mode sourced by a D4 spherical shell. Since the D4 brane is a magnetic source for $`F_4^1`$, the Bianchi identity for $`F`$ (3) gets modified to: $$dF_4=J_5,$$ (87) with $$J_5\alpha ^{3/2}\delta ^4(y)\delta (wr_0)dwd^4y,$$ (88) where the 2-sphere has radius $`r_0`$ in the $`w`$ 3-plane and the 4 transverse coordinates are denoted by $`y`$. Without doing the computation we can see by dimensional analysis that $$H_3^{\mathrm{normalizable}}\frac{\alpha ^{3/2}r_0^4}{r^7}dw^1dw^2dw^3.$$ (89) This gives $$<\lambda \lambda >m^4N^4\alpha ^{}.$$ (90) The exact numerical coefficient of the normalizable mode can be computed through a straightforward generalization of the procedure in . Nevertheless, since the condensate contains a combination of the fermion bilinear and its supersymmetric partners, the individual condensates are not known. It is harder to find the condensates in the confining vacuum. The ellipsoidal geometry complicates things. However, we can estimate the condensates by dimensional analysis $$<\lambda \lambda >1/\alpha ^k(mN)^{5/2}\alpha ^{1/4},$$ (91) where the power $`k>0`$ of the squashing ratio $`\alpha `$ cannot be determined by dimensional analysis alone. Note that the condensate disappears in the limit of an ellipsoid squashed to a line since the Ramond-Ramond charges on the opposite sides of the ellipsoid cancel. The condensate corresponding to the NS5 brane wrapped on a three sphere is simply $$<\lambda \lambda >(mN)^{5/2}\alpha ^{1/4}.$$ (92) ### 5.4 Baryon Vertices. It was argued in that the dual of a baryon vertex in a 4-dimensional $`𝒩=4`$ $`SU(N)`$ SYM is a D5 brane wrapped on the five sphere in the $`AdS_5\times S^5`$ geometry. In the 3-dimensional $`𝒩=8`$ theory one expects to find the dual of a baryon vertex in terms of a D6 brane of geometry $`R\times S^6`$. At a very high energy the non perturbed picture of the baryon vertex is still valid. We expect this picture to change when the energy becomes comparable to the fermion masses. In a confining vacuum, as a shrinking D6 brane crosses the NS5 brane shell, a D4 brane is created via the Hanany-Witten effect . This D4 brane fills the ball whose boundary is the NS5 three-sphere (or three-ellipsoid). A D4 brane ending on an NS5 brane is a source for the 1-form field strength living on the NS5 brane. The D2 branes dissolved in the NS5 brane become $`F_3`$ flux. The world-volume action of the NS5 brane contains a term : $$F_3B_2.$$ (93) We can see that in order for this term to be invariant under the gauge transformation $`\delta B_2=d\xi _1`$, $`N`$ fundamental strings should end on the D4-NS5 junction. Therefore this configuration represents a baryon vertex. The energy of this D4 vertex is $$\mu _4_{B^4}d^4xe^\mathrm{\Phi }G_{induced}^{1/2}=\frac{\mu _4}{g_s}_{B^4}d^4xZ^{1/2}.$$ (94) The integral is divergent near $`r_0`$ and has to be regularized at $`r_0\rho _c`$. The warped $`Z`$ is given in terms of the hypergeometric function (82). To get the leading contribution to the baryon vertex mass we use the small $`\rho `$ limit (83) to obtain: $`{\displaystyle \frac{\mu _4(2R^5/3\pi r_0^3)^{1/2}}{g_s}}{\displaystyle ^{r_0\rho _c}}2\pi ^2w^3𝑑w{\displaystyle \frac{1}{wr_0}}`$ (95) $`N\left({\displaystyle \frac{m^3Ng_s}{\sqrt{\alpha ^{}}}}\right)^{1/4}\mathrm{ln}(g_sN)\left({\displaystyle \frac{1}{3\sqrt{\pi }}}\sqrt{1+{\displaystyle \frac{1}{\sqrt{7}}}}\right)^{3/2}.`$ (96) It would be interesting to also understand baryon vertices in the other vacua. ### 5.5 Instantons. Our candidates for instantons, objects of space time dimension 0, are either fundamental strings wrapped around the D4 two-sphere or D2 branes wrapped on the NS5 three-ellipsoid. None of this configurations is stable in string theory: a wrapped string can lower its tension by attaching to the D4 brane and unwrapping. Similarly a wrapped D2 brane can lower its tension by attaching and sliding from the NS5 brane. ### 5.6 Glueballs and other states. A supergravity state localized on a two sphere has transverse momentum $`k_m1/r_0`$. Its energy is found by requiring $$G_{\mu \nu }k^\mu k^\nu G_{mn}k^mk^n,$$ (97) which gives the glueball mass $$k_\mu \left(\frac{r_0}{R}\right)^{5/2}\frac{1}{r_0}\frac{\alpha ^{1/4}m^{3/2}N}{g_s^{1/2}}.$$ (98) A typical string state, on the other hand has $$G_{\mu \nu }k^\mu k^\nu 1/\alpha ^{},$$ (99) which translates into the energy scale $$k_\mu \frac{\alpha ^{1/8}m^{5/4}N}{g_s^{1/4}}.$$ (100) We encounter the same puzzle as in . There is no sign of gauge theory states with masses given in the previous equations. ## 6 Flows to M theory. It is interesting to examine what happens to our configurations as we lift them to M-theory. Both the NS5 and the D2 branes correspond to unwrapped M5 and M2 branes respectively. Thus, the D2-NS5 configurations just become the M2-M5 configurations studied in . In the case of D2-D4 configurations, as we increase the radius of the 11’th dimension, the D2 branes remain localized, while the D4 branes reveal their true nature. They become M5 brane cylinders of geometry $`R^3\times S^2\times S^1`$. As the radius of $`S^1`$ grows beyond the radius of the $`S^2`$, the M5 brane discovers that it is energetically favorable to be wrapped on a 3-sphere (or 3-ellipsoid). Thus, there will be a phase transition between these configurations when $`g\sqrt{\alpha ^{}}r_0`$. We have not studied oblique configurations (NS5 branes with D4 flux) in this paper. If they exist, it would be interesting to see what happens to them in the M-theory limit. Intuitively one may expect them to become either a tilted M5 brane (although this configuration is not a ground state in M theory), or two M5 branes polarized in orthogonal directions. ## 7 Conclusions and future directions. From a general relativity viewpoint the most important conclusion of this paper is the absence of naked singularities in the gravity duals of a perturbed 3-dimensional gauge theory. We have shown the existence of supersymmetric Higgs and non supersymmetric confining vacua in a three-dimensional Yang Mills theory. We have found their IIA string theory duals in terms of D2 branes polarized into D4 and NS5 branes respectively. The Higgs vacuum of D2-D4 system is expected from the T duality of D3-D5 system studied by Polchinski and Strassler. We found a nontrivial realization of a baryon vertex as D4 filling the inside of the NS5 brane 3-sphere. Other important results of this paper are the calculation of tensions of flux tubes, an estimate of condensates and discussion of domain walls in Yang Mills theory with no supersymmetry. There are still some unsolved problems. First, it would be interesting to better understand oblique vacua (which may correspond to D2 branes polarized into NS5 branes with D4 charge). Second, the tension of domain walls and their shape could be obtained. Also, the nontrivial connection between the NS5 brane action and the supergravity equations of motion deserves further study. The warping mechanism responsible for the resolution of singularities awaits a fuller understanding. It is the hope of the authors that this paper is a step in that direction. ## 8 Acknowledgements We are very grateful to Joe Polchinski for suggestions and for reading the manuscript. We thank Alex Buchel, and Costin Popescu for numerous discussions. We also thank Marina Nudelman for help with proofreading. The work of I.B. was supported in part by NSF grant PHY97-22022 and the work of A.N. was supported in part by DOE contract DE-FG-03-91ER40618. ## 9 Appendix Some useful relations between four-tensor spherical harmonics $`T_4`$ $`=`$ $`{\displaystyle \frac{1}{4!}}T_{mnpk}dx^mdx^ndx^pdx^k,`$ $`V_{mnpk}`$ $`=`$ $`{\displaystyle \frac{x^q}{r^2}}(x^mT_{qnpk}+x^nT_{mqpk}+x^pT_{mnqk}+x^kT_{mnpq}),`$ $`V_4`$ $`=`$ $`{\displaystyle \frac{1}{4!}}V_{mnpk}dx^mdx^ndx^pdx^k,`$ $`S_3`$ $`=`$ $`{\displaystyle \frac{1}{3!}}T_{mnpk}x^mdx^ndx^pdx^k,`$ $`dS_3`$ $`=`$ $`4T_4,dV_4=4d(\mathrm{ln}r)T_4,d(\mathrm{ln}r)S_3=V_4,`$ $`dT_4`$ $`=`$ $`0d(r^pS_3)=r^p(4T_4+pV_4),`$ (101) where $`r^2=x^mx^m`$. The 11 dimensional Hodge duals can be converted into flat metric Hodge duals on the 7-dimensional space transverse to the D2 branes: $`F_4={\displaystyle \frac{1}{Z}}(_7F_4)dx^0dx^1dx^2`$ $`H_3={\displaystyle \frac{1}{Z^{1/2}}}(_7H_3)dx^0dx^1dx^2`$ (102) For completeness we also give the formulas from , which relate 3-tensor spherical harmonics. $`T_3`$ $`=`$ $`{\displaystyle \frac{1}{3!}}T_{mnp}dx^mdx^ndx^p,S_2={\displaystyle \frac{1}{2}}T_{mnp}x^mdx^ndx^p,`$ $`V_{mnp}`$ $`=`$ $`{\displaystyle \frac{x^q}{r^2}}(x^mT_{qnp}+x^nT_{mqp}+x^pT_{mnq}),V_3={\displaystyle \frac{1}{3!}}V_{mnp}dx^mdx^ndx^p,`$ $`dS_2`$ $`=`$ $`3T_3,d(\mathrm{ln}r)S_2=V_3,dV_3=3d(\mathrm{ln}r)T_3,`$ $`dT_3`$ $`=`$ $`0,d(r^pS_2)=r^p(3T_3+pV_3).`$ (103) where $`r^2=x^mx^m`$.
warning/0005/hep-th0005166.html	ar5iv	text	# 1 Introduction ## 1 Introduction In the last few years, a lot of progress has been made in the study of D-branes as probes of sub-stringy geometry. While the usual picture of space-time is supposed to be valid upto the string scale, it was argued in that D-branes are the natural candidates to probe the geometry of space-time beyond the string scale. As has been widely accepted by now, the standard concepts of space-time appears, from the D-brane perspective, as the vacuum moduli space of its world volume gauge theory. Indeed, it has been found that the space-time appearing in D-brane probes at ultra-short distance scales has qualitatively different features from that probed by bulk strings . Of particular interest has been the study of D-branes on Calabi-Yau manifolds. In ,, the moduli space of Calabi-Yau manifolds was studied using fundamental strings. A rich underlying geometric structure was discovered, and the moduli space was shown to contain topologically distinct geometric (Calabi-Yau) phases, as well as non-geometric (Landau-Ginzburg) phases. The Kähler sector of the vacuum moduli space was shown to consist of several domains (separated by singular Calabi-Yau spaces), whose union is isomorphic to the complex structure moduli space of the mirror manifold. It was also shown, that using mirror symmetry as in or the approach of gauged linear sigma models as in , one can interpolate smoothly between the topologically distinct points in the Kähler moduli space. In the fundamental string picture, the geometric and non-geometric phases of Calabi-Yau manifolds appear as two distinct possibilities, with the latter being thought of as an analytic continuation of the geometric phases in the presence of a non-zero theta angle in the language of . It was, however, shown in that the presence of additional open string sectors change this picture considerably, and that the D-brane linear sigma model probes only a part of the full linear sigma model vacuum moduli space, namely the geometric phase. It was explicitly demonstrated in that D-branes (in the particular examples of the orbifold $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟑}`$ and $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟓}`$) ‘project out’ the non-geometric phase, in agreement with the results obtained in . This was found to be generally true for Calabi-Yau orbifolds of the form $`𝐂^\mathrm{𝟑}/𝐙_𝐧`$ . Since the pioneering work of , which dealt with D-branes on Abelian orbifold singularities of the form $`𝐂^\mathrm{𝟐}/𝐙_𝐧`$ that was generalised to arbitrary non-Abelian singularities of the form $`𝐂^\mathrm{𝟐}/\mathrm{\Gamma }`$ ($`\mathrm{\Gamma }`$ being a subgroup of $`SU(2)`$), D-branes on Abelian and non-Abelian orbifold backgrounds have been extensively studied in ,, ,,,. D-branes at orbifolded conifold singularities have been considered in . (See also for an analysis of a blowup of the four-dimensional $`N=1`$, $`𝐙_\mathrm{𝟑}`$ orientifold). Of particular interest has been the applications of methods of toric geometry , to the study of D-brane gauge theories on such singularities. In the approach pioneered in , the matter content and the interactions of the D-brane gauge theory, which are specified by the D-terms and the F-terms of the gauge theory respectively, are treated on the same footing, and the gauge theory information can be encoded as algebraic equations of the toric variety. It is an interesting question to ask, if this procedure can be reversed, i.e, given a particular toric singularity, is there a way of consistently reading off the world volume gauge theory data for a D-brane that probes this singularity. A step in this direction has been recently taken in . The authors of have given an algorithm by which the geometrical data encoded in a toric diagram can be used to construct the matter content and superpotential interaction of the D-brane gauge theory that probes it, and have demonstrated this in the cases of the suspended pinch point (SPP) singularity of the blowup of the $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$ orbifold and the various partial resolutions of the $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟑}\times 𝐙_\mathrm{𝟑}`$ orbifold analysed in . In the present paper, we critically examine this procedure in several other examples. We first consider this algorithm, which we call the inverse toric procedure, for partial resolutions of the orbifold $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$. Further, we consider several blowups of the orbifolds $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟑}`$ and $`𝐂^\mathrm{𝟒}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$ and present results on the extraction of the D-brane gauge theory data from the given toric geometric data. The organisation of the paper is as follows. In section 2, we first recapitulate the essential details of the construction of the D-brane gauge theory on partial resolutions of the orbifold $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$ , that will also set the notation and conventions to be used in the rest of the paper and then demonstrate the inverse toric algorithm to construct the gauge theory of a D-brane probing the $`𝐙_\mathrm{𝟐}`$ and the conifold singularities obtained by partially resolving the $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$ orbifold. Section 3 deals with the orbifold $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟑}`$ and its various partial resolutions. In section 4, we consider the example of the orbifold $`𝐂^\mathrm{𝟒}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$ and consider some resolutions of the same. Section 5 contains discussions and directions for future work. ## 2 D-branes on $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$ and its blowups We begin this section by recapitulating the essential details of the application of toric methods to the construction of D3-branes on $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$ and some of its partial resolutions which will also set the notations and conventions to be followed in the rest of the paper. Recall that the world volume supersymmetric gauge theory of D-branes probing a toric singularity are characterised by its matter content and interactions. While the former are given by the D-term equations, the latter are specified, via the superpotential, as the F-term equations. These two sets of equations, in conjunction, define the moduli space of the theory. In order to define the D-brane moduli space by toric methods, the standard procedure is to re-express the F and D-term constraints in terms of a set of homogeneous variables $`p_\alpha `$, and a concatenation of these two sets of equations give rise to the toric data. Let us briefly review this construction for the example of D-3 branes on the orbifold $`𝐂^\mathrm{𝟑}/(𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐})`$. We start by considering the field theory of a set of D-3 branes at a $`𝐂^\mathrm{𝟑}/(𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐})`$ singularity. Generically, a single D-p brane of type II string theory at a point in the orbifold $`𝐂^\mathrm{𝟑}/\mathrm{\Gamma }`$ is constructed by considering the theory of $`\|\mathrm{\Gamma }\|`$ D-p branes in $`𝐂^\mathrm{𝟑}`$ and then projecting to $`𝐂^\mathrm{𝟑}/\mathrm{\Gamma }`$ where the projection is defined by a combined action of $`\mathrm{\Gamma }`$ on $`𝐂^\mathrm{𝟑}`$ as well as the D-brane Chan-Paton index. The unbroken supersymmetry in $`d=4`$ would then be $`N=2`$ in the closed string sector while the open string sector will have $`N=1`$ supersymmetry. In our example, a generator $`g`$ of the discrete group $`𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$ acts simultaneously on the coordinates of $`𝐂^\mathrm{𝟑}`$ and also the D-3 brane Chan-Paton factors. The surviving fields in the theory are the ones that are left invariant by a combination of these two actions. Specifically, denoting the action of the quotienting group on $`𝐂^\mathrm{𝟑}`$ by $`R(g)`$ and that on the Chan-Paton indices by $`\gamma (g)`$ (where $`\gamma `$ denotes a regular representation of $`g`$), the surviving components of the scalars $`X`$ that live on the brane world volume are those for which $`X=R(g)\gamma (g)X\gamma (g)^1`$, and the components of the gauge fields $`A`$ that survive the projection are those that satisfy $`A=\gamma (g)A\gamma (g)^1`$. After imposing these projections, the $`N=1`$ matter multiplet of the theory consists of the twelve fields $$(X_{13},X_{31},X_{24},X_{42}),(Y_{12},Y_{21},Y_{34},Y_{43}),(Z_{14},Z_{41},Z_{23},Z_{32})$$ (1) Where we have denoted the complex bosonic fields corresponding to the position coordinates tangential to $`𝐂^\mathrm{𝟑}`$ by $`X^1=X,X^2=Y,X^3=Z`$. After projection of the gauge field, the unbroken gauge symmetry in this case is the subgroup $`U(1)^4`$ of $`U(4)`$, and, apart from an overall $`U(1)`$, we have a $`U(1)^3`$ theory. The total charge matrix is given by $`d=\left(\begin{array}{cccccccccccc}1& 0& 0& 1& 1& 1& 0& 0& 1& 0& 1& 0\\ 0& 1& 1& 0& 1& 1& 0& 0& 0& 1& 0& 1\\ 0& 1& 1& 0& 0& 0& 1& 1& 1& 0& 1& 0\\ 1& 0& 0& 1& 0& 0& 1& 1& 0& 1& 0& 1\end{array}\right)`$ (2) With the last row of $`d`$ signifying the overall $`U(1)`$. The D-term constraints are now given by $$\underset{\mu }{}d_{(i)\mu }^\beta \|X_\beta ^\mu \|^2\zeta _i=0$$ (3) where $`d_{(i)\mu }^\beta `$ is the charge of $`X_\beta ^\mu `$ under the $`i`$’th $`U(1)`$ and $`\beta `$ signifies the matrix indices of the surviving components of $`X^\mu `$. $`_i\zeta _i=0`$ is the condition for the existence of supersymmetric vacua. The superpotential of the theory is given by the expression $$W=\text{Tr}[X^1,X^2]X^3$$ (4) With the vacuum satisfying $$\frac{W}{X^\mu }=0;\text{i.e}[X^\mu ,X^\nu ]=0$$ (5) Using the expressions for the $`X^\mu `$’s from eq.(1) in eq. (5), one can derive twelve constraints on the twelve fields $`X_{ij},Y_{ij},Z_{ij}`$, out of which six are seen to be independent, hence the twelve initial fields can be solved in terms of six independent fields. Denoting the twelve matter fields collectively by $`X_i,i=1\mathrm{}12`$, and the set of six independent fields (which we take to be $`X_{13},X_{24},Y_{21},Y_{34},Z_{14},Z_{32}`$) by $`x_b,(b=1\mathrm{}6)`$, the solution, which is of the form $$X_i=\underset{b=1}{\overset{6}{}}x_b^{K_{ib}}$$ (6) can be encoded in the columns of the matrix $`K`$, given by $`K=\left(\begin{array}{ccccccccccccc}& X_{13}& X_{24}& X_{31}& X_{42}& Y_{12}& Y_{21}& Y_{34}& Y_{43}& Z_{14}& Z_{23}& Z_{32}& Z_{41}\\ X_{13}& 1& 0& 0& 1& 1& 0& 0& 1& 0& 0& 0& 0\\ X_{24}& 0& 1& 1& 0& 1& 0& 0& 1& 0& 0& 0& 0\\ Y_{21}& 0& 0& 0& 0& 0& 1& 0& 1& 0& 1& 0& 1\\ Y_{34}& 0& 0& 0& 0& 1& 0& 1& 0& 0& 1& 0& 1\\ Z_{14}& 0& 0& 1& 1& 0& 0& 0& 0& 1& 1& 0& 0\\ Z_{32}& 0& 0& 1& 1& 0& 0& 0& 0& 0& 0& 1& 1\end{array}\right)`$ (7) The columns of the above matrix are vectors in the lattice $`N=𝐙^\mathrm{𝟔}`$, and define the edges of a cone $`\sigma `$. In order to make the description of the toric variety explicit, we consider the dual cone $`\sigma ^{}`$, which is defined in the dual lattice $`M=\text{Hom}(N,𝐙)`$. The dual cone $`\sigma ^{}`$ is defined to be the set of vectors in the dual lattice $`M`$ which have non-negative inner products with the vectors of $`\sigma `$, i.e $$\sigma ^{}=\{𝐦M:𝐦,𝐧0𝐧\sigma \}$$ (8) The dual cone corresponds to defining a new set of fields $`p_\alpha (\alpha =1\mathrm{}c)`$, in terms of which the variables $`x_b`$ are solved, in a way that is consistent with the relation (6). In this example, $`c=9`$ and the dual cone is generated by the columns of the matrix $`T=\left(\begin{array}{cccccccccc}& p_1& p_2& p_3& p_4& p_5& p_6& p_7& p_8& p_9\\ X_{13}& 1& 1& 0& 0& 0& 0& 0& 0& 1\\ X_{24}& 0& 1& 1& 0& 0& 0& 0& 0& 1\\ Y_{21}& 0& 0& 1& 1& 1& 0& 0& 0& 0\\ Y_{34}& 0& 0& 1& 0& 1& 1& 0& 0& 0\\ Z_{14}& 0& 0& 0& 0& 0& 1& 1& 0& 1\\ Z_{32}& 0& 0& 0& 0& 0& 1& 1& 1& 0\end{array}\right)`$ (9) The relationship between $`p_\alpha `$ and $`x_b`$ are then of the form $$x_b=\underset{\alpha }{}p_\alpha ^{T_{b\alpha }}$$ (10) Thus, there are nine variables $`p_\alpha `$ in terms of which we parametrize the six independent variables $`x_b`$. This parametrization, however, has redundancies, and in order to eliminate these, we introduce a new set of $`𝐂^{}`$ actions on the $`p_\alpha `$’s, with the condition that under these actions, the fields $`x_b`$ are left invariant. This can be equivalently described by the introduction of a new set of $`U(1)`$ gauge symmetries, such that the $`x_b`$’s are gauge invariant . In our example, since there are nine fields $`p_\alpha `$ and six independent variables $`x_b`$, we introduce a gauge group $`U(1)^3`$, such that the $`p_\alpha `$’s are charged with respect to these, with charges $`Q_{n\alpha }`$ for the $`n`$th $`U(1)`$. The gauge invariance condition demands that the charge matrix $`Q`$ must obey $$TQ^t=0$$ (11) where $`Q^t`$ is the transpose of $`Q`$. From the matrix $`T`$ given above, the charge matrix is given by $`Q=\left(\begin{array}{cccccccccc}& p_1& p_2& p_3& p_4& p_5& p_6& p_7& p_8& p_9\\ U\left(1\right)_1& 0& 0& 0& 1& 1& 1& 1& 0& 0\\ U\left(1\right)_2& 0& 1& 0& 0& 0& 0& 1& 1& 1\\ U\left(1\right)_3& 1& 1& 1& 0& 1& 0& 0& 0& 0\end{array}\right)`$ (12) We now consider the constraints imposed on the original fields $`X_i`$ via the D-flatness conditions, and impose these conditions in terms of the new fields $`p_\alpha `$. The charges of the fields $`x_b`$ under the original $`U(1)`$ charge assignments are given by the matrix $`V=\left(\begin{array}{cccccc}x_1& x_2& x_3& x_4& x_5& x_6\\ 1& 0& 1& 0& 1& 0\\ 0& 1& 1& 0& 0& 1\\ 1& 0& 0& 1& 0& 1\end{array}\right)`$ (13) From the relationship $`x_b=_\alpha p_\alpha ^{T_{b\alpha }}`$, it is clear that in order for the charge assignment $`q_{i\alpha }^p`$ for $`p_\alpha `$ (corresponding to the original set of $`(i=3)`$ $`U(1)`$ symmetries) to reproduce that for the $`x_b`$’s, we must have $`q^pT^t=V`$, and one possible choice for the matrix $`q^p`$ is $`q_{i\alpha }^p=(VU)_{i\alpha }`$ with the matrix $`U`$ satisfying $`TU^t=\text{Id}`$. The matrix $`U`$ is thus the inverse of the matrix $`T^t`$. However, since $`T^t`$ is a rectangular matrix the usual definition of its inverse (as in the case of a square matrix) is not applicable and there are several choices for $`U`$. In particular, for a rectangular matrix $`A`$, using the rules of single value decomposition, we can compute the Moore-Penrose inverse of $`A`$ which is defined in a way that the sum of the squares of all the entries in the matrix $`(AA^1\text{Id})`$ is minimised, where Id is the identity matrix in appropriate dimensions. For any choice of the inverse of $`T^t`$, the final toric data will of course be identical, and in this section and the next, we simply use the results of in defining the matrix $`U`$. For this case, the matrix $`U`$ is $`U=\left(T^t\right)^1=\left(\begin{array}{ccccccccc}1& 0& 0& 0& 0& 0& 0& 0& 0\\ 1& 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 1& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 0& 1& 0\end{array}\right)`$ Now, multiplying by $`V`$ and concatenating $`Q`$ and $`VU`$ we obtain the total charge matrix, $`Q_t=\left(\begin{array}{cccccccccc}0& 0& 0& 1& 1& 1& 1& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 1& 1& 1& 0\\ 1& 1& 1& 0& 1& 0& 0& 0& 0& 0\\ 1& 1& 0& 1& 0& 0& 0& 0& 1& \zeta _1\\ 1& 1& 0& 1& 0& 0& 0& 1& 0& \zeta _2\\ 1& 0& 0& 0& 0& 1& 1& 1& 0& \zeta _3\end{array}\right)`$ (14) where the $`\zeta _i`$ are the Fayet-Illiapoulos parameters for the original $`U(1)`$’s, there being three independent $`\zeta _i`$’s because of the constraint $`_{i=1}^4\zeta _i=0`$. The toric data can now be calculated from the cokernel of the transpose of $`Q_t`$, which (after a few row operations) become $`\stackrel{~}{T}=\left(\begin{array}{ccccccccc}0& 1& 0& 0& 1& 0& 1& 1& 1\\ 1& 1& 1& 0& 1& 0& 1& 0& 1\\ 1& 1& 1& 1& 1& 1& 1& 1& 1\end{array}\right)`$ (15) Now from the dual cone $`T`$, we can, via the matrix $`K`$, find expressions for all the twelve initial fields in terms of the fields $`p_\alpha `$. These are then used to define invariant variables in the following way $`X_{13}X_{31}=x`$ $`=`$ $`p_1p_2^2p_3p_8p_9`$ $`Y_{12}Y_{21}=y`$ $`=`$ $`p_1p_3p_4p_5^2p_6`$ $`Z_{14}Z_{41}=z`$ $`=`$ $`p_4p_6p_7^2p_8p_9`$ $`X_{13}Y_{34}Z_{41}=w`$ $`=`$ $`p_1p_2p_3p_4p_5p_6p_7p_8p_9`$ (16) which follow the relation $`xyz=w^2`$. The toric diagram of the orbifold $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$ is shown in figure 1(a). Let us now consider some examples of the possible partial resolutions of this space. We will discuss three possible partial blowups as in ,. The first example corresponds to blowing up the orbifold $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$ by a $`𝐏_1`$ parametrised by $`w^{}=\frac{w}{z}`$. This gives the suspended pinch point (SPP) singularity given by $`xy=zw^2`$. The SPP can be further blown up by introducing a second $`𝐏_1`$. One possibility is to parametrise this $`𝐏_1`$ by $`y^{}=\frac{y}{w^{}}`$ (or $`x^{}=\frac{x}{w^{}}`$) and the remaining singularity is the conifold $`x^{}y=zw^{}`$ (or $`xy^{}=zw^{}`$) of figure 1(b). One can also introduce a $`𝐏_1`$ parametrised by say $`x^{}=\frac{x}{z}`$ whence the resulting blowup is a $`𝐂^\mathrm{𝟐}/𝐙_\mathrm{𝟐}\times 𝐂`$. Two possible toric diagrams of the latter are shown in fig. 1(c) and fig. 1(d). Let us now come to the description of the inverse toric problem. Suppose we are given with a singularity depicted by a toric diagram of the form given in fig. 1(b), 1(c) or 1(d). We wish to extract the gauge theory data of a D-brane that probes such a singularity, starting from figure 1(a). This issue was addressed in , and let us review the basic steps in their algorithm. Recall that the toric data for a given singularity that a D-brane probes is given by the matrix $`\stackrel{~}{T}`$ which is the transpose of the kernel of the total charge matrix (obtained by concatenating the F and D-term constraints expressed in terms of the homogeneous coordinates $`p_\alpha `$). Hence, the solution to the inverse problem would imply the construction of an appropriately reduced total charge matrix $`Q_t^{red}`$, the cokernel of the transpose of which is the reduced matrix $`\stackrel{~}{T}_{red}`$ obtained by deleting the appropriate columns (corresponding to the fields that have been resolved) from $`\stackrel{~}{T}`$. However, although $`\stackrel{~}{T}_{red}`$ is known, there is no unique way of specifying the reduced charge matrix $`Q_t^{red}`$ from a knowledge of $`\stackrel{~}{T}_{red}`$. Further, there is no guarantee apriori that such a $`Q_t^{red}`$, if obtained, would continue to describe the world volume theory of a D-brane that probes the given singularity. The algorithm of gives a canonical method by which these issues can be addressed. It consists of determining the fields to be resolved by tuning the Fayet-Illiopoulos parameters appropriately (so that the resulting theory at the end is still a physical D-brane gauge theory) and then reading off the the matrix $`Q_t^{red}`$ which can be separated into the F and D-terms and thus can be used to obtain the matter content and the superpotential determining the gauge theory. Let us discuss this in some more details. The toric diagrams for the partially resolved singularities are obtained by the deletion of of a subset of nodes from the original one. Conversely, a given toric diagram of dimension $`k`$ can be embedded into a singularity of the form $`𝐂^k/\mathrm{\Gamma }(k,n)`$ where $`\mathrm{\Gamma }(k,n)=𝐙_𝐧\times 𝐙_𝐧\times \mathrm{}(k1\text{times})\mathrm{}𝐙_𝐧`$. In our case, the singularities mentioned above can be embedded into the toric diagram for the singularity $`𝐂^\mathrm{𝟑}/\mathrm{\Gamma }(3,2)`$. For example, the SPP, and the conifold singularity can both be embedded into the $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$ singularity, as is obvious from their toric diagrams. We shall henceforth refer to these singularities, in which we embed the partially resolved ones as the parent singularities. A toric diagram corresponding to a partial resolution can of course be embedded in more than one parent singularity and according to the algorithm of , one can choose the minimal embedding. We will return to this issue in a while. We now determine the fields to be resolved from the parent singularity in order to reach the given toric diagram, by appropriately choosing the values of $`\zeta _i`$. The difficulty that arises here is that given the toric diagram of the singularity that we are interested in, there is no unique way of knowing exactly which fields (in the toric data of the parent singularity) have to be resolved in order to reach the toric singularity of interest. Given a toric diagram for the parent singularity, a particular resolution, along with the elimination of one or more nodes, might also require a subset of fields from some other nodes to be resolved. Hence, we have to carefully tune the FI parameters for a consistent blowup of the parent singularity to the one that we are interested in. We now illustrate this procedure by the example of the blowup of the $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$ singularity to the $`𝐂^\mathrm{𝟐}/𝐙_\mathrm{𝟐}\times 𝐂`$ orbifold illustrated in fig. 1(c). This example is, in a sense, straightforward, because from the toric diagram of fig. 1(c), it is clear that this being the $`𝐙_\mathrm{𝟐}`$ singularity, we can directly read off the matter content and the interactions from the standard results for D3-branes on the $`𝐙_\mathrm{𝟐}`$ orbifold . However, note that this is the simplest example where the inverse procedure of calculating the D-brane gauge theory has to be valid. The $`𝐙_\mathrm{𝟐}`$ orbifold can be thought of as partial resolutions of the orbifolds $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$, or $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟑}`$, and the inverse toric procedure must give familiar results for both in agreement with ,. A lower dimensional toric orbifold can always be embedded in the toric diagram of one of higher dimension, and the inverse toric procedure must consistently reproduce the D-brane gauge theory on the former, starting from the latter data. As we will see, this is indeed the case. We proceed by embedding the $`𝐙_\mathrm{𝟐}`$ singularity in $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$. In order to determine the fields to be resolved, we perform Gaussian row reduction on the the matrix $`Q_t`$ to obtain $`\left(\begin{array}{cccccccccc}1& 0& 0& 0& 0& 1& 1& 0& 1& \zeta _1+\zeta _2+\zeta _3\\ 0& 1& 0& 0& 0& 0& 1& 0& 2& \zeta _1+\zeta _2\\ 0& 0& 1& 0& 0& 1& 1& 0& 1& \zeta _1\\ 0& 0& 0& 1& 0& 1& 0& 0& 0& \zeta _1+\zeta _3\\ 0& 0& 0& 0& 1& 2& 1& 0& 0& \zeta _1+\zeta _3\\ 0& 0& 0& 0& 0& 0& 0& 1& 1& \zeta _1+\zeta _2\end{array}\right)`$ The above set of equations imply the constraints $`x_1+x_6x_7+x_9`$ $`=`$ $`(\zeta _1+\zeta _2+\zeta _3)`$ $`x_2x_7+2x_9`$ $`=`$ $`(\zeta _1+\zeta _2)`$ $`x_3+x_6x_7+x_9`$ $`=`$ $`\zeta _1`$ $`x_4+x_6`$ $`=`$ $`(\zeta _1+\zeta _3)`$ $`x_5+2x_6x_7`$ $`=`$ $`(\zeta _1+\zeta _3)`$ $`x_8+x_9`$ $`=`$ $`(\zeta _1+\zeta _2)`$ (17) Where, for convenience of notation, we have labelled $`p_i^2=x_i`$. This set of equations is now solved in terms of the fields that we know would definitely get resolved. From the toric diagram, three such fields are $`p_7,p_8,p_9`$. However, from the relation between $`x_8`$ and $`x_9`$ above, it suffices to solve the above set of equations in terms of $`x_7`$ and $`x_9`$, say. This gives the solution set, $`[x_1,\mathrm{}x_9]`$ $`=`$ $`[x_1,(2x_9x_7\zeta _1\zeta _2),(x_1+\zeta _2+\zeta _3),(x_1+x_7x_9+\zeta _2),`$ (18) $`(2x_1+x_72x_9+\zeta _1+2\zeta _2+\zeta _3)(x_1+x_7x_9+\zeta _1+\zeta _2+\zeta _3),`$ $`x_7,(x_9\zeta _1\zeta _2),x_9]`$ Now, from the toric diagram of fig.1(c), we choose the fields $`[p_1,p_2,p_3,p_5]`$ to have zero vev i.e these fields continue to remain dynamical. Thus, in the above solution set, we impose $`x_1,x_2,x_3,x_5=0`$ to obtain $$[x_1,\mathrm{}x_9]=[0,0,0,(x_9\zeta _1),0,(x_7x_9+\zeta _1),x_7,(x_9\zeta _1\zeta _2),x_9]$$ (19) We further set $`x_9=\zeta _1`$ in order to make $`x_4=0`$. Hence, $`x_6=x_7=\zeta _1\zeta _2`$. Also, $`x_8=\zeta _2`$. Therefore, the fields to be resolved are $`[p_6,p_7,p_8,p_9]`$, all of which have positive vevs. We now need to obtain the reduced charge matrix for this choice of variables to be resolved. The method of doing this essentially consists of performing row operations on the full charge matrix $`Q_t`$ in eq. (14), such that the columns corresponding to the resolved fields $`[p_6,p_7,p_8,p_9]`$ are set to zero. The reduced charge matrix $`Q_t^{red}`$ thus obtained must be in the nullspace of the reduced toric matrix $`\stackrel{~}{T}_{red}`$, that can be directly evaluated from the full toric matrix in eq. (15) by removing the columns $`6,7,8,9`$. Let us first discuss the general formalism for achieving this, which can then be easily applied to our present example. Consider a general parent theory that has a total charge matrix $`Q_t`$ of dimensions $`(a\times b)`$, and suppose we wish to eliminate $`d`$ fields from the parent theory in order to reach the theory of interest. This implies that we have to perform row operations on $`Q_t`$ so that the given $`d`$ columns will have zero entries and thus be eliminated. The latter can be achieved by constructing the nullspace of the transpose of a submatrix that is formed by precisely the $`d`$ columns that need to be eliminated. Further, since the toric data $`\stackrel{~}{T}`$ was initially in the nullspace of $`Q_t`$, removal of the $`d`$ columns from $`Q_t`$ to obtain $`Q_t^{red}`$ would mean that the reduced toric data $`\stackrel{~}{T}_{red}`$ with the same $`d`$ columns removed would be in the nullspace of $`Q_t^{red}`$. Thus, the expression for $`Q_t^{red}`$ is given by $$Q_t^{red}=\left[\text{NullSpace}(Q_d^t)Q_r\right]$$ (20) where $`Q_d^t`$ is the transpose of the matrix obtained from $`Q`$ containing as its columns, the deleted columns of $`Q`$ and $`Q_r`$ is the remaining matrix. In our example, we evaluate $`Q_{red}`$ directly, by writing $`Q_d^t`$ as the transpose of the matrix that contains columns $`6,7,8,9`$ of the charge matrix $`Q_t`$, and $`Q_r`$ being the matrix containing the columns $`1,2,3,4,5,6,10`$ of $`Q_t`$. This implies that the reduced charge matrix is $`Q_t^{red}=\left(\begin{array}{cccccc}1& 1& 1& 0& 1& 0\\ 2& 1& 0& 0& 1& \zeta _2+\zeta _3\end{array}\right)`$ (21) From this, the reduced charge matrices corresponding to the F and D-terms are $$Q_{red}=\left(11101\right);(VU)_{red}=\left(21001\right)$$ (22) Denoting the kernel of the charge matrix by $`T_{red}`$, the transpose of its dual is the reduced $`K^t`$ matrix, $`K_{red}^t=\left(\begin{array}{ccccc}0& 1& 0& 1& 0\\ 1& 0& 0& 0& 0\\ 0& 1& 1& 0& 0\\ 0& 0& 1& 0& 1\end{array}\right)`$ Now, From eq. (6), the charge matrix for the initial fields $`X_i`$ are given, via the charge matrix for the independent fields $`x_b`$ by the relation $$\mathrm{\Delta }=VK^t$$ (23) Hence, in this case, making use of the identity $`U_{red}T_{red}^t=\text{Id}`$, the matter content for the D-brane gauge theory on the resolved space is obtained as $`\mathrm{\Delta }_{red}=(VU)_{red}(T_{red}^tK_{red}^t)`$. From which in this case we obtain the matrix $`d_{red}=\left(\begin{array}{cccccc}0& 1& 1& 1& 1& \\ 0& 1& 1& 1& 1\end{array}\right)`$ (24) where we have included an extra row corresponding to the redundant $`U(1)`$. Hence, the gauge group is seen to be $`U(1)^2`$, with four bifundamentals, in agreement with the result of . Also note that one adjoint field has appeared in our expression. This signals an ambiguity that exists in the toric procedure with regards to handling chargeless fields. Consider, for example, the superpotential interaction, which can be determined from the matrix $`K`$ that encodes the F-flatness conditions. The number of columns in the matrix $`K_{red}^t`$ is the number of fields that appear in the interactions and the terms in the superpotential are read off as linear relations between the columns of $`K_{red}^t`$. In this example, from the matrix $`K`$, we have the relation $`X_2X_5=X_3X_4`$, and from the $`d`$ matrix, we can see that each of these have charge zero. Hence, we try as an ansatz for the superpotential, $$W=\left(X_1+\varphi \right)\left(X_2X_5X_3X_4\right)$$ (25) where $`\varphi `$ is a field uncharged for both the $`U(1)`$’s. This agrees with after the identification $`x_{13}X_2;Y_{34}X_5;X_{24}X_3;X_{13}X_4;Z_{41}X_1;Z_{23}\varphi `$. As is clear, however, there is an ambiguity in writing the superpotential due to the presence of fields that are chargeless under both the $`U(1)`$’s . Let us also mention that we can carry out the inverse procedure just outlined for the case of the alternative blowup to $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐂`$ shown in figure 1(d). This gives exactly the same matter content and superpotential as the case just analysed, in the region of moduli space given by $`\zeta _1+\zeta _3=0`$, $`\zeta _2>>0`$. $``$ We now discuss the example of the blowup of the $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$ singularity to the conifold. Proceeding in the same way as above, we now find that in the region of the FI parameters determined by $`\zeta _2=0;\zeta _{1,3}>>0`$, the fields to be resolved are $`[p_3,p_5,p_6,p_7,p_9]`$, when the fields $`[p_1,p_2,p_4,p_8]`$ are chosen dynamical, i.e continue to have zero vev. Therefore, from eq. (20), we obtain the reduced charge matrix to be $$Q_t^{red}=\left(1111\zeta _2\right)$$ (26) Note that in this, case, there are no F-terms and we have $`Q_{red}=\mathrm{𝟎}_{1\times 4}`$, moreover, the kernel of this matrix is the $`4\times 4`$ identity matrix, i.e $`T_{red}=\text{Id}_{4\times 4}`$, which implies that (since the dual matrix of $`T_{red}`$ is also the $`4\times 4`$ identity matrix) $`\mathrm{\Delta }=Q_{red}`$, with the charge matrix given in this case (by adding the extra column corresponding to the overall $`U(1)`$) as $`d_{red}=\left(\begin{array}{cccc}1& 1& 1& 1\\ 1& 1& 1& 1\end{array}\right)`$ (27) In this case, the matrix $`Q_{red}`$ being zero, there are no F-terms and the entire gauge theory information is contained in the D-term. ## 3 D-branes on $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟑}`$ and its blowups Let us now consider D3-branes on the orbifold $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟑}`$. In this case, the bosonic fields corresponding to the D3-brane coordinates tangential to the $`𝐂^\mathrm{𝟑}`$ are $`6\times 6`$ matrices, and the components that survive the projection by the quotienting group constitute the $`N=1`$ matter multiplet, given by the set $`(X_{12},X_{23},X_{34},X_{45},X_{56},X_{61})`$ $`(Y_{13},Y_{24},Y_{35},Y_{46},Y_{51},Y_{62})`$ $`(Z_{14},Z_{25},Z_{36},Z_{41},Z_{52},Z_{63})`$ (28) The F-flatness condition (5) in this case imply $`18`$ constraints out of which $`8`$ are seen to be independent. Further, there are $`17`$ homogeneous coordinates $`p_\alpha ,(\alpha =1\mathrm{}17)`$, corresponding to the physical fields in terms of which the dual cone may be described. The total charge matrix of the homogeneous coordinates $`p_\alpha `$ is defined by introducing a set of $`9`$ $`𝐂^{}`$ actions that remove the redundancy in expressing the original independent fields in terms of the homogeneous coordinates and a further set of $`U(1)^5`$ charges expressing the D-term constraints in terms of the $`p_\alpha `$. (The $`6`$th $`U(1)`$ is redundant because of the relation $`_i\zeta _i=0`$). The total charge matrix can then be obtained by concatenating these two sets of charge matrices as in our earlier example, and is given (with the inclusion of a column specifying the $`5`$ independent Fayet-Illiopoulos parameters) by $`Q_t=\left(\begin{array}{cccccccccccccccccc}1& 0& 1& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& 0& 1& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 1& 0& 0& 0& 0& 0& 0& 1& 0& 0& 1& 0\\ 0& 0& 1& 1& 0& 0& 0& 0& 0& 1& 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 1& 0& 0& 0& 0& 0& 0& 0& 1& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0& 1& 0& 1& 0& 0& 0& 1& 1& 1& 0\\ 0& 0& 0& 0& 0& 1& 0& 0& 0& 1& 0& 0& 1& 0& 1& 1& 1& 0\\ 0& 0& 0& 0& 0& 0& 1& 0& 0& 1& 1& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1& 1& 0& 0& 1& 1& 0& 1& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& 0& 1& 1& 1& 0& 0\\ 0& 1& 0& 0& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& \zeta _1\\ 1& 1& 0& 0& 0& 1& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& \zeta _2\\ 1& 0& 1& 0& 1& 0& 0& 0& 1& 0& 0& 0& 0& 0& 0& 0& 0& \zeta _3\\ 0& 1& 0& 1& 1& 0& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& \zeta _4\\ 1& 1& 1& 0& 0& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& \zeta _5\end{array}\right)`$ The toric data is the co-kernel of the transpose of this charge matrix (when the FI parameters are set to zero) and is given by $`\stackrel{~}{T}=\left(\begin{array}{ccccccccccccccccc}1& 0& 1& 1& 0& 0& 1& 1& 0& 1& 1& 0& 0& 1& 1& 0& 0\\ 0& 0& 1& 1& 0& 1& 1& 0& 1& 0& 2& 0& 0& 1& 1& 0& 0\\ 1& 1& 1& 1& 1& 1& 1& 1& 1& 1& 1& 1& 1& 1& 1& 1& 1\end{array}\right)`$ (29) The invariant variables are defined by $`X_{12}X_{23}X_{34}X_{45}X_{56}X_{61}=x`$ $`=`$ $`p_1^4p_2p_3^6p_4^2p_5p_6^3p_7^2p_8^4p_9^3p_{10}^4p_{12}p_{13}p_{15}^2p_{16}p_{17}`$ $`Y_{24}Y_{46}Y_{62}=y`$ $`=`$ $`p_1p_2p_4^2p_5p_7^2p_8p_{10}p_{11}^3p_{12}p_{13}p_{15}^2p_{16}p_{17}`$ $`Z_{14}Z_{41}=z`$ $`=`$ $`p_2p_5p_6p_9p_{12}p_{13}p_{14}^2p_{16}p_{17}`$ $`X_{12}Y_{24}Z_{41}=w`$ $`=`$ $`p_1p_2p_3p_4p_5p_6p_7p_8p_9p_{10}p_{11}p_{12}p_{13}`$ (30) $`p_{14}p_{15}p_{16}p_{17}`$ From which it can be seen that $$xy^2z^3=w^6$$ (31) Let us now discuss some examples of blowups of this space that will illustrate the procedure outlined in section 2. We will use the notation $`x_i=p_i^2`$ in what follows. $``$ Our first example is illustrated in figure 2. From the diagram, we see that the field $`p_{14}`$ will definitely be resolved. Hence, after performing Gaussian row reduction on the charge matrix $`Q_t`$, which gives $`\left(\begin{array}{cccccccccccccccccc}1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& 0& 2& 0& 0& \zeta _12\zeta _2\zeta _42\zeta _5\\ 0& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& \zeta _1\zeta _2\zeta _3\\ 0& 0& 1& 0& 0& 0& 0& 0& 0& 0& 2& 0& 0& 0& 3& 0& 0& \zeta _12\zeta _2\zeta _42\zeta _5\\ 0& 0& 0& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& \zeta _2\zeta _5\\ 0& 0& 0& 0& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& \zeta _2\zeta _3\\ 0& 0& 0& 0& 0& 1& 0& 0& 0& 0& 1& 0& 0& 0& 1& 0& 1& \zeta _1\zeta _2\zeta _3\zeta _5\\ 0& 0& 0& 0& 0& 0& 1& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& \zeta _1\zeta _2\zeta _4\zeta _5\\ 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& 1& 0& 0& 0& 2& 0& 0& \zeta _2\zeta _5\\ 0& 0& 0& 0& 0& 0& 0& 0& 1& 0& 1& 0& 0& 0& 1& 0& 1& \zeta _2\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& 1& 0& 0& 0& 2& 0& 0& \zeta _1\zeta _2\zeta _4\zeta _5\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& 0& 0& 1& \zeta _4+\zeta _5\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& 0& 1& \zeta _4\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& 1& 0& 2& \zeta _3+\zeta _4+\zeta _5\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& 1& \zeta _3\end{array}\right)`$ (32) we use this matrix to solve for all the fields in terms of $`x_{14}`$, the solution set being given by $`\left[x_1\mathrm{}x_{17}\right]`$ $`=`$ $`[x_1,x_2,(2x_12x_2+x_{14}\zeta _1\zeta _3+\zeta _5),`$ (33) $`(2x_2x_{14}+2\zeta _1+\zeta _2+\zeta _3+\zeta _4),(x_2+\zeta _1),`$ $`(x_1x_2+x_{14}\zeta _1\zeta _3),(2x_2x_{14}+\zeta _1+\zeta _2+\zeta _3),`$ $`(x_1+\zeta _1+\zeta _2+\zeta _4+\zeta _5),(x_1x_2+x_{14}+\zeta _5),`$ $`(x_1+\zeta _2+\zeta _5),(x_1+4x_22x_{14}+3\zeta _1+2\zeta _2+2\zeta _3+\zeta _4),`$ $`(x_2+\zeta _1+\zeta _2+\zeta _3+\zeta _4+\zeta _5),(x_2+\zeta _1+\zeta _2+\zeta _3+\zeta _4),`$ $`x_{14},(2x_2x_{14}+2\zeta _1+2\zeta _2+\zeta _3+\zeta _4+\zeta _5),`$ $`(x_2+\zeta _1+\zeta _2),(x_2+\zeta _1+\zeta _2+\zeta _3)]`$ Now, we choose the fields $`[p_1,p_3,p_4,p_{11}]`$ as the ones that continue to have zero vev, and in eq. (33) set the values of these to zero. We thus obtain conditions on the FI parameters, $`\zeta _1+\zeta _2+\zeta _4=0,\zeta _2=2\zeta _5`$. In the new solution set obtained with these conditions, we choose $`\zeta _5=0`$ which makes $`x_8=0`$. This also imples that $`\zeta _2=0`$ (since $`\zeta _2=2\zeta _5`$). We also set $`x_2=x_{14}=\zeta _3\zeta _1`$ to obtain $`x_7,x_9,x_{15},x_{17}=0`$. Our choice of $`\zeta _2=\zeta _5=0`$ implies that $`x_{10}=0`$. This gives the solution set $`\left[x_1\mathrm{}x_{17}\right]`$ $`=`$ $`[0,(\zeta _3\zeta _1),0,0,\zeta _3,(\zeta _1\zeta _3),`$ (34) $`0,0,0,0,0,\zeta _1,\zeta _1,(\zeta _3\zeta _1),0,\zeta _3,0]`$ Thus the fields to be resolved are $`[p_2,p_5,p_6,p_{12},p_{13},p_{14},p_{16}]`$. Hence, from (20), we obtain the reduced charge matrix, $`Q_t^{red}=\left(\begin{array}{ccccccccccc}0& 0& 1& 0& 1& 0& 0& 1& 1& 0& 0\\ 0& 0& 0& 1& 0& 0& 1& 1& 1& 0& 0\\ 0& 1& 1& 0& 0& 0& 1& 1& 0& 0& 0\\ 1& 1& 0& 0& 0& 0& 0& 1& 1& 0& 0\\ 1& 0& 1& 0& 0& 1& 1& 1& 0& 1& \zeta _2\\ 0& 0& 1& 1& 0& 0& 0& 0& 0& 0& \zeta _1+\zeta _4\\ 1& 1& 1& 1& 0& 1& 1& 1& 0& 1& \zeta _5\end{array}\right)`$ In the usual way, we calculate the kernel of the charge matrix $`T_{red}`$, and the transpose of its dual matrix is $`K_{red}^t=\left(\begin{array}{cccccccc}0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 1& 0& 1& 0& 0\\ 0& 0& 1& 1& 1& 0& 0& 0\\ 0& 0& 0& 1& 0& 1& 1& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 1& 0& 1\end{array}\right)`$ from the matrices $`T_{red},K_{red}^t`$ and $`(VU)_{red}`$, we can read off the matter content and gauge group of the D-brane gauge theory from the matrix $`d_{red}=\left(\begin{array}{cccccccc}1& 1& 0& 1& 0& 0& 0& 1\\ 0& 0& 1& 1& 0& 0& 1& 1\\ 1& 1& 0& 0& 1& 1& 0& 0\\ 0& 0& 1& 0& 1& 1& 1& 0\end{array}\right)`$ This is an $`U(1)^4`$ gauge theory with $`8`$ matter fields, in agreement with ,. To calculate the superpotential, we first note that the relations between the columns of the matrix $`K^t`$, which can be written as $`X_3X_7=X_5X_6=X_4X_8`$. All these combinations are seen to be chargeless, as is the combination $`X_1X_2`$. Hence, our ansatz for the superpotential in this case is $$W=X_1X_2\left(X_3X_7X_5X_6\right)+\varphi _1\left(X_3X_7X_4X_8\right)+\varphi _2\left(X_5X_6X_4X_8\right)$$ (35) which agrees with that calculated in after the identification with their notation, $`Z_{25}X_1,Z_{52}X_2;X_{23}X_3;Y_{62}X_7;Y_{51}X_5;X_{45}X_6;YX_4;X_{61}X_8;Z_{36}\varphi _1;Z_{14}\varphi _2`$. In terms of the original variables of (28), this blowup corresponds to giving a vev to the fields $`Z_{63},Z_{41}`$. The invariant variables are in this case defined by $`x^{}`$ $`=`$ $`p_1^2p_3^3p_4p_7p_8^2p_9p_{10}^2p_{15};x=x^2z`$ $`y`$ $`=`$ $`p_1p_4^2p_7^2p_8p_{10}p_{11}^3p_{15}^2p_{17}`$ $`z`$ $`=`$ $`p_9p_{17}`$ $`w^{}`$ $`=`$ $`p_1p_3p_4p_7p_8p_{10}p_{11}p_{15};w=w^{}z`$ (36) In terms of these variables, the blown up space is given by $`x^{}y=zw^3`$. $``$ We now consider the second example which is depicted in figure 3(a). In this case, we start from the matrix (32), and in the solution set of eq. (33), keeping in mind the toric diagram of the blowup, we select the set of fields $`[x_7,x_9]`$ that remain dynamical, with zero vev. This gives the solution $`x_1=\zeta _1+\zeta _4`$. From the resulting solution set, we now select $`x_{17}=0`$, in order to obtain $`x_2=x_{14}=\zeta _1\zeta _2\zeta _3`$ and after imposing this condition, we make the choice $`\zeta _1+\zeta _4+\zeta _5=0`$ in order for $`p_9`$ to have zero vev. The final solution is given by $`\left[x_1\mathrm{}x_{17}\right]`$ $`=`$ $`[(\zeta _1+\zeta _4),(\zeta _1\zeta _2\zeta _3),(\zeta _1+\zeta _2+\zeta _4),(\zeta _1+\zeta _4),`$ (37) $`(\zeta _2\zeta _3),(\zeta _4\zeta _3),0,(\zeta _1+\zeta _2+\zeta _4),0,\zeta _2,0,\zeta _1,\zeta _4,`$ $`(\zeta _1\zeta _2\zeta _3),\zeta _2,\zeta _3,0]`$ Hence, the fields to be resolved are $`[p_1,p_2,p_3,p_4,p_5,p_6,p_8,p_{10},p_{12},p_{13},p_{14},p_{15},p_{16}]`$ and the dynamical fields are $`[p_7,p_9,p_{11},p_{17}]`$. We can thus directly evaluate the reduced charge matrix from (20) which is given by $$Q_t^{red}=\left(1111\zeta _1+\zeta _4+\zeta _5\right)$$ (38) Defining the invariant variables from eq. (28) in this case as $`x^{}`$ $`=`$ $`p_7p_9;x=x^2z`$ $`y^{}`$ $`=`$ $`p_{11}p_{17};y=y^{}w^2`$ $`z`$ $`=`$ $`p_9p_{17}`$ $`w^{}`$ $`=`$ $`p_7p_{11};w=w^{}z`$ (39) this blownup space can be seen to be the conifold $`x^{}y^{}=zw^{}`$ by using the redefined variables in (31). The gauge theory living on the D-brane is of course the same as in eq. (27). In terms of the original variables in eq. (28), this blowup can be shown to correspond to giving vevs to the fields $`X_{45},X_{23},Z_{63},Z_{41}`$. Since there are no F-terms in this case, all the information about the gauge theory is contained in the D-term constraint. $``$ Let us now consider the toric diagram given by figure 3(b). From the toric diagram of fig. 3(b), we choose the fields $`p_1,p_2,p_3,p_4,p_{11}`$ to have zero vev. These conditions imply that $`\zeta _1+\zeta _4=0`$ and also $`\zeta _2+\zeta _5=0`$. Substituting this in the solution set (33), we obtain $`\left[x_1\mathrm{}x_{17}\right]`$ $`=`$ $`[0,0,0,0,\zeta _1,\zeta _2,0,0,(\zeta _1+\zeta _3),0,0,\zeta _3,(\zeta _2+\zeta _3),`$ (40) $`(\zeta _1+\zeta _3\zeta _5),0,(\zeta _1+\zeta _2),(\zeta _1+\zeta _2+\zeta _3)]`$ Now, we have a choice. We put $`\zeta _1=0`$ and obtain $`x_5=0`$ which also implies that $`\zeta _4=0`$. Hence, from the final solution set, we see that the fields that continue to have zero vev and are dynamical are, in this case, given by the set $`[p_1,p_2,p_3,p_4,p_5,p_7,p_8,p_{10},p_{11},p_{15}]`$ and the fields that pick up non-zero vev and hence are resolved are $`[p_6,p_9,p_{12},p_{13},p_{14},p_{16},p_{17}]`$. Hence, according to (20), we obtain the reduced charge matrix as $`Q_t^{red}=\left(\begin{array}{ccccccccccc}0& 0& 0& 1& 0& 0& 1& 0& 1& 1& 0\\ 0& 0& 0& 0& 0& 1& 0& 1& 1& 1& 0\\ 0& 0& 1& 1& 0& 0& 0& 1& 1& 0& 0\\ 1& 0& 1& 0& 0& 0& 0& 0& 1& 1& 0\\ 2& 0& 1& 0& 0& 1& 0& 0& 0& 0& \zeta _2+\zeta _5\\ 0& 1& 0& 1& 1& 1& 0& 0& 0& 0& \zeta _4\\ 0& 1& 0& 0& 1& 0& 0& 0& 0& 0& \zeta _1\end{array}\right)`$ and the dual of the kernel of the charge matrix $`Q_{red}`$ is $`K_{red}^t=\left(\begin{array}{cccccccc}1& 0& 0& 0& 0& 1& 1& 0\\ 1& 0& 0& 0& 0& 1& 0& 1\\ 1& 0& 0& 0& 1& 0& 1& 0\\ 0& 0& 0& 1& 0& 1& 1& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\end{array}\right)`$ As before, in order to calculate the matter content, we use the matrix $`T_{red}^tK_{red}^t`$, from which the matrix $`\mathrm{\Delta }`$, which specifies the matter content and gauge group of the D-brane gauge theory is calculated to be $`d_{red}=\left(\begin{array}{cccccccc}1& 0& 0& 1& 0& 0& 1& 1\\ 1& 1& 1& 1& 1& 1& 0& 0\\ 0& 1& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 1& 1& 1\end{array}\right)`$ Hence, the gauge group is $`U(1)^4`$ and there are $`8`$ matter fields. Using the relations $`X_1X_4=X_5X_6=X_7X_8`$ with total charge zero and the chargeless combination $`X_2X_3`$, we write the superpotential in this case as $$W=X_2X_3\left(X_1X_4X_5X_6\right)+\varphi _1\left(X_1X_4X_7X_8\right)+\varphi _2\left(X_5X_6X_7X_8\right)$$ (41) where the fields $`\varphi _1,\varphi _2`$ are possible adjoints. The invariant variables (30) are defined by $`x^{}`$ $`=`$ $`p_1^2p_3^3p_4p_7p_8^2p_{10}^2p_{15}`$ $`y^{}`$ $`=`$ $`p_1p_4^2p_7^2p_8p_{10}p_{11}^3p_{15}^2`$ $`z`$ $`=`$ $`p_2p_5`$ $`w^{}`$ $`=`$ $`p_1p_3p_4p_7p_8p_{10}p_{11}p_{15}`$ (42) with the relations $`x=x^2z;y=y^{}z;w=w^{}z`$. The remaining singularity is seen from eq. (31) to be the equation $`x^{}y^{}=w^3`$ in $`𝐂^\mathrm{𝟑}`$, which is recognised as corresponding to blowing up the $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟑}`$ orbifold to the singularity $`𝐂^\mathrm{𝟐}/𝐙_\mathrm{𝟑}\times 𝐂`$. Let us also mention that in terms of the original variables in eq. (28), this blowup corresponds to giving vevs to the fields $`Y_{13},Z_{14},Z_{52},Z_{63}`$. $``$ We now come to the example shown in figure 3(c). In this case, from the Gaussian row-reduced matrix of eq. (33), we choose the set of fields $`[p_{10},p_7,p_{11}]`$ to have zero vev. This implies that $`\zeta _1+\zeta _2+\zeta _4+\zeta _5=0`$ and also that $`x_1=\zeta _1+\zeta _4`$. Substituting this in the solution set implies that in addition $`x_{15}=0`$, and further setting $`x_2=\zeta _3+\zeta _4+\zeta _5`$, so that $`x_{17}`$ is zero, we obtain the full solution set in this case as $`\left[x_1\mathrm{}x_{17}\right]`$ $`=`$ $`[(\zeta _2\zeta _5),(\zeta _3+\zeta _4+\zeta _5),(\zeta _2\zeta _5),(\zeta _2\zeta _5),`$ (43) $`(\zeta _2\zeta _3),(\zeta _4\zeta _5),0,(\zeta _2\zeta _5),\zeta _2,0,0,(\zeta _4+\zeta _5),\zeta _4,`$ $`(\zeta _3+\zeta _4+\zeta _5),0,\zeta _3,0]`$ Hence, the fields that are dynamical and have zero vev is given by the set $`[p_7,p_{10},p_{11},p_{15},p_{17}]`$ and the set to be resolved is given by $`[p_1,p_2,p_3,p_4,p_5,p_6,p_8,p_9,p_{12},p_{13},p_{14},p_{16}]`$. Using this, the reduced charge matrix is calculated from (20) to be $`Q_t^{red}=\left(\begin{array}{cccccc}1& 1& 1& 1& 0& 0\\ 0& 1& 1& 2& 0& \zeta _1+\zeta _2+\zeta _4+\zeta _5\end{array}\right)`$ From which the F and the D-terms as $$Q_{red}=\left(11110\right)(VU)_{red}=\left(01120\right)$$ (44) The dual of the kernel of $`Q_{red}`$ is given by $`K_{red}^t=\left(\begin{array}{ccccc}1& 0& 0& 0& 0\\ 0& 1& 1& 0& 0\\ 0& 1& 0& 1& 0\\ 0& 0& 1& 0& 1\end{array}\right)`$ and the matter content and gauge group of the D-brane gauge theory is obtained as $`d_{red}=\left(\begin{array}{ccccc}0& 1& 1& 1& 1\\ 0& 1& 1& 1& 1\end{array}\right)`$ The superpotential is written down, after noting that the field $`X_1`$ is chargeless, and the relations $`X_2X_5=X_3X_4`$ from the matrix $`K^t`$. It is given by $$W=\left(X_1+\varphi \right)\left(X_2X_5X_3X_4\right)$$ (45) where $`\varphi `$ is a possible adjoint. The matter content and the superpotential are exactly the same as in (24) and (25). This is of course as expected because both describe the world volume theory of a D3-brane on the $`𝐂^\mathrm{𝟐}/𝐙_\mathrm{𝟐}`$ singularity. This blowup corresponds to giving vevs to the fields $`X_{12},X_{23},X_{45}`$. The invariant variables (30) are in this case $`x^{}`$ $`=`$ $`p_7p_{10}^2p_{15}`$ $`y^{}`$ $`=`$ $`p_7p_{11}^2p_{15}`$ $`z`$ $`=`$ $`p_{17}`$ $`w^{}`$ $`=`$ $`p_7p_{10}p_{11}p_{15}`$ (46) where $`x=x^2z;y=y^{}w^{}z;w=w^{}z`$. The blown up space is, from eq. (31), given by the equation $`x^{}y^{}=w^2`$ in $`𝐂^\mathrm{𝟑}`$, which describes the blowup to $`𝐂^\mathrm{𝟐}/𝐙_\mathrm{𝟑}\times 𝐂`$. Let us pause to comment here that the toric diagram of figure 3(c) can be embedded both into that of $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟑}`$ and $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$; in fact, this diagram is identical to that in figure 1(d). The fact that either embedding gives correct results for the D-brane gauge theory data shows that the inverse toric procedure is indeed consistent. $``$ Finally, let us discuss the toric diagram of figure 3(d). which is the conifold obtained by partially resolving the $`𝐂^\mathrm{𝟐}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟑}`$ singularity. Following the same procedure as outlined in the previous examples, we start from the solution set of (33) and select the fields $`[p_1,p_3,p_6]`$ which we set to zero. From this, we obtain the solution $`x_2=\zeta _5`$ and $`x_{14}=(\zeta _1+\zeta _2+\zeta _5)`$. Substituting this in the solution set and further setting $`\zeta _5=0`$ so that $`x_2=0`$, we finally obtain the full solution set to be $`\left[x_1\mathrm{}x_{17}\right]`$ $`=`$ $`[0,0,0,(\zeta _1+\zeta _2+\zeta _4),\zeta _1,0,\zeta _2,(\zeta _1+\zeta _2+\zeta _3+\zeta _4),(\zeta _1+\zeta _3),`$ (47) $`\zeta _2,(\zeta _1+2\zeta _2+\zeta _4)(\zeta _1+\zeta _2+\zeta _3+\zeta _4),(\zeta _1+\zeta _2+\zeta _3+\zeta _4),`$ $`(\zeta _1+\zeta _3),(\zeta _1+2\zeta _2+\zeta _4),(\zeta _1+\zeta _2),(\zeta _1+\zeta _2+\zeta _3)]`$ Hence, the fields that are dynamical are given by the set $`[p_1,p_2,p_3,p_6]`$ while the rest acquire non-zero vev and hence are resolved, and in this case the reduced charge matrix is obtained from (20) to be $$Q=\left(1111\zeta _5\right)$$ (48) This blowup corresponds to the conifold. To see this, we define the invariant variables (30) in this case by $$x^{}=p_1p_3=\frac{w}{z};y=p_1p_2;z=p_2p_6;w^{}=p_3p_6$$ (49) In this case, from (31), we can see that the remaining singularity is given by the conifold $`x^{}y=zw^{}`$. Let us also mention that in terms of the original variables in eq. (28), this blowup corresponds to giving vevs to $`Y_{13},Y_{24},Z_{36}`$. ## 4 D-branes on $`𝐂^\mathrm{𝟒}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$ and its blowups We now consider D-branes on the orbifold $`𝐂^\mathrm{𝟒}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$ and some of its blowups . We will consider a D1-brane at this singularity, and the world volume theory is an $`N=(0,2)`$ SYM theory in two dimensions. The method of construction of the moduli space of the D-1 brane is similar to the method outlined in sections 2 and 3. The unbroken gauge symmetry in this case is (apart from a redundant $`U(1)`$), $`U(1)^7`$. The matter multiplet, after imposing the appropriate projections consist of $`32`$ surviving fields, $`(X_{18},X_{27},X_{36},X_{45},X_{54},X_{63},X_{72},X_{81}),(Y_{15},Y_{26},Y_{37},Y_{48},Y_{51},Y_{62},Y_{73},Y_{84})`$ $`(Z_{13},Z_{24},Z_{31},Z_{42},Z_{57},Z_{68},Z_{75},Z_{86}),(W_{12},W_{21},W_{34},W_{43},W_{56},W_{65},W_{78},W_{87})`$ (50) The F-term constraints of eq. (5), in this case imply that out of the $`32`$ fields, $`11`$ are independent, and the $`32`$ initial fields can be solved in terms of these $`11`$, and the solutions can be denoted as vectors in the lattice $`𝐙^{\mathrm{𝟏𝟏}}`$. The dual cone in this case has $`34`$ homogeneous coordinates (matter fields) $`p_\alpha `$, and the redundancy in the definition of the initial independent coordinates in terms of these can be eliminated by introducing a set of $`23`$ $`𝐂^{}`$ actions. The charge matrix $`Q`$ for the $`𝐂^{}`$ actions is given in . For the matrix denoting the original $`U(1)`$ charges expressed in terms of the homogeneous coordinates $`p_\alpha `$, we use the Moore-Penrose inverse of the matrix defining the dual cone, and concatenate these two matrices to obtain the total charge matrix given in the appendix. The toric diagram is given by the columns of the matrix $$\begin{array}{c}(\begin{array}{ccccccccccccccccc}p_1& p_2& p_3& p_4& p_5& p_6& p_7& p_8& p_9& p_{10}& p_{11}& p_{12}& p_{13}& p_{14}& p_{15}& p_{16}& p_{17}\\ 1& 0& 1& 0& 1& 2& 1& 0& 0& 1& 0& 1& 1& 1& 1& 1& 1\\ 1& 1& 1& 1& 0& 0& 0& 0& 0& 0& 2& 1& 0& 1& 1& 1& 1\\ 1& 0& 0& 1& 1& 0& 0& 1& 2& 1& 0& 0& 0& 1& 1& 1& 1\\ 0& 0& 0& 1& 0& 1& 1& 0& 1& 0& 1& 0& 1& 0& 0& 0& 0\end{array}\mathrm{}\mathrm{}\hfill \\ \\ \\ \mathrm{}\begin{array}{ccccccccccccccccc}p_{18}& p_{19}& p_{20}& p_{21}& p_{22}& p_{23}& p_{24}& p_{25}& p_{26}& p_{27}& p_{28}& p_{29}& p_{30}& p_{31}& p_{32}& p_{33}& p_{34}\\ 1& 0& 0& 0& 1& 1& 1& 1& 1& 1& 0& 1& 1& 1& 1& 1& 1\\ 1& 1& 0& 0& 1& 1& 1& 1& 1& 1& 1& 1& 1& 1& 1& 1& 1\\ 1& 1& 0& 1& 1& 1& 1& 1& 1& 1& 0& 1& 1& 1& 1& 1& 1\\ 0& 1& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0\end{array})\hfill \end{array}$$ The distinct columns in the toric data are shown in fig. 4. The labelling of the vectors and the corresponding $`p_\alpha `$ are as follows $`a`$ $`(2001)[p_6]`$ $`b`$ $`(1001)[p_7,p_{13}]`$ $`c`$ $`(1010)[p_5,p_{10}]`$ $`d`$ $`(1110)[p_1,p_{14},p_{15},p_{16},p_{17},p_{18},p_{22},p_{23},`$ $`p_{24},p_{25},p_{26},p_{27},p_{29},p_{30},p_{31},p_{32},p_{33},p_{34}]`$ $`e`$ $`(1100)[p_3,p_{12}]`$ $`f`$ $`(0001)[p_{20}]`$ $`g`$ $`(0010)[p_8,p_{21}]`$ $`h`$ $`(0021)[p_9]`$ $`i`$ $`(0111)[p_4,p_{19}]`$ $`j`$ $`(0201)[p_{11}]`$ $`k`$ $`(0100)[p_{11}]`$ (51) The invariant variables, in this case denoted by $`x=(X_{18}X_{81})`$; $`y=(Y_{15}Y_{51})`$; $`z=(Z_{13}Z_{31})`$; $`w=(W_{12}W_{21})`$; $`v=(X_{18}Y_{84}Z_{42}W_{21})`$, are $`x`$ $`=`$ $`p_1p_2p_7p_8p_{13}p_{14}p_{15}p_{16}p_{17}p_{18}p_{20}^2p_{21}p_{22}p_{23}p_{24}`$ $`p_{25}p_{26}p_{27}p_{28}p_{29}p_{30}p_{31}p_{32}p_{33}p_{34}`$ $`y`$ $`=`$ $`p_1p_2p_3p_4p_{11}^2p_{12}p_{14}p_{15}p_{16}p_{17}p_{18}p_{19}p_{22}p_{23}p_{24}`$ $`p_{25}p_{26}p_{27}p_{28}p_{29}p_{30}p_{31}p_{32}p_{33}p_{34}`$ $`z`$ $`=`$ $`p_1p_4p_5p_8p_9^2p_{10}p_{14}p_{15}p_{16}p_{17}p_{18}p_{19}p_{21}p_{22}p_{23}p_{24}`$ $`p_{25}p_{26}p_{27}p_{29}p_{30}p_{31}p_{32}p_{33}p_{34}`$ $`w`$ $`=`$ $`p_1p_3p_5p_6^2p_7p_{10}p_{12}p_{13}p_{14}p_{15}p_{16}p_{17}p_{18}p_{22}p_{23}p_{24}`$ $`p_{25}p_{26}p_{27}p_{29}p_{30}p_{31}p_{32}p_{33}p_{34}`$ $`v`$ $`=`$ $`p_1^2p_2p_3p_4p_5p_6p_7p_8p_9p_{10}p_{11}p_{12}p_{13}p_{14}^2p_{15}^2p_{16}^2p_{17}^2`$ (52) $`p_{18}^2p_{19}p_{20}p_{21}p_{22}^2p_{23}^2p_{24}^2p_{25}^2p_{26}^2p_{27}^2p_{28}p_{29}^2p_{30}^2p_{31}^2p_{32}^2p_{33}^2p_{34}^2`$ In terms of these variables, the space is defined by the surface $`xyzw=v^2`$ in $`𝐂^\mathrm{𝟓}`$. Two of its resolutions (corresponding to cases (Vb and VIb of ) are shown in figure 4. This singularity has to be analysed in the lines of in order to determine whether its partial resolutions are realised in the moduli space of the D-brane world volume gauge theory. We leave this issue for future work, and for the moment analyse the blowups of this singularity into lower dimensional orbifold singularities, using the inverse toric procedure. $``$ Let us consider the blowup illustrated in fig. 4(b). Performing Gaussian row reduction on the total charge matrix, we find that for the range of the FI parameters $`\zeta _1+\zeta _2=0;\zeta _3+\zeta _4=0;\zeta _5+\zeta _6=0`$, an appropriate initial choice determines the fields that retain zero vev as the set $`[p_2,p_4,p_8,p_9,p_{11},p_{19},p_{20},p_{21},p_{28}]`$, while all others have non-zero vev and are hence resolved. The reduced charge matrix is $`Q_t^{red}=\left(\begin{array}{cccccccccc}0& 1& 0& 1& 1& 1& 0& 0& 0& 0\\ 1& 0& 0& 0& 1& 0& 1& 0& 1& 0\\ 1& 0& 1& 1& 1& 0& 0& 1& 1& 0\\ 5& 5& 5& 0& 4& 1& 0& 1& 1& \zeta _5+\zeta _6\\ 5& 5& 5& 4& 0& 1& 0& 1& 1& \zeta _3+\zeta _4\\ 1& 5& 5& 4& 0& 1& 0& 1& 5& \zeta _1+\zeta _2\end{array}\right)`$ From the kernel of the charge matrix $`Q_{red}`$, the transpose of its dual is obtained as $`K^t=\left(\begin{array}{cccccccccccc}1& 1& 0& 1& 0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 1& 0& 0& 0& 1& 0& 1& 0& 0& 0\\ 1& 0& 0& 1& 0& 0& 0& 0& 1& 0& 1& 0\\ 0& 1& 1& 0& 1& 0& 0& 0& 0& 1& 0& 0\\ 0& 1& 0& 0& 0& 1& 0& 0& 0& 1& 0& 1\\ 0& 0& 1& 0& 1& 0& 1& 1& 0& 0& 0& 0\end{array}\right)`$ From these matrices, we can extract the gauge theory data which is now given by the matrix $`d_{red}=\left(\begin{array}{cccccccccccc}0& 1& 0& 1& 1& 0& 1& 0& 1& 0& 0& 1\\ 0& 0& 1& 1& 0& 1& 0& 1& 1& 1& 0& 0\\ 1& 1& 1& 0& 0& 0& 0& 1& 0& 0& 1& 1\\ 1& 0& 0& 0& 1& 1& 1& 0& 0& 1& 1& 0& \end{array}\right)`$ The gauge group is as expected, $`U(1)^4`$ with $`12`$ matter fields, and the superpotential is given by $`W`$ $`=`$ $`X_1X_5X_{12}X_3X_4X_{12}+X_2X_8X_9X_2X_7X_{11}+X_3X_6X_{11}`$ (53) $``$ $`X_5X_6X_9+X_{10}X_4X_7X_1X_8X_{10}`$ From the definition of the invariant variables in eq. (52), it can be seen that this space is $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$ defined by the equation $`xyz=v^2`$ in $`𝐂^\mathrm{𝟒}`$. In terms of the original fields, it corresponds to giving vevs to the fields $`W_{12}`$ and $`W_{21}`$. $``$ Finally, we come to the example of fig. 4(c). In this case, in the region of moduli space given by $`\zeta _2+\zeta _3+\zeta _6+\zeta _7=0`$, we find that the dynamical fields are given by the set $`[p_5,p_6,p_9,p_{10}]`$, while the rest acquire non-zero vev and are resolved. The reduced charge matrix is $`Q_t^{red}=\left(\begin{array}{ccccc}1& 1& 1& 1& 0\\ 5& 2& 2& 1& \zeta _2+\zeta _3+\zeta _6+\zeta _7\end{array}\right)`$ The matrix $`K^t`$ is in this case given by $`K^t=\left(\begin{array}{cccc}1& 1& 0& 0\\ 1& 0& 1& 0\\ 0& 1& 0& 1\end{array}\right)`$ Which specifies the matter content as a $`U(1)^2`$ gauge theory with the charge matrix given, as expected, by $`d_{red}=\left(\begin{array}{cccc}1& 1& 1& 1\\ 1& 1& 1& 1\end{array}\right)`$ With a superpotential $`W=\varphi \left(X_1X_4X_2X_3\right)`$ with $`\varphi `$ being a possible chargeless field. In terms of the original fields in the theory, this resolution corresponds to giving vevs to $`X_{18},X_{81},Y_{15},Y_{51}`$, and the resulting space is the $`𝐂^\mathrm{𝟐}/𝐙_\mathrm{𝟐}`$ orbifold $`zw=v^2`$ that we have discussed earlier. ## 5 Discussions In this paper, we have critically examined the inverse procedure of obtaining the world volume gauge theory data of D-branes probing certain toric singularities and their blowups from the geometric data of the resolution. We have shown that the algorithm of gives consistent results, for a given singularity that can be reached by partial resolutions of different parent theories. We have explicitly checked several partial resolutions of orbifolds of the form $`𝐂^\mathrm{𝟑}/\mathrm{\Gamma }`$ and $`𝐂^\mathrm{𝟒}/\mathrm{\Gamma }`$ and found the results to be in agreement with field theory calculations. However, as pointed out in , in the presence of chargeless fields in the matter multiplet, the procedure cannot be used to specify the superpotential interaction uniquely, because of the inherent problem in the handling of such fields by toric methods. We have treated the simplest examples of the resolution of the $`𝐂^\mathrm{𝟒}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$ singularity in this paper. It would be interesting to investigate the various partial resolutions of the moduli space of D-branes on this singularity the lines of and determine the gauge theory data thereof. Finally, as noted in , the ambiguities that exist in the inverse algorithm seem to imply that in some cases, different gauge theories, in the infra red limit flow to theories with identical moduli space. We have not dealt with this aspect in the present paper, and it would be an interesting direction for future work. Acknowledgements I would like to thank I. Biswas, S. Bhattacharjee, P. Chatterjee and D. Suryaramana for helpful discussions. Appendix Charge matrix for the partial resolutions of $`𝐂^\mathrm{𝟑}/𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}\times 𝐙_\mathrm{𝟐}`$: $$\begin{array}{c}Q_t=(\begin{array}{cccccccccccccccccc}p_1& p_2& p_3& p_4& p_5& p_6& p_7& p_8& p_9& p_{10}& p_{11}& p_{12}& p_{13}& p_{14}& p_{15}& p_{16}& p_{17}& \\ 0& 0& 0& 0& 1& 1& 0& 0& 1& 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 1& 0& 0& 0& 0& 1& 1& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 1& 0& 0& 1& 1& 0& 0& 0& 1& 1& 0& 0& 0\\ 1& 1& 0& 1& 1& 0& 1& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0\\ 1& 1& 1& 0& 1& 0& 0& 1& 0& 0& 0& 0& 0& 0& 0& 1& 0\\ 1& 1& 0& 1& 0& 1& 0& 1& 0& 0& 0& 0& 0& 0& 0& 0& 1\\ 1& 0& 1& 1& 1& 1& 0& 0& 1& 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0& 1& 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 1& 0& 0& 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 1& 1& 1& 0& 0& 0& 1& 0& 0& 0& 0\\ 2& 1& 1& 1& 1& 0& 0& 1& 1& 0& 1& 0& 1& 0& 0& 0& 0\\ 1& 0& 1& 1& 0& 0& 1& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 1& 1& 1& 0& 0& 0& 1& 0& 1& 0& 0& 0& 0& 0& 0& 0& 0\\ 1& 1& 1& 0& 0& 0& 0& 0& 0& 0& 1& 0& 1& 0& 0& 0& 0\\ 1& 0& 1& 1& 0& 0& 0& 1& 1& 0& 1& 0& 1& 0& 0& 0& 0\\ 2& 1& 1& 1& 1& 0& 1& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 1& 1& 0& 0& 0& 1& 0& 1& 0& 0& 0& 0\\ 1& 0& 1& 1& 1& 1& 1& 0& 0& 0& 0& 0& 1& 0& 0& 0& 0\\ 1& 1& 0& 1& 0& 1& 1& 1& 1& 0& 1& 0& 1& 0& 0& 0& 0\\ 1& 1& 1& 0& 1& 0& 1& 1& 1& 0& 1& 0& 1& 0& 0& 0& 0\\ 1& 1& 0& 1& 1& 0& 0& 0& 1& 0& 1& 0& 1& 0& 0& 0& 0\\ 1& 0& 0& 0& 1& 0& 1& 1& 0& 0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 1& 0& 1& 0& 1& 0& 1& 0& 0& 0& 0\\ 5& 2& 2& 2& 2& 0& 2& 2& 0& 2& 0& 2& 2& 1& 3& 3& 1\\ 1& 2& 2& 2& 2& 0& 2& 2& 0& 2& 0& 2& 2& 3& 1& 1& 5\\ 1& 2& 2& 2& 2& 0& 2& 2& 0& 2& 0& 2& 2& 3& 1& 1& 3\\ 3& 2& 2& 2& 2& 0& 2& 2& 0& 2& 0& 2& 2& 1& 3& 3& 1\\ 1& 2& 2& 2& 2& 0& 2& 2& 0& 2& 0& 2& 2& 5& 1& 1& 3\\ 3& 2& 2& 2& 2& 0& 2& 2& 0& 2& 0& 2& 2& 1& 5& 3& 1\\ 3& 2& 2& 2& 2& 0& 2& 2& 0& 2& 0& 2& 2& 1& 3& 5& 1\end{array}\mathrm{}\mathrm{}\hfill \\ \\ \\ \mathrm{}\begin{array}{ccccccccccccccccccc}p_{18}& p_{19}& p_{20}& p_{21}& p_{22}& p_{23}& p_{24}& p_{25}& p_{26}& p_{27}& p_{28}& p_{29}& p_{30}& p_{31}& p_{32}& p_{33}& p_{34}& p_{35}& \\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 0& 1& 0\\ 1& 2& 0& 2& 2& 3& 1& 1& 3& 2& 2& 1& 1& 3& 3& 1& 5& \zeta _1\\ 3& 2& 0& 2& 2& 1& 3& 3& 1& 2& 2& 3& 5& 1& 1& 3& 1& \zeta _2\\ 3& 2& 0& 2& 2& 1& 5& 5& 1& 2& 2& 3& 3& 1& 1& 3& 1& \zeta _3\\ 1& 2& 0& 2& 2& 5& 1& 1& 5& 2& 2& 1& 1& 3& 3& 1& 3& \zeta _4\\ 3& 2& 0& 2& 2& 1& 3& 3& 1& 2& 2& 3& 3& 1& 1& 5& 1& \zeta _5\\ 1& 2& 0& 2& 2& 3& 1& 1& 3& 2& 2& 1& 1& 3& 5& 1& 3& \zeta _6\\ 1& 2& 0& 2& 2& 3& 1& 1& 3& 1& 1& 1& 1& 5& 3& 1& 3& \zeta _7\end{array})\hfill \end{array}$$
warning/0005/hep-th0005116.html	ar5iv	text	# References IFT-UAM/CSIC-00-19 hep-th/0005116 Curved dilatonic brane-worlds and the cosmological constant problem Natxo Alonso-Alberca<sup>1</sup><sup>1</sup>1E-mail address: natxo.alonso@uam.es, Bert Janssen<sup>2</sup><sup>2</sup>2E-mail address: bert.janssen@uam.es and Pedro J. Silva<sup>3</sup><sup>3</sup>3E-mail address: psilva@delta.ft.uam.es Instituto de Física Teórica, C-XVI, Universidad Autónoma de Madrid E-28006 Madrid, Spain ABSTRACT We construct a model for dilatonic brane worlds with constant curvature on the brane, i.e. a non-zero four-dimensional cosmological constant, given in function of the dilaton coupling and the cosmological constant of the bulk. We compare this family of solutions to other known dilatonic domain wall solutions and apply a self-tuning mechanism to check the stability of our solutions under quantum fluctuations living on the brane. Recently the idea of brane worlds has received a lot of attention. In this picture our space-time has one (or more) extra non-compact spacial dimensions. In the scenario of Randall and Sundrum , a three-brane in a five-dimensional AdS-space was constructed, confining a Standard Model-like gauge theory on the brane. Due to the warped form of the metric of the solution, a different approach to the hierarchy problem was presented. But it also turns out that in this picture, perturbations of the metric on the brane have a five-dimensional profile which is normalisable and localised to the brane as if it were four-dimensional. Hence, an observer living in this brane world would see both gauge as gravity physics in the same way as an observer in a four-dimensional space-time. Different generalisations of this picture including a dilaton field were given in - while on the other hand, generalizations to non-flat brane worlds (without dilaton) appeared in , where a four-dimensional cosmological constant was introduced via constant curvature branes. Here it was shown that, in the brane world scenario, the four-dimensional cosmological constant is a geometrical property of the three-brane (namely its internal curvature), which in principle can be independent of the five-dimensional one. This shed new light on an old problem, namely why the observed cosmological constant in our universe is so small. One would expect that the non-zero vacuum expectation value of Standard Model fields would generate a non-zero vacuum energy, which would result in an effective non-zero cosmological constant.<sup>4</sup><sup>4</sup>4For a more general discussion on the cosmological constant problem, see and references therein. Now that it turns out that in the brane world picture the four-dimensional cosmological constant is a geometrical property of the brane, one should ask how this property is affected by quantum fluctuations of the field theory living on this brane. In it was observed that fluctuations of the brane tension, due to quantum corrections of the field theory living on a (flat) dilatonic brane, do not generate a four-dimensional cosmological constant. Via a self-tuning mechanism, the fluctuations in the brane tension can be absorbed into shifts of the dilaton, such that the quantum corrections do not curve the brane itself or the extra dimension. Thus, starting with flat brane worlds, no extra curvature, and hence no four-dimensional cosmological constant, is generated. Recent astronomical observations, however, point in the direction of a small positive, but non-zero cosmological constant. It is therefore interesting to look at brane-world models of non-zero curvature and see if there is a self-tuning mechanism working in these cases, which could protect this curvature from quantum corrections. The aim of this letter is to combine the results of - and , constructing dilatonic brane-worlds with a non-zero four-dimensional cosmological constant and test the self-tuning mechanism on these solutions. The organisation of this letter is as follows: we start with the construction of the curved dilatonic brane-world solutions and then give a small discussion of these solutions, comparing the obtained solutions with a class of domain wall solutions given in . Finally we will analyse the self-tuning mechanism in the case of curved dilatonic branes. As a starting point let us consider the action of five-dimensional dilatonic gravity coupled to a brane source in the presence of a (five-dimensional) cosmological constant: $`S`$ $`=`$ $`{\displaystyle \frac{1}{\kappa }}{\displaystyle d^4x𝑑y\sqrt{\|\widehat{g}\|}\left[\widehat{R}+\frac{4}{3}(\varphi )^2e^{\frac{\alpha }{3}\varphi }\mathrm{\Lambda }\right]}`$ (1) $`{\displaystyle d^4x\sqrt{\|\overline{g}\|}e^{{\scriptscriptstyle \frac{\beta }{3}}\varphi }V_0},`$ where $`V_0`$ is the tension of the brane source and $`\overline{g}_{mn}=\widehat{g}_{\mu \nu }\delta _m^\mu \delta _n^\nu `$ the induced metric on the brane. To solve the equations of motion, we propose the following curved brane-world Ansatz: $`ds^2`$ $`=`$ $`a(y)^2\stackrel{~}{g}_{mn}dx^mdx^ndy^2`$ $`\varphi `$ $`=`$ $`q\mathrm{log}a(y)+\varphi _0,`$ (2) with $`q`$ and $`\varphi _0`$ arbitrary constants. The internal brane metric $`\stackrel{~}{g}_{mn}`$ depends only on the internal coordinates $`x^m`$. Plugging this Ansatz in the equations of motion of the action (1) gives the following set of differential equations: $`\stackrel{~}{R}_{mn}\stackrel{~}{g}_{mn}[a\ddot{a}+3\dot{a}^2+\frac{1}{3}a^{\frac{\alpha q}{3}+2}e^{\frac{\alpha }{3}\varphi _0}\mathrm{\Lambda }]=+\frac{1}{6}\kappa V_0a^{\frac{\beta q}{3}+2}e^{\frac{\beta }{3}\varphi _0}\stackrel{~}{g}_{mn}\delta (y),`$ $`4a^1\ddot{a}+\frac{4}{3}q^2a^2\dot{a}^2+\frac{1}{3}a^{\frac{\alpha q}{3}}e^{\frac{\alpha }{3}\varphi _0}\mathrm{\Lambda }=\frac{2}{3}\kappa V_0a^{\frac{\beta q}{3}}e^{\frac{\beta }{3}\varphi _0}\delta (y),`$ (3) $`qa^1\ddot{a}+3qa^2\dot{a}^2\frac{\alpha }{8}a^{\frac{\alpha q}{3}}e^{\frac{\alpha }{3}\varphi _0}\mathrm{\Lambda }=+\frac{\beta }{8}\kappa V_0a^{\frac{\beta q}{3}}e^{\frac{\beta }{3}\varphi _0}\delta (y),`$ where $`\stackrel{~}{R}_{mn}`$ is the Ricci tensor of the internal metric $`\stackrel{~}{g}_{mn}`$ and a dot denotes derivative with respect to $`y`$. Analogous as in , the first equation of (3) can only be satisfied if the $`y`$-dependence vanishes: $$a\ddot{a}+3\dot{a}^2+\frac{1}{3}a^{\frac{\alpha q}{3}+2}e^{\frac{\alpha }{3}\varphi _0}\mathrm{\Lambda }+\frac{1}{6}\kappa V_0a^{\frac{\beta q}{3}+2}e^{\frac{\beta }{3}\varphi _0}\delta (y)=\stackrel{~}{\mathrm{\Lambda }},$$ (4) with $`\stackrel{~}{\mathrm{\Lambda }}`$ an abritrary integration constant which we interpret as the four-dimensional cosmological constant in the brane-world, since the first equation of (3) now translates into $`\stackrel{~}{R}_{mn}=\stackrel{~}{\mathrm{\Lambda }}\stackrel{~}{g}_{mn}`$. The solution to the equations (3)-(4) is given by $`ds^2`$ $`=`$ $`\left[1\frac{\alpha }{12}\sqrt{\mathrm{\Lambda }}\|y\|\right]^2e^{\frac{\alpha }{3}\varphi _0}\stackrel{~}{g}_{mn}dx^mdx^ndy^2,`$ $`\varphi `$ $`=`$ $`{\displaystyle \frac{6}{\alpha }}\mathrm{log}\left[1\frac{\alpha }{12}\sqrt{\mathrm{\Lambda }}\|y\|\right],`$ (5) where the coordinate $`y`$ runs from 0 to $`\frac{12}{\alpha \sqrt{\mathrm{\Lambda }}}`$ and the four-dimensional metric and cosmological constant satisfy $$\stackrel{~}{R}_{mn}=\stackrel{~}{\mathrm{\Lambda }}\stackrel{~}{g}_{mn},\stackrel{~}{\mathrm{\Lambda }}=\frac{16\alpha ^2}{48}\mathrm{\Lambda }e^{\frac{\alpha }{3}\varphi _0}.$$ (6) In the conformal frame, the above solution takes the form $`ds^2`$ $`=`$ $`\mathrm{exp}\left(\frac{\alpha }{6}\sqrt{\mathrm{\Lambda }}\|z\|\right)\left[e^{\frac{\alpha }{3}\varphi _0}\stackrel{~}{g}_{mn}dx^mdx^ndz^2\right],`$ $`\varphi `$ $`=`$ $`\frac{\sqrt{\mathrm{\Lambda }}}{2}\|z\|.`$ (7) The brane tension and dilaton coupling in the source term are given in terms of $`\alpha `$ and $`\mathrm{\Lambda }`$ by the jump equations: $$V_0=\frac{\alpha }{2\kappa }\sqrt{\mathrm{\Lambda }},\beta =\frac{8}{\alpha }.$$ (8) We thus see that the brane world are surfaces of positive (dS) or negative (AdS) constant curvature, depending on the dilaton coupling $`\alpha `$. It is remarkable that in the limit $`\stackrel{~}{\mathrm{\Lambda }}0`$ ($`\mathrm{\Lambda }0`$), we do not recover the general “flat” dilatonic RS theory , but only a particular case $`\alpha =\pm 4`$.<sup>5</sup><sup>5</sup>5In the case $`\alpha =4`$ was erroneously identified with the case $`\mathrm{\Lambda }=0`$ due to a coordinate singularity in the conformal frame. The holographic frame solution however is completely regular and can be seen as the limit of (5) for $`\stackrel{~}{\mathrm{\Lambda }}=0`$. Therefore, in contrast to “flat” dilatonic RS theory, it is no longer possible to make contact with the original RS scenario . Nor is it possible, in the limit $`\alpha 0`$ ($`\varphi `$ trivial), to make contact with the solutions of , as can clearly seen in the conformal frame (7). The solution for $`\mathrm{\Lambda }=0`$ ($`\stackrel{~}{\mathrm{\Lambda }}0`$) was given in . It is interesting to compare the solution (5) with the dilatonic domain wall solutions found in . It is clear that (5) can not be identified as one of the solutions given in . This is due to the fact that the Ansatz used in is different as our Ansatz (2). The main difference lays in the fact that the Ansatz of considers branes with constant spatial curvature, while the Ansatz (2) constructs branes with constant world volume (space-time) curvature. Still, there exists a particular solution that belongs to both classes, i.e. that can be obtained from each class as a special limit. For this case, the domain wall has to be flat, which implies $`\alpha =\pm 4`$ in equation (5) and $`M=k=0`$ in the Type II solutions of . Furthermore, in the Ansatz of the metric components $`g_{tt}`$ and $`g_{ij}`$ (the radial function in front of the spatial part of the brane world volume) should be identified. This particular solution is (in our notation, for $`\alpha =4`$ e.g.): $`ds^2`$ $`=`$ $`\left(1\frac{1}{3}\sqrt{\mathrm{\Lambda }}\|y\|\right)^2e^{\frac{4}{3}\varphi _0}\left[dt^2dx_i^2\right]dy^2,`$ $`e^\varphi `$ $`=`$ $`\left(1\frac{1}{3}\sqrt{\mathrm{\Lambda }}\|y\|\right)^{3/2}.`$ (9) An analysis, similar as the one done in reveals that this solution is a static domain wall, due to the fact that for this case the potential is constant. This particular solution (9) was first given in . A straightforward analysis of the perturbations of the metric (5) gives the following profile for the four-dimensional graviton $$\psi (y)=\left[1\frac{\alpha }{12}\sqrt{\mathrm{\Lambda }}\|y\|\right]^2,$$ (10) which is normalizable in the interval $`[0,\frac{12}{\alpha \sqrt{\mathrm{\Lambda }}}]`$ and localized around $`\|y\|=0`$. Note that the dependence on dilaton coupling $`\alpha `$ is much weaker than in the “flat” dilatonic RS scenario, where we had an exponential dependence on the coupling. Finally, note that, similar to other dilatonic brane world models, the solution (5)-(8) has a time-like naked singularity in $`y=\frac{12}{\alpha \sqrt{\mathrm{\Lambda }}}`$ ($`z=\mathrm{}`$). In the brane world picture, the four-dimensional cosmological constant $`\stackrel{~}{\mathrm{\Lambda }}`$ is, rather than a input parameter of the Lagrangian, a geometrical quantity, related to the curvature of the domain wall. In it was argued that for flat domain walls (i.e. $`\stackrel{~}{\mathrm{\Lambda }}=0`$), the quantum fluctuations of the gauge theory living on the brane can be absorbed in shifts of the dilaton, which turns out to be a symmetry of the solutions. This mechanism, called self-tuning, makes that no extra curvature is generated and that the four-dimensional cosmological constant remains zero, even after quantum corrections. Since recent astronomical observations seem to indicate that the cosmological constant might be very small, but non-zero and positive, a natural question to ask is how far this self-tuning mechanism can be extended, in particular to branes with non-zero curvature (i.e. solutions with non-zero four-dimensional cosmological constant). It turns out that no self-tuning mechanism is possible for the solutions we have constructed, due to the fact that the jump equation (8) does not depend explicitly on $`\varphi _0`$. In fact, $`V_0`$ only depends on the bulk parameters and an arbitrary integration constant (which, without lost of generality, we set equal to $`\frac{12}{\alpha \sqrt{\mathrm{\Lambda }}}`$), indicating the position of the naked singularity. Changes of $`V_0`$ due to quantum fluctuations could only be absorbed by this integration constant, but this would mean that the positions of the singularities should change, which seems to have no sense. This seems to indicate that the self-tuning mechanism of is only valid for the case of flat branes and does not provide a satisfying explanation of the cosmological constant problem in the light of the recent astronomical observations. The family of solutions we have constructed does not include the $`\alpha =0`$ case, as can be seen from (5) and (8). With the Ansatz $`\varphi =\varphi (y)`$, the equations of motion for $`\alpha =0`$ read: $$\begin{array}{ccc}& & \ddot{\varphi }+4a^1\dot{a}\dot{\varphi }=0,\hfill \\ & & a^1\ddot{a}+\frac{1}{3}\dot{\varphi }^2+\frac{1}{12}\mathrm{\Lambda }=0,\hfill \\ & & a^1\ddot{a}+3a^2\dot{a}^2+\frac{1}{3}\mathrm{\Lambda }=\stackrel{~}{\mathrm{\Lambda }}a^2.\hfill \end{array}$$ (11) The first of these equations gives $$\dot{\varphi }=\gamma a^4,$$ (12) where $`\gamma `$ is an arbitrary integration constant. Substituting (12) in (11), the last two equations reduce to a single one, given by: $$\dot{a}=ϵ\sqrt{\frac{\gamma ^2}{9}a^6\frac{\mathrm{\Lambda }}{12}a^2+\frac{\stackrel{~}{\mathrm{\Lambda }}}{3}},$$ (13) where $`ϵ=\pm 1`$ determines the branch of the square root chosen in the solution. The sign of the argument of the square root in (13) will depend on the values of $`\mathrm{\Lambda }`$ and $`\stackrel{~}{\mathrm{\Lambda }}`$. Both $`\mathrm{\Lambda }`$ and $`\stackrel{~}{\mathrm{\Lambda }}`$ are arbitrary and independent. The solution only makes sense when this argument is positive. Integrating (13) we obtain: $$^a\frac{ϵda^{}}{\sqrt{\frac{\gamma ^2}{9}a_{}^{}{}_{}{}^{6}\frac{\mathrm{\Lambda }}{12}a_{}^{}{}_{}{}^{2}+\frac{\stackrel{~}{\mathrm{\Lambda }}}{3}}}=y+y_0.$$ (14) The l.h.s. of (14) cannot be expressed in terms of any elementary function, so it is hard to see whether a self-tuning mechanism could work. As in , it could be possible to study the bounds on the brane cosmological constant depending on the signs and values of $`\mathrm{\Lambda }`$ and $`\stackrel{~}{\mathrm{\Lambda }}`$. We leave this question open for future research. Acknowledgments We would like to thank César Gómez, Patrick Meessen and Tomás Ortín for useful discussions. We also thank Jim Cline and Christophe Grojean for pointing out some problems in a previous version. The work of N.A.A. and B.J. has been supported by the TMR program FMRX-CT96-0012 on Integrability, non-perturbative effects, and symmetry in quantum field theory.
warning/0005/gr-qc0005060.html	ar5iv	text	# LASER INTERFEROMETER GRAVITATIONAL WAVE DETECTORS ## Acknowledgments This work was partially supported by the European Union, TMR Contract No. ERBFMBICT972771. ## References
warning/0005/hep-th0005094.html	ar5iv	text	# Untitled Document IC/2000/58 AdS/CFT Correspondence, Critical Strings and Stochastic Quantization Dimitri Polyakov polyakov@ictp.trieste.it The Abdus Salam International Centre for Theoretical Physics Strada Costiera 11 I-34014 Trieste, ITALIA Abstract In our previous paper we have shown that the NSR string sigma-model with the massless 5-form vertex operator in $`D=10`$ NSR string theory: $`V_5e^{3\varphi }\psi _0\psi _1\psi _2\psi _3\psi _t\overline{}X^te^{ik^{\|\|}x^{\|\|}}`$ ($`t=4,\mathrm{..9}`$) reproduces the correlators of the $`N=4`$ $`D=4`$ super Yang-Mills theory. In particular, this implies that the sigma-model with the $`V_5`$ operator in flat space-time should be the NSR analogue of the GS string theory on $`AdS_5\times S^5`$. This means that the $`V_5`$-operator plays the role of cosmological constant, curving flat ten-dimensional space-time into that of $`AdS_5\times S^5`$ In the present paper we show that dilaton beta-function equation in such a sigma-model has the form of stochastic Langevin equation with the non-Markovian noise. The worldsheet cutoff is identified with stochastic time and the $`V_5`$-operator plays the role of the noise. We derive the Fokker-Planck equation associated with this stochastic process and show that the Hamiltonian of the $`AdS_5`$ supergravity defines the distribution satisfying this Fokker-Planck equation. This means that the dynamical compactification of the space-time on $`AdS_5\times S^5`$ occurs as a result of the non-Markovian stochastic process, generated by the $`V_5`$-operator noise. This provides us with an insight into relations between holography principle and the concept of stochastic quantization from the point of view of critical string theory. May 2000 PACS: $`04.50.+h`$;$`11.25.Mj`$. Introduction The holography principle \[1,,2,,3,,4\] relates degrees of freedom of quantum field theories in D and $`D+1`$ dimensions. The remarkable example is the AdS/CFT correspondence which relates the AdS supergravity to supersymmetric conformal field theories on the AdS boundary. In particular, in case of $`D=4`$ this implies the correspondence between five-dimensional AdS supergravity and $`N=4`$ super Yang-Mills theory. Another, seemingly unrelated example connecting field theories in $`D`$ and $`D+1`$ dimensions is known from the concept of stochastic quantization. In the approach of stochastic quantization \[5\] (for comprehensive review see also \[6,,7\]) the dynamics of a given D-dimensional quantum field theory is studied by introducing the white noise to the field theory, or coupling it to infinite thermal reservoir. The interaction with the reservoir breaks the thermal equilibrium; as a result the field theory starts behaving like a brownian-like object, drifting in an $`additional`$ “fictitious” dimension, known as stochastic time $`\tau `$. In the limit of infinite stochastic time this non-equilibrium $`D+1`$-dimensional system with certain Fokker-Planck distribution returns to the initial equilibrium state; the logarithm of the equilibrium limit of the FP distribution is identified with the action of the original $`D`$-dimensional field theory. The dynamics of this $`D+1`$-dimensional stochastic field theory is described by the correlators $`<\varphi (x_1,\tau _1)\mathrm{}.\varphi (x_n,\tau _n)>`$ where $`\varphi (x,\tau )`$ are the solutions of the Langevin equation: $$\frac{d\varphi }{d\tau }=\frac{\delta S_E}{\delta \varphi }+\eta (x,\tau )$$ Here $`S_E`$ is the action of the equilibrium $`D`$-dimensional theory and $`\eta `$ is the white noise. To obtain correlation functions the product of these solutions (as functions of $`\eta `$, obtained by iterations) must be averaged over the white noise with the gaussian distribution. This gives rise to the $`D+1`$-dimensional stochastic diagrams which, in the limit of $`\tau \mathrm{}`$ reproduce the Feynman graphs of the D-dimensional theory. In the present paper we will try to show that these two approaches , the stochastic quantization and the holography principle, are in fact related to each other in a rather intriguing way and that such a relation may be observed and explored from string theory point of view. The key element of our construction is so called brane-like sigma-model which may be understood as the NSR analogue of the GS superstring theory on $`AdS_5\times S^5`$, as the analysis of its correlation functions shows \[8\]. The action for this model is given by: $$\begin{array}{cc}\hfill S_{b.l.}=d^2zX^\mu \overline{}X_\mu +\psi ^\mu \overline{}\psi _\mu +\overline{\psi }^\mu \overline{\psi }_\mu & \\ \hfill +\lambda _0ϵ_{a_1a_2a_3a_4}\frac{d^4k^{\|\|}}{k_{}^{\|\|}{}_{}{}^{4}}e^{3\varphi }\psi ^{a_1}\psi ^{a_2}\psi ^{a_3}\psi ^{a_4}\psi _t\overline{}X^te^{ik^{\|\|}X}+c.c.+ghosts& \end{array}$$ Here $`X^\mu (z,\overline{z})(\mu =0,\mathrm{},9)`$ are space-time coordinates, $`\psi _\mu `$ and $`\overline{\psi }_\mu `$ are NSR worldsheet fermions and $`\varphi `$ is bosonized superconformal ghost field. Furthermore, the space-time index $`\mu `$ is split in the $`4+6`$ way, so that $`a=0,\mathrm{},3`$ and $`t=4,\mathrm{},9`$. In other words, the sigma-model (2) is given by NSR string theory action in flat space-time with the potential term related to the closed-string brane-like vertex operator $$V_5(z,\overline{z},k^{\|\|})=\frac{\lambda _0}{k_{}^{\|\|}{}_{}{}^{4}}ϵ_{a_1\mathrm{}a_4}e^{3\varphi }\psi _{a_1}\mathrm{}\psi _{a_4}\psi _t\overline{}X^te^{ik^{\|\|}X}(z,\overline{z})$$ which we will also refer to as $`V_5`$-operator in the rest of the paper. The properties of the sigma-model with the $`V_5`$-operator have been studied in \[8\] where the relevance of this vertex operator to non-perturbative $`D3`$-brane dynamics has been shown. It is important that BRST invariance condition for the $`V_5`$-operator (3) requires that its propagation is confined to four dimensions. Namely, to insure the BRST invariance, the momentum $`k^{\|\|}`$ must be polarized along the $`a_1,\mathrm{}a_4`$ directions. Moreover, the condition of the worldsheet conformal invariance (preserving the conformally invariant form of the O.P.E. between two stress-energy tensors corresponding to the action (2) or the vanishing of the beta-function in the lowest order of string perturbation theory) requires that the space-time scalar field $`\lambda (k^{\|\|})`$, corresponding to the $`V_5`$-operator, should behave as $$\lambda (k^{\|\|})\frac{\lambda _0}{k_{}^{\|\|}{}_{}{}^{4}}$$ where $`\lambda _0`$ is constant. As has been argued in \[8\] the role of the $`V_5`$ operator is that it transforms the flat ten-dimensional space-time vacuum into that of $`AdS_5\times S^5`$, thus connecting two maximally supersymmetric backgrounds in ten dimensions. This is because adding the $`V_5`$-term to the sigma-model action (2) is in fact equivalent to introducing the $`D3`$-brane in the theory. As a result, one may explore string theory in the AdS background (and consequently, the large N limit of gauge theory) by means of the brane-like sigma-model (2) which technically lives in $`flat`$ ten-dimensional space-time. Using the correspondence between local gauge invariant operators in the large N Yang-Mills theory and massless vertex operators in string theory one may obtain the large N correlators in gauge theory by computing scattering amplitudes of the BRST invariant vertices in the sigma-model (2). For example, as the dilaton vertex operator $`V_\phi `$ corresponds to the $`TrF^2`$ field in gauge theory, the generating functional for various correlation functions of the $`TrF^2`$ operators is given by: $$\begin{array}{cc}\hfill Z(\lambda _0,\phi )=D[X]D[\psi ]D[ghosts]f(\mathrm{\Gamma },N)exp\{d^2zX^\mu \overline{}X_\mu +\psi ^\mu \overline{}\psi _\mu +\overline{\psi }^\mu \overline{\psi }_\mu & \\ \hfill +\lambda _0ϵ_{a_1a_2a_3a_4}\frac{d^4k^{\|\|}}{k_{}^{\|\|}{}_{}{}^{4}}e^{3\varphi }\psi ^{a_1}\psi ^{a_2}\psi ^{a_3}\psi ^{a_4}\psi _t\overline{}X^te^{ik^{\|\|}X}+d^{10}pV_\phi (p,z,\overline{z})\phi (p)\}& \\ \hfill +c.c.+ghosts& \end{array}$$ where $`\phi (p)`$ is ten-dimensional space-time dilaton field. The “measure function” $`f(\mathrm{\Gamma },N)(1+N^2\mathrm{\Gamma }^4)^1\times c.c.`$ ( $`\mathrm{\Gamma }`$ is picture-changing operator and $`N`$ corresponds to the gauge group parameter) needs to be introduced to the measure of integration to insure correct ghost number balance on the sphere and normalization of scattering amplitudes. The two-point dilaton correlation function, corresponding to the generating functional (5) is given by $$\begin{array}{cc}\hfill <V_\phi (p_1)V_\phi (p_2)>_{\sigma model}=\frac{\delta ^2Z(\lambda _0,\phi )}{\delta \phi (p_1)\delta \phi (p_2)}\|_{\phi =0}& \end{array}$$ To compute this correlator we have to expand the functional (5) in $`\lambda _0`$. The first non-trivial contribution has the order of $`\lambda _0^2`$ and it is given by $$\begin{array}{cc}\hfill A_{\lambda _0^2}(p_1,p_2)\lambda _0^2\frac{d^4k_1^{\|\|}}{k_{1}^{\|\|}{}_{}{}^{4}}\frac{d^4k_2^{\|\|}}{k_{2}^{\|\|}{}_{}{}^{4}}<V_\phi (p_1)V_\phi (p_2)V_5(k_1^{\|\|})V_5(k_2^{\|\|})>& \end{array}$$ where the four-point amplitude should be computed in the usual NSR string theory in flat space-time. In other words, this is just the usual four-point closed string Veneziano amplitude which has to be integrated over internal momenta of the $`V_5`$-vertices, i.e. over two out of three independent momenta. The straightforward computation of the four-point amplitude and the integration over the $`V_5`$ momenta has been performed in \[8\] and the answer is given by: $$\begin{array}{cc}\hfill A_{\lambda _0^2}\lambda _0^2(p_1^\|\|)^4log(p_1^{\|\|})^2d^2w\frac{log(\|log\|w\|\|)+log(\|log\|1w\|\|)}{\|1w\|^4}& \end{array}$$ where $`p_1^{\|\|}`$ is the longitudinal projection of the dilaton momentum to four longitudinal directions; $`(p^\|\|)^2=p_\alpha p^\alpha `$. It is remarkable that the amplitude (8) depends exclusively on four-dimensional longitudinal projection of the dilaton momentum; up to normalization it has the same form as the two-point correlator $`<TrF^2(p_1^{\|\|})TrF^2(p_1^{\|\|})>`$ in the $`N=4`$ super Yang-Mills theory in $`D=4`$, computed in the approximation of dilaton s-wave \[3\]. Fourier transforming the amplitude (8), one recovers the well-known expression for the two-point amplitude in the $`N=4`$ $`D=4`$ SYM theory in the four-dimensional coordinate space: $`TrF^2(x)TrF^2(y)\frac{1}{\|xy\|^8}`$. Furthermore, the momentum structure of amplitudes with more $`V_5`$ insertions agrees with the form of the $`<TrF^2TrF^2>`$ correlators computed at higher values of the dilaton angular momentum; in other words, expansion in the $`\lambda _0`$ parameter in the brane-like sigma-model (2) accounts for higher partial waves of the dilaton field in the $`AdS_5\times S^5`$ supergravity. Proceeding similarly, one can in principle compute higher point correlation functions from the generating functional (5) to show their agreement with the known expressions for 3 and 4-point correlators in the $`N=4,D=4`$ SYM theory. Regarding the four-point functions, one can recover the logarithmic singularities in the position space, observed previously in the $`AdS_5`$ supergravity computations \[9,,10\]. In this paper we will explore the mechanism of the dynamical compactication of flat ten-dimensional space-time on $`AdS_5\times S^5`$ caused by the presence of potential term with the $`V_5`$ vertex in the brane-like sigma-model (9). Our explanation of this mechanism is that the compactification on $`AdS_5\times S^5`$ occurs as a result of certain very special non-Markovian stochastic process. Namely, we will show that the $`V_5`$ insertion in the sigma-model (5) has a meaning of a “random force” term with the $`V_5`$-operator playing the role of a non-Markovian stochastic noise, which correlations are determined by the worldsheet beta-function associated with the $`V_5`$ vertex. We will see that the presence of the $`V_5`$-term in the sigma-model (2) leads to Langevin-type equation in the low-energy effective theory. The role of the stochastic time variable in this Langevin equation will be played by the worldsheet cutoff. This means that the $`V_5`$-term effectively describes the interaction of critical NSR string theory with some infinite thermal reservoir, breaking thermal equilibrium and causing the theory to drift stochastically, behaving like a brownian-type object. Moreover, the reservoir itself is not in the equilibrium, which is reflected in the non-Markovian properties of the $`V_5`$-noise. Therefore the $`V_5`$-noise is not the usual white noise but it propagates in the stochastic time. Such a propagation is defined by the worldsheet correlators of the $`V_5`$-vertices. As a result of the reservoir evolution, the state of thermal equilibrium, reached in the limit of infinite stochastic time in the Langevin equations, is $`different`$ from the initial equilibrium state. In other words, the overall scenario is the following: in the beginning we have a field theory in $`flat`$ space time, corresponding to the low-energy limit of usual NSR string theory. This corresponds to the situation when we are very far from the $`D3`$-brane, corresponding to the $`V_5`$-insertion. Next, the $`D3`$-brane, or the non-Markovian $`V_5`$-noise appears in the picture. This leads to the appropriate Langevin-type equations for the dilaton, associated with corresponding modified Fokker-Planck equation ( this FP equation should be modified due to non-trivial correlations of the noise). Next, we shall find that the distribution solving this modified FP equation is given by the exponent of the $`AdS_5`$ supergravity Hamiltonian, living on a constant time slice. Evolution in the stochastic time corresponds to the flow in the direction transverse to the $`D3`$-brane worldvolume. Finally, in the limit of the infinite stochastic time we reach $`another`$ equilibrium limit, corresponding to the theory living in the $`D3`$-brane worldvolume, i.e. the four-dimensional gauge theory. Of course, such a scenario also implies the correspondence between radial AdS coordinate and the stochastic time which in turn corresponds to the worldsheet cutoff in the model (2). This interpretation also implies that the expansion in the $`\lambda _0`$ parameter in the brane-like sigma-model (9) corresponds to the the “stochastic diagrams” which have to be averaged over the white noise distribution. It is possible that the consequence of such an averaging is reflected in the numerical factor given by the integral over the worldsheet coordinate $`w`$ appearing in the amplitude (8), the only worldsheet integration left after fixing Koba-Nielsen’s measure in the four-point correlator. Though the last suggestion seems rather heuristic and we will not provide any evidence for it in this paper, it nevertheless seems very tempting to ask whether the double logarithm appearing in the integral over $`w`$ in (8) has anything to do with the one appearing in the well-known law of the iterated logarithm which limits the maximum distance covered by a Brownian particle for a time interval t: $`\|x(t)x(0)\|<const\times \sqrt{t}log\|log\|t\|\|`$. Indeed, in this paper we will point out connections between the radial worldsheet coordinate (more precisely, the worldsheet cutoff), the stochastic time parameter and the radial $`AdS_5`$ coordinate. With all these intuitive ideas in mind, in the following section we will proceed to deriving the Langevin equation from the brane-like sigma-model (2) and the corresponding Fokker-Planck equation. Langevin and Fokker-Planck equations from the brane-like sigma-model Consider expansion of the generating functional (5) in $`\phi (X)`$ and $`\lambda _0`$. We get: $$\begin{array}{cc}\hfill Z(\lambda _0,\phi )D[X]D[\psi ]D[ghosts]e^{S_0^{NSR}}\{1+d^{10}p\phi (p)d^2zV_\phi (p,z,\overline{z})& \\ \hfill +\frac{1}{2}d^{10}p_1d^{10}p_2\phi (p_1)\phi (p_2)d^2z_1d^2z_2V_\phi (p_1,z_1,\overline{z_1})V_\phi (p_2,z_2,\overline{z_2})& \\ \hfill +\lambda _0\frac{d^4k^{\|\|}}{k_{}^{\|\|}{}_{}{}^{4}}d^2wV_5(k^{\|\|},w,\overline{w})& \\ \hfill +\frac{1}{2}\lambda _0^2\frac{d^4k_1^{\|\|}}{k_{1}^{\|\|}{}_{}{}^{4}}\frac{d^4k_2^{\|\|}}{k_{2}^{\|\|}{}_{}{}^{4}}d^2w_1d^2w_2V_5(k_1^{\|\|},w_1,\overline{w_1})V_5(k_2^{\|\|},w_2,\overline{w_2})& \\ \hfill +\frac{\lambda _0^3}{6}\frac{d^4k_1^{\|\|}}{k_{1}^{\|\|}{}_{}{}^{4}}\frac{d^4k_2^{\|\|}}{k_{2}^{\|\|}{}_{}{}^{4}}\frac{d^4k_3^{\|\|}}{k_{3}^{\|\|}{}_{}{}^{4}}d^2w_1d^2w_2d^2w_3V_5(k_1^{\|\|},w_1,\overline{w_1})& \\ \hfill \times V_5(k_2^{\|\|},w_2,\overline{w_2})V_5(k_3^{\|\|},w_3,\overline{w_3})+\mathrm{}\}f(\mathrm{\Gamma },N)& \end{array}$$ where $`S_0^{NSR}`$ is the free part of NSR superstring action in flat space-time and we only retained the terms relevant to our discussion. Consider first the term proportional to the square of the dilaton in the expansion (9). The O.P.E. between two dilaton vertex operators is given by $$\begin{array}{cc}\hfill L_\phi =d^{10}p_1d^{10}p_2d^2z_1d^2z_2V_\phi (z_1,\overline{z_1},p_1)V_\phi (z_2,\overline{z_2},p_2)& \\ \hfill d^{10}p_1d^{10}p_2d^2z_1d^2z_2[\frac{C_\phi (p_1,p_2)}{\|z_1z_2\|^2}V_\phi (\frac{1}{2}(z_1+z_2),\frac{1}{2}(\overline{z_1}+\overline{z_2}),p_1+p_2)+\mathrm{}]& \\ \hfill C_\phi (p_1,p_2)(p_1p_2)^2& \end{array}$$ Here $`C_\phi (p_1,p_2)`$ is the structure constant and we have dropped all the terms not contributing to the dilaton’s beta-function. The worldsheet integrals in (9), (10) are divergent as the logarithmic divergence originates from the region $`\|z_1z_2\|\mathrm{\Lambda }`$ where $`\mathrm{\Lambda }`$ is the worldsheet cutoff. The relevant cutoff dependent part of the O.P.E. in (10) is then given by $$\begin{array}{cc}\hfill L_\phi log\mathrm{\Lambda }d^{10}p_1d^{10}p_2d^2zC_\phi \phi (p_1)\phi (p_2)V_\phi (z,\overline{z},p_1+p_2)& \\ & \\ \hfill =log\mathrm{\Lambda }d^{10}pd^2zV_\phi (p,z,\overline{z})d^{10}qC_\phi (q)\phi (\frac{p+q}{2})\phi (\frac{pq}{2})& \end{array}$$ where we have changed the variables: $$\begin{array}{cc}\hfill z=z_1+z_2& \\ \hfill w=z_1z_2& \\ \hfill p=p_1+p_2& \\ \hfill q=p_1p_2& \end{array}$$ Combining this with the linear dilaton term in (9), it is easy to see that in the absence of the $`V_5`$ term the renormalization of the space-time dilaton field necessary to remove the logarithmic divergence in the partition function (9) would be given by $$\phi (p)\phi (p)log\mathrm{\Lambda }d^{10}qC_\phi (q)\phi (\frac{p+q}{2})\phi (\frac{pq}{2})$$ provided that one does not turn on the graviton excitation. The condition for the worldsheet conformal invariance leads then to low energy effective equations of motion for the dilaton field in flat space-time: $$\beta _\phi \frac{d\phi }{d(log\mathrm{\Lambda })}=d^{10}qC_\phi (q)\phi (\frac{p+q}{2})\phi (\frac{pq}{2})=0$$ The presence of the $`V_5`$ vertex, however, modifies the dilaton beta-function. The relevant part of the operator algebra is given by: $$V_5(w_1,\overline{w_1},k_1^{\|\|})V_5(w_2,\overline{w_2},k_2^{\|\|})\frac{C_{5\phi }(k_1^{\|\|},k_2^{\|\|})}{\|w_1w_2\|^2}V_\phi (\frac{1}{2}(w_1+w_2),\frac{1}{2}(\overline{w_1}+\overline{w_2});k_1^{\|\|}+k_2^{\|\|})$$ where $`C_{5\phi }(k_1^{\|\|}+k_2^{\|\|})^4`$ is again the appropriate structure constant; it is obtained easily by computing the three-point correlator of two $`V_5`$ states with the dilaton on the sphere. It is remarkable that this structure constant is proportional to the fourth power in the momentum and has no contributions quadratic in $`p`$, unlike the case of the graviton. Proceeding precisely as above and substituting the structure constant $`C_{5\phi }`$ we find that the logarithmically divergent contributions relevant to the dilaton’s beta function are given by: $$\begin{array}{cc}\hfill Z(\lambda _0,\phi )D[X]D[\psi ]D[ghosts]e^{S_0^{NSR}}\{1+d^{10}pd^2zV_\phi (z,\overline{z},p)\{\phi (p)& \\ \hfill +log\mathrm{\Lambda }[d^{10}qC_\phi (q)\phi (\frac{pq}{2})\phi (\frac{p+q}{2})& \\ \hfill +\lambda _0^2d^4k_1^{\|\|}\frac{1}{(p^{\|\|}k_1^{\|\|})^4}(1+\lambda _0\frac{d^4k_2^{\|\|}}{k_2^{\|\|}}_\mathrm{\Lambda }d^2wV_5(w,\overline{w},k_2^{\|\|}))]\}\}& \end{array}$$ where the integral over $`w`$ is taken outside the two-dimensional cutoff region. To analyze this cutoff dependence, we write $`w=re^{i\alpha }`$ and $$\begin{array}{cc}\hfill _\mathrm{\Lambda }d^2wV_5(w,\overline{w},k^{\|\|})=_0^{2\pi }𝑑\alpha _\mathrm{\Lambda }^{\mathrm{}}𝑑rrV_5(r,\alpha ,k^{\|\|})=_0^{2\pi }𝑑\alpha _0^{\mathrm{}}𝑑rrV_5(r+\mathrm{\Lambda },\alpha ,k^{\|\|})& \end{array}$$ The integral over $`k_1^{\|\|}`$in (16) is taken over the on-shell surface; it is computed by introducing the exponential regulator: $`\frac{d^4k^{\|\|}}{k_{}^{\|\|}{}_{}{}^{4}}\frac{d^4k^{\|\|}e^{i\alpha k}}{k_{}^{\|\|}{}_{}{}^{4}}`$ and then taking the limit $`\alpha 0`$ The value of this integral is equal to $`1+log\alpha ^2`$ which gives just 1 after the regularization. Therefore the dilaton’s beta-function equation following from the partition function expansion is given by $$\frac{d\phi (p)}{d(log\mathrm{\Lambda })}=d^{10}qC_\phi (q)\phi (\frac{pq}{2})\phi (\frac{p+q}{2})+\eta _5(p^{\|\|},\mathrm{\Lambda })$$ where $$\eta _5(p^{\|\|},\mathrm{\Lambda })\lambda _0^2(1+\lambda _0d^4k_2^{\|\|}_0^{2\pi }𝑑\alpha _0^{\mathrm{}}𝑑rrV_5(r+\mathrm{\Lambda },\alpha ,k^{\|\|}))$$ The dilaton beta-function equation (18) has the form of the Langevin equation (1), with the role of the stochastic noise term being played by the truncated worldsheet integral of the $`V_5`$-vertex (19). The logarithm of the worldsheet cutoff parameter plays the role of the stochastic time in the Langevin equation (18). The noise (19) is non-Markovian and it is generated by the $`V_5`$ operator (3), as was already noted above. Prior to computing the correlators of the $`V_5`$-noise (19), let us derive the Fokker-Planck equation associated with the Langevin equation (18). It should be stressed that in our derivation we will have to take into account that the noise is non-Markovian, i.e. its correlators depend non-trivially on the stochastic time. Let us assume the two-point noise correlation function is given by $$<\eta _5(p_1^{\|\|},\tau _1)\eta _5(p_2^{\|\|},\tau _2)>G(p_1^{\|\|},p_2^{\|\|},\tau _1,\tau _2)$$ where $`\tau log\mathrm{\Lambda }`$ and the noise Green’s function $`G(p_1^{\|\|},p_2^{\|\|},\tau _1,\tau _2)`$ is determined by properly regularized worldsheet correlators of the $`V_5`$ vertices in (19). The precise form of this Green’s function will be computed below. Let us furthermore assume that the Green’s function (19) corresponds to some local effective action $`I(\eta _{V_5})`$ which determines the propagation of the noise in stochastic time: $$I(\eta )d^4p𝑑\tau (\eta _5K\eta _5+V(\eta _5))$$ where $`K`$ is some differential operator and $`V`$ is the potential. To derive the Fokker-Planck equation, associated with the non-Markovian Langevin equation (18), consider an arbitrary $`\tau `$-dependent functional of the dilaton field $`f(\phi (p,\tau ),\tau )`$ Next, consider the average of the time derivative of this functional, taking the Langevin equation (18) into account: $$\begin{array}{cc}\hfill <\frac{df(\phi ,\tau )}{d\tau }>=d^4p<\frac{\delta f(\phi ,\tau )}{\delta \phi (p,\tau )}\frac{d\phi }{d\tau }>& \\ \hfill =d^4p\{<\frac{\delta f(\phi ,\tau )}{\delta \phi (p,\tau )}\frac{\delta S_E}{\delta \phi }>+<\frac{\delta f(\phi ,\tau )}{\delta \phi (p,\tau )}\eta _5(p,\tau )>\}& \end{array}$$ where $`S_E`$ is the dilaton low energy effective action in flat space-time (i.e. in the absence of the $`V_5`$-term and the graviton). Integrating by parts, we can write the left-hand side of (22) as $$<\frac{df(\phi ,\tau )}{d\tau }>=D\phi P_{FP}(\phi ,\tau )\frac{df(\phi ,\tau )}{d\tau }=D\phi \frac{dP_{FP}(\phi ,\tau )}{d\tau }f(\phi ,\tau )$$ where $`P_{FP}(\phi ,\tau )`$ is five-dimensional Fokker-Planck distribution associated with the Langevin equation (18). Next, the first correlator on the right hand side of (22), involving $`S_E`$ may be written as $$\begin{array}{cc}\hfill <\frac{\delta f(\phi ,\tau )}{\delta \phi (p,\tau )}\frac{\delta S_E}{\delta \phi }>=d^4pD\phi P_{FP}(\phi ,\tau )<\frac{\delta f(\phi ,\tau )}{\delta \phi (p,\tau )}\frac{\delta S_E}{\delta \phi }>& \\ \hfill =D\phi f(\phi ,\tau )d^4pd^4q\frac{\delta }{\delta \phi (p,\tau )}[\frac{\delta S_E}{\delta \phi (q)}P_{FP}(\phi ,\tau )]& \end{array}$$ given that partial integration was again performed. Finally, consider the second term in the right hand side of (22). We have: $$\begin{array}{cc}\hfill d^4p<\frac{\delta f(\phi ,\tau )}{\delta \phi (p,\tau )}\eta _5(p,\tau )>=D\eta _5d^4pe^{I(\eta _5)}\frac{\delta f(\phi ,\tau )}{\delta \phi (p,\tau )}\eta _5(p,\tau )& \\ \hfill =lim_{j0}d^4p\frac{\delta }{\delta j(p,\tau )}D\eta _5e^{I(\eta _5)+{\scriptscriptstyle d^4q𝑑tj(q,t)\eta _5(q,t)}}\frac{\delta f}{\delta \phi }(p,\tau )& \end{array}$$ where we introduced infinitezimal source term in the last expression. Let us perform the transformation: $$\eta _5\eta _0+\stackrel{~}{\eta }$$ where $`\eta _0`$ is the solution of equation $`\frac{\delta I}{\delta \eta }=j`$ The functional integral (25) becomes $$\begin{array}{cc}\hfill lim_{j0}\frac{\delta }{\delta j}d^4pD\stackrel{~}{\eta }e^{[I(\stackrel{~}{\eta })+{\scriptscriptstyle d^4k𝑑\tau {\scriptscriptstyle \frac{\delta I(\stackrel{~}{\eta })}{\delta \stackrel{~}{\eta }}}\eta _0(k,\tau )}+\mathrm{}]}\times \{\frac{\delta f}{\delta \phi }(p,\tau )& \\ \hfill +d^4qdt\eta _0(q,t)\frac{\delta }{\delta \stackrel{~}{\eta }(q,t)}(\frac{\delta f}{\delta \phi }(p,\tau ))+O(\eta _0^2)\}& \end{array}$$ In case when $`I(\eta _5)=\eta _5K\eta _5`$, where K is some self-adjoint operator, we have $`\eta _0(q,\tau )=d^4k𝑑tG(q,\tau ,k,t)j(k,t)`$ where $`G`$ is the noise two-point function. Substituting into (27) we get: $$\begin{array}{cc}\hfill lim_{j0}\frac{\delta }{\delta j}d^4pD\stackrel{~}{\eta }e^{[I(\stackrel{~}{\eta })+{\scriptscriptstyle d^4q𝑑\tau d^4k𝑑tj(q,\tau )G(q,k,\tau ,t)j(k,t)}]}\{\frac{\delta f}{\delta \phi }(p,\tau )& \\ \hfill +d^4qdt\frac{\delta }{\delta \stackrel{~}{\eta }(q,t)}(\frac{\delta f}{\delta \phi }(p,\tau ))j(q,t)G(p,\tau ,q,t)+O(j^2)\}& \\ \hfill =D\stackrel{~}{\eta }e^{I(\stackrel{~}{\eta })}d^4q𝑑td^4pG(p,\tau ,q,t)\frac{\delta }{\delta \stackrel{~}{\eta }(q,t)}\frac{\delta f}{\delta \phi }(p,\tau )& \\ \hfill =d^4q𝑑td^4pG(p,\tau ,q,t)<\frac{\delta }{\delta \eta (q,t)}\frac{\delta f}{\delta \phi (p,\tau )}>& \\ \hfill =d^4q𝑑td^4pG(p,\tau ,q,t)<\frac{\delta }{\delta \phi (q,t)}\frac{\delta f}{\delta \phi (p,\tau )}>& \end{array}$$ where we used the identity $$\frac{\delta }{\delta \eta (q,t)}=\frac{\delta }{\delta \phi (q,t)}$$ To derive it, note that the general solution to the Langevin equation (18) is given by: $$\phi (k,\tau )=𝑑te^{k^2(t\tau )}\theta (t\tau )\eta (k,t)$$ for which $`\frac{\delta \phi (k,\tau )}{\delta \eta (k,\tau )}=\theta (0)=1`$. \[6\] It is easy to see that in cases when $`I(\eta )`$ is not quadratic, the relation (28) is still true in the WKB limit. Next, writing the average (28) in terms of the functional integral in $`\phi `$ and integrating by parts as in (23) and (24) we obtain $$\begin{array}{cc}\hfill d^4p<\frac{\delta f(\phi ,\tau )}{\delta \phi (p,\tau )}>& \\ \hfill =D\phi f(\phi ,\tau )d^4pd^4q𝑑t\frac{\delta }{\delta \phi (q,t)}G(p,\tau ,q,t)\frac{\delta }{\delta \phi (p,\tau )}P_{FP}(\phi ,\tau )& \end{array}$$ Finally, putting together (23), (24), (31) and using the Langevin equation (18) we obtain the following Fokker-Planck equation for the non-Markovian stochastic process (18), (19): $$\begin{array}{cc}\hfill \frac{dP_{FP}}{d\tau }=d^4pd^4q\frac{\delta }{\delta \phi (q,\tau )}[\frac{\delta S_E}{\delta \phi (p)}P_{FP}(\phi ,\tau )]& \\ \hfill d^4pd^4q𝑑t\frac{\delta }{\delta \phi (q,t)}G(p,\tau ,q,t)\frac{\delta }{\delta \phi (p,\tau )}P_{FP}(\phi ,\tau )& \end{array}$$ To complete our derivation we need to evaluate the two-point noise correlator $`G(p,q,\tau ,t)`$ introduced in (20), which is determined by the worldsheet correlations of the $`V_5`$ vertex in (19). Firstly, let us compute the correlator $`<_{\mathrm{\Lambda }_1}d^2zV_5(z,\overline{z})_{\mathrm{\Lambda }_2}d^2wV_5(w,\overline{w})>`$ of the $`V_5`$-operators entering the definition (19) of the noise. We have: $$\begin{array}{cc}\hfill <\mathrm{\Gamma }^4_{\mathrm{\Lambda }_1}d^2zV_5(z,\overline{z})_{\mathrm{\Lambda }_2}d^2wV_5(w,\overline{w})>=6_{\mathrm{\Lambda }_1}d^2z_{\mathrm{\Lambda }_2}d^2w\frac{1}{\|zw\|^4}& \end{array}$$ where the fourth power of picture-changing operator $`\mathrm{\Gamma }^4`$ has been introduced in the correlator to insure the cancellation of ghost number anomaly on a sphere. The integrals are divergent and their regularized dependence on $`\mathrm{\Lambda }_1`$, $`\mathrm{\Lambda }_2`$ determines the form of noise correlations in stochastic time. Writing $`z=r_1e^{i\alpha _1},w=r_2e^{i\alpha _2}`$ we have: $$\begin{array}{cc}\hfill <_{\mathrm{\Lambda }_1}d^2zV_5(z,\overline{z})_{\mathrm{\Lambda }_2}d^2wV_5(w,\overline{w})>& \\ \hfill =6_{\mathrm{\Lambda }_1}^{\mathrm{}}𝑑r_1r_1_0^{2\pi }𝑑\alpha _1_{\mathrm{\Lambda }_2}^{\mathrm{}}𝑑r_2r_2_0^{2\pi }𝑑\alpha _2\frac{1}{(r_1^2+r_2^22r_1r_2cos(\alpha _1\alpha _2))^2}& \\ \hfill =3\pi ^2_{\mathrm{\Lambda }_1}^{\mathrm{}}𝑑r_1^2_{\mathrm{\Lambda }_2}^{\mathrm{}}𝑑r_2^2\{\frac{1}{(r_1^2r_2^2)^2}+\frac{2r_2^2}{(r_1^2r_2^2)^3}\}& \end{array}$$ The regularized value of the of the first integral gives: $$_{\mathrm{\Lambda }_1}^{\mathrm{}}𝑑r_1^2_{\mathrm{\Lambda }_2}^{\mathrm{}}𝑑r_2^2\frac{1}{(r_1^2r_2^2)^2}=log(\mathrm{\Lambda }_1^2\mathrm{\Lambda }_2^2)$$ while the second one is given by $$_{\mathrm{\Lambda }_1}^{\mathrm{}}𝑑r_1^2_{\mathrm{\Lambda }_2}^{\mathrm{}}𝑑r_2^2\frac{2r_2^2}{(r_1^2r_2^2)^3}=\frac{\mathrm{\Lambda }_2^2}{\mathrm{\Lambda }_1^2\mathrm{\Lambda }_2^2}log(\mathrm{\Lambda }_1^2\mathrm{\Lambda }_2^2)$$ Adding these contributions together we see that the logarithmic terms cancel. Recalling that $`log\mathrm{\Lambda }=\tau `$ where $`\tau `$ is stochastic time we get the following expression for the non-Markovian noise correlator: $$\begin{array}{cc}\hfill <\eta _5(p^{\|\|},\tau _1)\eta _5(q^{\|\|},\tau _2)>\lambda _0^4(1+\frac{\lambda _0^2}{e^{2(\tau _1\tau _2)}1})& \end{array}$$ Substituting this expression into the Green’s function in (32) concludes our derivation of the Fokker-Planck equation associated with the non-Markovian stochastic process (18) in the brane-like sigma-model. In the following section we will study the ansatz solutions to this equation showing their relevance to the $`AdS_5`$ supergravity Hamiltonian. $`AdS_5`$ Supergravity as a solution of the Fokker-Planck Equation Consider the bosonic part of the action for the $`AdS_5`$ supergravity: $$S_{AdS_5}N^2𝑑\lambda d^4x\frac{1}{\lambda ^3}(_\lambda \phi _\lambda \phi +_\mu \phi ^\mu \phi )$$ where $`\phi (\lambda ,x)`$ is the dilaton, $`(\lambda ,x^\mu )`$ are the transverse and longitudinal $`AdS_5`$ coordinates ( $`\mu =0,\mathrm{}3`$). The gauge-theoretic $`N`$ parameter entering (38) is related to the $`AdS_5`$ radius $`R`$ and ten-dimensional gravitational constant $`\kappa `$ through $`\frac{R^8}{\kappa ^2}N^2`$. As is easy to check, the action (38) corresponds to the $`AdS_5`$ metric with normalization $$ds^2=\frac{N^{\frac{4}{3}}}{\lambda ^2}(d\lambda d\lambda +dx_\mu dx^\mu )$$ Substituting $`t=N^{\frac{2}{3}}log\lambda `$ we can write the metric as $$ds^2=dtdt+N^{\frac{4}{3}}e^{2tN^{\frac{2}{3}}}dx_\mu dx^\mu $$ In the new coordinates corresponding to the temporal gauge, the gravity action becomes: $$S_{AdS_5}=N^{\frac{8}{3}}𝑑td^4xe^{4tN^{\frac{2}{3}}}(_t\phi _t\phi +N^{\frac{4}{3}}e^{2tN^{\frac{2}{3}}}_\mu \phi _\mu \phi )$$ The conjugate mimentum $`\pi `$, computed at a constant “time” t is given by $$\pi =\frac{1}{\sqrt{g}}\frac{\delta S_{AdS_5}}{\delta \phi }=_t\phi $$ The corresponding ADM-type Hamiltonian is given by \[11,,12\] $$H_{AdS_5}(t)=d^4x\sqrt{g}(\frac{1}{2}\pi ^2+L)=N^{\frac{8}{3}}d^4xe^{4tN^{\frac{2}{3}}}(\frac{1}{2}_t\phi _t\phi +N^{\frac{4}{3}}e^{2tN^{\frac{2}{3}}}_\mu \phi ^\mu \phi )$$ where $`L=_\mu \varphi ^\mu \varphi `$ is the local Lagrangian density. In the $`(\lambda ,x)`$ coordinates this “equal time slice” Hamiltonian is given by: $$H_{AdS_5}(\lambda )=N^{\frac{4}{3}}d^4x\frac{1}{\lambda ^2}(\frac{1}{2}_\lambda \phi _\lambda \phi +_\mu \phi ^\mu \phi )$$ We shall also need the expression for this Hamiltonian in the $`(\rho ,x)`$ coordinates with $`\rho =\lambda ^2`$: $$H_{AdS_5}(\rho )=N^{\frac{4}{3}}d^4x(2_\rho \phi _\rho \phi +\frac{1}{\rho }_\mu \phi ^\mu \phi )$$ Our strategy now is to relate the radial $`AdS_5`$ coordinate to the stochastic time parameter. This will allow us to interpret the Hamiltonian $`H_{AdS_5}`$ as a logarithm of the Fokker-Planck distribution $`P_{FP}e^{H_{FP}}`$ solving the Fokker-Planck equation (32). However, an important remark should be made here. We know from the holography principle that the radial $`\rho `$ coordinate of $`AdS_5`$ in (45) corresponds to the gauge theory cutoff; therefore if one tries to interpret the AdS/CFT correspondence in terms of the stochastic quantization and the radial $`AdS`$ coordinate as a stochastic time, it is the gauge theoretic cutoff that should enter the Langevin and Fokker-Planck equations solved by the $`AdS_5`$ supergravity Hamiltonian (43), (44). At the same time, however, the stochastic variable in the Langevin equation (18) corresponds to the worldsheet cutoff in string theory. It is clear that the cutoffs in gauge and string theories are not the same; they are related to each other in some non-trivial way. In fact, the knowledge of relation between the string theoretic cutoff to the gauge-theoretic one is necessary to make sense of any possible connections between strings and Yang-Mills theories. In other words, the $`AdS_5`$ gravity Hamiltonian (43), (44) must be a solution to the the Fokker-Planck equation (32) expressed in terms of the gauge theory cutoff (stochastic time) rather than the string-theoretic one. The crucial question is what precisely is the relation between $`\rho `$ and $`\mathrm{\Lambda }_{string}\mathrm{\Lambda }`$. Below we will show that choosing $$\rho =log\mathrm{\Lambda }$$ and expressing the Langevin equation (18), (19) in terms of $`\rho `$ (as well as the associate Fokker-Planck equation) one can see that $`P_{FP}=e^{H_{AdS_5}}`$ is the solution. Consider the Langevin equation (18). Substituting (46) for $`\mathrm{\Lambda }`$ we have: $$\frac{d\phi (p)}{d\rho }=\frac{\delta S_E}{\delta \phi }+\eta _5(p,e^\rho )$$ where $$\eta _5(p,e^{2\rho })\lambda _0^2(1+\lambda _0_{e^\rho }d^2zV_5(z,\overline{z}))$$ and the worldsheet integral is cut off at the scale $`e^{2\rho }`$. To cut the theory at the size of $`\rho `$ corresponding to the the gauge theory cutoff we have to conformally transform $`u=logz`$. Using the conformal covariance of the $`V_5`$ vertex operator we obtain $$\eta _5(p,\rho )\lambda _0^2(1+\lambda _0_\rho d^2uV_5(u,\overline{u}))$$ The Fokker-Planck equation associated with the stochastic process (47) is then given by: $$\begin{array}{cc}\hfill \frac{dP_{FP}}{d\rho }=d^4pd^4q\frac{\delta }{\delta \phi (p,\rho )}[\frac{\delta S_E}{\delta \phi (q,\rho )}P_{FP}(\phi ,\rho )]& \\ \hfill d^4pd^4q𝑑\omega \frac{\delta }{\delta \phi (q,\omega )}G(p,\rho ,q,\omega )\frac{\delta }{\delta \phi (p,\rho )}P_{FP}(\phi ,\rho )& \end{array}$$ where $$G(p,\rho ,q,\omega )=G(p,\rho ,q,\omega )=\lambda _0^4(1+\lambda _0^2\frac{\omega ^2}{\omega ^2\rho ^2})$$ , i.e. it has the same form as the correlator (37) expressed in terms of $`\mathrm{\Lambda }_1`$, $`\mathrm{\Lambda }_2`$. The integral over $`\omega `$ is taken from minus infinity to infinity. We are now ready to check that the Fokker-Planck equation (50) with noise propagator (51) is satisfied by the distribution with Hamiltonian (45). It is convenient to write the Hamiltonian (45) as $$H_{AdS_5}(\rho )=H_1(\rho )+H_2(\rho )$$ where $$\begin{array}{cc}\hfill H_1(\rho )=2N^{\frac{4}{3}}d^4x_\rho \phi _\rho \phi =2N^{\frac{4}{3}}d^4q_\rho \phi (q,\rho )_\rho \phi (q,\rho )& \\ \hfill H_2(\rho )=\frac{N^{\frac{4}{3}}}{\rho }d^4x_\mu \phi ^\mu \phi =\frac{N^{\frac{4}{3}}}{\rho }d^4q\phi (q,\rho )q^2\phi (q,\rho )& \end{array}$$ and the Fourier transform to the momentum space has been performed. We shall use the “constant time slice” variational calculus implying $`\frac{\delta \phi (q_1,\rho _1)}{\delta \phi (q_2,\rho _2)}=\delta ^{(4)}(q_1q_2)\theta _{\rho _1\rho _2}`$ where $`\theta _\rho `$ is the function which is 1 at $`\rho =0`$ and zero elsewhere. This function can be expressed in terms of the usual Heaviside step function as $`\theta _\rho =\theta (\rho )\theta (\rho )`$ Note that the product of two Heaviside functions is a well-defined object in the context of the Colombeau theory of multiplication of distributions \[13\] The useful identity for $`\theta _\rho `$ is $$_{\mathrm{}}^{\mathrm{}}𝑑\rho \frac{\theta _\rho }{\rho }=1$$ which can be easily checked using the integral representation for the Heaviside function: $`\theta (\rho )=_{\mathrm{}}^{\mathrm{}}𝑑x\frac{e^{ix\rho }}{x}`$. It is convenient to represent the r.h.s. of the Fokker-Planck equation (50) (apart from the drift term) as $`\widehat{I_1}+\widehat{I_2}+\mathrm{}+\widehat{I_6}P_{FP}`$ where schematically $$\begin{array}{cc}\hfill \widehat{I_1}=\frac{\delta }{\delta \phi }G\frac{\delta }{\delta \phi }H_1& \\ \hfill \widehat{I_2}=\frac{\delta }{\delta \phi }G\frac{\delta }{\delta \phi }H_2& \\ \hfill \widehat{I_3}=\frac{\delta H_1}{\delta \phi }G\frac{\delta H_1}{\delta \phi }& \\ \hfill \widehat{I_4}=\frac{\delta H_2}{\delta \phi }G\frac{\delta H_2}{\delta \phi }& \\ \hfill \widehat{I_5}=\frac{\delta H_1}{\delta \phi }G\frac{\delta H_2}{\delta \phi }& \\ \hfill \widehat{I_6}=\frac{\delta H_2}{\delta \phi }G\frac{\delta H_1}{\delta \phi }& \end{array}$$ First, let us consider $$\begin{array}{cc}\hfill \widehat{I_2}P_{FP}(\rho )=\{d^4qd^4p𝑑\omega \frac{\delta }{\delta \phi (q,\rho )}G(p,\rho ,q,\omega )\frac{\delta }{\delta \phi (p,\omega )}H_2(\rho )\}P_{FP}(\rho )& \\ \hfill =\{\lambda _0^4N^{\frac{4}{3}}d^4pd^4q\delta ^{(4)}(pq)𝑑\omega \theta _{\omega \rho }(1+\lambda _0^2\frac{p^2}{\omega }\frac{\omega ^2}{\omega ^2\rho ^2})\}P_{FP}(\rho )& \end{array}$$ Using (54) we can evaluate the $`\omega `$ integral in (55) and obtain: $$\widehat{I_1}P_{FP}(\rho )=\frac{N^{\frac{4}{3}}\lambda _0^6}{2}d^4qq^2P_{FP}(\rho )$$ Next, consider the piece $$\begin{array}{cc}\hfill \widehat{I_1}P_{FP}(\rho )=\{d^4qd^4p𝑑\omega \frac{\delta }{\delta \phi (q,\rho )}G(p,\rho ,q,\omega )\frac{\delta }{\delta \phi (p,\omega )}H_1(\rho )\}P_{FP}(\rho )& \end{array}$$ To compute this piece, consider first of all the variation: $$\begin{array}{cc}\hfill \frac{\delta H_1(\rho )}{\delta \phi (q,\rho )}P_{FP}(\rho )=N^{\frac{4}{3}}\frac{\delta }{\delta \phi (q,\rho )}d^4p𝑑u_u\phi (p,u)_u\phi (p,u)\delta (\rho u)P_{FP}(u)& \\ \hfill =N^{\frac{4}{3}}\frac{\delta }{\delta \phi (q,\rho )}d^4p𝑑u\phi (p,u)_u[\delta (\rho u)_u\phi (p,u)P_{FP}(u)]=& \\ \hfill N^{\frac{4}{3}}𝑑u\theta _{\rho u}_u[\delta (\rho u)_u\phi (q,u)P_{FP}(u)]& \\ \hfill =N^{\frac{4}{3}}𝑑u\phi (q,u)_u[\delta (\rho u)_u\theta _{\rho u}P_{FP}(u)]& \end{array}$$ Next, performing the second variation and integrating by parts we have: $$\begin{array}{cc}\hfill \widehat{I_1}P_{FP}(\rho )=N^{\frac{4}{3}}𝑑u𝑑\omega (\lambda _0^4+\lambda _0^6\frac{\omega ^2}{\omega ^2\rho ^2})\theta _{\omega u}_u[\delta (\rho u)_u\theta _{\rho u}P_{FP}(u)]& \\ \hfill =N^{\frac{4}{3}}𝑑u𝑑\omega (\lambda _0^4+\frac{\lambda _0^6\omega ^2}{\omega ^2\rho ^2})_u\theta _{\omega u}\delta (\rho u)_u\theta _{\rho u}P_{FP}(u)& \\ \hfill =N^{\frac{4}{3}}𝑑w(\lambda _0^4+\frac{\lambda _0^6\omega ^2}{\omega ^2\rho ^2})_\omega J(\rho ,\omega )& \end{array}$$ where we used $`_u\theta _{\omega u}=_\omega \theta _{\omega u}`$ and denoted $`J(\rho ,\omega )=𝑑u\theta _{\omega u}\delta (\rho u)_u\theta _{\rho u}P_{FP}(u)`$. Next, integrating by parts again and using the above identity for the derivatives of $`\theta `$ we obtain: $$\begin{array}{cc}\hfill J(\rho ,\omega )=𝑑u_u\theta _{\rho u}\theta _{\omega u}\delta (\rho u)P_{FP}(u)=𝑑u_\omega \theta _{\rho u}\theta _{\omega u}\delta (\rho u)P_{FP}(u)& \\ \hfill =_\rho 𝑑u\theta _{\rho u}\theta _{\omega u}\delta (\rho u)P_{FP}(u)+𝑑u\theta _{\rho u}\theta _{\omega u}(\delta (\rho u))^{}P_{FP}(u)& \\ \hfill =_\rho (\theta _{\omega \rho }P_{FP}(\rho ))_\rho (\theta _{\omega \rho }P_{FP}(\rho ))=2_\rho (\theta _{\omega \rho }P_{FP}(\rho ))& \end{array}$$ Substituting $`J(\rho ,\omega )`$ into $`\widehat{I_1}P_{FP}(\rho )`$ we obtain: $$\begin{array}{cc}\hfill \widehat{I_1}P_{FP}(\rho )=2N^{\frac{4}{3}}𝑑\omega (\lambda _0^4+\frac{\lambda _0^6\omega ^2}{\omega ^2\rho ^2})_\omega _\rho (\theta _{\omega \rho }P_{FP}(\rho ))& \\ \hfill =4N^{\frac{4}{3}}\lambda _0^6\rho ^2_\rho P_{FP}(\rho )𝑑\omega \frac{\omega \theta _{\omega \rho }}{(\rho ^2\omega ^2)^2}+4N^{\frac{4}{3}}\lambda _0^6\rho ^2P_{FP}(\rho )𝑑\omega \frac{\omega _\omega \theta _{\omega \rho }}{(\rho ^2\omega ^2)^2}& \\ \hfill =4N^{\frac{4}{3}}\lambda _0^6\rho ^2_\rho P_{FP}(\rho )_\rho (\frac{1}{4\rho })4N^{\frac{4}{3}}\lambda _0^6\rho ^2P_{FP}(\rho )𝑑\omega \theta _{\omega \rho }(\frac{1}{(\rho ^2\omega ^2)^2}+\frac{4\omega ^2}{(\rho ^2\omega ^2)^3})& \\ \hfill =N^{\frac{4}{3}}\lambda _0^6_\rho P_{FP}(\rho )+N^{\frac{4}{3}}\lambda _0^6\rho ^2(\frac{1}{2\rho ^3}+\frac{1}{2\rho ^3})P_{FP}(\rho )\lambda _0^6_\rho P_{FP}(\rho )& \end{array}$$ To summarize, $$\widehat{I_1}P_{FP}(\rho )=N^{\frac{4}{3}}\lambda _0^6P_{FP}(\rho )$$ Next, consider $$\begin{array}{cc}\hfill \widehat{I_3}P_{FP}(\rho )=d^4pd^4q𝑑\omega \frac{\delta H_1(\rho )}{\delta \phi (\rho )}G(p,q,\rho ,\omega )\frac{\delta H_1(\rho )}{\delta \phi (\omega )}P_{FP}(\rho )& \\ \hfill =\lambda _0^6N^{\frac{8}{3}}d^4pd^4q𝑑\omega \frac{\omega ^2}{\omega ^2\rho ^2}_\rho ^2\phi (p,\rho )_\omega (\theta _{\omega \rho }_\omega \phi (q,\rho ))P_{FP}(\rho )& \\ \hfill =\lambda _0^6N^{\frac{8}{3}}d^4pd^4q\{d\omega \frac{\omega ^2}{\omega ^2\rho ^2}\theta _{\omega \rho }_\rho ^2\phi (p,\rho )_\omega ^2\phi (q,\omega )& \\ \hfill d\omega \theta _{\omega \rho }_\omega (\frac{\omega ^2}{\omega ^2\rho ^2})_\rho ^2\phi (p,\rho )_\omega \phi (q,\omega )\}P_{FP}(\rho )& \\ \hfill =\lambda _0^6N^{\frac{8}{3}}d^4pd^4q(\frac{\rho }{2}_\rho ^2\phi (p,\rho )\frac{\rho }{2}_\rho ^2\phi (q,\rho ))_\rho ^2\phi (p,\rho )P_{FP}(\rho )=0& \end{array}$$ The next piece is given by $$\begin{array}{cc}\hfill \widehat{I_4}P_{FP}(\rho )=d^4pd^4q𝑑\omega (\lambda _0^4+\lambda _0^6\frac{\omega ^2}{\omega ^2\rho ^2})\frac{\delta H_2(\rho )}{\delta \phi (p,\rho )}\frac{\delta H_2(\rho )}{\delta \phi (q,\omega )}P_{FP}(\rho )& \\ \hfill =\lambda _0^6N^{\frac{8}{3}}d^4pd^4q\frac{p^2q^2}{\rho }\phi (p,\rho )𝑑\omega \frac{\theta _{\omega \rho }\omega \phi (q,\omega )}{(\omega \rho )(\omega +\rho )}P_{FP}(\rho )& \\ \hfill =\frac{\lambda _0^6N^{\frac{8}{3}}}{2\rho }d^4pd^4qp^2q^2\phi (p,\rho )\phi (q,\rho )P_{FP}(\rho )& \end{array}$$ Next, $$\begin{array}{cc}\hfill \widehat{I_5}P_{FP}(\rho )=d^4pd^4q𝑑\omega (\lambda _0^4+\lambda _0^6\frac{\omega ^2}{\omega ^2\rho ^2})\frac{\delta H_1(\rho )}{\delta \phi (\rho )}\frac{\delta H_2(\rho )}{\delta \phi (\omega )}P_{FP}(\rho )& \\ \hfill =\lambda _0^6N^{\frac{8}{3}}d^4pd^4q𝑑\omega \frac{\omega ^2}{\omega ^2\rho ^2}_\rho ^2\phi (p,\rho )\frac{q^2}{\omega }\phi (q,\omega )P_{FP}(\rho )& \\ \hfill =\frac{\lambda _0^6N^{\frac{8}{3}}}{2}d^4pd^4qq^2\phi (q,\rho )_\rho ^2\phi (p,\rho )P_{FP}(\rho )& \end{array}$$ At last, $$\begin{array}{cc}\hfill \widehat{I_6}P_{FP}(\rho )=d^4pd^4q𝑑\omega (\lambda _0^4+\lambda _0^6\frac{\omega ^2}{\omega ^2\rho ^2})\frac{\delta H_2(\rho )}{\delta \phi (\rho )}\frac{\delta H_1(\rho )}{\delta \phi (\omega )}P_{FP}(\rho )& \\ \hfill =N^{\frac{8}{3}}\lambda _0^6d^4pd^4q𝑑\omega \frac{p^2\phi (p,\rho )}{\rho }_\omega (\theta _{\omega \rho }_\omega \phi (q,\omega ))P_{FP}(\rho )=0& \end{array}$$ Finally, using (55)-(67) we find that substituting the $`AdS_5`$ supergravity Hamiltonian to the “kinetic” part of the r.h.s. of the Fokker-Planck equation (50) (i.e. the r.h.s. without a drift term) is given by: $$\begin{array}{cc}\hfill d^4pd^4qd\omega \frac{\delta }{\delta \phi (q,\omega )}G(p,\rho ,q,\omega )\frac{\delta }{\delta \phi (p,\rho )}P_{FP}(\rho )=N^{\frac{4}{3}}\lambda _0^6\{\frac{}{\rho }+\frac{1}{2}d^4qq^2& \\ \hfill +\frac{N^{\frac{4}{3}}}{2\rho }d^4pd^4qp^2\phi (p,\rho )q^2\phi (q,\rho )+\frac{N^{\frac{4}{3}}}{2}d^4pd^4q_\rho ^2\phi (p,\rho )q^2\phi (q,\rho )\}P_{FP}(\rho )& \end{array}$$ The next step is to evaluate the drift term of the Fokker-Planck equation (50). First of all, we have to point out the correct form and normaliaztion of the equilibrium action $`S_E`$. As the stochastic time $`\rho `$ corresponds to the radial $`AdS_5`$ coordinate, orthogonal to the worldvolume of the underlying D3-brane, the stochastic process (18), (19) describes the evolution of an “equal time slice” gravity Hamiltonian as it drifts from infinity towards the $`D3`$-brane worldvolume. The state of thermal equilibrium, reached in the limit of infinite stochastic time, is described by the free field scalar theory in four dimensions since in the region far from the D3-brane the space-time is flat. On the other hand, action in the drift term of the Langevin equation (18) is the one of the scalar field theory in $`ten`$ dimensions, i.e. the dilaton effective action (in the Einstein frame) in the absence of the graviton mode and the $`V_5`$ term. Note, however, that truncation of the $`\lambda _0`$ expansion in the brane-like sigma-model, performed in (9), corresponds to the dilaton s-wave approximation in the $`AdS_5\times S^5`$ picture, as can be shown by considering sigma-model correlation functions and comparing them to those of the $`D=4`$ super Yang-Mills theory refs. In the s-wave approximation the $`AdS_5\times S^5`$ supergravity is reduced to the $`AdS_5`$ one with the action (38) while the ten-dimensional free field dilaton action entering the Langevin equation (18), given by $$S_{10}\frac{1}{\kappa ^2}d^{10}x_\alpha \phi ^\alpha \phi $$ can be reduced to the five-dimensional one given by $$\begin{array}{cc}\hfill S_5\frac{R^5}{\kappa ^2}d^4x𝑑\rho (_\mu \phi ^\mu \phi +_\rho \phi _\rho \phi )& \\ \hfill =\frac{N^2}{R^3}d^4x𝑑\rho (_\mu \phi ^\mu \phi +_\rho \phi ^\rho \phi )& \end{array}$$ Of course, higher order $`\lambda _0`$ corrections (corresponding to dilaton partial waves with higher angular momenta) modify the kinetic part of the Fokker-Planck equation (50) and the full ten-dimensional equilibrium action must then be used. The ADM-type “equal time slice” four-dimensional Hamiltonian may be constructed out of this action precisely as in the $`AdS_5`$ case; in the limit $`\rho \mathrm{}`$ it is given by $$S_E=\frac{N^{\frac{4}{3}}}{R^2}d^4x_\mu \phi ^\mu \phi $$ This gives us the expression for the four-dimensional equilibrium distribution in the Fokker-Planck equation (50) obtained from the full ten-dimensional dilaton action in the Langevin equation (18) in the s-wave approximation. Physically, it describes the dilaton field living in the four-dimensional subspace infinitely far from the D3-brane. Knowing the equilibrium action $`S_E`$, we may now carry out the calculation of the drift term in (50): $$\begin{array}{cc}\hfill d^4pd^4q\frac{\delta }{\delta \phi (p,\rho )}(\frac{\delta S_E}{\delta \phi (q,\rho )}P_{FP}(\rho ))& \\ \hfill =d^4pd^4q\{\frac{\delta ^2S_E}{\delta \phi (p,\rho )\delta \phi (q,\rho )}P_{FP}(\rho )+\frac{\delta S_E}{\delta \phi (p,\rho )}\frac{\delta (H_1+H_2)}{\delta \phi (q,\rho )}P_{FP}(\rho )\}& \\ \hfill =\{N^{\frac{4}{3}}d^4qq^2+\frac{N^{\frac{8}{3}}}{\rho }d^4pd^4qp^2\phi (p,\rho )q^2\phi (q,\rho )& \\ \hfill +N^{\frac{8}{3}}d^4pd^4q_\rho ^2\phi (p,\rho )q^2\phi (q,\rho )\}P_{FP}(\rho )& \end{array}$$ Finally, substituting (68) and (72) into the right-hand side of the Fokker-Planck equation (50) we get: $$\begin{array}{cc}\hfill \widehat{H}P_{FP}(\rho )d^4pd^4q\frac{\delta }{\delta \phi (p,\rho )}[\frac{\delta S_E}{\delta \phi (q,\rho )}P_{FP}(\rho )]& \\ \hfill d^4pd^4q𝑑\omega \frac{\delta }{\delta \phi (q,\omega )}G(p,\rho ,q,\omega )\frac{\delta }{\delta \phi (p,\rho )}P_{FP}(\rho )& \\ \hfill =\lambda _0^6N^{\frac{4}{3}}_\rho P_{FP}(\rho )+\{(\frac{N^{\frac{4}{3}}}{R^2}\lambda _0^6N^{\frac{4}{3}})d^4pp^2& \\ \hfill (\lambda _0^6N^{\frac{8}{3}}\frac{N^{\frac{8}{3}}}{R^2})d^4pd^4q[p^2\phi (p,\rho )(q^2\phi (q,\rho )+_\rho ^2\phi (q,\rho ))]\}P_{FP}(\rho )& \end{array}$$ Choosing $`\lambda _0^6=R^2`$ (i.e. proportional to the curvature of the AdS space) and rescaling the stochastic time $`\rho `$ parameter as $`\rho =\lambda _0^6N^{\frac{4}{3}}s`$ we see that all the terms in the right hand side of (73) cancel out, except for the one proportional to the derivative of $`P_{FP}(\rho )e^{H_{AdS_5}(\rho )}`$: $$\widehat{H}P_{FP}(s)=\frac{}{s}P_{FP}(s)$$ This constitutes the proof that the $`AdS_5`$ gravity indeed defines the distribution satisfying the Fokker-Planck equation associated with the stochastic process (18) with the non-Markovian $`V_5`$-noise (19). Discussion Our results have the following geometrical interpretation. In the beginning ($`\rho \mathrm{}`$) we have an observer living in asymptotically flat four-dimensional space-time. At certain moment the D3-brane, or the $`V_5`$ vertex, located infinitely far from the observer, is introduced. The gravitational force of the brane causes the drift of the four-dimensional space-time slice, in which the observer lives, towards the D3-brane location. From the point of view of the observer, however, the process of the four dimensional space-time slice deformation is perceived as the non-Markovian stochastic process with the white noise propagating in the additional fictitious dimension. However, what perceived by the observer as a stochastic drift in the fictitious time dimension is in fact the gravitational drift in the $`real`$ dimension, that is, the coordinate transverse to the $`D3`$-brane worldvolume. The equilibrium state in the limit of $`\rho \mathrm{}`$ corresponds to reaching the D3-brane horizon by the observer. In this limit the $`V_5`$ noise propagator (51) becomes $$lim_\rho \mathrm{}G(p,\rho ,q,\omega )=\lambda _0^4+\lambda _0^6$$ and in this limit the kinetic term of the Fokker-Planck equation (50) vanishes. The only contribution comes then from the drift term and the Fokker-Planck equation (50) is reduced to $`p^2P_{FP}(\phi )0`$ or $$\frac{^2}{X_\mu X^\mu }P_{FP}0$$ in the position space. The structure of the last equation is similar to the regularized large N loop equation $$\begin{array}{cc}\hfill \widehat{L}W(C)=0& \end{array}$$ with $$\widehat{L}lim_{ϵ0}_ϵ^ϵ𝑑\alpha \frac{\delta ^2}{\delta X_\mu (s\alpha )\delta X^\mu (s+\alpha )}$$ Therefore in the limit of infinite stochastic time the Fokker-Planck distribution $`e^{H_{AdS_5}(\rho )}`$ solving the equation (50) is reduced to the Wilson loop of the large N gauge theory living in the D3-brane worldvolume. From point of view of the Langevin equation (18) reaching the equilibrium limit corresponds to restoring the worldsheet conformal invariance (or cutoff independence). Our final remark regards higher order $`\lambda _0`$ corrections in the brane-like sigma-model (9). As has been said before, this four-dimensional picture is only appropriate in the dilaton s-wave approximation. Higher order $`\lambda _0`$ corrections correspond to higher angular momentum modes of the dilaton on $`S^5`$, therefore the initial equilibrium distribution in (18) must be taken ten-dimensional. Just as described above, the stochastic process is non-Markovian and the equilibrium limit reached at infinite stochastic time is different from the initial one and it corresponds to the $`AdS_5\times S^5`$ gravity. Summing up to all orders of $`\lambda _0`$ one should obtain string theory compactified on the $`AdS_5\times S^5`$ space. Thus flat ten-dimensional string theory and the one in the $`AdS_5\times S^5`$ space-time appear as two $`different`$ thermodynamical limits of the stochastic process caused by the non-Markovian $`V_5`$-noise. This picture offers the dynamical explanation of how the $`AdS_5\times S^5`$ space-time geometry emerges in the brane-like sigma-model (5) due to the presence of the 5-form $`V_5`$ vertex operator. Another question of interest is the loop equation which open string theory in the $`AdS_5\times S^5`$ background should satisfy \[14\]. Though Wilson loops have been studied extensively in the context of AdS/CFT correspondence (for review see \[15\]), it is not clear that the functional of the $`AdS_5`$ supergravity satisfies the loop equation. So far the zig-zag invariance of the $`AdS_5`$ string has been proved in the WKB limit and for the special class of contours only \[14\]. On the other hand, it has been argued \[16\] that the loop equation may be understood as an equilibrium condition for the Fokker-Planck equation with the white noise. It seems hard, however, to draw any straightforward connection between Fokker-Planck and Schwinger-Dyson equations and, in particular, to interpret the loop equation as an infinite time limit of the $`<H_{FP}(\tau )W(C)>=0`$ relation discussed in \[16,,17,,18\]. The reason is that the loop equation is derived in the regularized theory while no corresponding regularization is performed in the Fokker-Planck counterpart. Therefore in order to to make sense of possible connections between Wheeler-de Witt and Schwinger-Dyson equations one has to obtain more information about string and gauge degrees of freedom and their interrelation. In particular, it is necessary to point out how string and gauge theory cutoffs correspond to each other. We hope that our approach, attempting to explore the holographic effects from brane-like sigma-model point of view, will be helpful in trying to find answers to these questions. Acknowledgements It is a pleasure to thank the Physics Department of Tel Aviv University and particularly J.Sonnenschein for their gracious hospitality. The author also gratefully acknowledges the very kind hospitality of the Erwin Schroedinger International Institute for Mathematical Physics as well as the organizers of the workshop “Dualities in String Theory” at the ESI in Vienna. In particular, it is a pleasure to thank H. Grosse and S. Theisen. I also would like to thank H.Huffel, B.Kol, M. Sheikh-Jabbari, J.Sonnenschein and S.Theisen for useful discussions. This work is partially supported by the EC contract no. ERBFMRX-CT 96-0090. 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warning/0005/cond-mat0005003.html	ar5iv	text	# Tunneling Dynamics of Bose-Einstein Condensates with Feshbach Resonances \[ ## Abstract We study tunneling dynamics of atomic pairs in Bose-Einstein condensates with Feshbach resonances. It is shown that the tunneling of the atomic pairs depends on not only the tunneling coupling between the atomic condensate and the molecular condensate, but also the inter-atomic nonlinear interactions and the initial number of atoms in these condensates. It is found that in addition to oscillating tunneling current between the atomic condensate and the molecular condensate, the nonlinear atomic-pair tunneling dynamics sustains a self-locked population imbalance: macroscopic quantum self-trapping effect. Influence of decoherence induced by non-condensate atoms on tunneling dynamics is investigated. It is shown that decoherence suppresses atomic-pair tunneling. \] The atomic Bose-Einstein condensates offer new opportunities for studying quantum-degenerate fluids. All the essential properties of atomic Bose-Einstein condensed systems are determined by the strength of the atomic interactions. In contrast with the situation of the traditional superfluids, the strength of the inter-particle interactions in the atomic Bose-Einstein condensate can vary over a wide range of values through changing external fields. Hence one can manipulate and control condensate properties by varying the strength of interactions. The Feshbach resonance approach is considered as an effective one to alter the inter-atomic interactions in Bose-Einstein condensates. The magnetic-field-induced Feshbach resonances in an atomic Bose condensate have already been observed experimentally . Theoretical studies of the ultracold atoms with Feshbach resonances showed that the interactions responsible for the Feshbach resonances produce a second condensate component, a molecular condensate, and predict that tunneling of atomic pairs occurs between the atomic condensate and the molecular condensate. Recently, the molecular Bose-Einstein condensate has been produced experimentally . The purpose of this paper is to study tunneling dynamics of atomic pairs between the atomic condensate and the molecular condensate. We show that in addition to oscillating tunneling current between the atomic condensate and the molecular condensate, the nonlinearity of the tunneling dynamics sustains a self-maintained population imbalance: macroscopic quantum self-trapping effect (MQST). We also discuss the influence of decoherence induced by non-condensate atoms on tunneling dynamics and find that decoherence suppresses atomic-pair tunneling. The binary atom Feshbach resonances are hyperfine-induced spin-flip processes that bring the colliding atoms to a bound molecular state of different spins, and then return an unbound state. These processes can be described by Hamiltonian ($`\mathrm{}=1`$) $$\widehat{H}_{FR}=\alpha 𝑑𝐫\widehat{\psi }_m^+(𝐫)\widehat{\psi }_a(𝐫)\widehat{\psi }_a(𝐫)+h.c.,$$ (1) where $`\widehat{\psi }_m(𝐫)`$, $`\widehat{\psi }_m^+(𝐫)`$ ($`\widehat{\psi }_a𝐫)`$, $`\widehat{\psi }_a^+(𝐫)`$) are the annihilation and creation field operators of the molecules (atoms), $`\alpha `$ stands for the coupling constant. The Hamiltonian $`\widehat{H}_{FR}`$ together with atomic, molecular, and atom-molecule interaction Hamiltonians ($`\mathrm{}=1`$) $`\widehat{H}_a`$ $`=`$ $`{\displaystyle 𝑑𝐫\widehat{\psi }_a^+(𝐫)[\frac{1}{2M}^2+V(𝐫)]\widehat{\psi }_a(𝐫)}`$ (3) $`+{\displaystyle \frac{\lambda _a^{}}{2}}{\displaystyle 𝑑𝐫\widehat{\psi }_a^+(𝐫)\widehat{\psi }_a^+(𝐫)\widehat{\psi }_a(𝐫)\widehat{\psi }_a(𝐫)},`$ $`\widehat{H}_m`$ $`=`$ $`{\displaystyle 𝑑𝐫\widehat{\psi }_m^+(𝐫)[\frac{1}{4M}^2+V(𝐫)+ϵ]\widehat{\psi }_m(𝐫)}`$ (5) $`+{\displaystyle \frac{\lambda _m^{}}{2}}{\displaystyle 𝑑𝐫\widehat{\psi }_m^+(𝐫)\widehat{\psi }_m^+(𝐫)\widehat{\psi }_m(𝐫)\widehat{\psi }_m(𝐫)},`$ $`\widehat{H}_{am}`$ $`=`$ $`\lambda ^{}{\displaystyle 𝑑𝐫\widehat{\psi }_a^+(𝐫)\widehat{\psi }_m^+(𝐫)\widehat{\psi }_a(𝐫)\widehat{\psi }_m(𝐫)},`$ (6) forms a total Hamiltonian $$\widehat{H}=\widehat{H}_a+\widehat{H}_m+\widehat{H}_{am}+\widehat{H}_{FR},$$ (7) which governs the dynamics of the system under our consideration. Here $`V(𝐫)`$ represents the trapped potential, $`\lambda _{a(m)}^{}=4\pi a_{a(m)}/((2)M)`$ with $`M`$ being the atomic mass and $`a_{a(m)}`$ the scattering length, $`\lambda ^{}`$ denotes the coupling constant of atom-molecule interaction. The detuning $`ϵ`$ linearly depends on the magnetic field $`ϵBB_0`$ with $`B_0`$ being the resonant magnetic field. For small atomic and molecular condensates , the atomic and molecular field operators can be approximated as $`\widehat{\psi }_a(𝐫)=\widehat{a}\varphi _a(𝐫)`$, $`\widehat{\psi }_m(𝐫)=\widehat{b}\varphi _b(𝐫)`$ where $`\varphi _a(𝐫)`$ and $`\varphi _b(𝐫)`$ are real normalized mode functions for the two condensates, and $`\widehat{a}`$ and $`\widehat{b}`$ are associated mode annihilation operators which satisfy the standard bosonic communtation relations. Then the total Hamiltonian becomes the two-mode Hamiltonian $`\widehat{H}`$ $`=`$ $`\widehat{H}_0+\widehat{H}^{},`$ (8) $`\widehat{H}_0`$ $`=`$ $`\omega _a\widehat{a}^+\widehat{a}+\omega _b\widehat{b}^+\widehat{b}+\lambda _a\widehat{a}^{+2}\widehat{a}^2+\lambda _b\widehat{b}^{+2}\widehat{b}^2+\lambda \widehat{a}^+\widehat{a}\widehat{b}^+\widehat{b},`$ (9) $`\widehat{H}^{}`$ $`=`$ $`\alpha (\widehat{b}^+\widehat{a}^2+\widehat{b}\widehat{a}^{+2}),`$ (10) where $`\alpha \widehat{b}^+\widehat{a}^2`$ describes the annihilation of a atomic pair in the atomic condensate and the creation of one molecule in the molecular condensate thereby transferring a pair of atoms from the atomic condensate to the molecular condensate with $`\alpha `$ being the corresponding tunneling coupling constant. The Hermitian conjugate part $`\alpha \widehat{b}\widehat{a}^{+2}`$ describes the reverse process. Hence, what the Feshbach Hamiltonian $`\widehat{H}^{}`$ describes is not only a three-body recombination process of the molecular formation from a pair of atoms but also a tunneling process between the atomic condensate and the molecular condensate. In general, the Hamiltonian (6) can not be exactly solved, but it can be perturbatively solved in the off-resonant regime with a large detuning $`ϵ0`$, by which we mean here that $`ϵ`$ greatly exceeds the Feshbach-resonant interaction energy, so we can treat $`H^{}`$ as a perturbation. Let $`\widehat{H}_0\|n,m=E_{n,m}^{(0)}\|n,m`$ where $`\|n,m`$ is an eigenstate of the number operators $`\widehat{a}^+\widehat{a}`$ and $`\widehat{b}^+\widehat{b}`$ defined by $`\widehat{a}^+\widehat{a}\|n,m=n\|n,m`$ and $`\widehat{b}^+\widehat{b}\|n,m=m\|n,m`$. It is easy to find that $`E_{n,m}^{(0)}`$ $`=`$ $`n\omega _a+m\omega _b+(n^2n)\lambda _a`$ (12) $`+(m^2m)\lambda _b+nm\lambda .`$ For simplicity, we consider the nondegenerate case and assume that $`\widehat{H}\|\psi _{n,m}=E_{n,m}\|\psi _{n,m}`$. Then the perturbative energy and eigenstate are found to be $`E_{n,m}`$ $``$ $`E_{n,m}^{(0)}+\alpha [a_{n,m}\sqrt{n(n1)(m+1)}`$ (14) $`+b_{n,m}\sqrt{(n+1)(n+2)m}],`$ $`\|\psi _{n,m}`$ $``$ $`A_{n,m}[\|n,m+a_{n,m}\|n2,m+1`$ (16) $`+b_{n,m}\|n+2,m1],`$ where the normalization constant $`A_{n,m}`$ and two coefficients $`a_{n,m}`$ and $`b_{n,m}`$ are given by $`A_{n,m}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{1+a_{n,m}^2+b_{n,m}^2}}},`$ (17) $`a_{n,m}`$ $`=`$ $`{\displaystyle \frac{\alpha \sqrt{n(n1)(m+1)}}{E_{n,m}^{(0)}E_{n2,m+1}^{(0)}}},`$ (18) $`b_{n,m}`$ $`=`$ $`{\displaystyle \frac{\alpha \sqrt{(n+1)(n+2)m}}{E_{n,m}^{(0)}E_{n+2,m1}^{(0)}}},`$ (19) Let the two condensates be initially in a state $`\|\mathrm{\Psi }(0)=_{n,m}C_{n,m}(0)\|\psi _{n,m}`$, the time evolution of the wave function is then given by the expression $$\|\mathrm{\Psi }(t)\underset{n,m}{}C_{n,m}(0)e^{iE_{n,m}t}\|\psi _{n,m}.$$ (20) In order to investigate tunneling dynamics, we introduce the population difference $$P(t)=n_a(t)2n_b(t),$$ (21) where $`n_{a(b)}(t)`$ is the number of atoms (molecules) in the atomic (molecular) condensate at time $`t`$. Now let us assume that initially both the atomic and molecular condensates are in a Fock state $`\|n_1,n_2`$. From Eqs. (11) and (12) we find that $$P(t)=P_0(n_1,n_2)+\underset{i=1}{\overset{3}{}}P_i(n_1,n_2)\mathrm{cos}[\omega _i(n_1,n_2)t],$$ (22) where the oscillation frequencies are given by $`\omega _1(n_1,n_2)`$ $`=`$ $`E_{n_1,n_2}E_{n_1+2,n_21},`$ (23) $`\omega _2(n_1,n_2)`$ $`=`$ $`E_{n_1,n_2}E_{n_12,n_2+1},`$ (24) $`\omega _3(n_1,n_2)`$ $`=`$ $`E_{n_1+2,n_21}E_{n_12,n_2+1},`$ (25) and the coefficients $`P_i(n_1,n_2)`$ are defined by $`P_0(n_1,n_2)`$ $`=`$ $`A_{n_1,n_2}^2[(n_12n_2)+(n_12n_24)a_{n_1,n_2}^2`$ (33) $`+(n_12n_2+4)b_{n_1,n_2}^2]`$ $`+A_{n_1+2,n_21}^2a_{n_1+2,n_21}^2[(n_12n_2+4)`$ $`+(n_12n_2)a_{n_1+2,n_21}^2`$ $`+(n_12n_2+8)b_{n_1+2,n_21}^2]`$ $`+A_{n_12,n_2+1}^2b_{n_12,n_2+1}^2[(n_12n_24)`$ $`+(n_12n_26)a_{n_12,n_2+1}^2`$ $`+(n_12n_2)b_{n_12,n_2+1}^2],`$ $`P_1(n_1,n_2)`$ $`=`$ $`2A_{n_1,n_2}A_{n_1+2,n_21}a_{n_1+2,n_21}`$ (36) $`\times [(n_12n_2)a_{n_1+2,n_21}`$ $`+(n_12n_2+4)b_{n_1,n_2}],`$ $`P_2(n_1,n_2)`$ $`=`$ $`2A_{n_1,n_2}A_{n_12,n_2+1}b_{n_12,n_2+1}`$ (39) $`\times [(n_12n_24)a_{n_1,n_2}`$ $`+(n_12n_2)b_{n_12,n_2+1}],`$ $`P_3(n_1,n_2)`$ $`=`$ $`2A_{n_1+2,n_21}A_{n_12,n_2+1}a_{n_1+2,n_21}^2`$ (41) $`\times b_{n_12,n_2+1}^2(n_12n_2).`$ From Eq. (12) we see that the population difference between the atomic condensate and the molecular condensate exhibits oscillating behaviors. Especially, we can obtain the nonzero time average of the population difference labeled by $`\overline{P}`$ given by $$\overline{P}=P_0(n_1,n_2),$$ (42) which implies that there is a self-locked population imbalance between the atomic condensate and the molecular condensate. This is the MQST which occurs in the usual two Bose condensate system as well. It is easy to check that the MQST vanishes when the nonlinearities in interactions are absent. Hence the MQST phenomenon is a nonlinear effect. Eq. (12) indicates that the population imbalance between the atomic condensate and the molecular condensate exhibits nonlinear oscillations with the time evolution. It is these oscillations that leads to a Josephson-like tunneling current between the atomic condensate and molecular condensate, which can be defined as $`I(t)=\dot{P}(t)/N`$ with $`N`$ being the total number of atoms in the atomic and molecular condensates. Making use of Eq. (13) it is easy to find that $$I(t)=\underset{i=1}{\overset{3}{}}\frac{P_i(n_1,n_2)\omega _i(n_1,n_2)}{n_1+2n_2}\mathrm{sin}[\omega _i(n_1,n_2)t].$$ (43) From Eqs. (7)-(16) and (18) we see that the tunneling of the atomic pairs depends on not only the tunneling coupling between the atomic condensate and the molecular condensate but also the inter-atomic nonlinear interactions and the initial number of atoms in these condensates. In order to further understand the influence of inter-atomic interactions on tunneling dymamics of the system under our consideration, let us specialize to the case of the atomic condensate initially being in a number state $`\|N`$, and the molecular condensate initially being unpopulated, i.e., $`\|\mathrm{\Psi }(0)=\|N,0`$. In this case, the population difference and the tunneling current between the two condensates given by the expressions $`P(t)`$ $`=`$ $`P_0(N,0)+P_2(N,0)\mathrm{cos}[\omega _2(N,0)t],`$ (44) $`I(t)`$ $`=`$ $`I_{am}\mathrm{sin}[\omega _2(N,0)t],`$ (45) where we have set $`I_{am}=P_2(N,0)\omega _2(N,0)/N`$. It is easy to see that the amplitude $`I_{am}`$ and the frequency $`\omega _2(N,0)`$ depend upon the initial number of atoms in the condensates, the tunneling coupling $`\alpha `$, the nonlinear interaction strengths $`\lambda _a`$, $`\lambda _b`$, and $`\lambda `$. To find out how the initial number of atoms in the condensates and the interaction strengths affect the tunneling current, in Figure 1 we plot the amplitude of the tunneling current as a function of interaction strengths for different initial number of atoms when the initial state is $`\|\mathrm{\Psi }(0)=\|N,0`$, $`\lambda _a=\lambda _b=\lambda `$, and $`\omega _a=2\omega _b`$. Figure 1 indicates that the amplitude of of the tunneling current is almost independent of the initial number of atoms, the tunneling coupling, and inter-atomic nonlinear interactions in the regime of weak (strong) tunneling coupling (nonlinear couplings) $`0<\alpha /\lambda <4`$. However, the tunneling coupling and the inter-atomic nonlinear interactions strongly affect the amplitude of the tunneling current in the regime of strong (weak) tunneling coupling (nonlinear couplings) $`\alpha /\lambda >4`$. From Figure 1 we can see that the amplitude of the tunneling current increases with increasing both the initial number of atoms and the tunneling coupling in the regime of strong (weak) tunneling coupling (nonlinear couplings) $`\alpha /\lambda >4`$. In Figure 2, we display the scaled frequency of the tunneling current, $`\omega /\lambda =\omega _2(N,0)/\lambda `$, as a function of interaction strengths for different initial number of atoms when the initial state is $`\|\mathrm{\Psi }(0)=\|N,0`$, $`\lambda _a=\lambda _b=\lambda `$, and $`\omega _a=2\omega _b`$. It is interesting to note that from Figure 2 we can see that there exists a zero-frequency point, labeled by D. From Eqs. (14) it is straight forward to see that the zero-frequency point is a degenerate point of energy of the system under our consideration, at which the nondegenerate perturbation theory is broken. From Eqs. (7), (8), (14), and (20) we can find that the degenerate point is given by the expression $`\alpha /\lambda =[(3N4)(7N20)/2(N2)(N3)]^{1/2}`$. Figure 2 indicates that on the left hand side of the degenerate point D the scaled frequency decreases with increasing the tunneling coupling and/or decreasing the nonlinear interaction strengths , and increases with increasing the initial number of atoms. On the other hand, on the left hand side of the degenerate point D the scaled frequency increases with increasing both the tunneling coupling and the initial number of atoms, and/or decreasing the nonlinear interaction strengths . However, it is customary to consider a Bose-Einstein condensate to be in a coherent state, associated with a macroscopic wave function with both amplitude and a phase, the presence of which is originated from Bose broken symmetry. Assume that the two condensate are initially in the coherent states $`\|\alpha `$ and $`\|\beta `$, which are eigenstates of $`\widehat{a}`$ and $`\widehat{b}`$, repectively, then we have $`C_{n,m}(0)`$ $`=`$ $`{\displaystyle \frac{A_{n,m}}{\mathrm{exp}(\|\alpha \|^2+\|\beta \|^2)}}[{\displaystyle \frac{\alpha ^n\beta ^m}{\sqrt{n!m!}}}`$ (48) $`+{\displaystyle \frac{a_{n,m}\alpha ^{n2}\beta ^{m+1}}{\sqrt{(n2)!(m+1)!}}}`$ $`+{\displaystyle \frac{b_{n,m}\alpha ^{n+2}\beta ^{m1}}{\sqrt{(n+2)!(m1)!}}}].`$ Making use of Eqs. (11), (12) and (21) we find the expression of the population difference $`P(t)`$ $`=`$ $`2{\displaystyle \underset{n,m}{}}\{{\displaystyle \frac{1}{2}}p_0(n,m)C(n,m;n,m)`$ (55) $`+p_1(n,m)C(n,m;n2,m+1)`$ $`\times \mathrm{cos}(E_{n,m}E_{n2,m+1})t`$ $`+p_2(n,m)C(n,m;n+2,m1)`$ $`\times \mathrm{cos}(E_{n,m}E_{n+2,m1})t`$ $`+p_3(n,m)C(n,m;n+4,m2)`$ $`\times \mathrm{cos}(E_{n,m}E_{n+4,m2})t\},`$ where we have introduced the notations $`p_0(n,m)`$ $`=`$ $`(n2m)+(n2m4)a_{n,m}^2`$ (57) $`+(n2m+4)b_{n,m}^2,`$ $`p_1(n,m)`$ $`=`$ $`(n2m4)a_{n,m},`$ (58) $`p_2(n,m)`$ $`=`$ $`(n2m+4)b_{n,m},`$ (59) $`p_3(n,m)`$ $`=`$ $`(n2m+4)a_{n+4,m2}b_{n,m},`$ (60) and $$C(n,m;n^{},m^{})=C_{n,m}(0)C_{n^{},m^{}}^{}(0).$$ (62) From Eq. (22) we can get the nonzero time average of the population difference labeled by $`\overline{P}`$ given by $$\overline{P}=\underset{n,m}{}C(n,m;n,m)p_0(n,m),$$ (63) which implies that there exists the MQST between the atomic condensate and the molecular condensate. From Eq. (22) we can obtain the Josephson-like tunneling current $`I(t)`$ $`=`$ $`{\displaystyle \underset{n,m}{}}[I_1(n,m)\mathrm{sin}(E_{n,m}E_{n2,m+1})t`$ (66) $`I_2(n,m)\mathrm{sin}(E_{n,m}E_{n+2,m1})t`$ $`I_3(n,m)\mathrm{sin}(E_{n,m}E_{n+4,m2})t],`$ with $`I_1(n,m)`$ $`=`$ $`2p_1(n,m)C(n,m;n2,m+1)`$ (68) $`\times (E_{n,m}E_{n2,m+1})/N,`$ $`I_2(n,m)`$ $`=`$ $`2p_2(n,m)C(n,m;n+2,m1)`$ (70) $`\times (E_{n,m}E_{n+2,m1})/N,`$ $`I_3(n,m)`$ $`=`$ $`2p_3(n,m)C(n,m;n+4,m2)`$ (72) $`\times (E_{n,m}E_{n+4,m2})/N.`$ where $`N=\|\alpha \|^2+2\|\beta \|^2`$ is the total number of the atoms in the two condensates. We now discuss the effect of the decoherence. In experiments on trapped Bose condensates of atomic gases, condensate atoms continuously interact with non-condensate atoms (environment). As is well known, interactions between a quantum system and environment cause two types of unwelcomed effects: dissipation and decoherence . The dissipation effect, which dissipates the energy of the quantum system into the environment, is characterized by the relaxation time scale $`\tau _r`$. In contrast, the decoherence effect is much more insidious because the coherence information leaks out into the environment in another time scale $`\tau _d`$, which is much shorter than $`\tau _r`$. Since macroscopic quantum phenomena in Bose-Einstein condensates mainly depend on $`\tau _d`$ rather than $`\tau _r`$, the discussions in present paper only focus on the decoherence problem rather than the dissipation effect. We use a reservoir consisting of an infinite set of harmonic oscillators to model environment of condensate atoms and molecules in a trap, and assume the total Hamiltonian to be $`\widehat{H}_T`$ $`=`$ $`\widehat{H}+{\displaystyle \underset{k}{}}\omega _k\widehat{b}_k^{}\widehat{b}_k+F(\{\widehat{S}\}){\displaystyle \underset{k}{}}c_k(\widehat{b}_k^{}+\widehat{b}_k)`$ (74) $`+F(\{\widehat{S}\})^2{\displaystyle \underset{k}{}}{\displaystyle \frac{c_k^2}{\omega _k^2}},`$ where the second term is the Hamiltonian of the reservoir. The last term is a renormalization term. The third term represents the interaction between the system and the reservoir with a coupling constant $`c_k`$, where $`\{\widehat{S}\}`$ is a set of linear operators of the system or their linear combinations in the same picture as that of $`\widehat{H}`$, $`F(\{\widehat{S}\})`$ is an operator function of $`\{\widehat{S}\}`$. In order to enable what the interaction between the system and environment describes is decoherence not dissipation, we require that the linear operator $`\widehat{S}`$ commutes with the the Hamiltonian of the system $`\widehat{H}`$. Then, the interaction term commutes with the Hamiltonian of the system. This implies that there is no energy transfer between the system and its environment. So that it does describe the decoherence. The concrete form of the function $`F(\{\widehat{S}\})`$, which may be considered as an experimentally determined quantity, may be different for different environment. The Hamiltonian $`\widehat{H}_T`$ can be exactly solved by using the unitary transformation $`\widehat{U}=\mathrm{exp}[\widehat{H}_k(c_k/\omega _k)(\widehat{b}_k^{}\widehat{b}_k)]`$. Corresponding to the Hamiltonian (29), the total density operator of the system plus reservoir can be expressed as $`\widehat{\rho }_T(t)=e^{i\widehat{H}t}\widehat{U}^1e^{it_k\omega _k\widehat{b}_k^{}\widehat{b}_k}\widehat{\rho }_T(0)\widehat{U}^1e^{it_k\omega _k\widehat{b}_k^{}\widehat{b}_k}\widehat{U}e^{i\widehat{H}t}`$. We assume that the system and reservoir are initially in thermal equilibrium and uncorrelated, so that $`\widehat{\rho }_T(0)=\widehat{\rho }(0)\widehat{\rho }_R`$, where $`\widehat{\rho }(0)`$ is the initial density operator of the system, and $`\widehat{\rho }_R`$ the density operator of the reservoir, which can be written as $`\widehat{\rho }_R=_k\widehat{\rho }_k(0)`$ with $`\widehat{\rho }_k(0)`$ is the density operator of the $`k`$-th harmonic oscillator in thermal equilibrium. After taking the trace over the reservoir, we can get the reduced density operator of the system, denoted by $`\widehat{\rho }(t)=tr_R\widehat{\rho }_T(t)`$, whose matrix elements in the eigenstate representation of $`\widehat{H}`$ are explicitly written as $`\rho _{(m^{},n^{})(m,n)}(t)`$ $`=`$ $`\|\rho _{(m^{},n^{})(m,n)}(0)\|e^{\gamma _{(m^{},n^{})(m,n)}(t)}`$ (76) $`\times e^{i\varphi _{(m^{},n^{})(m,n)(t)}},`$ where the damping factor and the phase shift are defined by $`\gamma _{(m^{},n^{})(m,n)}(t)`$ $`=`$ $`v_{}^2(m^{},n^{};m,n)Q_2(t),`$ (77) $`\varphi _{(m^{},n^{})(m,n)}(t)`$ $`=`$ $`v_+(m^{},n^{};m,n)v_{}(m^{},n^{};m,n)Q_1(t)`$ (79) $`+\theta _{(m^{},n^{})(m,n)},`$ where we have introduced the following notations: $`v_\pm (n,m;n^{},m^{})`$ $`=`$ $`F(\{S(n,m)\})\pm F(\{S(n^{},m^{})\}),`$ (80) $`\rho _{(m^{},n^{})(m,n)}(0)`$ $`=`$ $`\|\rho _{(m^{},n^{})(m,n)}(0)\|e^{i\theta _{(m^{},n^{})(m,n)}},`$ (81) and the two reservoir-dependent functions are given by $`Q_1(t)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑\omega J(\omega ){\displaystyle \frac{c^2(\omega )}{\omega ^2}}\mathrm{sin}(\omega t),`$ (82) $`Q_2(t)`$ $`=`$ $`2{\displaystyle _0^{\mathrm{}}}𝑑\omega J(\omega ){\displaystyle \frac{c^2(\omega )}{\omega ^2}}\mathrm{sin}^2({\displaystyle \frac{\omega t}{2}})\mathrm{coth}({\displaystyle \frac{\beta \omega }{2}}).`$ (83) Here we have taken the continuum limit of the reservoir modes: $`_k_0^{\mathrm{}}𝑑\omega J(\omega )`$, where $`J(\omega )`$ is the spectral density of the reservoir, $`c(\omega )`$ is the continuum expression for $`c_k`$, and $`\beta =1/k_BT`$ with $`k_B`$ and $`T`$ being the Boltzmann constant and temperature, respectively. Eq. (30) indicates that the interaction between the system and its environment induces a phase shift and a decaying factor in the reduced density operator of the system. We now consider the population difference between the atomic condensate and the molecular condensate in the presence of the decoherence defined by $`P(t)=Tr\widehat{\rho }(t)(\widehat{n}_a2\widehat{n}_b)`$. Making use of Eq. (30), We find that $`P(t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n,,m}{}}\{p_0(n,m)\|\rho _{(n,m)(n,m)}(0)\|\mathrm{cos}\theta _{(n,m)(n,m)}`$ (90) $`+p_1(n,m)\|\rho _{(n,m)(n2,m+1)}(0)\|`$ $`\times \mathrm{cos}\varphi _{(n,m)(n2,m+1)}(t)e^{\gamma _{(n,m)(n2,m+1)}(t)}`$ $`+p_2(n,m)\|\rho _{(n,m)(n+2,m1)}(0)\|`$ $`\times \mathrm{cos}\varphi _{(n,m)(n+2,m1)}(t)e^{\gamma _{(n,m)(n+2,m1)}(t)}`$ $`+p_3(n,m)\|\rho _{(n,m)(n+4,m2)}(0)\|`$ $`\times \mathrm{cos}\varphi _{(n,m)(n+4,m2)}(t)e^{\gamma _{(n,m)(n+4,m2)}(t)}\}.`$ Then, the tunneling current is given by $`I(t)`$ $`=`$ $`{\displaystyle \underset{n,m}{}}\{p_1(n,m)\|\rho _{(n,m)(n2,m+1)}(0)\|`$ (99) $`\times [\dot{\gamma }_{(n,m)(n2,m+1)}(t)+\dot{\varphi }_{(n,m)(n2,m+1)}(t)`$ $`\times \mathrm{sin}\varphi _{(n,m)(n2,m+1)}(t)]e^{\gamma _{(n,m)(n2,m+1)}(t)}`$ $`+p_2(n,m)\|\rho _{(n,m)(n+2,m1)}(0)\|`$ $`\times [\dot{\gamma }_{(n,m)(n+2,m1)}(t)+\dot{\varphi }_{(n,m)(n+2,m1)}(t)`$ $`\times \mathrm{sin}\varphi _{(n,m)(n+2,m1)}(t)]e^{\gamma _{(n,m)(n+2,m1)}(t)}`$ $`+p_3(n,m)\|\rho _{(n,m)(n+4,m2)}(0)\|`$ $`\times [\dot{\gamma }_{(n,m)(n+4,m2)}(t)+\dot{\varphi }_{(n,m)(n+4,m2)}(t)`$ $`\times \mathrm{sin}\varphi _{(n,m)(n+4,m2)}(t)]e^{\gamma _{(n,m)(n+4,m2)}(t)}\}.`$ From Eqs. (37) and (38) we can immediately draw one important qualitative conclusion: since $`\gamma _{(n,m)(n^{},m^{})}`$ is positive definite, the existence of the decoherence is always to tend to suppress the population difference and tunneling current between the atomic condensate and the molecular condensate. From Eqs. (31), (32), and (35)-(38) we see that all necessary information about the effects of the environment on the population difference and the tunneling current is contained in the spectral density of the reservoir. To procced further let us now specialize to the Ohmic case with the spectral distribution $`J(\omega )=[\eta \omega /c^2(\omega )]\mathrm{exp}(\omega /\omega _c)`$, where $`\omega _c`$ is the high frequency cut-off, $`\eta `$ is a positive characteristic parameter of the reservoir. With this choice, at low temperature the functions $`Q_1(t)`$ and $`Q_2(t)`$ are given by the expressions $`Q_1(t)=\eta \mathrm{tan}^1(\omega _ct)`$ and $`Q_2(t)=\eta \{\frac{1}{2}\mathrm{ln}[1+(\omega _ct)^2]+\mathrm{ln}[\frac{\beta }{\pi t}\mathrm{sinh}(\frac{\pi t}{\beta })]\}`$. In particular, At zero temperature and in the meaningful domain of time $`\omega _ct1`$, we have $`\dot{Q}_1(t)\eta /(\omega _ct^2)`$, and $`Q_2(t)\eta \mathrm{ln}(\omega _ct)`$. Then we find $`P(t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{n,,m}{}}\{p_0(n,m)\|\rho _{(n,m)(n,m)}(0)\|\mathrm{cos}\theta _{(n,m)(n,m)}`$ (106) $`+p_1(n,m)\|\rho _{(n,m)(n2,m+1)}(0)\|`$ $`\times \mathrm{cos}\varphi _{(n,m)(n2,m+1)}(t)(\omega _ct)^{(\eta _{1nm}^{})^2}`$ $`+p_2(n,m)\|\rho _{(n,m)(n+2,m1)}(0)\|`$ $`\times \mathrm{cos}\varphi _{(n,m)(n+2,m1)}(t)(\omega _ct)^{(\eta _{2nm}^{})^2}`$ $`+p_3(n,m)\|\rho _{(n,m)(n+4,m2)}(0)\|`$ $`\times \mathrm{cos}\varphi _{(n,m)(n+4,m2)}(t)(\omega _ct)^{(\eta _{3nm}^{})^2}\},`$ $`I(t)`$ $`=`$ $`{\displaystyle \underset{n,m}{}}\{(p_1(n,m)\|\rho _{(n,m)(n2,m+1)}(0)\|[{\displaystyle \frac{(\eta _{1nm}^{})^2}{t}}`$ (112) $`+{\displaystyle \frac{\eta _{1nm}^+}{\omega _ct^2}}\mathrm{sin}\varphi _{(n,m)(n2,m+1)}(t)](\omega _ct)^{(\eta _{1nm}^{})^2}`$ $`+p_2(n,m)\|\rho _{(n,m)(n+2,m1)}(0)\|[(\eta _{2nm}^{})^2`$ $`+{\displaystyle \frac{\eta _{2nm}^+}{\omega _c}}\mathrm{sin}\varphi _{(n,m)(n+2,m1)}(t)](\omega _ct)^{(\eta _{2nm}^{})^2}`$ $`+p_3(n,m)\|\rho _{(n,m)(n+4,m2)}(0)\|[(\eta _{3nm}^{})^2`$ $`+{\displaystyle \frac{\eta _{3nm}^+}{\omega _c}}\mathrm{sin}\varphi _{(n,m)(n+4,m2)}(t)](\omega _ct)^{(\eta _{3nm}^{})^2}\},`$ where we have used $`\eta _{1nm}^\pm =\sqrt{\eta }v_\pm (n,m;n2,m+1)`$, $`\eta _{2nm}^\pm =\sqrt{\eta }v_\pm (n,m;n+2,m1)`$, $`\eta _{3nm}^\pm =\sqrt{\eta }v_\pm (n,m;n+4,m2)`$, and $`\eta _{inm}^+=\eta _{inm}^+\eta _{inm}^{}`$. Eqs. (39) and (40) indicates that the tunneling current decays away according to the “power law”, where we have noted that the decaying factors can not be taken outside the summation on the r.h.s. of Eqs. (39) and (40). In summary, we have studied tunneling dynamics of atomic pairs in atomic Bose-Einstein condensates with Feshbach resonances, and shown that the tunneling of the atomic pairs depends on not only the tunneling coupling between the atomic condensate and the molecular condensate, but also the inter-atomic nonlinear interactions and the initial number of atoms in these condensates. especially, we have shown that the tunneling coupling and the inter-atomic nonlinear interactions strongly affect the tunneling of atomic pairs in the regime of strong (weak) tunneling coupling (nonlinear couplings) when the atomic condensate is in a number state and the molecular condensate in the vacuum state. This implies that the tunneling of atomic pairs between the atomic condensate and the molecular condensate can be manipulated and controlled by varying the tunneling coupling and/or inter-atomic nonlinear interaction strengths . We have revealed the existence of the MQST between the atomic condensate and the molecular condensate. The MQST is a kind of nonlinear effects which vanishes in the absence of the inter-atomic nonlinear interactions. We have also discussed the influence of decoherence induced by non-condensate atoms on the tunneling dynamics, and shown that decoherence suppresses the atomic-pair tunneling. Finally, it should be mentioned that inelastic collisions between the atomic and molecular condensates and reservoir may affect the atomic-pair tunneling between the atomic and molecular condensates. Influence of inelastic collisions between the atomic and molecular condensates can be taken account into through introducing an imaginary part in the interaction strengths $`\lambda `$ and $`\lambda _b`$ . Inelastic collisions between the system of the atomic-molecular condensates and reservoir lead to dissipation . A detailed investigation on the dissipation problem of the system of the atomic-molecular condensates is beyond the scope of the present paper, and will be given elsewhere. ACKNOWLEDGMENTS This work was supported in part by grants from Hong Kong Research Grants Council (RGC) and the Hong Kong Baptist University Faculty Research Grant (FRG). L.M.K. also acknowledges support from the Climbing Project of China and NSF of China, the Excellent Young-Teacher Foundation of the Educational Commission of China, ECF and STF of Hunan Province. The authors would like to thank Dr. J.H. Xiao for his useful discussions and help in preparing figures of the paper. Figure Captions FIG. 1. Amplitude of the tunneling current as a function of interaction strengths for different initial number of atoms when the initial state is $`\|\mathrm{\Psi }(0)=\|N,0`$, $`\lambda _a=\lambda _b=\lambda `$, and $`\omega _a=2\omega _b`$. Here we have set $`\omega =\omega _2(N,0)`$. FIG. 2. Scaled frequency of the tunneling current as a function of interaction strengths for different initial number of atoms when the initial state is $`\|\mathrm{\Psi }(0)=\|N,0`$, $`\lambda _a=\lambda _b=\lambda `$, and $`\omega _a=2\omega _b`$. Here we have set $`\omega =\omega _2(N,0)`$.
warning/0005/gr-qc0005124.html	ar5iv	text	# Maximal Acceleration Effects in Reissner-Nordström Space. <sup>a</sup> Dipartimento di Scienze Fisiche “E.R. Caianiello”, Università di Salerno, Italy. <sup>b</sup> Istituto Nazionale di Fisica Nucleare, Sezione di Napoli. <sup>c</sup> Facoltà d’Ingegneria, Università del Sannio, Benevento, Italy. <sup>d</sup> Department of Physics, University of Regina, Regina, Sask. S4S 0A2, Canada. <sup>e</sup> Canadian Institute for Theoretical Astrophysics, University of Toronto, Toronto, Ontario K7L 3N6, Canada <sup>f</sup> IIASS, Vietri sul Mare (SA), Italy The dynamics of a relativistic particle in a Reissner-Nordström background is studied using Caianiello model with maximal acceleration. The behaviour of the particle, embedded in a new effective geometry, changes with respect to the classical scenario because of the formation of repulsive potential barriers near the horizon. Black hole formation by accretion of massive particles is not therefore a viable process in the model. At the same time, the naked singularity remains largely unaffected by maximal acceleration corrections. PACS: 04.50.+h, 04.70.Bw Keywords: Reissner-Nordström metric, Quantum Geometry, Maximal Acceleration A model developed by Caianiello and collaborators provides quantum mechanics with a geometrical framework and delves into a number of fundamental issues such as the unification of general relativity and quantum mechanics and the regularization of field equations. In particular, the model interprets quantization as curvature of the eight-dimensional space-time tangent bundle TM, incorporates the Born reciprocity principle and the notion that the proper acceleration of massive particles has an upper limit $`𝒜_m`$. Classical and quantum arguments supporting the existence of a maximal acceleration (MA) have long been discussed in the literature . MA also appears in the context of Weyl space and of a geometrical analogue of Vigier’s stochastic theory . Its existence would rid black hole entropy of ultraviolet divergencies ,, and circumvent inconsistencies associated with the application of the point-like concept to relativistic quantum particles . Some authors regard $`𝒜_m`$ as a universal constant fixed by Planck’s mass ,, but a direct application of Heisenberg’s uncertainty relations , as well as the geometrical interpretation of the quantum commutation relations given by Caianiello, suggest that $`𝒜_m`$ be fixed by the rest mass of the particle itself according to $`𝒜_m=2mc^3/\mathrm{}`$. A limit on the acceleration also occurs in string theory. Here the upper limit manifests itself through Jeans-like instabilities which occur when the acceleration induced by the background gravitational field is larger than a critical value $`a_c=(m\alpha )^1`$for which the string extremities become causally disconnected . $`m`$ is the string mass and $`\alpha `$ is the string tension. Frolov and Sanchez have then found that a universal critical acceleration $`a_c=(m\alpha )^1`$ must be a general property of strings. While in all these instances the critical acceleration results from the interplay of the Rindler horizon with the finite extension of the particle ,, in the Caianiello model MA is a basic physical property of all massive particles which appears from the out-set in the physical laws. At the same time the model introduces an invariant interval in TM that leads to a regularization of the field equations that does not require a fundamental length as in and does therefore preserve the continuum structure of space-time. Applications of Caianiello’s model include cosmology , where the initial singularity can be avoided while preserving inflation, the dynamics of accelerated strings , the energy spectrum of a uniformly accelerated particle , the periodic structure as a function of momentum in neutrino oscillations and the expansion of the very early universe . The model also makes the metric observer–dependent, as conjectured by Gibbons and Hawking . The extreme large value that $`𝒜_m`$ takes for all known particles makes a direct test of the model very difficult. Nonetheless a direct test that uses photons in a cavity has also been suggested . More recently, we have worked out the consequences of the model for the classical electrodynamics of a particle , the mass of the Higgs boson and the Lamb shift in hydrogenic atoms . In the last instance the agreement between experimental data and MA corrections is very good for $`H`$ and $`D`$. For $`He^+`$ the agreement between theory and experiment is improved by $`50\%`$ when MA corrections are included. MA effects in muonic atoms also appear to be measurable . In all these works space-time is endowed with a causal structure obtained by means of an embedding procedure pioneered in and discussed at length in , . The procedure stipulates that the line element experienced by an accelerating particle, in the presence of gravity, is given by $$d\tau ^2=\left(1+\frac{g_{\mu \nu }\ddot{x}^\mu \ddot{x}^\nu }{𝒜_m^2}\right)g_{\alpha \beta }dx^\alpha dx^\beta \sigma ^2(x)g_{\alpha \beta }dx^\alpha dx^\beta ,$$ (1) where $`\ddot{x}^\mu =d^2x^\mu /ds^2`$ is the, in general, non–covariant acceleration of a particle along its worldline. As a consequence, the effective space-time geometry experienced by accelerated particles exhibits mass-dependent corrections, which in general induce curvature, give rise to a mass-dependent violation of the equivalence principle and vanish in the classical limit $`\left(𝒜_m\right)^1=\frac{\mathrm{}}{2mc^3}0`$. We have recently studied the modifications produced by MA in the motion of a test particle in a Schwarzschild field and found that these account for the presence of a spherical shell, external to the Schwarzschild sphere, that is forbidden to any classical particle and hampers the formation of a black hole. The analogous occurrence of a classically impenetrable shell was derived by Gasperini as a consequence of the breaking of the local $`SO(3,1)`$ symmetry . The shell remains impervious to quantum, scalar particles . Before proceeding, a few comments are in order . The effective theory presented is intrinsically non-covariant, as is the four-acceleration that appears in $`\sigma ^2(x)`$. In addition $`\sigma ^2(x)`$ could be eliminated from (1) by means of a coordinate transformation if one applied the principles of general relativity to this effective theory. On the contrary, the embedding procedure requires that $`\sigma ^2(x)`$ be present in (1) and that it be calculated in the same coordinates of the unperturbed gravitational background. Nonetheless the choice of $`\ddot{x}^\mu `$ is supported by the derivation of $`𝒜_m`$ from quantum mechanics, by special relativity and by the weak field approximation to general relativity. The model is not intended, therefore, to supersede general relativity, but rather to provide a way to calculate the quantum corrections to the structure of space-time implied by (1). Remarkably, the results of persist in isotropic coordinates . For convenience, the natural units $`\mathrm{}=c=G=1`$ are used below. The purpose of this paper is to extend the calculation of the MA corrections to the Reissner–Nordström metric. The problem is not trivial. This metric, of form $$ds^2=\left(1\frac{2M}{r}+\frac{e^2}{r^2}\right)dt^2\frac{dr^2}{1\frac{2M}{r}+\frac{e^2}{r^2}}r^2\left(d\vartheta ^2+\mathrm{sin}^2\vartheta d\phi ^2\right),$$ (2) does, in fact, contain two lengths, the mass of the central object $`M`$ and its charge $`e`$ . The global space structure therefore depends on the ratio of these two lengths. In addition, the solution does not pertain to vacuum and contains a naked singularity for certain values of the parameters. Apart from the trivial case of vanishing $`e`$, which yields the Schwarzschild metric, three cases must be distinguished. If $`\left\|e\right\|<M`$, the Schwarzschild horizon moves from $`2M`$ to $`r_A=M+\sqrt{M^2e^2}`$, while a new horizon forms at the origin. Its position is $`r_B=M\sqrt{M^2e^2}`$. In the region between these two horizons, the component $`g_{00}`$ of the metric tensor is negative and particles can only move towards the origin. Elsewhere, particles are allowed to move away from the origin. When $`\left\|e\right\|=M`$, the two horizons merge at $`r=M`$. Yet particles at this point cannot leave and the region inside the sphere of radius $`M`$ is still not accessible to an external observer. Finally, for charges larger than the critical value, no horizon exists, the metric (2) can be used in the whole space and particles can move in and out without restrictions. The solution then corresponds to a naked singularity. In order to calculate the corrections to the Reissner–Nordström field experienced by a particle initially at infinity and falling toward the origin along a geodesic, one must calculate the metric induced by the embedding procedure discussed in to first order in the parameter $`𝒜_m^2`$. On choosing $`\theta =\pi /2`$, one finds $$\sigma ^2\left(r\right)=1+\frac{1}{𝒜_m^2}\left[\left(1\frac{2M}{r}+\frac{e^2}{r^2}\right)\ddot{t}^2\left(1\frac{2M}{r}+\frac{e^2}{r^2}\right)^1\ddot{r}^2r^2\ddot{\phi }^2\right].$$ (3) The components of the four–acceleration can be easily obtained by following the same steps as in the Schwarzschild case . When these are substituted into (3), the conformal factor becomes $$\sigma ^2\left(r\right)=1+\frac{1}{𝒜_m^2}\{\frac{\left(\frac{M}{r^2}+\frac{e^2}{r^3}+\frac{\stackrel{~}{L}^2}{r^3}\frac{3ML^2}{r^4}+\frac{2e^2\stackrel{~}{L}^2}{r^5}\right)}{\left(1\frac{2M}{r}+\frac{e^2}{r^2}\right)}+$$ $$+[\stackrel{~}{E}^2\frac{\left(\frac{2M}{r^2}\frac{2e^2}{r^3}\right)^2}{\left(1\frac{2M}{r}+\frac{e^2}{r^2}\right)^3}\frac{4\stackrel{~}{L}^2}{r^4}][\stackrel{~}{E}^2(1\frac{2M}{r}+\frac{e^2}{r^2})(1+\frac{\stackrel{~}{L}^2}{r^2})]\}.$$ (4) As in , $`\stackrel{~}{E}`$ and $`\stackrel{~}{L}`$ are the energy and the angular momentum per unit of test particle mass. Bearing in mind that the magnitude of the energy–momentum four–vector is given by the rest mass $`m`$ of the particle according to $`g_{\alpha \beta }p^\alpha p^\beta =m^2`$ and using the modified metric (1), one can introduce an effective potential $`V_{eff}`$ by means of the equation $$\left(\frac{dr}{d\tau }\right)^2=\stackrel{~}{E}^2V_{eff}^2,$$ (5) where $$V_{eff}^2=\stackrel{~}{E}^2\frac{\stackrel{~}{E}^2}{\sigma ^4\left(r\right)}+\frac{1}{\sigma ^2\left(r\right)}\left(1\frac{2M}{r}+\frac{e^2}{r^2}\right)\left(1+\frac{\stackrel{~}{L}^2}{r^2\sigma ^2\left(r\right)}\right).$$ (6) In the limit $`𝒜_m\mathrm{}`$, or equivalently $`\sigma ^2\left(r\right)1`$, we recover the classical Reissner–Nordström potential. Moreover, as $`e0`$, (6) reduces to the effective potential of . In the following we introduce the adimensional variable $`\rho =r/M`$ and the parameters $`ϵ=\left(M𝒜_m\right)^1`$, $`\lambda =\stackrel{~}{L}/M`$ and $`\eta =e/M`$. Notice that $`V_{eff}^2(\rho )1`$ as $`\rho \mathrm{}`$ and $`V_{eff}^2(\rho )\stackrel{~}{E}^2`$ as $`\rho 0`$. The main features of (6) can be derived and compared with the corresponding corrections determined for the Schwarzschild case. a) $`\left\|\eta \right\|<1`$ and $`\lambda =0`$. Since the component $`g_{00}=\left(1\frac{2}{\rho }+\frac{\eta ^2}{\rho ^2}\right)`$ of the metric tensor vanishes at $`\rho _A=1+\sqrt{1\eta ^2}`$ and $`\rho _B=1\sqrt{1\eta ^2}`$, it follows from (4) that the conformal factor diverges as $`g_{00}^3`$ for the same values of $`\rho `$. $`\sigma ^2\left(\rho \right)`$ always appears in the denominators of the effective potential which thus equals $`\stackrel{~}{E}^2`$. Fig. 1 shows $`V_{eff}^2\left(\rho \right)`$ for a radially incoming particle ($`\lambda =0`$, $`\eta ^2=0.1`$) and different values of the energy. As in the Schwarzschild case , a shell is formed in $`\rho _A`$ that prevents external particles from falling into the no–return region. The shell is dynamical in origin, is impenetrable to classical particles and remains so at higher orders of approximation in $`𝒜_m^2`$. An expansion of (6) in the neighborhood of $`\rho _A=1+\sqrt{1\eta ^2}`$ shows that the potential barrier behaves like $$V_{eff}^2(\rho )=\stackrel{~}{E}^2+\frac{4f(\eta )}{\stackrel{~}{E}^4ϵ^2}\left(\rho \rho _A\right)^4+O\left((\rho \rho _A)^5\right),$$ (7) where $$f(\eta )=$$ $$\frac{\left(1\eta ^2\right)\left(1+\sqrt{1\eta ^2}\right)^5}{8\left(3272\eta ^2+54\eta ^415\eta ^6\right)+9\eta ^8+\sqrt{1\eta ^2}\left(4\eta ^2\right)\left(6496\eta ^2+36\eta ^4\eta ^6\right)}.$$ $`V_{eff}^2(\rho )`$ clearly has the minimum $`\stackrel{~}{E}^2`$ on the horizon $`\rho _A`$. On the left of the barrier, there is a divergence originated by a zero in $`\sigma ^2\left(\rho \right)`$. This divergence is always negative because the dominant term in (6), $`\frac{\stackrel{~}{E}^2}{\sigma ^4\left(\rho \right)}`$, is negative for the corresponding values of $`\rho `$. To the left of the divergence, the effective potential returns to values very close to those of the Reissner-Nordström potential. The shape of $`V_{eff}^2`$ is rather insensitive to variations of the parameters. In particular, the barrier and the following divergence never change. Near $`\rho _B`$ the situation is more complex since the conformal factor has two, or four more roots (depending on the values of $`\eta `$ and $`\stackrel{~}{E}`$), one on the left of $`\rho _B`$ and one, or three, on the right. The detailed behaviour of the corresponding divergences is, however, irrelevant in the present context because no classical particle can reach the region to the left of the potential barrier near $`\rho _A`$. b) $`\left\|\eta \right\|<1`$ and $`\lambda 0`$. As in a) above, the effective potential takes the value $`\stackrel{~}{E}^2`$ when $`\rho `$ tends to zero and to the horizons $`\rho _A`$ and $`\rho _B`$. On the right of $`\rho _A`$ the potential barrier is finite when the angular momentum $`\lambda `$ is lower than a critical value $`\lambda _c(ϵ,\stackrel{~}{E},\eta )`$, but infinite if $`\lambda >\lambda _c`$. This barrier prevents external particles from falling into the no–return region. Fig. 2 shows the behaviour of $`V_{eff}^2\left(\rho \right)`$ in the neighborhood of the horizon $`\rho _A`$ for two different values of $`\lambda `$, $`\lambda _1>\lambda _c`$ and $`\lambda _2<\lambda _c`$. The singularity structure on the left of $`\rho _A`$ is as complex as in case a), and likewise irrelevant. c) $`\left\|\eta \right\|=1`$. At this critical value, $`g_{00}`$ has a double root at $`\rho =1`$. Here, the effective potential forms the usual barrier as $`\sigma ^2\left(\rho \right)`$ diverges. This is shown in Fig. 3 (for $`\lambda =0`$) and in Fig. 4 (for $`\lambda 0`$). An expansion of (6) in the neighborhood of $`\rho _A=1`$ yields the potential barrier $$V_{eff}^2(\rho )=\stackrel{~}{E}^2+\frac{4}{\stackrel{~}{E}^4ϵ^2}\left(\rho \rho _A\right)^6+O\left((\rho \rho _A)^7\right)$$ which, clearly, has the minimum $`\stackrel{~}{E}^2`$ on the horizon $`\rho _A=1`$. Near the origin there still is the singularity produced by the zero of the conformal factor. Here too the sign is determined by the values of $`\stackrel{~}{E}`$ and $`\lambda `$; it becomes positive when $`\lambda ^2>\frac{\stackrel{~}{E}^2\rho _0^2}{\left(1\rho _0\right)^2}`$, where $`\rho _0`$ is the point at which $`\sigma ^2\left(\rho \right)`$ vanishes. d) $`\left\|\eta \right\|>1`$. $`g_{00}`$ is positive for all values of $`\rho `$. There are no more horizons and the barrier shrinks (Fig. 5) until it disappears. Although the naked singularity is no longer point-like, its behaviour remains unaffected by MA corrections. For a black hole with 3.3 solar masses, the critical value for the central charge is $`5.7\times 10^{20}`$ coulomb. Obviously, this value can never be reached in the normal case irrespective of the relative signs of the charges of $`m`$ and $`M`$. If, in fact, matter and black hole had charges of the same sign, then the electromagnetic repulsion would prevent the fall of matter. If the signs were opposite, then particles would be strongly attracted and would tend to neutralize the whole system. For this reason real black holes can only be “low charged” and should be described by the $`\left\|e\right\|<M`$ case. The corrections to the Reissner–Nordström metric induced by the presence of MA do indeed represent an interesting confirmation of the physics already discussed for the Schwarzschild case. Apart from the classical shift from $`\rho =2`$ to $`\rho =\rho _A`$, the possible presence of charge in the source essentially leaves the barrier at the external horizon unchanged. This means that all the remarks about the formation of black holes made in with regard to the Schwarzschild metric can be extended to the “low charged” Reissner–Nordström space. Classically, the presence of the barrier forbids the formation of a black hole by the accretion of matter, unless the latter is transformed first into massless particles and these are absorbed by the star at a rate higher than the corresponding re-emission rate. A recent calculation shows that even the quantum properties of matter do not alter this behaviour and that gravitational collapse would at least be slowed down. Nonetheless the shell would still retain the appearances of a very intensely radiating, charged, compact object ($`\rho _A2`$). Besides confirming what found in about black hole dynamics, the application of Caianiello model to the Reissner–Nordström metric offers other interesting aspects. Since the shell is impervious to charged matter of both signs, the charge of the source can not change. It is not possible, at present, to speculate on the behaviour of electron-positron pairs produced in the immediate vicinity of a source with sufficiently high charge. A barrier on a horizon is always accompanied by a singularity on the side where $`g_{00}`$ is negative. These divergences are even present in the curvature invariants, so they must be considered as physical singularities of the effective metric. They are however inaccessible to massive probes and only regard the behaviour of matter inside a black hole, if the latter somehow formed. It is however doubtful that the Reissner–Nordström solution would adequately describe this type of matter. For high values of the angular momentum the divergence in the effective potential becomes positive causing the formation of an infinite, repulsive potential barrier. It is important to remember that the Reissner–Nordström space–time allows causality violations in cases c) (and d)) and in the region $`\rho <\rho _B`$ of a) and b). The impenetrability of the shell renders inaccessible those regions with physically inadmissible properties that violate causality. Finally, MA alters the radius of the naked singularity, but not its nature and does not therefore embody any form of cosmic censorship. Acknowledgments Research supported by Ministero dell’Università e della Ricerca Scientifica of Italy, funds ex 40% and 60% DPR 382/80 and the Natural Sciences and Engineering Research Council of Canada.
warning/0005/hep-ph0005049.html	ar5iv	text	# 1 Introduction ## 1 Introduction Investigation of high-energy interactions of virtual photons with nucleons and nuclei gives information on the dynamics of high-density partonic systems as well as on the transition between perturbative and nonperturbative regimes in QCD. Experiments at HERA have found two extremely important properties of small-$`x`$ physics : a fast increase of parton densities as $`x`$ decreases and the existence of substantial diffractive production in deep inelastic scattering (DIS). From a theoretical point of view there are good reasons to believe that the fast increase of the $`\sigma _{\gamma ^{}p}`$ with energy in the HERA energy range will change to a milder (logarithmic) increase at much higher energies. This is due to unitarity effects, which are related to shadowing in highly dense systems of partons - with eventual “saturation” of densities. This problem has a long history (for reviews see ) and has been extensively discussed in recent years . It is closely connected to the problem of the dynamics of very high-energy heavy ion collisions . In this paper we will address the problem of unitarization in $`\gamma ^{}p`$-interactions using the reggeon calculus with a supercritical Pomeron ($`\mathrm{\Delta }_P\alpha _P(0)1>0`$). In this approach the unitarization effects mentioned above are described by multi-Pomeron exchanges. The problem of the Pomeron in QCD is not solved yet. The calculation based on QCD perturbation theory leads to the so-called BFKL-Pomeron and the problem of unitarization for such Pomeron has been considered by many authors - in the leading logarithmic approximation. However, recent calculations indicate that NLO corrections are large and modify substantially the picture of high-energy interactions. In particular the intercept of the leading Regge pole depends on nonperturbative effects . Here, we will use a more phenomenological approach and will assume that the Pomeron is a simple pole with an intercept $`\alpha _P(0)=1.2`$ determined from the analysis of high-energy hadronic interaction - when all multi-Pomeron contributions are taken into account . In ref. it was suggested that the increase of the effective intercept of the Pomeron, $`\alpha _{eff}=1+\mathrm{\Delta }_{eff}`$, as $`Q^2`$ increases from zero to several GeV<sup>2</sup> is mostly due to a decrease of shadowing effects with increasing $`Q^2`$. A parametrization of the $`Q^2`$ dependence of $`\mathrm{\Delta }_{eff}`$ such that $`\mathrm{\Delta }_{eff}0.1`$ for $`Q^20`$ (as in soft hadronic interactions) and $`\mathrm{\Delta }_{eff}0.2`$ (our input or bare Pomeron intercept) for $`Q^2`$ of the order of a few GeV<sup>2</sup>, gives a good description of all existing data on $`\gamma ^{}p`$ total cross-sections in the region of $`Q^2<5÷10`$ GeV<sup>2</sup> . At larger $`Q^2`$, effects due to QCD evolution become important. Using the above parametrization as initial condition in the QCD evolution equation, allows to describe the data in the whole region of $`Q^2`$ studied at HERA . In the reggeon calculus the amount of rescatterings is closely related to diffractive production. AGK-cutting rules allow to calculate the cross-section of inelastic diffraction if contributions of multi-Pomeron exchanges to the elastic scattering amplitude are known. Thus, it is very important for self-consistency of theoretical models to describe not only total cross-sections, but, simultaneously, inelastic diffraction. Indeed, in the reggeon calculus the variation of $`\mathrm{\Delta }_{eff}`$ with $`Q^2`$ is directly related to the corresponding variation of the ratio of diffractive to total cross-sections. In this paper we present an explicit model which leads to the pattern of energy behaviour of $`\sigma _{\gamma ^{}p}^{(tot)}(W^2,Q^2)`$ for different $`Q^2`$ described above. Moreover, it allows to describe simultaneously diffraction production by real and virtual photons. The plan of the paper is as follows. In Section 2 the model is described. Section 3 contains the expressions of the total cross-section $`\sigma _{\gamma ^{}p}^{tot}`$. Formulae for diffractive production are given in Section 4. Section 5 is devoted to the description of experimental data on $`\sigma _{\gamma ^{}p}^{(tot)}(W^2,Q^2)`$ and diffraction dissociation of the photon. In Section 6 we compare our results with results of other authors on this subject, discuss the problem of “saturation” and give predictions for future experiments. The parameters of the model and their numerical values are given in Appendix 1. ## 2 Description of the model The amplitude of elastic $`\gamma ^{}p`$ scattering at very high energy can be described in terms of the single Pomeron exchange (Fig. 1a) and multi-Pomeron exchanges (Figs. 1 b, 1c). The Pomeron exchange contribution to $`\sigma _{\gamma ^{}p}^{tot}`$ (or to the structure function $`F_2=\frac{Q^2}{4\pi ^2\alpha _{e.m}}\sigma _{\gamma ^{}p}^{(tot)}`$), increases as $`\left(\frac{1}{x}\right)^\mathrm{\Delta }`$ at small $`x=\frac{Q^2}{W^2+Q^2}`$. The contribution of $`n`$-Pomeron exchange behaves as $`\left(\frac{1}{x}\right)^{n\mathrm{\Delta }}`$. These contributions are important at very small-$`x`$ to restore unitarity. A resummation of multi-Pomeron contribution is provided by Gribov reggeon calculus . In this approach a two-Pomeron exchange diagram (Fig. 1b), for example, can be represented as a sum over all intermediate states (Fig. 2) which can be diffractively produced by a single Pomeron exchange. Thus the magnitude of the two Pomeron exchange in the elastic amplitude is related to the cross-section of diffractive production. The multi-Pomeron exchanges (Fig. 1c) are related by AGK-cutting rules to shadowing corrections to diffractive production. Let us first discuss the main mechanisms of diffraction dissociation of a proton and of a virtual photon. For a proton the main contribution to the diagram of Fig. 2 is due to the proton intermediate state. The dissociation of the proton will be taken into account in the quasi-eikonal approximation \- which works well in hadronic interactions. As for a virtual photon, it can diffractively produce vector mesons (Fig. 3a) (the so-called quasi-elastic processes) and large mass states. For $`M^2Q^2`$, the latter correspond to the triple-Regge diagrams of Fig. 3b with $`P`$ and $`f`$-exchanges in the $`t`$-channel. Models of this type allow to describe diffractive dissociation both at $`Q^2=0`$ and at large $`Q^2`$ . A simple quasi-eikonal ansatz for multi-Pomeron exchanges based on this type of model was used in ref. for the description of multiparticle production in DIS. However, the quasi-eikonal approximation for diffractive dissociation of a highly virtual photon is difficult to interpret in terms of reggeon diagrams - while an explicit calculation of these diagrams for multi-Pomeron exchanges leads to too many parameters for the matrix of transition between different channels. For these reasons, we adopt here a different approach, based on the quark-parton picture of high-energy interaction of a real or virtual photon. Taking into account that a photon dissociates into a $`q\overline{q}`$-pair and this pair can have multiple interactions with the target, we separate such pairs into two groups : “aligned” or asymmetric configurations, which will be denoted by $`A`$ and the rest or symmetric configurations denoted by $`S`$. The $`A`$-configuration was introduced by Bjorken and Kogut in the framework of the parton model in DIS and is characterized by a strongly asymmetric distribution of the relative longitudinal momenta of $`q`$ and $`\overline{q}`$ ($`z_q1`$, $`z_{\overline{q}}1`$ or viceversa). This separation into $`A`$ and $`S`$ components is very important at large $`Q^2`$, where the $`A`$-configuration has a large transverse size and a large interaction cross-section with the proton, while the $`S`$-component has small transverse size ($`1/Q`$) and a small cross-section. On the other hand the probability of the $`A`$-configuration decreases as $`1/Q^2`$ and both configurations lead to the same behaviour, $`\sigma _{\gamma ^{}p}^{(tot)}1/Q^2`$, at large $`Q^2`$ . What is important, however, is that these two components give rise to a different behavior for diffraction : the $`A`$-component gives the same $`1/Q^2`$ behaviour for total and diffractive cross-sections, while the $`S`$-component leads to a diffractive cross-section which decreases as $`(1/Q^2)^2`$ at large $`Q^2`$. These results are well known and have been extensively discussed in the literature in connection with HERA experiments. The $`q\overline{q}`$-state is described, in general, by non-perturbative dynamics and only for large $`Q^2`$ and small size component it is possible to use perturbative QCD. Finally, we stress that at present energies, due to phase space limitations, only diagrams with none or one triple reggeon interactions are important. In particular large mass diffraction corresponds to a triple Pomeron interaction diagram. However, to be more general we resum all fan-type diagrams - of the type shown in Fig. 4. This will be discussed in more detail in the next Section. ## 3 The total $`\gamma ^{}p`$ cross-section We formulate the expressions of cross-sections in the impact parameter space and take into account that $`A`$\- and $`S`$-components are diagonal : $$\sigma _{\gamma ^{}p}^{(tot)}(s,Q^2)=4d^2b\sigma _{\gamma ^{}p}^{(tot)}(b,s,q^2),$$ (1) $$\sigma _{\gamma ^{}p}^{(tot)}(b,s,Q^2)=\underset{i}{}g_i(Q^2)\sigma _i^{(tot)}(b,s,Q^2)i=A,S.$$ (2) Here $`g_i(Q^2)`$ is the probability of the $`i`$-th component, which according to the discussion above, has been chosen in the form $$g_A(Q^2)=\frac{g_A(0)}{1+\frac{Q^2}{m_A^2}};g_S(Q^2)=g_S(0)$$ (3) where $`g_A(0)`$, and $`m_A^2`$ are considered as phenomenological free parameters. The cross-sections $`\sigma _i^{(tot)}(b,s,q^2)`$ are written in the form $$\sigma _i^{(tot)}(b,s,Q^2)=\frac{1\mathrm{exp}(C\chi _i(b,s,Q^2)}{2C},$$ (4) $$\chi _i(s,b,Q^2)=\frac{\chi _{i0}^P(s,b,Q^2)}{1+a\chi _3(s,b,Q^2)}+\chi _{i0}^f(s,b,Q^2).$$ (5) The expressions for the eikonals $`\chi _{i0}`$ are : $$\chi _{i0}^k(s,b,Q^2)=C_i^k\frac{f_i(Q^2)}{\lambda _{0k}^i(Q^2,\xi )}\mathrm{exp}\left(\mathrm{\Delta }_k\xi \frac{b^2}{4\lambda _{0k}^i}\right)$$ (6) with $`k=P,f`$ and $`i=A,S`$ where $`\mathrm{\Delta }_k=\alpha _k(0)1;\xi =\mathrm{}n{\displaystyle \frac{s+Q^2}{s_0+Q^2}}`$ $`\lambda _{0k}^i=R_{0ki}^2(Q^2)+\alpha _k^{}\xi .`$ (7) The quantity $`\xi `$ is chosen in such a way as to behave as $`\mathrm{}n\frac{1}{x}`$ for large $`Q^2`$ and as $`\mathrm{}n\frac{s}{s_0}`$ for $`Q^2=0`$. The functions $`f_i(Q^2)`$ have been chosen in the formIn eqs. (3) and (8) we used the simplest parametrization of $`g_i(Q^2)`$ and $`f_i(Q^2)`$ consistent with our requirements. In general $`g_S(Q^2)`$ and $`f_A(Q^2)`$ can depend on $`Q^2`$. $$f_i(Q^2)=\{\begin{array}{cc}\frac{1}{1+Q^2/m_S^2}& ,i=S\hfill \\ & \\ 1& ,i=A.\hfill \end{array}$$ (8) The parametrization of the radius $`R_{0ki}^2`$ is given in Appendix 1. The slopes of the trajectories $`\alpha _P^{}=0.25`$ GeV<sup>-2</sup> and $`\alpha _f^{}=0.9`$ GeV<sup>-2</sup> were fixed at their values known from soft hadronic interactions. We turn next to the denominator in the first term of eq. (5). Here $`a=g_{pp}^P(0)r_{PPP}/16\pi `$ where $`g_{pp}^P(0)`$ and $`r_{PPP}(0)`$ are, respectively, the Pomeron-proton coupling and the triple Pomeron vertex, both at $`t=0`$. The expression of $`\chi _3(s,b,Q^2)`$ is given in Section 4. Note that eqs. (4)-(6) correspond to a quasi-eikonal model . For $`a=0`$, its Born term, eq. (5), is a sum of Pomeron and secondary $`(f)`$ reggeon exchanges. The latter is included in order to be able to use the model at energies as low as $`\sqrt{s}10`$ GeV. The coefficient $`C`$ takes into account the dissociation of the proton and is taken to be equal to 1.5 . To first order in $`a`$, the Born term, eq. (5), also contains the contribution of $`PPP`$ and $`PfP`$ triple regge terms (see Fig. 5). This can be easily seen by developing the first term of the r. h. s. of eq. (5) to first order in $`a`$, and using the expression of $`\chi _3`$ in eqs. (26) and (27). Taken to all orders in $`a`$, eq. (5) corresponds to the resummation of all fan-type diagrams of the type shown in Fig. 4, using the Schwimmer formula . Here, however, each new branching contains not only a Pomeron as in Fig. 4 but also a secondary $`f`$-reggeon exchange. These branchings are controlled by the parameters $`a\gamma _P`$ and $`a\gamma _f`$, respectively (see eq. (26))Eq. (5) is valid when these couplings are the same irrespectively of whether the branching takes place off a Pomeron or an $`f`$.. As discussed in the last paragraph of Section 2, diagrams with more than one branching (i. e. more than one triple reggeon interaction) are not important at present energies due to phase space limitations. Note that the triple Pomeron interaction introduces a large (gluonic) size component in our $`S`$-component. Actually, the separation between the two components is washed out at large $`Q^2`$ due to QCD evolution. Note also that the second term of eq. (5) ($`f`$-exchange) does not contain a denominator - as does the first one. Such a denominator would add to the $`f`$-exchange $`fPf`$ and $`fff`$ terms. However, such double reggeon $`fP`$ and $`ff`$ exchanges are already present in eq. (4) (terms $`\chi _{i0}^P`$ $`\chi _{i0}^f`$ and $`(\chi _{i0}^f)^2`$, respectively). In the above formulae we have neglected the real parts of the reggeon exchanges given by the signature factors. Since $`\alpha _P(0)`$ is not very different from unity, their effect is expected to be rather small - except when the contribution of the multi-Pomeron exchanges becomes very important. As for secondary ($`f`$) exchange, the signature factor enters only in terms with multiple $`f`$-exchange - which are not important at the comparatively high energies of the present experiments. Many parameters of the model such as $`R_{0ki}`$ are strongly constrained by the data on hadronic interactions and were fixed (see Appendix 1). Let us also note that the $`f`$ exchange in $`\sigma _{\gamma ^{}p}^{(tot)}`$ is related to the valence quarks contribution, which for small $`x`$ belongs to the $`A`$-component. Thus $`C_S^f=0`$. The list of all parameters and their numerical values is given in Appendix 1. ## 4 Diffractive production Let us now introduce the expressions for diffraction dissociation of a virtual photon. They can be obtained, using AGK-cutting rules, from the imaginary part of the $`\gamma ^{}p`$-elastic scattering, given by eqs. (1)-(8). Neglecting s-channel iterations of the triple-Pomeron graphs, the total diffraction dissociation cross-section can be written as a sum of three components $$\sigma _{\gamma ^{}p}^{(dif)}=\underset{i=A,S}{}\sigma _i^{(0)}+\sigma _{PPP}$$ (9) where $$\sigma _i^{(0)}=4g_i^2(Q^2)\left(\sigma _i^{(tot)}(b,s,Q^2)\right)^2d^2b$$ (10) and $$\sigma _{PPP}=2\underset{i}{}g_i^2(Q^2)\chi _{PPP}^i(b,s,Q^2)e^{2C\chi _i(s,b,Q^2)}d^2b$$ (11) where $`\chi _{PPP}^i(b,s,Q^2)=a\chi _i^P(s,b,Q^2)\chi _3(s,b,Q^2)`$, and $`\chi _i^P(s,b,Q^2)`$ is given by the first term in the r. h. s. of eq. (5). For the total diffractive production cross-section, which includes the diffraction dissociation of the proton, it is necessary to multiply expressions (10), (11) by the factor $`C=1.5`$, introduced in Section 3. Note that to first order in $`a`$, $`\sigma _{PPP}`$ corresponds to the sum of the two triple reggeon contributions in Fig. 5. Taken to all orders in $`a`$, $`\sigma _{PPP}`$ corresponds to the resummation of all fan-type diagrams \- obtained by adding to the diagrams in Fig. 5, multiple branchings with Pomeron and $`f`$-reggeon at each branching (see the discussion in Section 3). In the following, the sum of these two types of contributions will be referred to as the triple Pomeron ($`PPP`$) - although the diagram in Fig. 5b corresponds to an interference term. With gaussian forms of the Born term in impact parameter, eq. (6), the expression (10) can be also written as follows (see second paper of ) $$\sigma _i^{(0)}=\frac{g_i^2(Q^2)\sigma _{B_i}}{C}\left[f\left(\frac{Z_i}{2}\right)f(Z_i)\right]$$ (12) where $$f(Z)=\underset{n=1}{\overset{\mathrm{}}{}}\frac{(Z)^{n1}}{nn!}$$ (13) $$\sigma _{B_i}=2\chi _i(s,b,Q^2)d^2b$$ (14) $$Z_i=\frac{8C}{\sigma _{B_i}}\chi _i^2(s,b,Q^2)d^2b.$$ (15) In the next sections we will compare our model with differential diffractive cross-sections as functions of $`M^2`$ (square of the mass of the diffractively produced system) or of $`\beta =\frac{Q^2}{M^2+Q^2}`$ (which is convenient at large $`Q^2`$ and plays the same role for the Pomeron as the variable $`x`$ for the proton<sup>§</sup><sup>§</sup>§It should be noted that the “structure function of the Pomeron” $`F_P(\beta ,Q^2)`$ can be used only for the cases when multi-Pomeron exchanges are small, because in general it is impossible to separate out this amplitude.) or of the variable $`x_P`$ ($`x_p=\frac{x}{\beta }`$). Present experiments at HERA measure the differential cross-sections integrated over the variable $`t`$ (square of the momentum transfer between initial and final protons) and usually the function $`F_{2D}^{(3)}`$ is introduced : $$x_PF_{2D}^{(3)}=\frac{Q^2}{4\pi ^2\alpha _{e.m.}}x_P\frac{d\sigma }{dx_Pdt}𝑑t.$$ (16) In our model it can be written as a sum of three terms $$F_{2D}^{(3)}=\left(\underset{i}{}F_{2Di}^{(3)}(x,Q^2,\beta )+F_{2DPPP}^{(3)}(x,Q^2,\beta )\right)$$ (17) where $$F_{2Di}^{(3)}=F_{2Di}^{(3)B}\zeta _i.$$ (18) $`F_{2Di}^{(3)B}(x,Q^2,\beta )`$ are the lowest order (Born) approximations for these functions (their explicit forms are given below) while $`\zeta _i`$ are “suppression factors” due to higher order multi-Pomeron exchanges $$\zeta _i=\left[f\left(\frac{Z_i}{2}\right)f(Z_i)\right]\frac{8}{Z_i}.$$ (19) The functions $`F_{2Di}^{(3)}`$ satisfy the condition $$_{x_{P_{min}}}^{x_{P_{max}}}F_{2Di}^{(3)}𝑑x_P=\sigma _i^{(0)}(x,Q^2)\frac{Q^2}{4\pi ^2\alpha _{e.m.}}$$ (20) where $`x_{P_{min}}=\frac{x}{\beta _{max}}`$ ; $`x_{P_{max}}=0.1`$ ; $`\beta _{max}=\frac{Q^2}{M_{min}^2+Q^2}`$ : with $`M_{min}^2=4m_\pi ^2`$. The $`\beta `$-dependence of the $`A`$-component was chosen according to ref. , where the small-$`\beta `$-dependence was determined from the $`PPf`$-triple-Regge asymptotics with a simple ansatz for the behaviour at $`\beta 1`$ $$x_PF_{2DA}^{(3)B(P)}(x,Q^2,\beta )=\frac{Q^2g_A^2(Q^2)}{4\pi ^2\alpha _{e.m.}}d^2b\left(\chi _A^P(x,Q^2,b)\right)^2\frac{\stackrel{~}{\beta }^{2\mathrm{\Delta }_P\mathrm{\Delta }_f}(1\beta )^{n_P(Q^2)}}{_{\beta _{min}}^{\beta _{max}}\frac{d\beta }{\beta }\stackrel{~}{\beta }^{2\mathrm{\Delta }_p\mathrm{\Delta }_f}(1\beta )^{n_P(Q^2)}}$$ (21) where $`\stackrel{~}{\beta }=\frac{Q^2+S_0}{Q^2+M^2}=\beta \frac{\stackrel{~}{x}}{x}`$ ; $`\beta _{min}=\frac{x}{x_P^{max}}=10x`$. In the following all powers of $`1\beta `$ are chosen according to the arguments presented in refs. and : $$n_P(Q^2)=\frac{1}{2}+\frac{3}{2}\left(\frac{Q^2}{c+Q^2}\right)$$ (22) with $`c=3.5`$ GeV<sup>2</sup>. The $`S`$-component should be maximal at $`\beta `$ close to 1 and was chosen in the form $$x_PF_{2DS}^{(3)B}(x,Q^2,\beta )=\frac{Q^2g_S^2(Q^2)}{4\pi ^2\alpha _{e.m.}}d^2b\chi _S^2(x,Q^2,b)\frac{\stackrel{~}{\beta }}{_{\beta _{min}}^{\beta _{max}}\frac{d\beta }{\beta }\stackrel{~}{\beta }}.$$ (23) This component contains the contribution of comparatively small masses and should decrease as $`\beta 1`$ faster than $`PPf`$ contribution, eq. (21). Our choice satisfies this condition. Note that from perturbative QCD a behaviour $`\beta ^3`$ has been found . However, the contribution of this component at large $`Q^2`$, where perturbative QCD is applicable, is small $`(m_S^2/Q^2`$). Our results are rather insensitive to the power of $`\beta `$ in the range 1 to 3 . The triple-Pomeron ($`PPP`$ plus $`PfP`$) contribution $`F_{2DPPP}^{(3)}(x,Q^2,\beta )`$, is given by $$x_PF_{2DPPP}^{(3)}(x,Q^2,\beta )=x_PF_{2DPPP}^{(3)B}(x,Q^2,\beta )\frac{\sigma _{PPP}}{\sigma _{PPP}^B},$$ (24) where $`\sigma _{PPP}`$ is given by eq. (11), its Born term, $`\sigma _{PPP}^B`$, by the same equation with $`C=0`$, and $$x_PF_{2DPPP}^{(3)B}(x,Q^2,\beta )=\frac{Q^2}{4\pi ^2\alpha _{e.m.}}2ad^2b\underset{i}{}g_i^2(Q^2)\chi _i^P(b,s,Q^2)\chi _3(s,b,Q^2,\beta ),$$ (25) with $$\chi _3(s,b,Q^2,\beta )=\underset{k}{}\gamma _k\mathrm{exp}\left(\frac{b^2}{4\lambda _k\left(\frac{\stackrel{~}{\beta }}{\stackrel{~}{x}}\right)}\right)\left(\frac{\stackrel{~}{\beta }}{\stackrel{~}{x}}\right)^{\mathrm{\Delta }_k}\frac{(1\beta )^{n_P(Q^2)+4}}{\lambda _k\left(\frac{\stackrel{~}{\beta }}{\stackrel{~}{x}}\right)}.$$ (26) Here $`\gamma _P=1`$, $`\gamma _f`$ determines the relative strength of the $`PfP`$-contribution and $`\lambda _k=R_{1k}^2+\alpha _K^{}\mathrm{}n\left(\frac{\stackrel{~}{\beta }}{\stackrel{~}{x}}\right)`$. The function $`\chi _3(s,b,Q^2)`$, which enters in eqs. (5) and (11) is given by $$\chi _3(s,b,Q^2)=_{\beta _{min}}^{\beta ^{max}}\frac{d\beta }{\beta }\chi _3(s,b,Q^2,\beta ).$$ (27) Since the triple Pomeron formula is not valid for low masses, we use here $`M_{min}=1`$ GeV. Note that $`x_PF_{2DPPP}^{(3)B}`$ is obtained from $`\sigma _{PPP}^B`$ replacing $`\chi _3(s,b,Q^2)`$ by $`\chi _3(s,b,Q^2,\beta )`$. In this way $`𝑑x_PF_{2DPPP}^{(3)}𝑑x_P=\frac{Q^2}{4\pi ^2\alpha _{e.m.}}\sigma _{PPP}`$. ## 5 Comparison with experiment The model outlined above was used for a joint fit of the data on the structure function $`F_2(x,Q^2)`$ and diffractive structure function $`F_{2D}^{(3)}(x,Q^2,\beta )`$ in the region of small x ($`x<10^2`$) and $`Q^210`$ GeV<sup>2</sup>Actually, for $`F_2`$ only data with $`Q^23.5`$ GeV<sup>2</sup> were included in the fit.. The full list of parameters and their values either fixed or obtained in the fit are given in the Appendix 1. In Fig. 6, the dependence of $`\sigma _{\gamma ^{}p}`$ as a function of $`Q^2`$ for different energies is shown and Fig. 7 gives the $`x`$-dependence of $`F_2(x,Q^2)`$ for different $`Q^2`$. Our results are compared with experimental data. The description of the data is excellent. Note that the $`x`$-dependence and its variation with $`Q^2`$ is fixed in the model and is strongly correlated with the ratio of $`\sigma _{\gamma ^{}p}^{(dif)}/\sigma _{\gamma ^{}p}^{(tot)}`$ and its dependence on $`Q^2`$. The model confirms the $`Q^2`$-dependence of the effective intercept $`\mathrm{\Delta }_{eff}=\frac{d\mathrm{}nF_2(x,Q^2)}{d\mathrm{}n\left(\frac{1}{x}\right)}`$, assumed in ref. . The function $`\mathrm{\Delta }_{eff}`$ depends not only on $`Q^2`$, but also on $`x`$ and it is shown in Fig. 8 for two different intervals of $`x`$. The function $`\mathrm{\Delta }_{eff}`$ decreases as $`Q^2`$ or $`x`$ decreases due to the increase of shadowing effects. It is interesting to note that, for $`Q^2=0`$, $`\mathrm{\Delta }_{eff}`$ in the region of HERA energies is close to 0.13, i.e. substantially higher than in hadronic interactions. For $`s\mathrm{}`$ the cross-section $`\sigma _{\gamma ^{}p}^{(tot)}`$ has a Froissard-type behavior in $`\mathrm{}n^2\frac{s}{s_0}`$. In Figs. 9 and 10 the results of the model are compared with experimental data on diffractive structure function $`F_2^{D(3)}(x,\beta ,Q^2)`$. The $`\beta `$ dependence for fixed values of $`Q^2`$ and $`x_P`$ is shown in Fig. 9 and the $`x_P`$-dependence for different values of $`Q^2`$ and $`\beta `$ is presented in Fig. 10. The model reproduces the experimental data rather well. For comparison we limited ourselves to the data at low values of $`x_P10^2`$, where the effect of the non-diffractive RRP contribution (which is not included in the model) is small. Diffractive production at $`Q^2=0`$ is reasonably well reproduced as well (Fig. 11). The dependence of diffractive cross-sections on $`\mathrm{\Delta }_{eff}^{dif}\left(F_{2D}^{(3)}\left(\frac{1}{x_P}\right)^{2\mathrm{\Delta }_{eff}^D}\right)`$ is shown in Fig. 12 as a function of $`Q^2`$ (for fixed $`M`$). The experimental points are also shown. The model predicts a weak dependence of $`\mathrm{\Delta }_{eff}^{dif}`$ on $`Q^2`$ for $`Q^2>5÷10`$ GeV<sup>2</sup> and QCD-evolution does not change this result . ## 6 Discussion and conclusions The model we propose here for the description of total and diffractive cross-sections of interaction of a virtual photon with a proton is a natural generalization of models used for the description of high-energy hadronic interactions. The main parameter of the model - intercept of the Pomeron, $`\mathrm{\Delta }_P\alpha _P(0)1`$, was fixed from a phenomenological study of these interactions ($`\mathrm{\Delta }_P=0.2`$) and was found to give a good description of $`\gamma ^{}p`$-interactions in a broad range of $`Q^2`$ ($`0Q^2<10GeV^2`$). Note that a single Pomeron is present in the model. It should be noticed that at higher values of $`Q^2`$, QCD evolution is important and our results should be used as initial condition in DGLAP evolution (see refs ). In particular, due to QCD evolution, the effective Pomeron intercept in $`F_2`$ will reach values significantly larger than 1.2 at large $`Q^2`$. For diffractive process the situation is different. In this case, the effective Pomeron intercept at large $`Q^2`$ is not significantly modified as a result of QCD evolution. Moreover, the latter has rather small effect for intermediate values of $`\beta `$ . For these reasons our model, without QCD evolution, can be used here at comparatively larger values of $`Q^2`$. Another important parameter of the theory is the triple Pomeron vertex, which appears in the parameter $`a`$. It is quite remarkable that the same value of this parameter allows to describe diffractive data both at $`Q^2=0`$ and at a few GeV<sup>2</sup>. Thus, in our approach, the value of this parameter can be determined from soft diffraction data . The variation of the photon virtuality gives a unique possibility to study unitarization effects for different scales. Let us discuss the qualitative features of our model and its relation to the models introduced by other authors . It has been mentioned above that the different role of small and large size components of a virtual photon has been emphasized by many authors and perturbative calculations of unitarization effects for $`\sigma _{\gamma ^{}p}`$ have been carried out in ref. . In our opinion the perturbative calculations can be valid only for the small size component ($`S`$) at large $`Q^2`$. An interesting analysis of unitarization effects for both total $`\gamma ^{}p`$-cross-section and diffractive production was performed by Golec-Biernat and Wüsthoff . These authors, however, postulated an unconventional form for the elastic $`\gamma ^{}p`$-amplitude and it is not clear which diagrams lead to these unitarization effectsCorrelations in hadronic interactions have been studied for many years. It is known that they can be described in an eikonal model - where the Born term, at fixed impact parameter, is used as eikonal. In the model of Golec-Biernat and Wüsthoff , on the contrary, the Born term is integrated over impact parameter and then eikonalized. The unitarity effects obtained in this way are very large. It is unlikely that they can describe the observed features of correlations.. We agree with their conclusion that for the structure function $`F_2(x,Q^2)`$ the unitarization effects at HERA energies are important mostly in the region of rather small $`Q^2<2÷4`$ GeV<sup>2</sup>. Indeed, for $`Q^2>2÷4`$ GeV<sup>2</sup>, the experimental value of $`\mathrm{\Delta }_{eff}`$ is close to the one of the input (bare) Pomeron $`\mathrm{\Delta }_P=0.2`$. Within our approach this is only possible if the $`S`$-component (where unitarity corrections are higher twists) dominates. However, some unitarity corrections are present at large $`Q^2`$ due to the $`A`$-component and the triple Pomeron interaction. Moreover, they will become more important at higher energies (see the discussion in the next paragraph). However, we differ in many details from their predictions on the properties of diffractive production at large $`Q^2`$ and on the pattern of “saturation” of parton densities at large $`Q^2`$ and extremely small $`x`$. Let us now discuss the saturation patterns in our model, i.e. the properties of $`\sigma _{\gamma ^{}p}^{(tot)}(s,Q^2)`$ in the limit $`s\mathrm{}`$. It follows from formulae (1)-(8) that for the $`A`$-component, which has a large cross-section, $`\sigma _{\gamma ^{}p}^{(tot)}(s,b,Q^2)`$ will tend to a saturation limit $`1/2C`$ rather fast as energy increases. On the other hand, for the small size (at large $`Q^2`$) component ($`S`$) the situation is different. If we neglect the triple-Pomeron contribution (i.e. with $`a=0`$ in eq. (5)), $`\chi _S(s,b,Q^2)\frac{Cm_S^2}{Q^2}\mathrm{exp}\left(\mathrm{\Delta }_P\xi \frac{b^2}{4\lambda _{0P}^S}\right)`$ will increase with energy until the increase of $`\mathrm{exp}(\mathrm{\Delta }_P\xi )`$ does not overcompensate the smallness of $`\frac{m_S^2}{Q^2}`$. For such extremely large energies, cross-sections in the small impact parameter region will saturate to a $`Q^2`$-independent value and $`F_2(x,Q^2)Q^2`$. This is a usual picture of saturation in perturbative QCD -. However, the inclusion of the $`PPP`$-fan diagrams (large distance, nonperturbative effects) according to eqs. (5), (26), (27) changes the picture drastically. In this case the increase of $`\chi _3(s,b,Q^2)\mathrm{exp}(\mathrm{\Delta }\xi )`$ will compensate the increase of $`\chi _S(s,b,Q^2)`$ at very large $`\xi `$ and lead, at large $`Q^2`$, to a behaviour $`\sigma _{\gamma ^{}p}^{(tot)}\frac{1}{Q^2}f(\mathrm{}nQ^2)`$ even for $`x0`$. (Note that for large enough $`Q^2`$ this will happen in a region where $`\chi _S`$ is still small). The above effects are relatively small at present energies due to the smallness of the triple Pomeron vertex $`r_{PPP}`$ which determines the constant $`a`$. However, in the “saturation limit” the properties of $`\sigma _{\gamma ^{}p}^{(tot)}(s,Q^2)`$ in our model are quite different from those in other models . Our predictions for $`\sigma _{\gamma ^{}p}^{(tot)}(s,Q^2)`$ for energies higher than those accessible at HERA are shown in Figs. 13. In a forthcoming publication we introduce a different but closely related model with also two components for the interaction of the $`q\overline{q}`$ pair : a large-size component, parametrized in the same way as the $`A`$-component in the present work, and a small-size component computed using the $`q\overline{q}`$ perturbative QCD wave function - with its longitudinal and transverse components. The results are very similar to the ones obtained here. Again, a single Pomeron with intercept 1.2 allows to describe the data on both total $`\gamma ^{}p`$ cross-section and diffractive production. ## Appendix 1 In this Appendix we give the list of all the parameters in the formulae for $`\sigma _{\gamma ^{}p}^{(tot)}(s,Q^2)`$ (or $`F_2(x,Q^2)`$) and diffractive production, given in sections 3 and 4 respectively - together with their numerical values. The total cross-section of $`\gamma ^{}p`$ interactions is related to the structure function $`F_2`$ as follows $$\sigma _{\gamma ^{}p}^{(tot)}(s,Q^2)=\frac{4\pi ^2\alpha _{em}}{Q^2}F_2(x,Q^2),$$ $`(\mathrm{A}.1)`$ where $`s=W^2`$ and $`x=\frac{Q^2}{1+Q^2}`$. This cross-section is related by eq. (1) to the corresponding quantity, $`\sigma _{\gamma ^{}p}^{(tot)}(b,s,Q^2)`$, in the impact parameter space. The latter is written as the sum of two components (eq. (2)) with probabilities given eq. (3) where $`g_A(0)`$, $`g_S(0)`$ and $`m_A^2`$ are free parameters. Their values obtained from the fit are $$g_A(0)=0.15\alpha _{e.m.},g_S(0)=0.367\alpha _{e.m.},m_A^2=0.227\mathrm{GeV}^2.$$ $`(\mathrm{A}.2)`$ For $`a=0`$ (i.e. when triple reggeon vertices vanish ; see eq. (5)) the cross-sections $`\sigma _i^{(tot)}(b,s,Q^2)`$ are written in a quasi-eikonal form (eq. (4)) with Born terms given by the exchange of the Pomeron $`P`$ and secondary reggeon ($`f`$) trajectories). Their expressions contain the intercepts $`\alpha _k(0)=1+\mathrm{\Delta }_k`$ $`(k=P,f)`$ for which we take $$\mathrm{\Delta }_P=0.2,\mathrm{\Delta }_f=0.3.$$ $`(\mathrm{A}.3)`$ They also contain five other parameters treated as free ones : $`C_A^P`$, $`C_S^P`$, $`C_A^f`$, ($`C_S^f=0`$), $`s_0`$ and $`m_S^2`$. Their values obtained in the fit are: $$m_S^2=1.343\mathrm{GeV}^2,s_0=0.463\mathrm{GeV}^2,$$ $$C_A^P=0.830\mathrm{GeV}^2,C_S^P=0.807\mathrm{GeV}^2,C_A^f=14.28\mathrm{GeV}^2.$$ $`(\mathrm{A}.4)`$ The value of the Pomeron intercept $`\alpha _P(0)=1.2`$ has been found in soft hadronic interactions . It is also consistent with the intercept of the BFKL Pomeron at NLO. When treated as a free parameter in the fit its value turns out to be very close to the one in eq. (A.3). The value of the intercept of the $`f`$ has been allowed to change within some restricted limits. The value $`\alpha _f(0)=0.7`$ leads to better results than the more conventional value $`\alpha _f(0)=0.5`$. Our expressions contain also the triple Pomeron term $`PPP`$ whose strength is controlled by the parameter $`a`$ in eq. (5), and the interference term $`PfP`$ whose strength (relative to the $`PPP`$ term) is given by the parameter $`\gamma _f`$. We use $$a=0.052\mathrm{GeV}^2\gamma _f=8.$$ $`(\mathrm{A}.5)`$ Note that the value of $`a`$ obtained in the fit agrees with the value obtained from soft diffraction (see the discussion in Section 6). Finally we turn to the quantities $$\lambda _{0k}^i=R_{0ki}^2(Q^2)+\alpha _k^{}\xi $$ $`(\mathrm{A}.6)`$ and $$\lambda _k=R_{1k}^2+\alpha _k^{}\mathrm{}n(\beta /x)$$ $`(\mathrm{A}.7)`$ in eqs. (5) and (26). They are directly related to the $`t`$-dependence of the elastic $`\gamma ^{}p`$ amplitudes and of the triple reggeon vertices, respectively, and have the standard regge behavior as functions of $`\xi `$ \- or $`\beta /x`$. We use the conventional parameters of the reggeon slopes $$\alpha _P^{}=0.25\mathrm{GeV}^2,\alpha _f^{}=0.9\mathrm{GeV}^2.$$ $`(\mathrm{A}.8)`$ Our results for the diffractive cross-sections integrated over $`t`$ are not sensitive to the values of the parameters $`\lambda `$ and we have fixed them to the values obtained from soft interaction data and from vector dominance. We take $$\begin{array}{c}R_{0kA}^2(Q^2)=R_{vkA}^2+R_{pk}^2\hfill \\ \\ R_{0kS}^2(Q^2)=\frac{R_{vkS}^2}{1+\frac{Q^2}{m_\rho ^2}}+R_{pk}^2\hfill \end{array}$$ $`(\mathrm{A}.9)`$ with $$R_{pP}^2=R_{Pf}^2=2\mathrm{GeV}^2,R_{vki}^2=1\mathrm{GeV}^2$$ $`(\mathrm{A}.10)`$ and $$R_{1P}^2=R_{1f}^2=2.2\mathrm{GeV}^2.$$ $`(\mathrm{A}.11)`$ Note that the $`Q^2`$-dependence of the $`\gamma ^{}\gamma ^{}P`$ and $`\gamma ^{}\gamma ^{}f`$ vertices for the $`S`$ component have been chosen in such a way that the corresponding radii tend to zero for $`Q^2\mathrm{}`$ . Acknowledgments It is a pleasure to thank K. Boreskov, O. Kancheli, G. Korchemski, U. Maor and C. Merino for discussions. This work is partially supported by NATO grant OUTR.LG 971390. E. G. F. and C. A. S. thank Ministerio de Educación y Cultura of Spain for financial support. ## Figure captions Figure 1. Single Pomeron exchange (Fig. 1a) and multi-Pomeron exchanges (Figs. 1 b, 1c). Figure 2. Intermediate states that can be diffractively produced by a single Pomeron exchange. Figure 3. Diagram for diffractive vector mesons production (Fig. 3a) and triple-Regge diagram (Fig. 3b) with $`P`$ and $`f`$-exchanges in the $`t`$-channel ($`PPP`$ and $`PPf`$). Figure 4. Fan-type diagrams with Pomeron branchings. Figure 5. Triple Pomeron diagrams: a) triple Pomeron $`PPP`$ \- already represented in Fig. 3b; b) the interference term $`PfP`$. Figure 6. The $`\gamma ^{}p`$ cross section as a function of $`Q^2`$ for different energies compared to the following experimental data: H1 1995 (black points), ZEUS 1995 (black squares), E665 (black triangles) and ZEUSBPT97 (white circles). The data at $`Q^2=6.5`$ GeV<sup>2</sup> (H1 1994 -black stars- and ZEUS 1994 -white stars-) are from . Figure 7. $`F_2`$ as a function of $`x`$ for different values of $`Q^2`$. The experimental data are the same of Fig. 6. Figure 8. The function $`\mathrm{\Delta }_{eff}`$ as a function of $`Q^2`$ for 2 intervals of $`x`$. The solid line corresponds to the interval $`\stackrel{~}{x}=\frac{Q^2+s_0}{Q^2+s}=1.8510^4÷1.1610^5`$ and the dashed line to the interval $`\stackrel{~}{x}=7.4110^{14}÷4.6410^{15}`$. Figure 9. The diffractive structure function $`x_PF_2^{D(3)}`$ as a function of $`\beta `$ for fixed values of $`Q^2`$ and $`x_P`$. The experimental data (H1 1994) are from . Figure 10. The diffractive structure function $`x_PF_2^{D(3)}`$ as a function of $`x_P`$ for fixed values of $`Q^2`$ and $`\beta `$. The experimental data (H1 1994 and H1 1995) are from and . Figure 11. Diffractive production at $`Q^2=0`$ GeV for two different energies. The experimental data are from . Figure 12. The effective Pomeron slope $`\mathrm{\Delta }_{eff}^D`$ as a function of $`Q^2`$ for $`M_X=5`$ GeV. The experimental values are from ZEUS 1994 . Figure 13. Our predictions for the $`\gamma ^{}p`$ cross section at energies higher than those accessible at HERA. The experimental data are the same as in Figs. 6 and 7. Figure 1 Figure 2 Figure 3 Figure 4 Figure 5 Figure 6 Figure 7 Figure 8 Figure 9 Figure 10 Figure 11 Figure 12 Figure 13
warning/0005/math-ph0005009.html	ar5iv	text	# 1 Introduction ## 1 Introduction Let $`YX`$ be a smooth fibre bundle (throughout the paper, smooth manifolds are assumed to be real, finite-dimensional, Hausdorff, paracompact, and connected). We study cohomology of differential algebras of exterior forms on the infinite-order jet space $`J^{\mathrm{}}Y`$ of $`YX`$. This cohomology plays an important role in some physical models, e.g., in the calculus of variations in Lagrangian field theory \[1-4\] and in the field-antifield BRST formalism for constructing the descent equations \[5-8\]. Recall that the infinite-order jet space of a smooth fibre bundle $`YX`$ is defined as a projective limit $`(J^{\mathrm{}}Y,\pi _r^{\mathrm{}})`$ of the surjective inverse system $$X\stackrel{\pi }{}Y\stackrel{\pi _0^1}{}\mathrm{}J^{r1}Y\stackrel{\pi _{r1}^r}{}J^rY\mathrm{}$$ (1) of finite-order jet manifolds $`J^rY`$ . Provided with the projective limit topology, $`J^{\mathrm{}}Y`$ is a paracompact Fréchet (but not Banach) manifold . Given a bundle coordinate chart $`(\pi ^1(U_X);x^\lambda ,y^i)`$ on the fibre bundle $`YX`$, we have the coordinate chart $`((\pi ^{\mathrm{}})^1(U_X);x^\lambda ,y_\mathrm{\Lambda }^i)`$, $`0\|\mathrm{\Lambda }\|,`$ on $`J^{\mathrm{}}Y`$, together with the transition functions $$y_{}^{}{}_{\lambda +\mathrm{\Lambda }}{}^{i}=\frac{x^\mu }{x^\lambda }d_\mu y_\mathrm{\Lambda }^i,$$ (2) where $`\mathrm{\Lambda }=(\lambda _k\mathrm{}\lambda _1)`$, $`\|\mathrm{\Lambda }\|=k`$, is a multi-index, $`\lambda +\mathrm{\Lambda }`$ is the multi-index $`(\lambda \lambda _k\mathrm{}\lambda _1)`$ and $`d_\lambda `$ are the total derivatives $`d_\lambda =_\lambda +{\displaystyle \underset{\|\mathrm{\Lambda }\|=0}{}}y_{\lambda +\mathrm{\Lambda }}^i_i^\mathrm{\Lambda }.`$ In studying cohomology of differential algebras on the infinite-order jet manifold, the key point lies in the following fact (see Appendix A). LEMMA 1. A smooth fibre bundle $`Y`$ is a strong deformation retract of its infinite-order jet space $`J^{\mathrm{}}Y`$. Since $`J^{\mathrm{}}Y`$ is paracompact, it follows that there is an isomorphism $$H^{}(J^{\mathrm{}}Y,R)=H^{}(Y,R)$$ (3) of the cohomology groups of the infinite-order jet space $`J^{\mathrm{}}Y`$ with coefficients in the constant sheaf $`R`$ and those $`H^{}(Y,R)`$ of the fibre bundle $`Y`$. The goal is to construct a resolution of the constant sheaf $`R`$ on $`J^{\mathrm{}}Y`$ by acyclic sheaves of some differential algebra on $`J^{\mathrm{}}Y`$. Then the well-known theorem on cohomology of global sections of these sheaves (, Theorem 2.12.1) can be called into play. For the sake of convenience, we will agree to call it the general De Rham theorem. Given the surjective inverse system (1), we have the direct system $$𝔒_X^{}\stackrel{\pi ^{}}{}𝔒_Y^{}\stackrel{\pi _0^1^{}}{}𝔒_1^{}\stackrel{\pi _1^2^{}}{}\mathrm{}\stackrel{\pi _{r1}^r^{}}{}𝔒_r^{}\mathrm{}$$ (4) of ringed spaces $`(J^rY,𝔒_r^{})`$ whose structure sheaves $`𝔒_r^{}`$ are sheaves of differential $`R`$-algebras of exterior forms on finite-order jet manifolds $`J^rY`$, and $`\pi _{r1}^r^{}`$ are the pull-back morphisms. We follow the terminology of Ref. where by a sheaf is meant a sheaf bundle. The direct system (4) admits a direct limit $`𝔔_{\mathrm{}}^{}`$ which is a sheaf of differential exterior $`R`$-algebras on the infinite-order jet space $`J^{\mathrm{}}Y`$. Accordingly, we have the direct system $$𝒪^{}(X)\stackrel{\pi ^{}}{}𝒪^{}(Y)\stackrel{\pi _0^1^{}}{}𝒪_1^{}\stackrel{\pi _1^2^{}}{}\mathrm{}\stackrel{\pi _{r1}^r^{}}{}𝒪_r^{}\mathrm{}$$ (5) of the structure algebras $`𝒪_r^{}=\mathrm{\Gamma }(J^rY,𝔒_r^{})`$ of global sections of the sheaves $`𝔒_r^{}`$, i.e., $`𝒪_r^{}`$ are differential $`R`$-algebras of (global) exterior forms on finite-order jet manifolds $`J^rY`$. The direct limit $`(𝒪_{\mathrm{}}^{},\pi _r^{\mathrm{}})`$ of the direct system (5) is a differential $`R`$-algebra of all exterior forms on finite-order jet manifolds modulo the pull-back identification. Therefore, one usually thinks of elements of $`𝒪_{\mathrm{}}^{}`$ as being the pull-back onto $`J^{\mathrm{}}Y`$ of exterior forms on finite-order jet manifolds. It should be emphasized that the direct limit $`𝒪_{\mathrm{}}^{}`$ of the direct system (5) of structure algebras of sheaves $`𝔒_r^{}`$ fails to coincide with the structure algebra $`𝒬_{\mathrm{}}^{}=\mathrm{\Gamma }(J^{\mathrm{}}Y,𝔔_{\mathrm{}}^{})`$ of the direct limit $`𝔔_{\mathrm{}}^{}`$ of these sheaves. By definition, $`𝔔_{\mathrm{}}^{}`$ is the sheaf of germs of local exterior forms on finite-order jet manifolds. These local forms constitute a presheaf $`𝔒_{\mathrm{}}^{}`$ from which the sheaf $`𝔔_{\mathrm{}}^{}`$ is constructed. It means that, given a section $`\varphi \mathrm{\Gamma }(𝔔_{\mathrm{}}^{})`$ of $`𝔔_{\mathrm{}}^{}`$ over an open subset $`UJ^{\mathrm{}}Y`$ and any point $`qU`$, there exists a neighbourhood $`U_q`$ of $`q`$ such that $`\varphi \|_{U_q}`$ is the pull-back of a local exterior form on some finite-order jet manifold. However, $`𝔒_{\mathrm{}}^{}`$ does not coincide with the canonical presheaf $`\mathrm{\Gamma }(U,𝔔_{\mathrm{}}^{})`$ of sections of the sheaf $`𝔔_{\mathrm{}}^{}`$. There are obvious monomorphisms $`𝒪_{\mathrm{}}^{}𝒬_{\mathrm{}}^{}`$ and $`𝔒_{\mathrm{}}^{}\mathrm{\Gamma }(𝔔_{\mathrm{}}^{})`$. For short, we agree to call $`𝔔_{\mathrm{}}^{}`$ and $`𝒬_{\mathrm{}}^{}`$ the sheaf and algebra of locally pull-back exterior forms on $`J^{\mathrm{}}Y`$. Being restricted to a coordinate chart $`(\pi ^{\mathrm{}})^1(U_X)`$ on $`J^{\mathrm{}}Y`$, elements of $`𝒬_{\mathrm{}}^{}`$ can be written in the familiar coordinate form, where basic forms $`\{dx^\lambda \}`$ and contact 1-forms $`\{\theta _\mathrm{\Lambda }^i=dy_\mathrm{\Lambda }^iy_{\lambda +\mathrm{\Lambda }}^idx^\lambda \}`$ provide the local generators of the algebra $`𝒬_{\mathrm{}}^{}`$. There is the canonical splitting of the space of $`m`$-forms $`𝒬_{\mathrm{}}^m=𝒬_{\mathrm{}}^{0,m}𝒬_{\mathrm{}}^{1,m1}\mathrm{}𝒬_{\mathrm{}}^{m,0}`$ into spaces $`𝒬_{\mathrm{}}^{k,mk}`$ of $`k`$-contact forms. Accordingly, the exterior differential on $`𝒬_{\mathrm{}}^{}`$ is decomposed into the sum $`d=d_H+d_V`$ of horizontal and vertical differentials $`d_H:𝒬_{\mathrm{}}^{k,s}𝒬_{\mathrm{}}^{k,s+1},d_V:𝒬_{\mathrm{}}^{k,s}𝒬_{\mathrm{}}^{k+1,s}`$ which obey the nilpotency rule $$d_Hd_H=0,d_Vd_V=0,d_Vd_H+d_Hd_V=0.$$ (6) Traditionally, one has tried to introduce the algebra of locally pull-back forms on $`J^{\mathrm{}}Y`$ in a standard geometric way . The difficulty lies in the geometric interpretation of derivations of the $`R`$-ring $`𝒬_{\mathrm{}}^0`$ of locally pull-back functions as vector fields on the Fréchet manifold $`J^{\mathrm{}}Y`$ and their duals as differential forms on $`J^{\mathrm{}}Y`$. Therefore, one usually considers the subalgebra $`𝒪_{\mathrm{}}^{}`$ of pull-back exterior forms on $`J^{\mathrm{}}Y`$, given as the direct limit of the direct system (5). Accordingly, the infinite-order De Rham complex of these forms $$0R𝒪_{\mathrm{}}^0\stackrel{d}{}𝒪_{\mathrm{}}^1\stackrel{d}{}\mathrm{}$$ (7) is the direct limit of the De Rham complexes of exterior forms on finite-order jet manifolds. Hence, as was repeatedly proved, the cohomology groups $`H^{}(𝒪_{\mathrm{}}^{})`$ of the complex (7) are equal to the De Rham cohomology groups $`H^{}(Y)`$ of the fibre bundle $`Y`$ \[1-4\]. At the same time, $`d_H`$\- and $`d_V`$-cohomology of the differential algebra $`𝒪_{\mathrm{}}^{}`$ remains unknown. The algebraic Poincaré lemma on the local exactness of differentials $`d_H`$ and $`d_V`$ only has been repeatedly proved (see, e.g., \[13-15\]). Here, we show that the problem of cohomology of differential forms on the infinite-order jet space $`J^{\mathrm{}}Y`$ has a comprehensive solution by enlarging the algebra $`𝒪_{\mathrm{}}^{}`$ to the algebra $`𝒬_{\mathrm{}}^{}`$. From the physical viewpoint, it enables one also to study effective field theories whose Lagrangians involve derivatives of arbitrary high order. The key point is that the Fréchet manifold $`J^{\mathrm{}}Y`$ admits a partition of unity performed by elements of the ring $`𝒬_{\mathrm{}}^0`$ of locally pull-back functions on $`J^{\mathrm{}}Y`$ . It follows that the sheaf $`𝔔_{\mathrm{}}^{}`$ of $`𝒬_{\mathrm{}}^0`$-modules on $`J^{\mathrm{}}Y`$ is fine and, consequently, acyclic. Therefore, studying different resolutions performed by subsheaves of the sheaf $`𝔔_{\mathrm{}}^{}`$, one may hope to get a complete picture of cohomology of the differential algebra $`𝒬_{\mathrm{}}^{}`$. Given a smooth fibre bundle $`YX`$, we will show the following. * 1. The De Rham cohomology groups of the differential algebra $`𝒬_{\mathrm{}}^{}`$ are isomorphic to those of the fibre bundle $`Y`$. 2. Its $`d_V`$-cohomology groups are related to the cohomology groups of the fibre bundle $`Y`$ with coefficients in the sheaf $`𝔒_X^{}`$ of exterior forms on the base $`X`$. 3. The $`d_H`$-cohomology groups of contact elements $`\varphi 𝒬_{\mathrm{}}^{0<,}`$ of the algebra $`𝒬_{\mathrm{}}^{}`$ are trivial. 4. In degrees $`r<n=\mathrm{dim}X`$, the $`d_H`$-cohomology groups of its horizontal elements $`\varphi 𝒬_{\mathrm{}}^{0,}`$ coincide with the De Rham cohomology groups of the fibre bundle $`Y`$. Note that, as was mentioned above, the result (i) is also true for the differential algebra $`𝒪_{\mathrm{}}^{}`$. The result (iii) recovers that in Refs. , obtained by means of the Mayer-Vietoris sequence. Point out the particular case of an affine bundle $`YX`$, interesting for physical applications, e.g., to BRST theory. In this case, $`X`$ is a strong deformation retract of $`Y`$ and the cohomology of $`Y`$ under consideration is equal to that of $`X`$. Then the above results are reformulated as follows. * 1. Any closed form $`\varphi 𝒬_{\mathrm{}}^{}`$ is decomposed into the sum $`\varphi =\phi +d\xi `$ where $`\phi 𝒪^{}(X)`$ is a closed form on the base $`X`$. 2. Any $`d_V`$-closed form $`\varphi 𝒬_{\mathrm{}}^{}`$ is $`d_V`$-exact. 3. Any $`d_H`$-closed form $`\varphi 𝒬_{\mathrm{}}^{}`$ is decomposed into the sum $`\varphi =\phi +d_H\xi `$ where $`\phi 𝒪^{}(X)`$ is a closed form on the base $`X`$. The results (i) and (ii) are also true for the differential algebra $`𝒪_{\mathrm{}}^{}`$. ## 2 De Rham cohomology Let us start from De Rham cohomology of the differential algebra $`𝒬_{\mathrm{}}^{}`$ of locally pull-back exterior forms on the infinite-order jet manifold $`J^{\mathrm{}}Y`$. We consider the complex of sheaves of $`𝒬_{\mathrm{}}^0`$-modules $$0R𝔔_{\mathrm{}}^0\stackrel{d}{}𝔔_{\mathrm{}}^1\stackrel{d}{}\mathrm{}$$ (8) on the infinite-order jet space $`J^{\mathrm{}}Y`$ and the corresponding infinite-order De Rham complex $$0R𝒬_{\mathrm{}}^0\stackrel{d}{}𝒬_{\mathrm{}}^1\stackrel{d}{}\mathrm{}$$ (9) of algebras of locally pull-back exterior forms on $`J^{\mathrm{}}Y`$. Since locally pull-back exterior forms fulfill the Poincaré lemma, the complex of sheaves (8) is exact. Since the paracompact space $`J^{\mathrm{}}Y`$ admits a partition of unity performed by elements of $`𝒬_{\mathrm{}}^0`$, the sheaves $`𝔔_{\mathrm{}}^r`$ of $`𝒬_{\mathrm{}}^0`$-modules are fine for all $`r0`$. Then they are acyclic, i.e., the cohomology groups $`H^{>0}(J^{\mathrm{}}Y,𝔔_{\mathrm{}}^r)`$ of the paracompact space $`J^{\mathrm{}}Y`$ with coefficients in sheaves $`𝔔_{\mathrm{}}^r`$ vanish. Consequently, the exact sequence (8) is a fine resolution of the constant sheaf $`R`$ on $`J^{\mathrm{}}Y`$. Then, in accordance with the above mentioned general De Rham theorem, we have an isomorphism $$H^{}(𝒬_{\mathrm{}}^{})=H^{}(J^{\mathrm{}}Y,R)$$ (10) of the De Rham cohomology groups $`H^{}(𝒬_{\mathrm{}}^{})`$ of the differential algebra $`𝒬_{\mathrm{}}^{}`$ and the cohomology groups $`H^{}(J^{\mathrm{}}Y,R)`$ of the infinite-order jet space $`J^{\mathrm{}}Y`$ with coefficients in the constant sheaf $`R`$. Combining isomorphisms (3) and (10), we come to the manifested assertion. PROPOSITION 2. There is an isomorphism $`H^{}(𝒬_{\mathrm{}}^{})=H^{}(Y)`$ of the De Rham cohomology groups $`H^{}(𝒬_{\mathrm{}}^{})`$ of the differential algebra $`𝒬_{\mathrm{}}^{}`$ to the De Rham cohomology groups $`H^{}(Y)`$ of the fibre bundle $`Y`$. ## 3 Cohomology of $`d_V`$ Due to the nilpotency rule (6), the vertical and horizontal differentials $`d_V`$ and $`d_H`$ define the bicomplex of sheaves $$\begin{array}{ccccccccccccccccccc}& & & & \hfill _{d_V}& \text{}\hfill & & \hfill _{d_V}& \text{}\hfill & & & \hfill _{d_V}& \text{}\hfill & & & \hfill _{d_V}& \text{}\hfill & & \\ & & 0& & & 𝔔_{\mathrm{}}^{k,0}\hfill & \stackrel{d_H}{}& & 𝔔_{\mathrm{}}^{k,1}\hfill & \stackrel{d_H}{}& \mathrm{}& & 𝔔_{\mathrm{}}^{k,m}\hfill & \stackrel{d_H}{}& \mathrm{}& & 𝔔_{\mathrm{}}^{k,n}\hfill & & \\ & & & & & \mathrm{}\hfill & & & \mathrm{}\hfill & & & & \mathrm{}\hfill & & & & \mathrm{}\hfill & & \\ & & & & \hfill _{d_V}& \text{}\hfill & & \hfill _{d_V}& \text{}\hfill & & & \hfill _{d_V}& \text{}\hfill & & & \hfill _{d_V}& \text{}\hfill & & \\ 0& & R& & & 𝔔_{\mathrm{}}^0\hfill & \stackrel{d_H}{}& & 𝔔_{\mathrm{}}^{0,1}\hfill & \stackrel{d_H}{}& \mathrm{}& & 𝔔_{\mathrm{}}^{0,m}\hfill & \stackrel{d_H}{}& \mathrm{}& & 𝔔_{\mathrm{}}^{0,n}\hfill & & \\ & & & & \hfill _\pi ^{\mathrm{}}& \text{}\hfill & & \hfill _\pi ^{\mathrm{}}& \text{}\hfill & & & \hfill _\pi ^{\mathrm{}}& \text{}\hfill & & & \hfill _\pi ^{\mathrm{}}& \text{}\hfill & & \\ 0& & R& & & 𝔒_X^0\hfill & \stackrel{d}{}& & 𝔒_X^1\hfill & \stackrel{d}{}& \mathrm{}& & 𝔒_X^m\hfill & \stackrel{d}{}& \mathrm{}& & 𝔒_X^n\hfill & \stackrel{d}{}& 0\\ & & & & & \text{}\hfill & & & \text{}\hfill & & & & \text{}\hfill & & & & \text{}\hfill & & \\ & & & & & 0\hfill & & & 0\hfill & & & & 0\hfill & & & & 0\hfill & & \end{array}$$ (11) \[1-4,10,13-15\]. The rows and columns of these bicomplex are horizontal and vertical complexes. Moreover, the above mentioned algebraic Poincaré lemma is obviously extended to elements of $`𝒬_{\mathrm{}}^{}`$ and leads to the following. LEMMA 3. The columns and rows of the bicomplex (11) are exact sequences of sheaves. It follows that, since all sheaves $`𝔔_{\mathrm{}}^{k,m}`$ are fine, the columns and rows of the bicomplex (11) are fine resolutions of their first terms. Then the general De Rham theorem on a resolution of a sheaf on a paracompact manifold can be used again. Let us consider a vertical exact sequence of sheaves $$0𝔒_X^m\stackrel{\pi ^{\mathrm{}}}{}𝔔_{\mathrm{}}^{0,m}\stackrel{d_V}{}\mathrm{}\stackrel{d_V}{}𝔔_{\mathrm{}}^{k,m}\stackrel{d_V}{}\mathrm{},mn,$$ (12) and the corresponding complex of their structure algebras $$0𝒬^m(X)\stackrel{\pi ^{\mathrm{}}}{}𝒬_{\mathrm{}}^{0,m}\stackrel{d_V}{}\mathrm{}\stackrel{d_V}{}𝒬_{\mathrm{}}^{k,m}\stackrel{d_V}{}\mathrm{},mn.$$ (13) The exact sequence (12) is a resolution of the sheaf $`\pi ^{\mathrm{}}𝔒_X^m`$ of the pull-back onto $`J^{\mathrm{}}Y`$ of exterior forms on the base $`X`$. Then, by virtue of the general De Rham theorem, we have an isomorphism $$H^{}(d_V,m)=H^{}(J^{\mathrm{}}Y,\pi ^{\mathrm{}}𝔒_X^m)$$ (14) of the cohomology groups $`H^{}(d_V,m)`$ of the complex of differential algebras (13) to the cohomology groups $`H^{}(J^{\mathrm{}}Y,\pi ^{\mathrm{}}𝔒_X^m)`$ of the infinite-order jet manifold $`J^{\mathrm{}}Y`$ with coefficients in the sheaf $`\pi ^{\mathrm{}}𝔒_X^m`$. Since $`J^{\mathrm{}}Y`$ and $`Y`$ are homotopic, there is an isomorphism of cohomology groups $`H^{}(J^{\mathrm{}}Y,\pi ^{\mathrm{}}𝔒_X^m)=H^{}(Y,\pi ^{}𝔒_X^m).`$ Thus, it is stated the following. PROPOSITION 4. There is an isomorphism of the cohomology groups $`H^{}(d_V,m)=H^{}(Y,\pi ^{}𝔒_X^m).`$ In particular, if $`YX`$ is an affine bundle, we have $`H^{}(d_V,m)=H^{}(Y,\pi ^{}𝔒_X^m)=H^{}(X,𝔒_X^m)=0`$ because the sheaf $`𝔒_X^m`$ on $`X`$ is fine. This result also follows directly from the expression for the corresponding homotopy operator and, therefore, remains true for to the differential algebra $`𝒪_{\mathrm{}}^{}`$ (see Appendix B). ## 4 Cohomology of $`d_H`$ Turn now to the rows of the bicomplex (11) (excluding the bottom one which is obviously the De Rham complex on the base $`X`$). We have the exact sequences of sheaves $`0𝔔_{\mathrm{}}^{k,0}\stackrel{d_H}{}𝔔_{\mathrm{}}^{k,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}𝔔_{\mathrm{}}^{k,n},k>0,`$ (15) $`0R𝔔_{\mathrm{}}^0\stackrel{d_H}{}𝔔_{\mathrm{}}^{0,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}𝔔_{\mathrm{}}^{0,n}.`$ (16) Speaking rigorously, the exact sequences (15) and (16) fail to be fine resolutions of the sheaves $`𝔔_{\mathrm{}}^{k,0}`$ and $`R`$, respectively, because of their last terms. At the same time, following directly the proof of the above mentioned Theorem 2.12.1 in Ref. till these terms, one can show the following. PROPOSITION 5. The cohomology groups $`H^r(d_H,k)`$, $`r<n`$, of the complex $`0𝒬_{\mathrm{}}^{k,0}\stackrel{d_H}{}𝒬_{\mathrm{}}^{k,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}𝒬_{\mathrm{}}^{k,n}`$ are isomorphic to the cohomology groups $`H^r(J^{\mathrm{}}Y,𝔔_{\mathrm{}}^{k,0})`$ of $`J^{\mathrm{}}Y`$ with coefficients in the sheaf $`𝔔_{\mathrm{}}^{k,0}`$ and, consequently, are trivial because the sheaf $`𝔔_{\mathrm{}}^{k,0}`$ is fine. PROPOSITION 6. The cohomology groups $`H^r(d_H)`$, $`r<n`$, of the complex $`0R𝒬_{\mathrm{}}^0\stackrel{d_H}{}𝒬_{\mathrm{}}^{0,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}𝒬_{\mathrm{}}^{0,n}`$ are isomorphic to the cohomology groups $`H^r(J^{\mathrm{}}Y,R)`$ of $`J^{\mathrm{}}Y`$ with coefficients in the constant sheaf $`R`$ and, consequently, to the cohomology groups $`H^r(Y,R)`$. Note that one can also study the exact sequences of presheaves $`0𝔒_{\mathrm{}}^{k,0}\stackrel{d_H}{}𝔒_{\mathrm{}}^{k,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}𝔒_{\mathrm{}}^{k,n},k>0,`$ $`0R𝔒_{\mathrm{}}^0\stackrel{d_H}{}𝔒_{\mathrm{}}^{0,1}\stackrel{d_H}{}\mathrm{}\stackrel{d_H}{}𝔒_{\mathrm{}}^{0,n},`$ but comes again to the results of Propositions 4, 4. Because $`J^{\mathrm{}}Y`$ is paracompact, the cohomology groups $`H^{}(J^{\mathrm{}}Y,𝔔_{\mathrm{}}^,)`$ of $`J^{\mathrm{}}Y`$ with coefficients in a sheaf $`𝔔_{\mathrm{}}^,`$ and those $`H^{}(J^{\mathrm{}}Y,𝔒_{\mathrm{}}^,)`$ with coefficients in a presheaf $`𝔒_{\mathrm{}}^,`$ are isomorphic. It follows that the cohomology group $`H^0(J^{\mathrm{}}Y,𝔒_{\mathrm{}}^,)`$ of a presheaf $`𝔒_{\mathrm{}}^,`$ is isomorphic to the $`R`$-module $`𝒬_{\mathrm{}}^,=H^0(J^{\mathrm{}}Y,𝔔_{\mathrm{}}^,)`$, but not $`𝒪_{\mathrm{}}^,`$. ## 5 Appendix A Here, we construct a desired homotopy from the infinite order jet manifold $`J^{\mathrm{}}Y`$ to a fibre bundle $`Y`$ in an explicit form. Given a coordinate chart $`((\pi ^{\mathrm{}})^1(U_X);x^\lambda ,y_\mathrm{\Lambda }^i)`$ on $`J^{\mathrm{}}Y`$, let us consider the map $`[0,1]\times J^{\mathrm{}}Y(t;x^\lambda ,y^i,y_\mathrm{\Lambda }^i)(x^\lambda ,y^i,y_\mathrm{\Lambda }^i)J^{\mathrm{}}Y,0<\|\mathrm{\Lambda }\|,`$ $`y_\mathrm{\Lambda }^i=f_k(t)y_\mathrm{\Lambda }^i+(1f_k(t))\mathrm{\Gamma }_{(k)}{}_{\mathrm{\Lambda }}{}^{i}),\|\mathrm{\Lambda }\|=k>0,`$ where $`\mathrm{\Gamma }_{(k)}`$ is a section of the affine jet bundle $`J^kYJ^{k1}Y`$ and $`f_k(t)`$ is a continuous (smooth) real function on $`[0,1]`$ such that $`f_k(t)=\{\begin{array}{cc}0,\hfill & t12^k,\hfill \\ 1,\hfill & t12^{(k+1)}.\hfill \end{array}`$ A glance at the transition functions (2) shows that, given in a coordinate form, this map is well-defined. It is a desired homotopy from $`J^{\mathrm{}}Y`$ to $`Y`$ which is identified with its image under the global section $`\gamma :YJ^{\mathrm{}}Y,y_\mathrm{\Lambda }^i\gamma =\mathrm{\Gamma }_{(k)}{}_{\mathrm{\Lambda }}{}^{i}\mathrm{\Gamma }_{(k1)}\mathrm{}\mathrm{\Gamma }_1,\|\mathrm{\Lambda }\|=k.`$ ## 6 Appendix B We start from a remark that, Studying cohomology of exterior forms on the infinite-order jet space of an affine bundle, we can restrict our consideration to vector bundles $`YX`$ without loss of generality as follows. Let $`YX`$ be a smooth affine bundle modelled over a smooth vector bundle $`\overline{Y}X`$. A glance at the transformation law (2) shows that $`J^{\mathrm{}}YX`$ is an affine topological bundle modelled on the vector bundle $`J^{\mathrm{}}\overline{Y}X`$. This affine bundle admits a global section $`J^{\mathrm{}}s`$ which is the infinite-order jet prolongation of a global section $`s`$ of $`YX`$. With $`J^{\mathrm{}}s`$, we have a homeomorphism $`\widehat{s}_{\mathrm{}}:J^{\mathrm{}}Yqq(J^{\mathrm{}}s)(\pi ^{\mathrm{}}(q))J^{\mathrm{}}\overline{Y}`$ of the topological spaces $`J^{\mathrm{}}Y`$ and $`J^{\mathrm{}}\overline{Y}`$, together with an exterior algebra isomorphism $`\widehat{s}_{\mathrm{}}^{}:\overline{𝒪}_{\mathrm{}}^{}𝒪_{\mathrm{}}^{}`$. Moreover, it is readily observed that the pull-back morphism $`\widehat{s}_{\mathrm{}}^{}`$ commutes with the differentials $`d`$, $`d_V`$ and $`d_H`$. Therefore, the differential algebras $`𝒬_{\mathrm{}}^{}`$ and $`\overline{𝒬}_{\mathrm{}}^{}`$ (as like as $`𝒪_{\mathrm{}}^{}`$ and $`\overline{𝒪}_{\mathrm{}}^{}`$) have the same $`d`$-, $`d_V`$\- and $`d_H`$-cohomology. Let $`YX`$ be a vector bundle. Let us consider the vertical complex $`0𝒪^m(X)\stackrel{\pi ^{\mathrm{}}}{}𝒪_{\mathrm{}}^{0,m}\stackrel{d_V}{}\mathrm{}\stackrel{d_V}{}𝒪_{\mathrm{}}^{k,m}\stackrel{d_V}{}\mathrm{},mn,`$ of differential algebras of pull-back exteriror forms on $`J^{\mathrm{}}Y`$. Its local exactness on a coordinate chart $`((\pi ^{\mathrm{}})^1(U_X);x^\lambda ,y_\mathrm{\Lambda }^i)`$, $`0\|\mathrm{\Lambda }\|,`$ on $`J^{\mathrm{}}Y`$ follows from a version of the Poincaré lemma with parameters (see, e.g., ). We have the corresponding homotopy operator $`\sigma ={\displaystyle _0^1}t^k[\overline{y}\varphi (x^\lambda ,ty_\lambda ^i)]dt,\varphi \text{}^{k,m}_{\mathrm{}},`$ where $`\overline{y}=y_\mathrm{\Lambda }^i_i^\mathrm{\Lambda }`$. Since $`YX`$ is a vector bundle, it is readily observed that, given in a coordinate form, this homotopy operator is globally defined on $`J^{\mathrm{}}Y`$, and so is the exterior form $`\sigma `$.
warning/0005/astro-ph0005214.html	ar5iv	text	# Discovery of a Young Radio Pulsar in a Relativistic Binary Orbit ## 1 Introduction Relativistic binary pulsars have been celebrated as being unique laboratories for high-precision tests of general relativity and accurate determinations of neutron star masses (Taylor & Weisberg (1989); Taylor et al. (1992); Stairs et al. (1998)). To date, post-Keplerian general relativistic parameters have been measured for nine binary pulsar systems (see reviews by Taylor et al. (1992); Thorsett & Chakrabarty (1999); Kaspi (1999)). Five are double neutron star binaries (PSRs B1913+16, B1534+12, B2127+11C, J1518+4904, J1811$``$1736) and four (PSRs J1713+0747, B1802$``$07, B1855+09, B2303+46) have white dwarf companions. Through radio timing observations of PSR B1913+16, general relativity has been confirmed at the $``$0.3% level (Taylor (1992)). Timing observations of these pulsars show that the neutron star masses fall in a narrow range centered on 1.35 M (Thorsett & Chakrabarty (1999)). The lack of variation of these masses is surprising given the variety of binary evolution mechanisms by which they were formed. Of the above relativistic binary systems, all but one are likely to share a similar evolutionary history. The observed pulsar in most of these systems has a large characteristic age $`\tau _c`$ and low surface magnetic field $`B`$ relative to the bulk of the isolated pulsar population. This is because the observed neutron star, the first-born in the binary, was “recycled” by a phase of mass transfer from its companion as the latter ascended the giant branch (see Bhattacharya & van den Heuvel 1991 for details of neutron-star binary evolution). This results in spin up of the neutron star and the reduction of its magnetic field. The binary observed today has survived the supernova explosion of the secondary (although in the cases of PSRs B2127+11C and B1802$``$07, which are in globular clusters, the possibility of a more complicated evolutionary history cannot be discounted). The exception among the above binaries is PSR B2303+46 which shows no evidence for having been recycled, having $`B`$ and $`\tau _c`$ comparable to the bulk of the isolated pulsar population. PSR B2303+46, long thought to be part of a double neutron-star binary, has recently been shown to have a white-dwarf companion (van Kerkwijk & Kulkarni (1999)) which must have formed prior to the pulsar’s birth. Formation mechanisms for such a system have been suggested by Portegies Zwart & Yungelson (1999) and Tauris & Sennels (2000) and involve an initial primary, the white dwarf progenitor, which transferred sufficient matter onto the initial secondary that the latter underwent core collapse. We report here on PSR J1141$``$6545, a relativistic binary pulsar recently discovered as part of the Parkes Multibeam Pulsar survey. A preliminary report on this system was presented by Manchester et al. (2000). ## 2 Observations and Results ### 2.1 Discovery Observations The Parkes multibeam pulsar survey (Lyne et al. (2000)), currently underway at the 64-m radio telescope in Parkes, Australia, makes use of a multibeam receiver which operates at a wavelength of 20 cm and permits simultaneous observations of 13 separate regions of the sky (Staveley-Smith et al. (1996)). The survey, roughly halfway complete, has been extremely successful, having discovered over 500 new radio pulsars (Camilo et al. (2000); D’Amico et al. (2000)). The observing system used has been described by Lyne et al. (2000). Observations are centered on 1374 MHz, with 288 MHz of bandwidth in each of two orthogonal polarizations, for each of the 13 beams. The filtered and amplified signals are fed into a $`2\times 13\times 96\times 3`$ MHz filterbank system, then are square-law detected. The polarization pairs are summed, and the results are 1-bit digitized every 250 $`\mu `$s and recorded on tape. Observations are of 35-min duration. Offline, the data are dedispersed at many trial dispersion measures (DMs) and subjected to $`2^{23}`$-point Fast Fourier Transforms. PSR J1141$``$6545 was discovered in our standard analysis. A periodicity at 394 ms and at DM of 118 pc cm<sup>-3</sup> was reported as having signal-to-noise ratio 15.4 in data obtained on 1999 April 28. Immediately apparent in the discovery observation was the variation in the pulse period due to binary acceleration. Had the pulsar been isolated, its signal-to-noise ratio would have been 40. Confirmation observations were made at Parkes on 1999 July 10. A series of observations in the following week determined the binary parameters. ### 2.2 Timing Observations Regular timing observations of PSR J1141$``$6545 were carried out at numerous epochs from 1999 July 10 through 2000 January 19 using the central beam of the multibeam receiver. In the timing observations, the same data acquisition and online software system are used as in the regular survey observations. At 1374 MHz, we obtained 123 pulse arrival times at 3-MHz frequency resolution, and 31 pulse arrival times using a higher resolution $`2\times 512\times 0.5`$-MHz filterbank. Two observations were made at 600 MHz, on 1999 August 23 and 24, using a $`2\times 256\times 0.125`$-MHz filterbank and the standard data acquisition system. Timing observations were divided into 1-min sub-integrations which were partially de-dispersed into 8 sub-bands and folded at the predicted topocentric pulsar period using 1024 pulse phase bins. Typical total integration times were $``$10 min. Final dedispersion of the sub-bands and folding of the sub-integrations were done in a subsequent processing step. Each pulse profile was then cross-correlated with a high signal-to-noise template appropriate for the observing frequency and filterbank resolution. The template for the high-resolution 20-cm observations is shown in Figure 1. The other templates appear similar, including that at 600 MHz. The pulse times-of-arrival (TOAs) which result from the cross-correlation were then subject to a timing analysis using the TEMPO software package<sup>1</sup><sup>1</sup>1see http://pulsar.princeton.edu/tempo which converts them to barycentric TOAs at infinite frequency using the Jet Propulsion Laboratory DE200 solar-system ephemeris (Standish (1982)) and performs a multi-parameter least-squares fit for the pulsar spin and binary parameters. During the fitting, the pulsar’s position was held fixed at the value determined by interferometric observations (see §2.3). Also, the DM was determined using the two 600-MHz observations and six 1400-MHz observations obtained on 1999 August 22 and 23, and was held fixed for the rest of the timing analysis. The results of the timing analysis are presented in Table 1. The fit parameters were $`P`$, $`\dot{P}`$, the five Keplerian orbital parameters, and a rate of periastron advance, $`\dot{\omega }`$. Upper limits on other post-Keplerian parameters were determined individually. Post-fit timing residuals are shown in Figure 2. The uncertainties in the arrival times were determined empirically by demanding that the reduced $`\chi ^2`$ for observations over a short span (1–2 days) be unity. The formal reduced $`\chi ^2`$ for these residuals is 4.7 (145 degrees of freedom). This clearly indicates that there are unmodeled effects in the data. That the residuals as a function of orbital phase do not display systematic trends leads us to believe the orbital parameters are not biased by these effects. Indeed the two arrival times with the largest residuals, near epoch 1999.96, are at orbital phases separated by $``$0.5. The unmodeled features may be manifestations of standard “timing noise,” a stochastic process seen in many young pulsars like PSR J1141$``$6545 (e.g. Arzoumanian et al. (1994)). This may bias the measured spin parameters slightly. ### 2.3 Radio Interferometric Observations PSR J1141$``$6545 was observed on 1999 August 22 and 23 with the Australia Telescope Compact Array (ATCA) which is located in Narrabri, Australia. The array was in its “6D” configuration, which has maximum baseline 5878 m. The observations were done in pulsar gating mode at center frequencies of 1384 and 2496 MHz simultaneously, with 128 MHz of bandwidth in each of two linear polarizations for each frequency. The primary and secondary calibrators for the observations were the bright radio point sources 1934$``$638 and 1329$``$665. The data were processed using the MIRIAD processing package.<sup>2</sup><sup>2</sup>2http://www.atnf.csiro.au/computing/software/miriad/ After standard interference flagging and calibration were done, on- and off-pulse images were formed. The off-pulse emission was then subtracted from the UV data set leaving only on-pulse emission from the pulsar. A point source was then fitted to these UV data using the nominal position and a rough estimate of the flux density of the pulsar. A model of the resulting UV data was then fitted using trial positions and flux densities for the pulsar using the MIRIAD UVFIT task. Best-fit values and uncertainties were obtained for both the 1384- and 2496-MHz data sets. Since the pulsar is brighter at 1384 MHz, we use the more accurate fitted position from that data set in the timing analysis (see Table 1). The 1384-MHz pulsed flux ($``$3.3 mJy) is consistent with that estimated from the Parkes observations; the $`3\sigma `$ upper limit on point-like radio emission off-pulse is 0.66 mJy at this frequency. ### 2.4 Polarization We have measured the polarization properties of PSR J1141$``$6545, which are especially interesting given the possibility of significant precession of the pulsar (see §3.3 below). Observations were made on 2000 March 26 (MJD 51629) using the center beam of the multibeam receiver and the Caltech correlator (Navarro 1994). Instrumental phases and gains were calibrated using a pulsed noise signal injected at $`45^{}`$ to the two signal probes (see Navarro et al. 1997 for a full description). Two 25-min observations using 64-MHz bandwidth at orthogonal feed angles were summed to produce the polarization profiles shown in Figure 3. Position angles were corrected for Faraday rotation in the Earth’s ionosphere. The ionospheric contribution to the rotation measure (RM) was estimated to be $`3.2`$ rad m<sup>-2</sup> at the time of the observation using an ionospheric model, so the plotted angles are approximately $`9^{}`$ greater than those observed. The precision of the absolute position angles is limited by the uncertainty in this correction, which we estimate to be $`\pm 2^{}`$. On average, the pulse profile is 16% linearly polarized (see Figure 3). The position angle shows a complicated variation across the pulse with a clear orthogonal transition in the leading wing and a sharp drop toward the trailing edge of the pulse. The circular polarization is relatively strong, with a mean value of $`10`$%. There is a possible reversal in the sign of $`V`$ at a pulse phase slightly preceding the orthogonal transition in the linear polarization. These properties are not consistent with any simple model for the polarization. The shape of the total intensity profile suggests that the main peak is near the trailing edge of a conal beam. This may be supported by the location of the reversal of sign of $`V`$ which is usually (but not always) located near the center of the beam (Han et al. 1998). However, the observed position angle variation is certainly not consistent with a rotating vector model having a beam center on the leading wing. The interstellar component of the RM, estimated by computing the weighted mean position angle difference between the two halves of the observed band, is $`86\pm 3`$ rad m<sup>-2</sup>. This implies a mean line-of-sight interstellar field, weighted by the electron density along the path, of 0.9 $`\mu `$G directed away from us. RMs of neighboring pulsars are of mixed sign and magnitude (Han, Manchester and Qiao 1999). At the estimated distance of 3.2 kpc, as inferred from the DM and a model for the Galactic electron distribution (Taylor & Cordes (1993)), the path to the pulsar crosses the Carina arm; the structure of the Galactic magnetic field in this region is evidently complex. ## 3 Discussion PSR J1141$``$6545 is different from most of the other known relativistic binary pulsars in several important respects. The pulsar’s characteristic age, $`\tau _c=1.4\times 10^6`$ yr, and inferred surface magnetic field strength, $`B=1.3\times 10^{12}`$ G, are similar to those of the bulk of the isolated pulsar population. Thus, PSR J1141$``$6545 is unlikely to have ever been recycled. Like the other short-$`P_b`$, eccentric relativistic binaries, it may be a double neutron star binary, but if so, we must be observing the second-formed neutron star in its short-lived radio pulsar phase. In this case, its companion could still be an observable radio pulsar that was not detected in our Parkes observations, either because its short spin period varies too rapidly due to the binary orbit (a possibility we are checking, given our knowledge of the binary orbit) or because its radio beam does not intersect our line of sight. However, we show here that the companion is unlikely to be a second neutron star, given the system’s mass function, its rate of precession $`\dot{\omega }`$, and what is known about the neutron star mass distribution. The mass function is given by $$f(M_p)=\frac{4\pi ^2(a\mathrm{sin}i)^3}{GP_b^2}=\frac{M_c^3\mathrm{sin}^3i}{(M_p+M_c)^2},$$ (1) where $`M_p`$ and $`M_c`$ are the pulsar and companion masses, respectively, $`P_b`$ is the orbital period, and $`a\mathrm{sin}i`$ is the projected semi-major axis of the pulsar orbit. This provides a constraint on the minimum $`M_c`$ for any assumed $`M_p`$, by setting $`i`$, the inclination of the orbital angular momentum with respect to the line-of-sight, to $`90^{}`$. In addition, $`\dot{\omega }`$, under the assumption that it is due to general relativistic periastron precession (but see §3.2), is given by $$\dot{\omega }_{GR}=3\left(\frac{P_b}{2\pi }\right)^{5/3}(T_{}M)^{2/3}(1e^2)^1,$$ (2) where $`T_{}=4.925490947\times 10^6`$ s. This then determines the total system mass $`M=M_p+M_c`$. For $`\dot{\omega }=(5.32\pm 0.02)^{}`$yr<sup>-1</sup> (Table 1), $`M=(2.300\pm 0.012)`$ M. Figure 4 shows these constraints in $`M_pM_c`$ phase space. The shaded region is ruled out by the mass function. The intersection of the $`\dot{\omega }`$ straight line with the boundary of this region determines the maximum allowed pulsar mass: $`M_p<1.331`$ M ($`1\sigma `$) or $`M_p<1.348`$ M ($`3\sigma `$). Similarly, we set a lower limit $`M_c>0.974`$ M ($`1\sigma `$) or $`M_c>0.968`$ M ($`3\sigma `$). Thorsett & Chakrabarty (1999) showed that measurements of neutron star masses are consistent with all being drawn from a Gaussian distribution having mean and standard deviation 1.35 and 0.04 M, respectively. The parameters in the PSR J1141$``$6545 system are thus interesting: if $`M_p`$ is close to its maximum allowed value corresponding to $`i90^{}`$, as would be consistent with other neutron star masses, then its companion has mass much lower than those of all known neutron stars, and therefore is unlikely to be one. Of course on a posteriori statistical grounds, $`i90^{}`$ is improbable. On the other hand, a smaller $`i`$ implies a less massive pulsar. For the median value $`i=60^{}`$, $`M_p=1.17`$ M, and $`M_c=1.13`$ M, both significantly lower than for any other known neutron stars. Note that none of these values can be ruled out based on neutron star stability arguments, since masses of as little as $``$0.1 M are allowed (e.g. Colpi, Shapiro, & Teukolsky (1993)). However the formation mechanism for such low-mass neutron stars is unclear. ### 3.1 A Neutron Star/White Dwarf Binary? A more likely possibility is that the PSR J1141$``$6545 system is a neutron star/CO white-dwarf binary seen edge-on. (ONeMg white dwarfs have minimum mass $``$1.1 M\[Wanajo, Hashimoto & Nomoto 1999\] so this possibility cannot be ruled out.) The evolutionary history of the PSR J1141$``$6545 binary, if the companion is a massive white dwarf, is then clear (see e.g. Dewey & Cordes (1987), Portegies Zwart & Yungelson (1999), Tauris & Sennels (2000)): the system originated as a binary consisting of two main sequence stars having mass ratio near unity, but with neither sufficiently massive to independently form a neutron star. For example, the primary may have had mass 7 M and the secondary 5 M. The primary evolved first to form the white dwarf, in the process transferring sufficient matter onto the secondary for its mass to exceed that necessary to form a neutron star. After the white dwarf formed, it spiraled into the envelope of the now more massive secondary as the latter ascended the giant branch, greatly decreasing the orbital period, and ejecting the common envelope. The secondary, a helium star after the ejection of the envelope, then exploded in a supernova, which fortuitously did not disrupt the binary. The radio pulsar we see is therefore the second evolved star. The evolutionary scenario in which a massive white dwarf forms prior to the neutron star has been considered in population synthesis studies. Dewey & Cordes (1987) and Portegies Zwart & Yungelson (1999) showed that the birth rate of white dwarf/young pulsar binaries is comparable to that of double neutron star binaries. However, Portegies Zwart & Yungelson (1999) argued that the white dwarf should have mass $`\stackrel{>}{_{}}1.1`$ M, a result of their assumption that a neutron-star forming binary had to have an initial primary mass of $`>7`$ M. To form a $``$1 M white dwarf, an initial minimum primary mass of $``$6 M is required (S. Portegies Zwart, private communication). The work of Tauris & Sennels (2000) confirms that white dwarf/young neutron star systems should be observable, but they suggest that the birth rate of such systems is much higher than that of double neutron star systems. A lower allowed initial primary mass for forming white dwarfs in the Portegies Zwart & Yungelson (1999) analysis could reduce the disagreement. One prediction of the binary evolution theory is that the space velocity of PSR J1141$``$6545 will be greater than $``$150 km s<sup>-1</sup> (Tauris & Sennels (2000)). At a distance of 3.2 kpc, this implies a proper motion of $`>10`$ mas yr<sup>-1</sup>, which could be measurable by timing on a time scale of a few years. VLBI observations may also be able to detect it. In fact, given the pulsar’s timing age, Galactic latitude, distance estimate, and assuming that the pulsar was born in the Galactic plane, the component of the pulsar’s space velocity perpendicular to the plane must be $`150(d/3.2\mathrm{kpc})`$ km s<sup>-1</sup>, already roughly consistent with the prediction. If the companion is indeed a white dwarf, allowing for interstellar reddening, its B magnitude is likely to be in the range 25 – 26.5 and R in the range 24.5 – 26. From a Digital Sky Survey image of the field, it is clear that the field is relatively crowded, but not impossibly so. Excellent seeing conditions on the VLT would allow detection of the companion. If detected, then a follow-up temperature measurement using HST may provide a useful age constraint. ### 3.2 Classical Contribution to $`\dot{\omega }`$? The observed $`\dot{\omega }`$ may not be purely relativistic, in which case the above conclusions would have to be modified. Classical contributions to $`\dot{\omega }`$ can come either from tidal deformation of the companion by the neutron star (significant only if the companion is non-degenerate), or from a rotation-induced quadrupole moment in the companion, possibly relevant if it is a rapidly rotating white dwarf. We consider these possibilities in turn. A tide raised on the companion star gives it a quadrupole moment that results in a classical apsidal advance, $`\dot{\omega }_{tide}`$, given by $$\dot{\omega }_{tide}=3.44\times 10^6k_2\left(\frac{M_p}{M_c}\right)\left(\frac{M}{M_{}}\right)^{5/3}\left(\frac{R_c}{R_{}}\right)^5{}_{}{}^{}\mathrm{yr}_{}^{1},$$ (3) where $`R_c`$ is the companion radius and $`k_2`$ its apsidal constant, a measure of its internal density distribution (Roberts, Masters, & Arnett (1976); Smarr & Blandford (1976)). A hydrogen main sequence star companion could in principle fit in the orbit, but would result in a classical $`\dot{\omega }`$ $``$1000 times the relativistic value (see Masters & Roberts (1975)) so is certainly ruled out. However, a helium main sequence star has much smaller radius, and would result in a classical $`\dot{\omega }`$ that is comparable to the expected relativistic value. We have explored this possibility by considering a range of values for $`M_p`$ and $`M_c`$, and finding the expected $`\dot{\omega }=\dot{\omega }_{tide}+\dot{\omega }_{GR}`$ that match the observed value. For these calculations, we assumed $`k_2(R_c/R_{})^5=4\times 10^6(M_c/\mathrm{M}_{})^{4.59}`$, as derived by Roberts et al. (1976). In Figure 4, the dashed line shows the allowed locus in $`M_pM_c`$ space for a non- or slowly rotating helium star companion. Clearly this is possible only if $`M_p0.75`$ M. This seems unlikely given observed masses of other neutron stars. A rapidly rotating white dwarf would have a quadrupole moment that could result in a significant classical $`\dot{\omega }`$ (Roberts et al. 1976, Smarr & Blandford (1976)). However unlike the tidal quadrupole, the rotationally-induced quadrupole lies in a plane perpendicular to the white dwarf spin axis. Since the pulsar formed after the white dwarf, it is unlikely that the white dwarf spin axis is aligned with the orbital angular momentum, as this would demand a fortuitously symmetric supernova explosion. A misalignment of the angular momenta results in classical spin-orbit coupling, that is, a precession of the orbital plane, $`di/dt`$ as in the pulsar/B-star binary PSR J0045$``$7319 (Lai, Bildsten, & Kaspi (1995); Kaspi et al. (1996)). This precession would manifest itself as a time-variable projected semi-major axis $`xa\mathrm{sin}i`$. In general, the magnitude of $`di/dt`$ is comparable to that of the periastron precession $`\dot{\omega }`$ (Wex (1998)). For PSR J1141$``$6545, we find an upper limit on $`\dot{x}`$ that is much smaller than $`\dot{\omega }`$ (Table 1). This might suggest that spin-orbit coupling is unlikely to be occurring. However, $`\dot{x}`$ varies as $`\mathrm{cos}i`$; since $`i90^{}`$ (especially if $`\dot{\omega }_{GR}`$ is smaller than our observed value), we cannot rule out orbital plane precession, i.e. a large $`di/dt`$. The above reasoning also holds for a helium-star companion that is rotating rapidly. In this case, the dashed line in Figure 4 is inappropriate as it assumes no contribution from a rotation-induced quadrupole moment. Hence, we cannot presently rule out this possibility. However, a 1-M helium main sequence star would be considerably brighter than a white dwarf: assuming it is on the helium main sequence, it would have bolometric luminosity $`\mathrm{log}(L/L_{})=2.4`$ and $`T_{eff}=50,000`$ K (Keppenhahn & Weigert (1990)). Using bolometric corrections from Bessell, Castelli & Plez (1998) and given the distance and expected reddening, a helium star should have $`B17`$, and should be easily distinguishable from a white dwarf. ### 3.3 Future Relativistic Observables The prospects for a high-precision determination of both component masses in the PSR J1141$``$6545 binary system are excellent. Our upper limit on the combined time dilation and gravitational redshift parameter $`\gamma `$ is shown in Table 1. Note that this is largely independent of which model for the system is correct. Given that the fractional uncertainty in $`\gamma `$ scales as $`T^{3/2}`$, where $`T`$ is the observing span (Damour & Taylor (1992)), we expect a 3$`\sigma `$ detection by 2001, assuming $`\dot{\omega }`$ is purely relativistic (for which $`\gamma 0.7`$ ms for $`M_p=1.32`$ M). Furthermore, a general relativistic orbital period derivative, $`\dot{P_b}`$, should be measurable with interesting precision by 2004, given that its fractional uncertainty scales as $`T^{5/2}`$ and its expected value is $`3.5\times 10^{13}`$. Thus, continued high-precision timing observations of PSR J1141$``$6545 will eventually test our assumptions about a purely relativistic $`\dot{\omega }`$. Although the orbit is likely to be highly inclined, measurement of Shapiro delay will be difficult: the maximum expected delay is only $``$$`\mu `$s, well below our current timing precision. Although PSR J1141$``$6545 will coalesce in $``$1.5 Gyr due to gravitational-radiation-induced orbit decay, it, and coalescing neutron-star/white-dwarf binaries in general, will not be detected by LIGO. The detectability depends on the frequency of the emitted radiation near the time of coalescence. This depends on the size of the orbit at coalescence, which, to order of magnitude, is the white dwarf radius, $`10^9`$ cm. This implies an orbital frequency of $``$0.25 Hz at coalescence; the implied gravitational wave frequency, twice the orbital frequency, is well outside the LIGO band of 10–10,000 Hz. However, such signals would be within the 0.00001–1 Hz band of the proposed Laser Interferometric Space Antenna (LISA; Bender (1998)). General relativistic geodetic precession (Barker & O’Connell (1975)) is predicted to result in the precession of the pulsar’s spin axis at a rate of $`1.7^{}`$yr<sup>-1</sup> (assuming a purely relativistic $`\dot{\omega }`$). This is greater than that predicted for the two pulsars for which it has been observed, PSR B1913+16 (Weisberg, Romani, & Taylor (1989); Kramer (1998)) and PSR B1534+12 (Arzmoumanian (1995); Stairs et al. (2000)), although the measured amplitude will be suppressed by a possibly large factor that depends on the unknown line-of-sight geometry and spin-orbit misalignment. Given the narrow pulse profile seen in PSR J1141$``$6545 (Fig. 1, FWHM of 0.011$`P`$), this effect could be detected in a few years, barring unfortunate geometries. Thus far, we have detected no strong evidence for pulse profile changes, although a detailed analysis is beyond the scope of this paper. That the pulsar was not detected in past pulsar radio surveys of the region is intriguing. The Johnston et al. (1992) 20-cm survey of the Galactic plane had minimum detectable flux density of $``$1 mJy, below the pulsar’s 20-cm flux density of 3.3 mJy. At this frequency, the pulsar displays no evidence for scintillation. The Parkes all-sky survey at 70 cm may also have been able to detect it. That survey’s sensitivity reached 20 mJy for Galactic plane sources; extrapolating the flux densities in Table 1, the estimated flux density at 70 cm is $``$40 mJy. Scattering for the pulsar at that frequency is negligible. Thus, it is plausible that the pulsar’s beam has only recently precessed into our line-of-sight, a result of geodetic precession. If so, the pulse profile should be evolving rapidly. We thank M. Kramer for assistance at Parkes, B. Gaensler for help with the ATCA data, and V. Kalogera, M. van Kerkwijk, S. Portegies Zwart, F. Rasio and S. Thorsett for helpful discussions. VMK is supported by a National Science Foundation CAREER award (AST-9875897). FC is supported by NASA grant NAG 5-3229.
warning/0005/hep-th0005026.html	ar5iv	text	# A Noncommutative M–Theory Five–brane ## I Introduction The M–theory five–brane is a most mysterious object . Several puzzles remain concerning its effective action and the theory of coincident five–branes, see e.g. . Recently, there has been much interest in the noncommutative geometry that arises on D–brane worldvolumes from including a constant Neveu Schwarz two–form background potential . It is natural to ask oneself whether similar deformed geometries play a role in M–theory. It has been suggested that the five–brane of M–theory plays the role of a D–brane in M–theory where, instead of an open string ending on the D–brane, we have an open membrane ending on the M5–brane . We note that the open membrane ending on the five–brane can be reduced to the fundamental string ending on the D$`4`$–brane and by dimensional reduction one may therefore learn about the properties of D–branes by analysing the M5–brane, see for example . The M–theory origin of the Neveu–Schwarz two–form background potential is given by the three–form potential $`C`$ of eleven–dimensional supergravity. In order to investigate the occurrence of deformed geometries in M–theory one is therefore naturally led to consider an M–theory open membrane/five–brane system in the background of a constant $`C`$–field. The M5–brane in the background of a constant $`C`$–field has been investigated previously in the context of the AdS/CFT correspondence and supersymmetric n–cycles . It is known that, in order to derive the noncommutative properties of the D–brane, one may consider the interaction of open string end points with the Neveu–Schwarz two–form potential $`B_{\mathrm{NS}}`$ on the D–brane and perform a canonical quantisation . Analogously, to determine the effect of including the background $`C`$–field on the usual five–brane geometry, we examine how an open membrane couples to an eleven–dimensional supergravity background with background five–branes. A preliminary analysis can be found in . In this paper we will not perform the full canonical quantisation programme to the open membrane. Inspired by the case of D–branes , we will consider a certain decoupling limit in which (1) the bulk modes of the membrane decouple; (2) the dynamics of the boundary string in the five–brane is governed by the Wess–Zumino term of the M2–brane action; and (3) the worldvolume theory of the M–theory five–brane is described by a (2,0) supersymmetric field theory. Following this, we canonically quantise this Wess–Zumino boundary term only. The definition of the decoupling limit is more subtle than in the case of strings because: (i) there is no coupling constant parameter $`g_s`$ in M-theory; (ii) the M5–brane couples to the dual of the $`C`$–field; and (iii) the five–brane equations of motion involve a non–linear algebraic self–duality condition which must be taken into account . The canonical analysis of the boundary term leads to a noncommutative loop space on the five–brane worldvolume. This is a functional analogue of the noncommutative geometry that occurs for D–branes. Following this, we give a suggestion for a loop space star product. The decoupling limit mentioned above involves taking the field theory limit on the five–brane in the presence of a background field strength. We conjecture that the appropriate metric on the five–brane for describing this limit is not the usual induced closed membrane metric but rather an open membrane metric. This metric is analogous to the open string metric on a D–brane and indeed reduces to the latter upon double dimensional reduction. The open brane metrics are the natural metrics for describing the D–brane and the five–brane worldvolume theories in the presence of a background field strength in that it is these metrics rather than the induced ones that give the linearised mass shell condition on the worldvolume which in turn defines the field theory limits on the branes. The structure of the paper is as follows. In Section II we introduce the open membrane/five–brane system. In Section III we discuss some properties of open brane metrics. Next, in Section IV, we give the decoupling limit and discuss its relation to the decoupling limit of . In Section V we derive the deformed five–brane geometry. First we apply the Dirac canonical quantisation programme to the decoupled boundary Wess–Zumino term and derive the Dirac brackets of the five–brane coordinates. We then connect to the results in the literature for the D4–brane, examine the special case of the infinite momentum frame and finally give a star product for particular cases. Finally, in Section VI we give a discussion of our results and possible extensions to this work. ## II The open membrane/five–brane system Our starting point is an open membrane in an eleven–dimensional supergravity background with background five–branes. The membrane boundary is a string that is constrained to lie within one of the five–branes which is separated from the stack of remaining five–branes. The reason that we include the additional stack of five–branes in the background is that it effectively stiffens the separated five–brane so that the membrane is not able to deform this five–brane and can be treated as a test membrane, as we shall discuss in more detail in Section IV. The action for the open bosonic membrane is as follows (for the kappa-symmetric version of this action, see ): $$S=S_k+_{M^3}f_2^{}C+_{M^3}f_1^{}b,$$ (1) where the kinetic term can be written in Polyakov form as $$S_k=\frac{1}{2(\mathrm{}_p)^2}_{M^3}d^3\sigma \sqrt{det\gamma }\left(\gamma ^{\alpha \beta }_\alpha X^M_\beta X^N\widehat{g}_{MN}+\mathrm{}_p^2\right).$$ (2) Here $`\mathrm{}_p`$ is the $`D=11`$ Planck’s constant, $`\widehat{g}_{MN}`$ is the $`D=11`$ spacetime metric and $`\gamma _{\alpha \beta }`$ is the auxiliary worldvolume metric. The maps $`f_2`$ and $`f_1`$ denote the embedding of the membrane and its boundary into the spacetime and the five–brane, respectively. The worldvolume three–form $`f_2^{}C`$ is the pull–back of the $`D=11`$ three–form potential $`C`$ to the membrane worldvolume and, similarly, $`f_1^{}b`$ is the pull–back of the five–brane two–form potential $`b`$ to the boundary of the membrane. In terms of components, we write $$(f_2^{}C)_{\alpha \beta \gamma }=_\alpha X^M_\beta X^N_\gamma X^PC_{MNP},(f_1^{}b)_{ij}=_iX^\mu _jX^\nu b_{\mu \nu },$$ (3) where $`M=0,1,\mathrm{},9,11`$ are spacetime indices, $`\mu =0,1\mathrm{},5`$ are five–brane worldvolume indices, $`\alpha =0,1,2`$ are membrane worldvolume indices and $`i=0,1`$ are indices on the boundary of the membrane. Since the decoupling limit involves scaling Planck’s constant $`\mathrm{}_p`$ we need to carefully assign dimensions to all quantities. We are using units in which the spacetime metric $`\widehat{g}_{MN}`$, the worldvolume metric $`\gamma _{\alpha \beta }`$, all differential forms and the membrane worldvolume parameters $`\sigma ^\alpha `$ are dimensionless. The spacetime coordinates $`X^M`$ and the five–brane coordinates $`X^\mu `$ have dimension length. Note that the components $`C_{MNP}`$ and $`b_{\mu \nu }`$ have dimension (mass)<sup>3</sup> and (mass)<sup>2</sup>, respectively, in order for $`C`$ and $`b`$ to be dimensionless. The coupling of $`b`$ to the boundary of the membrane ensures that the open membrane action is invariant under the spacetime gauge transformations $`\delta C=d\mathrm{\Lambda }`$ provided that $`\delta b=f_5^{}\mathrm{\Lambda }`$, where $`f_5`$ denotes the embedding of the five–brane into spacetime. The two–form $`b`$ satisfies the five–brane field equations. These are equivalent to a non–linear self–duality condition on the following gauge invariant three–form field strength of $`b`$: $$=db+f_5^{}C.$$ (4) Here the last term is the pull–back of the spacetime three–form potential to the five–brane: $$(f_5^{}C)_{\mu \nu \rho }=_\mu x^M_\nu x^N_\rho x^PC_{MNP},$$ (5) where $`x^M(X^\mu )`$ are local embedding functions satisfying the five–brane equations of motion. The non–linear self–duality condition on the five–brane reads $$\frac{\sqrt{detg}}{6}ϵ_{\mu \nu \rho \sigma \lambda \tau }^{\sigma \lambda \tau }=\frac{1+K}{2}(G^1)_\mu {}_{}{}^{\lambda }_{\nu \rho \lambda }^{},$$ (6) where $`ϵ^{012345}=1`$ and the scalar $`K`$ and the tensor $`G_{\mu \nu }`$ are given by $`K`$ $`=`$ $`\sqrt{1+{\displaystyle \frac{\mathrm{}_p^6}{24}}^2},`$ (7) $`G_{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{1+K}{2K}}\left(g_{\mu \nu }+{\displaystyle \frac{\mathrm{}_p^6}{4}}_{\mu \nu }^2\right).`$ (9) We shall argue in the next section that the tensor $`G_{\mu \nu }`$ is the metric on the five–brane seen by an open membrane in the presence of a background three–form field strength $``$. It is also understood that in the above three equations the indices are contracted with the induced five–brane metric: $$g_{\mu \nu }=_\mu x^M_\nu x^N\widehat{g}_{MN}.$$ (10) We may parameterise the $`D=11`$ spacetime by local coordinates $`X^M=(X^\mu ,Y^m)`$, where $`m=6,7,8,9,11`$ refer to the directions perpendicular to the five–brane, such that the Dirichlet condition on the membrane embedding fields becomes $`Y^m(\sigma ^\alpha )\|_{M^3}=0`$. The remaining parallel embedding fields $`X^\mu (\sigma ^\alpha )`$ obey a mixed Dirichlet and Neumann boundary condition. By varying the action (1), using $`\delta Y^m\|_{M^3}=0`$ together with the definitions (4) and (10) one finds that $$\left[\frac{1}{\mathrm{}_p^2}\sqrt{det\gamma }n^\alpha _\alpha X^M\widehat{g}_{\mu M}+ϵ^{\alpha \beta \gamma }n_\alpha _{\mu \nu \rho }_\beta X^\nu _\gamma X^\rho \right]\|_{M^3}=0,$$ (11) where $`n_\alpha `$ is the normal vector at the boundary. The first term is non–local from the point of view of the boundary string. As we shall see, this term drops out in the decoupling limit leaving a local field equation for the embedding fields $`X^\mu (\sigma ^i)`$ of a closed string in six dimensions. The closed string boundary conditions are $$X^\mu (\tau ,\sigma +2\pi )=X^\mu (\tau ,\sigma ),$$ (12) where $`\tau \sigma ^0`$ is the time coordinate and $`\sigma \sigma ^1`$ is the spatial coordinate on the worldsheet. We shall consider backgrounds where $`_{\mu \nu \rho }`$ is constant. This is only consistent with (4) provided we require that the pull–back of the spacetime four–form field strength $`F=dC`$ to the five–brane vanishes, i.e. $`f_5^{}F=0`$. After some manipulations (given in appendix A) we end up with the following action: $$S=S_k+_{M^3}f_2^{}\stackrel{~}{C}+\frac{1}{3}_{M^3}d^2\sigma _{\mu \nu \rho }X^\mu \dot{X}^\nu X^\rho ,$$ (13) where $`\stackrel{~}{C}`$ is the part of the three–form potential which is perpendicular to the five–brane, i.e. $`f_5^{}\stackrel{~}{C}=0`$, and the dot and prime indicates differentiation with respect to $`\tau `$ and $`\sigma `$, respectively. It will be useful to introduce a specific parametrization of the solutions of the self–duality condition (6) as followsThis parameterisation has been derived independently by . (see Appendix B): $$_{\mu \nu \rho }=\frac{h}{\sqrt{1+\mathrm{}_p^6h^2}}ϵ_{\alpha \beta \gamma }v_\mu ^\alpha v_\nu ^\beta v_\rho ^\gamma +hϵ_{abc}u_\mu ^au_\nu ^bu_\rho ^c,$$ (14) $$G_{\mu \nu }=\frac{\left(1+\sqrt{1+h^2\mathrm{}_p^6}\right)^2}{4}\left(\frac{1}{1+h^2\mathrm{}_p^6}\eta _{\alpha \beta }v_\mu ^\alpha v_\nu ^\beta +\delta _{ab}u_\mu ^au_\nu ^b\right).$$ (15) Here $`h`$ is a real field of dimension (mass)<sup>3</sup> and $`(v_\mu ^\alpha ,u_\mu ^a)`$, $`\alpha =0,1,2`$, $`a=3,4,5`$, are sechsbein fields in the nine–dimensional coset $`SO(5,1)/SO(2,1)\times SO(3)`$ satisfying $$g^{\mu \nu }v_\mu ^\alpha v_\nu ^\beta =\eta ^{\alpha \beta },g^{\mu \nu }u_\mu ^av_\nu ^\beta =0,g^{\mu \nu }u_\mu ^au_\nu ^b=\delta ^{ab},$$ (16) $$g_{\mu \nu }=\eta _{\alpha \beta }v_\mu ^\alpha v_\nu ^\beta +\delta _{ab}u_\mu ^au_\nu ^b.$$ (17) ## III Open Brane Metrics The natural metric to examine the worldvolume of a D–brane with a non–trivial background Born–Infeld field strength $`_{\mu \nu }`$ is the so called open string metric. For a D$`p`$–brane this metric is defined as follows: $$G_{\mu \nu }=g_{\mu \nu }\alpha ^2_{\mu \rho }g^{\rho \lambda }_{\lambda \nu },\mu =0,1,\mathrm{},p,$$ (18) where $`g_{\mu \nu }`$ is the induced metric on the brane. The open string metric defines the mass shell condition for the open string modes propagating in the D–brane. The expression (18) for the open string metric can also be read off from the low energy effective Dirac–Born–Infeld action which yields massless field equations with D’Alembertians $`(G^1)^{\mu \nu }_\mu _\nu `$ and Dirac operators $`(G^1)^{\mu \nu }\mathrm{\Gamma }_\mu _\nu `$ (where $`_\mu `$ contains the Christoffel symbol of $`g_{\mu \nu }`$). The decoupling limit on a D$`p`$–brane as defined in is non–degenerate in the sense that it yields a $`(p+1)`$–dimensional, massless field theory on the brane if the leading components of the open string metric form a rank $`p+1`$ matrix. In analogy with how branes appear as solutions in supergravity, the open string is visible as a supersymmetric worldvolume soliton solution to the massless D–brane field equations. It is instructive to examine the geometry of these so called BIon solutions using the open string metric. A charged, static BIon, corresponding to the string ending on a D$`p`$–brane, excites one transverse scalar X and the time component of the Born-Infeld vector field $`A_\mu `$. The solution is given by $$(\alpha ^{})^{\frac{1}{2}}A_0=(\alpha ^{})^{\frac{1}{2}}X=H,$$ (19) where $`H`$ is a harmonic function on the transverse worldvolume space $`\text{𝔼}^p`$ given by: $$H=1+\frac{Q}{r^{p2}},p>2.$$ (20) After substituting this solution into the open string metric (18) we find the following two–block 0–brane line element: $$ds^2(G)=\frac{dt^2}{1+\alpha ^{}(H)^2}+(dy^m)^2,m=1,\mathrm{}p,$$ (21) where $`(H)^2=\delta ^{mn}_mH_nH`$ and $`(dy^m)^2`$ is the line element on $`\text{𝔼}^p`$. Taking the limit of large $`Q`$ and carrying out the coordinate transformation $`r^{2(p2)}=\alpha ^{}Q^2/u^2`$ we find the following metric: $$ds^2(G)=\frac{r^2}{(p2)^2}\left(\frac{dt^2+du^2}{u^2}+(p2)^2d\mathrm{\Omega }_{p1}^2\right).$$ (22) This metric is conformally $`AdS_2\times S^{p1}`$. The ordinary induced metric $`ds^2(g)`$ is of three–block form with an additional line element proportional to $`dr^2`$ and therefore not conformally $`AdS_2\times S^{p1}`$ in the near horizon region. The natural properties of the open string metric (18) in examining the open string spectrum, motivates the introduction of an analogous metric for studying an open membrane probing a five–brane in the presence of a background $``$ field. We propose that the conformal class of this so called open membrane metric is given by the tensor $`G_{\mu \nu }`$ given in (9). This metric has been shown to dimensionally reduce to the open string metric on the D4–brane up to a scale factor . Moreover, the massless field equations on the five–brane have kinetic terms involving the box operator $`(G^1)^{\mu \nu }_\mu _\nu `$ and the Dirac operator $`(G^1)^{\mu \nu }\mathrm{\Gamma }_\mu _\nu `$ . Further arguments in favour of the proposed open membrane metric are obtained by considering other excitations than the massless tensor multiplet, like e.g. the self–dual string soliton corresponding to the open membrane ending on the five–braneOnly in this section we use the index $`m`$ to indicate the five–brane worldvolume directions transverse to the self–dual string. In all other sections the index $`m`$ indicates the target space directions transverse to the five–brane.: $$4\mathrm{}_p^2b_{01}=\mathrm{}_p^1X=H,_{mnp}=\frac{\mathrm{}_p^2}{4}ϵ_{mnpq}_qH,m=2,3,4,5,$$ (23) where $`H`$ is a harmonic function on the transverse worldvolume space $`\text{𝔼}^4`$ given by: $$H=1+\frac{Q}{r^2}.$$ (24) Substituting this solution into (9) the line element of the open membrane metric has the structure of a two–block 1–brane solution: $$ds^2(G)=\frac{1}{4}\left(1+\sqrt{1+\mathrm{}_p^2(H)^2}\right)^2\left(\frac{dt^2+dx^2}{1+\mathrm{}_p^2(H)^2}+(dy^m)^2\right),$$ (25) whereas the induced metric has an additional line element proportional to $`dr^2`$. In the large $`Q`$ limit the open membrane metric, upon making the coordinate transformation $`r^4=\mathrm{}_p^2Q^2/u^2`$, becomes: $$ds^2(G)=\frac{u^2}{4}\left(\frac{dt^2+d\sigma ^2+du^2}{u^2}+4d\mathrm{\Omega }_3^2\right),$$ (26) which we identify as $`AdS_3\times S^3`$ with a conformal factor. ## IV The decoupling limit We will now consider a limit of the open membrane/five–brane system with the following three properties. Firstly, the bulk modes of the membrane can be neglected. Secondly the boundary string that lives in the five–brane is governed solely by the Wess–Zumino term. Thirdly, the field theory limit on the five–brane is taken such that the physics on the five–brane is described by the self–dual tensor field theory. In order to demonstrate the limiting procedure we require, it is instructive to first examine the the limit taken by Seiberg and Witten in deriving the noncommutative geometry of a D–brane from the open string. ### The open string case We study the decoupling limit at small string coupling $`g_s`$. The effective tensions $`\tau `$ of the string and the D$`p`$–brane behave like $`\tau _{\mathrm{F1}}1,\tau _{\mathrm{Dp}}1/g_s`$. Therefore, for small $`g_s`$, the string is much lighter than the D$`p`$–brane and can be treated as a test string probing the D$`p`$–brane. Furthermore, the effective gravitational couplings $`G_N\tau `$ (Newton’s constant times tension) behave like $`G_N\tau _{\mathrm{F1}}g_s^2,G_N\tau _{\mathrm{Dp}}g_s`$ and therefore we can assume that the spacetime background is approximately flat. The open string action reads $$S=\frac{1}{2\alpha ^{}}_{M^2}d^2\sigma \overline{}X^MX^N\eta _{MN}+\frac{1}{2}_{M^2}𝑑\tau _{\mu \nu }X^\mu \dot{X}^\nu ,$$ (27) where $`_{\mu \nu }`$ is the background field strength on the D$`p`$–brane. We assume that the only non–vanishing components of $``$ are $`_{r^{}s^{}}`$, where we have decomposed the worldvolume index $`\mu `$ as $`\mu =(r,r^{})`$ with $`r=0,1,\mathrm{},p\mathrm{rank}`$ and $`r^{}=p+1\mathrm{rank},\mathrm{},p`$. The string theory has closed string modes in the bulk, i.e. away from the D–brane. Their mass shell condition is governed by the closed string metric $`\eta _{MN}`$: $$\eta ^{MN}p_Mp_N=\frac{4}{\alpha ^{}}\left(N_\mathrm{L}1\right)=\frac{4}{\alpha ^{}}\left(N_\mathrm{R}1\right).$$ (28) Here $`N_{\mathrm{L},\mathrm{R}}`$ indicates the oscillator number of the left and right movers and $`p_M`$ is the bulk momentum of the closed string state. There are also open string modes propagating in the D–brane. An open string state with worldvolume momentum $`k_\mu `$ obeys the mass shell condition $$(G^1)^{\mu \nu }k_\mu k_\nu =\frac{1}{\alpha ^{}}\left(N^{\mathrm{open}}1\right),$$ (29) where $`N^{\mathrm{open}}`$ is the open string oscillator number and the tensor $`G_{\mu \nu }`$ is the open string metric defined in (18). From the conditions on $`_{\mu \nu }`$, it follows that $`G_{\mu \nu }`$ is given by $$G_{\mu \nu }=\{\begin{array}{cc}\eta _{\mu \nu }\hfill & \text{for}\mu ,\nu =r,s\hfill \\ \eta _{\mu \nu }\alpha ^2_{\mu \rho }\eta ^{\rho \sigma }_{\sigma \nu }\hfill & \text{for}\mu ,\nu =r^{},s^{}\hfill \end{array}$$ (30) In order to describe the limit we first decompose the target spacetime index $`M`$ as $`M=(r,r^{},m)`$ and split the spacetime metric $`\eta _{MN}`$ into parts that are parallel and perpendicular to $``$ and the D–brane as follows: $`\eta _{MN}\eta _{rs}\eta _{r^{}s^{}}\eta _{mn}`$, where $`\eta _{rs}`$ is perpendicular to $``$ and parallel to the D–brane, $`\eta _{r^{}s^{}}`$ is parallel to $``$, and $`\eta _{mn}`$ is perpendicular to the D–brane. The limit of Seiberg and Witten is obtained by taking $`ϵ0`$ such that<sup>\</sup><sup>\</sup>\In the equation below it is understood that the $`\eta _{r^{}s^{}}`$ and $`\alpha ^{}`$ occurring at the right–hand–side are $`ϵ`$–independent.: $$\eta _{r^{}s^{}}ϵ\eta _{r^{}s^{}},\alpha ^{}ϵ^{\frac{1}{2}}\alpha ^{},$$ (31) while keeping all other quantities fixed. The open string action scales as follows: $`S`$ $``$ $`{\displaystyle \frac{1}{2ϵ^{1/2}\alpha ^{}}}{\displaystyle _{M^2}}d^2\sigma \overline{}X^mX^n\eta _{mn}+{\displaystyle \frac{1}{2ϵ^{1/2}\alpha ^{}}}{\displaystyle _{M^2}}d^2\sigma \overline{}X^rX^s\eta _{rs}`$ (34) $`+{\displaystyle \frac{ϵ^{1/2}}{2\alpha ^{}}}{\displaystyle _\mathrm{\Sigma }}d^2\sigma \overline{}X^r^{}X^s^{}\eta _{r^{}s^{}}+{\displaystyle \frac{1}{2}}{\displaystyle _{M^2}}𝑑\tau _{r^{}s^{}}X^r^{}\dot{X}^s^{}`$ and the various string mass shell conditions become $`\mathrm{bulk}:`$ $`\{\begin{array}{c}\eta ^{mn}p_mp_n,\eta ^{rs}p_rp_s\frac{4}{ϵ^{1/2}\alpha ^{}}(N1),N=N_\mathrm{L}=N_\mathrm{R},\hfill \\ \eta ^{r^{}s^{}}p_r^{}p_s^{}\frac{4ϵ^{1/2}}{\alpha ^{}}(N1),\hfill \end{array}`$ (37) $`\mathrm{D}\mathrm{brane}:`$ $`(G^1)^{\mu \nu }k_\mu k_\nu {\displaystyle \frac{1}{ϵ^{1/2}\alpha ^{}}}(N^{\mathrm{open}}1),`$ (38) where the open string metric is finite in this limit and is given by the maximal rank matrix $$G_{\mu \nu }=\{\begin{array}{cc}\eta _{\mu \nu }\hfill & \text{for}\mu ,\nu =r,s\hfill \\ \alpha ^2_{\mu \rho }\eta ^{\rho \sigma }_{\sigma \nu }\hfill & \text{for}\mu ,\nu =r^{},s^{}\hfill \end{array}$$ (39) The limit (31) has the following three crucial properties that we wish to emulate for the open membrane/five–brane system: 1. The closed string bulk modes perpendicular to $``$ with momentum $`p_M=(p_r,0,0)`$ or $`p_M=(0,0,p_m)`$ are frozen out<sup>††</sup><sup>††</sup>†† The massless, perpendicular bulk modes remain but their effective bulk coupling constant goes to zero. thereby isolating the dynamics of the D–brane theory from the bulk. 2. The closed string bulk modes with momentum $`p_M=(0,p_r^{},0)`$ give vanishing contribution to the action. 3. On the D–brane all massive open string modes are frozen out and the decoupled, massless field theory on the brane is non–degenerate in the sense that it is $`(p+1)`$–dimensional and has finite effective coupling. The noncommutative nature of the D–brane arises from quantising the remaining Wess–Zumino term. ### The open membrane case We now would prefer to proceed by analogy to the open membrane/five–brane system. However, the analysis of the decoupling limit for this system requires a slight modification, due to three circumstances: the absence of an analog of the string coupling constant; the fact that the open membrane probe couples to the three–form potential which is dual to the potential of the background five–brane; and the non–linear self–duality condition on the field strength in the five–brane worldvolume. Thus, in M-theory, there is no sense in which the tension of an isolated five–brane can be said to be much larger than the tension of an isolated membrane. Neither can we assume an approximately flat spacetime background around a background five–brane. To prevent the membrane from deforming the background geometry, we therefore consider a $`D=11`$ background consisting of a large stack of parallel five–branes, given by the solution $$ds^2(\widehat{g})=H^{1/3}(dx^\mu )^2+H^{2/3}(dy^m)^2,H=1+\frac{N_5\mathrm{}_p^3}{r^3},F=N_5ϵ_4,$$ (40) where $`\mu =0,1,\mathrm{},5;m=6,7,8,9,11,N_5`$ is the number of stacked five–branes and $`ϵ_4`$ is the volume form on the transverse $`S^4`$. We let the open membrane end on one of these five–branes removed from the stack and placed at radius $`r_0`$. If $`N_5>>1`$ and $`r_0`$ is small, then the interactions between the stack and the separated five–brane effectively stiffens the latter so that the membrane can probe it without deforming it<sup>‡‡</sup><sup>‡‡</sup>‡‡ By analysing the Nambu–Goto action for the five–brane scalars, it can be shown that for large $`N_5`$ the characteristic length scale of five–brane deformations with large gradients is $`N_5^{\frac{1}{3}}\mathrm{}_p`$.. Under these conditions the induced metric on the five–brane (10) is given by: $$g_{\mu \nu }=H^{1/3}(r_0)\eta _{\mu \nu }.$$ (41) Moreover, from (40) it follows that the $`D=11`$ background four–form field strength satisfies $`f_5^{}F=0`$. From the discussion in Section II it follows that we may consider an open membrane action given by $`S`$ $`=`$ $`{\displaystyle \frac{1}{2(\mathrm{}_p)^2}}{\displaystyle _{M^3}}d^3\sigma \sqrt{det\gamma }\left(H^{1/3}\gamma ^{\alpha \beta }_\alpha X^\mu _\beta X^\nu \eta _{\mu \nu }H^{2/3}\gamma ^{\alpha \beta }_\alpha Y^m_\beta Y^n\delta _{mn}+\mathrm{}_p^2\right)`$ (43) $`+N_5{\displaystyle _{M^3}}f_2^{}\stackrel{~}{C}+{\displaystyle \frac{1}{3}}{\displaystyle _{M^3}}d^2\sigma _{\mu \nu \rho }X^\mu \dot{X}^\nu X^\rho ,`$ where the $`D=11`$ background three–form potential $`\stackrel{~}{C}`$ obeys $`d\stackrel{~}{C}=ϵ_4`$ and $`f_5^{}\stackrel{~}{C}=0`$ and the background three–form field strength $`_{\mu \nu \rho }`$ on the five–brane is constant. Unlike the case of the string ending on a D–brane, where the Neveu–Schwarz two–form potential can be tuned independently of the Ramond–Ramond potentials, the membrane couples to the dual potential $`N_5\stackrel{~}{C}`$ of the five–brane six–form potential, which becomes large in the limit of large $`N_5`$. However, $`\stackrel{~}{C}`$ affects only the perpendicular closed membrane bulk modes and does not prevent these modes from decoupling in the limit from the dynamics on the five–brane. We will therefore, from now on, drop the term in the open membrane action (43) containing $`\stackrel{~}{C}`$. Another difference from the string case is that we cannot with impunity take a scaling limit for $``$, since the field equations for the five–brane impose the non–linear algebraic constraints (6). In order for our limit to be guaranteed to obey (6) we will use the explicit solution for $``$ given by (14) and (17). We propose the decoupling limit obtained by taking $`ϵ0`$ such that $`\mathrm{}_p`$ $``$ $`ϵ\mathrm{}_p,`$ (44) $`{\displaystyle \frac{N_5}{r_0^3}}`$ $``$ $`ϵ^{3\delta }{\displaystyle \frac{N_5}{r_0^3}},`$ (45) $`h`$ $``$ $`ϵ^\lambda h.`$ (46) For simplicity we shall assume that $`\delta >1`$<sup>\</sup><sup>\</sup>\ One can also consider the case $`0<\delta 1`$. This does not lead to any qualitative changes in the final result. such that we may drop the $`1`$ from the harmonic function in the metric (40) <sup>\†</sup><sup>\†</sup>\†We point out that the limit considered here differs from the Maldacena limit since energies are not kept fixed in the near horizon region. . It then follows from (41) that the induced five–brane metric and the sechsbein fields in (17) scales as $`g_{\mu \nu }ϵ^{\delta 1}g_{\mu \nu },u_\mu ^aϵ^{\frac{1}{2}(\delta 1)}u_\mu ^a,v_\mu ^\alpha ϵ^{\frac{1}{2}(\delta 1)}v_\mu ^\alpha .`$ (47) Furthermore, we assume that $`\lambda 3`$. This implies that $`h\mathrm{}_p^3`$ remains finite which enables us to keep the three–form field strength and the open membrane metric non–degenerate in the limit. Thus we find that the open membrane action (43) scales as $`S`$ $``$ $`ϵ^\mathrm{\Delta }[{\displaystyle \frac{1}{2\mathrm{}_p^2}}{\displaystyle _{M^3}}d^3\sigma \sqrt{\gamma }(ϵ^{\mathrm{\Delta }+\delta 3}H^{\frac{1}{3}}\gamma ^{\alpha \beta }_\alpha X^\mu _\beta X^\nu \eta _{\mu \nu }ϵ^{\mathrm{\Delta }2\delta }H^{\frac{2}{3}}\gamma ^{\alpha \beta }_\alpha Y^m_\beta Y^n\delta _{mn}+ϵ^\mathrm{\Delta }\mathrm{}_p^2)`$ (49) $`+{\displaystyle \frac{1}{3}}{\displaystyle _{M^3}}d^2\sigma _{\mu \nu \rho }X^\mu \dot{X}^\nu X^\rho ],`$ where we have defined $$\mathrm{\Delta }=\lambda \frac{3}{2}(\delta 1).$$ (50) We now make our requirements, as for the string case: 1. The perpendicular bulk modes $`Y^m`$ are frozen out provided $`\mathrm{\Delta }2\delta <0`$. 2. The action for the parallel bulk modes vanishes if $`\mathrm{\Delta }+\delta 3>0`$. (This amounts to the vanishing of the first term in (11) such that (11) turns into a local field equation on the string worldsheet.) 3. We require a non–degenerate field theory limit on the five–brane. As discussed in Section III, we conjecture that the relevant metric is the open membrane metric (9). Thus we require the distances $`ds^2(G)`$ measured in the open membrane metric to scale more slowly than $`\mathrm{}_p^2`$ and the leading part of the open membrane metric to be non–degenerate in the limit<sup>\‡</sup><sup>\‡</sup>\‡This differs from the limit proposed in where the rank of the open membrane metric is reduced from six to five in the limit.: $$G_{\mu \nu }ϵ^\beta \left(\overline{G}_{\mu \nu }+𝒪(ϵ)\right),$$ (51) where $$\beta <2,det\overline{G}_{\mu \nu }0.$$ (52) From (15) it follows that the second condition in (52) requires $`h\mathrm{}_p^3`$ to remain finite in the limit which implies that $`\lambda 3`$. The first condition in (52) amounts to $`\delta <3`$. Combining our assumption that $`\delta >1`$ with the conditions obtained under (i)-(iii) we find the following restrictions on our parameters: $$\mathrm{\Delta }+\frac{3}{2}(\delta 1)3<\mathrm{\Delta }+\delta ,\mathrm{\Delta }<2\delta ,1<\delta <3.$$ (53) These conditions are solved by $`(\mathrm{\Delta },\delta )`$ in a finite size region. For instance, $`\delta =\frac{5}{3}`$, $`\lambda =3`$ and $`\mathrm{\Delta }=2`$ leads to a decoupled five–brane theory in a background with a non–linearly self–dual field strength, while $`\delta =\frac{4}{3}`$, $`\lambda =2`$ and $`\mathrm{\Delta }=2`$ yields a linearly self–dual field strength. A noteworthy feature is that (53) implies $`\mathrm{\Delta }>0`$, such that there is necessarily an overall scaling of the action in (49). Such a scaling was not required in the string case. A crucial difference between the string and the membrane is that only the string action (27) has a microscopic interpretation. On the other hand, the membrane action should be seen as an effective action. One interpretation of the scaling (49) with $`\mathrm{\Delta }>0`$ is that actually we are taking a semiclassical limit. Summarizing, in order to understand the geometry of the five–brane worldvolume we are led to study the quantisation of the following action: $$S=\frac{1}{3}_{M^3}d^2\sigma _{\mu \nu \rho }X^\mu \dot{X}^\nu X^\rho ,$$ (54) with $``$ constant. ### A comparison with the limits for a lifted D–brane Before proceeding with the canonical analysis in the next section, we wish to compare the decoupling limit we propose for the M-theory five–brane with the one discussed in . By applying the usual relations between string theory and M-theory $$g_s=(R/\mathrm{}_p)^{3/2},\alpha ^{}=\frac{\mathrm{}_p^3}{R},_{\mu \nu }=R_{\mu \nu 5},\mu ,\nu =0,1,2,3,4,$$ (55) one may lift the limits taken for the D4–brane to the case of the M5–brane. We assume that the rank of $``$ is $`4`$. This motivates the following limit for the M–theory five–brane on $`𝐑^6`$ : $`g_{IJ}`$ $``$ $`ϵg_{IJ},g_{55}ϵ^2g_{55},g_{00}1,\mathrm{}_pϵ^{1/2}\mathrm{}_p,I=1,2,3,4,`$ (56) $`_{012}`$ $``$ $`ϵ^1\sqrt{1+ϵ}h_0,_{034}=ϵ^1h,_{125}h_0,_{345}\sqrt{1+ϵ}h,`$ (57) with $`h_0^2h^2+\mathrm{}_p^6h_0^2h^2=0`$ and $`h_0<0`$. The presented constant three–form solution is a special case of (14) where a diagonalised $`_{012}_{345}`$ split background has been boosted in an $`ϵ`$ dependent way in the $`05`$ direction which turns on $`_{034}`$ and $`_{125}`$ components. This is required in order for the reduction along the fifth direction to give a rank $`4`$ constant two–form solution on the D4–brane. Equivalently, we can reduce along a skew spacelike direction in the $`05`$ plane instead of reducing along the fifth direction to obtain a rank $`4`$ solution on the D4–brane. By construction, the dimensional reduction of this solution and limit will give the appropriate limit for the D4–brane. However, under this limit the open membrane metric (15) behaves like $`G_{55}\mathrm{}_p^20`$. This means that our condition (iii) for the field theory limit (involving the open membrane metric) on the five–brane to be valid is not satisfied. Obviously, this is not a problem for the compactified theory where the excitations in the fifth direction are suppressed. ## V Canonical Analysis In this section we will canonically quantise the action (54) with constant field strength $`_{\mu \nu \rho }`$. In the first part of this section we will assume that the field strength can be diagonalised as follows (see Appendix B for details) : $$_{012}=\frac{h}{\sqrt{1+\mathrm{}_p^6h^2}},_{345}=h.$$ (58) where the dimensionless tensor multiplet “coupling” $`h\mathrm{}_p^3`$ is non–vanishing provided the decoupling limit (46) has been taken with $`\lambda =3`$. For $`\lambda <3`$ the limit results in a linear tensor multiplet. At the end of this section we discuss the case in which the field strength cannot be brought into the above form (see Appendix C for details). In the parameterisation (58) the action (54) splits into two independent Lagrangians for the two sets of coordinates $`X^{0,1,2}`$ and $`X^{3,4,5}`$: $$S=\frac{h}{3\sqrt{1+\mathrm{}_p^6h^2}}_{M^3}d^2\sigma ϵ_{\alpha \beta \gamma }X^\alpha \dot{X}^\beta X^\gamma +\frac{h}{3}_{M^3}d^2\sigma ϵ_{abc}X^a\dot{X}^bX^c,$$ (59) where $`\alpha =0,1,2`$ and $`a=3,4,5`$. The action is invariant under worldsheet reparameterisations: $$\delta _\xi X^\alpha =\xi ^i_iX^\alpha ,\delta _\eta X^a=\eta ^i_iX^a,i=0,1.$$ (60) Note that, due to the absence of a worldsheet metric, there is no need to identify the vector fields $`\xi `$ and $`\eta `$. The equations of motion are: $$ϵ_{\alpha \beta \gamma }\dot{X}^\beta X^\gamma =0,ϵ_{abc}\dot{X}^bX^c=0.$$ (61) This means that the embedding of the worldsheet in the 0,1,2 directions is a one–dimensional submanifold of that space. And similarly the embedding in the 3,4,5 directions is also one–dimensional. A special feature of our system is that both the equations of motion and the gauge transformations are first order in derivatives. Thus, before we start the canonical analysis, it is instructive to first analyse the solutions to the equations of motion. We assume that the one–dimensional embedding of the worldsheet in the 0,1,2 directions is timelike. This means that we can fix the following static gauge for the $`\xi `$ symmetry: $$X^0=\tau .$$ (62) The field equation then implies that $$X^{\mathrm{\hspace{0.17em}1},2}=0.$$ (63) There are two inequivalent sectors of the theory which differ by how the one–dimensional embedding in the remaining 3,4,5 directions takes place. 1. The string sector, for which the gauge fixed $`X^a`$ field equation is $$\dot{X}^a=0.$$ (64) 2. The particle sector, for which the gauge fixed $`X^a`$ field equation is $$X^a=0.$$ (65) Hence, in the case of (i) the embedded worldsheet is a two–dimensional surface which is unconstrained in the $`X^{1,2}`$ directions and fixed in the $`X^{3,4,5}`$ directions. In the case of (ii) the worldsheet is contracted into a one–dimensional worldline which is unconstrained in all five spatial directions. (By unconstrained we mean that the coordinate is pure gauge; this will become apparent after the analysis described below.) In the case of strings ending on D–branes the analogous field equations are much simpler, namely: $`\dot{X}^i=0,i=1,\mathrm{},p`$. This leads to the notion of a zero–brane in the D–brane. Here however we find a membrane in the five–brane due to the additional Dirichlet conditions (64). The boundary string is positioned in the space transverse to this membrane. Let us continue by analysing the phase space dynamics of the three Euclidean coordinates $`\stackrel{}{X}=(X^3,X^4,X^5)`$. The canonical momenta are given by $$\mathrm{\Pi }_a(\sigma ):=\frac{\delta S}{\delta \dot{X}^a(\sigma )}=\frac{h}{3}ϵ_{abc}X^bX^{}{}_{}{}^{c}.$$ (66) The non–trivial canonical Poisson brackets are: $$\{X^a(\sigma ),\mathrm{\Pi }_b(\sigma ^{})\}:=\delta _b^a\delta (\sigma \sigma ^{}).$$ (67) Since the Lagrangian is first order in time derivatives there is one primary constraint $`\varphi _a(\sigma )`$ for each canonical momentum, given by $$\varphi _a:=\mathrm{\Pi }_a+\frac{h}{3}ϵ_{abc}X^bX^{}{}_{}{}^{c}0.$$ (68) The canonical Hamiltonian $`H:=\stackrel{}{\mathrm{\Pi }}\dot{\stackrel{}{X}}L`$ vanishes. Instead, one introduces a generalised Hamiltonian $`H_{\mathrm{gen}}:=𝑑\sigma \lambda ^a(\sigma )\varphi _a(\sigma )`$ where $`\lambda ^a(\sigma )`$ are three Lagrange multipliers. To proceed with the canonical analysis we study the consistency conditions $$\dot{\varphi }_a(\sigma ):=\{H_{\mathrm{gen}},\varphi _a(\sigma )\}=\lambda ^b(\sigma )M_{ba}(\sigma )0,$$ (69) where $$\{\varphi _a(\sigma ),\varphi _b(\sigma ^{})\}=M_{ab}(\sigma )\delta (\sigma \sigma ^{}),M_{ab}=hϵ_{abc}X^{}{}_{}{}^{c}.$$ (70) Clearly, (69) imposes no further phase space constraints. It simply sets to zero the Lagrange multipliers in the directions where the matrix $`M_{ab}`$ is non–degenerate. These directions correspond to second class constraints. In a direction where $`M_{ab}`$ is degenerate the Lagrange multiplier remains undetermined and such a direction corresponds to a first class constraint. Hence, if $`\|\stackrel{}{X}^{}\|0`$ then there are three first class constraints and the phase space is trivial. On the other hand, provided that $`\|\stackrel{}{X}^{}\|0`$ the matrix $`M_{ab}`$ has only one zero eigenvector, given by $`\lambda ^a(\sigma )=\lambda (\sigma )X^{}^a`$. The matrix $`M_{ab}`$ is therefore non–degenerate in the two–dimensional subspace orthogonal to $`\stackrel{}{X}^{}`$. Thus one introduces a projection onto this subspace as follows ($`I=1,2`$): $`P_I{}_{}{}^{a}(\sigma )P_J{}_{}{}^{b}(\sigma )\delta _{ab}`$ $`=`$ $`\delta _{IJ},`$ (71) $`\delta ^{IJ}P_I{}_{}{}^{a}(\sigma )P_J{}_{}{}^{b}(\sigma )`$ $`=`$ $`\delta ^{ab}{\displaystyle \frac{X^aX^b}{\|\stackrel{}{X}^{}\|^2}},`$ (72) $`ϵ^{IJ}P_I{}_{}{}^{a}(\sigma )P_J{}_{}{}^{b}(\sigma )`$ $`=`$ $`{\displaystyle \frac{ϵ^{abc}X^c}{\|\stackrel{}{X}^{}\|}}.`$ (73) The three constraints $`\varphi _a`$ now split into the two second class constraints $$\chi _I:=P_I{}_{}{}^{a}\varphi _{a}^{},$$ (74) with the now non–degenerate matrix $$\{\chi _I(\sigma ),\chi _J(\sigma ^{})\}:=M_{IJ}(\sigma )\delta (\sigma \sigma ^{}),M_{IJ}=P_I{}_{}{}^{a}P_{J}^{}{}_{}{}^{b}M_{ab}^{},$$ (75) and one first class constraint $$\varphi :=X^{}{}_{}{}^{a}\varphi _{a}^{}X^a\mathrm{\Pi }_a,$$ (76) which acts as the generator of $`\sigma `$–reparameterisations. The resulting generalised Hamiltonian, $$H_{\mathrm{gen}}=𝑑\sigma \lambda (\sigma )\varphi (\sigma ),$$ (77) leads to the canonical field equations $$\dot{X}^a(\sigma )=\lambda (\sigma )X^{}{}_{}{}^{a}(\sigma ),$$ (78) which are indeed consistent with the Lagrangian field equations (61) when $`\|\stackrel{}{X}^{}\|0`$. We remark that the first class nature of the constraint $`\varphi `$ means that the gauge invariant dynamics is trivially realised (i.e. $`\dot{V}=0`$) on reparameterisation invariant functionals $`V[\stackrel{}{X}]`$ satisfying $$𝑑\sigma \eta (\sigma )X^a(\sigma )\frac{\delta V}{\delta X^a(\sigma )}0.$$ (79) The analysis of the $`X^\alpha `$ coordinates is similar, but we require this system to admit the solution $`X^0=\tau `$. This implies, from the equations of motion, that we must consider the trivial sector $`X^\alpha 0`$. (In this sector $`X^0=\tau `$ is a gauge choice for the constraint $`\varphi _00`$.) To summarise, the two admissible sectors of the theory given above are: (i) $`X^\alpha =0`$ and $`\|\stackrel{}{X}^{}\|=0`$; and (ii) $`X^\alpha =0`$ and $`\|\stackrel{}{X}^{}\|0`$. Let us continue by deriving the symplectic structure of the phase space. In the 0,1,2 directions there are no second class constraints and we find vanishing brackets between the $`X^{0,1,2}`$ coordinates. We next turn to the $`X^{3,4,5}`$ coordinates. In the string sector, the presence of the second class constraints leads us to define a Dirac bracket as follows: $$[A,B]^D:=\{A,B\}𝑑\sigma \{A,\chi _I(\sigma )\}(M^1)^{IJ}(\sigma )\{\chi _J(\sigma ),B\}.$$ (80) Here $`A`$ and $`B`$ are general phase space variables and $`(M^1)^{IJ}`$ is the ordinary matrix inverse of the matrix $`M_{IJ}`$ defined in (75). The basic Dirac brackets between the coordinates $`X^a`$ and the first class constraint $`\varphi `$ are: $$[X^a(\sigma ),X^b(\sigma ^{})]^D=\frac{1}{h}\frac{ϵ^{abc}X^{}{}_{}{}^{c}(\sigma )}{\|\stackrel{}{X}^{}(\sigma )\|^2}\delta (\sigma \sigma ^{}),$$ (81) $$[X^a(\sigma ),\varphi (\sigma ^{})]^D=X^a(\sigma )\delta (\sigma \sigma ^{})$$ (82) $$[\varphi (\sigma ),\varphi (\sigma ^{})]^D=2\varphi (\sigma )\delta ^{}(\sigma \sigma ^{})+\varphi ^{}(\sigma )\delta (\sigma \sigma ^{}).$$ (83) From the point of view of the five–brane the fields $`X^a`$ are worldvolume coordinates. The Dirac bracket (81) therefore describes a noncommutative loop space on the five–brane. Finally, we discuss three special topics: dimensional reduction, the canonical analysis in the infinite momentum frame and the Moyal quantisation in case of a compact direction. ### Double Dimensional Reduction We wrap the five–brane and the membrane around a spacetime circle of radius $`R`$. In the limit of small $`R`$ the five–brane becomes a D4–brane and the wrapped membrane becomes a fundamental string. Thus one recovers the case of a string ending on a D4–brane in the background of constant Neveu–Schwarz two–form potential. Explicitly, taking winding numbers into account, we find for a five–brane winding $`M`$ times in spacetime and a string winding $`N`$ times in the five–brane $$X^5(\sigma )=NR\sigma ,_{\mu \nu 5}=\frac{M}{R}_{\mu \nu },$$ (84) where the index $`\mu `$ now refers to the D4–brane worldvolume and $``$ is normalised such that it obeys $`d=H^{NS}`$ where $`H^{NS}`$ is the $`D=10`$ Neveu–Schwarz three–form field strength. The membrane winds $`MN`$ times in spacetime and the membrane action (1) reduces to $`MNS_{F1}`$, where $`S_{F1}`$ is the $`D=10`$ string action (27). For illustrative purposes, let us choose $`_{345}=h`$ and $`_{012}=h/\sqrt{1+\mathrm{}_p^6h^2}`$. For this choice we find a rank two magnetic fieldThis result can be generalized to a rank $`4`$ magnetic field by either reducing over a skew spacelike direction in the $`05`$ plane or boosting the constant $``$ background in the $`05`$ direction turning on $`_{034}`$ and $`_{125}`$ components, and then reduce over the fifth direction. $$_{34}=\frac{hR}{M},$$ (85) leading to a noncommutative D–brane worldvolume with $$[X^3,X^4]=\frac{1}{MN}(^1)^{34}.$$ (86) Alternatively, the double dimensional reduction can be performed directly at the level of the algebra (81)-(83) by fixing the gauge $`X^5(\sigma )=NR\sigma `$. The gauge fixed Dirac bracket $`[,]^D^{}`$ reads $$[X^3(\sigma ),X^4(\sigma ^{})]^D^{}=\frac{1}{hNR}\delta (\sigma \sigma ^{}).$$ (87) This bracket reduces to (86) provided we make use of (85) and identify $`X^{3,4}`$ with the zero–modes $`X^{3,4}=𝑑\sigma X^{3,4}(\sigma )`$ of the doubly reduced string. ### Canonical Analysis in the infinite momentum frame In the infinite momentum frame discussed in Appendix C the field strength $``$ is degenerate along a null direction: $$_{\mu \nu }=0.$$ (88) The remaining components of the field strength are given by (see Eq. (C9) in Appendix C): $$_{+pq}=F_{pq}^{}=\frac{1}{2}ϵ_{pqrs}F_{rs}^{},_{pqr}=0,p=1,2,3,4.$$ (89) Thus the lightcone coordinate $`X^{}`$ is decoupled from the remaining fields $`X^+`$ and $`X^p`$ in the string action (54). The action has the reparameterisation invariance: $$\delta _\eta X^+=\eta ^i_iX^+,\delta _\eta X^p=\eta ^i_iX^p.$$ (90) In the non–trivial sector the embedding of the worldsheet is a two–dimensional surface and the gauge fixed solutions to the field equations can be taken to be $$X^{}=\sigma ^{},_{}X^+=_{}X^p=0,$$ (91) where $`\sigma ^\pm =\frac{1}{\sqrt{2}}(\tau \pm \sigma )`$. To give a phase space formulation, we introduce canonical momenta with respect to two–dimensional lightcone velocities as follows: $$\mathrm{\Pi }_\mu (\sigma ^+):=\frac{\delta S}{\delta _{}X^\mu (\sigma ^+)},\mu =+,p.$$ (92) The first order nature of the Lagrangian implies that there are five primary constraints $$\varphi _\mu :=\mathrm{\Pi }_\mu +\frac{1}{3}_{\mu \nu \rho }X^\nu _+X^\rho 0,\mu =+,p.$$ (93) Since $`F_{pq}^{}`$ has rank four, it follows that the bracket $`\{\varphi _\mu ,\varphi _\nu \}`$ has rank four and zero direction given by $`_+X^\mu `$. Hence there is one first class constraint, that we can choose to be $$\varphi :=_+X^\mu \varphi _\mu _+X^\mu \mathrm{\Pi }_\mu ,\mu =+,p,$$ (94) and four second class constraints $`\chi _I`$, $`I=1,2,3,4`$. Provided that $`_+X^+0`$, we can choose these to be $$\chi _I=\delta _I^p\varphi _p.$$ (95) The basic Dirac brackets are found to be $`[X^p(\sigma _1^+),X^q(\sigma _2^+)]^D`$ $`=`$ $`{\displaystyle \frac{[(F^{})^1]^{pq}}{_+X^+}}\delta (\sigma _1^+\sigma _2^+),`$ (96) $`[X^+(\sigma _1^+),X^\mu (\sigma _2^+)]^D`$ $`=`$ $`0,\mu =+,p,`$ (98) and the analogs of (82) and (83). In a background with a compact light–like direction of radius $`R`$, the gauge choice $`X^+=NR\sigma ^+`$ leads to a local, field independent bracket analogous to (87): $$[X^p(\sigma _1^+),X^q(\sigma _2^+)]^D=\frac{[(F^{})^1]^{pq}}{NR}\delta (\sigma _1^+\sigma _2^+).$$ (99) A double dimensional reduction along the compact direction gives the Dirac bracket for the rank $`4`$ noncommutative D4–brane with Born-Infeld field strength $`_{\mu \nu }=\frac{R}{M}\delta _\mu ^p\delta _\nu ^qF_{pq}^{}`$ where $`M`$ is the number of times the five–brane winds around the compact direction. ### Moyal quantisation Usually after deriving the noncommutative structure one determines the star product. This cannot be done in a direct way for the bracket (81). We will discuss a proposal for how this might be done in the following discussion section. However, in the special case where the five–brane has a compact spacelike direction, the gauge fixed Dirac brackets (87) are straightforward to quantise by introducing the following Moyal product $`(I,J=3,4)`$: $$F(X)G(X)=\mathrm{exp}\left[\frac{1}{2hNR}𝑑\sigma ϵ^{IJ}\frac{\delta }{\delta X^I(\sigma )}\frac{\delta }{\delta Y^J(\sigma )}\right]F(X)G(Y)\|_{X=Y},$$ (100) where $`F`$ and $`G`$ are functionals of $`X^{3,4}`$. Similarly, in a background with a compact light–like direction of radius $`R`$, the Moyal quantised star–product, based upon the brackets (99), is given by ($`p=1,2,3,4)`$: $$F[X]G[X]=\mathrm{exp}\left[\frac{1}{2NR}𝑑\sigma ^+[(F^{})^1]^{pq}\frac{\delta }{\delta X^p(\sigma ^+)}\frac{\delta }{\delta Y^q(\sigma ^+)}\right]F[X]G[Y]\|_{X=Y},$$ (101) where $`F`$ and $`G`$ are functionals of $`X^p`$. ## VI Discussions In this paper we have studied the effect of having a constant $`C`$–field on the M–theory five–brane using an open membrane probe. We proposed a specific decoupling limit in which the open membrane metric remains non–degenerate. An unconventional feature of this decoupling limit, not encountered in the case of a string probing a D–brane, is that the open membrane action scales with an overall factor, as was discussed in Section IV. We hope that our work provides the motivation and initial steps to study noncommutative loop algebras. A natural question to ask in this context is whether we can define a star product like in the case of D–branes. The utility of the star product is that to incorporate the effects of the noncommutativity of the space on fields one simply replaces ordinary products with star products to give the deformed theory. In the case of a point particle moving in a finite dimensional Poisson manifold $`M`$ with Poisson structure $`\mathrm{\Omega }^{ij}`$, one may in principle always find local canonical coordinates in which $`\mathrm{\Omega }^{ij}`$ is a constant matrix and the star product then becomes simple. The reparameterisation independent definition of the star product was first given by Kontsevich and has recently been cast into a more physical context by the work of . The latter work makes use of a path integral representation of the star product. There is a natural extension of the definition of star product given in that applies to the case of the open membrane/five–brane system. A physically intuitive definition of the star product of two loop functionals $`F[X]`$ and $`G[X]`$ is given by the following path integral expression $$(FG)[X(\sigma )]=\underset{\underset{\tau \pm \mathrm{}}{lim}X(\tau ,\sigma )=X(\sigma )}{}𝒟XF[X(\tau =1,\sigma )]G[X(\tau =0,\sigma )]\mathrm{exp}\frac{i}{\mathrm{}}S[X],$$ (102) with the action S given by (54). In case the five–brane winds around a compact direction in the background we expect (102) to reproduce the Moyal products given in (100) and (101). It would be interesting to see whether the proposed path integral representation of the deformed five–brane geometry will facilitate the construction of a loop space analog of the star product of ordinary fields. It will be desirable to represent the loop algebra also on ordinary fields such as those of the massless $`(2,0)`$ tensor multiplet on the five–brane. One possible direction for such constructions may be to introduce so called loop space covariant derivatives , which naturally incorporates the two–form potential in the tensor multiplet. The main motivation for introducing a noncommutative loop algebra in this work was provided by a study of the M–theory five–brane. However, there are also other reasons to study the same system. For instance, the action (54) is formally the same as the action for a string interacting with a linear background Neveu–Schwarz two–form. Here we identify $`B_{\mu \nu }^{\mathrm{NS}}=\frac{1}{3}_{\mu \nu \rho }X^\rho `$ such that the Neveu–Schwarz field strength, $`H^{\mathrm{NS}}=dB^{\mathrm{NS}}`$ is a constant. However, despite the resemblance of the terms in the action, the conditions on the background fields differ. In the previous case $``$ must obey the field equations of the five–brane, in the case of the string with a background Neveu–Schwarz field strength, the background fields must obey the Einstein’s equations for the appropriate supergravity. One can construct a solution with constant Neveu–Schwarz field strength, but the space is necessarily curved and one must again take care when discussing limits. It is interesting that the action (54) also arises from the dimensional reduction of the higher order Chern–Simons term that appears in $`11`$–dimensional supergravity $`S=_{M^{11}}CHH`$. More precisely, take $`M^{11}=M^2\times M^9`$ and then reduce on $`M^9`$ as follows: Let $`C=X^I\gamma _I`$ where $`\gamma _IH^3(M^9),I=1\mathrm{}b^3(M^9)`$ and $`X^I`$ are scalars on $`M^2`$, one then identifies $`H`$ as the triple intersection form on $`M^9`$, $`H_{IJK}=_{M^9}\gamma _I\gamma _J\gamma _K`$. For this case the only constraints on the background field $`H`$ is that it is a triple intersection form. The kinetic terms of the supergravity will introduce period matrices of $`\{\gamma _I\}`$ that will play the role of a metric for $`X^I`$. In $`5`$–dimensional supergravity there is an analogous treatment of the $`A_1F_2F_2`$ term. It would be interesting to see whether (compactified) supergravity theories allow limits in which only the Wess–Zumino terms survive. It is worth remarking that in the decoupling limit we defined in this work, it appears possible to take seriously a quantisation program for the membrane. This is because one has decoupled the bulk modes and so one is left with only the boundary string in the five–brane. It would be interesting to see how far one can carry out the quantisation of such a system. In doing this one should use the full supersymmetric branes. Finally, it is known that for the case of the D–brane the noncommutative U(1) theory is non–abelian. It is our hope that the development of the noncommutative M5–brane might shed light on the as yet unknown non–abelian structure of the five–brane worldvolume theory. Note Added: After this paper appeared, we received the work of in which related ideas are discussed from a different point of view. ## Acknowledgements We would like to thank Martin Cederwall, Chong-Sun Chu, Christiaan Hofman, Whee Ky Ma and Ergin Sezgin for discussions. D.B. would like to thank the organisers and participants of the Leiden conference on Noncommutative Geometry. This work is supported by the European Commission TMR program ERBFMRX-CT96-0045, in which E.B., D.B. and J.P.v.d.S. are associated to the University of Utrecht. The work of J.P.v.d.S. and P.S. is part of the research program of the “Stichting voor Fundamenteel Onderzoek der Materie” (FOM). ## A The Wess–Zumino term of the membrane In this appendix we derive the form of the Wess–Zumino term of the membrane as given in (13). Assume that $`f_5^{}F=0`$ and $``$ is constant. Then we can write $`f_5^{}C=dc`$ and $`=d(b+c)`$, at least locally. The two Wess–Zumino terms in (1) can therefore be written as a Wess–Zumino term for the bulk modes perpendicular to the five–brane plus a string Wess–Zumino term as follows: $$_{M^3}f_2^{}\stackrel{~}{C}+_{M^3}f_1^{}(b+c),$$ (A1) where $`\stackrel{~}{C}`$ obeys $`f_5^{}\stackrel{~}{C}=0`$. To see this, one writes $`C=\stackrel{~}{C}+dC_2`$ where $`f_5^{}C_2=c`$. One then applies Stokes’ theorem to $`_{M^3}C`$ and uses the identity $$f_2^{}C_2\|_{M^3}=f_1^{}f_5^{}C_2=f_1^{}c.$$ (A2) Finally, one uses the fact that, since $``$ is constant, we have the identity $$(b+c)_{\mu \nu }=\frac{1}{3}_{\mu \nu \rho }X^\rho .$$ (A3) Combining everything then leads to the form (13) of the Wess–Zumino term for the string in the five–brane. ## B The $`SO(5,1)/SO(2,1)\times SO(3)`$ parameterisation of $`_{\mu \nu \rho }`$ In this appendix we parameterise a generic solution to the non–linear self–duality condition (6) using a real scale factor and a coset element in $`SO(5,1)/SO(2,1)\times SO(3)`$. The special solutions corresponding to an infinite momentum frame are described in the next appendix. The non–linear self–duality condition (6) is equivalent to a linear self–duality condition $`_{\mu \nu \rho }`$ $`=`$ $`4(m^1)_\mu {}_{}{}^{\lambda }h_{\nu \rho \lambda }^{},h_{\mu \nu \rho }={\displaystyle \frac{\sqrt{g}}{6}}ϵ_{\mu \nu \rho \sigma \lambda \tau }h^{\sigma \lambda \tau },`$ (B1) $`m_{\mu \nu }`$ $`=`$ $`g_{\mu \nu }2\mathrm{}_p^6k_{\mu \nu },k_{\mu \nu }=h_{\mu \lambda \rho }h_\nu {}_{}{}^{\lambda \rho }.`$ (B2) It can be shown that the matrix $`k_{\mu \nu }`$ is traceless and that its square is proportional to $`\eta _{\mu \nu }`$. Hence $`k_{\mu \nu }`$ can be written as $$k_{\mu \nu }=\lambda _+v_\mu ^+v_\nu ^++\lambda _{}v_\mu ^{}v_\nu ^{}+\lambda \underset{I=1}{\overset{2}{}}(e_\mu ^Ie_\nu ^I+e_\mu ^{I+2}e_\nu ^{I+2}),$$ (B3) where the three real parameters $`\lambda _\pm `$ and $`\lambda `$ obey $`\lambda _+\lambda _{}=\lambda ^2`$ and $`(v_\mu ^\pm ,e_\mu ^p)`$, $`p=1,2,3,4`$, is an element of the coset $`SO(5,1)/SO(1,1)\times SO(4)`$ defined by $$v^{\pm \mu }v_\mu ^\pm =\eta ^{\pm \pm },v^{\pm \mu }e_\mu ^p=0,e^{p\mu }e_\mu ^q=\delta ^{pq},$$ (B4) $$\eta ^+=\eta ^+=1,\eta ^{++}=\eta ^{}=0,ϵ^{+pqrs}=ϵ_{+pqrs}=ϵ^{pqrs}.$$ (B5) If $`\lambda _\pm 0`$ (the cases $`\lambda _+=\lambda =0`$ and $`\lambda _{}=\lambda =0`$ correspond to the infinitely boosted solutions discussed in the next appendix) we can use a boost to set $`\lambda _+=\lambda _{}=\lambda `$ such that $$k_{\mu \nu }=\lambda \eta _{\alpha \beta }v_\mu ^\alpha v_\nu ^\beta +\lambda \delta _{ab}u_\mu ^au_\nu ^b,$$ (B6) where $`(v_\mu ^\alpha ,u_\mu ^a)`$, $`\alpha =0,1,2`$, $`a=3,4,5`$, is an element of the coset $`SO(5,1)/SO(2,1)\times SO(3)`$ obeying (17). Keeping the coset element fixed, any three–form can be expanded in this basis as follows $$\omega _{\mu \nu \rho }=Au_{\mu \nu \rho }^3+B_\alpha ^a(u^2v)_a^\alpha {}_{\mu \nu \rho }{}^{}+C_a^\alpha (uv^2)_\alpha ^a{}_{\mu \nu \rho }{}^{}+Dv_{\mu \nu \rho }^3,$$ (B7) where the coefficients together make up twenty real components and we have defined the following basis elements for three–forms: $`u_{\mu \nu \rho }^3`$ $`=`$ $`ϵ_{abc}u_\mu ^au_\nu ^bu_\rho ^c,(u^2v)_a^\alpha {}_{\mu \nu \rho }{}^{}=ϵ_{abc}u_{[\mu }^bu_\nu ^cv_{\rho ]}^\alpha ,`$ (B8) $`(uv^2)_\alpha ^a_{\mu \nu \rho }`$ $`=`$ $`ϵ_{\alpha \beta \gamma }u_{[\mu }^av_\nu ^\beta v_{\rho ]}^\gamma ,v_{\mu \nu \rho }^3=ϵ_{\alpha \beta \gamma }v_\mu ^\alpha v_\nu ^\beta v_\rho ^\gamma .`$ (B9) For our choice of Lorentz basis, the Hodge $``$ acts as follows $$u^3=v^3,(uv^2)^{a,\alpha }=(u^2v)^{a,\alpha }.$$ (B10) Hence, we can write a general self–dual tensor in the form (B7) using the ten coefficients $`A=D`$ and $`B^{a,\alpha }=C^{a,\alpha }`$. In particular, we can apply this expansion to $`h_{\mu \nu \rho }`$ itself, which yields the following expression for $`k_{\mu \nu }`$: $`k_{\mu \nu }`$ $`=`$ $`2A^2(u_{\mu \nu }^2v_{\mu \nu }^2)+{\displaystyle \frac{2}{9}}\left[\left(B^2\delta _{ab}2B_{ab}^2\right)u_\mu ^au_\nu ^b+\left(B^2\eta _{\alpha \beta }+2B_{\alpha \beta }^2\right)v_\mu ^\alpha v_\nu ^\beta \right]`$ (B12) $`+{\displaystyle \frac{4}{9}}ϵ_{abc}ϵ_{\alpha \beta \gamma }B^{a,\alpha }B^{b,\beta }u_{(\mu }^cv_{\nu )}^\gamma ,`$ where we have written $`B_{\alpha \beta }^2=\delta ^{ab}B_{a,\alpha }B_{b,\beta }`$, $`B_{ab}^2=\eta ^{\alpha \beta }B_{a,\alpha }B_{b,\beta }`$ and $`B^2=\delta ^{ab}B_{ab}^2=\eta ^{\alpha \beta }B_{\alpha \beta }^2`$. However, by assumption the matrix $`k_{\mu \nu }`$ must be diagonal in this basis. This implies that there must exist real numbers $`X`$ and $`Y`$ such that $$B_{ab}^2=X\delta _{ab},B_{\alpha \beta }^2=Y\eta _{\alpha \beta }.$$ (B13) This implies that $`B_{a,\alpha }`$ has to vanish since it is an irreducible representation of $`SO(3)\times SO(2,1)`$. This can be verified explicitly by first taking the determinant of the equations in (B13) which shows that $`X=Y<0`$, and then deducing from $`B_{1,1}^2=B_{2,2}^2=Y<0`$ that $`B_{a,1}=B_{a,2}=0`$ and hence $`X=Y=0`$. It then follows from $`B_{0,0}^2=0`$ that also $`B_{a,0}=0`$. Substituting $`B_{a,\alpha }=C_{a,\alpha }=0`$ into Eq. (B7) gives $$h_{\mu \nu \rho }=A\left(u_{\mu \nu \rho }^3+v_{\mu \nu \rho }^3\right).$$ (B14) From this is follows that the matrix $`m_{\mu \nu }`$ in Eq. (B2) is given by $$m_{\mu \nu }=(14\mathrm{}_p^6A^2)u_{\mu \nu }^2+(1+4\mathrm{}_p^6A^2)v_{\mu \nu }^2.$$ (B15) After application of (B2) one finally obtains the expressions (14) and (15) for the non–linearly self–dual three–form and the open membrane metric provided one makes the following identification: $$h=\frac{4A}{14\mathrm{}_p^6A^2}.$$ (B16) ## C The decoupling limit in the infinite momentum frame In this appendix we discuss a modified version of the decoupling limit of Section IV where a boost parameter is scaled such that it becomes infinite in the limit. This results in a decoupled string action of the form (54) where the constant background field strength is given by applying the infinite boost to the generic solution (14). The infinitely boosted field strength is most naturally described in terms of an $`SO(5,1)/SO(1,1)\times SO(4)`$ coset element (B4) and a self–dual two–form in the four–dimensional Euclidean space perpendicular to the boost direction. ### The $`SO(5,1)/SO(1,1)\times SO(4)`$ parameterisation of $`_{\mu \nu \rho }`$ Given a coset element (B4) an arbitrary three–form $`h_{\mu \nu \rho }`$ can be parameterised in terms of two anti–symmetric matrices $`F_{pq}`$ and $`F_{pq}^{}`$ and two vectors $`G_p`$ and $`G_p^{}`$ by giving its components with the respect to the frame as follows ($`p=1,2,3,4)`$: $$h_{+pq}=F_{pq}^{}+F_{pq}^+,h_{pq}=F_{pq}^++F_{pq}^{},$$ (C1) $$h_{+p}=G_p+G_p^{},h_{pqr}=ϵ_{pqrs}(G_sG_s^{}),$$ (C2) where $`F_{pq}^\pm =\frac{1}{2}(F_{pq}\pm ϵ_{pqrs}F_{rs})`$ and $`F_{pq}^\pm =\frac{1}{2}(F_{pq}^{}\pm ϵ_{pqrs}F_{rs}^{})`$. A self–dual three–form is obtained by setting $`F_{pq}^{}=G_p^{}=0`$. By an $`SO(5,1)`$ rotation we can fix a coset frame such that $`G_p=0`$. In this adapted frame the self–dual three–form then has the expansion $$h_{\mu \nu \rho }=3v_{[\mu }^+e_\nu ^pe_{\rho ]}^qF_{pq}^{}+3v_{[\mu }^{}e_\nu ^pe_{\rho ]}^qF_{pq}^+=3v_{[\mu }^+F_{\nu \rho ]}^{}+3v_{[\mu }^{}F_{\nu \rho ]}^+,$$ (C3) where $`F_{\mu \nu }^\pm =e_\mu ^pe_\nu ^qF_{pq}^\pm `$. The local $`SO(1,1)\times SO(4)`$ symmetry can be fixed by taking $`v_\mu ^\pm =(1,\pm \widehat{n})`$, where $`\widehat{n}`$ is a unit vector in $`𝐑^5`$, and choosing $`e_\mu ^p`$ such that $`F_{pq}^\pm =f_\pm (i\sigma ^2(\pm i\sigma ^2))_{pq}`$, where $`f_\pm `$ are two real parameters. This makes a total of ten independent degrees of freedom. From (B2) and (9) it follows that the non–linearly self–dual three–form and the open membrane metric are given in this parameterisation by $$_{\mu \nu \rho }=\frac{12}{14\mathrm{}_p^{12}(F^+)^2(F^{})^2}\left[v_{[\mu }^+F_{\nu \rho ]}^{}+v_{[\mu }^{}F_{\nu \rho ]}^+2\mathrm{}_p^6(F^{})^2v_{[\mu }^+F_{\nu \rho ]}^+2\mathrm{}_p^6(F^+)^2v_{[\mu }^{}F_{\nu \rho ]}^{}\right],$$ (C4) $`G_{\mu \nu }={\displaystyle \frac{1}{\left(14\mathrm{}_p^{12}(F^+)^2(F^{})^2\right)^2}}`$ $`[(1+4\mathrm{}_p^{12}(F^+)^2(F^{})^2)g_{\mu \nu }`$ (C5) $`+`$ $`4\mathrm{}_p^6v_\mu ^+v_\nu ^+(F^{})^2+4\mathrm{}_p^6v_\mu ^{}v_\nu ^{}(F^+)^216\mathrm{}_p^6F_{\mu \rho }^{}F_\nu ^+{}_{}{}^{\rho }].`$ (C6) where we have defined $`(F^\pm )^2=F^{\pm \mu \nu }F_{\mu \nu }^\pm =F^{\pm pq}F_{pq}^\pm `$. Provided that $`F_{\mu \nu }^\pm 0`$ this parameterisation is equivalent to the parameterisation in (14). An infinite boost along the direction $`\widehat{n}`$ is obtained by setting $`F_{\mu \nu }^+=0`$ (corresponding to the special case $`\lambda _{}=\lambda =0`$ mentioned under Eq. (B5)) and results in the expressions $`_{\mu \nu \rho }`$ $`=`$ $`12v_{[\mu }^+F_{\nu \rho ]}^{},`$ (C7) $`G_{\mu \nu }`$ $`=`$ $`g_{\mu \nu }+4\mathrm{}_p^6v_\mu ^+v_\nu ^+(F^{})^2.`$ (C8) We notice that (C7) is actually linearly self–dual and the second term in (C8) drops out of the non–linear condition (6). More explicitly, in the infinite momentum frame the components of $``$ are $$_{+pq}=4F_{pq}^{},_{pq}=0,_{pqr}=_{+p}=0.$$ (C9) ### The Infinite Momentum Decoupling Limit In order to define a decoupling limit leading to (54) with the field strength given by (C7), we first boost (14) along the $`5`$ direction: $$\left(\begin{array}{c}v^0\hfill \\ u^5\hfill \end{array}\right)=\left(\begin{array}{cc}\sqrt{1+a^2}& a\\ a& \sqrt{1+a^2}\end{array}\right)\left(\begin{array}{c}\stackrel{~}{v}^0\hfill \\ \stackrel{~}{u}^5\hfill \end{array}\right),$$ (C10) and then keep $`\stackrel{~}{v}^0`$ and $`\stackrel{~}{u}^5`$ fixed, while scaling according to (45) and $$aϵ^\gamma a,\gamma >0.$$ (C11) The Wess–Zumino term now has weight $`\gamma \mathrm{\Delta }`$, such that in the rescaled action (where the Wess–Zumino term is fixed) the kinetic energy of the parallel modes has weight $`2+\delta 1+\gamma +\mathrm{\Delta }>0`$ and the kinetic energy of the perpendicular modes has weight $`2+\gamma +\mathrm{\Delta }<0`$. The open membrane metric becomes $`G_{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{\left(1+\sqrt{1+h^2\mathrm{}_p^6}\right)^2}{4}}[\delta _{ab}u_\mu ^au_\nu ^b+{\displaystyle \frac{1}{1+h^2\mathrm{}_p^6}}\eta _{\alpha \beta }v_\mu ^\alpha v_\nu ^\beta `$ (C12) $`+`$ $`{\displaystyle \frac{a^2h^2\mathrm{}_p^6}{1+h^2\mathrm{}_p^6}}(v_\mu ^0v_\nu ^0+2\sqrt{1+{\displaystyle \frac{1}{a^2}}}v_{(\mu }^0u_{\nu )}^5+u_\mu ^5u_\nu ^5)],`$ (C13) where we have dropped the tildes on $`v^0`$ and $`u^5`$. Requiring the leading components of this matrix to form a rank six matrix, we see that $`a^2h^2\mathrm{}_p^6`$ has to be finite, i.e. $`\gamma +\mathrm{\Delta }+\frac{3}{2}(\delta 1)3`$. Thus the conditions on $`\gamma +\mathrm{\Delta }`$ are the same as the conditions on $`\mathrm{\Delta }`$ given in (53). Since $`h\mathrm{}_p^30`$, the result of the limit is the decoupled boundary Wess–Zumino term (54) with background field strength now given by the infinitely boosted solution (C7) provided we identify $$v_\mu ^+=\frac{1}{\sqrt{2}}(v_\mu ^0+u_\mu ^5),F_{\mu \nu }^{}=\frac{ah}{\sqrt{2}}(v_{[\mu }^1v_{\nu ]}^2+u_{[\mu }^3u_{\nu ]}^4).$$ (C14)
warning/0005/hep-th0005008.html	ar5iv	text	# 1 Introduction ## 1 Introduction Recent developments in non-perturbative string theory and M theory have led to new insights into the relation between low energy field theory and string theory: it has been argued in particular that non-perturbative dynamics takes place in low energy field theory and higher gauge symmetries emerge when compactifying Calabi-Yau and $`K_3`$ manifolds degenerate and some of their homology cycles vanish. For instance, when $`K_3`$ surface develops an $`A`$-$`D`$-$`E`$ singularity, there appears an enhanced $`A`$-$`D`$-$`E`$ gauge symmetry in 6-dimensions. Similarly, when a family of $`𝐏^1`$’s shrinks to zero size along a rational curve in Calabi-Yau threefold, we obtain a 4-dimensional $`SU(2)`$ gauge theory with $`𝒩=2`$ supersymmetry. In this article we would like to study the dynamics of SUSY gauge theories in 5-dimensions by compactifying the M theory on degenerate Calabi-Yau manifolds. We construct an effective action of an $`SU(2)`$ gauge theory on a 5-dimensional space $`M_4\times S^1`$ where $`M_4`$ is the Minkowki space and $`S^1`$ is a circle of radius $`R`$. In the limit of $`R0`$ this theory reproduces the standard $`𝒩=2`$ SUSY gauge theory in 4-dimensions, i.e. Seiberg-Witten theory . Thus our 5-dimensional model gives an M-theoretic generalization of Seiberg-Witten theory by incorporating Kaluza-Klein excitations. In the opposite limit $`R\mathrm{}`$ our model reduces to the gauge theory in uncompactified 5-dimensions $`M_5`$. Characteristic features of this theory have been studied using the brane-probe picture and also from the point of view of classical geometry of collapsing del Pezzo surfaces and the behavior of the low-energy effective gauge coupling has been determined exactly. In this paper we would like to propose an exact solution of the 5-dimensional theory on $`M_4\times S^1`$ which reproduces the known results at both limits $`R0`$ and $`R\mathrm{}`$. As it turns out, our prepotential follows directly from that of the type II string theory compactified on singular Calabi-Yau manifolds using the method of local mirror symmetry . In the case of pure $`SU(2)`$ gauge theory without matter our result is obtained from the type II theory compactified on the local $`𝐅_2`$, i.e. Hirzebruch surface $`𝐅_2`$ lying inside a Calabi-Yau threefold which is the canonical bundle over $`𝐅_2`$. Similarly, $`SU(2)`$ gauge theories with $`N_f`$ matter in vector representations are also obtained from the type II theory compactified on the local $`𝐅_2`$ blown up at $`N_f`$ points ($`0N_f4`$). We will find that our model at $`R=\mathrm{}`$ has an infinite bare coupling constant and yields a non-trivial interacting field theory in the infra-red limit. The local mirror symmetry is a method of mirror symmetry adapted in the case of non-compact Calabi-Yau manifolds. Suppose, for instance, we are given a compact Calabi-Yau threefold which is an elliptic fibration over $`𝐅_n`$. One considers the limit of the size of the fiber $`t_E`$ going to $`\mathrm{}`$. Then the resulting non-compact manifold is modeled by the local $`𝐅_n`$, i.e. $`𝐅_n`$ inside a Calabi-Yau with the normal bundle being given by the canonical bundle of $`𝐅_n`$. The limit of $`t_E\mathrm{}`$ may also be considered as the limit of shrinking $`𝐅_n`$ with the size of the fiber kept fixed. $`𝐅_n`$ ($`n=0,1,2`$) and its various blow ups are the del Pezzo surfaces and these are in fact the type of manifolds which featured in the geometrical interpretation of 5-dimensional gauge theory . ## 2 $`SU(2)`$ Gauge Theory without Matter Let us start from the case of 5-dimensional gauge theory without matter. We consider the local $`𝐅_2`$ model which is described by the toric data given in the Appendix. Following the standard procedure one obtains a curve in the B-model given by $$P=a_0x+a_1x^2+a_2\zeta +a_3+a_4\frac{1}{\zeta }=0.$$ (2.1) Introducing a new variable $`y=a_2\zeta a_4/\zeta `$ the curve is rewritten as $$y^2=(a_1x^2+a_0x+a_3)^24a_2a_4.$$ (2.2) Complex moduli of the B-model are defined by $$z_F=\frac{a_1a_3}{a_0^2},z_B=\frac{a_2a_4}{a_3^2}.$$ (2.3) If we choose $`a_1=a_3=1,a_0=s,4a_2a_4=K^4`$, we find $$y^2=(x^2+sx+1)^2K^4,z_F=\frac{1}{s^2},z_B=\frac{K^4}{4}.$$ (2.4) (2.4) is in fact the curve proposed by Nekrasov for the description of 5-dimensional gauge theory and its properties have been studied in in detail. If we introduce a variable $`U`$ which is the analogue of $`u`$ of the Seiberg-Witten solution, the parameter $`s`$ is written as $$s=2R^2U$$ (2.5) where $`R`$ is the radius of $`S^1`$. In terms of $`U`$ and $`R`$ the curve reads as $$y^2=(x^2R^4(U^2\frac{1}{R^4}))^2K^4.$$ (2.6) By comparing (2.6) with the Seiberg-Witten curve $`y^2=(x^2u)^2\mathrm{\Lambda }^4`$, we find the correspondence between the parameters $`U`$ and $`u`$ $$uR^2(U^2\frac{1}{R^4}).$$ (2.7) As we see from (2.7), $`U`$ variable describes two copies of the $`u`$-plane. Strong coupling region $`u0`$ maps to $`U\pm 1/R^2`$ and thus the two copies are separated by a distance of order $`1/R^2`$. In the brane-probe interpretation of Seiberg-Witten solution , $`u`$-plane is identified as a local region around one of the four fixed points (O-7 planes) in type I’ theory compactification to 8-dimensions on $`T^2`$ ($`u=0`$ is identified as the location of the O-7 plane). Then the curve (2.6) describes a theory which contains two of these O-7 planes. Since the fixed plane acts like a reflecting mirror, D3-brane probe will possess an infinite number of mirror images when inserted into a background of two orientifold planes. These mirror images are separated by distances $`n/R^2,n𝐙`$ and open strings connecting them generate Kaluza-Klein modes of supersymmetric gauge fields. Thus the curve (2.6) effectively describes a theory on a 5-dimensional manifold $`M_4\times S^1`$ with $`R`$ being the radius of $`S^1`$. In the limit of $`R0`$ one of the $`u`$-planes moves off to $`\mathrm{}`$ and the model (2.6) is expected to reduce to the Seiberg-Witten theory. Periods of the B-model of local $`𝐅_2`$ (2.4) is determined by solving differential equations (Gelfand-Kapranov-Zelevinskij (GKZ) system) associated with the toric data (see Appendix). Differential operators are given by $`_1=z_F(4\theta _{z_F}^2+2\theta _{z_F})+\theta _{z_F}(2\theta _{z_B}\theta _{z_F}),`$ (2.8) $`_2=z_B(2\theta _{z_B}\theta _{z_F}+1)(2\theta _{z_B}\theta _{z_F})\theta _{z_B}^2,`$ (2.9) where $$\theta _{z_F}z_F\frac{}{z_F},\theta _{z_B}z_B\frac{}{z_B}.$$ (2.10) These operators have a regular singular point at $`z_F=z_B=0`$: they possess “single-log” solutions $`\omega _F,\omega _B`$ behaving as $`\omega _F=\mathrm{log}z_F+\mathrm{}`$ and $`\omega _B=\mathrm{log}z_B+\mathrm{}`$ at $`z_F,z_B0`$. There also exists a “double-log” solution $`\mathrm{\Omega }`$ behaving as $`\mathrm{\Omega }=(\mathrm{log}z_F)^2+(\mathrm{log}z_F)(\mathrm{log}z_B)+\mathrm{}`$. We identify the two single-log solutions as the Kähler parameters $`t_F,t_B`$ of the A-model: $`t_F`$ represents the size of the $`𝐏^1`$ fiber of $`𝐅_2`$ and $`t_B`$ the size of its base $`𝐏^1`$. $`t_F`$ is given by $`t_F\omega _F=\mathrm{log}z_F+\left[{\displaystyle \underset{n1,m0}{}}{\displaystyle \frac{2(2n1)!}{(n2m)!n!m!^2}}z_F^nz_B^m{\displaystyle \underset{m1}{}}{\displaystyle \frac{(2m1)!}{m!^2}}z_B^m\right],`$ $`=2\mathrm{log}\left({\displaystyle \frac{1}{\sqrt{4z_F}}}+\sqrt{{\displaystyle \frac{1}{4z_F}}1}\right)`$ $`+\left[{\displaystyle \underset{n1,m1}{}}{\displaystyle \frac{2(2n1)!}{(n2m)!n!m!^2}}z_F^nz_B^m{\displaystyle \underset{m1}{}}{\displaystyle \frac{(2m1)!}{m!^2}}z_B^m\right].`$ (2.11) Similarly, $`t_B`$ is given by $`t_B\omega _B=2\mathrm{log}\left({\displaystyle \frac{1}{\sqrt{4z_B}}}+\sqrt{{\displaystyle \frac{1}{4z_B}}1}\right).`$ (2.12) In our interpretation as the 5-dimensional gauge theory, the size of the fiber $`t_F`$ is identified as the vacuum expectation value $`A`$ of the scalar field in the vector-multiplet $$t_F=4RA.$$ (2.13) On the other hand, the size of the base $`t_B`$ is related to the dynamical mass parameter $`\mathrm{\Lambda }`$ as $$e^{t_B}=4R^4\mathrm{\Lambda }^4.$$ (2.14) Then the mirror transformation (2.12) becomes $$K=\frac{2R\mathrm{\Lambda }}{\sqrt{1+4R^4\mathrm{\Lambda }^4}}.$$ (2.15) Note that (2.13) and (2.14) are in fact the identification of variables suggested in . One can invert the relations (2.11),(2.12) perturbatively and express the B-model parameters $`z_F`$,$`z_B`$ in terms of $`t_F,t_B`$. In the case of local mirror symmetry the holomorphic solution of GKZ system is a constant and the mirror transformation is simpler than in the compact Calabi-Yau case. One may then represent the double-log solution in terms of the Kähler parameters $`\mathrm{\Omega }=t_F^2+t_Ft_B+4{\displaystyle \underset{n=1}{}}{\displaystyle \frac{1}{n^2}}q_F^n+4q_B\left({\displaystyle \underset{n=1}{}}n^2q_F^n\right)`$ $`+q_B^2\left(q_F^2+36q_F^3+260q_F^4+1100q_F^5+\mathrm{}\right)+𝒪(q_B^3),`$ (2.16) where $$q_F=e^{t_F},q_B=e^{t_B}.$$ (2.17) We identify the double-log solution as $`A_D`$, the dual of the variable $`A`$ $$\mathrm{\Omega }=8\pi iRA_D.$$ (2.18) First two terms of (2.16) represent the classical intersection numbers of the Calabi-Yau manifold and the remaining terms represent the contribution of world-sheet instantons. According to ref the double-log solution has a generic form $$\mathrm{\Omega }=\underset{i,j=1}{\overset{2}{}}t_it_jJ^iJ^j+\underset{k=1}{}\underset{n,m0}{}(\underset{i}{}x^i_{t_i})d_{n,m}\frac{q_1^{kn}q_2^{km}}{k^3},$$ (2.19) for a local model of a surface $`S`$ (with two Kähler parameters). $`J_i`$ denotes the Kähler classes of $`S`$ and $`J_iJ_j`$ their intersection numbers. Numerical coefficients $`x^i`$ are defined by $$c_1(S)=\underset{i}{}x^iJ_i,$$ (2.20) where $`c_1(S)`$ is the first Chern class of $`S`$. $`d_{n,m}`$ gives the number of rational holomorphic curves in the homology class $`nJ_1+mJ_2`$. Sum over $`k`$ in (2.19) represents the multiple-cover factor. By comparing (2.16) and (2.19) (set $`q_1=q_F,q_2=q_B`$) we find $`d_{1,0}=2,d_{n,0}=0,n>1,d_{n,1}=2n,`$ $`d_{1,2}=d_{2,2}=0,d_{3,2}=6,d_{4,2}=32,\mathrm{}`$ (2.21) By integrating $`A_D`$ over $`A`$ we have the prepotential $`F`$ for the local $`𝐅_2`$ model $`F={\displaystyle \frac{1}{32\pi iR^2}}[{\displaystyle \frac{t_F^3}{3}}{\displaystyle \frac{t_F^2t_B}{2}}+4{\displaystyle \underset{n=1}{}}{\displaystyle \frac{1}{n^3}}q_F^n+4q_B\left({\displaystyle \underset{n=1}{}}nq_F^n\right)`$ $`+q_B^2({\displaystyle \frac{1}{2}}q_F^2+12q_F^3+65q_F^4+220q_F^5+\mathrm{})+𝒪(q_B^3)].`$ (2.22) ## 3 Small and Large Radius Limits Let us next examine the small and large radius limits $`R0,\mathrm{}`$ of (2). We first consider the 4-dimensional limit $`R0`$. Due to the relations (2.13), (2.14) small $`R`$ corresponds to the base of $`𝐅_2`$ becoming large ($`t_B\mathrm{}`$) and the fiber becoming small ($`t_F0`$) while the ratio $`e^{t_B}/t_F^4`$ is kept fixed. At small $`R`$, we have $`K2R\mathrm{\Lambda }`$ and $`U\mathrm{cosh}(2RA)`$. As explained by Katz, Klemm and Vafa , quantum parts of $`F`$ are suppressed because of the powers of $`q_BR^4`$ and the surviving contributions come from the divergent parts of the series $`d_{n,m}q_F^n`$ over $`q_F`$ as $`q_F=e^{4RA}1`$. At small $`R`$ the gauge coupling $`\tau =^2F/A^2`$ behaves as $`\tau ={\displaystyle \frac{4i}{\pi }}RA{\displaystyle \frac{i}{2\pi }}\mathrm{log}(4R^4\mathrm{\Lambda }^4){\displaystyle \frac{2i}{\pi }}{\displaystyle \underset{n=1}{}}{\displaystyle \frac{q_F^n}{n}}{\displaystyle \frac{8i}{\pi }}R^4\mathrm{\Lambda }^4{\displaystyle \underset{n=1}{}}n^3q_F^n+\mathrm{}`$ $`{\displaystyle \frac{2i}{\pi }}\mathrm{log}(2\sqrt{2}{\displaystyle \frac{A}{\mathrm{\Lambda }}}){\displaystyle \frac{3i}{16\pi }}{\displaystyle \frac{\mathrm{\Lambda }^4}{A^4}}+\mathrm{}`$ (3.1) (3.1) reproduces the one-loop beta function and the one-instanton contribution to the Seiberg-Witten solution . We can check the agreement with Seiberg-Witten theory for higher instanton terms. In (3.1) the world-sheet instanton expansion of type II theory is converted into the space-time instanton expansion of gauge theory in the $`R0`$ limit. Coefficients of the $`m`$-instanton amplitudes of gauge theory are determined by the asymptotic behavior of the number of holomorphic curves $`d_{n,m}`$ as $`n\mathrm{}`$ with fixed $`m`$. Let us next examine the $`R\mathrm{}`$ limit of the uncompactified 5-dimensional gauge theory. It is known that the gauge theory on $`M_5`$ has no quantum corrections and its gauge coupling is simply expressed in terms of classical intersection numbers of Calabi-Yau manifold. In fact by taking $`R\mathrm{}`$ world-sheet instanton terms disappear and we have $$\underset{R\mathrm{}}{lim}\frac{\tau }{2\pi iR}\tau _5=\frac{2}{\pi ^2}A,$$ (3.2) (we have rescaled $`\tau `$ so that $`\tau _5`$ corresponds to the gauge coupling of 5-dimensional theory, $`\tau _5=1/g_5^2`$). In the next section we will discuss local $`𝐅_2`$ model blown up at $`N_f`$ points. We then find that the above formula is generalized as $$\tau _5=\frac{2}{\pi ^2}\left(A\underset{i=1}{\overset{N_f}{}}\frac{1}{16}\|AM_i\|\underset{i=1}{\overset{N_f}{}}\frac{1}{16}\|A+M_i\|\right).$$ (3.3) (3.3) is exactly the behavior of gauge coupling of $`SU(2)`$ theory with $`N_f`$ matter in vector representations (at infinite bare coupling) . Thus we reproduce correct results also in the uncompactified limit $`R\mathrm{}`$. We should note that the local models of $`𝐅_0`$ and $`𝐅_1`$ also reproduce Seiberg-Witten solution in the $`R0`$ limit since the asymptotic behavior of the number of holomorphic curves $`d_{n,m}`$ are the same for all models $`𝐅_i,i=0,1,2`$. $`𝐅_0,𝐅_1`$, however, have a different classical topology from $`𝐅_2`$ and do not reproduce (3.2) at $`R=\mathrm{}`$. Thus the $`𝐅_2`$ model is singled out as the unique candidate for the description of 5-dimensional gauge theory on $`M_4\times S^1`$. ## 4 $`SU(2)`$ theory with matter Let us next consider the case of $`SU(2)`$ gauge theory coupled to $`N_f`$ matter ($`1N_f4`$) in vector representations. We first discuss the $`N_f=1`$ case. Relevant geometry is given by the local $`𝐅_2`$ with one-point blown up. Corresponding curve is given by (see Appendix) $$y^2=(x^2+sx+1)^2K^3(xw+\frac{1}{w}).$$ (4.1) Comparison with the $`N_f=1`$ Seiberg-Witten curve suggests the identification $$w=e^{RM},$$ (4.2) where $`M`$ denotes the bare mass of the matter multiplet. Complex moduli of the B-model are given by (we denote $`z_{FE}`$ as $`z_F`$ and $`z_{F+E}`$ as $`z_F^{}`$ for notational simplicity) $$z_F=\frac{w^2}{s},z_F^{}=\frac{1}{sw^2},z_B=\frac{K^3}{4w}.$$ (4.3) In this case the GKZ system is given by five partial differential equations in three variables (see Appendix) and is of considerable complexity. Here we content ourselves with the analysis of prepotential at the tree and one-loop level ignoring the space-time instantons. This suffices for our purpose of extracting the $`R\mathrm{}`$ behavior of the theory. The small $`R`$ behavior has already been studied in and argued to reproduce Seiberg-Witten theory (we have also verified that the 1-instanton term is correctly reproduced). When we ignore instantons, single log solutions are given by $`t_F\omega _F=\mathrm{log}(z_F)+{\displaystyle \underset{n=1}{}}{\displaystyle \frac{(2n1)!}{n!^2}}(z_Fz_F^{})^n,`$ (4.4) $`t_F^{}\omega _F^{}=\mathrm{log}(z_F^{})+{\displaystyle \underset{n=1}{}}{\displaystyle \frac{(2n1)!}{n!^2}}(z_Fz_F^{})^n,`$ (4.5) $`t_B\omega _B=\mathrm{log}(z_B).`$ (4.6) Identification with the variables of the gauge theory is given by $$t_F=2R(AM),t_F^{}=2R(A+M),e^{t_B}=\frac{2R^3\mathrm{\Lambda }^3}{w}.$$ (4.7) Then by inverting relations (4.4),(4.5) we find $$z_F=\frac{e^{2R(AM)}}{1+e^{4RA}},z_F^{}=\frac{e^{2R(A+M)}}{1+e^{4RA}}.$$ (4.8) The double log solution $`\mathrm{\Omega }=t_F^2+{\displaystyle \frac{1}{2}}t_F^{}^2+2t_Ft_F^{}+t_B(t_F+t_F^{})+{\displaystyle \underset{n=1}{}}{\displaystyle \frac{4(2n1)!}{n!^2}}{\displaystyle \underset{j=1}{\overset{n}{}}}{\displaystyle \frac{1}{j}}(z_Fz_F^{})^n`$ $`+\left({\displaystyle \underset{n>\mathrm{}}{}}+{\displaystyle \underset{n<\mathrm{}}{}}\right){\displaystyle \frac{(n+\mathrm{}1)!}{n!\mathrm{}!}}{\displaystyle \frac{1}{n\mathrm{}}}z_F^nz_F^{}^{\mathrm{}},`$ (4.9) can then be re-expressed as $`\mathrm{\Omega }=t_F^2+{\displaystyle \frac{1}{2}}t_F^{}^2+2t_Ft_F^{}+t_B(t_F+t_F^{})`$ $`+4{\displaystyle \underset{n=1}{}}{\displaystyle \frac{1}{n^2}}e^{4nRA}{\displaystyle \underset{n=1}{}}{\displaystyle \frac{1}{n^2}}e^{2nR(A+M)}{\displaystyle \underset{n=1}{}}{\displaystyle \frac{1}{n^2}}e^{2nR(AM)},`$ (4.10) using (4.8). When we take the derivative in $`A`$, (4.10) gives $$\tau =\frac{i}{\pi }\left[2\mathrm{log}\mathrm{sinh}(2RA)\frac{1}{4}\mathrm{log}\mathrm{sinh}(R(A+M))\mathrm{sinh}(R(AM))\right].$$ (4.11) In the 4-dimensional limit (4.11) becomes $$\tau \frac{i}{\pi }(2\frac{1}{2})\mathrm{log}A$$ (4.12) which gives the 1-loop beta function of Seiberg-Witten theory. On the other hand in the 5-dimensional limit we find $$\tau _5=\underset{R\mathrm{}}{lim}\frac{\tau }{2\pi iR}=\frac{2}{\pi ^2}\left[A\frac{1}{16}\|AM\|\frac{1}{16}\|A+M\|\right]$$ (4.13) as we have claimed before. We note that in formulas (3.2) and (4.13) $`\tau _5`$ does not contain an additive constant which is identified as the bare coupling constant $`\tau _{5,B}=1/g_{5,B}^2`$. Thus the theory sits at the infinite bare coupling constant limit $`g_{5,B}^2=\mathrm{}`$ which yields non-trivial 5-dimensional theory in the infra-red regime . We can similarly study the system with more matter up to $`N_f=4`$ using the local mirror symmetry. Curves of the B-model are given by (see Appendix) $$y^2=(x^2+sx+1)^2K^{4N_f}\underset{i=1}{\overset{N_f}{}}(w_ix+\frac{1}{w_i}),$$ (4.14) where parameters $`w_i`$ correspond to the bare masses of the matter multiplets $$w_i=e^{RM_i},i=1,\mathrm{},4$$ (4.15) (4.14) agrees with the curve suggested by and (at $`N_f=4`$ the factor $`K^{4N_f}`$ should be replaced by a dimensionless parameter $`q`$). In the cases $`N_f2`$ the analysis of GKZ system becomes further involved: we will instead use a simpler method based on the Picard-Fuchs (PF) equation derived for the elliptic curves (4.14). It turns out that the PF equation of the elliptic curve is an ordinary differential equation of third order in the variable $`A`$ and can be studied relatively easily. This is the method used in ref . This system, however, is not complete unlike that of GKZ case. The periods are determined only up to integration constants and one can not precisely fix the mirror transformation. This ambiguity, however, affects only the quantum part of the computation and one still obtains precise results for the prepotential at the tree and one-loop level. First we note that the quadratic curve $`y^2=ax^4+4bx^3+6cx^2+4dx+e`$ is transformed into the Weierstrass form $$y^2=4x^3g_2xg_3,$$ (4.16) by the relation $$g_2=ae4bd+3c^2,g_3=ace+2bcdad^2b^2ec^3.$$ (4.17) We regard $`g_2,g_3`$ as functions of the parameter $`s`$. Periods $`\omega `$ of the elliptic curve (4.16) obey the PF equation $$\frac{d^2\omega }{ds^2}+c_1(s)\frac{d\omega }{ds}+c_0(s)\omega =0,$$ (4.18) where $`c_1={\displaystyle \frac{d}{ds}}\mathrm{log}\left({\displaystyle \frac{3}{2\mathrm{\Delta }}}(2g_2{\displaystyle \frac{dg_3}{ds}}3{\displaystyle \frac{dg_2}{ds}}g_3)\right),`$ (4.19) $`c_0={\displaystyle \frac{1}{12}}c_1{\displaystyle \frac{d}{ds}}\mathrm{log}\mathrm{\Delta }+{\displaystyle \frac{1}{12}}{\displaystyle \frac{d^2}{ds^2}}\mathrm{log}\mathrm{\Delta }{\displaystyle \frac{1}{16\mathrm{\Delta }}}\left(g_2({\displaystyle \frac{dg_2}{ds}})^212({\displaystyle \frac{dg_3}{ds}})^2\right)`$ (4.20) and $`\mathrm{\Delta }=(g_2)^327(g_3)^2`$ denotes the discriminant of the curve. Since $`dA/ds`$ is one of the periods of the curve, it satisfies the PF equation. Regarding $`s`$ as a function of $`A`$, we obtain $$\frac{ds}{dA}\frac{d^3s}{dA^3}3(\frac{d^2s}{dA^2})^2+c_1(\frac{ds}{dA})^2\frac{d^2s}{dA^2}c_0(\frac{ds}{dA})^4=0.$$ (4.21) This determines $`s`$ in terms of $`A`$. Similarly $`dA_D/ds=dA/dsd^2F/dA^2`$ satisfies the PF equation and we obtain $$(\frac{ds}{dA})^1\frac{d^4F}{dA^4}3(\frac{ds}{dA})^2\frac{d^2s}{dA^2}\frac{d^3F}{dA^3}+c_1\frac{d^3F}{dA^3}=0.$$ (4.22) This equation can be integrated once and we find $$\frac{d^3F}{dA^3}=\frac{const.}{\mathrm{\Delta }}\left(2g_2\frac{dg_3}{ds}3\frac{dg_2}{ds}g_3\right)\left(\frac{ds}{dA}\right)^3.$$ (4.23) Solving (4.23) determines the prepotential (const. in the right-hand-side is fixed by a suitable normalization of $`F`$). In the case $`N_f=2`$, we find the gauge coupling at the classical and one-loop level as $$\tau =\frac{i}{\pi }\left[2\mathrm{log}\mathrm{sinh}(2RA)\frac{1}{4}\underset{i=1}{\overset{2}{}}\mathrm{log}\mathrm{sinh}(R(A+M_i))\mathrm{sinh}(R(AM_i))\right].$$ (4.24) In the 4-dimensional limit (4.24) reads as $$\tau \frac{2i}{\pi }(2\frac{1}{2}2)\mathrm{log}A,$$ (4.25) which gives the beta function of $`N_f=2`$ Seiberg-Witten theory. On the other hand, in the five-dimensional limit we have $$\tau _5=\frac{2}{\pi ^2}\left(A\underset{i=1}{\overset{2}{}}\frac{1}{16}\|AM_i\|\underset{i=1}{\overset{2}{}}\frac{1}{16}\|A+M_i\|\right)$$ (4.26) in agreement with (3.3). We have checked that the local model of $`𝐅_2`$ blown up at 3 and 4 points (see Appendix) also reproduce (3.3). Thus we have obtained a model which appears to describe correctly the physics of 5-dimensional theory up to the number of flavors $`N_f=4`$ making use of the mirror symmetry. ## 5 Discussions Since $`𝐅_2`$ is obtained from $`𝐅_1`$ which is a one-point blow up of $`𝐏^2`$, we effectively have up to 5-point blow ups of $`𝐏^2`$ describing the gauge theory. Unfortunately, beyond the 5-point blow up mirror symmetry of del Pezzo surfaces can not be described by toric geometry and we can not apply our analysis for these cases. In fact in the range $`N_f5`$ something drastic must happen, since in this asymptotically non-free region Seiberg-Witten solution does not exist and we should not have a smooth 4-dimensional limit. On the other hand, in this range the 5-dimensional gauge theories are expected to possess $`E_n,n=6,7,8`$ global symmetries and are of particular interests. It is a challenging problem to clarify the physics of gauge theory with $`E_n`$ symmetry. It may shed some light on the nature of asymptotically non-free field theories in 4-dimensions. Our result shows that from the point of view of the type II/M theory compactified on Calabi-Yau manifold, low energy 5-dimensional theory on $`M_4\times S^1`$ emerges most naturally with its prepotential being exactly the same as that of the string theory. On the other hand, 4-dimensional Seiberg-Witten theory appears only in the fine-tuned limit of the Kähler parameters. It seems possible that our 5-dimensional model is an example of an “M-theoretic” lift of various 4-dimensional quantum field theories. We may imagine most of the 4-dimensional SUSY field theories in fact have a lift to 5-dimensions where the quantum effects of the loops and instantons are replaced by purely geometrical effects of world-sheet instantons. It will be also quite interesting to see if there is a further lift of quantum field theory to 6-dimensions as suggested by the duality between F- and M-theory. We would like to thank M. Jinzenji, M. Naka and Y. Ohta for discussions. Researches of T.E. and H.K. are supported by the fund for the Special Priority Area No.707, Japan Ministry of Education. We also acknowledge the stimulating atmosphere at Summer Institute ’99 where this research was started. ## Appendix ### A.1 $`N_f=0`$ For pure Yang-Mills case we take the Hirzebruch surface $`𝐅_2`$ lying inside a Calabi-Yau manifold. Its toric diagram has five vertices; $$\nu _0=(0,0),\nu _1=(1,0),\nu _2=(0,1),\nu _3=(1,0),\nu _4=(2,1).$$ (A.1) This is a two dimensional reflexive polyhedra no.4 in Fig.1 of . The charge vectors satisfying the linear relation $$\underset{i}{}\mathrm{}_i^{(k)}\nu _i=0$$ (A.2) are given by $$\mathrm{}^{(F)}=(2;1,0,1,0),\mathrm{}^{(B)}=(0;0,1,2,1).$$ (A.3) The constraint among B-model variables $`Y_i,i=0,1,\mathrm{},4`$ $$\underset{\mathrm{}_i^{(k)}>0}{}Y_i^{\mathrm{}_i^{(k)}}=\underset{\mathrm{}_i^{(k)}<0}{}Y_i^{\mathrm{}_i^{(k)}}$$ (A.4) gives $$Y_0^2=Y_1Y_3,Y_2Y_4=Y_3^2,$$ (A.5) and we have a solution $$Y=(sx,x^2,\zeta ,s^2,\frac{s^4}{\zeta }).$$ (A.6) Set $`s=1`$. Then we obtain the following curve in the B model side $`P_{N_f=0}`$ $`=`$ $`{\displaystyle a_iY_i},`$ (A.7) $`=`$ $`a_0x+a_1x^2+a_2\zeta +a_3+a_4{\displaystyle \frac{1}{\zeta }}=0.`$ Introducing a new variable $`y=a_2\zeta (a_4/\zeta )`$, we can rewrite the curve as $$y^2=(a_1x^2+a_0x+a_3)^24a_2a_4.$$ (A.8) For each charge vector $`\mathrm{}^{(k)}`$, a corresponding complex structure modulus is given by $$z_k=\underset{i}{}a_i^{\mathrm{}_i^{(k)}}.$$ (A.9) In the present case we have $$z_F=\frac{a_1a_3}{a_0^2},z_B=\frac{a_2a_4}{a_3^2}.$$ (A.10) By setting $$a_1=a_3=1,a_0=s,4a_2a_4=K^4,$$ (A.11) we have $$z_F=\frac{1}{s^2},z_B=\frac{K^4}{4}.$$ (A.12) GKZ system is defined by a system of differential operators $$\underset{\mathrm{}_i^{(k)}>0}{}\left(\frac{}{a_i}\right)^{\mathrm{}_i^{(k)}}=\underset{\mathrm{}_i^{(k)}<0}{}\left(\frac{}{a_i}\right)^{\mathrm{}_i^{(k)}},$$ (A.13) In the $`𝐅_2`$ case it is given by $`_1=z_F(4\theta _{z_F}^2+2\theta _{z_F})+\theta _{z_F}(2\theta _{z_B}\theta _{z_F}),`$ (A.14) $`_2=z_B(2\theta _{z_B}\theta _{z_F}+1)(2\theta _{z_B}\theta _{z_F})\theta _{z_B}^2,`$ (A.15) $`\theta _{z_i}z_i{\displaystyle \frac{}{z_i}},i=F,B.`$ ### A.2 $`N_f=1`$ It is known that matters in the vector representation are generated by blowing up the manifold (see, for intance ). We expect that a blow up of the local $`𝐅_2`$ provides the description of gauge theory with $`N_f=1`$ matter. The corresponding reflexive polyhedra is no.6 of and given by the vertices $`\nu _0`$ $`=`$ $`(0,0),\nu _1=(1,0),\nu _2=(0,1),`$ $`\nu _3`$ $`=`$ $`(1,1),\nu _4=(1,0),\nu _5=(1,1).`$ (A.16) We see that the charge vectors are $`\mathrm{}^{(B)}=(0;0,0,1,2,1),`$ $`\mathrm{}^{(FE)}=(1;0,1,1,1,0),`$ $`\mathrm{}^{(E)}=(1;1,1,1,0,0).`$ (A.17) The constraint $$Y_3Y_5=Y_4^2,Y_0Y_3=Y_2Y_4,Y_0Y_2=Y_1Y_3,$$ (A.18) is solved by $$Y=(sx,sx^2,\frac{tx}{\zeta },\frac{s}{\zeta },s,\zeta s).$$ (A.19) Setting $`s=1`$ gives the curve $$P_{N_f=1}=a_0x+a_1x^2+a_2\frac{x}{\zeta }+a_3\frac{1}{\zeta }+a_4+a_5\zeta ,$$ (A.20) or $$y^2=(a_1x^2+a_0x+a_4)^24a_5(a_2x+a_3).$$ (A.21) If we set $$a_1=a_4=1,a_0=s,4a_5=K^3,a_2=w,a_3=w^1,$$ (A.22) we obtain the curve $$y^2=(x^2+sx+1)^2K^3(wx+\frac{1}{w})$$ (A.23) Complex moduli are given by $$z_B=\frac{a_3a_5}{a_4^2}=\frac{K^3}{4w},z_{FE}=\frac{a_2a_4}{a_0a_3}=\frac{w^2}{s},z_{F+E}=\frac{a_1a_3}{a_0a_2}=\frac{1}{sw^2}.$$ (A.24) Complete set of differential equations is given by $`_1=\theta _B(\theta _B\theta _{FE}+\theta _{F+E})z_B(2\theta _B\theta _{FE})(2\theta _B\theta _{FE}+1),`$ (A.25) $`_2=(\theta _{FE}\theta _{F+E})(\theta _{FE}2\theta _B)z_{FE}(\theta _{FE}+\theta _{F+E})(\theta _{FE}\theta _B\theta _{F+E}),`$ (A.26) $`_3=\theta _{F+E}(\theta _{F+E}+\theta _B\theta _{FE})z_{F+E}(\theta _{F+E}+\theta _{FE})(\theta _{F+E}\theta _{FE}),`$ (A.27) $`_4=(\theta _{FE}\theta _{F+E})\theta _Bz_Bz_{FE}(\theta _{FE}+\theta _{F+E})(2\theta _B\theta _{FE}),`$ (A.28) $`_5=\theta _{F+E}(\theta _{FE}2\theta _B)z_{FE}z_{F+E}(\theta _{FE}+\theta _{F+E})(\theta _{FE}+\theta _{F+E}+1).`$ (A.29) We note that the last two operators are necessary to obtain a unique double log solution $`\mathrm{\Omega }`$. ### A.3 $`N_f=2`$ We choose the following vertices $`\nu _0`$ $`=`$ $`(0,0),\nu _1=(1,0),\nu _2=(0,1),\nu _3=(1,1),`$ $`\nu _4`$ $`=`$ $`(1,0),\nu _5=(1,1),\nu _6=(1,1).`$ (A.30) This is a reflexive polyhedra no.8 of . The curve is given by $$P_{N_f=2}=a_0x+a_1x^2+a_2\frac{x}{\zeta }+a_3\frac{1}{\zeta }+a_4+a_5\zeta +a_6\frac{x^2}{\zeta }=0,$$ (A.31) or $$y^2=(a_1x^2+a_0x+a_4)^24a_5(a_6x^2+a_2x+a_3).$$ (A.32) Substituting the relation $`a_1`$ $`=`$ $`a_4=1,a_0=s,4a_5=K^2,`$ $`a_6`$ $`=`$ $`w_1w_2,a_2=\left({\displaystyle \frac{w_2}{w_1}}+{\displaystyle \frac{w_1}{w_2}}\right),a_3=(w_1w_2)^1,`$ (A.33) we obtain the complex structure moduli $`z_B`$ $`=`$ $`{\displaystyle \frac{a_3a_5}{a_4^2}}={\displaystyle \frac{K^2}{4w_1w_2}},z_{F_1}={\displaystyle \frac{a_2a_4}{a_0a_3}}={\displaystyle \frac{w_1^2+w_2^2}{s}},`$ $`z_{F_2}`$ $`=`$ $`{\displaystyle \frac{a_1a_3}{a_0a_2}}={\displaystyle \frac{1}{s(w_1^2+w_2^2)}},z_{F_3}={\displaystyle \frac{a_1a_2}{a_0a_6}}={\displaystyle \frac{w_1^2+w_2^2}{sw_1^2w_2^2}}.`$ (A.34) ### A.4 $`N_f=3`$ Polyhedron for the 3-point blow up is obtained by adding a vertex $`\nu _7=(0,1)`$ to (A.30) (no.12 of ). Elliptic curve is given by $$P_{Nf=3}=a_0x+a_1x^2+a_2\frac{x}{\zeta }+a_3\frac{1}{\zeta }+a_4+a_5\zeta +a_6\frac{x^2}{\zeta }+a_7x\zeta =0.$$ (A.35) or $$y^2=(a_1x^2+a_0x+a_4)^24(a_6x^2+a_2x+a_3)(a_7x+a_5).$$ (A.36) By choosing the variables as $`a_0=s,a_1=1,a_2={\displaystyle \frac{K}{4}}({\displaystyle \frac{w_1}{w_2}}+{\displaystyle \frac{w_2}{w_1}}),a_3={\displaystyle \frac{K}{4w_1w_2}},a_4=1,a_5={\displaystyle \frac{1}{w_3}},`$ $`a_6={\displaystyle \frac{Kw_1w_2}{4}},a_7=w_3,`$ (A.37) we find the complex moduli $`z_B={\displaystyle \frac{a_3a_5}{a_4^2}}={\displaystyle \frac{K}{4w_1w_2w_3}},z_{F_1}={\displaystyle \frac{a_2a_4}{a_0a_3}}={\displaystyle \frac{1}{s}}(w_1^2+w_2^2),z_{F_2}={\displaystyle \frac{a_1a_3}{a_0a_2}}={\displaystyle \frac{1}{s}}{\displaystyle \frac{1}{w_1^2+w_2^2}},`$ $`z_{F_3}={\displaystyle \frac{a_1a_2}{a_0a_6}}={\displaystyle \frac{1}{s}}({\displaystyle \frac{1}{w_1^2}}+{\displaystyle \frac{1}{w_2^2}}),z_{F_4}={\displaystyle \frac{a_4a_7}{a_0a_5}}={\displaystyle \frac{w_3^2}{s}}.`$ ### A.5 $`N_f=4`$ Polyhedron for the 4-point blow up is obtained by further adding a vertex $`\nu _8=(1,1)`$ to the polyhedron of $`N_f=3`$ (no.15 of ). Elliptic curve is given by $$P_{N_f=4}=a_0x+a_1x^2+a_2\frac{x}{\zeta }+a_3\frac{1}{\zeta }+a_4+a_5\zeta +a_6\frac{x^2}{\zeta }+a_7x\zeta +a_8x^2\zeta =0.$$ (A.39) By eliminating $`\zeta `$ it becomes $$y^2=(a_1x^2+a_0x+a_4)^24(a_6x^2+a_2x+a_3)(a_8x^2+a_7x+a_5).$$ (A.40) By choosing the variables as $`a_0=s,a_1=1,a_2=({\displaystyle \frac{w_1}{w_2}}+{\displaystyle \frac{w_2}{w_1}})q,a_3={\displaystyle \frac{1}{w_1w_2}}q,a_4=1,a_5={\displaystyle \frac{1}{w_3w_4}},`$ $`a_6=w_1w_2q,a_7=({\displaystyle \frac{w_3}{w_4}}+{\displaystyle \frac{w_4}{w_3}}),a_8=w_3w_4`$ (A.41) we find the complex moduli $`z_B={\displaystyle \frac{a_3a_5}{a_4^2}}={\displaystyle \frac{q}{w_1w_2w_3w_4}},z_{F_1}={\displaystyle \frac{a_2a_4}{a_0a_3}}={\displaystyle \frac{1}{s}}(w_1^2+w_2^2),z_{F_2}={\displaystyle \frac{a_1a_3}{a_0a_2}}={\displaystyle \frac{1}{s}}{\displaystyle \frac{1}{w_1^2+w_2^2}},`$ $`z_{F_3}={\displaystyle \frac{a_1a_2}{a_0a_6}}={\displaystyle \frac{1}{s}}({\displaystyle \frac{1}{w_1^2}}+{\displaystyle \frac{1}{w_2^2}}),z_{F_4}={\displaystyle \frac{a_4a_7}{a_0a_5}}={\displaystyle \frac{w_3^2+w_4^2}{s}},z_{F_5}={\displaystyle \frac{a_1a_7}{a_0a_8}}={\displaystyle \frac{1}{s}}({\displaystyle \frac{1}{w_3^2}}+{\displaystyle \frac{1}{w_4^2}}).`$
warning/0005/cond-mat0005298.html	ar5iv	text	# 1 Introduction ## 1 Introduction In the present paper we consider the six-vertex model (ice model) on the two-dimensional square lattice with the special boundary conditions of “domain wall” . First the six-vertex model was solved in the papers , where the partition function and bulk free energy were found for the case of periodically boundary conditions. An explicit expression for the partition function with domain wall boundary conditions on the lattice of finite volume was found by A. G. Izergin . Izergin’s representation contains a determinant of $`N\times N`$ matrix (where $`N^2`$ is the number of vertices of the lattice). Thus, the problem of the thermodynamic limit of the partition function can be reduced to the study of the asymptotic behavior of the determinant at $`N\mathrm{}`$. The goal of this paper is to obtain new representation for the partition function, appropriate for the asymptotic analysis. A Fredholm determinant representation seems to be suitable for this goal. Let us explain in brief the main idea of the method suggested. Let a matrix $``$ of the size $`N`$ is given, whose entries $`_{jk}(t)`$ depend on some parameter $`t`$. Determinants of this matrix $`det`$ generate a sequence of functions $`f_N(t)`$ for different $`N`$. As a rule the question on the properties of the limiting function $`f_{\mathrm{}}(t)=lim_N\mathrm{}f_N(t)`$ is rather complicated. In spite of $`det`$ is a finite sum of finite products, nevertheless such “explicit” expression for large $`N`$ becomes extremely complicated for analysis. Suppose however that we can replace $`det`$ with a Fredholm determinant $$det=det(I+V)$$ of an integral operator, whose kernel parametrically depends on $`t`$ and $`N`$: $`V(x,y)=V_N(x,y\|t)`$. The Fredholm determinant is an infinite sum, and each term of the last one can be presented as a multiple integral, for example, $$det(I+V)=\mathrm{exp}\left(\underset{k=1}{\overset{\mathrm{}}{}}\frac{(1)^{k+1}}{k}V_N(x_1,x_2\|t)\mathrm{}V_N(x_k,x_1\|t)d^kx\right).$$ While at finite $`N`$ the computations of the mentioned integrals are very complicated, at $`N\mathrm{}`$ the existence of large parameter in the integrands often allows to estimate not only each integral, but the series as a whole. Thus we have the possibility to obtain an estimate for the determinant. Hence, we can say that for $`N\mathrm{}`$ representations in terms determinants of finite size matrices become “less explicit”, while Fredholm determinants turn to be “more explicit”. Transformations of finite determinants into Fredholm ones were found to be useful, for instance, in asymptotic analysis of correlation functions of quantum integrable models, asymptotic behavior of orthogonal polynomials, the spectrum of random matrices etc. We hope that in the case under consideration the Fredholm determinant representation will permit eventually to solve the problem of the thermodynamic limit of the partition function of the six-vertex model<sup>1</sup><sup>1</sup>1Recently the thermodynamic limit of the six-vertex model was studied by use of Toda-chain equation and matrix models methods . The approach suggested in the present paper is completely different from the mentioned ones, and we hope that it is also worth attention.. The content of the paper is as follows. In the next section the necessary information on the six-vertex model and Izergin’s formula for the partition function are given. In the third section we transform this formula to the Fredholm determinant representation. In the last section we discuss the outlook for application of the obtained formula to the asymptotic analysis. ## 2 The six-vertex model Consider a statistical system on a square lattice of the size $`N\times N`$, associating with each edge some variables (classical spin), taking values $`\pm 1`$. Usually these variables are denoted by arrows, identifying $`(,)`$ with $`1`$, and $`(,)`$ with $`1`$. Thus, in general case there are $`16`$ different configurations of arrows in each vertex. The six-vertex is specified by the condition that the number of arrows entering and leaving each vertex coincide. The types of possible vertices are shown on Fig. 1. With each configuration we associate a statistical weights. Generically the last ones can depend on the position of the vertex. We consider the model, which is symmetric with respect to simultaneous reversal of all arrows. Thus, we have $`3N^2`$ statistical weights $`w_{jk}=\{a_{jk},b_{jk},c_{jk}\}`$, where indices $`j,k=0,1,\mathrm{},N1`$ numerate vertices. The domain wall boundary conditions correspond to the entering arrows on the lower and upper boundaries, and leaving arrows on the left and right ones. The partition function is given by $$Z_N=\underset{\{C\}}{}\underset{j,k=0}{\overset{N1}{}}w_{jk}.$$ (2.1) Here the sum is taken with respect to all possible configurations $`C`$ on the whole lattice. The six-vertex model is exactly solvable, if $`3N^2`$ weights $`\{a_{jk},b_{jk},c_{jk}\}`$ are parameterized by $`2N+1`$ variables $`\lambda _j`$, $`\xi _j`$, $`(j=0,1,\mathrm{},N1)`$ and $`\eta `$: $$a_{jk}=\phi (\lambda _j\xi _k+\eta ),b_{jk}=\phi (\lambda _j\xi _k),c_{jk}=\phi (\eta ).$$ (2.2) Parameters $`\lambda _j`$ are associated with horizontal lines, $`\xi _j`$—with vertical ones (see. Fig. 2), $`\eta `$ is an arbitrary complex constant. Function $`\phi (x)`$ is defined up to normalization factor and is equal to $$\phi (x)=\mathrm{sinh}(x),\text{or}\phi (x)=x.$$ (2.3) In the first case the model is equivalent to the quantum $`XXZ`$ magnet, in the second—to the $`XXX`$ magnet. The Izergin’s formula for the partition function on the lattice of finite volume with the domain wall boundary conditions has the form $$Z_N=\frac{\underset{a,b=0}{\overset{N1}{}}\phi (\lambda _a\xi _b+\eta )\phi (\lambda _a\xi _b)}{\underset{N>a>b0}{}\phi (\lambda _a\lambda _b)\phi (\xi _b\xi _a)}det\left(\frac{\phi (\eta )}{\phi (\lambda _j\xi _k)\phi (\lambda _j\xi _k+\eta )}\right).$$ (2.4) We would like to draw attention of the reader that here and further we numerate indices of matrices starting with zero, but not with one: $`j,k=0,1,\mathrm{},N1`$. We are interesting in homogeneous limit, when the statistical weights do not depend on the position of the vertex. In other words $`\lambda _j=\lambda ,\xi _j=\xi ,j=0,1\mathrm{},N1`$. The corresponding limit in (2.4) can be obtained via $$\underset{\genfrac{}{}{0pt}{}{\lambda _j\lambda }{\xi _k\xi }}{lim}\frac{det[\mathrm{\Phi }(\lambda _j,\xi _k)]}{\mathrm{\Delta }(\lambda )\mathrm{\Delta }(\xi )}=det\left[\frac{1}{j!k!}\frac{^j}{\lambda ^j}\frac{^k}{\xi ^k}\mathrm{\Phi }(\lambda ,\xi )\right].$$ (2.5) Here $`\mathrm{\Phi }(\lambda ,\xi )`$ is some enough smooth two-variable function, $`\mathrm{\Delta }(\lambda )`$ and $`\mathrm{\Delta }(\xi )`$ are Van der Monde determinants of variables $`\{\lambda \}`$ and $`\{\xi \}`$ respectively. In the rational case, when $`\phi (x)=x`$ the direct use of (2.5) gives $$Z_N^{(XXX)}=\frac{\left[\nu (\nu +\eta )\right]^{N^2}}{\underset{m=0}{\overset{N1}{}}(m!)^2}det\left[(j+k)!\left(\nu ^{jk1}(\nu +\eta )^{jk1}\right)\right],$$ (2.6) where $`\nu =\lambda \xi `$. In the trigonometric case it convenient first to use variables $`u_j=e^{2\lambda _j}`$, $`v_j=e^{2\xi _j}`$ and $`q=e^\eta `$. Then (2.4) takes the form $$Z_N^{(XXZ)}=\underset{a=0}{\overset{N1}{}}\left(e^{\lambda _a\xi _a}\right)\frac{\underset{a,b=0}{\overset{N1}{}}\mathrm{sinh}(\lambda _a\xi _b+\eta )(u_av_b)}{\underset{N>a>b0}{}(u_au_b)(v_bv_a)}det\left(\frac{1}{u_jv_k}\frac{q}{qu_jq^1v_k}\right).$$ (2.7) Now the equation (2.5) can be easily applied, and after simple algebra we obtain $$Z_N^{(XXZ)}=e^{N\nu }\frac{\left[\mathrm{sinh}(\nu )\mathrm{sinh}(\nu +\eta )\right]^{N^2}}{\underset{m=0}{\overset{N1}{}}(m!)^2}det\left[\frac{(j+k)!}{\mathrm{sinh}^{j+k+1}(\nu )}\frac{q^{jk+1}(j+k)!}{\mathrm{sinh}^{j+k+1}(\nu +\eta )}\right].$$ (2.8) Here as in (2.6) $`\nu =\lambda \xi `$, and $`q=e^\eta `$. Below we shall focus on the trigonometric case (2.8), as the most general one. As usual, the answer for the partition function $`Z_N^{(XXX)}`$ can be obtained in the limit $`\nu =ϵ\nu ,\eta =ϵ\eta ,ϵ0`$. Our goal is to transform (2.8) to the expression, suitable for the asymptotic analysis at $`N\mathrm{}`$. There are factors with trivial asymptotics in the equation (2.8): $`e^{N\nu }`$ and $`\left[\mathrm{sinh}(\nu )\mathrm{sinh}(\nu +\eta )\right]^{N^2}`$. It is also easy to evaluate the behavior of the factor $`_{m=0}^{N1}(m!)^2`$ at $`N\mathrm{}`$. The main problem is the determinant. In the next section we consider some special transforms of (2.8), which result finally to the new representation for the partition function, containing the determinant of an integral operator. ## 3 Fredholm determinant As the first step, it is convenient to extract, for example, $`\mathrm{sinh}^{jk1}(\nu )`$ out off the determinant (2.8). Then (2.8) becomes $$Z_N^{(XXZ)}=e^{N\nu }\frac{\mathrm{sinh}^{N^2}(\nu +\eta )}{\underset{m=0}{\overset{N1}{}}(m!)^2}det[AqtQTATQ^1].$$ (3.1) Here $$A_{jk}=(j+k)!,T_{jk}=\delta _{jk}t^j,Q_{jk}=\delta _{jk}q^j,t=\frac{\mathrm{sinh}(\nu )}{\mathrm{sinh}(\nu +\eta )}.$$ (3.2) Note. Similarly we could extract terms $`\mathrm{sinh}^{jk1}(\nu +\eta )`$. Then, up to the replacements $`tt^1`$, $`qq^1`$ and trivial common factor, the arising expression coincides with (3.1). Now we can transform (3.1) as follows: $$det[AqtQTATQ^1]=detAdet[IU],\text{where}U=qtA^1QTATQ^1.$$ (3.3) The advantage of this method is that $`detA`$ cancels the product of factorials in the denominator, since $$detA=\underset{m=0}{\overset{N}{}}(m!)^2,$$ (3.4) (see Appendix 1). Thus our problem reduces to the derivation of the inverse matrix $`A^1`$ and calculation of the product $`U=qtA^1QTATQ^1`$. The inverse matrix $`A^1`$ can be easily written in terms of Laguerre polynomials, therefore below we recall their basic properties. The Laguerre polynomials $`L_n(x)`$ are defined by $$L_n(x)=\frac{e^x}{n!}\frac{d^n}{dx^n}(x^ne^x)$$ (3.5) and generate the orthogonal system on the half-axis $`R_+`$ with the weight $`e^x`$: $$_0^{\mathrm{}}L_n(x)L_m(x)e^x𝑑x=\delta _{nm}.$$ (3.6) The key object, determining the matrix $`A^1`$, is the kernel $$K_n(x,y)=\underset{k=0}{\overset{n}{}}L_k(x)L_k(y).$$ (3.7) One more useful representation for $`K_n(x,y)`$ follows from Christoffel–Darboux formula: $$K_n(x,y)=\frac{n+1}{xy}\left(L_{n+1}(x)L_n(y)L_{n+1}(y)L_n(x)\right).$$ (3.8) Obviously, for arbitrary polynomial $`\pi _m(x)`$ of degree $`m`$, less or equal to $`n`$, holds $$_0^{\mathrm{}}\pi _m(x)K_n(x,y)e^x𝑑x=\pi _m(y),0mn.$$ (3.9) To prove (3.9), it is enough to substitute to the l.h.s. the expansion of $`\pi _m(x)`$ with respect to Laguerre polynomials $`\pi _m(x)=_{k=0}^mc_kL_k(x)`$, and to use the definition (3.7). Being a polynomial of variables $`x`$ and $`y`$, the kernel $`K_n(x,y)`$ can be presented in the form $$K_n(x,y)=\underset{j,k=0}{\overset{n}{}}K_{jk}^{(n)}x^jy^k,$$ (3.10) where $$K_{jk}^{(n)}=\frac{1}{j!k!}\frac{^{j+k}}{x^jy^k}K_n(x,y)\|_{x=y=0}.$$ (3.11) Now we are in position to formulate the theorem on the inverse matrix. ###### Theorem 1 Let $`A_{jk}=(j+k)!`$$`j,k=0,1\mathrm{},N1`$. Then the inverse matrix has the entries $$(A^1)_{jk}=K_{jk}^{(N1)}.$$ (3.12) Proof. Using the integral representation for $`(j+k)!`$, we have $$\underset{l=0}{\overset{N1}{}}A_{jl}K_{lk}^{(N1)}=\frac{1}{k!}\frac{^k}{y^k}\underset{l=0}{\overset{N1}{}}\frac{1}{l!}\frac{^l}{x^l}_0^{\mathrm{}}s^{j+l}e^sK_{N1}(x,y)𝑑s\|_{x=y=0}.$$ (3.13) The sum with respect to $`l`$ is the Taylor series (in fact, since $`K_{N1}(x,y)`$ is $`(N1)`$-degree polynomial of $`x`$, the mentioned series turns to the finite sum). Hence, $$\underset{l=0}{\overset{N1}{}}A_{jl}K_{lk}^{(N1)}=\frac{1}{k!}\frac{^k}{y^k}_0^{\mathrm{}}s^je^sK_{N1}(s,y)𝑑s\|_{y=0}.$$ (3.14) Due to (3.9) the integral in (3.14) is equal to $`y^j`$. Thus, we arrive at $$\underset{l=0}{\overset{N1}{}}A_{jl}K_{lk}^{(N1)}=\frac{1}{k!}\frac{^k}{y^k}y^j\|_{y=0}=\delta _{jk}.$$ (3.15) The theorem is proved. Note. For generalization of the relationship between Hankel matrices and Christoffel–Darboux kernels see . The obtained explicit expression for $`A^1`$ allows us to compute the matrix $`U`$ (3.3). Indeed, $$U_{jk}=\frac{1}{k!}\frac{^k}{y^k}\underset{l=0}{\overset{N1}{}}t^{l+k+1}q^{lk+1}\frac{1}{l!}\frac{^l}{x^l}_0^{\mathrm{}}s^{j+l}e^sK_{N1}(x,y)𝑑s\|_{x=y=0}.$$ (3.16) Similarly to the proof of the Theorem the sum with respect to $`l`$ easily can be computed and it gives $$U_{jk}=\frac{qt}{k!}\frac{^k}{y^k}_0^{\mathrm{}}x^j(q^1t)^ke^xK_{N1}(xqt,y)𝑑x\|_{y=0}.$$ (3.17) Thus, the partition function $`Z_N^{(XXZ)}`$ turns to be proportional to the determinant of the matrix $`\delta _{jk}U_{jk}`$ (3.17). To compute the last one it is convenient to use the equation $$det(IU)=\mathrm{exp}\{tr\mathrm{log}(IU)\}=\mathrm{exp}\left\{\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n}trU^n\right\}.$$ (3.18) Hence, we need to find the traces of powers of the matrix $`U`$. It is easy to see that $$(U^n)_{jk}=\frac{(qt)^n}{k!}\frac{^k}{y^k}_0^{\mathrm{}}x_1^j(q^1t)^kK_{N1}(x_1qt,x_2q^1t)\mathrm{}K_{N1}(x_nqt,y)e^{_{m=1}^nx_m}d^nx\|_{y=0}.$$ (3.19) Indeed, for $`n=1`$ (3.19) coincides with (3.17). Assuming that (3.19) holds for $`U^n`$, we immediately obtain $`{\displaystyle \underset{l=0}{\overset{N1}{}}}(U^n)_{jl}U_{lk}={\displaystyle \frac{(qt)^{n+1}}{k!}}{\displaystyle \frac{^k}{y^k}}{\displaystyle _0^{\mathrm{}}}x_1^j(q^1t)^kK_{N1}(x_1qt,x_2q^1t)\mathrm{}`$ $`\mathrm{}K_{N1}(x_nqt,x_{n+1}q^1t)K_{N1}(x_{n+1}qt,y)e^{_{m=1}^{n+1}x_m}d^{n+1}x\|_{y=0}.`$ (3.20) (Here the sum with respect to $`l`$ again turns to the Taylor series). Hence, $$tr(U^n)=(qt)^n_0^{\mathrm{}}K_{N1}(x_1qt,x_2q^1t)\mathrm{}K_{N1}(x_nqt,x_1q^1t)e^{_{m=1}^nx_m}d^nx.$$ (3.21) The obtained expression (3.21) for the trace of $`U^n`$ allows to replace $`det(\delta _{jk}U_{jk})`$ with Fredholm determinant. In order to do this, we introduce integral operator $`IV`$, acting on $`R_+`$ $$[(IV)f](x)=f(x)_0^{\mathrm{}}V(x,y)f(y)𝑑y,$$ (3.22) with the kernel $$V(x,y)=qtK_{N1}(xqt,yq^1t)e^{(x+y)/2}=qt\underset{k=0}{\overset{N1}{}}L_k(xqt)L_k(yq^1t)e^{(x+y)/2}.$$ (3.23) Due to (3.8), $`V(x,y)`$ also can be presented in the form $$V(x,y)=\frac{qN}{qxq^1y}\left(L_{N1}(xqt)L_N(yq^1t)L_N(xqt)L_{N1}(yq^1t)\right)e^{(x+y)/2}.$$ (3.24) Then due to (3.21) $$trU^n=trV^n,$$ (3.25) and we finally obtain $$det(\delta _{jk}U_{jk})=det(IV).$$ (3.26) ## 4 Discussions Thus, we have found new representation for the partition function of the six-vertex model, containing the Fredholm determinant of the integral operator with the kernel (3.23), (3.24): $$Z_N^{(XXZ)}=e^{N\nu }\left[\mathrm{sinh}(\nu +\eta )\right]^{N^2}det(IV).$$ (4.1) Since the kernel $`V`$ is degenerated, we always can turn back to the determinant of the finite size matrix $$det\left(Iqt\underset{k=0}{\overset{N1}{}}L_k(xqt)L_k(yq^1t)e^{(x+y)/2}\right)=det\left(\delta _{jk}qt_0^{\mathrm{}}L_j(xqt)L_k(xq^1t)e^x𝑑x\right).$$ (4.2) However, for the asymptotic analysis at $`N\mathrm{}`$ the Fredholm determinant representation seems to be preferable due to several reasons. First, the equation (3.9) shows that the kernel $`K_N(x,y)`$ acts as identity operator on the subspace of the polynomials of $`N`$-th degree. Then one can think that for $`N`$ going to infinity, $`K_N(x,y)`$ goes to identity operator (delta-function). Of course, this limiting procedure requires a serious analysis of convergency. It is remarkable however, that for certain values of $`q`$ and $`t`$ such a naive interpretation of the kernel (3.23) does holds, and the Fredholm determinant can be evaluated explicitly in the limit $`N\mathrm{}`$. Second, similar integral operators with the kernels, depending on Laguerre polynomials, often appear in the theories of random matrices and random permutations . Some methods of the analysis of such operators can be directly applied to the case under consideration. Besides, at large $`N`$ the Laguerre polynomials can be approximated by Bessel functions, what also leads us to the integral operator, analogous to ones, considered in the papers . Third, in particular case of the partition function $`Z_N^{(XXX)}`$, when $`q=1`$, the kernel $`V(x,y)`$ takes the form $$V(x,y)=\frac{N}{xy}\left(L_{N1}(xt)L_N(yt)L_N(xt)L_{N1}(yt)\right)e^{(x+y)/2}.$$ (4.3) The integral operator with the kernel (4.3) belongs to the class of integrable integral operators . The corresponding Fredholm determinant turn to be $`\tau `$-function of classical exactly solvable equation and can be evaluated via the matrix Riemann–Hilbert problem methods. We are planning to present the detailed asymptotic analysis of the partition function of the six-vertex model with domain wall boundary conditions in forthcoming publications. The author thanks A. R. Its for numerous and useful discussions. The work was supported in parts by RFBR, Grant 99-01-00151 and INTAS-99-1782. ## Appendix A Calculation of the determinant One of the method to proof (3.4) is to use well known formula for Cauchy determinant $$\frac{\underset{N>a>b0}{}(\alpha _a\alpha _b)(\beta _b\beta _a)}{\underset{a,b=0}{\overset{N1}{}}(\alpha _a\beta _b)}=det\left(\frac{1}{\alpha _j\beta _k}\right).$$ (A.1) Here $`\{\alpha \}`$ and $`\{\beta \}`$ are arbitrary complex. Dividing both sides of (A.1) by Van der Monde determinants $`\mathrm{\Delta }(\alpha )`$ and $`\mathrm{\Delta }(\beta )`$ and taking homogeneous limit $`\alpha _j\alpha ,\beta _j\beta ,j=0,1,\mathrm{},N1`$, we obtain due to (2.5) $$(\alpha \beta )^{N^2}=\underset{m=0}{\overset{N1}{}}(m!)^2det\left(\frac{(j+k)!}{(\alpha \beta )^{j+k+1}}\right),$$ (A.2) what immediately implies (3.4). Another method to compute $`detA`$ is based on the application of Laguerre polynomials. Below we present this method, since it allows to find $`A^1`$ as well as to prove (3.4). Using the integral representation for factorial, we have $$detA=det\left(_0^{\mathrm{}}e^ss^{j+k}𝑑s\right).$$ (A.3) Denote integration variable in the $`j`$-th row as $`s_j`$. Then we come to the multiple integral $$detA=_0^{\mathrm{}}𝑑s_0\mathrm{}𝑑s_{N1}det\left(e^{s_j}s_j^{j+k}\right)=_0^{\mathrm{}}𝑑s_0\mathrm{}𝑑s_{N1}\underset{j=0}{\overset{N1}{}}\left(e^{s_j}s^j\right)\mathrm{\Delta }(s),$$ (A.4) where $`\mathrm{\Delta }(s)`$ is Van der Monde determinant. Now we can use well known identity $$\mathrm{\Delta }(s)=det[p_k(s_j)],$$ (A.5) where $`p_k(s)`$ is a system of arbitrary polynomials with the highest coefficient equal to $`1`$: $$p_k(s)=s^k+o(s).$$ (A.6) In particular we can choose the Laguerre polynomials $$p_k(s)\stackrel{~}{L}_k(s)=(1)^ke^s\frac{d^k}{ds^k}(s^ke^s).$$ (A.7) The normalization in (A.7) differs from one in (3.5) by factor $`(1)^kk!`$ in order to provide (A.6). Therefore we denote the polynomials (A.7) as $`\stackrel{~}{L}_k`$. The orthogonality condition then takes the form $$_0^{\mathrm{}}\stackrel{~}{L}_m(s)\stackrel{~}{L}_n(s)e^s𝑑s=(n!)^2\delta _{nm}.$$ (A.8) Thus, the equation (A.4) becomes $$detA=_0^{\mathrm{}}𝑑s_0\mathrm{}𝑑s_{N1}\underset{j=0}{\overset{N1}{}}\left(e^{s_j}s^j\right)det[\stackrel{~}{L}_k(s_j)].$$ (A.9) Integrating inside the determinant, we obtain $$detA=det\left(_0^{\mathrm{}}e^ss^j\stackrel{~}{L}_k(s)𝑑s\right).$$ (A.10) Due to orthogonality of the polynomials $`\stackrel{~}{L}_k(s)`$, the entries in (A.10) vanish, if $`j<k`$. Thus we deal with the determinant of triangular matrix and, hence, $$detA=\underset{m=0}{\overset{N1}{}}\left(_0^{\mathrm{}}e^ss^m\stackrel{~}{L}_m(s)𝑑s\right).$$ (A.11) It remains to observe that expansion of $`s^m`$ with respect to $`\stackrel{~}{L}_m(s)`$ has the form $$s^m=\stackrel{~}{L}_m(s)+\text{polynomials of lower degree}$$ therefore due to (A.8) $$detA=\underset{m=0}{\overset{N1}{}}(m!)^2,$$ (A.12) which ends the proof. The described method allows to compute not only $`detA`$, but the minors of the size $`N1`$ as well. This provides the possibility to find $`A^1`$.
warning/0005/gr-qc0005027.html	ar5iv	text	# Relative velocity and relative acceleration induced by the torsion in (𝐿_𝑛,𝑔)- and 𝑈_𝑛-spaces ## 1 Introduction ### 1.1 Differential geometry and space-time geometry In the last years, the evolution of the relations between differential geometry and space-time geometry has made some important steps toward applications of more comprehensive differential-geometric structures in the models of space-time than these used in (pseudo) Riemannian spaces without torsion ($`V_n`$-spaces). 1. It has been proved recently that every differentiable manifold with one affine connection and metrics \[$`(L_n,g)`$-space\] could be used as a model for a space-time. In it the equivalence principle (related to the vanishing of the components of an affine connection at a point or on a curve in a manifold) holds $`÷`$ , . Even if the manifold has two different (not only by sign) connections for tangent and co-tangent vector fields \[$`(\overline{L}_n,g)`$-space\] the principle of equivalence is fulfilled at least for one of the two types of vector fields . On this grounds, every free moving spinless test particle in a suitable basic system (frame of reference) will fulfil an equation identical with the equation for a free moving spinless test particle in the Newtons mechanics or in the special relativity. In $`(\overline{L}_n,g)`$\- and $`(L_n,g)`$-spaces, this equation could be the autoparallel equation \[different from the geodesic equation in contrast to the case of (pseudo) Riemannian spaces without torsion ($`V_n`$-spaces)\]. 2. There are evidences that $`(L_n,g)`$\- and $`(\overline{L}_n,g)`$-spaces can have similar structures as the $`V_n`$-spaces for describing dynamical systems and the gravitational interaction. In such type of spaces one could use Fermi-Walker transports conformal transports and different frames of reference . All these notions appear as generalizations of the corresponding notions in $`V_n`$-spaces. For instance, in $`(L_n,g)`$\- and $`(\overline{L}_n,g)`$-spaces a proper non-rotating accelerated observer’s frame of reference could be introduced by analogy of the same type of frame of reference related to a Fermi-Walker transport in the Einstein theory of gravitation (ETG). 3. All kinematic characteristics related to the notion of relative velocity as well as the kinematic characteristics related to the notion of relative acceleration have been worked out for $`(L_n,g)`$\- and $`(\overline{L}_n,g)`$-spaces without changing their physical interpretation in $`V_n`$-spaces . Necessary and sufficient conditions as well as only necessary an only sufficient conditions for vanishing shear, rotation and expansion acceleration induced by the curvature are found . The last results are related to the possibility of building a theoretical basis for construction of gravitational wave detectors on the grounds of gravitational theories over $`(L_n,g)`$ and $`(\overline{L}_n,g)`$–spaces. Usually, in the gravitational experiments the measurements of two basic objects are considered : (a) The relative velocity between two particles (or points) related to the rate of change of the length (distance) between them. The change of the distance is supposed to be caused by the gravitational interaction. (b) The relative accelerations between two test particles (or points) of a continuous media. These accelerations are related to the curvature of the space-time and supposed to be induced by a gravitational force. Together with the accelerations induced by the curvature in $`V_n`$-spaces accelerations induced by the torsion would appear in $`U_n`$-spaces, as well as by torsion and by non-metricity in $`(L_n,g)`$\- and $`(\overline{L}_n,g)`$-spaces. In particular, in other models of a space-time \[different from the (pseudo) Riemannian spaces without torsion\] the torsion has to be taken into account if we consider the characteristics of the space-time. 4. On the one hand, until now, there are a few facts that the torsion could induce some very small and unmesasurable effects in quantum mechanical systems considered in spaces with torsion . At the same time, there are no evidences that the model’s descriptions of interactions on macro level should include the torsion as a necessary mathematical tool. From this (physical) point of view the influence of the torsion on dynamical systems could be ignored since it could not play an important role in the description of physical systems in the theoretical gravitational physics. On the other hand, from mathematical point of view (as we will try to show in this paper), the role of the torsion in new theories for description of dynamical systems could be important and not ignorable. ### 1.2 Problems and results In this paper the influence of the torsion on the relative velocity and the relative acceleration between particles moving in $`(L_n,g)`$\- and $`U_n`$-spaces ($`n=4`$) is considered. In Sec. 2. the notion of relative velocity as well as the related to it notions of shear, rotation and expansion velocities, induced by the torsion, are determined. The change of the length of a vector field (deviation vector field) related to two moving particles is given in an explicit form for $`(L_n,g)`$\- and $`U_n`$-spaces. It is shown that the existence of deformation, shear, rotation and expansion (velocities), induced by the torsion, can compensate (under given conditions) the action of these induced by external forces. The necessary and sufficient conditions as well as only necessary and only sufficient conditions for vanishing deformation, shear, rotation and expansion are found. The vanishing of all these quantities could be caused by the torsion. In Sec. 3. the notion of relative acceleration and the related to it notions of shear, rotation and expansion accelerations induced by the torsion are found. It is shown that the existence of kinematic characteristics, induced by the torsion, can compensate (under given conditions) the action of these induced by the curvature or by external forces. The change of the rate of change of the length of a deviation vector field is given in explicit form for $`(L_n,g)`$\- and $`U_n`$-spaces. Concluding remarks comprise the final Sec. 4. The most considerations are given in details (even in full details) for those readers who are not familiar with the considered problems. ## 2 Relative velocity and its kinematic characteristics induced by the torsion Let us now recall some well known facts from differential geometry. Every contravariant vector field $`\xi T(M)`$ over a differentiable manifold $`M`$ can be written by means of its projection along and orthogonal to a contravariant non-null (nonisotropic) vector field $`u`$ in two parts - one collinear to $`u`$ and one - orthogonal to $`u`$, i.e. $$\xi =\frac{l}{e}.u+h^u[g(\xi )]\text{ }=\frac{l}{e}.u+\overline{g}[h_u(\xi )]\text{ ,}$$ (1) where $`\overline{g}[h_u(\xi )]:=\xi _{}=g^{ik}h_{kl}\xi ^le_i`$ is called deviation vector field and $$\begin{array}{c}l=g(\xi ,u)=g_{ij}\xi ^iu^j\text{ , }\\ \text{ }h^u=\overline{g}\frac{1}{e}.uu\text{ , }\xi =\xi ^i._i=\xi ^k.e_k\text{ , }h^u=h^{ij}.e_i.e_j\text{ ,}\\ \overline{g}(h_u)\overline{g}=h^u\text{ , }h_u(\overline{g})(g)=h_u\text{ , }h^u(g)(\overline{g})=h^u\text{ , }g(h^u)g=h_u\text{ , }\\ h_u=g\frac{1}{e}g(u)g(u)\text{ ,}\\ \text{ }\overline{g}=g^{ij}e_i.e_j\text{}e_i.e_j=\frac{1}{2}(e_ie_j+e_je_i)\text{ , }g^{ij}=g^{ji}\text{ ,}\\ g=g_{ij}e^i.e^j\text{ , }e^i.e^j=\frac{1}{2}(e^ie^j+e^je^i)\text{ , }g_{ij}=g_{ji}\text{ .}\end{array}$$ (2) Therefore, the covariant derivative $`_u\xi `$ of the contravariant vector field $`\xi `$ along a non-null vector field $`u`$ \[as a result of the action of the contravariant differential operator $`_u:\xi _u\xi T(M)`$, $`_u\xi =\xi _{;j}^iu^j_i`$, $`\xi ^i{}_{;j}{}^{}=\xi ^i{}_{,j}{}^{}+\mathrm{\Gamma }_{lj}^i\xi ^l`$, $`\mathrm{\Gamma }_{lj}^i`$ are the components of the affine connection $`\mathrm{\Gamma }`$ in a $`(L_n,g)`$-space\] can be written in the form $$_u\xi =\frac{\overline{l}}{e}u+\overline{g}[h_u(_u\xi )]\text{ , }\overline{l}=g(_u\xi ,u)\text{ .}$$ (3) ### 2.1 Relative velocity in $`(L_n,g)`$-spaces The notion relative velocity vector field (relative velocity) $`{}_{rel}{}^{}v`$ can be defined in $`(L_n,g)`$-spaces (regardless of its physical interpretation) as the projection \[orthogonal to a non-null (nonisotropic) contravariant vector field $`u`$\] of the first covariant derivative $`_u\xi `$ (along the same non-null vector field $`u`$) of (another) contravariant vector field $`\xi `$, i.e. $${}_{rel}{}^{}v=\overline{g}(h_u(_u\xi ))=g^{ij}h_{jk}\xi ^k\text{ }_{;l}u^le_i\text{,}e_i=_i\text{ (in a co-ordinate basis),}$$ (4) where (the indices in a co-ordinate and in a non-co-ordinate basis are written in both cases as Latin indices instead of Latin and Greek indices) $$\begin{array}{c}h_u=g\frac{1}{e}.g(u)g(u)\text{ },\text{ }h_u=h_{ij}e^i.e^j\text{ , }\overline{g}=g^{ij}e_i.e_j,\\ _u\xi =\xi ^i\text{ }_{;j}u^je_i\text{ , }\xi ^i\text{ }_{;j}=e_j\xi ^i+\mathrm{\Gamma }_{kj}^i\xi ^k\text{}\mathrm{\Gamma }_{kj}^i\mathrm{\Gamma }_{jk}^i\text{ },\\ e=g(u,u)=g_{ij}u^iu^j=u_iu^i0\text{ , }g(u)=g_{ik}u^ke^i=u_ie^i\text{ ,}\\ \text{ }h_u(_u\xi )=h_{ij}\xi ^j\text{ }_{;k}u^ke^i\text{ , }h_{ij}=g_{ij}\frac{1}{e}u_iu_j\text{ ,}\\ h_u(_u\xi )=h_{ij}\xi ^j\text{ }_{;k}u^ke^i\text{ .}\end{array}$$ (5) In a co-ordinate basis $`e_j\xi ^i=\xi ^i`$ $`{}_{,j}{}^{}=_j\xi ^i=\xi ^i/x^j`$, $`e^j=dx^j`$, $`e_i=_i=/x^i`$,$`u=u^i_i`$, $`e_{,k}=e_ke=_ke`$. Using the relation between the Lie derivative $`\mathrm{\pounds }_\xi u`$ and the covariant derivative $`_\xi u`$ $$\mathrm{\pounds }_\xi u=_\xi u_u\xi T(\xi ,u)\text{ , }T(\xi ,u)=T_{ij}^k.\xi ^i.u^j.e_k\text{ ,}$$ (6) one can write $`_u\xi `$ in the form $$_u\xi =(k)g(\xi )\mathrm{\pounds }_\xi u\text{ }=k[g(\xi )]\mathrm{\pounds }_\xi u\text{,}$$ (7) or, taking into account the above expression for $`\xi `$, in the form $`_u\xi =k[h_u(\xi )]+\frac{l}{e}.a\mathrm{\pounds }_\xi u`$, where $$\begin{array}{c}k[g(\xi )]=_\xi uT(\xi ,u)\text{ , }k=(u^i\text{ }_{;l}T_{lk}^i.u^k).g^{lj}.e_ie_j=k^{ij}.e_ie_j\text{ ,}\\ k[g(u)]=k(g)u=k^{ij}.g_{jk}.u^k.e_i=a=_uu=u^i\text{ }_{;j}.u^j.e_i\text{ .}\end{array}$$ (8) $`T_{ij}^k`$ are the components of the torsion tensor $`T`$: $$\begin{array}{c}T_{ij}^k=T_{ji}^k=\mathrm{\Gamma }_{ji}^k\mathrm{\Gamma }_{ij}^kC_{ij}\text{ }^k\text{ (in a non-co-ordinate basis }\{e_k\}\text{) , }\\ [e_i,e_j]=\mathrm{\pounds }_{e_i}e_j\text{ }=C_{ij}\text{ }^k.e_k\text{ ,}\\ \text{ }T_{ij}^k=\mathrm{\Gamma }_{ji}^k\mathrm{\Gamma }_{ij}^k\text{ (in a co-ordinate basis }\{_k\}\text{ ) ,}\end{array}$$ For $`h_u(_u\xi )`$, it follows $$h_u(_u\xi )=h_u(\frac{l}{e}a\mathrm{\pounds }_\xi u)+h_u(k)h_u(\xi )\text{ ,}$$ (9) where $`h_u(k)h_u(\xi )=h_{ik}k^{kl}h_{lj}\xi ^je^i`$, $`h_u(u)=0`$,$`u(h_u)=0`$, $`h_u(k)h_u(u)=0`$, $`(u)h_u(k)h_u=0`$. If we introduce the abbreviation $$d=h_u(k)h_u=h_{ik}k^{kl}h_{lj}e^ie^j=d_{ij}e^ie^j\text{ ,}$$ (10) the expression for $`{}_{rel}{}^{}v`$ can take the form $${}_{rel}{}^{}v=\overline{g}[h_u(_u\xi )]=\overline{g}(h_u)(\frac{l}{e}a\mathrm{\pounds }_\xi u)+\overline{g}[d(\xi )]=$$ $$=[g^{ik}h_{kl}(\frac{l}{e}a^l\mathrm{\pounds }_\xi u^l)+g^{ik}d_{kl}\xi ^l]e_i\text{ }=_{rel}v^ie_i\text{ ,}$$ (11) or $$g(_{rel}v)=h_u(_u\xi )=h_u(\frac{l}{e}a\mathrm{\pounds }_\xi u)+d(\xi )\text{ .}$$ (12) For the special case, when the vector field $`\xi `$ is orthogonal to $`u`$, i.e. $`\xi =\overline{g}[h_u(\xi )]`$, and the Lie derivative of $`u`$ along $`\xi `$ is zero, i.e. $`\mathrm{\pounds }_\xi u=0`$, then the relative velocity can be written in the form $`g(_{rel}v)=d(\xi )`$ or in the form $${}_{rel}{}^{}v=\overline{g}[d(\xi )]\text{.}$$ Remark. All further calculations leading to a useful representation of $`d`$ are quite straightforward. The problem here was the finding out a representation of $`h_u(_u\xi )`$ in the form (9) which is not a trivial task. ### 2.2 Deformation velocity, shear velocity, rotation velocity and expansion velocity The covariant tensor field $`d`$ is a generalization for $`(L_n,g)`$-spaces of the well known deformation velocity tensor for $`V_n`$-spaces , . It is usually represented by means of its three parts: the trace-free symmetric part, called shear velocity tensor (shear), the anti-symmetric part, called rotation velocity tensor (rotation) and the trace part, in which the trace is called expansion velocity (expansion) invariant. After some more complicated as for $`V_n`$-spaces calculations, the deformation velocity tensor $`d`$ can be given in the form $$d=h_u(k)h_u=h_u(k_s)h_u+h_u(k_a)h_u=\sigma +\omega +\frac{1}{n1}\theta h_u\text{ ,}$$ (13) where $`k_s=_sk^{ij}e_i.e_j`$, $`{}_{s}{}^{}k_{}^{ij}=\frac{1}{2}(k^{ij}+k^{ji})`$, $`{}_{a}{}^{}k=_ak^{ij}e_ie_j`$, $`{}_{a}{}^{}k_{}^{ij}=\frac{1}{2}(k^{ij}k^{ji})`$, $`e_ie_j=\frac{1}{2}(e_ie_je_je_i)`$ The tensor $`\sigma `$ is the shear velocity tensor (shear) , $$\begin{array}{c}\sigma =_sE_sP=EP\frac{1}{n1}\overline{g}[EP]h_u=\sigma _{ij}e^i.e^j=\\ =EP\frac{1}{n1}(\theta _o\theta _1)h_u\text{ ,}\end{array}$$ (14) where $$\begin{array}{c}{}_{s}{}^{}E=E\frac{1}{n1}\overline{g}[E]h_u\text{ , }\overline{g}[E]=g^{ij}E_{ij}=\theta _o\text{ ,}\\ E=h_u(\epsilon )h_u\text{ , }k_s=\epsilon m\text{ , }\epsilon =\frac{1}{2}(u_{\text{ };l}^ig^{lj}+u_{\text{ };l}^jg^{li})e_i.e_j\text{ ,}\\ m=\frac{1}{2}(T_{lk}^iu^kg^{lj}+T_{lk}^ju^kg^{li})e_i.e_j\text{ .}\end{array}$$ (15) The tensor $`{}_{s}{}^{}E`$ is the torsion-free shear velocity tensor, $`{}_{s}{}^{}P`$ is the shear velocity tensor induced by the torsion, $$\begin{array}{c}{}_{s}{}^{}P=P\frac{1}{n1}\overline{g}[P]h_u\text{ , }\overline{g}[P]=g^{kl}P_{kl}=\theta _1\text{}P=h_u(m)h_u\text{ ,}\\ \theta _1=T_{kl}^ku^l\text{ , }\theta _o=u^n\text{ }_{;n}\frac{1}{2e}(e_{,k}u^kg_{kl;m}u^mu^ku^l)\text{ , }\theta =\theta _o\theta _1\text{ . }\end{array}$$ (16) The invariant $`\theta `$ is the expansion velocity, $`\theta _o`$ is the torsion-free expansion velocity, $`\theta _1`$ is the expansion velocity induced by the torsion. The tensor $`\omega `$ is the rotation velocity tensor (rotation velocity), $$\begin{array}{c}\omega =h_u(k_a)h_u=h_u(s)h_uh_u(q)h_u=SQ\text{ ,}\\ s=\frac{1}{2}(u^k\text{ }_{;m}g^{ml}u^l\text{ }_{;m}g^{mk})e_ke_l\text{ ,}\\ q=\frac{1}{2}(T_{mn}^kg^{ml}T_{mn}^lg^{mk})u^ne_ke_l\text{ , }S=h_u(s)h_u\text{ , }Q=h_u(q)h_u\text{ .}\end{array}$$ (17) The tensor $`S`$ is the torsion-free rotation velocity tensor, $`Q`$ is the rotation velocity tensor induced by the torsion. By means of the expressions for $`\sigma `$, $`\omega `$ and $`\theta `$ the deformation velocity tensor $`d`$ can be decomposed in two parts: $`d_0`$ and $`d_1`$ $$d=d_od_1\text{ , }d_o=_sE+S+\frac{1}{n1}.\theta _o.h_u\text{ , }d_1=_sP+Q+\frac{1}{n1}.\theta _1.h_u\text{ ,}$$ (18) where $`d_o`$ is the torsion-free deformation velocity tensor and $`d_1`$ is the deformation velocity tensor induced by the torsion. For the case of $`V_n`$-spaces $`d_1=0`$ ($`{}_{s}{}^{}P=0`$, $`Q=0`$, $`\theta _1=0`$). Remark. The shear velocity tensor $`\sigma `$ and the expansion velocity $`\theta `$ can also be written in the forms $$\begin{array}{c}\sigma =\frac{1}{2}\{h_u(_u\overline{g}\mathrm{\pounds }_u\overline{g})h_u\frac{1}{n1}(h_u[_u\overline{g}\mathrm{\pounds }_u\overline{g}])h_u\}\text{ }=\\ =\frac{1}{2}\{h_{ik}(g^{kl}\text{ }_{;m}u^m\mathrm{\pounds }_ug^{kl})h_{lj}\frac{1}{n1}h_{kl}(g^{kl}\text{ }_{;m}u^m\mathrm{\pounds }_ug^{kl})h_{ij}\}e^i.e^j\text{ ,}\end{array}$$ (19) $$\theta =\frac{1}{2}h_u[_u\overline{g}\mathrm{\pounds }_u\overline{g}]=\frac{1}{2}h_{ij}(g^{ij}\text{ }_{;k}u^k\mathrm{\pounds }_ug^{ij})\text{ .}$$ The physical interpretation of the velocity tensors $`d`$, $`\sigma `$, $`\omega `$, and of the invariant $`\theta `$ for the case of $`V_4`$-spaces , , can also be extended for $`(L_4,g)`$-spaces. In this case the torsion plays an equivalent role in the velocity tensors as the covariant derivative. It is easy to be seen that the existence of some kinematic characteristics ($`{}_{s}{}^{}P`$, $`Q`$, $`\theta _1`$) depends on the existence of the torsion tensor field. They vanish if it is equal to zero (e.g. in $`V_n`$-spaces). On the other side, the kinematic characteristics, induced by the torsion, can compensate the result of the action of the torsion-free kinematic characteristics. If $`d=0`$, $`\sigma =0`$, $`\omega =0`$, $`\theta =0`$, then we could have the relations $`d_0=d_1`$, $`{}_{s}{}^{}E=_sP`$, $`S=Q`$, $`\theta _0=\theta _1`$ respectively leading to vanishing the relative velocity $`{}_{rel}{}^{}v`$ under the additional conditions $`g(u,\xi )=l=0`$ and $`\mathrm{\pounds }_\xi u=0`$. Since $`_u\xi =\frac{\overline{l}}{e}u+_{rel}v`$ and $`\overline{l}=g(_u\xi ,u)=ulg(\xi ,a)(_ug)(\xi ,u)=g(\xi ,a)(_ug)(\xi ,u)`$ for $`l=0`$, $`_u\xi =(_ug)(\xi ,u)`$ can be seen under the above conditions as a measure for the non-metricity ($`_ug`$) of the space-time if $`a=0`$ and $`{}_{rel}{}^{}v=0`$. ### 2.3 Special contravariant vector fields with vanishing kinematic characteristics related to the relative velocity The explicit forms of the quantities $`d`$, $`\sigma `$, $`\omega `$, and $`\theta `$ related to the relative velocity can be used for finding conditions for existence of special types of contravariant vector fields with vanishing characteristics induced by the relative velocity. #### 2.3.1 Contravariant vector fields with vanishing deformation velocity $`(d=0)`$ If we consider the explicit form for $`d`$ $$d:=h_u(k)h_u$$ we can prove the following propositions: Proposition 1. The necessary and sufficient condition for the existence of a non-null contravariant vector field $`u`$ with vanishing deformation velocity $`(d=0)`$ is the condition $$k=\frac{1}{e}\{au+u[g(u)](k)\}\frac{1}{2e^2}[ue(_ug)(u,u)]uu\text{ ,}$$ or in a co-ordinate basis $`k^{ij}`$ $`=`$ $`{\displaystyle \frac{1}{e}}(a^iu^j+u^ik^{lj}g_{lm}u^m)`$ $`{\displaystyle \frac{1}{2e^2}}(e_{,k}u^kg_{kl;m}u^mu^ku^l)u^iu^j\text{ .}`$ Proof. 1. Necessity. From $`d=h_u(k)h_u`$, after writing the explicit form of $`h_u`$, it follows that $`d`$ $`=`$ $`g(k)g{\displaystyle \frac{1}{e}}g(u)[g(u)](k)g{\displaystyle \frac{1}{e}}g(k)[g(u)]g(u)`$ $`+{\displaystyle \frac{1}{e^2}}[g(u)](k)[g(u)]uu\text{ .}`$ Since $`(k)[g(u)]=a=_uu`$, it follows further that $$[g(u)](k)[g(u)]=[g(u)](a)=g(u,a)=\frac{1}{2}[ue(_ug)(u,u)]\text{ .}$$ Therefore, $`d`$ $`=`$ $`0:g(k)g={\displaystyle \frac{1}{e}}\{g(u)[g(u)](k)g+g(a)g(u)\}`$ $`{\displaystyle \frac{1}{e^2}}g(u,a)g(u)g(u)\text{ .}`$ From $`\overline{g}(g(k)g)\overline{g}=k`$, we obtain $`g(k)g`$ $`=`$ $`{\displaystyle \frac{1}{e}}\{au+u[g(u)](k)\}`$ $`{\displaystyle \frac{1}{2e^2}}[ue(_ug)(u,u)]uu\text{ .}`$ 2. Sufficiency. From the explicit form of $`k`$, it follows that $`g(k)g`$ $`=`$ $`{\displaystyle \frac{1}{e}}g(u)[g(u)](k)g+{\displaystyle \frac{1}{e}}g(k)[g(u)]g(u)`$ $`{\displaystyle \frac{1}{e^2}}g(u,a)g(u)g(u)`$ which is identical to $`h_u(k)h_u=d=0`$. Special case: $`_uu=a:=0`$, $`_\xi g:=0`$ for $`\xi T(M)`$ ($`U_n`$-space), $`ue=0:e=`$const.$`0`$ ($`u`$ is normalized, non-null contravariant vector field). $$d=0:k=\frac{1}{e}u[g(u)](k)\text{ ,}$$ $$(k)[g(\xi )]=\frac{1}{e}u[g(u)](k)[g(\xi )]=\frac{1}{e}[g(u)](k)[g(\xi )]u\text{ ,}$$ $`(k)[g(\xi )]`$ $`=`$ $`(u_{;l}^iT_{lk}^iu^k)g^{lm}g_{mj}\xi ^j_i=_\xi uT(\xi ,u)=`$ $`=`$ $`_u\xi \mathrm{\pounds }_\xi u\text{ ,}`$ $${}_{rel}{}^{}v=\overline{g}[h_u(_u\xi )]=\overline{g}(h_u)(\mathrm{\pounds }_\xi u)\text{ for }\xi T(M)\text{ .}$$ (20) Proposition 2. A sufficient condition for the existence of a non-null contravariant vector field with vanishing deformation velocity ($`d=0`$) is the condition $$k=0\text{ ,}$$ equivalent to the condition $$_\xi u=T(\xi ,u)\text{ for }\xi T(M)\text{ ,}$$ or in a co-ordinate basis $$k^{ij}=0:u_{;j}^i=T_{jk}^iu^k\text{ .}$$ Proof. From $`k=0`$ and $`(k)[g(\xi )]=_\xi uT(\xi ,u)`$ for $`\xi T(M)`$, it follows that $`_\xi uT(\xi ,u)=0`$ or in a co-ordinate basis $`u_{;j}^iT_{jk}^iu^k=0`$. In this case $`\mathrm{\pounds }_\xi u=_\xi u_u\xi T(\xi ,u)=_u\xi `$. Corollary. A deformation free contravariant vector field $`u`$ with $`k=0`$ is an auto-parallel contravariant vector field. Proof. It follows immediately from the condition $`_\xi u=T(\xi ,u)`$ and for $`\xi =u`$ that $`_uu=a=0`$. Proposition 3. The necessary condition for the existence of a deformation free contravariant vector field $`u`$ with $`k=0`$ is the condition $$[R(u,v)]\xi =[\mathrm{\pounds }\mathrm{\Gamma }(u,v)]\xi \text{ for }\xi ,vT(M)\text{ ,}$$ or in a co-ordinate basis $$R_{ilj}^ku^l=\mathrm{\pounds }_u\mathrm{\Gamma }_{ij}^k\text{ .}$$ Proof. By the use of the explicit form of the curvature operator $`R(u,v)`$ acting on a contravariant vector field $`\xi `$ $$[R(u,v)]\xi =_u_v\xi _v_u\xi _{\mathrm{\pounds }_uv}\xi \text{ , }\xi \text{}v\text{}uT(M)\text{ ,}$$ and the explicit form of the deviation operator $`\mathrm{\pounds }\mathrm{\Gamma }(u,v)`$ acting on a contravariant vector field $`\xi `$ $$[\mathrm{\pounds }\mathrm{\Gamma }(u,v)]\xi =\mathrm{\pounds }_u_v\xi _v\mathrm{\pounds }_u\xi _{\mathrm{\pounds }_uv}\xi $$ we obtain under the condition $`_\xi u=T(\xi ,u)`$ (equivalent to the condition $`\mathrm{\pounds }_u\xi =_u\xi `$) $`[\mathrm{\pounds }\mathrm{\Gamma }(u,v)]\xi `$ $`=`$ $`\mathrm{\pounds }_u_v\xi _v\mathrm{\pounds }_u\xi _{\mathrm{\pounds }_uv}\xi =`$ $`=`$ $`_u_v\xi _{_v\xi }uT(u,_v\xi )_v_u\xi _{\mathrm{\pounds }_uv}\xi =`$ $`=`$ $`_u_v\xi _v_u\xi _{\mathrm{\pounds }_uv}\xi [_{_v\xi }u+T(u,_v\xi )]=`$ $`=`$ $`[R(u,v)]\xi [_{_v\xi }u+T(u,_v\xi )]\text{ .}`$ Since $$_{_v\xi }u+T(u,_v\xi )=0\text{ for }v,\xi T(M)\text{ , }$$ we have $$[R(u,v)]\xi =[\mathrm{\pounds }\mathrm{\Gamma }(u,v)]\xi \text{ for }v,\xi T(M)\text{ . }$$ The last condition appears as the integrability condition for the equation for $`u`$ $$_\xi u=T(\xi ,u)\text{ for }\xi T(M)\text{ .}$$ Proposition 4. A deformation-free contravariant non-null vector field $`u`$ with $`k=0`$ is an auto-parallel non-null shear-free ($`\sigma =0`$), rotation-free ($`\omega =0`$) and expansion-free ($`\theta =0`$) contravariant vector field with vanishing deformation acceleration ($`A=0`$) . Proof. If $`k=k^{ij}_i_j=0`$ and $`_uu=a=0`$, then $`k_s=k^{(ij)}_i._j=\frac{1}{2}(k^{ij}+k^{ji})_i._j=0`$, and $`k_a=k^{[ij]}_i_j=\frac{1}{2}(k^{ij}k^{ji})_i_j=0`$. Therefore, $`\sigma =h_u(k_s)h_u\frac{1}{n1}\overline{g}[h_u(k_s)h_u]h_u=0`$, $`\theta =\overline{g}[h_u(k_s)h_u]=0`$, and $`\omega =h_u(k_a)h_u=0`$. From the explicit form of the deformation acceleration $`A`$, it follows that $`A=0`$. From the identity for the Riemannian tensor $`R_{jkl}^i`$ $`R_{jkl}^i+R_{ljk}^i+R_{klj}^i`$ $``$ $`T_{jk}^i{}_{;l}{}^{}+T_{lj}^i{}_{;k}{}^{}+T_{kl}^i{}_{;j}{}^{}+`$ $`+T_{jk}^mT_{ml}^i+T_{lj}^mT_{mk}^i+T_{kl}^mT_{mj}^i\text{ ,}`$ after contraction with $`g_i^l`$ and summation over $`l`$ we obtain $`R_{jk}R_{kj}+R_{ijk}^i`$ $``$ $`T_{jk}^i{}_{;i}{}^{}+T_{ij}^i{}_{;k}{}^{}T_{ik}^i{}_{;j}{}^{}+`$ $`+T_{jk}^mT_{mi}^i+T_{ij}^mT_{mk}^iT_{ik}^mT_{mj}^i\text{ .}`$ If we introduce the abbreviations $${}_{a}{}^{}R_{jk}^{}:=\frac{1}{2}(R_{jk}R_{kj})\text{ , }T_{ji}^i:=T_j\text{ , }$$ where $`T_{ik}^i=T_{ki}^i=T_k`$, then the last expression for $`R_{jk}`$ can be written in the form $`2_aR_{jk}`$ $``$ $`R_{ijk}^i+T_{jk}^i{}_{;i}{}^{}+T_{ij}^i{}_{;k}{}^{}T_{ik}^i{}_{;j}{}^{}+`$ $`+T_{jk}^mT_m+T_{ij}^mT_{mk}^iT_{ik}^mT_{mj}^i\text{ .}`$ Therefore, $`{}_{a}{}^{}R_{ij}^{}u^j`$ can be written in the form $`2_aR_{ij}u^j+R_{ijk}^iu^j`$ $`=`$ $`T_{k;j}u^jT_{j;k}u^j+T_{jk}^i{}_{;i}{}^{}u^j+`$ $`+T_mT_{jk}^mu^j+T_{ij}^mu^jT_{mk}^iT_{ik}^mT_{mj}^iu^j\text{ .}`$ From the other side, from $`u_{;j}^i=T_{jl}^iu^l`$ and $`a^i=u_{;j}^iu^j=0`$, we have $$u_{;j;k}^i=T_{jl}^i{}_{;k}{}^{}u^l+T_{jm}^iT_{kl}^mu^l\text{ ,}$$ $`u_{;j;k}^iu_{;k;j}^i`$ $`=`$ $`u^lR_{ljk}^i+T_{jk}^mT_{ml}^iu^l=`$ $`=`$ $`T_{jl}^i{}_{;k}{}^{}u^l+T_{jm}^iT_{kl}^mu^l`$ $`T_{kl}^i{}_{;j}{}^{}u^lT_{km}^iT_{jl}^mu^l\text{ ,}`$ $`u^lR_{ljk}^i`$ $`=`$ $`T_{jk}^mT_{ml}^iu^l+T_{kl}^i{}_{;j}{}^{}u^lT_{jl}^i{}_{;k}{}^{}u^l+`$ $`+T_{km}^iT_{jl}^mu^lT_{jm}^iT_{kl}^mu^l\text{ ,}`$ $$R_{lj}u^l=T_{l;j}u^lT_{jl}^i{}_{;i}{}^{}u^lT_mT_{jl}^mu^l\text{ ,}$$ $$R_{lj}u^lu^j=_sR_{lj}u^lu^j=I=T_{l;j}u^lu^j=(T_iu^i)_{;j}u^j=\dot{\theta }_1\text{ . }$$ By the use of the decompositions $`R_{ij}=_aR_{ij}+_sR_{ij}`$, $`R_{ij}u^j=_aR_{ij}u^j+_sR_{ij}u^j`$, and the above expression for $`R_{lj}u^l`$, we can find the following relations $$2_aR_{jk}u^j=T_{k;j}u^jT_{j;k}u^j2T_{kj}^i{}_{;i}{}^{}u^j2T_mT_{kj}^mu^j\text{ ,}$$ $$R_{ijk}^iu^j=(T_{kj}^i{}_{;i}{}^{}+T_mT_{kj}^m+T_{ij}^mT_{mk}^iT_{ik}^mT_{mj}^i)u^j\text{ ,}$$ $$2_sR_{jk}u^j=(T_{j;k}+T_{k;j})u^j\text{ }.$$ It follows that in a $`(\overline{L}_n,g)`$-space the projections of the symmetric part of the Ricci tensor on the non-null contravariant vector field $`u`$ with $`k=0`$ is depending on the covariant derivatives of $`T_i`$ (respectively on the torsion $`T_{ik}^l`$) and not on the torsion $`T_{ik}^l`$ itself. #### 2.3.2 Contravariant non-null (nonisotropic) vector fields with vanishing shear velocity $`(\sigma =0)`$ If we consider the explicit form of the shear velocity tensor (shear velocity, shear) $$\sigma =h_u(k_s)h_u\frac{1}{n1}\overline{g}[h_u(k_s)h_u]h_u$$ we can prove the following propositions: Proposition 5. The necessary and sufficient condition for the existence of a shear-free non-null contravariant vector field is the condition $`k_s`$ $`=`$ $`{\displaystyle \frac{1}{2e}}\{ua+au+u[g(u)](k)+[g(u)](k)u`$ $`{\displaystyle \frac{1}{e}}[ue(_ug)(u,u)]uu\}+`$ $`+{\displaystyle \frac{1}{n1}}\theta h^u\text{,}`$ or in a co-ordinate basis $`h_s^{ij}`$ $`=`$ $`{\displaystyle \frac{1}{2e}}\{u^ia^j+u^ja^i+u^ig_{\overline{m}\overline{n}}u^nk^{mj}+u^jg_{\overline{m}\overline{n}}u^nk^{mi}`$ $`{\displaystyle \frac{1}{e}}[e_{,k}u^kg_{km;n}u^nu^{\overline{k}}u^{\overline{m}}]u^iu^j\}+{\displaystyle \frac{1}{n1}}\theta h^{ij}\text{ .}`$ Proof. 1. Necessity. From $`\sigma =0`$, it follows that $`h_u(k_s)h_u=\frac{1}{n1}\overline{g}[h_u(k_s)h_u]h_u=\frac{1}{n1}\theta h_u`$. Further, from the explicit form of $`h_u`$ and $`k_s`$, it follows that $`h_u(k_s)h_u`$ $`=`$ $`{\displaystyle \frac{1}{n1}}\theta h_u=g(k_s)g{\displaystyle \frac{1}{e}}\{g(u)[g(u)](k_s)g+g(k_s)[g(u)]g(u)\}+`$ $`+{\displaystyle \frac{1}{e^2}}[g(u)](k_s)[g(u)]g(u)g(u)\text{ ,}`$ or $`k_s`$ $`=`$ $`{\displaystyle \frac{1}{e}}\{u[g(u)](k_s)+(k_s)[g(u)]u\}{\displaystyle \frac{1}{e^2}}[g(u)](k_s)[g(u)]uu+`$ $`+{\displaystyle \frac{1}{n1}}\theta \overline{g}(h_u)\overline{g}\text{ .}`$ Since $`[g(u)](k_s)=(k_s)[g(u)]`$, $`(k_s)[g(u)]=\frac{1}{2}\{(k)[g(u)]+[g(u)](k)\}`$, $`(k)[g(u)]=a`$, $`(k_s)[g(u)]=\frac{1}{2}\{a+[g(u)](k)\}`$, $`[g(u)](k_s)[g(u)]=[g(u)](k)[g(u)]=g(u,a)=\frac{1}{2}[ue(_ug)(u,u)]`$, $`\theta =\overline{g}[h_u(k_s)h_u]=\overline{g}[h_u(k)h_u]`$, and $`\overline{g}(h_u)\overline{g}=h^u`$, the explicit form of $`k_s`$ can be found as $`k_s`$ $`=`$ $`{\displaystyle \frac{1}{2e}}\{ua+au+u[g(u)](k)+[g(u)](k)u`$ $`{\displaystyle \frac{1}{e}}[ue(_ug)(u,u)]uu\}+`$ $`+{\displaystyle \frac{1}{n1}}\theta h^u\text{ .}`$ 2. Sufficiency. From the last expression and the above relations, it follows that $`h_u(k_s)h_u=\frac{1}{n1}\theta h_u`$, and therefore $`\sigma =0`$. Proposition 6. A sufficient condition for the existence of a non-null vector field with vanishing shear velocity ($`\sigma =0`$) and expansion velocity ($`\theta =0`$) is the condition $$h_u(k_s)h_u=0\text{ ,}$$ identical with the condition $`k_s`$ $`=`$ $`{\displaystyle \frac{1}{2e}}\{ua+au+u[g(u)](k)+[g(u)](k)u`$ $`{\displaystyle \frac{1}{e}}[ue(_ug)(u,u)]uu\}`$ Proof. If $`h_u(k_s)h_u=0`$, then $`\theta =\overline{g}[h_u(k_s)h_u]=0`$. Therefore, $`\sigma =h_u(k_s)h_u\frac{1}{n1}\theta h_u=0`$. Corollary. If $`h_u(k_s)h_u=0`$, then $$g[k_s]=\frac{1}{2e}[ue(_ug)(u,u)]\text{ .}$$ Proof. It follows from the above proposition that $`\theta `$ $`=`$ $`g[k_s]{\displaystyle \frac{1}{e}}g(u,a)=0\text{ ,}`$ $`g[k_s]`$ $`=`$ $`{\displaystyle \frac{1}{e}}g(u,a)={\displaystyle \frac{1}{2e}}[ue(_ug)(u,u)]\text{ .}`$ #### 2.3.3 Contravariant non-null vector fields with vanishing rotation velocity $`(\omega =0)`$ If we consider the explicit form for $`\omega `$ $$\omega =h_u(k_a)h_u$$ we can prove the following propositions: Proposition 7. The necessary and sufficient condition for the existence of a contravariant non-null vector field with vanishing rotation velocity $`(\omega =0)`$ is the condition $$k_a=\frac{1}{e}\{u[g(u)](k_a)[g(u)](k_a)u\}\text{ ,}$$ (21) or in a co-ordinate basis $$k_a^{ij}=\frac{1}{e}g_{\overline{m}\overline{n}}u^n(u^ik_a^{mj}u^jk_a^{mi})\text{ .}$$ (22) Proof. 1. Necessity. Form $`h_u(k_a)h_u=0`$, it follows that $`h_u(k_a)h_u`$ $`=`$ $`0=g(k_a)g{\displaystyle \frac{1}{e}}g(u)[g(u)](k_a)g{\displaystyle \frac{1}{e}}g(k_a)[g(u)]g(u)+`$ $`+{\displaystyle \frac{1}{e^2}}[g(u)](k_a)[g(u)]g(u)g(u)\text{ .}`$ Since $$[g(u)](k_a)[g(u)]=g_{\overline{i}\overline{m}}u^mk_a^{ij}g_{\overline{j}\overline{n}}u^n=g_{\overline{i}\overline{m}}u^mk_a^{ij}g_{\overline{j}\overline{n}}u^n\text{ ,}$$ we have $`[g(u)](k_a)[g(u)]=0`$. Therefore, $$g(k_a)g=\frac{1}{e}\{g(u)[g(u)](k_a)g+g(k_a)[g(u)]g(u)\}\text{ .}$$ From the last expression and from the relation $`\overline{g}[g(k_a)g]\overline{g}=k_a`$, it follows that $`k_a`$ $`=`$ $`{\displaystyle \frac{1}{e}}\{u[g(u)](k_a)+(k_a)[g(u)]u\}=`$ $`=`$ $`{\displaystyle \frac{1}{e}}\{u[g(u)](k_a)[g(u)](k_a)u\}\text{ ,}`$ because of $`(k_a)[g(u)]=`$ $`[g(u)](k_a)`$. In a co-ordinate basis we obtain (22). 2. Sufficiency. From (\[10.1a\]) we have $$g(k_a)g=\frac{1}{e}\{g(u)[g(u)](k_a)g+g(k_a)[g(u)]g(u)\}\text{ ,}$$ which is identical to $`h_u(k_a)h_u=0`$. On the other hand, after direct computations, it follows that $$(k_a)[g(u)]=\frac{1}{2}\{(k)[g(u)][g(u)](k)\}\text{ .}$$ Since $`(k)[g(u)]=a`$, we have the relation $$(k_a)[g(u)]=\frac{1}{2}\{a[g(u)](k)\}\text{ .}$$ Then $$k_a=\frac{1}{2e}\{auua+u[g(u)](k)[g(u)](k)u\}\text{ .}$$ Proposition 8. A sufficient condition for the existence of a contravariant non-null vector field with vanishing rotation velocity $`(\theta =0)`$ is the condition $$k_a=0\text{ .}$$ Proof. If $`k_a=0`$, then it follows directly from $`\omega =h_u(k_a)h_u`$ that $`\omega =0`$. In a co-ordinate basis $`k_a`$ is equivalent to the expression $$u_{;l}^ig^{lj}u_{;l}^jg^{li}=(T_{lm}^ig^{lj}T_{lm}^jg^{li})u^m\text{ .}$$ On the other side, after multiplying the last expression with $`g_{\overline{j}\overline{k}}u^k`$ and summarizing over $`j`$, we obtain $$a^i=u_{;k}^iu^k=g_{\overline{j}\overline{k}}u^k(u_{;l}^jT_{lm}^ju^m)g^{li}=g_{\overline{j}\overline{k}}u^kk^{ji}\text{ ,}$$ or in a form $$a=[g(u)](k)\text{ .}$$ Proposition 9. The necessary condition for $`k_a=0`$ is the condition $$a=[g(u)](k)\text{.}$$ Proof. From $`k_a=0`$ and $`(k_a)[g(u)]=\frac{1}{2}\{a[g(u)](k)\}`$, it follows that $`a=[g(u)](k)`$. If the rotation velocity $`\omega `$ vanishes $`(\omega =0)`$ for an auto-parallel $`(_uu=0)`$ contravariant non-null vector field $`u`$, then the rotation acceleration tensor $`W`$ will have the form $$W=\frac{1}{2}[h_u(_u\overline{g})\sigma \sigma (_u\overline{g})h_u]\text{ .}$$ From the last expression it is obvious that the nonmetricity $`(_ug0)`$ in a $`(\overline{L}_n,g)`$-space is responsible for nonvanishing the rotation acceleration $`W`$. ### 2.4 Relative velocity and change of the length of a contravariant vector field over a $`(L_n,g)`$-space Let us now consider the influence of the kinematic characteristics related to the relative velocity and respectively to the relative velocity, induced by the torsion, upon the change of the length of a contravariant vector field. Let $`l_\xi =g(\xi ,\xi )^{\frac{1}{2}}`$ be the length of a contravariant vector field $`\xi `$. The rate of change $`ul_\xi `$ of $`l_\xi `$ along a contravariant vector field $`u`$ can be expressed in the form $`\pm \mathrm{\hspace{0.17em}2}.l_\xi .(ul_\xi )=(_ug)(\xi ,\xi )+2g(_u\xi ,\xi )`$. By the use of the projections of $`\xi `$ and $`_u\xi `$ along and orthogonal to $`u`$ we can find the relations $$\begin{array}{c}2g(_u\xi ,\xi )=2.\frac{l}{e}.g(_u\xi ,u)+2g(_{rel}v,\xi _{})\text{ ,}\\ (_ug)(\xi ,\xi )=(_ug)(\xi _{},\xi _{})+2.\frac{l}{e}.(_ug)(\xi _{},u)+\frac{l^2}{e^2}.(_ug)(u,u)\text{ .}\end{array}$$ Then, it follows for $`\pm \mathrm{\hspace{0.17em}2}.l_\xi .(ul_\xi )`$ the expression $$\begin{array}{c}\pm \mathrm{\hspace{0.17em}2}.l_\xi .(ul_\xi )=(_ug)(\xi _{},\xi _{})+2.\frac{l}{e}.[(_ug)(\xi _{},u)+g(_u\xi ,u)]+\\ +\frac{l^2}{e^2}.(_ug)(u,u)+2g(_{rel}v,\xi _{})\text{ ,}\end{array}$$ (23) where $$g(_{rel}v,\xi _{})=\frac{l}{e}.h_u(a,\xi _{})+h_u(\mathrm{\pounds }_u\xi ,\xi _{})+d(\xi _{},\xi _{})\text{ ,}$$ (24) $$d(\xi _{},\xi _{})=\sigma (\xi _{},\xi _{})+\frac{1}{n1}.\theta .l_\xi _{}^2\text{ .}$$ (25) For finding out the last two expressions the following relations have been used: $$g(\overline{g}(h_u)a,\xi _{})=h_u(a,\xi _{})\text{ , }g(\overline{g}(h_u)(\mathrm{\pounds }_u\xi ),\xi _{})=h_u(\mathrm{\pounds }_u\xi ,\xi _{})\text{ ,}$$ (26) $$g(\overline{g}[d(\xi )],\xi _{})=d(\xi _{},\xi _{})\text{ , }d(\xi )=d(\xi _{})\text{ .}$$ (27) Special case: $`g(u,\xi )=l:=0:\xi =\xi _{}`$. $$\pm \mathrm{\hspace{0.17em}2}.l_\xi _{}.(ul_\xi _{})=(_ug)(\xi _{},\xi _{})+2g(_{rel}v,\xi _{})\text{ .}$$ (28) Special case: $`V_n`$-spaces: $`_\eta g=0`$ for $`\eta T(M)`$ ($`g_{ij;k}=0`$), $`g(u,\xi )=l:=0:\xi =\xi _{}`$. $$\pm l_\xi _{}.(ul_\xi _{})=g(_{rel}v,\xi _{})\text{ .}$$ (29) In $`(L_n,g)`$-spaces the covariant derivative $`_ug`$ of the metric tensor field $`g`$ along $`u`$ can be decomposed in its trace free part $`{}_{}{}^{s}_{u}^{}g`$ and its trace part $`\frac{1}{n}.Q_u.g`$ as $$_ug=^s_ug+\frac{1}{n}.Q_u.g\text{ , }dimM=n\text{ ,}$$ where $`\overline{g}[^s_ug]=0`$, $`Q_u=\overline{g}[_ug]=g^{kl}.g_{kl;j}.u^j=Q_j.u^j`$, $`Q_j=g^{kl}.g_{kl;j}`$. Remark. The covariant vector $`\overline{Q}=\frac{1}{n}.Q=\frac{1}{n}.Q_j.dx^j=\frac{1}{n}.Q_\alpha .e^\alpha `$ is called Weyl’s covector field. The operator $`{}_{}{}^{s}_{u}^{}=_u\frac{1}{n}.Q_u`$ is called trace free covariant operator. If we use now the decomposition of $`_ug`$ in the expression for $`\pm \mathrm{\hspace{0.17em}2}.l_\xi .(ul_\xi )`$ we find the relation $$\begin{array}{c}\pm \mathrm{\hspace{0.17em}2}.l_\xi .(ul_\xi )=(^s_ug)(\xi ,\xi )+\frac{1}{n}.Q_u.l_\xi ^2+2g(_u\xi ,\xi )=\\ =(^s_ug)(\xi _{},\xi _{})+\\ +\frac{l}{e}.[2.(^s_ug)(\xi _{},u)+2.g(_u\xi ,u)+\frac{l}{e}.(^s_ug)(u,u)]+\\ +\frac{1}{n}.Q_u.(l_\xi _{}^2+\frac{l^2}{e})+2.g(_{rel}v,\xi _{})\text{ ,}\end{array}$$ (30) where $`l_\xi _{}^2=g(\xi _{},\xi _{})`$, $`l=g(\xi ,u)`$. For $`l_\xi 0:`$ $$ul_\xi =\pm \frac{1}{2.l_\xi }(^s_ug)(\xi ,\xi )\pm \frac{1}{2.n}.Q_u.l_\xi \pm \frac{1}{l_\xi }.g(_u\xi ,\xi )\text{ .}$$ (31) In the case of a parallel transport ($`_u\xi =0`$) of $`\xi `$ along $`u`$ the change $`ul_\xi `$ of the length $`l_\xi `$ is $$ul_\xi =\pm \frac{1}{2.l_\xi }(^s_ug)(\xi ,\xi )\pm \frac{1}{2.n}.Q_u.l_\xi \text{ . }$$ (32) Special case: $`_u\xi =0`$ and $`{}_{}{}^{s}_{u}^{}g=0`$. $$ul_\xi =\pm \frac{1}{2.n}.Q_u.l_\xi \text{ . }$$ (33) Special case: $`g(u,\xi )=l:=0:\xi =\xi _{}`$. $$\pm 2.l_\xi _{}.(ul_\xi _{})=(^s_ug)(\xi _{},\xi _{})+\frac{1}{n}.Q_u.l_\xi _{}^2+2.g(_{rel}v,\xi _{})\text{ .}$$ $$ul_\xi _{}=\pm \frac{1}{2.l_\xi _{}}.(^s_ug)(\xi _{},\xi _{})\pm \frac{1}{2n}.Q_u.l_\xi _{}\pm \frac{1}{l_\xi _{}}.g(_{rel}v,\xi _{})\text{ , }l_\xi _{}0\text{ .}$$ (34) Special case: Quasi-metric transports: $`_ug:=2.g(u,\eta ).g`$, $`u`$, $`\eta T(M)`$. $$\pm 2.l_\xi .(ul_\xi )=2.g(u,\eta ).(l_\xi _{}^2+\frac{l^2}{e})+2.[\frac{l}{e}.g(_u\xi ,u)+g(_{rel}v,\xi _{})]\text{ .}$$ (35) ### 2.5 Change of the cosine between two contravariant vector fields and the relative velocity The cosine between two contravariant vector fields $`\xi `$ and $`\eta `$ has been defined as $`g(\xi ,\eta )=l_\xi .l_\eta .\mathrm{cos}(\xi ,\eta )`$. The rate of change of the cosine along a contravariant vector field $`u`$ can be found in the form $$\begin{array}{c}l_\xi .l_\eta .\{u[\mathrm{cos}(\xi ,\eta )]\}=(_ug)(\xi ,\eta )+g(_u\xi ,\eta )+g(\xi ,_u\eta )\\ [l_\eta .(ul_\xi )+l_\xi .(ul_\eta )].\mathrm{cos}(\xi ,\eta )\text{ .}\end{array}$$ Special case: $`_u\xi =0`$, $`_u\eta =0`$, $`{}_{}{}^{s}_{u}^{}g=0`$. $$l_\xi .l_\eta .\{u[\mathrm{cos}(\xi ,\eta )]\}=\frac{1}{n}.Q_u.g(\xi ,\eta )[l_\eta .(ul_\xi )+l_\xi .(ul_\eta )].\mathrm{cos}(\xi ,\eta )\text{ .}$$ Since $`g(\xi ,\eta )=l_\xi .l_\eta .\mathrm{cos}(\xi ,\eta )`$, it follows from the last relation $$l_\xi .l_\eta .\{u[\mathrm{cos}(\xi ,\eta )]\}=\{\frac{1}{n}.Q_u.l_\xi .l_\eta [l_\eta .(ul_\xi )+l_\xi .(ul_\eta )]\}.\mathrm{cos}(\xi ,\eta )\text{ .}$$ Therefore, if $`\mathrm{cos}(\xi ,\eta )=0`$ between two parallel transported along $`u`$ vector fields $`\xi `$ and $`\eta `$, then the right angle between them \[determined by the condition $`\mathrm{cos}(\xi ,\eta )=0`$\] does not change along the contravariant vector field $`u`$. In the cases, when $`\mathrm{cos}(\xi ,\eta )0`$, the rate of change of the cosine of the angle between two vector fields $`\xi `$ and $`\eta `$ is linear to $`\mathrm{cos}(\xi ,\eta )`$. By the use of the definitions and the relations: $$_{rel}v_\xi :=\overline{g}[h_u(_u\xi )]=_{rel}v\text{ , }_{rel}v_\eta :=\overline{g}[h_u(_u\eta )]\text{ ,}$$ (36) $$\begin{array}{c}g(_u\xi ,\eta )=\frac{1}{e}.g(u,\eta ).g(_u\xi ,u)+g(_{rel}v_\xi ,\eta )\text{ ,}\\ g(_u\eta ,\xi )=\frac{1}{e}.g(u,\xi ).g(_u\eta ,u)+g(_{rel}v_\eta ,\xi )\text{ ,}\end{array}$$ (37) $$(_ug)(\xi ,\eta )=(^s_ug)(\xi ,\eta )+\frac{1}{n}.Q_u.g(\xi ,\eta )\text{ ,}$$ (38) $$\begin{array}{c}(^s_ug)(\xi ,\eta )=(^s_ug)(\xi _{},\eta _{})+\frac{l}{e}.(^s_ug)(u,\eta _{})+\frac{\stackrel{~}{l}}{e}.(^s_ug)(\xi _{},u)+\\ +\frac{l}{e}.\frac{\stackrel{~}{l}}{e}.(^s_ug)(u,u)\text{ , }\stackrel{~}{l}=g(u,\eta )\text{ , }\eta _{}=\overline{g}[h_u(\eta )]\text{ , }l=g(u,\xi )\text{ ,}\end{array}$$ (39) $$\begin{array}{c}(_ug)(\xi ,\eta )=(^s_ug)(\xi ,\eta )+\frac{1}{n}.Q_u.g(\xi ,\eta )=\\ =(^s_ug)(\xi _{},\eta _{})+\frac{l}{e}.(^s_ug)(u,\eta _{})+\frac{\stackrel{~}{l}}{e}.(^s_ug)(\xi _{},u)+\\ +\frac{l}{e}.\frac{\stackrel{~}{l}}{e}.(^s_ug)(u,u)+\frac{1}{n}.Q_u.[\frac{l.\stackrel{~}{l}}{e}+g(\xi _{},\eta _{})]\text{ ,}\end{array}$$ (40) the expression of $`l_\xi .l_\eta .\{u[\mathrm{cos}(\xi ,\eta )]\}`$ follows in the form $$\begin{array}{c}l_\xi .l_\eta .\{u[\mathrm{cos}(\xi ,\eta )]\}=(^s_ug)(\xi _{},\eta _{})+\frac{l}{e}.[(^s_ug)(u,\eta _{})+g(_u\eta ,u)]+\\ +\frac{\stackrel{~}{l}}{e}.[(^s_ug)(\xi _{},u)+g(_u\xi ,u)]+\frac{l.\stackrel{~}{l}}{e^2}.(^s_ug)(u,u)+\\ +\frac{1}{n}.Q_u.[\frac{l.\stackrel{~}{l}}{e}+g(\xi _{},\eta _{})]+g(_{rel}v_\xi ,\eta )+g(_{rel}v_\eta ,\xi )\\ [l_\eta .(ul_\xi )+l_\xi .(ul_\eta )].\mathrm{cos}(\xi ,\eta )\text{ .}\end{array}$$ (41) Special case: $`g(u,\xi )=l:=0`$, $`g(u,\eta )=\stackrel{~}{l}:=0:\xi =\xi _{}`$, $`\eta =\eta _{}`$. $$\begin{array}{c}l_\xi _{}.l_\eta _{}.\{u[\mathrm{cos}(\xi _{},\eta _{})]\}=(^s_ug)(\xi _{},\eta _{})+\frac{1}{n}.Q_u.l_\xi _{}.l_\eta _{}.\mathrm{cos}(\xi _{},\eta _{})+\\ +g(_{rel}v_\xi _{},\eta _{})+g(_{rel}v_\eta _{},\xi _{})[l_\eta _{}.(ul_\xi _{})+l_\xi _{}.(ul_\eta _{})].\mathrm{cos}(\xi _{},\eta _{})\text{ ,}\end{array}$$ (42) where $`g(\xi _{},\eta _{})=l_\xi _{}.l_\eta _{}.\mathrm{cos}(\xi _{},\eta _{})`$. Special case: $`{}_{}{}^{s}_{u}^{}g:=0:_ug=\frac{1}{n}.Q_u.g`$ (Weyl’s space with torsion). $$\pm 2.l_\xi .(ul_\xi )=2.\frac{l}{e}.g(_u\xi ,u)+\frac{1}{n}.Q_u.(l_\xi _{}^2+\frac{l^2}{e})+2.g(_{rel}v,\xi _{})\text{ ,}$$ (43) $$g(_u\xi ,u)=ul\frac{1}{n}.Q_u.lg(\xi ,a)\text{ , }a=_uu\text{ ,}$$ (44) $$\begin{array}{c}l_\xi .l_\eta .\{u[\mathrm{cos}(\xi ,\eta )]\}=\frac{l}{e}.g(_u\eta ,u)+\frac{\stackrel{~}{l}}{e}.g(_u\xi ,u)+\\ +\frac{1}{n}.Q_u.[\frac{l.\stackrel{~}{l}}{e}+g(\xi _{},\eta _{})]+g(_{rel}v_\xi ,\eta )+g(_{rel}v_\eta ,\xi )\\ [l_\eta .(ul_\xi )+l_\xi .(ul_\eta )].\mathrm{cos}(\xi ,\eta )\text{ ,}\end{array}$$ (45) Special case: $`{}_{}{}^{s}_{u}^{}g:=0:_ug=\frac{1}{n}.Q_u.g`$, $`g(u,\xi )=l:=0`$, $`g(u,\eta )=\stackrel{~}{l}:=0`$ (Weyl’s space with torsion and orthogonal to $`u`$ vector fields $`\xi _{}`$ and $`\eta _{}`$). $$\pm 2.l_\xi _{}.(ul_\xi _{})=\frac{1}{n}.Q_u.l_\xi _{}^2+2.g(_{rel}v,\xi _{})\text{ ,}$$ (46) $$\begin{array}{c}l_\xi _{}.l_\eta _{}.\{u[\mathrm{cos}(\xi _{},\eta _{})]\}=\frac{1}{n}.Q_u.g(\xi _{},\eta _{})+\\ +g(_{rel}v_\xi _{},\eta _{})+g(_{rel}v_\eta _{},\xi _{})[l_\eta _{}.(ul_\xi _{})+l_\xi _{}.(ul_\eta _{})].\mathrm{cos}(\xi _{},\eta _{})\text{ }=\\ =g(_{rel}v_\xi _{},\eta _{})+g(_{rel}v_\eta _{},\xi _{})\\ [l_\eta _{}.(ul_\xi _{})+l_\xi _{}.(ul_\eta _{})+\frac{1}{n}.Q_u.l_\xi _{}.l_\eta _{}].\mathrm{cos}(\xi _{},\eta _{})\end{array}$$ (47) Special case: $`U_n`$-spaces: $`_\xi g=0`$ for $`\xi T^{}(M)`$: $$\begin{array}{c}l_\xi _{}.l_\eta _{}.\{u[\mathrm{cos}(\xi _{},\eta _{})]\}=\\ =g(_{rel}v_\xi _{},\eta _{})+g(_{rel}v_\eta _{},\xi _{})[l_\eta _{}.(ul_\xi _{})+l_\xi _{}.(ul_\eta _{})].\mathrm{cos}(\xi _{},\eta _{}).\end{array}$$ ### 2.6 Rate of change of the length of a vector field $`\xi _{}`$ connecting two particles (points) in the space-time as a $`U_n`$-space The distance between a particle (as a basic point) and an other particle (as an observed point) lying in a neighborhood of the first one can be determined by the use of the length of the vector field $`\xi _{}`$. The rate of change of the length $`l_\xi _{}`$ of the vector field $`\xi _{}`$ (along the vector field $`u`$) in a $`U_n`$-space can be given in the form $$ul_\xi _{}=\pm \frac{1}{l_\xi _{}}.g(_u\xi _{},\xi _{})\text{ , }l_\xi _{}0\text{ .}$$ If $`u`$ is tangential vector field along a congruence of curves with parameter $`s`$, i. e. if $`u=\frac{d}{ds}`$, then $`ul_\xi _{}`$ will have the form $$\frac{dl_\xi _{}}{ds}=\pm \frac{1}{l_\xi _{}}.g(_u\xi _{},\xi _{})=\pm \frac{1}{l_\xi _{}}.g_{ij}.\xi _{}^i{}_{;k}{}^{}.u^k.\xi _{}^j\text{ .}$$ By the use of the decomposition of $`_u\xi _{}`$, $$_u\xi _{}=\frac{\overline{l}}{e}.u+_{rel}v\text{}\overline{l}=g(_u\xi _{},u)\text{ ,}$$ under the condition for the orthogonality between $`u`$ and $`\xi _{}:g(u,\xi _{})=l=0`$, we obtain $$\frac{dl_\xi _{}}{ds}=\pm \frac{1}{l_\xi _{}}.g(_{rel}v,\xi _{})\text{ .}$$ Under the assumption $`\mathrm{\pounds }_\xi _{}u=\mathrm{\pounds }_u\xi _{}=0`$, it follows that $$\frac{dl_\xi _{}}{ds}=\pm \frac{1}{l_\xi _{}}.g(_{rel}v,\xi _{})=\pm \frac{1}{l_\xi _{}}.d(\xi _{},\xi _{})\text{ .}$$ (48) If we use the explicit form for $`d`$ and the fact that $`\omega (\xi _{},\xi _{})=0`$, then we can find the expression for the rate of change of $`l_\xi _{}`$ in the form $$\frac{dl_\xi _{}}{ds}=\pm \frac{1}{l_\xi _{}}.d(\xi _{},\xi _{})=\pm \frac{1}{l_\xi _{}}.\sigma (\xi _{},\xi _{})\pm \frac{1}{n1}.\theta .l_\xi _{}\text{ .}$$ (49) Remark. The sign $`\pm `$ depends on the sign of the metric $`g`$ (for $`n=4`$, sign $`g=\pm 2`$). Since the existence of the torsion could cause the condition $`d=0`$, respectively $`\sigma =0`$, and $`\theta =0`$, the length $`l_\xi _{}`$ would not change along $`u`$ and therefore, along the proper time $`\tau =\frac{s}{c}`$, if we consider $`\tau `$ as the proper time of the basic particle (point). ## 3 Relative acceleration and its kinematic characteristics induced by the torsion ### 3.1 Relative acceleration in $`(L_n,g)`$-spaces The notion relative acceleration vector field (relative acceleration) $`{}_{rel}{}^{}a`$ in $`(L_n,g)`$-spaces can be defined (in analogous way as $`{}_{rel}{}^{}v`$) (regardless of its physical interpretation) as the orthogonal to a non-null contravariant vector field $`u`$ \[$`g(u,u)=e0`$\] projection of the second covariant derivative (along the same non-null vector field $`u`$) of (another) contravariant vector field $`\xi `$, i.e. $$_{rel}a=\overline{g}(h_u(_u_u\xi ))=g^{ij}.h_{jk}.(\xi ^k\text{ }_{;l}.u^l)_{;m}.u^m.e_i\text{ .}$$ (50) $`_u_u\xi `$ $`=(\xi ^i`$ $`{}_{;l}{}^{}.u^l)_{;m}.u^m.e_i`$ is the second covariant derivative of a vector field $`\xi `$ along the vector field $`u`$. It is an essential part of all types of deviation equations in $`V_n`$\- and ($`L_n,g`$)-spaces , . If we take into account the expression for $`_u\xi :_u\xi =k[g(\xi )]\mathrm{\pounds }_\xi u`$, and differentiate covariant along $`u`$, then we obtain $$_u_u\xi =\{_u[(k)g]\}(\xi )+(k)(g)(_u\xi )_u(\mathrm{\pounds }_\xi u)\text{ .}$$ By means of the relations $$\begin{array}{c}k(g)\overline{g}=k\text{ , }_u[k(g)]=(_uk)(g)+k(_ug)\text{ ,}\\ \text{ }\{_u[k(g)]\}\overline{g}=_uk+k(_ug)\overline{g}\text{ ,}\end{array}$$ $`_u_u\xi `$ can be written in the form $$_u_u\xi =\frac{l}{e}.H(u)+B(h_u)\xi k(g)\mathrm{\pounds }_\xi u_u(\mathrm{\pounds }_\xi u)$$ (51) \[compare with $`_u\xi =\frac{l}{e}.a+k(h_u)\xi \mathrm{\pounds }_\xi u`$\], where $$H=B(g)=(_uk)(g)+k(_ug)+k(g)k(g)\text{ ,}$$ $$B=_uk+k(g)k+k(_ug)\overline{g}=_uk+k(g)kk(g)(_u\overline{g})\text{ .}$$ The orthogonal to $`u`$ covariant projection of $`_u_u\xi `$ will have therefore the form $$h_u(_u_u\xi )=h_u[\frac{l}{e}H(u)k(g)\mathrm{\pounds }_\xi u_u\mathrm{\pounds }_\xi u]+[h_u(B)h_u](\xi )\text{ .}$$ (52) In the special case, when $`g(u,\xi )=l=0`$ and $`\mathrm{\pounds }_\xi u=0`$ , the above expression has the simple form $$g(_{rel}a)=h_u(_u_u\xi )=[h_u(B)h_u](\xi )=A(\xi )\text{ ,}$$ (53) \[compare with $`h_u(_u\xi )=[h_u(k)h_u](\xi )=d(\xi )`$\]. The explicit form of $`H(u)`$ follows from the explicit form of $`H`$ and its action on the vector field $`u`$ $$H(u)=(_uk)[g(u)]+k(_ug)(u)+k(g)(a)=_u[k(g)(u)]=_ua\text{ .}$$ (54) Now $`h_u[_u_u\xi ]`$ can be written in the form $$h_u(_u_u\xi )=h_u[\frac{l}{e}._uak(g)(\mathrm{\pounds }_\xi u)_u(\mathrm{\pounds }_\xi u)]+A(\xi )$$ (55) \[compare $`h_u(_u\xi )=h_u(\frac{l}{e}.a\mathrm{\pounds }_\xi u)+d(\xi )`$\]. The explicit form of $`A=h_u(B)h_u`$ can be found in an analogous way as the explicit form for $`d=h_u(k)h_u`$ in the expression for $`{}_{rel}{}^{}v`$. ### 3.2 Deformation acceleration, shear acceleration, rotation acceleration and expansion acceleration The covariant tensor $`A`$, named deformation acceleration tensor can be represented as a sum, containing three terms: a trace-free symmetric term, an antisymmetric term and a trace term $$A=_sD+W+\frac{1}{n1}.U.h_u$$ (56) where $$D=h_u(_sB)h_u\text{ , }W=h_u(_aB)h_u\text{ , }U=\overline{g}[_sA]=\overline{g}[D]$$ (57) $${}_{s}{}^{}B=\frac{1}{2}(B^{ij}+B^{ji}).e_i.e_j\text{ , }_aB=\frac{1}{2}(B^{ij}B^{ji}).e_ie_j\text{,}$$ (58) $${}_{s}{}^{}D=D\frac{1}{n1}.\overline{g}[D].h_u=D\frac{1}{n1}.U.h_u\text{ .}$$ (59) The tensor $`{}_{s}{}^{}D`$ is the shear acceleration tensor (shear acceleration), $`W`$ is the rotation acceleration tensor (rotation acceleration) and $`U`$ is the expansion acceleration invariant (expansion acceleration). Furthermore, every one of these quantities can be divided into three parts: torsion- and curvature-free acceleration, acceleration induced by the torsion and acceleration induced by the curvature. Let us now consider the representation of every acceleration quantity in its essential parts connected with its physical interpretation. The deformation acceleration tensor $`A`$ can be written in the following forms $$A=_sD+W+\frac{1}{n1}.U.h_u=A_0+G=_FA_0_TA_0+G\text{ ,}$$ (60) $$A=_sD_0+W_0+\frac{1}{n1}.U_0.h_u+_sM+N+\frac{1}{n1}.I.h_u\text{ ,}$$ (61) $$\begin{array}{c}A=_{sF}D_0+_FW_0+\frac{1}{n1}._FU_0.h_u(_{sT}D_0+_TW_0+\frac{1}{n1}._TU_{0.}h_u)+\\ +_sM+N+\frac{1}{n1}.I.h_u\text{ ,}\end{array}$$ (62) where $`A_0`$ is curvature-free deformation acceleration tensor, $`G`$ is the deformation acceleration tensor induced by the curvature, $`{}_{F}{}^{}A_{0}^{}`$ is the torsion-free and curvature-free deformation acceleration tensor, $`{}_{T}{}^{}A_{0}^{}`$ is the deformation acceleration tensor induced by the torsion. Every of these tensors could be represented in its three parts (the symmetric trace-free part, the antisymmetric part and the trace part) determining the corresponding shear acceleration, rotation acceleration and expansion acceleration: $$A_0=_FA_0_TA_0=_sD_0+W_0+\frac{1}{n1}.U_0.h_u\text{ ,}$$ (63) $${}_{F}{}^{}A_{0}^{}=_{sF}D_0+_FW_0+\frac{1}{n1}._FU_0.h_u\text{ , }_FA_0(\xi )=h_u(_\xi _{}a)\text{ , }$$ (64) $${}_{T}{}^{}A_{0}^{}=_{sT}D_0+_TW_0+\frac{1}{n1}._TU_0.h_u\text{ ,}$$ (65) $$G=_sM+N+\frac{1}{n1}.I.h_u=h_u(K)h_u\text{ ,}$$ (66) $$h_u([R(u,\xi )]u)=h_u(K)h_u(\xi )\text{ for }\text{ }\xi T(M)\text{ , }$$ (67) $$[R(u,\xi )]u=_u_\xi u_\xi _uu_{\mathrm{\pounds }_u\xi }u\text{ ,}$$ (68) $$K=K^{kl}.e_ke_l\text{ , }K^{kl}=R^k\text{ }_{mnr}.g^{rl}.u^m.u^n\text{ ,}$$ (69) The components $`R^k`$ <sub>mnr</sub> are the components of the contravariant Riemannian curvature tensor, $$K_a=K_a^{kl}.e_ke_l\text{}K_a^{kl}=\frac{1}{2}(K^{kl}K^{lk})\text{ , }K_s=K_s^{kl}.e_k.e_l\text{ , }K_s^{kl}=\frac{1}{2}(K^{kl}+K^{lk})\text{ ,}$$ (70) $${}_{s}{}^{}D=_sD_0+_sM\text{ , }W=W_0+N=_FW_0_TW_0+N\text{ ,}$$ (71) $$U=U_0+I=_FU_0_TU_0+I\text{ ,}$$ (72) $${}_{s}{}^{}M=M\frac{1}{n1}.I.h_u\text{ , }M=h_u(K_s)h_u\text{ , }I=\overline{g}[M]=g^{\overline{i}\overline{j}}.M_{ij}\text{ ,}$$ (73) $$N=h_u(K_a)h_u\text{ ,}$$ (74) $${}_{s}{}^{}D_{0}^{}=_{sF}D_0_{sT}D_0=_FD_0\frac{1}{n1}._FU_0.h_u(_TD_0\frac{1}{n1}._TU_0.h_u)\text{ ,}$$ (75) $${}_{s}{}^{}D_{0}^{}=_{sF}D_0_TD_0\frac{1}{n1}(_FU_0_TU_0)h_u\text{ ,}$$ (76) $${}_{s}{}^{}D_{0}^{}=D_0\frac{1}{n1}.U_0.h_u\text{ ,}$$ (77) $${}_{sF}{}^{}D_{0}^{}=_FD_0\frac{1}{n1}._FU_0.h_u\text{ , }_FD_0=h_u(b_s)h_u\text{ ,}$$ (78) $$b=b_s+b_a\text{ , }b=b^{kl}.e_ke_l\text{ , }b^{kl}=a^k\text{ }_{;n}.g^{nl}\text{ ,}$$ (79) $$a^k=u^k\text{ }_{;m}.u^m\text{ , }b_s=b_s^{kl}.e_k.e_l\text{ , }b_s^{kl}=\frac{1}{2}(b^{kl}+b^{lk})\text{ ,}$$ (80) $$b_a=b_a^{kl}.e_ke_l\text{ , }b_a^{kl}=\frac{1}{2}(b^{kl}b^{lk})\text{ ,}$$ (81) $${}_{F}{}^{}U_{0}^{}=\overline{g}[_FD_0]=g[b]\frac{1}{e}.g(u,_ua)\text{ , }g[b]=g_{kl}.b^{kl}\text{ ,}$$ (82) $${}_{sT}{}^{}D_{0}^{}=_TD_0\frac{1}{n1}._TU_0.h_u=_{sF}D_0_sD_0\text{ , }_TD_0=_FD_0D_0\text{ ,}$$ (83) $$U_0=\overline{g}[D_0]=_FU_0_TU_0\text{ , }_TU_0=\overline{g}[_TD_0]\text{ ,}$$ (84) $${}_{F}{}^{}W_{0}^{}=h_u(b_a)h_u\text{ , }_TW_0=_FW_0W_0\text{ .}$$ (85) Under the conditions $`\mathrm{\pounds }_\xi u=0`$ , $`\xi =\xi _{}=\overline{g}(h_u(\xi ))`$ , ($`l=0`$), the expression for $`h_u(_u_\upsilon \xi )`$ can be written in the forms $$g(_{rel}a)=h_u(_u_u\xi _{})=A(\xi _{})=A_0(\xi _{})+G(\xi _{})\text{ ,}$$ (86) $$g(_{rel}a)=h_u(_u_u\xi _{})=_FA_0(\xi _{})_TA_0(\xi _{})+G(\xi _{})\text{ ,}$$ (87) $$g(_{rel}a)=h_u(_u_u\xi _{})=(_{sF}D_0+_FW_0+\frac{1}{n1}._FU_0.g)(\xi _{})$$ $$(_{sT}D_0+_TW_0+\frac{1}{n1}._TU_0.g)(\xi _{})+(_sM+N+\frac{1}{n1}.I.g)(\xi _{})\text{ ,}$$ (88) which enable one to find a physical interpretation of the quantities $`{}_{s}{}^{}D`$,$`W`$,$`U`$ and of the quantities$`{}_{sF}{}^{}D_{0}^{}`$, $`{}_{F}{}^{}W_{0}^{}`$, $`{}_{F}{}^{}U_{0}^{}`$, $`{}_{sT}{}^{}D_{0}^{}`$, $`{}_{T}{}^{}W_{0}^{}`$, $`{}_{T}{}^{}U_{0}^{}`$, $`{}_{s}{}^{}M`$, $`N`$, $`I`$ contained in their structure . After the above consideration the following proposition can be formulated: Proposition 10. The covariant vector field $`g(_{rel}a)=h_u(_u_u\xi )`$ can be written in the form $$g(_{rel}a)=h_u[\frac{l}{e}._ua_{\mathrm{\pounds }_\xi u}u_u(\mathrm{\pounds }_\xi u)+T(\mathrm{\pounds }_\xi u,u)]+A(\xi )\text{ ,}$$ with $`A(\xi )=_sD(\xi )+W(\xi )+\frac{1}{n\mathrm{\hspace{0.17em}\hspace{0.17em}1}}.U.h_u(\xi )`$. For the case of affine symmetric (Levi-Civita) connection \[$`T(w,v)=0`$ for $``$ $`w,vT(M)`$ , $`T_{ij}^k=0`$, $`\mathrm{\Gamma }_{ij}^k=\mathrm{\Gamma }_{ji}^k`$\] and Riemannian metric ($`_vg=0`$ for $`vT(M)`$, $`g_{ij;k}=0`$), kinematic characteristics are obtained in $`V_n`$-spaces connected with the notion relative velocity , and the relative acceleration . For the case of affine non-symmetric connection \[$`T(w,v)0`$ for $``$ $`w,vT(M)`$ , $`\mathrm{\Gamma }_{jk}^i\mathrm{\Gamma }_{kj}^i`$\] and Riemannian metric kinematic characteristics are obtained in $`U_n`$-spaces . The shear, rotation and expansion accelerations induced by the torsion can be expressed by means of the kinematic characteristics of the relative velocity : (a) Shear acceleration tensor induced by the torsion $`{}_{sT}{}^{}D_{0}^{}`$ $${}_{sT}{}^{}D_{0}^{}=_TD_0\frac{1}{n1}._TU_0.h_u$$ $$\begin{array}{c}{}_{T}{}^{}D_{0}^{}=\frac{1}{2}[_sP(\overline{g})\sigma +\sigma (\overline{g})_sP]+\frac{1}{2}[Q(\overline{g})\omega +\omega (\overline{g})Q]+\\ +\frac{1}{n1}(\theta _1.\sigma +\theta ._sP)+\frac{1}{n1}(\theta _1^.+\frac{1}{n1}.\theta _1.\theta )h_u+_u(_sP)+\\ +\frac{1}{2}[_sP(\overline{g})\omega \omega (\overline{g})_sP]+\frac{1}{2}[Q(\overline{g})\sigma \sigma (\overline{g})Q]+\\ +\frac{1}{2e}[h_u(a)(g(u))(m+q)h_u+h_u((g(u))(m+q))h_u(a)]+\\ +\frac{1}{e}[_sP(a)g(u)+g(u)_sP(a)]+\\ +\frac{1}{2}[h_u(_u\overline{g})_sP+_sP(_u\overline{g})h_u]+\frac{1}{2}[h_u(_ug)QQ(_ug)h_u]\text{ .}\end{array}$$ (89) (b) Expansion acceleration induced by the torsion $`{}_{T}{}^{}U_{0}^{}`$ ($`\theta _1^.:=u\theta _1`$) $$_TU_0=\overline{g}[_sP(\overline{g})\sigma ]+\overline{g}[Q(\overline{g})\omega ]+\theta _1^.+\frac{1}{n1}.\theta _1.\theta +\frac{1}{e}.g(u,T(a,u))\text{ .}$$ (90) (c) Rotation acceleration tensor induced by the torsion $`{}_{T}{}^{}W_{0}^{}`$ $$\begin{array}{c}{}_{T}{}^{}W_{0}^{}=\frac{1}{2}[_sP(\overline{g})\sigma \sigma (\overline{g})_sP]+\frac{1}{2}[Q(\overline{g})\omega \omega (\overline{g})Q)]+\\ +\frac{1}{n1}(\theta _1.\omega +\theta .Q)+_uQ+\frac{1}{2}[_sP(\overline{g})\omega +\omega (\overline{g})_sP]+\\ +\frac{1}{2}[Q(\overline{g})\sigma +\sigma (\overline{g})Q]+\\ +\frac{1}{2e}[h_u(a)(g(u))(m+q)h_uh_u((g(u))(m+q))h_u(a)]+\\ +\frac{1}{e}[Q(a)g(u)g(u)Q(a)]+\\ +\frac{1}{2}[h_u(_u\overline{g})_sP_sP(_u\overline{g})h_u]+\frac{1}{2}[h_u(_u\overline{g})Q+Q(_u\overline{g})h_u]\text{ .}\end{array}$$ (91) The kinematic characteristics related to the notion of relative acceleration can be used in finding out their influence on the rate of change of the length of a contravariant vector field. ### 3.3 Relative acceleration and change of the change of the length of a contravariant vector field over a $`(L_n,g)`$-space The rate of change $`ul_\xi `$ of the length $`l_\xi =g(\xi ,\xi )^{\frac{1}{2}}`$ of a contravariant vector field $`\xi `$ along a contravariant vector field $`u`$ can be written in the form $`\pm 2.l_\xi .(ul_\xi )=(_ug)(\xi ,\xi )+2g(_u\xi ,\xi )`$. After a further differentiation of the last expression along the vector field $`u`$ we find the relation $$\begin{array}{c}\pm (ul_\xi )^2\pm l_\xi .(u(ul_\xi ))=\frac{1}{2}.(_u_ug)(\xi ,\xi )+g(_u_u\xi ,\xi )+\\ +2.(_ug)(_u\xi ,\xi )+g(_u\xi ,_u\xi )\text{ .}\end{array}$$ (92) Our aim is to represent the terms of the right side of the last expression by the use of the projective metrics of the vector field $`u`$ and the decomposition of the metric tensor $`g`$ in its trace free and trace parts. Since we already have the representation of this type for $`ul_\xi `$ (see the above section), we will find the corresponding representation for $`u(ul_\xi )=uul_\xi `$. After some computations, the following relations can be found: (a) Representation of $`(_u_ug)(\xi ,\xi )`$. By the use of the expressions: $$^s_uQ_u=uQ_u\frac{1}{n}.Q_u^2\text{ ,}$$ (93) $$_u_ug=^s_u^s_ug+\frac{1}{n}.[(uQ_u).g+2.Q_u.(^s_ug)]\text{ ,}$$ (94) $$(_u_ug)(\xi ,\xi )=(_u_ug)(\xi _{},\xi _{})+\frac{l^2}{e^2}.(_u_ug)(u,u)+2.\frac{l}{e}.(_u_ug)(u,\xi _{})\text{ ,}$$ (95) $$(_u_ug)(\xi _{},\xi _{})=(^s_u^s_ug)(\xi _{},\xi _{})+\frac{2}{n}.Q_u.(^s_ug)(\xi _{},\xi _{})+\frac{1}{n}.(uQ_u).l_\xi _{}^2\text{ ,}$$ (96) $$\frac{l^2}{e^2}.(_u_ug)(u,u)=\frac{l^2}{e^2}.(^s_u^s_ug)(u,u)+\frac{2}{n}.\frac{l^2}{e^2}.Q_u.(^s_ug)(u,u)+\frac{1}{n}.\frac{l^2}{e}.(uQ_u)\text{ ,}$$ (97) $$(_u_ug)(u,\xi _{})=(^s_u^s_ug)(u,\xi _{})+\frac{2}{n}.Q_u.(^s_ug)(u,\xi _{})\text{ ,}$$ (98) the expression for $`(_u_ug)(\xi ,\xi )`$ follows in the form $$\begin{array}{c}(_u_ug)(\xi ,\xi )=(^s_u^s_ug)(\xi _{},\xi _{})+2.\frac{l}{e}.(^s_u^s_ug)(u,\xi _{})+\\ +\frac{l^2}{e^2}.(^s_u^s_ug)(u,u)+\frac{2}{n}.Q_u.[(^s_ug)(\xi _{},\xi _{})+2.\frac{l}{e}.(^s_ug)(u,\xi _{})]+\\ +\frac{1}{n}.(uQ_u).(l_\xi _{}^2+\frac{l^2}{e})+\frac{2}{n}.\frac{l^2}{e^2}.Q_u.(^s_ug)(u,u)\text{ .}\end{array}$$ (99) (b) Representation of $`g(_u_u\xi ,\xi )`$: $$g(_u_u\xi ,\xi )=g(_{rel}a,\xi _{})+\frac{l}{e}.g(u,_u_u\xi )\text{ ,}$$ (100) (c) Representation of $`(_ug)(_u\xi ,\xi )`$. By the use of the expressions $$\begin{array}{c}(_ug)(_u\xi ,\xi )=\frac{l}{e^2}.(_ug)(u,u).g(_u\xi ,u)+\frac{1}{e}.(_ug)(u,\xi _{}).g(_u\xi ,u)+\\ +\frac{l}{e}.(_ug)(_{rel}v,u)+(_ug)(_{rel}v,\xi _{})\text{ , }_{rel}v=\overline{g}[h_u(_u\xi )]\text{ ,}\end{array}$$ (101) and the representation of $`_ug`$ by the use of $`{}_{}{}^{s}_{u}^{}g`$ and $`Q_u`$, we obtain for $`(_ug)(_u\xi ,\xi )`$ $$\begin{array}{c}(_ug)(_u\xi ,\xi )=\frac{1}{e}.(^s_ug)(u,\xi _{}).g(_u\xi ,u)+\frac{l}{e^2}.(^s_ug)(u,u).g(_u\xi ,u)+\\ +(^s_ug)(_{rel}v,\xi _{})+\frac{l}{e}.(^s_ug)(_{rel}v,u)+\frac{1}{n}.Q_u.[g(_{rel}v,\xi _{})+\frac{l}{e}.g(_u\xi ,u)]\text{ .}\end{array}$$ (102) (d) Representation of $`g(_u\xi ,_u\xi )`$: $$g(_u\xi ,_u\xi )=g(_{rel}v,_{rel}v)+\frac{1}{e}.[g(_u\xi ,u)]^2\text{ ,}$$ (103) where $$g(_u\xi ,_u\xi )=l_{_u\xi }^2\text{ , }g(_{rel}v,_{rel}v)=l_{{}_{rel}{}^{}v}^2\text{ , }e=g(u,u)=l_u^20\text{ ,}$$ (104) $$g(_u\xi ,u)=ul\frac{1}{n}.Q_u.l(^s_ug)(\xi ,u)g(\xi ,a)\text{ , }a=_uu\text{ .}$$ (105) By the use of all main results from (a) - (d) we find the explicit form for $`u(ul_\xi )`$ $$\begin{array}{c}\pm l_\xi .(u(ul_\xi ))=(ul_\xi )^2+\frac{1}{2}.[(^s_u^s_ug)(\xi _{},\xi _{})+\frac{l^2}{e^2}.(^s_u^s_ug)(u,u)]+\\ +\frac{l}{e}.(^s_u^s_ug)(u,\xi _{})+\frac{1}{n}.Q_u.[(^s_ug)(\xi _{},\xi _{})+2.\frac{l}{e}.(^s_ug)(u,\xi _{})+\\ +\frac{l^2}{e^2}.(^s_ug)(u,u)+2.g(_{rel}v,\xi _{})+2.\frac{l}{e}.g(_u\xi ,u)]+\frac{1}{2n}.(uQ_u).(l_\xi _{}^2+\frac{l^2}{e})+\\ +\frac{2}{e}.(^s_ug)(u,\xi _{}).g(_u\xi ,u)+2.\frac{l}{e^2}.(^s_ug)(u,u).g(_u\xi ,u)+\\ +2.(^s_ug)(_{rel}v,\xi _{})+2.\frac{l}{e}.(^s_ug)(_{rel}v,u)+g(_{rel}v,_{rel}v)+\frac{1}{e}.[g(_u\xi ,u)]^2+\\ +g(_{rel}a,\xi _{})+\frac{l}{e}.g(_u_u\xi ,u)\text{ .}\end{array}$$ (106) The last expression can help us to investigate the behavior of the length of a contravariant vector filed $`\xi `$ when transported along a non-null contravariant vector field $`u`$. There are several physical interpretation of such type of a transport. If we interpret the vector field $`u`$ as a time like vector field and as the 4-velocity of an observer in a $`(L_n,g)`$-space considered as a model of the space-time, then $`u(ul_\xi )=(d^2l_\xi )/(ds^2)`$ represents the acceleration acting on the length $`l_\xi `$ of the vector $`\xi `$ along the observers trajectory $`x^i(s)`$. The vector field $`u`$ can also be considered as a space like vector field with different from the above physical interpretation. Special case: $`g(u,\xi )=l:=0:\xi =\xi _{}`$. $$\begin{array}{c}\pm l_\xi _{}.(u(ul_\xi _{}))=(ul_\xi _{})^2+\frac{1}{2}.(^s_u^s_ug)(\xi _{},\xi _{})+\\ +\frac{1}{n}.Q_u.[(^s_ug)(\xi _{},\xi _{})+2.g(_{rel}v,\xi _{})]+\frac{1}{2n}.(uQ_u).l_\xi _{}^2+\\ +\frac{2}{e}.(^s_ug)(u,\xi _{}).g(_u\xi _{},u)+2.(^s_ug)(_{rel}v,\xi _{})+\\ +g(_{rel}v,_{rel}v)+\frac{1}{e}.[g(_u\xi _{},u)]^2+g(_{rel}a,\xi _{})\text{ .}\end{array}$$ (107) Special case: $`g(u,\xi )=l:=0:\xi =\xi _{}`$, $`_\xi g=0`$ for $`\xi T(M)`$ \[$`g_{ij;k}:=0`$\] ($`V_n`$\- or $`U_n`$-spaces). $$\pm l_\xi _{}.(u(ul_\xi _{}))=(ul_\xi _{})^2+g(_{rel}v,_{rel}v)+\frac{1}{e}.[g(_u\xi _{},u)]^2+g(_{rel}a,\xi _{})\text{ .}$$ (108) Special case: $`{}_{}{}^{s}_{u}^{}g:=0:_ug=\frac{1}{n}.Q_u.g`$ (Weyl’s space with torsion). $$\begin{array}{c}\pm l_\xi .(u(ul_\xi ))=(ul_\xi )^2+\frac{2}{n}.Q_u.[g(_{rel}v,\xi _{})+\frac{l}{e}.g(_u\xi ,u)]+\\ +\frac{1}{2n}.(uQ_u).(l_\xi _{}^2+\frac{l^2}{e})+g(_{rel}v,_{rel}v)+\frac{1}{e}.[g(_u\xi ,u)]^2+g(_{rel}a,\xi _{})+\\ +\frac{l}{e}.g(_u_u\xi ,u)\text{ .}\end{array}$$ (109) Special case: $`{}_{}{}^{s}_{u}^{}g:=0:_ug=\frac{1}{n}.Q_u.g`$, $`g(u,\xi )=l:=0`$, $`g(u,\eta )=\stackrel{~}{l}:=0`$ (Weyl’s space with torsion and orthogonal to $`u`$ vector fields $`\xi _{}`$ and $`\eta _{}`$). $$\begin{array}{c}\pm l_\xi _{}.(u(ul_\xi _{}))=(ul_\xi _{})^2+\frac{2}{n}.Q_u.g(_{rel}v,\xi _{})+\\ +\frac{1}{2n}.(uQ_u).l_\xi _{}^2+g(_{rel}v,_{rel}v)+\frac{1}{e}.[g(_u\xi _{},u)]^2+g(_{rel}a,\xi _{})\text{ .}\end{array}$$ (110) The representations of $`ul_\xi `$, $`u(ul_\xi )`$ and $`u[\mathrm{cos}(\xi ,\eta )]`$ can be useful tools for the considerations of the motion of physical systems with given dimensions in $`(L_n,g)`$-spaces. The induced by the torsion kinematic characteristics of the relative velocity $`{}_{rel}{}^{}v`$ and the relative acceleration $`{}_{rel}{}^{}a`$ have an equal position to the other kinematic characteristics induced by the curvature or by external forces. The relative acceleration between two (neighbor) points (particles) in an $`U_n`$-space $`(n=4)`$ is also related to quantities induced by non-autoparallel motions (under external forces), by the curvature and by the torsion. ### 3.4 Rate of change of the rate of change of the length of the vector field $`\xi _{}`$ over an $`U_n`$-space From the relation in $`U_n`$-spaces $$ul_\xi _{}=\pm \frac{1}{l_\xi _{}}.g(_u\xi _{},\xi _{})\text{ , }l_\xi _{}0\text{ ,}$$ after covariant differentiation along the vector field $`u`$, we obtain for $`U_n`$-spaces the relation $$u(ul_\xi _{})=\frac{1}{l_\xi _{}^2}.(ul_\xi _{}).g(_u\xi _{},\xi _{})\pm \frac{1}{l_\xi _{}}.g(_u_u\xi _{},\xi _{})\pm \frac{1}{l_\xi _{}}.g(_u\xi _{},_u\xi _{})\text{ ,}$$ where $$\frac{1}{l_\xi _{}}.g(_u\xi _{},\xi _{})=ul_\xi _{}\text{ , }\frac{1}{l_\xi _{}^2}.(ul_\xi _{}).g(_u\xi _{},\xi _{})=\frac{1}{l_\xi _{}^3}.[g(_{rel}v,\xi _{})]^2\text{ .}$$ By the use of the decomposition of $`_u_u\xi _{}`$, $$_u_u\xi _{}=\frac{\widehat{l}}{e}.u+_{rel}a\text{ , }\widehat{l}=g(_u_u\xi _{},u)\text{ ,}$$ under the condition $`g(u,\xi _{})=l=0`$, we obtain $$u(ul_\xi _{})=\frac{1}{l_\xi _{}^3}.[g(_{rel}v,\xi _{})]^2\pm \frac{1}{l_\xi _{}}.g(_{rel}a,\xi _{})\pm \frac{1}{l_\xi _{}}.\frac{\overline{l}^2}{e}\pm \frac{1}{l_\xi _{}}g(_{rel}v,_{rel}v)\text{ ,}$$ (111) where $$\text{ }\overline{l}=g(_u\xi _{},u)=_u[g(\xi _{},u)](_ug)(\xi _{},u)g(\xi _{},_uu)\text{ .}$$ From the last relation, it follows for $`U_n`$-spaces $`(_ug=0)`$ and for $`l=g(u,\xi _{})=0`$ that $$\text{ }\overline{l}=g(_u\xi _{},u)=g(\xi _{},_uu)=g(\xi _{},a)\text{ , }_uu=a\text{ .}$$ For $`(L_n,g)`$-spaces and $`l=0`$, it follows that $$\text{ }\overline{l}=g(_u\xi _{},u)=(_ug)(\xi _{},u)g(\xi _{},_uu)\text{ .}$$ Then we obtain for $`U_n`$-spaces $$u(ul_\xi _{})=\pm \frac{1}{l_\xi _{}}.\{g(_{rel}a,\xi _{})\frac{1}{l_\xi _{}^2}[g(_{rel}v,\xi _{})]^2+\frac{1}{e}.[g(\xi _{},a)]^2+g(_{rel}v,_{rel}v)\}\text{ . }$$ (112) Therefore, for $`u=\frac{d}{ds}`$ the relation follows $$\frac{d^2l_\xi _{}}{ds^2}=\pm \frac{1}{l_\xi _{}}.\{g(_{rel}a,\xi _{})+g(_{rel}v,_{rel}v)+\frac{1}{e}.[g(a,\xi _{})]^2\frac{1}{l_\xi _{}^2}[g(_{rel}v,\xi _{})]^2\}\text{ . }$$ (113) If $`\mathrm{\pounds }_u\xi _{}=0`$, then $`{}_{rel}{}^{}a=\overline{g}[A(\xi _{})]`$, $`{}_{rel}{}^{}v=\overline{g}[d(\xi _{})]`$. It is obvious that the existing kinematic terms induced by the torsion in $`A`$ and in $`d`$ could compensate the action of the kinematic terms induced by the curvature or by external forces in such a way that no change of rate of the rate of change of $`l_\xi _{}`$ could be registered. If $`{}_{rel}{}^{}a=0`$, $`{}_{rel}{}^{}v=0`$ and $`g(a,\xi _{})=0`$, then $`(d^2l_\xi _{}/ds^2)=0`$. ## 4 Conclusion The recent development of the mathematical models of the space-time shows the possible use of spaces with affine connections and metrics. Such type of spaces have the torsion as an intrinsic characteristic. Even in (pseudo) Riemannian spaces with torsion ($`U_n`$-spaces) the tensor of the torsion and the kinematic characteristics related to it could influence under certain conditions the effects in the space-time caused by the curvature tensor and by external forces. Especially, the torsion could lead up to possible new gravitational experiments. In such cases a consideration of the influence of the torsion as a characteristic of the space-time appears as a necessary step to a better understanding of the properties of the space-time related to the gravitational interaction between physical systems. Acknowledgments This work is supported in part by the National Science Foundation of Bulgaria under Grant No. 642.
warning/0005/gr-qc0005048.html	ar5iv	text	# Lemma: ## Lemma: A conformally flat metric with zero scalar curvature is analytic. ## Proof: By $`=0`$ and conformal flatness there holds $$D_k_{ij}=D_i_{kj}.$$ (32) Taking $`D^k`$ of Equ. (32), commuting derivatives and using $$_{kji\mathrm{}}=2_{i[k}g_{j]\mathrm{}}2_{\mathrm{}[k}g_{j]i}$$ (33) and the Bianchi identity we obtain $$\mathrm{\Delta }_{ij}=3_i{}_{}{}^{k}_{jk}^{}+g_{ij}_k\mathrm{}^k\mathrm{}.$$ (34) Now recalling that the Ricci tensor of a metric $`g_{ij}`$, when written in harmonic coordinates, gives an elliptic operator for $`\gamma _{ij}`$, we see that $$_{ij}=\sigma _{ij},$$ (35) together with (34) yields an elliptic system for $`(\gamma _{ij},\sigma _{ij})`$ with analytic coefficients. Thus, by , the Lemma is proved. In the case at hand we have that $``$ vanishes and, furthermore, that $`\gamma _{ij}`$ is conformally flat. In fact, from (29) we can directly obtain $$_{i[j;k]}=0.$$ (36) Combining the analyticity of $`\gamma _{ij}`$ with the system (30,31) we obtain analyticity of the pair $`(\mathrm{\Omega },\gamma _{ij})`$. The relativistic situation, given by Equ.’s (13,14), is similar, but more complicated. Instead of $`\mathrm{\Omega }=\stackrel{~}{U}^2/M^2`$ it is more convenient to take $$\mathrm{\Omega }=\frac{4}{M^2}\mathrm{sinh}^2\frac{\stackrel{~}{U}}{2}$$ (37) as conformal factor. As a field variable it is useful to take, instead of $`\mathrm{\Omega }`$, the quantity $`\rho `$, defined by $$\rho =\mathrm{tanh}^2\frac{\stackrel{~}{U}}{2}.$$ $`(37^{})`$ Then, again, it turns out that $`=0`$ and $$\mathrm{\Delta }\rho =\sigma :=\frac{3}{2}\frac{\rho _i\rho ^i}{\rho },$$ (38) whereas (28) gets replaced by $$\rho \left(1\rho \right)_{ij}=\rho _{ij}\frac{1}{3}\gamma _{ij}\mathrm{\Delta }\rho .$$ (39) Clearly the metric $`\gamma _{ij}`$ can not any longer be conformally flat. In fact, taking the “curl” of (39), Equ. (36) of the Newtonian theory gets replaced by $$\left(1\rho \right)_{i[j;k]}=2_{i[j}\rho _{k]}+\gamma _{i[j}_{k]\mathrm{}}\rho ^{\mathrm{}}.$$ (40) As a side remark we add,that, from the above equation it is not hard to show that $`\gamma _{ij}`$ is conformally flat iff $`(\mathrm{\Omega },\gamma _{ij})`$ comes from the Schwarzschild solution. For $`\sigma `$ we obtain, instead of (31), the relation $$\mathrm{\Delta }\sigma =3\rho \left(1\rho \right)^2+3_{ij}\rho ^i\rho ^j.$$ (41) Apropos Equ. (41) we wish to point out that, with the substitution $$\kappa :=\frac{\sigma ^{1/4}}{(1\rho )^{1/2}}$$ (42) and using (40), there results the identity $$\mathrm{\Delta }\kappa =\frac{1}{2}\kappa ^5+\frac{9}{16}\kappa ^7_{i[j;k]}^{i[j;k]},$$ (43) which plays an important role in black hole uniqueness theory (see e.g. ). In any case, we can now consider the system consisting of (35), (38), (41) and the equation of the form $$\mathrm{\Delta }\sigma _{ij}=\mathrm{}$$ (44) obtained by multiplying (40) by $`(1\rho )^1`$ (which is non-zero near $`\mathrm{\Lambda }`$) and taking $`D^k`$. This furnishes an elliptic system with analytic coefficients for the variables $`(\rho ,\gamma _{ij},\sigma ,\sigma _{ij})`$ (that-is-to-say if harmonic coordinates are used in Equ. (35)). We conclude that $`(\mathrm{\Omega },\gamma _{ij})`$ is analytic. In the Newtonian case, the multipole moments of the solution were nothing but the Taylor coefficients at $`\mathrm{\Lambda }`$ of the unphysical potential $`U`$ in the standard CC. In particular this implies that these moments determine the physical solution uniquely (which however already follows from the expansion (18)). In the relativistic case, and using the CC above based on $`\mathrm{\Omega }=\stackrel{~}{U}^2/M^2`$, it is not immediately clear whether the moments determine the solution. In fact, the expansions in (24,25) use a specific coordinate condition. In the unphysical picture one would like a “gauge invariant” definition of multipole moments. Luckily, Geroch in came up with such a definition. It goes as follows: Consider the following recursively defined set of tensor fields built from $`(U,\mathrm{\Omega },\gamma _{ij})`$ $`P`$ $`=`$ $`U`$ $`P_i`$ $`=`$ $`U_i`$ $`P_{ij}`$ $`=`$ $`U_{ij}{\displaystyle \frac{1}{3}}\gamma _{ij}\mathrm{\Delta }U`$ (45) $`P_{i_1\mathrm{}i_{s+1}}`$ $`=`$ $`𝒞\left[D_{i_{s+1}}P_{i_1\mathrm{}i_s}{\displaystyle \frac{s(2s1)}{2}}_{i_1i_2}P_{i_3\mathrm{}i_{s+1}}\right],`$ where $`𝒞`$ denotes the operation of taking the tracefree symmetric part. Now define $$M_{i_1\mathrm{}i_s}:=P_{i_1\mathrm{}i_s}\|_\mathrm{\Lambda }.$$ (46) Clearly, in the Newtonian case and using the standard CC for which $`_{ij}=0`$, this definition coincides with the standard one. As for the gauge dependence, when $$\overline{\mathrm{\Omega }}=\omega \mathrm{\Omega },\overline{U}=\omega ^{1/2}U,\overline{\gamma }_{ij}=\omega ^2\gamma _{ij},$$ (47) with $`\omega \|_\mathrm{\Lambda }=1`$ it turns out that $$\overline{M}_{i_1\mathrm{}i_s}=M_{i_1\mathrm{}i_s}+𝒞\underset{r=0}{\overset{s1}{}}\left(\begin{array}{c}s\\ r\end{array}\right)(2s1)\mathrm{}(2r+1)(1)^{rs}M_{i_1\mathrm{}i_s}b_{i_{r+1}}\mathrm{}b_{i_s},$$ (48) where $`b_i:=\frac{1}{2}\omega _i\|_\mathrm{\Lambda }`$. Note that only first derivatives of $`\omega `$ at $`\mathrm{\Lambda }`$ enter the transformation law (48). Furthermore the dependence of $`M_{i_1\mathrm{}i_s}`$ on $`b_i`$ is exactly the same as the change of the Newtonian moments under a translation of the physical Euclidean space by the vector $`b_i`$. We can now come back to the conformal gauge with $`\mathrm{\Omega }=\stackrel{~}{U}^2/M^2`$. Here it follows from (45) that $`UM`$ so that the moments are equivalent to $`M`$, together with $`𝒞`$ of the derivatives of $`_{ij}`$ at $`\mathrm{\Lambda }`$. But it is not difficult to infer from the above equations that these data determine $`(\mathrm{\Omega },\gamma _{ij})`$ uniquely whence $`(\stackrel{~}{U},\stackrel{~}{\gamma }_{ij})`$ is also determined. We remark, finally, that the $`M_{i_1\mathrm{}i_s}`$-terms in the physical expansion, which we have written out for $`k=2`$ in (24,25), coincide with the Geroch moments from above. There are two problems which the above analysis leaves open. The first one is to give an analyticity proof in the case where $`M`$ is zero. Although the “physical” expansion of $`(\stackrel{~}{U},\stackrel{~}{\gamma }_{ij})`$ seems to be completely insensitive to whether $`M`$ is nonzero or not, the unphysical situation, because of the choice of conformal factor makes vital use of the nonvanishing of $`M`$. The second open issue is that of convergence. Namely, given a set of moments $`M_{i_1\mathrm{}i_s}`$, $`s=0,1,\mathrm{}`$: what is the condition on the behaviour for large $`s`$ so that a static vacuum solution having these as multipole moments exists? The Newtonian situation is clear: the moments have to be such that the Taylor series having them as coefficients converges. In the relativistic case, the Taylor coefficients at $`\mathrm{\Lambda }`$ of the quantities $`(\mathrm{\Omega },\gamma _{ij})`$, say in Riemannian normal coordinates, depend on these moments in a nonlinear fashion. In particular it is not even obvious whether a solution exists for which only finitely many moments are nonzero.
warning/0005/hep-th0005180.html	ar5iv	text	# Untitled Document HUTP-00/A016 hep-th/0005180 Unifying Themes in Topological Field Theories Cumrun Vafa Jefferson Physical Laboratory Harvard University Cambridge, MA 02138, USA We discuss unifying features of topological field theories in 2, 3 and 4 dimensions. This includes relations among enumerative geometry (2d topological field theory) link invariants (3d Chern-Simons theory) and Donaldson invariants (4d topological theory). (Talk presented in conference on Geometry and Topology in honor of M. Atiyah, R. Bott, F. Hirzebruch and I. Singer, Harvard University, May 1999). May 2000 1. Introduction There has been many exciting interactions between physics and mathematics in the past few decades. Many of these developments on the physics side are captured by certain field theories, known as topological field theories. The correlation function of these theories compute certain mathematical invariants. Even though the original motivation for introducing topological field theories was to gain insight into these mathematical invariants, topological field theories have been found to be important for answers to many questions of interest in physics as well. The aim of my talk here is to explain certain connections that have been discovered more recently among various topological field theories. I will first briefly review what each one is, and then go on to explain some of the connections which has been discovered between them. The main examples of topological field theories that have been proposed appear in dimension two known as topological sigma models, in dimension three known as Chern-Simons theory and in dimension four known as topological Yang-Mills theory. The 2d and the 4d topological theory are related to an underlying supersymmetric quantum field theory, and there is no difference between the topological and standard version on the flat space. The difference between conventional supersymmetric theories and topological ones in these cases only arise when one considers curved spaces. In such cases the topological version, is a modified version of the supersymmetric theory on flat space where some of the fields have different Lorentz transformations properties (compared to the conventional choice). This modification of Lorentz transformation properties is also known as twisting, and is put in primarily to preserve supersymmetry on curved space. In particular this leads to having at least one nilpotent supercharge $`Q`$ as a scalar quantity, as opposed to a spinor, as would be in the conventional spin assignments. The physical observables of the topological theory are elements of the $`Q`$ cohomology. The path integral is localized to field configurations which are annihilated by $`Q`$ and this typically leads to some moduli problem which lead to mathematical invariants. In these theories the energy momentum tensor is $`Q`$ trivial, i.e., $$T_{\mu \nu }=\{Q,\mathrm{\Lambda }_{\mu \nu }\}$$ which (modulo potential anomalies) leads to the statement that the correlation functions are all independent of the metric on the curved space, thus leading to the notion of topological field theories (i.e. metric independence). The case of the 3d topological theory, is somewhat different. In this case, namely the case of Chern-Simons theory, one starts from an action which is manifestly independent of the metric on the 3 manifold, and thus topological nature of the field theory is manifest. The organization of this paper is as follows: In section 2 I briefly review each of the three classes of topological theories and discuss how in each case one goes about computing the correlation functions. In section 3 I discuss relations between 2d and 4d topological theories. In section 4 I discuss relations between 2d and 3d topological theories. 2. A Brief Review of Topological Field theories In this section I give a rather brief review of topological field theories in dimensions 2, 3 and 4. 2.1. TFT in $`d=2`$: Topological Sigma Models Topological sigma models are based on $`(2,2)`$ supersymmetric theories in 2 dimensions. These typically arise by considering supersymmetric sigma models on Kahler manifolds. In other words, we consider maps from 2 dimensional Riemann surfaces $`\mathrm{\Sigma }`$ to target spaces $`M`$ which are Kahler manifolds (together with fermionic degrees of freedom on the Riemann surface which map to tangent vectors on the Kahler manifold). The topological theory in this case localizes on holomorpic maps from Riemann surfaces to the target: $$X:\mathrm{\Sigma }M$$ $$\overline{}X=0$$ If we get a moduli space of such maps we have to evaluate an appropriate class over it. This class is determined by the topological theory one considers (for precise mathematical definitions see ). Also there are two versions of this topological theory: coupled or uncoupled to gravity. Coupling to gravity in this case means allowing the complex structure of $`\mathrm{\Sigma }`$ to be arbitrary and looking for holomorphic curves over the entire moduli space of curves. The case coupled to gravity is also sometimes referred to as ‘topological strings’. A particularly interesting class of sigma modelds both for the physics as well as for mathematics, corresponds to choosing $`M`$ to be a Calabi-Yau threefold, and considering topological strings on $`M`$. In this case the virtual dimension of the moduli space of holomorphic maps is zero. If this space is given by a number of points, the topological string amplitude just counts how many such points there are, weighted by $`e^{k(.)}`$ where $`k(.)`$ is the area of the holomorphic map (pullback of the Kahler form integrated over the surface) times $`\lambda ^{2g2}`$, where $`g`$ denotes the genus of the Riemann surface and $`\lambda `$ denotes the string coupling constant. More generally the space of holomorphic maps will involve a moduli space. This space comes equipped with a bundle with the same dimension as the tangent bundle (the existence of this bundle and the fact that its dimensions is the same as the tangent bundle follows from the fact that the relevent index is zero). Topological string computes the top Chern class of such bundles again weighted by $`e^{k(.)}\lambda ^{2g2}`$. These have to be defined carefully, due to singularities and issues of compactifications, and lead in general to rational numbers. The sum of these numbers for a given class $`vH_2(M,𝐙)`$ and fixed genus $`g`$, which we will denote by $`r_{g,v}`$, is known as Gromov-Witten invariant. We thus have the full partition function of topological string given by $$F(\lambda ,k)=\underset{vH_2(M,𝐙)}{}r_{g,v}e^{k(v)}\lambda ^{2g2}$$ here $`k`$ denotes the Kahler class of $`M`$. Even though the numbers $`r_{g,v}`$ are not integers, it has been shown, by physical arguments that $`F`$ can also be expressed in terms of other integral invariants . These integral invariants are related to certain aspects of cohomology classes of moduli of holomorphic curves together with flat bundles. These invariants associate for each $`vH_2(M,𝐙)`$ and eash positive (including 0) integer $`s`$ a number $`N_{v,s}`$ which denotes the ‘net’ number of BPS membranes with charge in class $`v`$ and ‘spin’ $`s`$ (for precise definitions see ). Then we have $$F(\lambda ,k)=\underset{n>0,vH_2(M,𝐙)}{}\frac{1}{n}N_{v,s}e^{nk(v)}[2Sin(n\lambda /2)]^{2s2}$$ For all cases checked thus far the Gromov-Witten invariants $`r_{g,v}`$ has been shown to be captured by these simpler integral invariants $`N_{v,s}`$ through the above map. In particular the checks made for constant maps and for contribution of isolated genus $`g`$ curves to all loops as well as some low genus computations for non-trivial CY 3-folds all support the above identification. Let us illustrate the above results in the case of a simple non-compact Calabi-Yau threefold, which we will later use in this paper. Consider the total space of the rank 2 vector bundle $`O(1)+O(1)𝐏^1`$. This space has vanishing $`c_1`$, and is a non-compact CY 3-fold. In this case the only BPS state is a membrane wrapping $`𝐏^1`$ once. This state has spin $`s=0`$. If we denote the area of $`𝐏^1`$ by $`t`$, then we have from (2.1) $$F=\underset{n>0}{}\frac{1}{n[2Sin(n\lambda /2)]^2}e^{nt}$$ For this particular case this has also been derived using the direct definition of topological strings in . 2.2. Topological Field Theory in 3d: Chern-Simons Theory The 3d topological theory we consider is Chern-Simons theory, which is given by the Chern-Simons action for a gauge field $`A`$: $$S_{CS}=\frac{k}{4\pi }_MTr[AdA+\frac{2}{3}A^3]$$ where $`M`$ is a 3-manifold and $`k`$ is an integer which is quantized in order for $`exp(iS)`$ to be well defined. As is clear from the definition of the above action, $`S`$ does not depend on any metric on $`M`$ and in this sense the theory is manifestly topological (i.e. metric independent)<sup>1</sup> At the quantum level there is a metric dependence which can be captured by a gravitational Chern-Simons term .. Thus the partition function of Chern-Simons theory gives rise to topological invariants for 3-manifolds for each group $`G`$. In other words $$Z_M(G)=exp(F_M(G))=𝒟Aexp[iS_{CS}]$$ where $`A`$ is a connection on $`M`$ for the gauge group $`G`$ and the above integral is over all inequivalent $`G`$-connections on $`M`$. The simplest way to compute such invariants is to use the relation between Hilbert space of Chern-Simons theory on a Riemann surface $`\mathrm{\Sigma }`$ and the chiral blocks of WZW model on $`\mathrm{\Sigma }`$ with group $`G`$ and level $`k`$. For example the partition function on $`S^3`$ can be computed by viewing $`S^3`$ as a sum of two solid 2-tori, which are glued along $`T^2`$ by an order 2 element of $`SL(2,𝐙)`$ on $`T^2`$. In this way the partition function gets identified with $$Z_{S^3}(G,k)=S_{00}(G,k)$$ where $`S_{00}=0\|S\|0`$ is a particular element of the order 2 operation of $`SL(2,𝐙)`$ on chiral characters, and is well studied in the context of WZW models. In particular for $`G=SU(N)`$ it is given by: $$Z_{S^3}(SU(N),k)=exp(F)=e^{i\pi N(N1)/8}\frac{1}{(N+k)^{N/2}}\sqrt{\frac{N+k}{N}}\underset{j=1}{\overset{N1}{}}(2\mathrm{sin}\frac{j\pi }{N+k})^{Nj}.$$ One can also consider knot invariants: Consider a knot $`\gamma `$ in $`M`$ and choose a representation $`R`$ of the group $`G`$ and consider the character of the holonomoy of $`A`$ around the knot $`\gamma `$, i.e. $$P[\gamma ,R]=Tr_RP\mathrm{exp}(i_\gamma A)$$ By the equation of motion for Chern-Simons theory, which leads to flatness of $`A`$, we learn that the above operator only depends on the choice of the knot type and not the actual knot<sup>2</sup> In the quantum theory one also needs to choose a framing for the knot.. One then obtains a knot invariant by computing the correlation function $$<\underset{i}{}P[\gamma _i,R_i]>=𝒟A\underset{i}{}P[\gamma _i,R_i]\mathrm{exp}(iS_{CS})$$ Again these quantities can be computed by the braiding properties of chiral blocks in 2 dimensional WZW models and leads in particular to HOMFLY polynomial invariants for the knots. 2.3. Topological Field Theories in 4 Dimensions If one consider $`N=2`$ supersymmetric Yang-Mills, with an unconventional spin assignments, one finds a topological field theory. The partition function is localized on the moduli space of instantons and the observables of this theory are given by intersection theory on the moduli space of instantons. More precisely each $`d`$-cycle on the four manifold $`M`$ will lead to a $`4d`$ cohomology element on the moduli space of instantons (obtained by integrating out $`\mathrm{Tr}FF`$ over the corresponding cycle on the universal moduli space of instantons), and the wedging of the cohomology classes gives rise to the observables in Donaldson theory. This does not depend on the metric in $`M`$ (except when $`b_2^+(M)=1`$) but will depend on the choice of smooth structure on $`M`$. The computations in this case can be done for many choices of $`M`$ by finding an equivalence of this theory and a simpler abelian theory. In this case studying the moduli space of non-abelian instantons gets replaced with the study of an abelian system known as the Seiberg-Witten equation. The relevant geometry for the case of $`SU(N)`$ Yang-Mills is captured by a certain geometric data related to a Jacobian variety over an $`N1`$ dimensional family of genus $`N1`$ Riemann surface, known as Seiberg-Witten geometry . For topological field theory aspects and how the Seiberg-Witten geometry leads to computation of the topological correlation functions see . There is another topological theory in 4 dimensions which has been studied and is related to twisting the maximal supersymmetric gauge theory in 4 dimensions. This theory computes the Euler characteristic of moduli space of instantons. In particular for each group $`G`$ and each complex parameter $`q`$ one considers $$Z_M(G)=q^{c(M,G)}\underset{k}{}q^k\chi (_k)$$ for some universal constant $`c`$ (depending on $`M`$ and $`G`$), where $`k`$ denotes the instanton number and $`\chi (_k)`$ denotes the euler characteristics (of a suitable resolution and compactification) of $`_k`$, the moduli space of anti-self dual $`G`$-connections with instanton number $`k`$ on $`M`$. Moreover, according to Montonen-Olive duality conjecture one learns that the above partition function is expected to be modular with respect to some subgroup of $`SL(2,𝐙)`$ acting in the standard way on $`\tau `$ where $`q=exp(2\pi i\tau )`$. For certain $`M`$ (such as $`K3`$ ) the above partition function has been computed and is shown to be modular in a striking way. For recent mathematical discussion on this see and references therein. 3. Connections between 2d $``$ 4d TFT’s There are three different links between 4 dimensional and 2 dimensional TFT’s that I would like to discuss. In all three links the common theme is that the moduli space of instantons are mapped to moduli space of holomoprhic curves on appropriate spaces. 3.1. Topological Reduction of $`4d`$ to $`2d`$ The simplest link between the two theories involves studying the 4d TFT on a geometry involving the product of two Riemann surfaces $`\mathrm{\Sigma }_1\times \mathrm{\Sigma }_2`$, which was studied in . In the limit where $`\mathrm{\Sigma }_1`$ is small compared to $`\mathrm{\Sigma }_2`$ one obtains an effective theory on $`\mathrm{\Sigma }_2`$ which is the topological sigma model with target space given by moduli space of flat connections on $`\mathrm{\Sigma }_1`$, in case one considers $`N=2`$ topological field theories in 4 dimensions or the Hitchin space associated with $`\mathrm{\Sigma }_1`$ if one considers $`N=4`$ topological field theories. This is natural to expect because studying light supersymmetric modes in either case gives rise to the corresponding space of solutions, which thus behaves from the viewpoint of the space $`\mathrm{\Sigma }_2`$ as a target space. In particular the moduli space of 4d instantons get mapped to moduli space of holomorphic maps for these target spaces. Thus quantum cohomology rings of moduli of flat connections on a Riemann surface, which are encoded in 2d topological correlation functions capture the corresponding topological correlation functions of the 4 dimensional $`N=2`$ theory. Similarly in the $`N=4`$ case the reduction to 2 dimensions yields a sigma model on the Hitchin space (which can also be viewed as a Jacobian variety). In this context the Montonen-Olive duality of $`N=4`$ theory gets mapped to mirror symmetry of this 2d sigma model (by a fiberwise application of T-duality to Jacobian fibers). 3.2. A more subtle 2d $``$ 4d link For the $`N=2`$ topologically twisted theory, an important role is played by the Seiberg-Witten geometry, which is an abelian simplification of the non-abelian gauge theory. This geometry is a quantum deformation of the classical one, due to pointlike four dimensional instantons. This geometry was first conjectured based on consistency with various properties of $`N=2`$ quantum field theories and its deformation to $`N=1`$ quantum field theories with mass gap, where plausible properties of $`N=1`$ theories were assumed. With the recent advances in our understanding of string theory, the same 4 dimensional gauge theories have been obtained by considering particular geometries where strings propagate in. This procedure is known as geometric engineering of QFT’s (see and references therein). These geometries involve a non-compact Calabi-Yau threefold geometry which is a blow up of a geometry with some loci of A-D-E singularities (locally modelled by $`𝐂^2/G`$ where $`G`$ is a discrete subgroup of $`SU(2)`$), giving rise to the corresponding gauge theory in 4 dimensions. Depending on the detailed structure of singularities one can obtain various interesting gauge groups and various matter representations. It turns out that in this description of gauge theory, the guage theory instantons are mapped to stringy instantons, which are just worldsheet instantons. Thus being able to compute worldsheet instantons, i.e. counting of holomorphic curves in these target geometries, captures the geometry of 4 dimensional gauge theory instantons. Counting of holomrophic curves is precisely what the (A-model) topological string computes and thus in this way the geometry of vacua of 4 dimensional gauge theory gets mapped to solving topological amplitudes in 2d. This in turn can be done by using (local) mirror symmetry. For a physical derivation of mirror symmetry and some references on this subject see the recent work . In this way, not only the Seiberg-Witten geometry has been rederived, but also other geometries which describe other $`N=2`$ systems with various kinds of gauge groups and intricate matter representations have been obtained . 3.3. $`N=4`$ Yang-Mills on elliptic surfaces and $`2d`$ topological theories If we consider an $`N=4`$ supersymmetric $`SU(N)`$ topological theory on an elliptic surface, with base $`B`$, the stable bundles get mapped to spectral covers of $`B`$ on a dual elliptic surface $`M`$ (where the Kahler class of the elliptic fiber is inverted). This uses the fact that in the limit of small tori, the stable bundles become flat fiberwise and flat bundles on tori are related to points on the dual tori. See for a discussion of how this arises. In particular a rank $`N`$ stable bundle with instanton number $`k`$ gets mapped to a spectral curve which is a holomorphic curve wrapping the base $`N`$ times and the elliptic fiber $`k`$ times. Thus the topological $`N=4`$ amplitude on $`M`$, denoted by $`Z_M(SU(N))`$ which computes the Euler characteristic of moduli space of $`SU(N)`$ instantons on $`M`$ gets mapped to computing Euler characteristic of moduli space of holomorphic curves (together with a flat bundle) which in turn is captured by genus zero topological string amplitudes, and can be computed using mirror symmetry. This idea has been implemented in great detail for the case of rational elliptic surface (also known as “half $`K3`$”) . The results for the case of rank $`2`$ and its implications for the Euler characteristic of moduli space of instantons on rational elliptic surface has been confirmed using rigorous mathematical methods in . 4. Connections between 2d $``$ 3d TFTs Over two decades ago ’t Hooft conjectured that $`SU(N)`$ gauge theories with large $`N`$ look alot like string theories. In particular the partition function for these theories can be organized in terms of Riemann surfaces where each Riemann surface is weighted with $`N^\chi `$ where $`\chi `$ denotes the Euler characteristic of the Riemann surface. In particular the low genera dominate in the large $`N`$ limit. The weight factor $`N^\chi `$ follows simply from the combinatorics of Feynman diagrams, and the Riemann surface can be identified with the Feynamn diagrams where the would be holes have been filled. The main difficulty in the conjecture of ‘t Hooft is to identify precisely which string theory one obtains. In the past few years for serveral interesting gauge theories and in particular some in 4 dimensions the corresponding string theory has been identified . Even though it has not been possible to actually compute the string theory amplitudes in these cases, due to the complicated background strings propagate in, there has been mounting evidence for the validity of the identification. One would like to have a similar conjecture in a setup which is more computable. An ideal setup for this is topological guage theories, and in particular the topological Chern-Simons theory. If we consider $`SU(N)`$ Chern-Simons theory on $`S^3`$ in the limit of large $`N`$, one could hope to get a string theory. It has been conjectured in that this is indeed the case. In particular it has been conjectured that $`SU(N)`$ Chern-Simons theory at level $`k`$ on $`S^3`$ is equivalent to topological string with target being a non-compact Calabi-Yau threefold which is the total space of $`O(1)+O(1)𝐏^1`$, where the (complexified) size of $`𝐏^1`$ is given by $`t=2\pi iN/(k+N)`$ and the string coupling constant $`\lambda =\frac{2\pi i}{N+k}`$. This is a natural conjecture in the following sense: The Chern-Simons theory on $`S^3`$ can itself be viewd as an open string theory with target $`T^{}S^3`$ . By open string we mean considering Riemann surfaces with boundaries, where the boundaries are mapped to $`S^3`$. The geometry $`O(1)+O(1)𝐏^1`$ can be obtained from the $`T^{}S^3`$ geometry by shrinking $`S^3`$ to zero size and blowing $`𝐏^1`$ instead. This kind of transition is also very similar to what is observed to happen in the other cases where large $`N`$ string theory description was discovered . In fact one can determine the map of the parameters $`t`$ and $`\lambda `$ given above using this picture (and recalling the metric dependence anomaly in Chern-Simons theory). This conjecture has been checked at the level of the partition function (which we have briefly reviewed for both the Chern-Simons theory on $`S^3`$ and for $`O(1)+O(1)𝐏^1`$ in section 2). The implications of this conjecture for knot invariants has been explored in and provides a reformulation of knot invariants in terms of integral invariants which again capture the degeneracy of spectrum of (BPS) particles in the corresponding string theory. This involves considering a Largrangian submanifold which intersects $`T^{}S^3`$ along the knot and following it through the transition to $`O(1)+O(1)𝐏^1`$ where it corresponds to a Lagrangian submanifold. The corresponding computation on the topological string side will now involove Riemann surfaces with boundaries, where the boundary can lie on this Lagrangian submanifold in $`O(1)+O(1)𝐏^1`$. The results for the unknot as well as the integrality properties of the torus knots are in perfect agreement with the conjecture. 5. 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warning/0005/nucl-th0005032.html	ar5iv	text	# Non-Abelian Energy Loss at Finite Opacity ## Abstract A systematic expansion in opacity, $`L/\lambda `$, is used to clarify the non-linear behavior of induced gluon radiation in quark-gluon plasmas. The inclusive differential gluon distribution is calculated up to second order in opacity and compared to the zeroth order (factorization) limit. The opacity expansion makes it possible to take finite kinematic constraints into account that suppress jet quenching in nuclear collisions below RHIC ($`\sqrt{s}=200`$ AGeV) energies. PACS numbers: 12.38.Mh; 24.85.+p; 25.75.-q Introduction. The production of high transverse momentum jets in QCD is always accompanied by gluon showers. For jets produced inside nuclei via $`e+A`$ or quark-gluon plasmas via $`A+A`$, final state interactions of jet and radiated gluons induce further radiation and also broaden the gluon shower. The non-abelian radiative energy loss in a medium is expected to be observable as “jet quenching”. The observable consequences in nuclear collisions should be seen as a suppression of the high $`p_{}`$ tails of single hadron distributions and a broadening of the jet cone . Non-abelian energy loss in pQCD can be calculated analytically in two limits. In both cases, the energy of the leading parton is assumed to be large enough that its angular deflection can be neglected. One analytic limit applies to thin plasmas where the mean number, $`\overline{n}=L/\lambda `$, of jet scatterings is small . The other limit is the thick plasma one , where $`\overline{n}1`$. This mean number is a measure of the opacity or geometrical thickness of the medium: $$\overline{n}=\frac{N\sigma _{el}}{A_{}}=𝑑z(\sigma _{el}\rho )\frac{dN}{dy}\frac{\sigma _{el}}{2\pi R_G^2}\mathrm{log}\frac{R_G}{\tau _0}.$$ (1) For a constant (box) density $`\rho =N/(LA_{})`$, $`L`$ is the target thickness and $`\lambda =1/(\rho \sigma _{el})`$ is the average mean free path. For a sharp cylinder geometry we can interpret $`L=1.2A^{1/3}R_s`$ as the nuclear radius. For a more realistic 3+1D expanding Gaussian cylinder with $`\rho (𝐱_{},\tau )=(\tau _0/\tau )\rho _0\mathrm{exp}\left((𝐱_{}^2+\mathrm{\Delta }\tau ^2)/R_G^2\right)I_0(2\|𝐱_{}\|\mathrm{\Delta }\tau /R_G^2)`$, $`L`$ is replaced by the equivalent rms Gaussian transverse radius $`R_G=0.75A^{1/3}`$ fm. The rightmost form in (1) is obtained by averaging over the expanding Gaussian cylinder. Here $`\tau _0`$ is the formation time of the plasma. At RHIC energies ($`\sqrt{s}200`$ AGeV) the expected rapidity density of the gluons is $`dN/dy1000`$ for $`A=200`$. With an elastic cross section $`\sigma _{el}2`$ mb and a plasma formation time $`0.5`$ fm/c, we obtain $`\overline{n}4`$. This suggests that neither analytic approaches may be strictly applicable in practice. However, because of the non-abelian analog of the Landau-Pomeranchuk-Migdal (LPM) effect, the radiation intensity angular distribution and total energy loss are controlled by the combined effect of the number of scatterings $`\overline{n}=L/\lambda `$ and a formation probability $`p_f=L/l_f=L\mu ^2/(2xE)`$ of the gluon in the medium. In the induced energy loss was shown to be nonlinear in the nuclear thickness, $`\mathrm{\Delta }E\mu ^2L^2/\lambda `$ assuming the $`\overline{n}1`$ limit. In , we showed that the angular pattern was even more strongly nonlinear in the exclusive (tagged target) case. We show below that in the inclusive case, important features of the radiation pattern are in fact governed by the product $`\overline{n}p_f=L^2(\mu ^2/k^+\lambda )`$. Because of this, even the first order in opacity reproduces the $`L^2`$ dependence of the energy loss and the range of applicability of the finite opacity expansion derived below may extend to realistic targets. Another important motivation for our approach is the apparent absence of jet quenching observed at SPS energies . As we find below, the opacity expansion shows that finite kinematic constraints suppress greatly the non-abelian energy loss at SPS energies. The work reported in this Letter extends Ref. by including virtual corrections necessary to compute inclusive (vs. tagged) jet quenching. Detailed derivation of these results is given in Ref. . Hard Gluon Distribution. At zeroth order in opacity, the gluon emission from the hard production vertex is approximately given by $$x\frac{dN^{(0)}}{dxd𝐤_{}^2}=\frac{C_R\alpha _s}{\pi }\left(1x+\frac{x^2}{2}\right)\frac{1}{𝐤_{}^2},$$ (2) where $`x=k^+/E^+\omega /E`$, and $`C_R`$ is the Casimir of the (spin $`1/2`$) jet in the $`d_R`$ dimensional color representation. For a spin 1 jet the gluon splitting function must be substituted above. The differential energy distribution outside a cone defined by $`𝐤_{}^2>\mu ^2`$ is given by $$\frac{dI^{(0)}}{dx}=\frac{2C_R\alpha _s}{\pi }\left(1x+\frac{x^2}{2}\right)E\mathrm{log}\frac{\|𝐤_{}\|_{\mathrm{max}}}{\mu },$$ (3) where the upper kinematic limit is $$𝐤_{\mathrm{max}}^2=\mathrm{min}[4E^2x^2,4E^2x(1x)].$$ (4) The energy loss outside the cone is then given by $$\mathrm{\Delta }E^{(0)}=\frac{4C_R\alpha _s}{3\pi }E\mathrm{log}\frac{E}{\mu }$$ (5) in the leading $`\mathrm{log}(E/\mu )`$ approximation (LLA). While this overestimates the radiative energy loss in the vacuum (self-quenching), it is important to note that $`\mathrm{\Delta }E^{(0)}/E50\%`$ is typically much larger than the medium induced energy loss. Opacity Expansion. To compute the induced radiation, we assume as in that the quark gluon plasma can be modeled by $`N`$ well-separated, i.e. $`\lambda 1/\mu `$, color-screened Yukawa potentials. In contrast to , we consider here the inclusive gluon distribution induced by a finite medium which remains unobserved. Without target tagging, we must add double Born (Virtual) amplitudes $`V_i`$ to the real (Direct) amplitudes $`D_i`$ derived in . We denote by $`G_0`$ the basic eikonal hard emission amplitude $$G_0=2ig_sϵ_{}𝐇e^{i\omega _0z_0}c,$$ (6) where we use the shorthand notation of , $`𝐇𝐤_{}/𝐤_{}^2,\omega _0𝐤_{}^2/2xE`$. The hard parton originates at $`z_0`$, and $`c`$ denotes the color matrix $`T_c`$ in its $`d_R`$ dimensional representation. In Ref. we developed an iterative procedure to generate from this amplitude the sum of all amplitudes to any order in opacity including both direct and virtual terms. To first order in opacity it is sufficient to consider the classes of amplitudes denoted $`G_0D_1`$ (sum of three graphs) and $`G_0V_1`$ (sum of four graphs) resulting from direct and virtual terms due to a scattering at position $`z_1`$. The cross section for the induced radiation consists of $`3^2`$ real and $`2\times 4`$ double Born contributions that sum to a simple result $`\mathrm{Tr}(G_0D_1D_1^{}G_0^{}+(G_0V_1G_0^{}+\mathrm{h}.\mathrm{c}.))`$ (7) $`=4g_s^2C_RC_Ad_R\left[(2𝐂_1𝐁_1)\left(1\mathrm{cos}(\omega _1\mathrm{\Delta }z_1)\right)\right],`$ (8) with $`𝐂_m`$ and $`\omega _m`$ obtained from $`𝐇`$ and $`\omega _0`$ through the substitution $`𝐤_{}𝐤_{}𝐪_m`$ and $`𝐁_m𝐇𝐂_m`$. Note that the interaction of the radiated gluon in the medium brought in a factor $`C_A=N_c`$. Unlike the tagged case , the above first order correction to the hard (factorization) distribution in Eq.(2) has no simple classical cascade limit. In the tagged case, the result contained terms with color factors $`C_R^2d_R`$ as well as $`C_RC_Ad_R`$ that were easily interpretable in terms of jet rescattering and induced radiation. In the inclusive case, the virtual corrections lead to long range non-local interference effects that have no classical analog. On the other hand, a remarkable “color triviality” is found at the inclusive level where all jet rescattering effects cancel . As show in detail in Ref. , the color factor at order $`n=2`$ is simply $`C_RC_A^2d_R`$, in contrast to the vastly more complicated color structure of the tagged case . The non-abelian LPM effect is seen in Eq.(8) as arising from the gluon formation factor $$\mathrm{\Phi }(\mathrm{\Delta }z_1)=1\mathrm{cos}(\omega _1\mathrm{\Delta }z_1),$$ (9) where $`\mathrm{\Delta }z_1=z_1z_0`$. This must also be averaged over the longitudinal target profile, $`n(z)`$. For a box density of thickness $`L=1.2A^{1/3}`$, $`n=\theta (Lz)/L`$. For analytic simplicity we take instead an exponential form $`n_e(z)=e^{z/L_e}/L_e`$. In order to compare results for the two cases, we must require identical mean target depths, i.e. $`L_e=L/2`$. With $`n_e(z)`$ the ensemble averaged formation factor is $`{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{2d\mathrm{\Delta }z_1}{L}}e^{\frac{2\mathrm{\Delta }z_1}{L}}\mathrm{\Phi }(\mathrm{\Delta }z_1)={\displaystyle \frac{(𝐤𝐪_1)_{}^4L^2}{16x^2E^2+(𝐤𝐪_1)_{}^4L^2}}.`$ (10) (11) The formation probability in this case is controlled by simple Lorentzian factors. Averaging over the momentum transfer $`𝐪_\mathrm{𝟏}`$ via the color Yukawa potential leads finally to the gluon double differential distribution $`x{\displaystyle \frac{dN^{(1)}}{dxd𝐤_{}^2}}`$ $`=`$ $`x{\displaystyle \frac{dN^{(0)}}{dxd𝐤_{}^2}}{\displaystyle \frac{L}{\lambda _g}}{\displaystyle _0^{q_{\mathrm{max}}^2}}d^2𝐪_1{\displaystyle \frac{\mu _{eff}^2}{\pi (𝐪_1^2+\mu ^2)^2}}`$ (13) $`{\displaystyle \frac{2𝐤_{}𝐪_1(𝐤𝐪_1)_{}^2L^2}{16x^2E^2+(𝐤𝐪_1)_{}^4L^2}}.`$ where the opacity factor $`L/\lambda _g=N\sigma _{el}^{(g)}/A_{}`$ arises from the sum over the $`N`$ distinct targets. Note that the radiated gluon mean free path $`\lambda _g=(C_A/C_R)\lambda `$ appears rather than the jet mean free path. It is the color triviality of Eq. (8) that allows us to absorb $`C_R`$ factor into the hard distribution, Eq. (2), and $`C_A`$ factor into $`\lambda _g`$. The upper kinematic bound on the momentum transfer $`q_{\mathrm{max}}^2=s/43E\mu `$, ($`1/\mu _{eff}^2=1/\mu ^21/(\mu ^2+q_{\mathrm{max}}^2)`$). For SPS and RHIC energies, this finite limit cannot be ignored as we show below. The second order contribution in opacity requires a more complex calculation involving the sum of $`7^2`$ direct and $`2\times 96`$ virtual terms as discussed in . The result is $`\mathrm{Tr}(G_0D_1D_2D_2^{}D_1^{}G_0^{}+(G_0V_1D_2D_2^{}G_0^{}+\mathrm{h}.\mathrm{c}.)`$ (14) $`+(G_0D_1V_2D_1^{}G_0^{}+\mathrm{h}.\mathrm{c}.)+(G_0V_1V_2G_0^{}+\mathrm{h}.\mathrm{c}.))`$ (15) $`=\mathrm{\hspace{0.17em}4}g_s^2C_RC_A^2d_R[\mathrm{\hspace{0.17em}2}𝐂_1𝐁_1(1\mathrm{cos}(\omega _1\mathrm{\Delta }z_1))`$ (16) $`+\mathrm{\hspace{0.17em}2}𝐂_2𝐁_2(\mathrm{cos}(\omega _2\mathrm{\Delta }z_2)\mathrm{cos}(\omega _2(\mathrm{\Delta }z_1+\mathrm{\Delta }z_2))`$ (17) $`\mathrm{\hspace{0.17em}2}𝐂_{(12)}𝐁_2\left(\mathrm{cos}(\omega _2\mathrm{\Delta }z_2)\mathrm{cos}(\omega _{(12)}\mathrm{\Delta }z_1+\omega _2\mathrm{\Delta }z_2)\right)`$ (18) $`\mathrm{\hspace{0.17em}2}𝐂_{(12)}𝐁_{2(12)}(1\mathrm{cos}(\omega _{(12)}\mathrm{\Delta }z_1))],`$ (19) where with $`𝐂_{(mn)}`$ and $`\omega _{(mn)}`$ obtained from $`𝐇`$ and $`\omega _0`$ through the substitution $`𝐤_{}𝐤_{}𝐪_m𝐪_n`$ and $`𝐁_{m(nl)}𝐂_m𝐂_{(nl)}`$. Note the color triviality of this inclusive result in contrast to the tagged target case in . Both Eqs. (8,19) actually hold for either quark or gluon jets in the small $`x`$ approximation. We note further that in the two important limits $`\|𝐤_{}\|0`$ and $`\|𝐤_{}\|\mathrm{}`$ the distributions vanish for any fixed $`x`$ due to azimuthal angular averaging $$x\frac{dN^{(n)}}{dxd𝐤_{}^2}_0^{2\pi }\frac{d\varphi }{2\pi }𝐤_{}𝐪_m=\mathrm{\hspace{0.17em}0}.$$ (20) This is clearly seen in Fig. 1. The average over the scattering points, $`z_1,z_2`$ in the second order case is performed with the exponential density profile as follows: $$\mathrm{}=_0^{\mathrm{}}\frac{3d\mathrm{\Delta }z_1}{L}_0^{\mathrm{}}\frac{3d\mathrm{\Delta }z_2}{L}e^{\frac{3(\mathrm{\Delta }z_1+\mathrm{\Delta }z_2)}{L}}\mathrm{},$$ (21) where $`\mathrm{\Delta }z_i=z_iz_{i1}`$. We must use $`L_e=L/(n+1)`$ for n<sup>th</sup> order in opacity in order to insure that the first moment of the m<sup>th</sup> scattering center is identical to that in a box distribution ($`z_m_n=mL/(n+1)`$) as discussed in . Numerical results comparing the first and second order in opacity corrections to the hard distribution Eq. (2) are illustrated in Figs. 1 and 2. We consider a 50 GeV quark jet in a medium with $`\lambda _g=1`$ fm. A screening scale $`\mu =0.5`$ GeV and $`\alpha _s=0.3`$ were assumed. The “angular” distribution in Fig. 1 shows that the first order angular distribution is wider than the medium independent hard ($`1/𝐤_{}^2`$) distribution. The second order corrections redistributes the gluons further. In Fig. 2 the $`𝐤_{}`$ integrated contributions to the gluon intensity, $`dI/dx`$, are shown for the same conditions as in Fig. 1. The induced intensity is concentrated at small $`x`$ in contrast to the relatively constant intensity originating from the hard “self-quenching” term (3). The long $`1/x`$ tail contributes however a logarithmic factor $`\mathrm{log}(E/\mu )`$ as we discuss further below. The second order correction suppresses the intensity at small $`x`$ and enhances it somewhat at higher $`x0.2`$. We note that these two effects tend to cancel in the integrated energy loss as seen in Fig. 3. Radiation Intensity and Energy loss. To gain analytic insight into the above numerical results, we consider the first order induced radiation intensity $`dI^{(1)}/dx`$ in the approximation that $`𝐤_{\mathrm{max}}^2=\mathrm{}`$. This allows us to change variables $`𝐪_{}^{}𝐤_{}𝐪_1`$ in Eq. (13) and express the integrand in the azimuthal $`\varphi `$ integral as a partial derivative with respect to $`𝐤_{}^2`$. The remaining $`𝐪_{}^{}`$ integral can be performed then analytically, resulting in $$\frac{dI^{(1)}}{dx}=\frac{C_R\alpha _s}{\pi }\left(1x+\frac{x^2}{2}\right)E\frac{L}{\lambda _g}f(\gamma ,\delta ),$$ (22) where $`\gamma =L\mu ^2/(4xE)`$ is a formation probability and $`\delta =\mu /E`$ is a measure of the importance of finite kinematics. The formation function is given by $$f(\gamma ,\delta )=\frac{\gamma \left(2\mathrm{tan}^1\frac{\gamma }{\delta }+\gamma \mathrm{log}\frac{\gamma ^2+\delta ^2}{1+\delta ^2}\right)}{(1+\gamma ^2)},$$ (23) which in the $`\delta 0`$ and $`\gamma 1`$ limit reduces to a simple form $`f(\gamma ,0)\pi \gamma L`$. It is this latter limit that leads to the characteristic quadratic dependence on $`L`$: FIG. 1. The medium induced gloun differential multiplicity normalized by the factorization distribution eq.(2) is plotted vs. $`𝐤_{}^2/\mu ^2`$ for opacity $`L/\lambda _g=1,2,3`$. Solid curves show first order (13) and dashed including second order in opacity (19) for a gluon with $`x=0.2`$. (Quark jet of $`E_{jet}=50`$ GeV and $`\mu =0.5`$ GeV). FIG. 2. The contributions to the induced radiation intensity to first and second order in opacity are compared to the medium independent intensity (3). The same conditions as in Fig. 1 are assumed. $$\frac{dI^{(1)}}{dx}\frac{C_R\alpha _s}{4}\frac{1x+\frac{x^2}{2}}{x}\frac{L^2\mu ^2}{\lambda _g}.$$ (24) This formula breaks down at both $`x0`$ and $`x1`$ because $`\|𝐤_{}\|_{\mathrm{max}}`$ cannot be approximated by $`\mathrm{}`$ and because the small $`x`$ approximations used above break down as $`x1`$. The total radiative energy loss is then given by $$\mathrm{\Delta }E^{(1)}=\frac{C_R\alpha _s}{N(E)}\frac{L^2\mu ^2}{\lambda _g}\mathrm{log}\frac{E}{\mu },$$ (25) FIG. 3. The radiated energy loss of a quark jet with energy $`E_{jet}=5,50,500`$ GeV (at SPS, RHIC, LHC) is plotted as a function of the opacity $`L/\lambda _g`$. ($`\lambda _g=1`$ fm, $`\mu =0.5`$ GeV). Solid curves show first order, while dashed curves show results up to second order in opacity. The energy loss (solid triangles) from Ref. (with $`\stackrel{~}{v}=2.5`$) is shown for comparison. with $`N(E)=4`$ if the kinematic bounds are ignored as in Eq. (24). In practice, we emphasize that finite kinematic constraints cause $`N(E)`$ to deviate considerably from the asymptotic value 4. We find that $`N(E)=7.3,\mathrm{\hspace{0.17em}10.1},\mathrm{\hspace{0.17em}24.4}`$ for $`E=500,\mathrm{\hspace{0.17em}50},\mathrm{\hspace{0.17em}5}`$ GeV. Together with the logarithmic dependence of energy, these kinematic effects suppress greatly the energy loss at lower (SPS) energies as seen in Fig. 3.. This is in contrast to the approximately energy independent result in Ref. where the finite kinematic bounds were neglected. Conclusions. We calculated the effect of final state interactions on the induced gluon differential distribution up to second order in opacity for hard jets produced in nuclear reactions. This work generalizes Ref. by taking into account virtual corrections to calculate inclusive rates and provides an complementary analytic approach to clarify nonlinear jet quenching effects predicted in Refs. . The inclusive $`xdN/dxd𝐤_{}^2`$, $`dI/dx`$, and $`\mathrm{\Delta }E`$, was studied as a function of nuclear thickness for jet energies in the SPS, RHIC and LHC range. One of our main results is the demonstration that the second order contribution (19) to the integrated energy loss remains surprisingly small up to realistic nuclear opacities $`L/\lambda _g5`$ ( except at low SPS energies). The rapid convergence of the opacity expansion even for realistic opacities results from the fact that the effective expansion parameter is actually the product of the opacity and the gluon formation probability $`L\mu ^2/xE`$ (see (23)). The leading quadratic dependence of the energy loss on nuclear thickness therefore arises from the simple first order term (8) in this approach. The detailed pattern of angular broadening and the $`x`$ dependence is also dominated by the first order contribution. We note that our first order (effectively power-law) $`\|𝐤_{}\|`$ and $`x`$ distributions, however, differ considerably from the Gaussian form obtained in the eikonal resummation approach of that applies to thick targets. Surprisingly, on the other hand, the magnitude of the integrated $`\mathrm{\Delta }E(L)`$ is rather similar at RHIC energies. At SPS energies kinematic effects suppress greatly the energy loss relative to . Our estimates provide a natural explanation for the absence of jet quenching in $`Pb+Pb`$ at 160 AGeV that has been a puzzle up to now . The short duration of the dense phase further limits the effective opacity at the SPS. A duration of $`L/\lambda _g2`$ for example leads to a total energy loss only $`100`$ MeV, which is much too small to be observable in soft multiple scattering background . At RHIC energies, on the other hand, a significant nonlinear (in $`A`$) pattern of suppression of high $`p_{}`$ hadrons relative to scaled $`pp`$ data should be observable to enable a direct test of non-abelian energy loss mechanisms in dense matter. We note finally that the simplicity of the first and second order results (8,19) will make it possible to improve significantly Monte Carlo event simulations of jet quenching . We thank A. Mueller and U. Wiedemann for discussion on virtual corrections. This work was supported by the DOE Research Grant under Contract No. De-FG-02-93ER-40764, partly by the US-Hungarian Joint Fund No.652 and OTKA No. T029158. \[Note added in proof: U. Wiedemann informed us of an opacity expansion derived independently in .\]
warning/0005/hep-ph0005153.html	ar5iv	text	# Isospin Breaking in the Extraction of Isovector and Isoscalar Spectral Functions from 𝑒⁺⁢𝑒⁻→ℎ⁢𝑎⁢𝑑⁢𝑟⁢𝑜⁢𝑛⁢𝑠 ## Introduction If we define the flavor $`ab`$ ($`a,b=3,8`$) vector current correlators via $$id^4xe^{iqx}0\|T(V_\mu ^a(x)V_\nu ^b(0)^{})\|0(g_{\mu \nu }q^2+q_\mu q_\nu )\mathrm{\Pi }^{ab}(q^2),$$ (1) the corresponding EM current-current correlator is given by $$\mathrm{\Pi }^{EM}=\mathrm{\Pi }^{33}+\frac{2}{\sqrt{3}}\mathrm{\Pi }^{38}+\frac{1}{3}\mathrm{\Pi }^{88},$$ (2) which reduces to $`\mathrm{\Pi }^{33}+\mathrm{\Pi }^{88}/3`$ in the isospin limit. $`\mathrm{\Pi }^{EM}`$ is related to the hadronic electroproduction cross-section by $$\sigma (e^+e^{}hadrons)=\frac{16\pi ^2\alpha ^2}{s}\mathrm{Im}\mathrm{\Pi }^{EM}(s).$$ (3) In the isospin limit, the $`33`$ and $`88`$ contributions associated with $`n\pi `$ states can be classified by $`G`$-parity, making an experimental determination of the separate isovector ($`33`$) and isoscalar ($`88`$) components of the EM spectral function, $`\rho ^{EM}(s)`$, possible. Isospin breaking (IB), however, complicates this separation because of the presence of a non-zero $`38`$ spectral component. One contribution to the $`38`$ component, associated with the process $`e^+e^{}\omega \pi ^+\pi ^{}`$, is conventionally removed by hand in determining the $`33`$ spectral function, but other $`38`$ contributions also exist which are not so subtracted. If we define, in obvious notation, the isovector $`(3)`$ and isoscalar $`(8)`$ components of the EM decay constant of vector meson, $`V`$, via $$F_V^{EM}=F_V^3+\frac{1}{\sqrt{3}}F_V^8$$ (4) then, for example, the intermediate $`\rho `$ contribution to $`\rho ^{EM}(s)`$ contains not only a $`33`$ component, proportional to $`[F_\rho ^3]^2`$, but also, to leading order in IB, a $`38`$ component proportional to $`F_\rho ^3F_\rho ^8`$. The $`\omega `$ contribution to $`\rho ^{EM}(s)`$ similarly contains a $`38`$ component proportional to $`F_\omega ^3F_\omega ^8`$. The presence of these $`38`$ “contaminations” in the conventionally extracted $`33`$ and $`88`$ spectral functions creates potential problems for * CVC tests * the determination of $`m_s`$ using Narison’s $`3388`$ “tau-decay-like” sum rulenarisonms * The inverse weighted chiral $`3388`$ sum rule for the $`6^{th}`$ order ChPT LEC, $`Q_V`$, which governs flavor/isospin breaking in vector current correlatorsgkqvsr . In the latter two cases, the IB effect can be particularly important (1) because the cancellation between the nominal isovector and isoscalar integrals is typically rather strong (to the $`10\%`$ level) and, (2) because the $`38`$ component of the $`\rho `$ and $`\omega `$ EM decay constants induced by $`\rho `$-$`\omega `$ mixing has the opposite sign for the $`\rho `$ and $`\omega `$, the $`38`$ corrections add constructively when one forms the nominal isovector-isoscalar difference. It is easy to see that the relative importance of this effect should be much larger for the isovector “contamination” of the $`\omega `$ decay constant than for the isoscalar “contamination” of the $`\rho `$ decay constant. Indeed, taking $`SU(3)_F`$ for the vector current vacuum-to-vector-meson matrix elements, and assuming ideal mixing of the vector meson sector, $`F_\rho ^{(0)}3F_\omega ^{(0)}`$, where the superscript $`(0)`$ indicates the isospin-unmixed states. With $`ϵ`$ the $`\rho `$-$`\omega `$ mixing angle, one then has $`F_\rho ^{EM}`$ $`=`$ $`F_\rho ^{(0)}ϵF_\omega ^{(0)}F_\rho ^{(0)}\left(1{\displaystyle \frac{ϵ}{3}}\right)`$ $`F_\omega ^{EM}`$ $`=`$ $`F_\omega ^{(0)}+ϵF_\rho ^{(0)}F_\omega ^{(0)}\left(1+3ϵ\right).`$ (5) The fractional contribution of the “wrong-isospin” current to the EM decay constant of the $`\omega `$ should thus be $`9`$ times that of the corresponding IB contribution to $`F_\rho ^{EM}`$, in the limit that the effect is dominated by $`\rho `$-$`\omega `$ mixing. In this paper, we employ a FESR analysis of the IB vector current correlator, $`\mathrm{\Pi }^{38}`$, in order to extract the product, $`F_V^3F_V^8`$, of IB and isospin-conserving (IC) decay constants for the low-lying vector mesons. This allows us to determine the $`38`$ contamination of the resonance contributions to the nominally extracted $`33`$ and $`88`$ spectral functions, and hence to make the corrections required to extract the actual versions thereof from data. ## FESR’s and the Determination of the IB Vector Meson Decay Constants The analyticity structure of $`\mathrm{\Pi }^{38}`$ is such that, for any function $`w(s)`$ analytic in the region of the contour, one has, using Cauchy’s theorem, the FESR relation $$_{s_{th}}^{s_0}𝑑sw(s)\rho ^{38}(s)=\frac{1}{2\pi i}_{\|s\|=s_0}𝑑sw(s)\mathrm{\Pi }^{38}(s),$$ (6) where $`\rho ^{38}(s)=\frac{1}{\pi }\mathrm{Im}\mathrm{\Pi }^{38}(s)`$. Choosing $`s_0`$ large enough that $`\mathrm{\Pi }^{38}`$ can be represented by the OPE on the RHS, but small enough that $`\rho ^{38}`$ on the LHS is still dominated by the known isovector and isoscalar vector mesons, one obtains a relation between the products $`f_V=F_V^3F_V^8`$ which govern the sizes of resonance contributions to $`\rho ^{38}`$ and the parameters ($`\alpha _s`$ and vacuum condensates) entering the OPE. Since the OPE is well known this allows us to constrain the resonance parameters. In order to make the analysis suggested above reliable, we will employ a form of FESR tested in the isovector vector channel and shown to be very accurately satisfied therekmfesr . These “pinch-weighted” FESR’s are those corresponding to weights satisfying $`w(s_0)=0`$. In the isovector vector channel, where they can be tested, they are known to be very accurately satisfied down to rather low scales ($`2\mathrm{GeV}^2`$, even when more general FESR’s, which do not suppress contributions from the region near the timelike real axis, are poorly satisfiedkmfesr ). It is convenient to work with two families of pinched weights, $`w_s(s)=[1s/s_0][1+As/s_0]`$ and $`w_d(s)=[1s/s_0]^2[1+As/s_0]`$, where in each case $`A`$ is a free parameter, used to vary the weight profile. It is worth noting that, using the FESR’s resulting from these two weight families to fit the decay constants of the first three $`\rho `$ resonances, one obtains a determination of the $`\rho (770)`$ decay constant in terms of OPE parameters which is accurate to within experimental errorskmfesr ; kma0 . On the hadronic side of our FESR’s we assume a sum of Breit-Wigner resonance contributions for the $`\rho `$, $`\omega `$, and $`\varphi `$ with PDG values of the masses and widths. In the $`\rho ^{}`$-$`\omega ^{}`$ region, a single effective resonance contribution with the average values of the masses and widths is employed since the individual $`\rho ^{}`$ and $`\omega ^{}`$ masses and widths are very similar, making it impossible for our sum rule to discriminate between them. The necessity of including a possible $`\varphi `$ term has been discussed in Ref. kmold38 , while the importance of including the resonance widths has been stressed in Ref. ijl ). The only remaining unknowns in the spectral ansatz are the products $`f_V=F_V^3F_V^8`$, which determine the heights of the resonance peaks. The $`D=0,2,4,6`$ terms in the OPE representation of $`\mathrm{\Pi }^{38}`$ are given by $`\left[\mathrm{\Pi }_{1\gamma E}^{38}\right]_{D=0}={\displaystyle \frac{\alpha }{16\pi ^3}}{\displaystyle \frac{1}{4\sqrt{3}}}ln(Q^2)`$ $`\left[\mathrm{\Pi }^{38}(Q^2)\right]_{D=2}={\displaystyle \frac{3}{2\pi ^2Q^2}}{\displaystyle \frac{1}{4\sqrt{3}}}\left[(m_d^2m_u^2)(Q^2)\right][1+{\displaystyle \frac{8}{3}}a(Q^2)`$ $`+({\displaystyle \frac{17981}{432}}+{\displaystyle \frac{62}{27}}\zeta (3){\displaystyle \frac{1045}{54}}\zeta (5))a^2(Q^2)]`$ $`\left[\mathrm{\Pi }^{38}(Q^2)\right]_{D=4}={\displaystyle \frac{2\left(m_u\overline{u}um_d\overline{d}d\right)}{4\sqrt{3}Q^4}}\left[1+{\displaystyle \frac{1}{3}}a(Q^2)+{\displaystyle \frac{11}{2}}a^2(Q^2)\right]`$ $`\left[\mathrm{\Pi }^{38}(Q^2)\right]_{D=6}={\displaystyle \frac{112\pi }{81\sqrt{3}Q^6}}\rho _{red}\gamma (\rho \alpha _s\overline{q}q^2),`$ (7) where $`\alpha `$ is the usual EM coupling, $`a(Q^2)=\alpha _s(Q^2)/\pi `$, $`\zeta (n)`$ is the Riemann zeta function, $`\gamma =[<\overline{d}d>/<\overline{u}u>]1`$, and $`\rho _{red}`$ is the ratio of the violation of the vacuum saturation approximation (VSA) estimate for the $`D=6`$ condensate combination in the $`38`$ channel to that in the flavor-diagonal $`33`$ vector channel (see Ref. narisonrho for the latter value). Numerically, the $`D=4`$ term is the largest, owing to the smallness of $`\alpha `$ and $`m_{u,d}`$. It can be rewritten in terms of the IB mass ratio $`r(m_dm_u)/(m_d+m_u)`$ (for which we use the value from Leutwyler’s ChPT analysisleutwylerqm ) and $`<(m_u+m_d)\overline{u}u>f_\pi ^2m_\pi ^2`$ (where we have used the GMOR relation). We employ the 4-loop versions of the running mass and coupling. Lack of knowledge of $`\rho _{red}`$ would normally limit the accuracy of the extraction of the $`f_V`$. Using both the $`w_s`$ and $`w_d`$ families of FESR’s, however, it turns out to be possible to determine $`\rho _{red}`$ self-consistently. In Figure 1 we illustrate the dependence of $`f_\rho `$ on $`\rho _{red}`$ for both the $`w_s`$ and $`w_d`$ weight families. Obviously only one value of $`\rho _{red}`$ provides a consistent determination of $`f_\rho `$. Checking analogous figures for the other $`f_V`$, one finds that the values of $`\rho _{red}`$ for which they become consistent agree with that in Figure 1 to better than $`1\%`$. This clearly demonstrates that this “self-consistency” determination of $`\rho _{red}`$ is physically meaningful. It is interesting to note that the degree of VSA violation is very similar for the flavor-diagonal isovector and IB flavor $`38`$ correlators. Further details of the calculation, including OPE input values, uncertainties, and a description of the contour-improved implementation of the OPE integrals cipt , may be found in Ref. kmcewv3v8 . Having extracted $`f_V`$, and knowing the EM decay constants, $`F_V^{EM}`$, one may then separately determine $`F_V^3`$ and $`F_V^8`$ for $`V=\rho ,\omega ,\varphi `$. To understand the impact of the IB decay constants on the extraction of the $`33`$ and $`88`$ spectral functions, it is convenient to quote the results in terms of the squared ratios, $`r_V`$, shown below. In each case, the numerator in the ratio of decay constants represents the IC decay constant of the meson in question (whose square is relevant to the contribution to either the $`33`$ or $`88`$ spectral function), while the denominator represents the nominal value of the numerator, obtained from EM data neglecting the presence of the “wrong-isospin” IB contributions discussed above. The deviations of the $`r_V`$ from $`1`$ reflect the presence of the IB $`F_\rho ^8`$, $`F_\omega ^3`$ and $`F_\varphi ^3`$ contributions to $`F_{\rho ,\omega ,\varphi }^{EM}`$. The ratios have been defined in such a way that the true resonance contributions to the $`33`$ (or $`88`$) spectral functions are obtained by multiplying the nominal results, obtained in the usual analysis, by $`r_V`$. The results of the analysis outlined above are then $`r_\rho `$ $`=`$ $`\left[{\displaystyle \frac{F_\rho ^3}{F_\rho ^{EM}}}\right]^2=0.982\pm 0.0021`$ $`r_\omega `$ $`=`$ $`\left[{\displaystyle \frac{F_\omega ^8}{\sqrt{3}F_\omega ^{EM}}}\right]^2=1.154\pm 0.017`$ $`r_\varphi `$ $`=`$ $`\left[{\displaystyle \frac{F_\varphi ^8}{\sqrt{3}F_\varphi ^{EM}}}\right]^2=1.009\pm 0.001,`$ (8) where the errors are dominated by the uncertainty in the quark mass ratio, $`r`$, and hence are strongly correlated. We note the following features of the above results: * The effect on the $`\rho `$ contribution to the $`33`$ spectral function is small, and hence has no impact on CVC tests, given the current experimental errors on the EM cross-sections. * As expected from the argument above, the effect is much larger for the $`\omega `$ contribution to the $`88`$ spectral function. Our result corresponds to a $`7.5\%`$ isovector contribution to the physical $`\omega `$ decay constant. * The ratio of $`\rho `$ and $`\omega `$ corrections is $`8.69`$, suggesting dominance by $`\rho `$-$`\omega `$ mixing. The impact of these corrections on the $`33`$-$`88`$ sum rules discussed above is as follows. * Using $`\tau `$ decay data for the isovector part, and the above corrections for the EM isoscalar data, the Narison $`m_s`$ sum rule implies $$m_s(1\mathrm{GeV}^2)=154(\pm 50)\mathrm{MeV},$$ (9) compatible with value $`157\pm 16\pm 15\pm 16`$ obtained in a recent analysis based on flavor breaking in hadronic $`\tau `$ decays jkkm00 , and to be compared to the central value $`197`$ MeV obtained neglecting the presence of IB correctionsnarisonms . * The correction to the inverse chiral $`3388`$ sum rule extraction of the $`6^{th}`$ order ChPT low-energy constant $`Q_V`$ is also large. The FESR solution above implies, via an inverse moment $`38`$ sum rule, the value kmcewv3v8 $$Q_V(m_\rho ^2)=(3.3\pm 0.4)\times 10^5,$$ (10) to be compared to the value from a recent analysis jksd99 of the vector current $`\mathrm{\Pi }^{ud}`$-$`\mathrm{\Pi }^{us}`$ difference, employing $`\tau `$ decay spectral data (which does not need IB corrections) $$Q_V(m_\rho ^2)=(2.8\pm 1.3)\times 10^5.$$ (11) The agreement of the results of Eq. (10) with the independent determination of Eq. (11) provides further strong support for the values of the IB vector meson decay constants obtained in the present analysis.
warning/0005/quant-ph0005039.html	ar5iv	text	# CU-TP972RBRC-89 A New Method to Derive Low-Lying N-dimensional Quantum Wave Functions by Quadratures Along a Single TrajectoryWork supported in part by the U.S. Department of Energy ## 1. Introduction In this paper, we begin with the generalization to arbitrary dimensions of a new method in which the low-lying quantum wave function of a particle in a potential $`+V`$ can be obtained by quadratures along the trajectory of a corresponding classical mechanical problem with $`V`$ as its potential. Let $`𝐪=(q_1,q_2,\mathrm{},q_N)`$ describes the $`N`$-dimensional coordinates of the particle, and $`\mathrm{\Phi }(𝐪)`$ is the ground-state (or low-lying) quantum wave function that satisfies $`H\mathrm{\Phi }(𝐪)=E\mathrm{\Phi }(𝐪),`$ (1.1) where $`H={\displaystyle \frac{1}{2}}^2+V(𝐪)`$ (1.2) is the Hamiltonian of a non-relativistic particle of unit mass, and $`^2={\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{^2}{q_i^2}}.`$ (1.3) We assume $`V(𝐪)`$ to be bounded from below; therefore we can choose $`V(𝐪)0.`$ (1.4) Its minimum $`V(𝐪)=0`$ may occur at more than one point, each with a non-vanishing curvature. As in , we introduce a scale factor $`g^2`$ by writing $`V(𝐪)=g^2v(𝐪),`$ (1.5) and consider the case of large $`g`$. We express $`\mathrm{\Phi }(𝐪)=e^{g𝐒(𝐪)}`$ (1.6) in terms of a formal expansion: $`g𝐒(𝐪)=g𝐒_0(𝐪)+𝐒_1(𝐪)+g^1𝐒_2(𝐪)+g^2𝐒_3(𝐪)+\mathrm{}`$ (1.7) with a similar expansion for $`E`$, $`E=gE_0+E_1+g^1E_2+g^2E_3+\mathrm{}.`$ (1.8) Substituting (1.6) - (1.8) into the Schroedinger equation (1.1) and equating the coefficients of $`g^n`$ on both sides, we derive $`(𝐒_0)^2`$ $`=`$ $`2v,`$ (1.9) $`𝐒_0𝐒_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}^2𝐒_0E_0,`$ (1.10) $`𝐒_0𝐒_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}[^2𝐒_1(𝐒_1)^2]E_1,`$ (1.11) $`𝐒_0𝐒_3`$ $`=`$ $`{\displaystyle \frac{1}{2}}[^2𝐒_22(𝐒_1)(𝐒_2)]E_2,`$ (1.12) etc., where $``$ is the $`N`$-dimensional gradient vector whose components are $`_i=/q_i`$ and $`i=1,2,\mathrm{},N.`$ (1.13) The underlying reason for considering the expressions (1.7) and (1.8) is that for large $`g`$ the ground-state quantum wave function has its amplitude centered around the minimum of $`V(𝐪)`$; near the minimum, $`V(𝐪)`$ may be approximated by a quadratic function of $`𝐪`$, with the curvature of the corresponding potential surface proportional to $`g^2`$. Thus, for large $`g`$, the ground-state energy is $`O(g)`$ and so is the exponent of the ground-state wave function, as suggested by the ground-state energy and the corresponding wave function of a harmonic oscillator. The parameter $`g^1`$ serves as a measure of anharmonicity. It is useful to write the first line of (1.9) as $`{\displaystyle \frac{1}{2}}(𝐒_0)^2v(𝐪)=e,`$ (1.14) where $`e=0+.`$ (1.15) Equation (1.14) is equivalent to the corresponding Hamilton-Jacobi equation in classical mechanics, with $`V(𝐪)`$ as the potential and $`e`$ the total energy. To solve $`(\text{1.14})`$, we pick a minimum of $`V(𝐪)`$, say $`𝐪=\mathrm{𝟎}`$ with $`V(\mathrm{𝟎})=g^2v(\mathrm{𝟎})=0.`$ (1.16) Consider a trajectory $`𝐪(t)`$ which begins at $`𝐪=\mathrm{𝟎}`$ when $`t=0`$ and ends at $`𝐪_T`$ when $`t=T`$; i.e., $`𝐪(0)=\mathrm{𝟎}\mathrm{𝖺𝗇𝖽}𝐪(T)=𝐪_𝐓.`$ (1.17) Hamilton’s action-integral is given by the minimum of the path integral $`𝐒_0(𝐪_𝐓,e)={\displaystyle \underset{0}{\overset{T}{}}}[{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}\underset{i}{\overset{}{q}}(t)^2(v(𝐪(t)))]𝑑t+Te,`$ (1.18) where $`T=({\displaystyle \frac{𝐒_0}{e}})_{𝐪_𝐓}`$ (1.19) is the total time of the trajectory from the initial point 0 to the final point $`𝐪_𝐓`$. Since the potential in the classical problem $`=v(𝐪)0`$, the integrand in (1.18) is positive everywhere. Hence, given any final point $`𝐪_𝐓`$ and any total energy $`e>0`$, the minimum of (1.18) exists, so is therefore the classical trajectory. (For clarity, assume the minimum to be unique, a restriction that will be removed later.) As $`e0+`$, the initial velocity $`\stackrel{}{𝐪}`$ is proportional to $`e^{1/2}`$ which also approaches zero; the total time $`T`$ of the classical trajectory is $`e^{1/2}\mathrm{}`$, yielding $`\underset{e=0+}{lim}Te=0.`$ (1.20) The same limit $`e=0+`$ of (1.18) gives $`𝐒_0`$ of (1.14) - (1.15). To go back to the first line of (1.9), we designate the end-point $`𝐪_𝐓`$ of the classical trajectory simply as $`𝐪`$, the point of interest for the quantum solution $`\mathrm{\Phi }(𝐪)`$. Setting $`𝐪_𝐓=𝐪,`$ (1.21) the solution $`𝐒_0(𝐪)`$ of the first line in (1.9) can now be written as $`𝐒_0(𝐪)={\displaystyle \underset{0}{\overset{T}{}}}[{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}\underset{i}{\overset{2}{\stackrel{}{q}}}(v(𝐪))]𝑑t,`$ (1.22) where the integral is along the trajectory $`q_i(t)`$ that satisfies the classical equations of motion $`\underset{i}{\overset{}{q}}=v/q_i,`$ (1.23) the energy conservation $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}\underset{i}{\overset{}{q}}(t)^2v(𝐪(t))=0+`$ (1.24) and the initial and final configurations (1.17) and (1.21). The solution $`𝐒_0(𝐪)`$ thus obtained is positive and continuous everywhere; it is an increasing function along the direction of the classical trajectory. In the case that the classical potential $`v(𝐪)`$ has other absolute maxima besides $`0`$, say $`v(𝐚_\mathrm{𝟏})=v(𝐚_\mathrm{𝟐})=\mathrm{}=v(𝐚_𝐧)=0,`$ (1.25) $`𝐒_0(𝐪)`$ is not analytic at these isolated points $`𝐪=𝐚_\mathrm{𝟏},𝐚_\mathrm{𝟐},\mathrm{},𝐚_𝐧.`$ (1.26) For a total energy $`e>0`$, along the classical trajectory the velocity must continue in its direction when the trajectory passes through any of the other maxima of $`v(𝐪)`$, say $`𝐪=𝐚_𝐧`$. Thus, as $`e0+`$, $`𝐒_0`$ becomes zero at $`𝐪=𝐚_𝐧`$, but its sign along the trajectory-direction remains the same before or after $`𝐚_𝐧`$; this forces $`𝐒_0`$ to develop a kink at $`𝐪=𝐚_𝐧`$. At the initial point $`𝐪=\mathrm{𝟎}`$, while $`v(\mathrm{𝟎})=0`$ and therefore $`𝐒_0=0`$ as $`e=0+`$, $`𝐒_0(𝐪)`$ is analytic, since different trajectories emanating from $`𝐪=\mathrm{𝟎}`$ have to go along different directions. At infinity, it is easy to see that $`𝐒_0(𝐪)=\mathrm{}`$, and therefore $`\mathrm{\Phi }_0e^{g𝐒_0}`$ is zero. To solve the second equation in (1.9), we require $`𝐒_1(𝐪)`$ to be also analytic at $`𝐪=\mathrm{𝟎}`$; this yields $`E_0={\displaystyle \frac{1}{2}}(^2𝐒_0)_{\mathrm{at}𝐪=\mathrm{𝟎}}.`$ (1.27) It is convenient to consider the surface $`𝐒_0(𝐪)=\mathrm{𝖼𝗈𝗇𝗌𝗍𝖺𝗇𝗍};`$ (1.28) its normal is along the corresponding classical trajectory passing through $`𝐪`$. Characterize each classical trajectory by the $`𝐒_0`$-value along the trajectory and a set of $`N1`$ angular variables $`\alpha =(\alpha _1(𝐪),\alpha _2(𝐪),\mathrm{},\alpha _{N1}(𝐪)),`$ (1.29) so that each $`\alpha `$ determines one classical trajectory with $`\alpha _j𝐒_0=0,`$ (1.30) where $`j=1,2,\mathrm{},N1.`$ (1.31) (As an example, we note that as $`𝐪\mathrm{𝟎}`$, $`v(𝐪)\frac{1}{2}\underset{i}{}\omega _i^2q_i^2`$ and therefore $`𝐒_0\frac{1}{2}\underset{i}{}\omega _iq_i^2`$. Consider the ellipsoidal surface $`𝐒_0=`$ constant $`\frac{1}{2}l^2`$. For $`l`$ sufficiently small, each classical trajectory is normal to this ellipsoidal surface. A convenient choice of $`\alpha `$ could be simply any $`N1`$ orthogonal parametric coordinates on the surface.) Each $`\alpha `$ designates one classical trajectory, and vice versa. Every $`(𝐒_0,\alpha )`$ is mapped into a unique set $`(q_1,q_2,\mathrm{},q_N)`$ with $`𝐒_00`$ by construction. Depending on the problem, the converse may or may not be true; i.e. $`(q_1,q_2,\mathrm{},q_N)(𝐒_0,\alpha )`$ could be either one to one, or one to many. In the latter case, we regard the points in the $`𝐪`$-space as specified by the coordinates $`(𝐒_0,\alpha )`$, rather than the original coordinates $`(q_1,q_2,\mathrm{},q_N)`$. Write $`𝐒_1(𝐪)=𝐒_1(𝐒_0,\alpha ),`$ (1.32) the second line of (1.9) becomes $`(𝐒_0)^2({\displaystyle \frac{𝐒_1}{𝐒_0}})_\alpha ={\displaystyle \frac{1}{2}}^2𝐒_0E_0,`$ (1.33) and therefore $`𝐒_1(𝐪)=𝐒_1(𝐒_0,\alpha )={\displaystyle \underset{0}{\overset{𝐒_0}{}}}{\displaystyle \frac{d𝐒_0}{(𝐒_0)^2}}[{\displaystyle \frac{1}{2}}^2𝐒_0E_0],`$ (1.34) where the integration is taken along the classical trajectory of constant $`\alpha `$. Likewise, the third, fourth and other lines of (1.9) lead to $`E_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}[^2𝐒_1(𝐒_1)^2]_{\mathrm{at}𝐪=\mathrm{𝟎}},`$ (1.35) $`𝐒_2(𝐪)=𝐒_2(𝐒_0,\alpha )`$ $`=`$ $`{\displaystyle \underset{0}{\overset{𝐒_0}{}}}{\displaystyle \frac{d𝐒_0}{(𝐒_0)^2}}\{{\displaystyle \frac{1}{2}}[^2𝐒_1(𝐒_1)^2]E_1\},`$ (1.36) $`E_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}[^2𝐒_22(𝐒_1)(𝐒_2)]_{\mathrm{at}𝐪=\mathrm{𝟎}},`$ (1.37) $`𝐒_3(𝐪)=𝐒_3(𝐒_0,\alpha )`$ $`=`$ $`{\displaystyle \underset{0}{\overset{𝐒_0}{}}}{\displaystyle \frac{d𝐒_0}{(𝐒_0)^2}}\{{\displaystyle \frac{1}{2}}[^2𝐒_22(𝐒_1)(𝐒_2)]E_2\},`$ (1.38) etc. These solutions give the convenient normalization convention at $`𝐪=\mathrm{𝟎}`$, $`𝐒(\mathrm{𝟎})=0`$ (1.39) $`\mathrm{and}`$ (1.40) $`\mathrm{\Phi }(\mathrm{𝟎})=1.`$ (1.41) The above expressions for $`E_0,E_1,E_2,\mathrm{}`$ can be derived more directly by introducing $`\sigma (𝐪)`$ $``$ $`g𝐒(𝐪)g𝐒_0(𝐪)`$ (1.42) $`=`$ $`𝐒_1(𝐪)+g^1𝐒_2(𝐪)+g^2𝐒_3(𝐪)+\mathrm{}.`$ (1.43) Thus, the ground-state wave function $`\mathrm{\Phi }(𝐪)`$ is $`\mathrm{\Phi }=e^{g𝐒_0\sigma },`$ (1.44) and the corresponding Schroedinger equation (1.1) gives $`g𝐒_0\sigma ={\displaystyle \frac{1}{2}}[g^2𝐒_0+^2\sigma (\sigma )^2]E.`$ (1.45) At $`𝐪=\mathrm{𝟎}`$, $`v(𝐪)=0`$ and $`𝐒_0=0`$, on account of $`(𝐒_0)^2=2v`$. Assuming $`\sigma `$ to be regular, we have $`E={\displaystyle \frac{1}{2}}[g^2𝐒_0+^2\sigma (\sigma )^2]_{\mathrm{𝖺𝗍}𝐪=\mathrm{𝟎}}`$ (1.46) which leads to (1.27), (1.35), (1.34), $`\mathrm{}`$ for $`E_0`$, $`E_1`$, $`E_2`$, $`\mathrm{}`$. Example. As an example, consider the special case $`V(𝐪)={\displaystyle \frac{g^2}{2}}(q_1^2+q_2^2+\mathrm{}+q_N^2).`$ (1.47) From (1.9), one can readily show that the first line gives $`𝐒_0(𝐪)={\displaystyle \frac{g}{2}}(q_1^2+q_2^2+\mathrm{}+q_N^2)`$ (1.48) and (1.27) leads to $`E_0=N/2;`$ (1.49) the other equations (1.34) - (1.35), etc., yield $`E_1=E_2=\mathrm{}=0`$ and $`S_1=S_2=\mathrm{}=0`$. The result is the well known exact answer $`\mathrm{\Phi }(𝐪)=exp[{\displaystyle \frac{g}{2}}(q_1^2+q_2^2+\mathrm{}+q_N^2)]`$ (1.50) and $`E={\displaystyle \frac{g}{2}}N.`$ (1.51) The organization of the paper is as follows: In Section 2, we discuss the case of a separable potential, and show that the solution of the wave function is a product of one-dimensional functions, as expected. The particular one-dimensional example of $`V(x)={\displaystyle \frac{g^2}{2}}(x^2a^2)^2`$ (1.52) has been analyzed in ref. using our method. In this problem, besides the anharmonicity parameter $`g^1`$, there is also a much smaller barrier penetration parameter $`e^{\frac{4}{3}ga^3}.`$ (1.53) (The same problem has also been extensively studied in the literature\[2-9\], using mostly Feynman’s path integration method. Quite often, the two parameters $`g`$ and $`a`$ in (1.52) are related by $`2ag=1`$, which makes $`V=\frac{1}{2}y^2(1gy)^2`$ with $`y=x+a`$. The dimensionless anharmonicity parameter $`1/ga^3`$ then becomes $`8g^2`$ and the barrier penetration parameter (1.53) becomes $`e^{1/6g^2}`$. In our $`g^1`$-expansion, we keep the parameter $`a`$ fixed.) The ground-state wave function $`\psi _{even}(x)`$ is even in $`x`$, and the first excited state is an odd function $`\psi _{odd}(x)`$. Let $`E_{even}`$ and $`E_{odd}`$ be their corresponding eigenvalues. We can expand both their sum and difference as the following double series: $`{\displaystyle \frac{1}{2}}(E_{odd}+E_{even})=ga{\displaystyle \underset{m,n}{}}C_{mn}(ga^3)^me^{\frac{4}{3}nga^3}`$ (1.54) and $`{\displaystyle \frac{1}{2}}(E_{odd}E_{even})`$ $`=`$ $`4({\displaystyle \frac{2}{\pi }}g^3a^5)^{\frac{1}{2}}e^{\frac{4}{3}ga^3}`$ (1.56) $`{\displaystyle \underset{m,n}{}}c_{mn}(ga^3)^me^{\frac{4}{3}(n1)ga^3}`$ where $`C_{mn}`$ and $`c_{mn}`$ are numerical coefficients, which can be explicitly expressed as definite integrals, with the results $`c_{m0}`$ $`=`$ $`0\mathrm{𝖿𝗈𝗋}\mathrm{𝖺𝗅𝗅}m,`$ (1.57) $`C_{00}`$ $`=`$ $`c_{01}=1,C_{10}={\displaystyle \frac{1}{4}},`$ (1.58) $`C_{20}`$ $`=`$ $`{\displaystyle \frac{9}{2^6}},C_{30}={\displaystyle \frac{89}{2^9}},\mathrm{𝖾𝗍𝖼}.`$ (1.59) The terms with $`n=1`$ represent the one-instanton contributions\[2-9\], $`n=2`$ the two-instanton contributions, etc. As one can see from the derivation, an important tool, in addition to the $`g^1`$ expansion (1.7) - (1.9) discussed above, is the use of the Sturm-Liouville type of Green’s functions. In this paper, we will show how to generalize such Green’s function technique to an N-dimensional problem. In Section 3, we introduce a compact notation which will pave the way towards the generalization. Section 4 contains the main body of this paper. We depart from the initial approach based on the Hamilton-Jacobi equation; instead, we consider two N-dimensional Hamiltonians: $`H={\displaystyle \frac{1}{2}}^2+V(𝐪)`$ (1.60) with, as before, $`V(𝐪)=g^2v(𝐪)0,`$ (1.61) and $`=H+ϵU(𝐪).`$ (1.62) Assuming that the ground-state wave function $`\mathrm{\Phi }(𝐪)=e^{g𝐒(𝐪)}`$ of the unperturbed Hamiltonian $`H`$ is known, we set out to derive the perturbation series expansion for the corresponding ground-state wave function $`\mathrm{\Psi }(𝐪)`$ of $``$ with $`ϵU`$ as the perturbation. Again, as we shall see, to each order in $`ϵ`$, the perturbed $`\mathrm{\Psi }(𝐪)`$ can be obtained by quadrature along a single trajectory connecting $`\mathrm{𝟎}`$, the point that the unperturbed $`\mathrm{\Phi }(𝐪)`$ is maximum (i.e. $`𝐒`$ is minimum), to the point $`𝐪`$ of interest; the trajectory is determined by being the normal to the $`𝐒`$ = constant surfaces along the entire course of the trajectory. (It may be called the quantum trajectory, in contrast to the classical trajectory determined by the Hamilton-Jacobi equation.) The form of the new perturbation series differs from the usual expression; the result is of course the same, as will be illustrated by a few examples. The new perturbation series formula will also be shown to cover the previous analysis in Sections 1-3 based on the classical trajectory derived from the Hamilton-Jacobi equation as special cases. In the same Section 4, the Green’s function $`G`$ of the N-dimensional Hamiltonian $`H`$ will be derived, with $`G`$ satisfying $`(HE)G=𝖨`$ (1.63) and $`E`$ being the ground-state energy of $`H`$. Usually, the Green’s function and the perturbation series expansion to each order $`ϵ^n`$, require either an infinite sum over all excited levels of $`H`$, or equivalently, a sum over all possible paths through Feynman’s path integration method. The Green’s function and the perturbation series formula derived here are different; they depend only on quadratures along a single trajectory. In Section 5, we give examples to illustrate the applications of the new perturbation series, as well as to clarify the properties of the Green’s function. In Section 6, we show how the trajectory-quadrature formulation can be extended to derive excited states, and in Section 7 we discuss how generalization can be made to allow for singular potentials, such as an attractive Coulomb potential with perturbations. As we shall see, with our new approach we can evaluate , e.g., the Stark effect to any finite order of the perturbation by quadratures along the radial trajectory. Although our methods are elementary, they become surprisingly powerful for the following reason. For example, in (1.38) by writing the wave function in the form of $`e^{g𝐒_0\sigma }`$, where $`g`$ is large, we cause the cross-term $`g𝐒_0\sigma `$ to be much larger than $`\frac{1}{2}[^2\sigma (\sigma )^2]`$. Therefore the differential equation for $`\sigma `$ can be written order by order in such a way that $`\frac{1}{2}[^2\sigma (\sigma )^2]`$ is known from previous orders, and $`𝐒_0\sigma =(𝐒_0)^2\frac{\sigma }{𝐒_0}`$ is inferred from the equation. In effect, to each order, $`\sigma `$ now satisfies a first order ordinary differential equation (along the trajectory of constant $`\alpha `$), not a second order partial differential equation. This makes possible a direct solution through quadratures. With variations, this theme runs through all the sections to be presented in this paper (sometimes $`e^\sigma `$ is written $`e^\tau `$, or $`\chi `$, and sometimes the known exponent is called $`g𝐒`$, rather than $`g𝐒_0`$). ## 2. Separable Potential For clarity, we consider the special case of dimension $`N=2`$. The discussions below for a separable potential can be trivially extended to any $`N>2`$. Let the two dimensional coordinate $`𝐪`$ in the previous section be designated $`𝐪=(x,y)`$ (2.1) and the quantum mechanical potential $`V(𝐪)`$ and $`^2`$ in (1.2) be $`V(𝐪)=g^2v(𝐪)=g^2[v_x(x)+v_y(y)]`$ (2.2) and $`^2=(^2/x^2)+(^2/y^2).`$ (2.3) As in (1.4) and (1.16), we assume $`v_x(x)0,v_y(y)0`$ (2.4) and at the point $`x=y=0`$, $`v_x(0)=v_y(0)=0.`$ (2.5) The classical Lagrangian density for the corresponding potential $`v(𝐪)`$ is $`{\displaystyle \frac{1}{2}}[\stackrel{}{x}(t)^2+\stackrel{}{y}(t)^2]+v_x(x)+v_y(y),`$ (2.6) where the dots denote time-derivatives, as before. The equations of motion are $`\stackrel{}{x}(t)={\displaystyle \frac{d}{dx}}v_x(x)`$ (2.7) $`\mathrm{and}`$ (2.8) $`\stackrel{}{y}(t)={\displaystyle \frac{d}{dy}}v_y(y).`$ (2.9) As in (1.17), we consider a trajectory that begins at $`x(0)=y(0)=0`$ (the maximum, or one of the maxima, of $`v`$) at $`t=0`$ and ends at $`x(T)=x_T\mathrm{and}y(T)=y_T`$ (2.10) when $`t=T>0`$, in which, $`x_T`$ and $`y_T`$ are variables independent of T. The action integral along the trajectory is given by the sum of $`A_x(x_T,T)={\displaystyle \underset{0}{\overset{T}{}}}[{\displaystyle \frac{1}{2}}\stackrel{2}{\stackrel{}{x}}(v_x(x))]𝑑t`$ (2.11) $`\mathrm{and}`$ (2.12) $`A_y(y_T,T)={\displaystyle \underset{0}{\overset{T}{}}}[{\displaystyle \frac{1}{2}}\stackrel{2}{\stackrel{}{y}}(v_y(y))]𝑑t`$ (2.13) The minimum of $`A_x(x_T,T)`$ and $`A_y(y_T,T)`$, keeping $`x_T`$, $`y_T`$ and $`T`$ fixed, gives the Lagrangian equations of motion (2.7), and the derivatives of $`A_x`$ and $`A_y`$ with respect to the final variables $`x_T`$, $`y_T`$ and $`T`$ are $`{\displaystyle \frac{}{x_T}}A_x(x_T,T)=\stackrel{}{x}(T)`$ (2.14) $`{\displaystyle \frac{}{y_T}}A_y(y_T,T)=\stackrel{}{y}(T)`$ (2.15) $`{\displaystyle \frac{}{T}}A_x(x_T,T)=e_x`$ (2.16) $`\mathrm{and}`$ (2.17) $`{\displaystyle \frac{}{T}}A_y(y_T,T)=e_y,`$ (2.18) where $`e_x`$ and $`e_y`$ are the energies in the $`x`$ and $`y`$ directions, given by $`e_x={\displaystyle \frac{1}{2}}\stackrel{2}{\stackrel{}{x}}v_x(x)`$ (2.19) $`\mathrm{and}`$ (2.20) $`e_y={\displaystyle \frac{1}{2}}\stackrel{2}{\stackrel{}{y}}v_y(y).`$ (2.21) For a given final configuration characterized by $`x_T`$, $`y_T`$ and $`T`$, both $`e_x`$, $`e_y`$, and therefore also their sum $`e=e_x+e_y`$ (2.22) are determined. Following (1.18), we transform the independent variables from $`x_T`$, $`y_T`$ and $`T`$ to $`x_T`$, $`y_T`$ and $`e`$. We write $`S_0(x_T,y_T,e)=A_x(x_T,T)+A_y(y_T,T)+Te,`$ (2.23) where $`T`$ is a dependent variable, determined by $`T=T(x_T,y_T,e)=({\displaystyle \frac{S_0}{e}})_{x_T,y_T}.`$ (2.24) Given $`x_T`$, $`y_T`$ and $`e`$, the classical trajectory is given by the minimum of (2.23), from which $`T`$, $`e_x`$, $`e_y`$ are also determined. Thus, we can regard $`e_x=e_x(x_T,y_T,e),`$ (2.25) $`\mathrm{and}`$ (2.26) $`e_y=e_y(x_T,y_T,e).`$ (2.27) Hamilton’s action integral (1.18) becomes $`S_0(x_T,y_T,e)=X_0(x_T,e_x)+Y_0(y_T,e_y),`$ (2.28) where $`X_0(x_T,e_x)=A_x(x_T,T)+Te_x`$ (2.29) $`\mathrm{and}`$ (2.30) $`Y_0(y_T,e_y)=A_y(y_T,T)+Te_y,`$ (2.31) in which $`T`$ is given by (2.24). Combining (2.14) and (2.29), we can readily show that $`({\displaystyle \frac{S_0}{x_T}})_{y_T,e}=({\displaystyle \frac{X_0}{x_T}})_{e_x}=({\displaystyle \frac{A_x}{x_T}})_T=\stackrel{}{x}(T)`$ (2.32) $`({\displaystyle \frac{S_0}{y_T}})_{x_T,e}=({\displaystyle \frac{Y_0}{y_T}})_{e_y}=({\displaystyle \frac{A_y}{y_T}})_T=\stackrel{}{y}(T)`$ (2.33) $`\mathrm{and}`$ (2.34) $`({\displaystyle \frac{X_0}{e_x}})_{x_T}=({\displaystyle \frac{Y_0}{e_y}})_{y_T}=({\displaystyle \frac{S_0}{e}})_{x_T,y_T}=T.`$ (2.35) Along each classical trajectory, when the final time variable changes from $`T`$ to $`T+dT`$, the end point moves from $`(x_T,y_T)`$ to ($`x_T+\stackrel{}{x}(T)dT`$, $`y_T+\stackrel{}{y}(T)dT`$); correspondingly, $`dS_0=({\displaystyle \frac{X_0}{x_T}})_{e_x}\stackrel{}{x}(T)dT+({\displaystyle \frac{Y_0}{y_T}})_{e_y}\stackrel{}{y}(T)dT+edT`$ (2.36) Combining the first two lines of (2.32) with (2.36) and taking the limit $`e=0+`$, we find $`dS_0`$ $`=`$ $`[({\displaystyle \frac{X_0}{x_T}})^2+({\displaystyle \frac{Y_0}{y_T}})^2]dT`$ (2.37) $`=`$ $`(S_0)^2dT.`$ (2.38) Along the trajectory we also have $`dT=dx_T/\stackrel{}{x}(T)=dy_T/\stackrel{}{y}(T),`$ (2.39) which, together with (2.32), leads to $`dT=dx_T/({\displaystyle \frac{X_0}{x_T}})_{e_x}=dy_T/({\displaystyle \frac{Y_0}{y_T}})_{e_y}.`$ (2.40) Substituting (2.37) and (2.40) into (1.34) and labeling $`x_T`$, $`y_T`$ simply as $`x`$, $`y`$, we obtain $`S_1(x,y)=X_1(x)+Y_1(y),`$ (2.41) where $`X_1(x)={\displaystyle \underset{0}{\overset{x}{}}}𝑑x({\displaystyle \frac{X_0}{x}})^1[{\displaystyle \frac{1}{2}}{\displaystyle \frac{^2X_0}{x^2}}E_{0,x}],`$ (2.42) $`Y_1(x)={\displaystyle \underset{0}{\overset{y}{}}}𝑑y({\displaystyle \frac{Y_0}{y}})^1[{\displaystyle \frac{1}{2}}{\displaystyle \frac{^2Y_0}{y^2}}E_{0,y}]`$ (2.43) with $`E_{0,x}={\displaystyle \frac{1}{2}}({\displaystyle \frac{^2X_0}{x^2}})_{\mathrm{at}x=0}`$ (2.44) $`\mathrm{and}`$ (2.45) $`E_{0,y}={\displaystyle \frac{1}{2}}({\displaystyle \frac{^2Y_0}{y^2}})_{\mathrm{at}y=0}`$ (2.46) and $`E_0=E_{0,x}+E_{0,y},`$ (2.47) consistent with (1.27). Likewise, $`𝐒_2`$, $`𝐒_3`$, $`\mathrm{}`$ can be written as sums of the form $`X_2(x)+Y_2(y)`$, $`X_3(x)+Y_3(y)`$, $`\mathrm{}`$, with $`X_2(x)`$ $`=`$ $`{\displaystyle \underset{0}{\overset{x}{}}}𝑑x({\displaystyle \frac{X_0}{x}})^1\{{\displaystyle \frac{1}{2}}[{\displaystyle \frac{^2X_1}{x^2}}({\displaystyle \frac{X_1}{x}})^2]E_{1,x}\},`$ (2.48) $`Y_2(y)`$ $`=`$ $`{\displaystyle \underset{0}{\overset{y}{}}}𝑑y({\displaystyle \frac{Y_0}{y}})^1\{{\displaystyle \frac{1}{2}}[{\displaystyle \frac{^2Y_1}{y^2}}({\displaystyle \frac{Y_1}{y}})^2]E_{1,y}\},`$ (2.49) $`X_3(x)`$ $`=`$ $`{\displaystyle \underset{0}{\overset{x}{}}}𝑑x({\displaystyle \frac{X_0}{x}})^1\{{\displaystyle \frac{1}{2}}[{\displaystyle \frac{^2X_2}{x^2}}2{\displaystyle \frac{X_1}{x}}{\displaystyle \frac{X_2}{x}}]E_{2,x}\},`$ (2.50) $`Y_3(y)`$ $`=`$ $`{\displaystyle \underset{0}{\overset{y}{}}}𝑑y({\displaystyle \frac{Y_0}{y}})^1\{{\displaystyle \frac{1}{2}}[{\displaystyle \frac{^2Y_2}{y^2}}2{\displaystyle \frac{Y_1}{y}}{\displaystyle \frac{Y_2}{y}}]E_{2,y}\},`$ (2.51) where $`E_{1,x}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[{\displaystyle \frac{^2X_1}{x^2}}({\displaystyle \frac{X_1}{x}})^2]_{\mathrm{𝖺𝗍}x=0},`$ (2.52) $`E_{1,y}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[{\displaystyle \frac{^2Y_1}{y^2}}({\displaystyle \frac{Y_1}{y}})^2]_{\mathrm{𝖺𝗍}y=0},`$ (2.53) $`E_{2,x}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[{\displaystyle \frac{^2X_2}{x^2}}2{\displaystyle \frac{X_1}{x}}{\displaystyle \frac{X_2}{x}}]_{\mathrm{𝖺𝗍}x=0},`$ (2.54) $`E_{2,y}`$ $`=`$ $`{\displaystyle \frac{1}{2}}[{\displaystyle \frac{^2Y_2}{y^2}}2{\displaystyle \frac{Y_1}{y}}{\displaystyle \frac{Y_2}{y}}]_{\mathrm{𝖺𝗍}y=0},`$ (2.55) etc.. The corresponding wave function factors $`e^{g𝐒_0}`$, $`e^{𝐒_1}`$, $`e^{g^1𝐒_2}`$, $`\mathrm{}`$ are all products $`e^{g(X_0+Y_0)}`$, $`e^{(X_1+Y_1)}`$, $`e^{g^1(X_2+Y_2)}`$, etc.. ## 3. A Compact Notation In this section we shall go over completely to the variables $`𝐒_0`$, $`\alpha `$ and express the $`g^1`$ expansion ((1.7) and (1.34)-(1.35)) described in Section 1 as iteration of an integral operator along a classical trajectory, that is for variable $`𝐒_0`$ and fixed $`\alpha `$. As we shall see, our results can be written in a compact form by introducing an implicit matrix notation in which the indices are two values of $`𝐒`$. In terms of the function $`\sigma (𝐪)`$, introduced by (1.42), the wave function $`\mathrm{\Phi }(𝐪)`$ is $`\mathrm{\Phi }(𝐪)=e^{g𝐒_0(𝐪)\sigma (𝐪)}.`$ (3.1) Because $`(𝐒_0)^2`$ $`=`$ $`2v,`$ (3.2) the Schroedinger equation (1.1) is equivalent to the following equation for $`e^\sigma `$: $`[g(𝐒_0)+T]e^{\sigma (𝐪)}=[{\displaystyle \frac{g}{2}}^2𝐒_0+E]e^{\sigma (𝐪)},`$ (3.3) where $`T={\displaystyle \frac{1}{2}}^2.`$ (3.4) Its solution is given by (1.42) and (1.46), which can be written in a more compact form, as we shall see. Consider the coordinate transformation $`q_1,q_2,q_3,\mathrm{},q_N𝐒_0,\alpha _1,\alpha _2,\mathrm{},\alpha _{N1}`$ (3.5) with $`\alpha =(\alpha _1,\alpha _2,\mathrm{},\alpha _{N1})`$ denoting the set of $`N1`$ orthogonal angular coordinates introduced in (1.29) - (1.31). Define the $`\theta _0`$-function: $`\theta _0(𝐒_0\overline{𝐒}_0)=\{\begin{array}{cc}1& \mathrm{𝗂𝖿}0\overline{𝐒}_0<𝐒_0\\ 0& \mathrm{𝗂𝖿}0𝐒_0<\overline{𝐒}_0\end{array}`$ (3.8) where $`𝐒_0`$ and $`\overline{𝐒}_0`$ vary from $`0`$ to $`\mathrm{}`$. For $`𝐒_0>0`$, its derivative gives Dirac’s $`\delta `$-function: $`{\displaystyle \frac{}{𝐒_0}}\theta _0(𝐒_0\overline{𝐒}_0)=\delta (𝐒_0\overline{𝐒}_0).`$ (3.9) Let us define $`C_0(\alpha )`$ as a square matrix in the $`𝐒_0`$-space, with its matrix element $`(𝐒_0\|C_0(\alpha )\|\overline{𝐒}_0)=g^1\theta _0(𝐒_0\overline{𝐒}_0)/(\overline{𝐒}_0)^2.`$ (3.10) An important feature is that each matrix element of $`C_0(\alpha )`$ connects two points $`𝐪(𝐒_0,\alpha )`$ and $`\overline{𝐪}(\overline{𝐒}_0,\alpha )`$ along the same classical trajectory that satisfies (1.23) - (1.24), and therefore having the same angular variables $`\alpha _1,\alpha _2,\mathrm{},\alpha _{N1}`$. Throughout the paper, the variables $`𝐒_0`$ and $`\overline{𝐒}_0`$ vary between $`0`$ and $`\mathrm{}`$, with the end points $`0`$ and $`\mathrm{}`$ treated as boundaries. All gradients refer to differentiations with respect to the Cartesian coordinates; i.e., $`_j={\displaystyle \frac{}{q_j}}\mathrm{𝖺𝗇𝖽}^2={\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{^2}{q_j^2}},`$ (3.11) as before. On account of (1.30) and (3.8), for $`𝐒_0>0`$, $`𝐒_0\theta _0(𝐒_0\overline{𝐒}_0)=(𝐒_0)^2\delta (𝐒_0\overline{𝐒}_0),`$ (3.12) where the $`j^{\mathrm{𝗍𝗁}}`$ component of $`\theta _0(𝐒_0\overline{𝐒}_0)`$ is $$_j\theta _0(𝐒_0\overline{𝐒}_0)=\frac{}{q_j}\theta _0(𝐒_0\overline{𝐒}_0)$$ with $`q_j=q_j(𝐒_0,\alpha )`$. From (3.10), we also have $`g𝐒_0(𝐒_0\|C_0(\alpha )\|\overline{𝐒}_0)=\delta (𝐒_0\overline{𝐒}_0).`$ (3.13) Like $`C_0(\alpha )`$, $`\theta _0`$ also denotes a square matrix whose matrix elements are $`(𝐒_0\|\theta _0\|\overline{𝐒}_0)=\theta _0(𝐒_0\overline{𝐒}_0).`$ (3.14) Matrix multiplications follow the usual rule; e.g., $`(𝐒_0\|C_0\theta _0\|\overline{𝐒}_0)={\displaystyle _0^{\mathrm{}}}𝑑\overline{\overline{𝐒}}_0(𝐒_0\|C_0\|\overline{\overline{𝐒}}_0)(\overline{\overline{𝐒}}_0\|\theta _0\|\overline{𝐒}_0).`$ (3.15) Introduce $`h_0^2(\alpha )`$ to be a diagonal matrix whose elements are $`(𝐒_0\|h_0^2(\alpha )\|\overline{𝐒}_0)=\delta (𝐒_0\overline{𝐒}_0)/(\overline{𝐒}_0)^2.`$ (3.16) Thus, $`(𝐒_0\|\theta _0h_0^2(\alpha )\|\overline{𝐒}_0)=\theta _0(𝐒_0\overline{𝐒}_0)/(\overline{𝐒}_0)^2,`$ (3.17) but $`(𝐒_0\|h_0^2(\alpha )\theta _0\|\overline{𝐒}_0)=\theta _0(𝐒_0\overline{𝐒}_0)/(𝐒_0)^2.`$ (3.18) Equations (3.10) and(3.13) can be written in their matrix forms $`C_0=g^1\theta _0h_0^2`$ (3.19) and $`g𝐒_0C_0=1,`$ (3.20) in which the $`\alpha `$-dependences of $`C_0`$ and $`h_0^2`$ are suppressed. Theorem 1. The solution of the wave equation (3.3) - (3.4) for $`e^\sigma `$ satisfies $`e^{\sigma (𝐒_0,\alpha )}`$ $`=`$ $`1+{\displaystyle _0^{\mathrm{}}}𝑑\overline{𝐒}_0(𝐒_0\|[1+C_0(\alpha )T]^1C_0(\alpha )\|\overline{𝐒}_0)`$ (3.22) $`[{\displaystyle \frac{g}{2}}^2\overline{𝐒}_0+E]e^{\sigma (\overline{𝐒}_0,\alpha )}.`$ Proof (i) We first establish the existence of the inverse matrix $`[1+C_0(\alpha )T]^1`$. Assuming the opposite, there would be a column matrix $`f(\alpha )`$ which satisfies $`[1+C_0(\alpha )T]f(\alpha )=0.`$ (3.23) Applying $`𝐒_0`$ from the left, using (3.20) and on account of $`T=\frac{1}{2}^2`$, we have for the matrix element $`f(𝐒_0,\alpha )`$ of $`f(\alpha )`$: $`g𝐒_0f(𝐒_0,\alpha ){\displaystyle \frac{1}{2}}^2f(𝐒_0,\alpha )=0.`$ (3.24) Hence, $`[e^{2g𝐒_0}f(𝐒_0,\alpha )]=0.`$ (3.25) Multiplying this equation by $`f(𝐒_0,\alpha )`$ and integrating over all space, we derive, through partial integration, $`{\displaystyle d^N𝐪e^{2g𝐒_0}[f(𝐒_0,\alpha )]^2}=0,`$ (3.26) provided $`e^{2g𝐒_0}ff=0\mathrm{𝖺𝗍}𝐒_0=\mathrm{}.`$ (3.27) It is clear that (3.26) implies $`f(𝐒_0,\alpha )=0`$; i.e., $`f(𝐒_0,\alpha )=`$ constant, which makes $`\frac{1}{2}^2f=0`$. It follows then from (3.23), $`f(𝐒_0,\alpha )=0.`$ (3.28) (ii) From (3.3), $`[g(𝐒_0)+T](e^\sigma 1)=[{\displaystyle \frac{g}{2}}^2𝐒_0+E]e^\sigma .`$ (3.29) Next, applying $`g(𝐒_0)`$ onto $`1+C_0T`$ and using (3.20), we find $`(g(𝐒_0))[1+C_0T]=g(𝐒_0)+T;`$ (3.30) therefore, (3.29) becomes $`[g(𝐒_0)][1+C_0T](e^\sigma 1)=[{\displaystyle \frac{g}{2}}^2𝐒_0+E]e^\sigma .`$ (3.31) On the other hand, assuming (3.22) we have $`e^\sigma 1=(1+C_0T)^1C_0[{\displaystyle \frac{g}{2}}^2𝐒_0+E]e^\sigma `$ (3.32) which satisfies (3.31). In deriving the above, we see that for $`f=e^\sigma 1`$, (3.27) is satisfied, which ensures the applicability of $`(1+C_0T)^1`$. The theorem is then proved. According to (3.19), $`C_0=O(g^1)`$, and therefore, for $`g^2>>1`$, we may expand $`(1+C_0T)^1C_0=C_0+(C_0T)C_0+(C_0T)^2C_0+\mathrm{}.`$ (3.33) Neglecting $`O(g^1)`$, $`\sigma 𝐒_1`$ and approximating $`(1+C_0T)^1C_0`$ by $`C_0`$, (3.32) reduces to $`e^{𝐒_1}1=C_0[{\displaystyle \frac{g}{2}}^2𝐒_0+E]e^{𝐒_1}.`$ (3.34) Taking its derivative through $`g𝐒_0`$, we have, as expected, $`g[𝐒_0𝐒_1]e^{𝐒_1}=[{\displaystyle \frac{g}{2}}^2𝐒_0E]e^{𝐒_1}`$ (3.35) which, upon the approximation $`EgE_0`$, gives the second line of (1.9). Likewise, we can derive (1.34) - (1.35) for $`𝐒_1`$, $`E_1`$, $`𝐒_2`$, $`E_2`$, $`\mathrm{}`$ from (3.32). ## 4. Green’s Functions and Perturbation Series As we have seen, the $`g^1`$ expansion enables us to probe the ground-state wave function $`\mathrm{\Phi }(𝐪)`$ along a classical trajectory; to each order $`g^n`$, the $`N`$-dimensional quantum wave function can be calculated through quadratures, by performing a finite number of definite integrals along the same trajectory. To the same order of accuracy, the corresponding energy is also determined by a finite number of differentiations at a single point where the potential has an absolute minimum, say at $`𝐪=\mathrm{𝟎}`$. This result may seem unfamiliar, since the Schroedinger wave function $`\mathrm{\Phi }(𝐪)`$ and its eigenvalue $`E`$ depend sensitively on the boundary conditions at $`\mathrm{}`$. But in our approach, the determination of $`\mathrm{\Phi }(𝐪)`$ is based only on a single trajectory connecting $`\mathrm{𝟎}`$ and $`𝐪`$. The boundary condition at infinity appears almost automatically, because the classical action (i.e., Hamilton’s action integral) satisfies the property: $`𝐒_0(𝐪)\mathrm{}\mathrm{𝖺𝗌}𝐪\mathrm{}.`$ (4.1) Since the classical action is calculated for a potential $`V(𝐪)`$, which has its maximum at $`0`$, this condition (4.1) imposes only a general restriction on the potential; in any case, it is considered to be an input in the new method. To explore further the essential features of the underlying mechanism, we shall in this section generalize our formalism without any explicit use of the Hamilton-Jacobi equation. In this section we shall derive several new forms of the Green’s function for the N-dimensional Schroedinger equation; these are all based on quadratures along a single specified trajectory. With the Green’s function we will be able to arrive at a new perturbation series expansion, different from the usual familiar ones. As we shall see (near the end of the section), the material discussed in previous sections 1-3 will appear as special cases. As mentioned in the Introduction, the new Green’s function will also enable us to extend the multi-instanton calculations developed for the one-dimensional problem in ref. to arbitrary dimensions. In what follows, we shall first give a number of definitions, then for clarity, arrange the main results in the form of three theorems (Theorem 2, 3 and 4 below). Examples of the new perturbation series and the Green’s function will be given in Section 5. Consider two $`N`$-dimensional Hamiltonians: $`H={\displaystyle \frac{1}{2}}^2+V(𝐪)`$ (4.2) with $`V(𝐪)0`$, as before, and $`=H+ϵU(𝐪).`$ (4.3) Their ground-state wave functions are $`\mathrm{\Phi }(𝐪)`$ and $`\mathrm{\Psi }(𝐪)`$ respectively: $`H\mathrm{\Phi }(𝐪)=E\mathrm{\Phi }(𝐪)`$ (4.4) and $`\mathrm{\Psi }(𝐪)=\mathrm{\Psi }(𝐪).`$ (4.5) Let $`ϵ\mathrm{\Delta }E.`$ (4.6) Write, as before, $`V(𝐪)=g^2v(𝐪),`$ (4.7) $`\mathrm{\Phi }(𝐪)=e^{g𝐒(𝐪)}`$ (4.8) and $`\mathrm{\Psi }(𝐪)=e^{g𝐒(𝐪)\tau (𝐪)}.`$ (4.9) We have, from (4.4) - (4.9), $`{\displaystyle \frac{g^2}{2}}(𝐒)^2{\displaystyle \frac{g}{2}}^2𝐒g^2v+E=0`$ (4.10) and $`g𝐒\tau +{\displaystyle \frac{1}{2}}[(\tau )^2^2\tau ]=ϵ(U\mathrm{\Delta }).`$ (4.11) In the subsequent discussions in this section, we assume that $`𝐒(𝐪)`$ is known, and our objective is to compute $`\tau (𝐪)`$, assuming that $`\|ϵ\|<<1.`$ (4.12) We assume further that $`𝐒(𝐪)`$ has an overall behavior similar to the Hamilton-Jacobi action $`𝐒_0(𝐪)`$ given by (1.22); in particular, $`𝐒(𝐪)`$ has an absolute minimum at $`𝐪=\mathrm{𝟎}`$. (This convention will be adopted here on.) Therefore $`𝐒=\mathrm{𝟎}\mathrm{𝖺𝗍}𝐪=\mathrm{𝟎}.`$ (4.13) (Note that from here on the absolute minimum of $`V(𝐪)`$ may or may not be at the same point.) Throughout this section, the role of the Hamilton-Jacobi action integral $`𝐒_0`$ in previous sections will be replaces by $`𝐒`$. As in (1.29) - (1.31), we introduce a set of $`N1`$ angular variables $`\beta =(\beta _1(𝐪),\beta _2(𝐪),\mathrm{},\beta _{N1}(𝐪)),`$ (4.14) which satisfy $`\beta _j𝐒=0`$ (4.15) with $`j=1,2,\mathrm{},N1.`$ (4.16) Each point $`𝐪`$ in the $`N`$-dimensional space will now be designated by $`(𝐒,\beta _1,\beta _2,\mathrm{},\beta _{N1}),`$ (4.17) instead of $`(q_1,q_2,q_3,\mathrm{},q_N)`$. At $`𝐪`$, let $`\stackrel{}{𝐒},\underset{1}{\overset{}{\beta }},\underset{2}{\overset{}{\beta }},\mathrm{},\underset{N1}{\overset{}{\beta }}`$ form a set of $`N`$ orthonormal unit vectors, which are the normals of the $`𝐒=`$ constant, $`\beta _1=`$ constant, $`\mathrm{}`$, $`\beta _{N1}=`$ constant surfaces at that point. Correspondingly, a line element $`d\stackrel{}{𝐪}`$ can be written as $`d\stackrel{}{𝐪}=\stackrel{}{𝐒}h_𝐒d𝐒+{\displaystyle \underset{j=1}{\overset{N1}{}}}\underset{j}{\overset{}{\beta }}h_jd\beta _j;`$ (4.18) the gradient is given by $`=\stackrel{}{𝐒}{\displaystyle \frac{1}{h_𝐒}}{\displaystyle \frac{}{𝐒}}+{\displaystyle \underset{j=1}{\overset{N1}{}}}\underset{j}{\overset{}{\beta }}{\displaystyle \frac{1}{h_j}}{\displaystyle \frac{}{\beta _j}},`$ (4.19) and $`T={\displaystyle \frac{1}{2}}^2`$ (4.20) can be decomposed into two parts: $`T=T_𝐒+T_\beta ,`$ (4.21) with $`T_𝐒={\displaystyle \frac{1}{2h_𝐒h_\beta }}{\displaystyle \frac{}{𝐒}}({\displaystyle \frac{h_\beta }{h_𝐒}}{\displaystyle \frac{}{𝐒}}),`$ (4.22) in which $`h_\beta ={\displaystyle \underset{j=1}{\overset{N1}{}}}h_j,`$ (4.23) and $`T_\beta ={\displaystyle \frac{1}{2h_𝐒h_\beta }}{\displaystyle \underset{j=1}{\overset{N1}{}}}{\displaystyle \frac{}{\beta _j}}({\displaystyle \frac{h_𝐒h_\beta }{h_j^2}}{\displaystyle \frac{}{\beta _j}}).`$ (4.24) It follows from (4.18), $`h_𝐒^2=[(𝐒)^2]^1,h_1^2=[(\beta _1)^2]^1,\mathrm{},h_j^2=[(\beta _j)^2]^1,\mathrm{}.`$ (4.25) The volume element in the $`𝐪`$-space is $`d^N𝐪=h_𝐒h_\beta d𝐒d\beta `$ (4.26) where $`d\beta ={\displaystyle \underset{j=1}{\overset{N1}{}}}d\beta _j.`$ (4.27) The trajectory of a given $`\beta `$ defines a continuous curve normal to the $`𝐒=`$ constant surfaces; along any such trajectory we will introduce the following matrices, $`\theta `$, $`C(\beta )`$, $`D(\beta )`$ and $`G(\beta )`$, each of which has its matrix element connecting two points $`(𝐒,\beta )`$ and $`(\overline{𝐒},\beta )`$ along the same trajectory. The matrix element of $`\theta `$ is given by $`(𝐒\|\theta \|\overline{𝐒})=\theta (𝐒\overline{𝐒})=\{\begin{array}{cc}1& \mathrm{𝗂𝖿}0\overline{𝐒}<𝐒\\ 0& \mathrm{𝗂𝖿}0𝐒<\overline{𝐒},\end{array}`$ (4.30) We now define $`C`$, $`D`$ and $`G`$ in matrix notations by (cf. (3.15)-(3.19)) $`C`$ $``$ $`g^1\theta h_𝐒^2,`$ (4.31) $`D`$ $``$ $`2e^{g𝐒}\theta e^{2g𝐒}{\displaystyle \frac{h_𝐒}{h_\beta }}\theta e^{g𝐒}h_𝐒h_\beta `$ (4.32) and $`GD(1+T_\beta D)^1=(1+DT_\beta )^1D,`$ (4.33) in which the $`\beta `$-dependence of the matrices is not exhibited explicitly and $`e^{\pm g𝐒}`$, $`h_𝐒`$, $`h_\beta `$ are all diagonal matrices. (The conditions for the existence of inverse matrices, like $`(1+T_\beta D)^1`$ will be given later, after (4.70).) For example, the matrix element of $`C`$ is $`(𝐒\|C\|\overline{𝐒})=g^1\theta (𝐒\overline{𝐒})/(\overline{𝐒})^2.`$ (4.34) Similar to (3.15) and (3.17)-(3.18), the product of any two matrices, say $`A`$ and $`B`$, defined along the same trajectory (i.e., sharing the same $`\beta `$) has its matrix element given by the usual rule $`(𝐒\|AB\|\overline{𝐒})={\displaystyle _0^{\mathrm{}}}𝑑\overline{\overline{𝐒}}(𝐒\|A\|\overline{\overline{𝐒}})(\overline{\overline{𝐒}}\|B\|\overline{𝐒}).`$ (4.35) . The matrix element of its derivative $`\frac{}{𝐒}AB`$ is given by $`(𝐒\|{\displaystyle \frac{}{𝐒}}AB\|\overline{𝐒})={\displaystyle _0^{\mathrm{}}}𝑑\overline{\overline{𝐒}}{\displaystyle \frac{}{𝐒}}(𝐒\|A\|\overline{\overline{𝐒}})(\overline{\overline{𝐒}}\|B\|\overline{𝐒}),`$ (4.36) with $`\frac{}{𝐒}`$ operating only on the left index $`𝐒`$, whereas that of $`\frac{}{\beta _j}AB`$ is $`(𝐒\|{\displaystyle \frac{}{\beta _j}}AB\|\overline{𝐒})=(𝐒\|({\displaystyle \frac{A}{\beta _j}})B+A({\displaystyle \frac{B}{\beta _j}})\|\overline{𝐒}).`$ (4.37) Theorem 2. $`1.g𝐒C=𝖨,`$ (4.38) $`2.D\mathrm{𝖺𝗇𝖽}G\mathrm{𝗌𝖺𝗍𝗂𝗌𝖿𝗒}`$ (4.39) $`D=e^{g𝐒}(1+CT_𝐒)^1Ce^{g𝐒}=e^{g𝐒}C(1+T_𝐒C)^1e^{g𝐒}`$ (4.40) $`\mathrm{𝖺𝗇𝖽}`$ (4.41) $`G=e^{g𝐒}(1+CT)^1Ce^{g𝐒}=e^{g𝐒}C(1+TC)^1e^{g𝐒}.`$ (4.42) 3. Furthermore, $`D`$ and $`G`$ are the Green’s functions of $`T_𝐒+VE`$ and $`T+VE`$: $`(T_𝐒+VE)D`$ $`=`$ $`𝖨,`$ (4.43) $`(HE)G`$ $`=`$ $`(T+VE)G=𝖨,`$ (4.44) where I denotes the unit matrix whose matrix element is $`\delta (𝐒\overline{𝐒})`$. Proof 1. For $`𝐒>0`$, $`{\displaystyle \frac{}{𝐒}}\theta (𝐒\overline{𝐒})=\delta (𝐒\overline{𝐒}).`$ (4.45) On account of the orthogonality $`𝐒\beta _j=0`$, we have $`𝐒=(𝐒)^2{\displaystyle \frac{}{𝐒}}.`$ (4.46) Eq.(4.38) follows. 2. Define $`\overline{D}e^{g𝐒}De^{g𝐒}\mathrm{𝖺𝗇𝖽}\overline{G}e^{g𝐒}Ge^{g𝐒}.`$ (4.47) Then, from (4.31) we see that $`\overline{D}=2\theta e^{2g𝐒}{\displaystyle \frac{h_𝐒}{h_\beta }}\theta e^{2g𝐒}h_𝐒h_\beta ,`$ (4.48) $`{\displaystyle \frac{\overline{D}}{𝐒}}=2e^{2g𝐒}{\displaystyle \frac{h_𝐒}{h_\beta }}\theta e^{2g𝐒}h_𝐒h_\beta ,`$ (4.49) $`{\displaystyle \frac{}{𝐒}}[{\displaystyle \frac{h_\beta }{h_𝐒}}{\displaystyle \frac{\overline{D}}{𝐒}}]=4ge^{2g𝐒}\theta e^{2g𝐒}h_𝐒h_\beta 2h_𝐒h_\beta .`$ (4.50) and, because of (4.22), $`T_𝐒\overline{D}={\displaystyle \frac{2g}{h_𝐒h_\beta }}e^{2g𝐒}\theta e^{2g𝐒}h_𝐒h_\beta +1.`$ (4.51) Multiplying by $`C`$ on the left, we find $`CT_𝐒\overline{D}=2\theta {\displaystyle \frac{h_𝐒}{h_\beta }}e^{2g𝐒}\theta e^{2g𝐒}h_𝐒h_\beta +C=\overline{D}+C,`$ (4.52) and therefore $`(1+CT_𝐒)\overline{D}=C,`$ (4.53) or $`\overline{D}=(1+CT_𝐒)^1C=C(1+T_𝐒C)^1.`$ (4.54) Since according to (4.43) $`D=e^{g𝐒}\overline{D}e^{g𝐒},`$ (4.55) (4.37) is proved. By using (4.31), we derive $`{\displaystyle \frac{D}{𝐒}}=gD2e^{g𝐒}{\displaystyle \frac{h_𝐒}{h_\beta }}\theta e^{g𝐒}h_𝐒h_\beta ,`$ (4.56) $`{\displaystyle \frac{h_\beta }{h_𝐒}}{\displaystyle \frac{D}{𝐒}}=g{\displaystyle \frac{h_\beta }{h_𝐒}}D2e^{g𝐒}\theta e^{g𝐒}h_𝐒h_\beta ,`$ (4.57) and $`{\displaystyle \frac{}{𝐒}}({\displaystyle \frac{h_\beta }{h_𝐒}}{\displaystyle \frac{D}{𝐒}})=g\{[{\displaystyle \frac{}{𝐒}}({\displaystyle \frac{h_\beta }{h_𝐒}})]g{\displaystyle \frac{h_\beta }{h_𝐒}}\}D2h_𝐒h_\beta .`$ (4.58) On the other hand, because $`\mathrm{\Phi }=e^{g𝐒}`$ satisfies $`(T_𝐒+VE)\mathrm{\Phi }=0`$ and since $`T_𝐒\mathrm{\Phi }={\displaystyle \frac{1}{2h_𝐒h_\beta }}{\displaystyle \frac{}{𝐒}}({\displaystyle \frac{h_\beta }{h_𝐒}}{\displaystyle \frac{}{𝐒}})e^{g𝐒}={\displaystyle \frac{1}{2h_𝐒h_\beta }}\{g[{\displaystyle \frac{}{𝐒}}({\displaystyle \frac{h_\beta }{h_𝐒}})]+g^2{\displaystyle \frac{h_\beta }{h_𝐒}}\}e^{g𝐒},`$ (4.59) we have $`{\displaystyle \frac{g}{2h_𝐒h_\beta }}\{[{\displaystyle \frac{}{𝐒}}({\displaystyle \frac{h_\beta }{h_𝐒}})]g{\displaystyle \frac{h_\beta }{h_𝐒}}\}=V+E.`$ (4.60) Therefore (4.58) leads to $`{\displaystyle \frac{1}{2h_𝐒h_\beta }}{\displaystyle \frac{}{𝐒}}({\displaystyle \frac{h_\beta }{h_𝐒}}{\displaystyle \frac{D}{𝐒}})+(VE)D=1;`$ (4.61) i.e., $`D`$ satisfies (4.43) and is the Green’s function of $`T_𝐒+VE`$. From (4.33) and the identity $`1=[1T_\beta D(1+T_\beta D)^1](1+T_\beta D)`$ (4.62) we have $`(T_𝐒+VE)D=[1T_\beta D(1+T_\beta D)^1](1+T_\beta D),`$ (4.63) i.e., $`(T_𝐒+VE)D(1+T_\beta D)^1=1T_\beta D(1+T_\beta D)^1.`$ (4.64) From (4.33), $`D(1+T_\beta D)^1`$ is $`G`$; we derive $`(T+VE)G=1,`$ (4.65) which establishes (4.40). From $`\overline{G}=e^{g𝐒}Ge^{g𝐒}`$ and since $`e^{g𝐒}`$ commutes with $`T_\beta `$, it follows that $`\overline{G}=\overline{D}(1+T_\beta \overline{D})^1.`$ (4.66) By using $`\overline{D}=C(1+T_𝐒C)^1`$, we have $`1+T_\beta \overline{D}`$ $`=`$ $`1+T_\beta C(1+T_𝐒C)^1`$ (4.67) $`=`$ $`(1+T_𝐒C+T_\beta C)(1+T_𝐒C)^1.`$ (4.68) Its inverse is $`(1+T_\beta \overline{D})^1=(1+T_𝐒C)(1+TC)^1`$ (4.69) which leads to $`\overline{G}=\overline{D}(1+T_𝐒C)(1+TC)^1=C(1+TC)^1,`$ (4.70) and that gives (4.38). From the above expression, we see that $`\overline{G}^1=C^1+T.(4.38)^{}`$ (4.71) Likewise, from (4.49) and (4.61), we can write $`\overline{D}^1=C^1+T_𝐒(4.37)^{}`$ (4.72) and $`\overline{G}^1=\overline{D}^1+T_\beta .(4.31)^{}`$ (4.73) These relations make transparent the connections between the original equations (4.38), (4.37) and (4.31). To complete the proof, we have to examine the conditions, under which $`1+CT`$, $`1+CT_𝐒`$ and $`1+DT_\beta `$ have inverses. As in (3.23), we assume $`f_1=f_1(𝐒,\beta )`$ to satisfy $`(1+CT)f_1=0.`$ (4.74) Operating on its left by $`2g𝐒`$, we derive $`2g𝐒f_1+^2f_1=0.`$ (4.75) Hence $`(e^{2g𝐒}f_1)=0.`$ (4.76) Multiplying by $`f_1d^N𝐪`$ and integrating over all space, we have $`{\displaystyle e^{2g𝐒}(f_1)^2d^N𝐪}=0,`$ (4.77) provided $`e^{2g𝐒}f_1f_1=0\mathrm{𝖺𝗍}\mathrm{}.`$ (4.78) From (4.77), we see that the only solution is $`f_1=0`$; i.e., $`f_1=`$ constant. However, from the original equation (4.74), it follows then $`f_1=0`$. Consequently, $`(1+CT)^1`$ exists provided (4.78) holds. Next, we examine the equation $`(1+CT_𝐒)f_2=0.`$ (4.79) Operating $`g𝐒`$ on its left, we find $`g𝐒f_2+T_𝐒f_2=0.`$ (4.80) Multiplying (4.80) by $`f_2h_𝐒h_\beta d𝐒`$ and integrating from $`𝐒=0`$ to $`\mathrm{}`$, we have $`{\displaystyle _0^{\mathrm{}}}f_2{\displaystyle \frac{}{𝐒}}({\displaystyle \frac{h_\beta }{h_𝐒}}e^{2g𝐒}{\displaystyle \frac{f_2}{𝐒}})𝑑𝐒=0.`$ (4.81) Assuming $`(f_2{\displaystyle \frac{1}{h_𝐒}}{\displaystyle \frac{f_2}{𝐒}})e^{2g𝐒}h_\beta =0\mathrm{𝖺𝗍}𝐒=0\mathrm{𝖺𝗇𝖽}𝐒=\mathrm{},`$ (4.82) (which is the $`𝐒`$-component equivalence of (4.78)), we derive $`{\displaystyle _0^{\mathrm{}}}e^{2g𝐒}({\displaystyle \frac{1}{h_𝐒}}{\displaystyle \frac{f_2}{𝐒}})^2h_\beta h_𝐒𝑑𝐒=0.`$ (4.83) Recognizing $`d^N𝐪=h_𝐒h_\beta d𝐒d\beta `$ is positive, it follows then $`f_2/𝐒=0`$. Substituting this back to (4.79), we find $`f_2=0.`$ (4.84) Hence, $`(1+CT_𝐒)^1`$ exists, provided (4.82) holds. Lastly, we consider $`(1+DT_\beta )f_3=0.`$ (4.85) Applying $`T_𝐒+VE`$ from the left and using (4.43), we obtain $`(T_𝐒+T_\beta +VE)f_3=0.`$ (4.86) Thus, the inverse of $`(1+DT_\beta )`$ exists, provided we restrict ourselves to a Hilbert space which excludes the ground-state of $`H`$. Theorem 2 is then proved. It is important to note that the existence of $`(1+CT)^1`$, and therefore $`G`$, as defined by (4.38), is free from the last restriction. We return now to (4.3)-(4.6). The following theorem will enable us to derive the ground-state wave function $`\mathrm{\Psi }=e^{g𝐒\tau }`$ of $`H+ϵU`$ in terms of $`e^{g𝐒}`$, the ground-state wave function of $`H`$: Theorem 3 $`e^\tau =1+(1+CT)^1Cϵ(U+\mathrm{\Delta })e^\tau ,`$ (4.87) or equivalently, $`\mathrm{\Psi }`$ $`=`$ $`e^{g𝐒}+(1+DT_\beta )^1Dϵ(U+\mathrm{\Delta })\mathrm{\Psi }`$ (4.88) $`=`$ $`e^{g𝐒}+Gϵ(U+\mathrm{\Delta })\mathrm{\Psi }.`$ (4.89) Proof. Eq.(4.11) can also be written as $`(g𝐒+T)e^\tau =ϵ(U+\mathrm{\Delta })e^\tau ,`$ (4.90) where $`T=\frac{1}{2}^2`$, as before. Since from (4.38) $`g𝐒C=1`$, and since $`(e^\tau 1)=e^\tau `$, the above expression gives $`g𝐒(1+CT)(e^\tau 1)=ϵ(U+\mathrm{\Delta })e^\tau .`$ (4.91) Hence, $`e^\tau 1=(1+CT)^1Cϵ(U+\mathrm{\Delta })e^\tau ,`$ (4.92) as can be derived by substituting (4.92) into (4.91). Thus (4.87) follows. Multiplying (4.92) by $`e^{g𝐒}`$ on the left gives $`\mathrm{\Psi }=e^{g𝐒}+e^{g𝐒}(1+CT)^1Cϵ(U+\mathrm{\Delta })e^\tau `$ (4.93) which leads to (4.88), on account of (4.33) and (4.38). Theorem 3 is thereby established. From (4.87), we may expand $`e^\tau `$ as a power series of $`ϵ`$: $`e^\tau =1+ϵ(1+CT)^1C(U+\mathrm{\Delta })+ϵ^2[(1+CT)^1C(U+\mathrm{\Delta })]^2`$ (4.94) $`+\mathrm{}+ϵ^n[(1+CT)^1C(U+\mathrm{\Delta })]^n+\mathrm{}.`$ (4.95) Likewise, from (4.88), we derive the perturbation series $`\mathrm{\Psi }=\mathrm{\Psi }_0+ϵ\mathrm{\Psi }_1+ϵ^2\mathrm{\Psi }_2+\mathrm{}+ϵ^n\mathrm{\Psi }_n+\mathrm{},`$ (4.96) where $`\mathrm{\Psi }_0=e^{g𝐒},`$ (4.97) and for $`n1`$, $`\mathrm{\Psi }_n=[G(U+\mathrm{\Delta })]^ne^{g𝐒},`$ (4.98) with $`G=(1+DT_\beta )^1D`$, or $`e^{g𝐒}(1+CT)^1Ce^{g𝐒}`$. From (4.98), we have $`\mathrm{\Psi }_n=G(U+\mathrm{\Delta })\mathrm{\Psi }_{n1}.`$ (4.99) Eqs. (4.88) and (4.98) can be directly established by using $`(HE)G=1`$; i.e., $`G`$ being the Green’s function of $`(HE)`$. Likewise, since $`(T_𝐒+VE)e^{g𝐒}=0`$ (4.100) and $`(T_𝐒+VE)\mathrm{\Psi }=[T_\beta +ϵ(U+\mathrm{\Delta })]\mathrm{\Psi },`$ (4.101) $`\mathrm{\Psi }`$ also satisfies $`\mathrm{\Psi }=e^{g𝐒}+D(T_\beta +ϵ(U+\mathrm{\Delta })]\mathrm{\Psi },`$ (4.102) on account of (4.43), $`D`$ being the Green’s function of $`(T_𝐒+VE)`$. The new perturbation series (4.96)-(4.98) is most effective, when it is complemented by an additional expansion in powers of $`g^1`$, the anharmonicity parameter. From $`C=g^1\theta h_𝐒^2`$ and (4.37)-(4.38), we see that the three Green’s functions $`C`$, $`D`$ and $`G`$ are all $`O(g^1)`$. The n<sup>th</sup> order perturbative wave function $`ϵ^n\mathrm{\Psi }_n`$ is given by (4.98). By using (4.32) for $`G`$, $`\mathrm{\Psi }_n`$ can be written as a power series in $`D`$, and therefore also in $`g^1`$: $`\mathrm{\Psi }_n=\{[DDT_\beta D+(DT_\beta )^2D\mathrm{}+()^m(DT_\beta )^mD+\mathrm{}](U+\mathrm{\Delta })\}^ne^{g𝐒},`$ (4.103) in which each $`D`$ assumes the form given by (4.31); i.e., $$D=2e^{g𝐒}\theta e^{2g𝐒}\frac{h_𝐒}{h_\beta }\theta e^{g𝐒}h_𝐒h_\beta ,$$ a double definite integral along the trajectory. Through (4.38), $`\mathrm{\Psi }_n`$ can also be expressed as a different power series in $`g^1`$ in terms of $`C`$: $`\mathrm{\Psi }_n=e^{g𝐒}\{[CCTC+\mathrm{}+()^m(CT)^mC+\mathrm{}](U+\mathrm{\Delta })\}^n,`$ (4.104) in which each $`C`$ is a single definite integral along the same trajectory. The equivalence between these two expressions rests on (4.37), which gives $$D=e^{g𝐒}[CCT_𝐒C+\mathrm{}+()^m(CT_𝐒)^mC+\mathrm{}]e^{g𝐒}.$$ The identity between this power series expression of $`D`$ and its double integral form given above has its origin in the Sturm-Liouville construction of the Green’s function for a one-dimensional Schroedinger equation. Recalling that $`e^{g𝐒}`$ is a solution of the second order ordinary differential equation in $`𝐒`$, $`(T_𝐒+VE)e^{g𝐒}=0,`$ (4.105) we can form an irregular solution $`F`$ which satisfies the same equation $`(T_𝐒+VE)F=0`$ (4.106) with $`F(𝐒)=e^{g𝐒}{\displaystyle \underset{0}{\overset{𝐒}{}}}e^{2g\overline{𝐒}}({\displaystyle \frac{h_𝐒}{h_\beta }})_{\mathrm{𝖺𝗍}\overline{𝐒}}𝑑\overline{𝐒},`$ (4.107) in which the angular variable $`\beta `$ is kept fixed and not exhibited explicitly in the argument. The Green’s function $`D`$ is related to the Wronskian-type expression by $`(𝐒\|D\|\overline{𝐒})=2[e^{g𝐒}F(\overline{𝐒})F(𝐒)e^{g\overline{𝐒}}](h_𝐒h_\beta )_{\overline{𝐒},\beta }\theta (𝐒\overline{𝐒}).`$ (4.108) In ref., we used a one-dimensional example to demonstrate how this form of the Green’s function can enable us to evaluate the multi-instanton expansion. Here we have extended the application of a one-dimensional Green’s function to an $`N`$-dimensional problem. Theorem 4. The energy shift $`ϵ\mathrm{\Delta }`$ is given by $`ϵ\mathrm{\Delta }={\displaystyle \frac{e^{g𝐒}ϵU\mathrm{\Psi }d^N𝐪}{e^{g𝐒}\mathrm{\Psi }d^N𝐪}}`$ (4.109) where $`d^N𝐪=h_𝐒h_\beta d𝐒d\beta `$, as in (4.26) and the integration is extended over all space. Proof From (4.31) or (4.91), the matrix element of $`D(\beta )`$ is $`(𝐒\|D(\beta )\|\overline{𝐒})=2e^{g𝐒}{\displaystyle _0^𝐒}𝑑\overline{\overline{𝐒}}({\displaystyle \frac{h_𝐒}{h_\beta }})_{\overline{\overline{𝐒}},\beta }\theta (\overline{\overline{𝐒}}\overline{𝐒})e^{2g\overline{\overline{𝐒}}}e^{g\overline{𝐒}}(h_𝐒h_\beta )_{\overline{𝐒},\beta },`$ (4.110) where the subscripts $`\overline{\overline{𝐒}},\beta `$ and $`\overline{𝐒},\beta `$ indicate the arguments of the functions inside the parenthesis. Let $`\chi `$ denotes the column matrix $`\chi (𝐒,\beta )[T_\beta +ϵ(U+\mathrm{\Delta })]\mathrm{\Psi }(𝐒,\beta ).`$ (4.111) Eq.(4.102) becomes $`\mathrm{\Psi }=e^{g𝐒}+D\chi .`$ (4.112) As $`𝐒\mathrm{}`$, the left side should be zero, but the right side is, on account of (4.110), controlled by the $`\overline{\overline{𝐒}}`$-integration in the upper region when $`\overline{\overline{𝐒}}`$ is near $`𝐒`$; i.e., writing $`D\chi `$ as the product of $`2e^{g𝐒}`$ times $`{\displaystyle _0^𝐒}𝑑\overline{\overline{𝐒}}e^{2g\overline{\overline{𝐒}}}({\displaystyle \frac{h_𝐒}{h_\beta }})_{\overline{\overline{𝐒}},\beta }{\displaystyle _0^{\overline{\overline{𝐒}}}}𝑑\overline{𝐒}e^{g\overline{𝐒}}(h_𝐒h_\beta )_{\overline{𝐒},\beta }\chi (\overline{𝐒},\beta ),`$ (4.113) we see that the factor to the right of (and including) the $`\overline{𝐒}`$integration sign approaches $`e^{2g𝐒}`$ times $`{\displaystyle _0^{\mathrm{}}}𝑑\overline{𝐒}e^{g\overline{𝐒}}(h_𝐒h_\beta )_{\overline{𝐒},\beta }\chi (\overline{𝐒},\beta ),`$ (4.114) which would make $`D\chi \mathrm{}`$ unless (4.114) is zero. Thus, the convergence of $`\mathrm{\Psi }`$ at $`\mathrm{}`$ requires $`{\displaystyle _0^{\mathrm{}}}𝑑𝐒e^{g𝐒}h_𝐒h_\beta [T_\beta +ϵ(U+\mathrm{\Delta })]\mathrm{\Psi }=0`$ (4.115) at all $`\beta `$. Integrating over the angular variables $`d\beta =_{j=1}^{N1}d\beta _j,`$ and because of $`h_𝐒h_\beta T_\beta ={\displaystyle \frac{1}{2}}{\displaystyle \underset{j=1}{\overset{N1}{}}}{\displaystyle \frac{}{\beta _j}}{\displaystyle \frac{h_𝐒h_\beta }{h_{\beta _j}^2}}{\displaystyle \frac{}{\beta _j}}`$ (4.116) and $$𝑑\beta h_𝐒h_\beta T_\beta \mathrm{\Psi }=0$$ we derive (4.109). This result is well known. (Since $`(H+ϵU)\mathrm{\Psi }=(E+ϵ\mathrm{\Delta })\mathrm{\Psi }`$, the multiplication of $`e^{g𝐒}`$ on the left and integrating over all space gives immediately (4.109).) We have demonstrated that the same result can also be derived, though somewhat awkwardly, by the Green’s function method developed in this paper. The usual perturbation series requires to each order $`n`$, either a product of infinite sums over all excited levels of the unperturbed Hamiltonian $`H`$, or a sum over all possible paths through Feynman’s path integration method. The perturbation series formula derived here is different; it depends only on quadratures along a single trajectory, as will be illustrated by examples in the next section. Remarks. Before leaving this section, we wish to address the connection between the perturbation series developed in this section and the previous series expansion given by (3.32)-(3.33). In sections 1-3, we discuss the relation between the Schroedinger wave function $$e^{g𝐒}=e^{g𝐒_0\sigma }$$ and the classical Hamilton’s action integral $`𝐒_0`$, which satisfies $$g^2(𝐒_0)^2=2V=2g^2v.$$ Thus formally we may regard $`e^{g𝐒_0}`$ as the solution of a Schroedinger-like wave equation whose eigenvalue happens to be $`0`$: $`({\displaystyle \frac{1}{2}}^2+V_c)e^{g𝐒_0}=0`$ (4.117) where $`V_c=g^2v{\displaystyle \frac{1}{2}}g^2𝐒_0.`$ (4.118) Since $`e^{g𝐒}`$ satisfies $`({\displaystyle \frac{1}{2}}^2+V_c+{\displaystyle \frac{1}{2}}g^2𝐒_0)e^{g𝐒}=Ee^{g𝐒},`$ (4.119) we may treat (4.117) as the unperturbed Schroedinger equation, and (4.119) as the perturbed one, with $`\frac{1}{2}g^2𝐒_0`$ as the perturbative Hamiltonian and $`E`$ as the energy shift. The corresponding small $`ϵ`$-parameter that characterize the perturbation becomes $`g^1`$, since $`\frac{1}{2}g^2𝐒_0`$ is $`O(g^1)`$ times $`V_c`$. This explains the similarity between (3.32) for $`e^\sigma `$ and (4.87) for $`e^\tau `$; it also identifies the $`g^1`$-expansion with the perturbative $`ϵ`$-expansion. ## 5. Examples From (4.33) and (4.38), the Green’s function $`G`$ can be either expressed in terms of the single integral form $`C`$, or equivalently, the double integral form $`D`$: $`G=e^{g𝐒}(1+CT)^1Ce^{g𝐒}`$ (5.1) or $`G=(1+DT_\beta )^1D.`$ (5.2) These two complement each other in their applications. As in (4.8) - (4.9), let $`e^{g𝐒}`$ be the ground-state wave function of $`H`$. The corresponding ground-state wave function $`\mathrm{\Psi }=e^{g𝐒\tau }`$ (5.3) of $`H+ϵU`$ can also be expressed either in terms of $`C`$ (eq.(4.87)): $`e^\tau =1+(1+CT)^1Cϵ(U+\mathrm{\Delta })e^\tau ,`$ (5.4) or in terms of $`D`$ (eq.(4.88)): $`\mathrm{\Psi }=e^{g𝐒}+(1+DT_\beta )^1Dϵ(U+\mathrm{\Delta })\mathrm{\Psi }.`$ (5.5) As shown in (4.94) and (4.103), both can be used for the perturbation series expansion. In this section, we shall examine several examples which illustrate the different merits of these two formulations. For clarity, consider the one-dimensional case, in which (5.2) and (5.5) reduce to $`G=D`$ and $`\mathrm{\Psi }=e^{g𝐒}+Dϵ(U+\mathrm{\Delta })\mathrm{\Psi }.`$ (5.6) As in (4.47), introduce $`\overline{D}e^{g𝐒}De^{g𝐒},`$ (5.7) which, in accordance with (4.53) and since $`T_𝐒=T`$ in the one-dimensional case, is related to $`C`$ by $`(1+CT)\overline{D}=C.`$ (5.8) In terms of $`\overline{D}`$, (5.6) becomes $`e^\tau =1+\overline{D}ϵ(U+\mathrm{\Delta })e^\tau .`$ (5.9) In order that the above expression is equivalent to the alternative form (5.4) in terms of $`C`$, we need to convert (5.8) into $`\overline{D}=(1+CT)^1C.`$ (5.10) According to (4.78), this requires $`e^{2g𝐒}f_1f_1=0`$ (5.11) at the boundary, where, for (5.9) , the function $`f_1`$ is $`f_1=\overline{D}(U+\mathrm{\Delta })e^\tau ,`$ (5.12) and in one-dimension, the boundary consists of $`x=\pm \mathrm{}.`$ (5.13) We recall further that $`C`$ satisfies (4.38) which is now an ordinary first order differential equation $`g{\displaystyle \frac{d𝐒}{dx}}{\displaystyle \frac{dC}{dx}}=1;`$ (5.14) this determines $`C`$ up to an integration constant, which can be set by requiring $`C=0\mathrm{𝖺𝗍}𝐒=0.`$ (5.15) Now, $`D`$ is the solution of the second order differential equation $`({\displaystyle \frac{1}{2}}{\displaystyle \frac{d^2}{dx^2}}+VE)D=1,`$ (5.16) and therefore contains two independent integration constants. In view of (5.10) and (5.15), we require $`\overline{D}=0\mathrm{𝖺𝗍}𝐒=0`$ (5.17) and write for the right hand side of (5.12) $`\overline{D}(U+\mathrm{\Delta })e^\tau =2{\displaystyle \underset{0}{\overset{x}{}}}e^{2g𝐒(y)}𝑑y{\displaystyle \underset{\mathrm{}}{\overset{y}{}}}e^{2g𝐒(z)}(U+\mathrm{\Delta })e^{\tau (z)}𝑑z`$ (5.18) so that (5.11) holds at the boundary $`x=\mathrm{}`$. (For convenience, we choose $`𝐒=0`$ at $`x=0`$.) In order to satisfy (5.11) at the other boundary $`x=\mathrm{}`$, we set $`{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}e^{2g𝐒(x)\tau (x)}(U+\mathrm{\Delta })𝑑x=0;`$ (5.19) otherwise, (5.18) would $`e^{2g𝐒}`$ at $`x=\mathrm{}`$, violating (5.11) - (5.12). Thus, the energy shift $`\mathrm{\Delta }`$ is also determined, and (5.19) is a special case of (4.109). On the other hand, if we use (5.1), expressing $`G`$ in terms of the single integral form $`C`$, then $`\mathrm{\Delta }`$ would be determined by differentiations at $`𝐒=\mathrm{𝟎}`$, similar to (1.35) and (1.34). In the following, we shall solve two examples by using $`G`$, first in terms of $`C`$, then in terms of $`D`$. Example 1. Consider a one-dimensional harmonic oscillator with $`H={\displaystyle \frac{1}{2}}{\displaystyle \frac{d^2}{dx^2}}+{\displaystyle \frac{g^2}{2}}x^2`$ (5.20) and a perturbative potential $`ϵU=ϵx^{2p},`$ (5.21) where $`p`$ is a positive integer. The unperturbed ground-state wave function is $`\mathrm{\Phi }=e^{gS}=e^{\frac{g}{2}x^2}.`$ (5.22) Since $`S=\frac{1}{2}x^2`$, we have $`h_S^2=(dS/dx)^2=x^2=(2S)^1.`$ (5.23) For $`n>0`$ and in accordance with (5.15), $`Cx^{2n}={\displaystyle _0^S}{\displaystyle \frac{dS}{g2S}}x^{2n}={\displaystyle \frac{1}{g(2n)}}x^{2n}.`$ (5.24) Lemma 1 $`(1+CT)^1C[x^{2n}\mathrm{\Gamma }_{1n}]={\displaystyle \underset{m=1}{\overset{n}{}}}\mathrm{\Gamma }_{mn}x^{2m}.`$ (5.25) where $`\mathrm{\Gamma }_{nn}`$ $`=`$ $`{\displaystyle \frac{1}{g2n}},\mathrm{\Gamma }_{n1n}={\displaystyle \frac{2n1}{2g^2(2n2)}},`$ (5.26) $`\mathrm{\Gamma }_{mn}`$ $`=`$ $`{\displaystyle \frac{(2n1)(2n3)\mathrm{}(2m+1)}{m(2g)^{nm+1}}}\mathrm{𝖿𝗈𝗋}1mn1`$ (5.27) and in particular $`\mathrm{\Gamma }_{1n}={\displaystyle \frac{(2n1)!!}{(2g)^n}}.`$ (5.28) Proof Since $`T=\frac{1}{2}\frac{d^2}{dx^2}`$, $$(1+CT)^1C=CCTC+(CT)^2C+\mathrm{}+(CT)^mC+\mathrm{}$$ and $$TCx^{2n}=\frac{1}{2g}(2n1)x^{2n2},$$ we find $`CTCx^{2n}={\displaystyle \frac{2n1}{2g^2(2n2)}}x^{2n2}=\mathrm{\Gamma }_{n1n}x^{2n2},`$ (5.29) $`(CT)^2Cx^{2n}={\displaystyle \frac{(2n1)(2n3)}{2^2g^3(2n4)}}x^{2n4}=\mathrm{\Gamma }_{n2n}x^{2n4},`$ (5.30) $`\mathrm{}`$, until $`(CT)^{n1}Cx^{2n}={\displaystyle \frac{(2n1)!!}{(2g)^n}}x^2=\mathrm{\Gamma }_{1n}x^2.`$ (5.31) But, $`(CT)^nCx^{2n}=C\mathrm{\Gamma }_{1n}`$, which by itself would be $`\mathrm{}`$. On the other hand, $`(CT)^nCx^{2n}C\mathrm{\Gamma }_{1n}=0;`$ (5.32) Lemma 1 is then proved. As in (4.9), the eigenstate of $`={\displaystyle \frac{1}{2}}{\displaystyle \frac{d^2}{dx^2}}+{\displaystyle \frac{g^2}{2}}x^2+ϵx^{2p}`$ (5.33) is $`e^{\frac{1}{2}gx^2\tau }`$. Write $`e^\tau =1+{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}a_lx^{2l}.`$ (5.34) From (4.87) $`{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}a_lx^{2l}=(1+CT)^1Cϵ(x^{2p}+\mathrm{\Delta })(1+{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}a_lx^{2l}).`$ (5.35) By using (5.25) and (5.34), we derive $`\mathrm{\Delta }=\mathrm{\Gamma }_{1p}+{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}a_l(\mathrm{\Gamma }_{1l+p}\mathrm{\Delta }\mathrm{\Gamma }_{1l})`$ (5.36) and for $`n1`$ $`a_n=ϵ\mathrm{\Gamma }_{np}ϵ{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}a_l(\mathrm{\Gamma }_{nl+p}\mathrm{\Delta }\mathrm{\Gamma }_{nl}).`$ (5.37) where $`\mathrm{\Gamma }_{mn}=0`$ for $`m>n`$; otherwise it is given by (5.26). From (5.36),we see that $`ϵ\mathrm{\Delta }=a_1.`$ (5.38) Expand $`ϵ\mathrm{\Delta }=ϵ\mathrm{\Delta }(1)+ϵ^2\mathrm{\Delta }(2)+ϵ^3\mathrm{\Delta }(3)+\mathrm{}`$ (5.39) $`\mathrm{𝖺𝗇𝖽}`$ (5.40) $`a_n=ϵa_n(1)+ϵ^2a_n(2)+ϵ^3a_n(3)+\mathrm{};`$ (5.41) it follows then $`\mathrm{\Delta }(1)=\mathrm{\Gamma }_{1p},a_n(1)=\mathrm{\Gamma }_{np},`$ (5.42) $`\mathrm{\Delta }(2)={\displaystyle \underset{m=1}{\overset{p}{}}}\mathrm{\Gamma }_{mp}(\mathrm{\Gamma }_{1m+p}+\mathrm{\Gamma }_{1m}\mathrm{\Gamma }_{1p})`$ (5.43) $`a_n(2)={\displaystyle \underset{m=1}{\overset{p}{}}}\mathrm{\Gamma }_{mp}(\mathrm{\Gamma }_{nm+p}\mathrm{\Gamma }_{nm}\mathrm{\Gamma }_{1p}),`$ (5.44) etc.. (Because of (5.37), $`a_n(2)=0`$ for $`n>2p`$.) As a special case, for $`ϵU=ϵx^4`$ (5.45) and therefore $`p=2`$, we find $`ϵ\mathrm{\Delta }(1)={\displaystyle \frac{3ϵ}{4g^2}},ϵ^2\mathrm{\Delta }(2)={\displaystyle \frac{21ϵ^2}{8g^5}},`$ (5.46) the same as the usual perturbation results. Example 2 Using the same one-dimensional harmonic oscillator Hamiltonian (5.20) as the unperturbed $`H`$, we consider now an odd power perturbation potential $`ϵU=ϵx^{2p+1}.`$ (5.47) The analysis given below is parallel to that in the first example, but with some changes, as will be indicated. For $`n0`$, $`Cx^{2n+1}={\displaystyle _0^S}{\displaystyle \frac{dS}{g2S}}x^{2n+1}={\displaystyle \frac{1}{g(2n+1)}}x^{2n+1}.`$ (5.48) It is straightforward to establish the following: Lemma 2 $`(1+CT)^1Cx^{2n+1}={\displaystyle \underset{m=0}{\overset{n}{}}}\gamma _{mn}x^{2m+1}`$ (5.49) where $`\gamma _{nn}={\displaystyle \frac{1}{g(2n+1)}},\gamma _{n1n}={\displaystyle \frac{n}{g^2(2n1)}},`$ (5.50) $`\gamma _{mn}={\displaystyle \frac{n(n1)\mathrm{}(m+1)}{g^{nm+1}(2m+1)}}\mathrm{𝖿𝗈𝗋}0m<n`$ (5.51) and in particular $$\gamma _{0n}=\frac{n!}{g^{n+1}}.$$ Instead of (5.34), we now set $`e^\tau =1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}b_nx^n.`$ (5.52) From (4.87), $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}b_nx^n=(1+CT)^1Cϵ(x^{2p+1}+\mathrm{\Delta })(1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}b_nx^n).`$ (5.53) It is convenient to extend the definitions of $`\gamma _{mn}`$ and $`\mathrm{\Gamma }_{mn}`$ (as in (5.26)), by defining $`\gamma _{mn}=\mathrm{\Gamma }_{mn}=0\mathrm{𝖿𝗈𝗋}m>n`$ (5.54) and $`\mathrm{\Gamma }_{0n}=0.`$ (5.55) Using the two Lemmas, we derive $`\mathrm{\Delta }={\displaystyle \underset{l}{}}(b_{2l+1}\mathrm{\Gamma }_{1l+p+1}\mathrm{\Delta }b_{2l}\mathrm{\Gamma }_{1l}),`$ (5.56) $`b_{2n}=ϵ{\displaystyle \underset{l}{}}(b_{2l+1}\mathrm{\Gamma }_{nl+p+1}\mathrm{\Delta }b_{2l}\mathrm{\Gamma }_{nl})`$ (5.57) and $$b_{2n+1}=ϵ\gamma _{np}ϵ\underset{l}{}(b_{2l}\gamma _{nl+p}\mathrm{\Delta }b_{2l+1}\gamma _{nl}).$$ Here and in the following, whenever the range of the summation index is not exhibited, it reads $`l(\mathrm{𝗈𝗋}m,n)=0,1,2,\mathrm{},\mathrm{}.`$ (5.58) By comparing the first two equations in (5.56), we have $`ϵ\mathrm{\Delta }=b_2.`$ (5.59) It can be readily verified that the following power series expansions hold: $`ϵ\mathrm{\Delta }=ϵ^2\mathrm{\Delta }(2)+ϵ^4\mathrm{\Delta }(4)+ϵ^6\mathrm{\Delta }(6)+\mathrm{}`$ (5.60) $`b_{2n}=ϵ^2b_{2n}(2)+ϵ^4b_{2n}(4)+ϵ^6b_{2n}(6)+\mathrm{},`$ (5.61) $`b_{2n+1}=ϵb_{2n+1}(1)+ϵ^3b_{2n+1}(3)+ϵ^5b_{2n+1}(5)+\mathrm{},`$ (5.62) in which $`b_{2n+1}(1)=\gamma _{np},\mathrm{\Delta }(2)={\displaystyle \underset{l}{}}\gamma _{lp}\mathrm{\Gamma }_{1l+p+1},`$ (5.63) $`b_{2n}(2)={\displaystyle \underset{l}{}}\gamma _{lp}\mathrm{\Gamma }_{nl+p+1},`$ (5.64) $`b_{2n+1}(3)={\displaystyle \underset{l,m}{}}\gamma _{lp}\mathrm{\Gamma }_{ml+p+1}\gamma _{nm+p}\mathrm{\Delta }(2){\displaystyle \underset{m}{}}\gamma _{nm}\gamma _{mp},`$ (5.65) $`\mathrm{\Delta }(4)=\mathrm{\Delta }(2){\displaystyle \underset{l,m}{}}(\gamma _{lp}\gamma _{mp}\mathrm{\Gamma }_{1l+p+1}+\mathrm{\Gamma }_{lm+p+1}\gamma _{mp}\mathrm{\Gamma }_{1l})`$ (5.66) $`{\displaystyle \underset{l,m,n}{}}\gamma _{lp}\mathrm{\Gamma }_{ml+p+1}\gamma _{nm+p}\gamma _{1n+p+1},`$ (5.67) etc. In the special case of $`p=0`$, we have $`ϵU=ϵx,`$ (5.68) the above expressions lead to $`ϵ^2\mathrm{\Delta }(2)={\displaystyle \frac{ϵ^2}{2g^2}}`$ (5.69) and all other $`\mathrm{\Delta }(2n)=0`$. This confirms the exact result, since $`V+ϵU={\displaystyle \frac{g^2}{2}}x^2+ϵx={\displaystyle \frac{g^2}{2}}(x+{\displaystyle \frac{ϵ}{g^2}})^2{\displaystyle \frac{ϵ^2}{2g^2}}.`$ (5.70) The lowest eigenvalue of $`V`$ is $`\frac{g}{2}`$, and that of $`V+ϵU`$ is $`\frac{g}{2}\frac{ϵ^2}{2g^2}`$. Likewise, (5.60)-(5.63) yield $`b_1=ϵb_1(1)=ϵ\gamma _{00}={\displaystyle \frac{ϵ}{g}},`$ (5.71) $`b_2=ϵ^2b_2(2)=ϵ^2\gamma _{00}\mathrm{\Gamma }_{11}={\displaystyle \frac{ϵ^2}{2g^2}},`$ (5.72) $`b_3=ϵ^3b_3(3)=ϵ^3\gamma _{00}\gamma _{11}\mathrm{\Gamma }_{11}={\displaystyle \frac{ϵ^3}{6g^3}},`$ (5.73) etc., so that $`e^\tau =1+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}b_nx^n=e^{ϵ/g}`$ (5.74) which leads to the expected expression $`e^{g𝐒\tau }e^{\frac{g}{2}(x+\frac{ϵ}{g^2})^2}.`$ (5.75) So far we have only used the Green’s function $`G`$ in terms of the single integral form $`C`$. Next, we will repeat the same examples, but using the double integral form $`D`$. Example 1 (repeat) We return to the same simple harmonic oscillator problem (5.20) - (5.21), with $$H=\frac{1}{2}\frac{d^2}{dx^2}+\frac{g^2}{2}x^2$$ and the perturbation $`ϵU=ϵx^{2p}.`$ (5.76) In order to derive the perturbation series (5.34) - (5.42) for $`e^\tau `$ and $`\mathrm{\Delta }`$ using $`D`$, we have to deal with integrals of the type $`\overline{D}x^{2n}=2{\displaystyle \underset{0}{\overset{x}{}}}e^{gy^2}𝑑y{\displaystyle \underset{\mathrm{}}{\overset{y}{}}}e^{gz^2}z^{2n}𝑑z.`$ (5.77) Introduce $`\xi =\sqrt{g}x,\eta =\sqrt{g}y,\zeta =\sqrt{g}z.`$ (5.78) (5.77) becomes $`\overline{D}x^{2n}={\displaystyle \frac{2}{g^{n+1}}}{\displaystyle \underset{0}{\overset{\xi }{}}}e^{\eta ^2}𝑑\eta {\displaystyle \underset{\mathrm{}}{\overset{\eta }{}}}e^{\zeta ^2}\zeta ^{2n}𝑑\zeta .`$ (5.79) Let $`H_l(\zeta )`$ be the usual Hermite polynomials, obtained by the generating function $`e^{t^2+2t\zeta }={\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{t^l}{l!}}H_l(\zeta ).`$ (5.80) The integrand in (5.79) can be expressed in terms of $`H_l(\zeta )`$ through $`\zeta ^{2n}={\displaystyle \underset{m=0}{\overset{n}{}}}{\displaystyle \frac{(2n)!}{2^{2n}m!(2n2m)!}}H_{2n2m}(\zeta ).`$ (5.81) For $`l1`$, $`e^{\xi ^2}H_l(\xi )={\displaystyle \frac{d}{d\xi }}[e^{\xi ^2}H_{l1}(\xi )];`$ (5.82) therefore $`{\displaystyle \underset{0}{\overset{\xi }{}}}e^{\eta ^2}𝑑\eta {\displaystyle \underset{\mathrm{}}{\overset{\eta }{}}}e^{\zeta ^2}H_l(\zeta )𝑑\zeta ={\displaystyle \underset{0}{\overset{\xi }{}}}H_{l1}(\eta )𝑑\eta `$ (5.83) $`={\displaystyle \frac{1}{2l}}[H_l(\xi )H_l(0)]`$ (5.84) and $`\overline{D}H_l(\sqrt{g}x)={\displaystyle \frac{2}{g}}{\displaystyle \underset{0}{\overset{\xi }{}}}e^{\eta ^2}𝑑\eta {\displaystyle \underset{\mathrm{}}{\overset{\eta }{}}}e^{\zeta ^2}H_l(\zeta )𝑑\zeta `$ (5.85) $`={\displaystyle \frac{1}{lg}}[H_l(\xi )H_l(0)].`$ (5.86) Since $`H_0(\zeta )=1`$, by separating out in (5.81) the term $`m=n`$ in the sum, we have $`\zeta ^{2n}=g^n\mathrm{\Gamma }_{1n}+{\displaystyle \underset{m=0}{\overset{n1}{}}}{\displaystyle \frac{(2n)!}{2^{2n}m!(2n2m)!}}H_{2n2m}(\zeta ),`$ (5.87) where $`\mathrm{\Gamma }_{1n}={\displaystyle \frac{(2n1)!!}{(2g)^n}},`$ (5.88) the same as (5.26). From (5.85) and (5.87), it follows that $`\overline{D}(x^{2n}\mathrm{\Gamma }_{1n})={\displaystyle \frac{(2n)!}{(4g)^n}}{\displaystyle \frac{1}{g}}{\displaystyle \underset{l=1}{\overset{n}{}}}{\displaystyle \frac{1}{(nl)!(2l)!2l}}[H_{2l}(\sqrt{g}x)H_{2l}(0)].`$ (5.89) The identity between the above expression and $`(1+CT)^1C(x^{2n}\mathrm{\Gamma }_{1n})`$ given by (5.25) can be established by first applying $`T`$ onto (5.85). Through the well known formula $`{\displaystyle \frac{d^2}{d\xi ^2}}H_l(\xi )=2\xi {\displaystyle \frac{d}{d\xi }}H_l(\xi )2lH_l(\xi ),`$ (5.90) we find $`T\overline{D}H_l(\sqrt{g}x)={\displaystyle \frac{1}{l}}[x{\displaystyle \frac{d}{dx}}H_l(\sqrt{g}x)+lH_l(\sqrt{g}x)].`$ (5.91) Next, apply $`C`$ onto (5.91): $`CT\overline{D}H_l(\sqrt{g}x)`$ $`=`$ $`{\displaystyle \frac{1}{g}}{\displaystyle \underset{0}{\overset{x}{}}}{\displaystyle \frac{dx}{x}}T\overline{D}H_l(\sqrt{g}x)`$ (5.92) $`=`$ $`{\displaystyle \frac{1}{lg}}[H_l(\sqrt{g}x)H_l(0)]+CH_l(\sqrt{g}x).`$ (5.93) Combining (5.85) and (5.92), we obtain $`(1+CT)\overline{D}H_l(\sqrt{g}x)=CH_l(\sqrt{g}x).`$ (5.94) Hence, the multiplication of $`(1+CT)`$ and (5.89) leads to $`(1+CT)\overline{D}(x^{2n}\mathrm{\Gamma }_{1n})=C(x^{2n}\mathrm{\Gamma }_{1n}),`$ (5.95) which gives the equality between (5.89) and (5.25). Thus, both forms of $`G`$, (5.1) and (5.2), yield the same perturbation series expansion (5.38) - (5.42). It is instructive to note that in (5.25), $$(1+CT)^1C[x^{2n}\mathrm{\Gamma }_{1n}]=\underset{m=1}{\overset{n}{}}\mathrm{\Gamma }_{mn}x^{2m},$$ each term, $`(1+CT)^1Cx^{2n}`$ or $`(1+CT)^1C\mathrm{\Gamma }_{1n}`$, is by itself $`\mathrm{}`$, which necessitates the combination that appears on the left-hand side. Using the single-integral form $`C`$ for $`G`$, the energy shift $`\mathrm{\Delta }`$ is required to make the integral finite at $`x=0`$. On the other hand, $`\overline{D}x^{2n}`$, or $$\overline{D}\mathrm{\Gamma }_{1n}=2\underset{0}{\overset{x}{}}e^{gy^2}𝑑y\underset{\mathrm{}}{\overset{y}{}}e^{gz^2}\mathrm{\Gamma }_{1n}𝑑z,$$ is well-defined at any finite $`x`$. As $`x\mathrm{}`$, either $`\overline{D}x^{2n}`$ or $`\overline{D}\mathrm{\Gamma }_{1n}2[\mathrm{\Gamma }_{1n}{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}e^{gz^2}𝑑z]xe^{gx^2}`$ (5.96) approaches $`\mathrm{}`$, which, when multiplied by $`e^{\frac{1}{2}gx^2}`$, would lead to a perturbed wave function $`e^{g𝐒\tau }`$ divergent at $`x=\mathrm{}`$. As shown in (5.18), this requires the energy shift $`\mathrm{\Delta }`$ to be determined by (5.19); in this example, it is equivalent to have the subtraction $`\overline{D}\mathrm{\Gamma }_{1n}`$ in (5.89). Example 2 (repeat). To complete the one-dimensional analysis, we continue with the same simple harmonic oscillator $`H`$ as the unperturbed Hamiltonian, but with the perturbation given by $`ϵU=ϵx^{2p+1}.`$ (5.97) Since by repeated partial integrations, $`{\displaystyle \underset{\mathrm{}}{\overset{y}{}}}e^{gz^2}z^{2n+1}𝑑z`$ $`=`$ $`{\displaystyle \frac{1}{2g}}[y^{2n}+{\displaystyle \frac{n}{g}}y^{2(n1)}+`$ (5.99) $`{\displaystyle \frac{n(n1)}{g^2}}y^{2(n2)}+\mathrm{}+{\displaystyle \frac{n!}{g^n}}]e^{gy^2},`$ we find $`\overline{D}x^{2n+1}`$ $`=`$ $`2{\displaystyle \underset{0}{\overset{x}{}}}e^{gy^2}𝑑y{\displaystyle \underset{\mathrm{}}{\overset{y}{}}}e^{gz^2}z^{2n+1}𝑑z`$ (5.100) $`=`$ $`{\displaystyle \underset{m=0}{\overset{n}{}}}\gamma _{mn}x^{2m+1},`$ (5.101) the same as (5.49) for $`(1+CT)^1Cx^{2n+1}`$, with $`\gamma _{mn}`$ given by (5.50). It is straight forward to see that using either $`\overline{D}`$, or $`(1+CT)^1C`$, we can arrive at the same perturbation series (5.60) and (5.63). ## 6. Excited States In this section, we apply the new approach to excited states. Let $`e^{g𝐒}`$ be the ground-state of (1.1) $$He^{g𝐒}=Ee^{g𝐒}$$ where $$H=\frac{1}{2}^2+g^2v$$ and, as in (1.4)–(1.5), $`v(𝐪)0.`$ (6.1) In this section, we assume $`e^{g𝐒}`$ to be already known; our purpose is to derive the excited states $`\mathrm{\Psi }_{ex}`$, which satisfies $`H\mathrm{\Psi }_{ex}=(E+)\mathrm{\Psi }_{ex}`$ (6.2) Write $`\mathrm{\Psi }_{ex}=\chi (𝐪)e^{g𝐒}.`$ (6.3) Since $$^2\mathrm{\Psi }_{ex}=(^2\chi 2g\chi 𝐒)e^{g𝐒}+\chi (^2e^{g𝐒}),$$ (6.1)-(6.3) lead to $`g𝐒\chi {\displaystyle \frac{1}{2}}^2\chi =\chi .`$ (6.4) The expansion (1.7), $`g𝐒=g𝐒_0+𝐒_1+g^1𝐒_2+g^2𝐒_3+\mathrm{},`$ (6.5) will now be accompanied by similar expansions for $`\chi `$ and $``$: $`\chi =\chi _0+g^1\chi _1+g^2\chi _2+\mathrm{}`$ (6.6) and $`=g_0+_1+g^1_2+\mathrm{}.`$ (6.7) Substituting (6.5) - (6.7) into (6.4) and equating the coefficients of $`g^n`$ on both sides, we obtain the following first order partial differential equations for $`\chi _0,\chi _1,\chi _2,\mathrm{}`$: $`(𝐒_0_0)\chi _0`$ $`=`$ $`0,`$ (6.8) $`(𝐒_0_0)\chi _1`$ $`=`$ $`(𝐒_1+{\displaystyle \frac{1}{2}}^2+_1)\chi _0,`$ (6.9) $`(𝐒_0_0)\chi _2`$ $`=`$ $`(𝐒_1+{\displaystyle \frac{1}{2}}^2+_1)\chi _1+(𝐒_2+_2)\chi _0,`$ (6.10) etc. To see how these equations can be solved we first give a simple example and then address the general solution. Example Take the example $`v(𝐪)={\displaystyle \frac{1}{2}}(\nu _1^2q_1^2+\nu _2^2q_2^2+\mathrm{}+\nu _N^2q_N^2),`$ (6.11) and correspondingly, $`𝐒_0(𝐪)`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\nu _1q_1^2+\nu _2q_2^2+\mathrm{}+\nu _Nq_N^2)`$ (6.12) $`𝐒_1(𝐪)`$ $`=`$ $`𝐒_2(𝐪)=\mathrm{}=0,`$ (6.13) $`E={\displaystyle \frac{g}{2}}(\nu _1+\nu _2+\mathrm{}+\nu _N).`$ (6.14) Introduce, as in (1.26)-(1.28), the variables $$(𝐒_0,\alpha )$$ with $`\alpha =(\alpha _1,\alpha _2,\mathrm{},\alpha _{N1})`$ and $`𝐒_0\alpha _j=0.`$ (6.15) In terms of these new coordinates $`𝐒_0`$ and $`\alpha `$, (6.8) becomes $`(𝐒_0)^2({\displaystyle \frac{\mathrm{ln}\chi _0}{𝐒_0}})_\alpha =_0,`$ (6.16) which will be integrated along the constant-$`\alpha `$ trajectory. The constant-$`\alpha `$ trajectory is determined by the solution of the Hamilton-Jacobi equation (1.11); it describes a classical trajectory with $`v`$ as the potential and a positive infinitesimal energy, starting from $`𝐪=\mathrm{𝟎}`$ at time near $`\mathrm{}`$, and $`𝐪=(q_1,q_2,\mathrm{},q_N)=(𝐒_0,\alpha )`$ at time $`t`$. As $`t`$ increases to $`t+dt`$, the end of the trajectory moves from $`𝐪=(𝐒_0,\alpha )`$ to $`(q_1+dq_1,q_2+dq_2,\mathrm{},q_N+dq_N)=(𝐒_0+d𝐒_0,\alpha )`$ (6.17) keeping $`\alpha `$ constant. In accordance with (2.20) - (2.22), we have $`dt={\displaystyle \frac{d𝐒_0}{(𝐒_0)^2}}={\displaystyle \frac{dq_1}{\nu _1q_1}}={\displaystyle \frac{dq_2}{\nu _2q_2}}=\mathrm{}={\displaystyle \frac{dq_N}{\nu _Nq_N}}.`$ (6.18) When $`𝐒_00`$, each $`q_i`$ must also $`0`$ along the trajectory. Since $`\chi _0`$ is a single-valued function of $`𝐪`$, we can classify $`\chi _0`$ by its power dependence on $`q_i`$, as $`q_i0`$. Take this to be $`\chi _0q_1^{n_1}q_2^{n_2}\mathrm{}q_N^{n_N}.`$ (6.19) Since $`({\displaystyle \frac{\mathrm{ln}\chi _0}{𝐒_0}})_\alpha ={\displaystyle \underset{i=1}{\overset{N}{}}}({\displaystyle \frac{\mathrm{ln}q_i}{𝐒_0}})_\alpha {\displaystyle \frac{\mathrm{ln}\chi _0}{\mathrm{ln}q_i}}.`$ (6.20) From (6.17), we see that along the classical (constant $`\alpha `$ ) trajectory $`({\displaystyle \frac{\mathrm{ln}q_i}{𝐒_0}})_\alpha ={\displaystyle \frac{\nu _i}{(𝐒_0)^2}},`$ (6.21) and consequently $`({\displaystyle \frac{\mathrm{ln}\chi _0}{𝐒_0}})_\alpha ={\displaystyle \frac{1}{(𝐒_0)^2}}{\displaystyle \underset{i=1}{\overset{N}{}}}\nu _i{\displaystyle \frac{\mathrm{ln}\chi _0}{\mathrm{ln}q_i}}.`$ (6.22) As $`𝐒_00`$, under the assumption (6.18), we have $`(𝐒_0)^2({\displaystyle \frac{\mathrm{ln}\chi _0}{𝐒_0}})_\alpha {\displaystyle \underset{i=1}{\overset{N}{}}}n_i\nu _i;`$ (6.23) therefore (6.15) yields $`_0={\displaystyle \underset{i=1}{\overset{N}{}}}n_i\nu _i`$ (6.24) and $`(𝐒_0)^2({\displaystyle \frac{\mathrm{ln}\chi _0}{𝐒_0}})_\alpha ={\displaystyle \underset{i=1}{\overset{N}{}}}n_i\nu _i`$ (6.25) along the entire trajectory for all $`𝐒_0>0`$. (It is reassuring that (6.23) is independent of $`\alpha `$.) It follows then $`\chi _0=q_1^{n_1}q_2^{n_2}\mathrm{}q_N^{n_N}`$ (6.26) at all $`𝐪`$. In this example, because $`𝐒_1=0`$, (6.9) becomes $`(𝐒_0)^2({\displaystyle \frac{\chi _1}{𝐒_0}})_\alpha _0\chi _1=_1\chi _0+{\displaystyle \frac{1}{2}}^2\chi _0,`$ (6.27) which, on account of (6.15) and (6.25), can also be written as $`\chi _0(𝐒_0)^2[{\displaystyle \frac{}{𝐒_0}}({\displaystyle \frac{\chi _1}{\chi _0}})]_\alpha =_1\chi _0+{\displaystyle \frac{1}{2}}[{\displaystyle \frac{n_1(n_11)}{q_1^2}}+{\displaystyle \frac{n_2(n_21)}{q_2^2}}+\mathrm{}+{\displaystyle \frac{n_N(n_N1)}{q_N^2}}]\chi _0.`$ (6.28) Thus, keeping $`\alpha `$ fixed, $`\chi _1=\chi _0{\displaystyle \frac{d𝐒_0}{(𝐒_0)^2}\{_1+\frac{1}{2}[\frac{n_1(n_11)}{q_1^2}+\frac{n_2(n_21)}{q_2^2}+\mathrm{}+\frac{n_N(n_N1)}{q_N^2}]\}},`$ (6.29) where the integration is along the (classical) trajectory normal to the $`𝐒_0=`$ constant surfaces. We leave (6.27) in the indefinite integral form, since the integration constant can be eliminated by the transformation $`\chi _1\chi _1+\mathrm{𝖼𝗈𝗇𝗌𝗍𝖺𝗇𝗍}\chi _0`$ (6.30) which affects only the overall normalization factor of $`\chi `$. (This convention will be adopted below for the general case as well.) For any $`0`$, and any partition $`_1=m_1\nu _1+m_2\nu _2+\mathrm{}+m_N\nu _N`$ (6.31) (where $`m_i`$ can be an arbitrary number) gives $`{\displaystyle \frac{d𝐒_0}{(𝐒_0)^2}_1}=m_1\mathrm{ln}q_1+m_2\mathrm{ln}q_2+\mathrm{}+m_N\mathrm{ln}q_N.`$ (6.32) This leads to an inadmissible wave function, and consequently $`_1=0.`$ (6.33) Hence, by using (6.17), (6.27) and (6.29), we find for this example $`\chi _1={\displaystyle \frac{\chi _0}{4}}[{\displaystyle \frac{n_1(n_11)}{\nu _1q_1^2}}+{\displaystyle \frac{n_2(n_21)}{\nu _2q_2^2}}+\mathrm{}+{\displaystyle \frac{n_N(n_N1)}{\nu _Nq_N^2}}].`$ (6.34) Recalling that the Hermite polynomial $`H_n(x)`$ is $`H_n(z)=(2z)^n[1{\displaystyle \frac{n(n1)}{4z^2}}+\mathrm{}],`$ (6.35) the sum $`\chi =\chi _0+g^1\chi _1+\mathrm{}`$ (6.36) can be shown to be, apart from an overall normalization factor, $`H_{n_1}(\sqrt{g\nu _1}q_1)H_{n_2}(\sqrt{g\nu _2}q_2)\mathrm{}H_{n_N}(\sqrt{g\nu _N}q_N).`$ (6.37) Returning to the general case, because $`𝐪=\mathrm{𝟎}`$ is the minimum of $`v(𝐪)`$, as $`𝐪\mathrm{𝟎}`$, $`v(𝐪)`$ depends quadratically on $`𝐪`$. Write $`v(𝐪){\displaystyle \frac{1}{2}}(\nu _1^2q_1^2+\nu _2^2q_2^2+\mathrm{}+\nu _N^2q_N^2);`$ (6.38) correspondingly, $`𝐒_0(𝐪){\displaystyle \frac{1}{2}}(\nu _1q_1^2+\nu _2q_2^2+\mathrm{}+\nu _Nq_N^2).`$ (6.39) Again, we shall classify $`\chi _0(𝐪)`$ according to its behavior as $`𝐪\mathrm{𝟎}`$, by using (6.18). As in this example, this leads to $`_0=n_1\nu _1+n_2\nu _2+\mathrm{}+n_N\nu _N,`$ (6.40) where, as before, $`n_1,n_2,\mathrm{},n_N`$ are positive integers. From (6.8), it follows that $`\mathrm{ln}\chi _0=_0{\displaystyle \frac{d𝐒_0}{(𝐒_0)^2}}`$ (6.41) with the integration taken along the trajectory of constant $`\alpha `$, and the integration constant is determined by the normalization condition (6.18), as $`𝐒_00`$. It is convenient to characterize the limit $`𝐒_00`$, by introducing along the trajectory an overall scale factor $`\lambda `$ (of the dimension $`[q_i]`$). Because of (6.17), for sufficiently small $`𝐒_0`$, we set $`𝐒_0\lambda ^2,q_i\lambda ,`$ (6.42) and therefore, on account of (6.17), $`{\displaystyle \frac{d𝐒_0}{(𝐒_0)^2}}={\displaystyle \frac{1}{\nu _1}}d\mathrm{ln}q_1={\displaystyle \frac{1}{\nu _2}}d\mathrm{ln}q_2=\mathrm{}={\displaystyle \frac{1}{\nu _N}}d\mathrm{ln}q_Nd\mathrm{ln}\lambda .`$ (6.43) We now turn to (6.9); as in (6.26)-(6.27), it can be written as $`\chi _0(𝐒_0)^2[{\displaystyle \frac{}{𝐒_0}}({\displaystyle \frac{\chi _1}{\chi _0}})]_\alpha =({\displaystyle \frac{1}{2}}^2𝐒_1)\chi _0+_1\chi _0,`$ (6.44) and therefore $`\chi _1=\chi _0{\displaystyle \frac{d𝐒_0}{(𝐒_0)^2}[\frac{1}{\chi _0}(\frac{1}{2}^2𝐒_1)\chi _0+_1]}.`$ (6.45) Near $`𝐒_0=0`$, we may use (6.36) to expand the first term inside the square bracket as a power series in $`\lambda `$: $`{\displaystyle \frac{1}{\chi _0}}({\displaystyle \frac{1}{2}}^2𝐒_1)\chi _0`$ $`=`$ $`b_2\lambda ^2+b_1\lambda ^1+b_0+`$ (6.47) $`b_1\lambda +b_2\lambda _2+\mathrm{},`$ where $`b_2,b_1,b_0,b_1,\mathrm{}`$ are constants. Because of (6.37), in order that $`\chi _1`$ be analytic at $`𝐪=\mathrm{𝟎}`$, we must require $`_1=b_0;`$ (6.48) otherwise, $`\chi _1`$ would have a term proportional to $`\chi _0\mathrm{ln}\lambda `$, which is not admissible. The function $`\chi _1`$ is given by the integral (6.38) along the $`\alpha `$-constant trajectory. In a similar way, we can determine $`_2,_3,\mathrm{}`$ and obtain the solutions $`\chi _2,\chi _3,\mathrm{}`$, in terms of quadrature along the classical (constant $`\alpha `$) trajectory. ## 7. Perturbation Around An Attractive Coulomb Potential Let $`H_c`$ be the Hamiltonian for a Coulomb potential: $`H_c={\displaystyle \frac{1}{2}}^2{\displaystyle \frac{g^2}{r}},`$ (7.1) where $`g^2=Ze^2`$ (7.2) with $`^2`$ denoting the three-dimensional Laplacian, $`r`$ the radius, $`Ze`$ the nuclear charge and $`e`$ the electronic charge. Consider a problem in which there is an additional perturbation $`ϵU`$; the corresponding Hamiltonian is $`H=H_c+ϵU,`$ (7.3) where $`U`$ is not singular at the origin. Let $`\psi _c`$ and $`\psi `$ be the ground-states of $`H_c`$ and $`H`$; i.e., $`H_c\psi _c=E_c\psi _c`$ (7.4) and $`H\psi =E\psi .`$ (7.5) 7.1 Isotropic Case We first discuss the case that $`U(r)`$ depends only on $`r`$. As we shall see, with modification our method can be adapted to derive $`\psi `$ by quadratures along the radial trajectory. The solution to the Coulomb problem is well known: $`\psi _c=e^{g^2r}\mathrm{𝖺𝗇𝖽}E_c={\displaystyle \frac{1}{2}}g^4.`$ (7.6) This suggests that instead of (1.7)-(1.8), a different $`g`$-power expansion is needed, one that should conform to the form (7.6) of the Coulomb wave function. Write $`\psi =e^S.`$ (7.7) Expand $`S=g^2S_0+S_1+g^2S_2+g^4S_3+\mathrm{}+g^{(2n2)}S_n+\mathrm{}`$ (7.8) and $`E=g^4E_0+g^2E_1+E_2+g^2E_3+\mathrm{}+g^{(2n4)}E_n+\mathrm{}.`$ (7.9) Since $$^2\psi =[(S)^2^2S]\psi ,$$ (7.5) becomes $`{\displaystyle \frac{1}{2}}(S)^2+{\displaystyle \frac{1}{2}}^2S{\displaystyle \frac{g^2}{r}}+ϵU=E.`$ (7.10) Substituting (7.8)-(7.9) into (7.10) and equating the coefficients of $`g^4,g^2,g^0,\mathrm{},g^{2m},\mathrm{}`$ on both sides, we obtain $`(S_0)^2`$ $`=`$ $`2E_0`$ (7.11) $`S_0S_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}^2S_0{\displaystyle \frac{1}{r}}E_1`$ (7.12) $`S_0S_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}(S_1)^2+{\displaystyle \frac{1}{2}}^2S_1+ϵUE_2,`$ (7.13) $`S_0S_3`$ $`=`$ $`S_1S_2+{\displaystyle \frac{1}{2}}^2S_2E_3,`$ (7.14) etc.; for $`n>1`$, $`S_0S_{2n}={\displaystyle \underset{m=1}{\overset{n1}{}}}S_mS_{2nm}{\displaystyle \frac{1}{2}}(S_n)^2+{\displaystyle \frac{1}{2}}^2S_{2n1}E_{2n}`$ (7.15) and for $`n1`$ $`S_0S_{2n+1}={\displaystyle \underset{m=1}{\overset{n}{}}}S_mS_{2n+1m}+{\displaystyle \frac{1}{2}}^2S_{2n}E_{2n+1}.`$ (7.16) For the ground-state, these $`S_m(r)`$ are all radial functions. Since $`U(r)`$ is regular at $`r=0`$, we can set $`U(0)=0.`$ (7.17) We adopt the same normalization condition (1.36), with $`\psi (0)=1\mathrm{𝖺𝗇𝖽}S(0)=0.`$ (7.18) From (7.11), it follows that $`{\displaystyle \frac{S_0}{r}}=\sqrt{2E_0}`$ (7.19) $`\mathrm{𝖺𝗇𝖽}`$ (7.20) $`S_0=\sqrt{2E_0}r.`$ (7.21) Substituting these expressions into (7.12), we find $`\sqrt{2E_0}{\displaystyle \frac{S_1}{r}}={\displaystyle \frac{1}{r}}(\sqrt{2E_0}1)E_1.`$ (7.22) Because $`S_1`$ should be regular at $`r=0`$, $`\sqrt{2E_0}1=0`$ and therefore $`E_0={\displaystyle \frac{1}{2}}\mathrm{𝖺𝗇𝖽}S_0=r`$ (7.23) confirming (7.6); (7.22) then leads to $`{\displaystyle \frac{S_1}{r}}=E_1`$ (7.24) $`\mathrm{𝖺𝗇𝖽}`$ (7.25) $`S_1=E_1r.`$ (7.26) Since $`^2r={\displaystyle \frac{2}{r}},`$ (7.27) we see that (7.13) becomes $`{\displaystyle \frac{S_2}{r}}={\displaystyle \frac{1}{2}}E_1^2{\displaystyle \frac{E_1}{r}}+ϵU(r)E_2.`$ (7.28) In order that $`S_2(r)`$ be regular at $`r=0`$, $`E_1=0,`$ (7.29) and consequently, on account of (7.24), $`S_1(r)=0`$ (7.30) and $`{\displaystyle \frac{S_2}{r}}=ϵU(r)E_2.`$ (7.31) By using (7.27), we see that as $`r0`$ $`^2S_2(^2S_2)_0={\displaystyle \frac{2}{r}}({\displaystyle \frac{S_2}{r}})_0={\displaystyle \frac{2}{r}}[ϵU(0)E_2].`$ (7.32) From (7.14), we see that in order to have $`S_3`$ also regular at $`r=0`$, $`(\frac{S_2}{r})_0`$ must be $`0`$; this yields on account of (7.17), $`E_2`$ $`=`$ $`ϵU(0)=0`$ (7.33) $`{\displaystyle \frac{S_2}{r}}`$ $`=`$ $`ϵU(r)`$ (7.34) and $`{\displaystyle \frac{S_3}{r}}={\displaystyle \frac{ϵ}{2r^2}}{\displaystyle \frac{}{r}}(r^2U)E_3.`$ (7.35) As $`r0`$, (7.14) implies $`{\displaystyle \frac{S_3}{r}}({\displaystyle \frac{S_3}{r}})_0=[{\displaystyle \frac{ϵ}{2r^2}}{\displaystyle \frac{}{r}}(r^2U)]_0E_3`$ (7.36) and, on account of (7.27), $`^2S_3{\displaystyle \frac{2}{r}}({\displaystyle \frac{S_3}{r}})_0.`$ (7.37) Hence by using (7.15) for $`S_4`$ and in order that $`S_4`$ should be regular at $`r=0`$, $`({\displaystyle \frac{S_3}{r}})_0=0,`$ (7.38) and therefore $`E_3={\displaystyle \frac{ϵ}{2}}[{\displaystyle \frac{1}{r^2}}{\displaystyle \frac{}{r}}(r^2U)]_0.`$ (7.39) The same reasoning yields for all $`m1`$, at $`r=0`$, $`({\displaystyle \frac{S_m}{r}})_0=0.`$ (7.40) From (7.39), we see that for $`ϵU(r)`$ $`=`$ $`ϵr^l`$ (7.41) $`E_3`$ $`=`$ $`\{\begin{array}{cc}\frac{3}{2}ϵ\hfill & \mathrm{𝗂𝖿}l=1\hfill \\ 0\hfill & \mathrm{𝗂𝖿}l>1.\hfill \end{array}`$ (7.44) Example To carry the analysis further, let us consider the case $`l=2`$; i.e., $`ϵU(r)=ϵr^2.`$ (7.45) In this example, the above results give $`E_1=E_2=E_3=0`$ (7.46) $`S_1=0,S_2={\displaystyle \frac{1}{3}}ϵr^3`$ (7.47) and $$S_3=ϵr^2.$$ Set $`n=2`$ in (7.15), we have $`{\displaystyle \frac{S_4}{r}}={\displaystyle \frac{1}{2}}ϵ^2r^4+3ϵE_4.`$ (7.48) Because of (7.40), at $`r=0`$, $`\frac{S_4}{r}=0`$; hence, $`E_4=3ϵ`$ (7.49) and $`S_4={\displaystyle \frac{1}{10}}ϵ^2r^5.`$ (7.50) In a similar way, we derive $`E_5=0,`$ $`S_5={\displaystyle \frac{7}{8}}ϵ^2r^4,`$ (7.51) $`E_6=0,`$ $`S_6={\displaystyle \frac{43}{12}}ϵ^2r^3+{\displaystyle \frac{1}{14}}ϵ^3r^7,`$ (7.52) $`E_7=0,`$ $`S_7={\displaystyle \frac{43}{4}}ϵ^2r^2+{\displaystyle \frac{13}{12}}ϵ^3r^6,`$ (7.53) $`E_8={\displaystyle \frac{129}{4}}ϵ^2,`$ $`\mathrm{𝖾𝗍𝖼}.`$ (7.54) Putting together, for $`ϵU=ϵr^2`$, we find $`S`$ $`=`$ $`g^2r+{\displaystyle \frac{ϵ}{3g^2}}r^3+{\displaystyle \frac{ϵ}{g^4}}r^2{\displaystyle \frac{ϵ^2}{10g^6}}r^5`$ (7.57) $`{\displaystyle \frac{7ϵ^2}{8g^8}}r^4+{\displaystyle \frac{1}{g^{10}}}({\displaystyle \frac{43}{12}}ϵ^2r^3+{\displaystyle \frac{1}{14}}ϵ^3r^7)`$ $`+{\displaystyle \frac{1}{g^{12}}}({\displaystyle \frac{43}{4}}ϵ^2r^2+{\displaystyle \frac{13}{12}}ϵ^3r^6)+O({\displaystyle \frac{ϵ^3}{g^{14}}})`$ and $`E={\displaystyle \frac{g^4}{2}}+{\displaystyle \frac{3ϵ}{g^4}}{\displaystyle \frac{129}{4g^{12}}}ϵ^2+O({\displaystyle \frac{ϵ^3}{g^{20}}}).`$ (7.58) In the usual perturbation series, to each order of the perturbation, the derivation of the perturbed Coulomb wave function requires summations over an infinite number of excited bound-states, plus the continuum. Here, to each order $`ϵ^m`$, the perturbed wave function can be obtained in closed form by quadratures along the radial trajectory. The perturbed energy can also be derived in an alternative way by using the integral form (4.92): $`E={\displaystyle \frac{1}{2}}g^4+{\displaystyle \frac{\underset{0}{\overset{\mathrm{}}{}}e^{g^2rS}ϵUr^2𝑑r}{\underset{0}{\overset{\mathrm{}}{}}e^{g^2rS}r^2𝑑r}}.`$ (7.59) To first order in $`ϵ`$, we can approximate the factor $`e^S`$ by $`e^{g^2r}`$ in the integrals. For $`U=r^2`$, $`E`$ $`=`$ $`{\displaystyle \frac{1}{2}}g^4+{\displaystyle \frac{ϵ\underset{0}{\overset{\mathrm{}}{}}e^{2g^2r}r^4𝑑r}{\underset{0}{\overset{\mathrm{}}{}}e^{2g^2r}r^2𝑑r}}+O(ϵ^2)`$ (7.60) $`=`$ $`{\displaystyle \frac{1}{2}}g^4+{\displaystyle \frac{3ϵ}{g^4}}+O(ϵ^2).`$ (7.61) To calculate $`E`$ to the accuracy of $`O(ϵ^2)`$, we need the wave function to $`O(ϵ)`$. From (7.57), $`e^S=e^{g^2r}[1ϵ({\displaystyle \frac{r^3}{3g^2}}+{\displaystyle \frac{r^2}{g^4}})+O(ϵ^2)].`$ (7.62) Substituting this expression into (7.59), we find the second term on its right-hand side to be: $`{\displaystyle \frac{ϵ\underset{0}{\overset{\mathrm{}}{}}e^{g^2r}r^4[1ϵ(\frac{r^3}{3g^2}+\frac{r^2}{g^4})]𝑑r}{\underset{0}{\overset{\mathrm{}}{}}e^{g^2r}r^2[1ϵ(\frac{r^3}{3g^2}+\frac{r^2}{g^4})]𝑑r}},`$ (7.63) confirming (7.58). 7.2 Stark Effect Consider the anisotropic example, in which the perturbation is $`ϵU=ϵr\mathrm{cos}a,`$ (7.64) where $`a`$ is the polar angle; i.e., $`r^2=x^2+y^2+z^2\mathrm{𝖺𝗇𝖽}z=r\mathrm{cos}a.`$ (7.65) Replacing $`ϵU(r)`$ by $`ϵr\mathrm{cos}a`$, we see that (7.1)-(7.33) remain intact with $`E_1=E_2=S_1=0`$; (7.29) becomes $`{\displaystyle \frac{S_2}{r}}=ϵr\mathrm{cos}a,`$ (7.66) and therefore $`S_2={\displaystyle \frac{1}{2}}ϵr^2\mathrm{cos}a,`$ (7.67) and $`^2S_2=2ϵ\mathrm{cos}a.`$ (7.68) Thus, (7.14) gives $`{\displaystyle \frac{S_3}{r}}=ϵ\mathrm{cos}aE_3,`$ (7.69) which yields, on account of (7.49), $`S_3=ϵzE_3r`$ (7.70) and $`^2S_3={\displaystyle \frac{2E_3}{r}}.`$ (7.71) Since $`S_4`$ satisfies (7.15) for $`n=2`$, $`{\displaystyle \frac{S_4}{r}}={\displaystyle \frac{1}{2}}(S_2)^2+{\displaystyle \frac{1}{2}}^2S_3E_4.`$ (7.72) In order that $`S_4`$ not have a $`\mathrm{ln}r`$ singularity, $`E_3=0`$ (7.73) and therefore $`S_3=ϵz.`$ (7.74) In order that $`^2S_4`$ not have a $`1/r`$ singularity, $`E_4=0.`$ (7.75) Thus, $`S_4={\displaystyle \frac{1}{24}}ϵ^2r^3(1+3\mathrm{cos}^2a).`$ (7.76) Likewise, it is straightforward to derive $`E_5`$ $`=`$ $`0,S_5={\displaystyle \frac{7}{16}}ϵ^2r^2(1+\mathrm{cos}^2a),`$ (7.77) $`E_6`$ $`=`$ $`{\displaystyle \frac{9}{4}}ϵ^2,S_6={\displaystyle \frac{1}{16}}ϵ^3r^4\mathrm{cos}a(1+\mathrm{cos}^2a),`$ (7.78) $`E_7`$ $`=`$ $`0,S_7={\displaystyle \frac{13}{48}}ϵ^3r^3\mathrm{cos}a(3+\mathrm{cos}^2a),`$ (7.79) $`E_8`$ $`=`$ $`0,S_8={\displaystyle \frac{53}{16}}ϵ^3r^2\mathrm{cos}a{\displaystyle \frac{1}{128}}ϵ^4r^5(1+10\mathrm{c}\mathrm{o}\mathrm{s}^2a+5\mathrm{c}\mathrm{o}\mathrm{s}^4a),`$ (7.80) $`E_9`$ $`=`$ $`0,S_9={\displaystyle \frac{53}{8}}ϵ^3r\mathrm{cos}a{\displaystyle \frac{99}{512}}ϵ^4r^4(1+6\mathrm{c}\mathrm{o}\mathrm{s}^2a+\mathrm{cos}^4a),`$ (7.81) $`E_{10}`$ $`=`$ $`0,S_{10}={\displaystyle \frac{761}{384}}ϵ^4r^3(1+3\mathrm{c}\mathrm{o}\mathrm{s}^2a)+O(ϵ^5),`$ (7.82) $`E_{11}`$ $`=`$ $`0,S_{11}={\displaystyle \frac{3131}{256}}ϵ^4r^2(1+\mathrm{cos}^2a)+O(ϵ^5),`$ (7.83) $`E_{12}`$ $`=`$ $`{\displaystyle \frac{3555}{64}}ϵ^4,`$ (7.84) etc. Combining these results together, for the potential $`\frac{g^2}{r}+ϵr\mathrm{cos}a`$, we find that, in powers of $`ϵ`$, the wave function $`e^S`$ and the energy are given by $`S`$ $`=`$ $`g^2r+{\displaystyle \frac{ϵr}{g^4}}\mathrm{cos}a(1+{\displaystyle \frac{1}{2}}g^2r){\displaystyle \frac{ϵ^2r^2}{g^8}}[{\displaystyle \frac{7}{16}}(1+\mathrm{cos}^2a)+{\displaystyle \frac{1}{24}}g^2r(1+3\mathrm{c}\mathrm{o}\mathrm{s}^2a)]`$ (7.85) $`+`$ $`{\displaystyle \frac{ϵ^3r}{g^{16}}}\mathrm{cos}a[{\displaystyle \frac{53}{8}}(1+{\displaystyle \frac{1}{2}}g^2r)+{\displaystyle \frac{13}{48}}(g^2r)^2(3+\mathrm{cos}^2a)+{\displaystyle \frac{1}{16}}(g^2r)^3(1+\mathrm{cos}^2a)]`$ (7.86) $`+`$ $`O(ϵ^4)`$ (7.87) and $`E={\displaystyle \frac{1}{2}}g^4{\displaystyle \frac{9}{4}}{\displaystyle \frac{ϵ^2}{g^8}}{\displaystyle \frac{3555}{64}}{\displaystyle \frac{ϵ^4}{g^{20}}}+O(ϵ^6).`$ (7.88) The $`\frac{9}{4}\frac{ϵ^2}{g^8}`$ is known. Our method gives closed expressions for both the wave function and the energy to any finite order of $`ϵ`$. ## Acknowledgment One of us (T. D. Lee) wishes to thank T. K. Lee for discussions, and for a private communication which outlines how our method can be extended to include the hydrogen atom. We also wish to give our thanks to RIKEN, Brookhaven National Laboratory and to the U.S. Department of Energy \[DE-AC02-98CH10886\] for providing the facilities essential for the completion of this work.
warning/0005/cond-mat0005101.html	ar5iv	text	# Wave functions in disordered wires in a weak magnetic field ## I Introduction Localization of all states in one-dimensional chains and quasi-one-dimensional wires with a finite thickness is a very well known phenomenon in the theory of disordered metals. Following the first works where the localization was predicted, a number of analytical methods has been developed to treat both the chains and wires quantitatively. A detailed discussion of the localization in the chains and wires can be found, e.g. in Refs. The localization of electron waves in weakly disordered chains and wires occurs due to destructive interference and can be destroyed at finite temperatures by inelastic scattering. The phenomenon of the electron localization in thick wires is richer than in chains. The localization length $`L_c`$ in chains is of the order of the mean free path $`l`$. Electrons move at distances smaller than the localization length ballistically without being scattered by impurities and get localized at distances exceeding $`L_c`$. In contrast, the localization length $`L_c`$ in wires is larger than the mean free path $`l`$ by a factor proportional to the number of channels of the transversal quantization $`N`$. At distances larger than $`l`$ but smaller than $`L_c`$, the electrons diffuse and become localized only at distances larger than $`L_c`$. One more important difference between chains and wires manifests itself if an external magnetic field $`H`$ is applied. In chains the magnetic field cannot influence the electron motion, whereas the motion in thick wires is quite sensitive to it. It turns out that the magnetic field does not destroy the localization of the wave functions but changes the localization length. Remarkably, the localization length $`L_c`$ in thick wires exactly doubles when applying a sufficiently strong magnetic field. For analytic calculations for thick wires a supermatrix $`\sigma `$-model proved to be a powerful tool that allows one to reduce the calculation of kinetic quantities at arbitrary frequencies $`\omega `$ to solving of a system of differential equations. The localization length $`L_c`$ is obtained from an exponential decay of the density-density correlation function. The most interesting limit $`\omega 0`$ can be investigated explicitly and one obtains relatively simple formulae for different physical quantities proving the localization of all states. Calculation of the transmittance of a finite sample with ideal leads is the basis of many numerical investigations of the localization properties. The equivalence of different methods of calculation of the localization length $`L_c`$ is traditionally guaranteed by the Borland conjecture stating a complete independence of the localization properties of the boundary conditions. The limiting cases of a strong magnetic field and no magnetic field, the unitary and the orthogonal ensembles, respectively, are relatively well studied in quasi 1D, whereas the crossover between these two limits has not been well understood. It was assumed that the localization length changed smoothly between $`L_c`$ and $`2L_c`$ and one had to calculate a curve connecting $`L_c`$ and $`2L_c`$ when increasing the magnetic field from $`0`$ to $`\mathrm{}`$. At the same time, investigation of the crossover can have quite interesting applications now because a systematic experimental study of the localization has been performed recently. In the experiment, a large number of submicrometer-wide wires prepared from Si $`\delta `$-doped GaAs were connected in parallel to insure statistical averaging. This enabled accurate measurements of the conductivity as a function of the temperature $`T`$ and the magnetic field $`H`$ applied perpendicular to the wires. The regimes of the weak and strong localization were observed and the experimental data in the regime of strong localization were used to extract the dependence of the localization length $`L_c`$ on $`T`$ and $`H`$. The localization length $`L_c`$ was shown to double in a strong magnetic field as compared with the limit of zero magnetic field $`H`$, which is in a complete agreement with the theory. As the measurements are done at an arbitrary magnetic field, study of the crossover between the orthogonal and unitary ensembles may become very important. In a recent publication, we discovered that already an arbitrarily weak magnetic field drastically changes the tails of the averaged amplitudes of the wave functions. At sufficiently large distances depending on the magnetic field, the exponential decay with the localization length $`L_c`$ changes to a decay with the length $`2L_c`$. This result is of interest not only from the theoretical point of view but it might also be important for the transport measurements, since the localization length $`L_c`$ is relevant for the hopping conductivity at low temperature. However, a numerical study of the averaged logarithm of the transmittance performed for sufficiently long wires did not manifest any two-scale behavior showed a smooth variation of the localization length between $`L_c`$ and $`2L_c`$. No sharp change in the far tail of the logarithm of the transmittance was observed in a weak magnetic field. As a possible explanation of the contradiction several reasons were suggested. As one of the possibilities, the difference between the logarithm of the averaged wave functions, the quantity we have considered, and the averaged logarithm of the wave functions was mentioned. In this paper, we give details of the calculations presented in our short publication and discuss the question in what situations our effect of the changing of the tails of the wave functions by a weak magnetic field can be identified. To understand better the shape of the wave functions we calculate not only the density-density correlation function but also its moments. The changes in the tails of the wave functions are seen in all moments of the density-density correlation functions and all of them are proportional at large distances $`x`$ to $`\mathrm{exp}(\|x\|/4L_{\mathrm{cu}})`$, where $`L_{\mathrm{cu}}`$ is the localization length for the unitary ensemble. The independence of the exponential decay of the asymptotics of the number of the moment signals that they are formed by strong rare splashes of the wave functions. The situation with the averaged logarithm of the wave function is more delicate. The large splashes contribute also to the averaged logarithm of the wave functions at finite distances. As a result, some broad distribution of logarithms occurs. Unfortunately, we are not able to predict how this distribution function develops in the limit of infinitely long distances. However, our results provide us reasons enough to suggest that the logarithmically normal distribution implied within the Borland conjecture is not universally applicable for length of the sample $`xL_c`$ for the description of the crossover region. Existing numerical studies might be interpreted as the head of some logarithmically normal distribution changes smoothly with $`H`$. Our results do not exclude a smooth behavior of the average logarithm of the wave function. We show, however, that another length scale, $`x_HL_c\mathrm{ln}(1/H)`$ enters the theory of localization, up to which the distribution should differ from the standard logarithmically normal one. Thus, the whole distribution function of the logarithm is of interest. For calculations we use the supersymmetry technique that makes possible consideration of disordered wires with any kind of boundary conditions. In particular, one can study the transmittance of the disordered wires with metallic leads. An external magnetic field is naturally accounted for within supersymmetry formalism allowing study of the crossover between the unitary and orthogonal ensembles. The localization of the wave functions discussed in the present paper is extracted from the moments of the density-density correlation function $$p_{\mathrm{}}^{(n)}\left(xx^{}\right)=\underset{\alpha }{}\left\|\psi _\alpha \left(x\right)\right\|^{2n}\left\|\psi _\alpha \left(x^{}\right)\right\|^{2n}\delta \left(\epsilon \epsilon _\alpha \right)$$ (1) where $`\psi _\alpha \left(x\right)`$ and $`\epsilon _\alpha `$ are the eigenfunction and eigenenergy of a state $`\alpha `$, $`x`$ and $`x^{}`$ are the coordinates along the wire, and the angle brackets stand for the averaging over the disorder. This quantity was previously studied for a closed finite multichannel wire, when the wave functions were taken at the ends of the wire. Generalizing the expressions for arbitrary $`n`$, the complete distribution was also obtained. Naturally, the very presence of hard boundaries or ideal leads inevitably changes localized states near the ends of the disordered sample. These states should differ from those localized in the bulk. Such a difference has been found indeed for the density of states of a single chain. The localization properties extracted from the transmittance of an open system and from the density-density correlator of a closed one, manifest a generic universality: All moments of the both quantities decay with the same rate, $`\mathrm{exp}(\|xx^{}\|/4L_c)`$, and the logarithm of the both quantities is normally distributed. The fact that the decay rate is the same for all moments shows that an important contribution comes from rare strong splashes of the wave functions. Due to these splashes the moments decay slower with the length $`4L_c`$ and not with the localization length $`L_c`$. At the same time, the logarithm of physical quantities like the transmittance or the product of the amplitudes of the wave functions at different points is a self-averaging quantity and approaches the value $`\|xx^{}\|/L_c`$ for $`\|xx^{}\|\mathrm{}`$. This quantity characterizes typical wave functions. The question what happens with the wave functions if a weak magnetic field is switched on has not been discussed. In the subsequent chapters we present details of our analysis. We use an approach based on the Lebedev-Kontorovich transformation to find the exact distribution of any quantity of interest in quasi-1D wires for the orthogonal and unitary ensemble. This allows us to calculate the correlator, Eq. (1), and the Landauer-like conductivity. Exact expressions valid at arbitrary distances are obtained for all moments for the pure orthogonal and unitary ensemble. Using the formulae for the moments we derive the entire distribution function. At a weak magnetic field, we show that a peculiar two-scale behavior characterizes any moment of the wave functions, Eq. (1). The paper is organized as follows: In Sec. II we introduce the quantities to be studied. A transfer-matrix formalism used for the calculations is described in Sec. III. Reduction to an effective Schrödinger equation is carried out in Sec. IV. Then, in Sec. V we calculate the moments and the distribution functions in the limiting cases of the orthogonal and unitary ensembles. In Sec. VI we present results for the wire in a weak magnetic field. Using a standard qualitative picture for localized states we discuss the physics involved in Sec. VII. Our arguments allow us to obtain without any calculation the characteristic length scales of the problem and reach a better understanding of the localization in the magnetic field. We discuss the results and make conclusions in Sec. VIII. ## II Correlation Functions Localization in disordered metals can be well characterized by different correlation and distribution functions. For an infinite sample one can consider the density-density correlation function and its moments. For a finite wire connected to metallic leads the transmittance and its moments are of interest. These quantities can be efficiently studied using the supersymmetry method. Following this approach one should reduce calculation of the correlation functions of interest to computation of correlation functions of a supermatrix $`\sigma `$-model. In the standard formulation the $`\sigma `$-model contains $`8\times 8`$ supermatrices $`Q`$. Fluctuations of the supermatrices $`Q`$ are strongly influenced by a magnetic field and one half of the “degrees of freedom” is suppressed if the magnetic field is sufficiently strong. This corresponds to reducing the size of the supermatrices $`Q`$ to $`4\times 4`$. Keeping this size of the supermatrices is sufficient to calculate such physical quantities as conductivity or density-density correlation function. In order to calculate higher moments of these quantities or distribution functions one has to increase the size of the supermatrices, which would make the theory very complicated. Fortunately, calculation of moments of certain quantities does not demand increasing the size of the supermatrices, which allows to compute them explicitly. In all such cases the free energy functional $`F\left[Q\right]`$ describing fluctuations of the $`8\times 8`$ supermatrices $`Q`$ can be written as $$F[Q]=\frac{\pi \nu }{8}\mathrm{Str}\left[D\left(_𝐫Q(𝐫)\frac{ie}{\mathrm{}c}𝐀[Q(𝐫),\tau _3]\right)^2+2i\omega \mathrm{\Lambda }Q(𝐫)\right]d𝐫.$$ (2) The $`\sigma `$-model, Eq. (2), appears after averaging over disorder and we are to carry out computations with the regular model. In Eq. (2) $`D`$ is the classical diffusion coefficient, $`\nu `$ is the density of states, and $`\omega `$ is the external frequency. The supermatrix $`Q`$ obeys the constraint $`Q^2=1`$, $`𝐀`$ is the vector potential corresponding to an external magnetic field $`H`$, and the standard notations for the supertrace $`\mathrm{Str}`$ and matrices $`\mathrm{\Lambda }`$, $`\tau _3`$ are used. Any quantity of interest in the present formulation is expressed as a functional integral over $`Q(𝐫)`$ with the weight $`\mathrm{exp}(F[Q])`$ and a combination of elements of the supermatrix $`Q`$ in the pre-exponential. We study properties of wave functions calculating their moments written in Eq. (1). As the first step, we express the moments of the wave functions in terms of retarded $`G^R`$ and advanced $`G^A`$ Green functions that can be written as $$G_\epsilon ^{R,A}(x,x^{})=\underset{\alpha }{}\frac{\psi _\alpha \left(x\right)\psi _\alpha ^{}\left(x^{}\right)}{\epsilon \epsilon _\alpha },$$ (3) where $`\psi _a`$ and $`\epsilon _\alpha `$ are eigenfunctions and eigenenergies of the electron states in the disordered system. Using the spectral expansion, Eq. (3), assuming that all states are localized, such that the spectrum is discrete, and considering the most divergent terms at vanishing frequency $`\omega 0`$ one can come to the following relation $$\left[G_\epsilon ^R(x,x^{})G_{\omega +\epsilon }^A(x^{},x)\right]^n=\frac{2\pi i(2n2)!(1)^{n1}}{\left((n1)!\right)^2\omega ^{2n1}}\underset{\alpha }{}\|\psi _\alpha (x)\|^{2n}\|\psi _\alpha (x^{})\|^{2n}\delta (\epsilon \epsilon _\alpha )$$ (4) Although Eq. (4) contains an arbitrary power of the Green functions we can express the LHS in terms of $`8`$-component supervectors. This possibility is due to the fact that the Green functions are taken at two different points only. The next step is the averaging over disorder and the decoupling the effective “interaction” by integration over the supermatrices $`Q`$. A functional integral over $`Q`$ is simplified in the saddle-point approximation and we come after standard manipulations to the following expression $$\left[G_\omega ^{}^R(x,x^{})G_{\omega +\omega ^{}}^A(x^{},x)\right]^n=(\pi \nu )^{2n}[Q_{33}^{12}(x)]^n[Q_{33}^{21}(x^{})]^n_F.$$ (5) In the LHS and RHS of Eq. (5), averaging over the impurity potential and averaging over the free energy, Eq. (2), is implied, respectively. The superscripts and subscripts of the supermatrices $`Q`$ stand for certain matrix elements. In Section V, exact expressions for the moments of the density-density correlation function, Eq. (1), as well as its entire distribution function $$𝒫_\psi (t)=\underset{\alpha }{}\delta \left(t\left\|\psi _\alpha \left(x\right)\right\|^2\left\|\psi _\alpha \left(x^{}\right)\right\|^2\right)\delta \left(\epsilon \epsilon _\alpha \right).$$ (6) will be obtained. Hereinafter we use dimensionless expression for the wave functions $`\psi SL_c\psi `$, where $`S`$ is the cross-section of the wire and $`L_c`$ is the localization length. To calculate the distribution function $`𝒫_\psi (t)`$, Eq. (6), we introduce an auxiliary correlation function $$𝒫_{a,b}(t_1,t_2)=\delta \left(t_1aQ_{33}^{12}(x_1)\right)\delta \left(t_2bQ_{33}^{21}(x_2)\right)_F,$$ (7) where $`a`$ and $`b`$ are some parameters. Once the correlator Eq. (7) is known, any other quantity of interest can be found by integrations of the type $`_{\mathrm{}}^{\mathrm{}}𝒫_{a,b(a)}𝑑a`$ using a proper choice of $`a`$, $`b(a)`$. As an example, we will calculate the function $`𝒫_\psi `$ and the distribution $$𝒫_Q=\delta \left(t+(i\omega L_c\nu S)^2Q_{33}^{12}(x)Q_{33}^{21}(x^{})\right)_F,$$ (8) characterizing the Landauer-type conductivity. Hereinafter, all the correlators are studied in the limit $`\omega 0`$. To express the distribution Eq. (6) in terms of a functional integral over the supermatrices $`Q`$, we relate the numerical coefficients of all moments of the distribution functions $`𝒫_\psi `$ and $`𝒫_Q`$, using the identity $$\frac{(2n2)!(1)^n}{\left((n1)!\right)^2}=\frac{\mathrm{d}}{\mathrm{d}\beta }\|_{\beta =1}_0^1dp(\beta p)^{n1}(\beta (1p))^{n1}.$$ (9) Thus, we reduce the entire distribution $`𝒫_\psi `$ to $`𝒫_Q`$ $$𝒫_\psi (t)=\frac{\mathrm{d}}{\mathrm{d}\beta }\|_{\beta =1}_0^1\frac{\mathrm{d}p}{\beta p(1p)}\delta (t+(i\omega L_cS)^2\beta ^2p(1p)Q_{33}^{12}(x)Q_{33}^{21}(x^{}))_F$$ (10) Eq. (10) can be proven taking $`𝒫_\psi \left(t\right)`$ from Eq. (6) and expanding both sides in $`\psi _\alpha `$ and $`Q`$ respectively. Then, one should use Eqs. (4, 5) and this gives finally Eq. (10) with the functional $`F`$ from Eq. (2). Eq. (10) allows us to calculate the distribution function of the amplitudes of the wave functions at two different points of an infinite sample. Although we calculate here the distribution functions of the quantities written at two different points $`x`$ and $`x^{}`$, our procedure can be extended to many-point functions. In quasi-1D case, one can calculate the functional integral, Eq. (7), exactly using the transfer-matrix technique. Within this method the functional integration is reduced to solving an effective “Schrödinger equation” in the space of variables parametrizing $`Q`$-matrices. Solving the “Schrödinger equation” and calculating certain “matrix elements” that are definite integrals over $`Q`$ one can obtain any physical quantity of interest. For instance, the moments corresponding to the distribution function $`𝒫_{a,b}`$, Eq. (7), read $$T_{mn}=t_1^nt_2^m_{𝒫_Q}=a^nb^m\mathrm{\Psi }(Q)\left(Q_{33}^{12}\right)^n\mathrm{\Gamma }(x,x^{};Q,Q^{})\left(Q_{33}^{\mathrm{\hspace{0.17em}21}}\right)^m\mathrm{\Psi }(Q^{})dQdQ^{},$$ (11) where the function $`\mathrm{\Psi }\left(Q\right)`$ is the partition functions of the parts of the wire to the left from the point $`x`$ and to the right from the point $`x^{}`$, and the function $`\mathrm{\Gamma }`$ represents the partition functions of the segment between $`x`$ and $`x^{}`$. The meaning of the functions $`\mathrm{\Psi }`$ and $`\mathrm{\Gamma }`$ becomes clear after the reduction of the functional integrals to the effective Schrödinger equation. The function $`\mathrm{\Psi }\left(Q\right)`$ appears to be the ground state of the effective Hamiltonian $`_Q`$ acting in the space of the elements of the $`Q`$-matrix whereas the function $`\mathrm{\Gamma }`$ is the Green function in this space. They satisfy the following equations $`_Q\mathrm{\Psi }\left(Q\right)`$ $`=`$ $`0`$ (12) $`\left({\displaystyle \frac{}{x}}+_Q\right)\mathrm{\Gamma }(x,x^{};Q,Q^{})dx`$ $`=`$ $`\delta \left(xx^{}\right)`$ (13) The fact that the energy of the ground state is zero is a consequence of the supersymmetry of the initial $`\sigma `$-model, Eq. (2). The explicit form of the Hamiltonian $`_Q`$ depends on the parametrization of the $`Q`$-matrix and its choice depends on the quantities calculated and the physics involved. The supermatrices $`Q`$ consist of compact and non-compact blocks containing circular and hyperbolic functions, respectively. In the cases of zero and strong magnetic fields only the non-compact sector survives in the limit $`\omega 0`$, which greatly simplifies the whole analysis. If a non-compact variable $`\lambda _1`$ is relevant only, the function $`\mathrm{\Psi }`$ takes in the limit of small $`\omega `$ the well-known form $$\mathrm{\Psi }=2z^{1/2}K_1(2z^{1/2}),\text{ }z=i\omega \lambda _1$$ (14) valid both for chains and thick wires. Using the modified Bessel functions $`K_{i\rho }`$ one can write in this limit the solution for the function $`\mathrm{\Gamma }`$ as $$\mathrm{\Gamma }(x,x^{};z,z^{})=\frac{8}{\pi ^2}_0^{\mathrm{}}\mathrm{e}^{\|xx^{}\|(1+\rho ^2)/4L_c}z^{1/2}K_{i\rho }(2z^{1/2})z^{1/2}K_{i\rho }(2z^{1/2})\mathrm{sinh}(\pi \rho )\rho d\rho .$$ (15) where $`L_c`$ is the localization length. Eq. (15) will be proven in Sec. V. Actually, it follows from the orthonormality of the modified Bessel functions $`K_{i\rho }`$ (see Eq. (43)) and the fact that the functions $`2z^{1/2}K_{i\rho }(2z^{1/2})`$ are eigenfunctions of the “heat” equation with the Laplacian in the space of the matrices $`Q`$. The expressions for $`\mathrm{\Psi }`$ and $`\mathrm{\Gamma }`$, Eqs. (14, 15), are sufficient to calculate all the moments, Eq. (11), for the orthogonal and unitary ensembles. This way is more straightforward and more economical, as compared to the usual way of representing the moments as solutions of differential equations obtained within the transfer matrix technique. At the same time, the conventional approach is more general and is applicable also in the crossover regime. It is worth mentioning that our analysis of the localization properties in infinite wires is somewhat easier than that for conductance and wave functions in finite wires. This is because only large values of the non-compact variables $`\lambda _1`$ are important for the calculations in our case. This is this simplification that allows us to calculate all moments as well as the entire distribution functions for the ‘pure’ (orthogonal and unitary) ensembles (see Sec. V). ## III Transfer-Matrix Technique At not very high frequencies, the zero transversal space harmonics of the supermatrix $`Q`$ in Eq. (2) gives the main contribution. In order to neglect the non-zero harmonics one should choose the London gauge for the vector potential $`𝐀`$. Writing $`𝐀=(Hy,0,0),`$ where $`H`$ is the external magnetic field perpendicular to the wire we represent the free energy, Eq. (2), for the quasi-1D sample in the form $$F[Q]=\frac{L_{\mathrm{cu}}}{16}\mathrm{Str}\left[\left(_xQ\right)^2+\frac{X^2}{16L_{\mathrm{cu}}^2}[Q,\tau _3]^2+\frac{2i\omega \mathrm{\Lambda }Q(x)}{D}\right]dx,$$ (16) where $`L_{\mathrm{cu}}=2\pi S\nu D`$ is the localization length for the unitary ensemble. The crossover parameter $`X`$ depends on the geometry of the sample, $`X=2\pi \varphi /\varphi _0`$, $`\varphi _0=hc/e`$ is the flux quantum, and $`\varphi =HL_{\mathrm{cu}}y^2_{\mathrm{sec}}^{1/2}`$ is the magnetic flux through the area limited by the localization length. The brackets imply averaging across the wire over the coordinate $`y`$ which is chosen to be perpendicular both to the direction of the wire and the magnetic field. This averaging gives $`y^2^{1/2}=d/\sqrt{12}`$, where $`d`$ is the width of the wire, for “flat” wires made on the basis of a 2D gas. For wires with a circular cross-section, $`y^2^{1/2}=d/4`$, where $`d`$ is the diameter. To determine the functions $`\mathrm{\Psi }`$ and $`\mathrm{\Gamma }`$ in Eq. (11), we discretize the wire by introducing sites on which the $`Q`$-supermatrix varies. The function $`\mathrm{\Psi }\left(Q\right)`$ is the partition function of a semi-infinite sample and therefore it does not depend on the coordinate $`x`$. This allows us to write immediately the a recurrence equation. Relating the function $`\mathrm{\Psi }(Q)`$ on one site and its value $`\mathrm{\Psi }\left(Q^{}\right)`$ on the neighboring one we obtain $$\mathrm{\Psi }(Q)=N(Q,Q^{})Z_0(Q^{})\mathrm{\Psi }(Q^{})dQ^{}.$$ (17) where the function $`N(Q,Q^{})`$ describing the coupling between the sites $$N(Q,Q^{})=\mathrm{exp}\left(\frac{L_{\mathrm{cu}}}{4\mathrm{\Delta }l}\mathrm{Str}QQ^{}\right),$$ (18) originates from the kinetic term $`(Q)^2`$ in the free energy, Eq. (16). The function $`Z_0\left(Q\right)`$, $$Z_0(Q)=\mathrm{exp}\left(\frac{i\omega L_{\mathrm{cu}}\mathrm{\Delta }l}{4D}\mathrm{Str}\mathrm{\Lambda }Q+\frac{X^2\mathrm{\Delta }l}{16L_{\mathrm{cu}}}\mathrm{Str}[Q,\tau _3]^2\right)$$ (19) describes the remaining terms in the free energy functional $`F`$, Eq. (16), and should be integrated on each site. In order to reduce the integral Eq. (17) to the differential “Schrödinger equation” for the functions $`\mathrm{\Psi }`$ (and similarly for $`\mathrm{\Gamma }`$), one should take the continuous limit $`\mathrm{\Delta }l0`$, where $`\mathrm{\Delta }l`$ is the length of the separation between the sites. In this continuous limit, $`\mathrm{\Delta }l`$ drops out from final expressions and therefore can be chosen to be arbitrarily short. Since the kernel $`\mathrm{\Gamma }`$ enters equations for each moment only in a combination with $`(Q^{21})^n\mathrm{\Psi }(Q)`$, it is advantageous to introduce a matrix function $`P_k^{(n)}`$ such that $$P_k^{(n)}(Q)=\underset{x}{}\mathrm{exp}[i(xx^{})k]\mathrm{\Gamma }(k;Q,Q^{})(Q^{\mathrm{\hspace{0.17em}21}})^n\mathrm{\Psi }(Q^{})dQ^{},$$ (20) where $`\mathrm{\Gamma }(k;Q,Q^{})`$ is the Fourier transform of $`\mathrm{\Gamma }(xx^{};Q,Q^{})`$. The matrix function $`P_k^{(n_)}(Q)`$ on a site is related to its value $`P_k^{(n_)}(Q^{})`$ on the neighboring site as $$P_k^{(n)}(Q)\mathrm{exp}(ik\mathrm{\Delta }l)N(Q,Q^{})Z_0(Q^{})P_k^{(n)}(Q^{})dQ^{}=(Q^{21})^n\mathrm{\Psi }(Q).$$ (21) The moments of the distribution $`𝒫_Q`$, expressed in terms of the functions $`\mathrm{\Psi }`$ and $`P_k`$, are given by $$T_{mn}=\mathrm{\Psi }(Q)(Q_{33}^{21})^n(P_{k\mathrm{\hspace{0.17em}33}}^{(m)}+P_{k\mathrm{\hspace{0.17em}33}}^{(m)})dQ,$$ (22) Solving Eqs. (20-21) and calculating the integral in Eq. (22) one can find, in principle, all the moments. However, further simplifications are necessary to obtain the results explicitly. A spectral expansion is convenient for the analysis of the function $`P_k^{(n)}`$. So, we expand the functions $`P_k^{(n)}`$ in eigenfunctions $`\varphi _{}(Q)`$, introduced as $$N(Q,Q^{})Z_0(Q^{})\varphi _{}(Q^{})dQ^{}=\varphi _{}(Q).$$ (23) The matrix functions $`\varphi _{}(Q)`$ have the same structure as $`Q^{21}`$. Their orthogonality and the normalization properties can be written in the form $$[\varphi _{^{}}^+(Q)]_{33}[\varphi _{}(Q)]_{33}dQ=\delta (^{}).$$ (24) Expanding $`P_k^{(m)}`$ in the complete set of the functions $`\varphi _{}`$, Fourier transforming it to the coordinate representation, and substituting the result into Eq. (11), we arrive at $$T_{mn}=\underset{}{}\mathrm{exp}(\|xx^{}\|)(Q_{33}^{12})^m[\varphi _{}(Q)]_{33}\mathrm{\Psi }(Q)dQ(Q_{33}^{21})^n[\varphi _{}(Q)]_{33}\mathrm{\Psi }(Q)dQ.$$ (25) Eq. (25) contains the sum over all eigenstates that should be found from Eq. (23) and establishes a link between the eigenvalue problem in the space of parameters of the $`Q`$-matrix and the localization properties. The lowest non-zero eigenenergy corresponding to the first excited state determines the exponential decay of the correlation functions at large distances. ## IV Effective Hamiltonians Eqs. (17, 23, 25) form a closed system completely solving the problem. To perform explicit calculations a parametrization of the $`Q`$-matrices should be chosen. For the problem involved one can use either the “standard” or the “magnetic” parametrization. The former parametrization has been used, e.g., for calculation of the density-density correlation function in disordered wires for the orthogonal and unitary ensembles, while the latter one helped to obtain explicit formulae for the crossover between the ensembles in the zero-dimensional situation. For the problem of the localization in wires the standard parametrization allows to calculate any quantity of interest for the orthogonal and unitary ensembles and it will be used in Sec. V for calculations of all moments as well as of the distribution function of the density-density correlator. At the same time, the standard parametrization is of little help in the crossover regime between the ensembles and we will use in this case the magnetic parametrization. The basic equations and their solutions for the orthogonal and unitary ensembles can be found in Refs. and we do not write them here. Instead, let us present explicit differential equations that can be obtained within the magnetic parametrization from Eqs. (17, 23, 25). The $`Q`$-matrix in this parametrization is constructed from two sets of variables, cooperon ones being cut by the magnetic field, and diffuson ones insensitive to $`H`$. An explicit form of the supermatrices $`Q`$ in these two parametrizations can be found elsewhere. The supermatrix $`Q`$ in this parametrization is represented as $$Q=V_dQ_c\overline{V}_d\text{}$$ (26) where the supermatrix $`Q_c`$, $`Q_c^2=1`$, contains cooperon “degrees of freedom” and $`V_d`$, $`\overline{V}_dV_d=1`$, diffuson ones. To make the reading easier, we present here the essential structure of the block $`Q^{12}`$ in the magnetic parametrization, keeping only those terms that contribute to the correlators considered $$\begin{array}{c}Q^{12}=u_di(1\widehat{\lambda }_d^2)^{1/2}[\widehat{\lambda }_c+2(\eta _c\eta _c^{}\kappa _c\kappa _c^{})(\lambda _{1c}\lambda _c)]\overline{v}_d,\end{array}$$ (27) where $`\widehat{\lambda }_{c,d}=\mathrm{diag}(\mathrm{cos}\theta _{c,d},\mathrm{cosh}\theta _{1c,\mathrm{\hspace{0.17em}1}d})`$, $`v_d`$ and $`u_d`$ are the standard $`4\times 4`$ unitary matrices in the $`V_d`$ (diffuson) block, and $`\eta _c`$, $`\eta _c^{}`$, $`\kappa _c`$, $`\kappa _c^{}`$ are the Grassmann variables from the cooperon block. An advantageous feature of the magnetic parametrization is that the term containing the magnetic field $`H`$ in the free energy has a simple form $$\mathrm{Str}[Q,\tau _3]^2=16(\lambda _{1c}^2\lambda _c^2).$$ (28) In contrast, the term with the frequency $`\omega `$ is quite complex since it contains the Grassmann variables $$\mathrm{Str}(\mathrm{\Lambda }Q)=4[\lambda _{1d}\lambda _{1c}\lambda _d\lambda _c+2(\eta _c\eta _c^{}\kappa _c\kappa _c^{})(\lambda _{1c}\lambda _c)(\lambda _{1d}\lambda _d)].$$ (29) The functions $`\mathrm{\Psi }`$ and $`\varphi _{}`$ correspond to states having respectively the “angular momentum” zero and one with respect to the unitary rotations $`V_d`$. The form of the functions $`\varphi _{}`$ is determined by the RHS of Eq. (21) and, as a result, it should be searched for in the class of functions with the same structure as the matrix $`Q^{12}`$ $$\varphi _{}=v_d\left(\begin{array}{cc}(1\lambda _d^2)^{1/2}f& 0\\ 0& i(\lambda _{1d}^21)^{1/2}f_1\end{array}\right)\overline{u}_d.$$ (30) where $`v_d`$ and $`u_d`$ are unitary matrices containing Grassmann variables only. The different forms of the functions $`\mathrm{\Psi }`$ and $`\varphi _\text{ }`$ lead to different forms of the effective Hamiltonians determining the functions $`\mathrm{\Psi }`$ and $`\varphi _{}`$. The functions $`f`$ and $`f_1`$ do not contain the diffuson Grassmann variables but may contain all other variables. The Schrödinger equations corresponding to Eqs. (17) and (23) are obtained considering slow variations of $`Q`$ on neighboring sites. The matrix $`Q^{}`$ is expanded near $`Q`$ and integrated over small variations of the matrix elements. The expansion up to quadratic variations is sufficient and therefore, the quadratic form obtained is identical to the one describing the square of the elementary length in the space of the supermatrices $`Q`$. This elementary length determines the Berezinian corresponding to a chosen parametrization. Correspondingly, the “Laplacian” entering the effective Schrödinger equation contains the Berezinian of the transformation. The Berezinian of the magnetic parametrization reads $$J=J_cJ_dJ_{cd},\text{ }J_{c,d}=\frac{1}{2^6\pi ^2(\lambda _{1c,d}\lambda _{c,d})^2},\text{ }J_{cd}=\frac{4\lambda _c^2}{(\lambda _{1c}+\lambda _c)^2},$$ (31) The parts $`J_d`$ and $`J_{cd}`$ of the Berezinian Eq. (31) originate from the diffuson Grassmann variables $`\eta _d`$, $`\kappa _d`$, whereas $`J_c`$ comes from the cooperon ones $`\eta _c`$, $`\kappa _c`$. The diffuson Grassmann variables do not enter Eqs. (28) and (29). In some sense, the diffuson variables decouple from other variables and this simplifies calculations. The situation with the cooperon variables is more difficult. Since the free energy, Eqs. (16, 29), contains explicitly the cooperon Grassmann variables $`\eta _c`$ and $`\kappa _c`$, so do the functions $`\mathrm{\Psi }`$ and $`\varphi _{}`$. In order to write equations for these functions explicitly one should expand them in powers of the Grassmann variables. Proceeding in this way one obtains a system of differential equations for the functions entering the expansion. Another complication occurs due to a non-trivial angular dependence on $`Q_c`$. This matrix is surrounded by the diffuson modes and couples to them. As a result, one comes to extremely cumbersome equations that can hardly be solved analytically for arbitrary $`X`$. The solution can be found for $`X1`$ but this limit corresponds to the well studied unitary ensemble. Fortunately, an important information about correlation functions can be extracted in the limit $`X1`$. To understand why this limit helps to study the problem we remind the reader that the possibility to obtain closed expressions for correlation functions in the orthogonal and unitary ensembles is due to the fact that the main contribution to integrals comes in the limit of small frequencies $`\omega `$ from large values of the non-compact variables (in the unitary ensemble $`\lambda _{1d}1/\omega `$). For large values of $`\lambda _{1d\text{ }}`$ the partial differential equations for the functions $`\mathrm{\Psi }`$ and $`\varphi _{}`$ are sufficiently simple and can be solved. In the crossover regime when $`X`$ is finite, the equations are still very complicated even in the limit $`\omega 0`$. However, in the limit $`X1`$, some quantities of interest are determined by large $`\lambda _{1c}1/X1`$. This leads to an additional simplification of the equations and the possibility to make estimates. At the same time, not any correlation function can be calculated in this way, which limits the applicability of the trick. In the continuous limit, Eq. (17) can be reduced to the form $$_X\mathrm{\Psi }=0$$ (32) Due to the presence of the Grassmann variables in the free energy functional, Eqs. (16, 29), the solution $`\mathrm{\Psi }`$ must contain them too, which is in contrast to the solutions for the orthogonal and unitary ensembles. The function $`\mathrm{\Psi }`$ can be represented as $$\mathrm{\Psi }=\mathrm{\Psi }_0(\lambda _{1d},\lambda _{1c})+\mathrm{\Psi }_1(\lambda _{1c},\lambda _{1d})\left(\eta _c\eta _c^{}\kappa _c\kappa _c^{}\right)+\mathrm{\Psi }_2(\lambda _{1c},\lambda _{1d})\eta _c\eta _c^{}\kappa _c\kappa _c^{}$$ (33) Then, one obtains a system of partial differential equations for the functions $`\mathrm{\Psi }_0`$, $`\mathrm{\Psi }_1`$, $`\mathrm{\Psi }_2`$ $`\left(_X^{(0)}+\stackrel{~}{\omega }x_1\right)\mathrm{\Psi }_0+2{\displaystyle \frac{\lambda _{1c}\lambda _c1}{\left(\lambda _{1c}\lambda _c\right)^2}}\mathrm{\Psi }_1=0`$ $$\left(_X^{(0)}+\stackrel{~}{\omega }x_1\right)\mathrm{\Psi }_1\frac{\lambda _{1c}\lambda _c1}{\left(\lambda _{1c}\lambda _c\right)^2}\mathrm{\Psi }_2+\stackrel{~}{\omega }x_2\mathrm{\Psi }_0=0$$ (34) $`\left(_X^{(0)}+\stackrel{~}{\omega }x_1\right)\mathrm{\Psi }_22\stackrel{~}{\omega }x_2\mathrm{\Psi }_1=0`$where $$_X^{(0)}=\underset{i=d,1d}{}\frac{1}{J_d}\frac{}{\lambda _i}J_d\|1\lambda _i^2\|\frac{}{\lambda _i}\underset{i=c,1c}{}\frac{1}{J_cJ_{cd}}\frac{}{\lambda _i}J_cJ_{cd}\|\lambda _i^21\|\frac{}{\lambda _i}+X^2\left(\lambda _{1c}^2\lambda _c^2\right).$$ (35) $`\stackrel{~}{\omega }=2\pi ^2\left(\nu S\right)^2D\omega `$and $$x_1=\lambda _{1d}\lambda _{1c}\lambda _d\lambda _c,\text{ }x_2=2\left(\lambda _{1c}\lambda _c\right)\left(\lambda _{1d}\lambda _d\right)$$ (36) In the limit $`\lambda _{1d}`$, $`\lambda _{1c}1,`$ the difference between equations for the functions $`\mathrm{\Psi }`$ and $`\varphi _{}`$ is negligible because it comes from the averages of the type $`\mathrm{\Delta }\eta _i\mathrm{\Delta }\eta _i^{}1/\lambda _{1i}`$, $`\mathrm{\Delta }\kappa _i\mathrm{\Delta }\kappa _i^{}1/\lambda _{1i}`$, $`i=c`$, $`d`$. In this limit, Eqs. (34, 35, 36) take a simpler form $`\left(_X^{(0)}+\stackrel{~}{\omega }x_1\right)\mathrm{\Psi }_0=0`$ $$\left(_X^{(0)}+\stackrel{~}{\omega }x_1\right)\mathrm{\Psi }_1+\stackrel{~}{\omega }x_2\mathrm{\Psi }_0=0$$ (37) $`\left(_X^{(0)}+\stackrel{~}{\omega }x_1\right)\mathrm{\Psi }_22\stackrel{~}{\omega }x_2\mathrm{\Psi }_1=0`$with $$_X^{(0)}=\lambda _{1d}^2\frac{^2}{\lambda _{1d}^2}\lambda _{1c}^2\frac{^2}{\lambda _{1c}^2}+2\lambda _{1c}\frac{}{\lambda _{1c}}+X^2\lambda _{1c}^2$$ (38) $`x_1=x_2/2=\lambda _{1c}\lambda _{1d}`$The same equations can be written in the limit $`\lambda _{1d}`$, $`\lambda _{1d}1`$ for the functions $`\varphi _{}`$. It is clear that Eqs. (33, 36) cannot be solved analytically for an arbitrary $`X`$ and the only hope is to solve Eqs. (37, 38). Of course, one is restricted then by calculation of those quantities for which large $`\lambda _{1d},\lambda _{1c}`$ give the main contribution. In the opposite limiting case $`X\mathrm{}`$ but $`\omega 0`$ all the cooperon variables are frozen and one comes to the effective Hamiltonian $$=\lambda _1^2\frac{^2}{\lambda _1^2}\frac{2i\omega L_{\mathrm{cu}}^2}{D}\lambda _1$$ (39) After a proper change of the localization length the same Hamiltonian can be written in the orthogonal case using the standard parametrization. For completeness let us remind that the localization lengths for the orthogonal and unitary ensembles, $`L_{\mathrm{co}}`$ and $`L_{\mathrm{cu}}`$ respectively, are given by the following expressions $$\begin{array}{cc}L_{\mathrm{cu}}=2\pi \nu SD,& L_{\mathrm{co}}=\pi \nu SD,\end{array}$$ (40) In these cases the function $`\mathrm{\Psi }`$ does not contain Grassmann variables. The equations to be solved take the form $$\mathrm{\Psi }=0,\varphi _{}=L_c\varphi _{}$$ (41) where the localization length $`L_c`$ equals either $`L_{\mathrm{co}}`$ or $`L_{\mathrm{cu}}`$ depending on the ensemble. In the next Section, solving Eqs. (41) we will calculate all moments and the entire distribution function of the density-density correlations. This will help us to come to certain conclusions about properties of the electron wave functions in the orthogonal and unitary ensembles. ## V Calculation of Moments and Distributions for the Orthogonal and Unitary Ensembles. Calculation of the moments and of the distribution functions is important for the understanding of the structure of the wave functions and properties of different averages. Higher moments and the distribution function for infinitely long wires have not been considered yet although the results are known for finite wires. So, calculation of the moments and of the distribution function can be interesting on its own but, what is more important, this information will also help us to make certain conclusions about properties of the wave functions in a weak magnetic field where the possibility of exact calculations is more limited. This will help us to extend the results of Ref. where only the averaged density-density correlation function was calculated. Solutions $`\mathrm{\Psi }`$ and $`\varphi _{}`$ of Eqs. (39, 41) are well known (see e.g. ) and can be written as $$\begin{array}{ccc}\mathrm{\Psi }=2z^{1/2}K_1(2z^{1/2}),& \varphi _{}=2z^{1/2}K_{i\rho }(2z^{1/2}),& =(1+\rho ^2)/4L_c,\end{array}$$ (42) where $`z=2i\lambda _1\omega L_c^2/D`$. Comparing Eqs. (25) and (11), we can understand that the propagator has the form of Eq. (15). The numerical coefficient in Eq. (15) is found using the Lebedev-Kontorovich transformation applicable for an arbitrary function $`\phi (y)`$ $$\begin{array}{cc}\stackrel{~}{\phi }(\rho )=_0^{\mathrm{}}\phi (y)K_{i\rho }(y)\frac{\mathrm{d}y}{y},& \phi (y)=\frac{2}{\pi ^2y}_0^{\mathrm{}}\stackrel{~}{\phi }(\rho )K_{i\rho }(y)\rho \mathrm{sinh}(\pi \rho )d\rho .\end{array}$$ (43) Now we can calculate all the moments of the two-point correlator, Eq. (11). Substituting Eqs. (42) into Eq. (25), and calculating integrals over all elements of the supermatrices $`Q`$ except $`\lambda _1`$, yields $$T_{mn}=t_1^nt_2^m_{𝒫_Q}=\frac{32}{\pi ^2}_0^{\mathrm{}}d\rho \rho \mathrm{sinh}(\pi \rho )e^{\frac{u}{4}\left(\rho ^2+1\right)}M_n\left(\rho \right)M_m\left(\rho \right).$$ (44) where $`M_n={\displaystyle dt_1\left(\frac{t_1}{2}\right)^{2n1}n^2K_1(t_1)K_{i\rho }(t_1)}.`$and $`u=\|xx^{}\|/L_c`$. The structure of Eq. (44) is not complicated and the integration over the variable $`t_1`$ can be performed exactly. As a result, we come to general expressions for all moments of the density-density correlations $`p_{\mathrm{}}^{\left(n\right)}\left(u\right)`$ introduced in Eq. (1) $$p_{\mathrm{}}^{(n)}\left(u\right)=\frac{2\nu }{\pi }\frac{(2n2)!}{\left((n1)!\right)^2}_0^{\mathrm{}}d\rho \rho \mathrm{sinh}(\pi \rho )\mathrm{e}^{\frac{u}{4}\left(\rho ^2+1\right)}\left[P^{(n)}(\rho )\right]^2,$$ (45) $`P^{(n)}(\rho )={\displaystyle \frac{n^2}{(2n1)!}}\mathrm{\Gamma }\left(n+{\displaystyle \frac{1}{2}}+{\displaystyle \frac{i\rho }{2}}\right)\mathrm{\Gamma }\left(n+{\displaystyle \frac{1}{2}}{\displaystyle \frac{i\rho }{2}}\right)\mathrm{\Gamma }\left(n{\displaystyle \frac{1}{2}}+{\displaystyle \frac{i\rho }{2}}\right)\mathrm{\Gamma }\left(n{\displaystyle \frac{1}{2}}{\displaystyle \frac{i\rho }{2}}\right).`$For $`n=1`$, Eq. (45) reduces to the expression obtained by Gogolin for a single chain. Eq. (45) is applicable at any $`x`$ including the values smaller and of the order of $`L_c`$ for both the unitary and orthogonal ensembles. Comparing the moments, Eqs. (45), with those obtained for finite wires , we see that their structure is similar, although they are not completely identical. At large distances, $`\|xx^{}\|L_c`$, essential contribution to the integral over $`\rho `$ in Eq. (45) comes from small $`\rho 1`$ and the integral can be easily calculated. We see that the moments of all quantities decay similarly as in infinite wires, Eqs. (45), (44), finite closed wires , as well as in open chains $$p_{\mathrm{}}^{(n)}(x)=C_n\left(\frac{1}{\|x\|}\right)^{3/2}\mathrm{exp}\left(\frac{\|x\|}{4L_c}\right).$$ (46) where $`C_n`$ is a constant depending on $`n`$. The exponential decay of all moments seen from Eq. (46) manifests the localization of the wave functions. A very important feature of the moments $`p_{\mathrm{}}^{(n)}(x)`$, Eq. (46), is that the dependence on the coordinate is the same for all $`n`$. This does not allow to interpret $`p_{\mathrm{}}^{\left(n\right)}\left(x\right)`$, Eq. (46), as functions characterizing the typical shape of the wave functions. The exact representation of the moments of the wave functions, Eq. (44), allows to obtain the entire distribution function of the density-density correlations. Of the crucial importance is the fact that Eq. (44) is valid for all $`n`$ including the limit $`n\mathrm{}`$. It is clear, if the $`n`$th moment of a quantity can be represented in a form of an integral of a function multiplied by the $`n`$th power of the variable of integration, then this function is the distribution of this quantity. In the case under consideration, Eq. (44) contains, besides the necessary power of the variable $`t`$, the factor $`n^2`$, originating from the expansion of the element $`(Q_{33}^{12})^n`$ in the Grassmann variables. This does not make the calculation of the distribution function more difficult because the presence of this factor leads merely to additional derivatives, and we obtain for the distribution function $`𝒫(t_1,t_2)`$, Eq. (7) $$𝒫(t_1,t_2)=\frac{2}{\pi ^2}_0^{\mathrm{}}d\rho \rho \mathrm{sinh}(\pi \rho )\mathrm{e}^{\frac{u}{4}\left(\rho ^2+1\right)}\mathrm{\Pi }_\rho (t_1)\mathrm{\Pi }_\rho (t_2),$$ (47) where $$\mathrm{\Pi }_\rho (t)=\left(\left[K_1(2t^{1/2})K_{i\rho }(2t^{1/2})\right]^{}t\right)^{}.$$ (48) and stands for the derivative over $`t`$. The two-point distribution function $`𝒫(t_1,t_2)`$ can be used to obtain any other correlator of interest. For instance, the distribution of the Landauer-type conductivity $`𝒫_Q(t)`$, Eq. (8) can be represented as the convolution $$𝒫_Q=d\tau \delta (t\tau i\omega \nu L_cSQ_{33}^{12}(x_1))\delta (\tau i\omega \nu L_cSQ_{33}^{21}(x_2))_F.$$ (49) Then, we obtain for the function $`𝒫_Q\left(t\right)`$ $$𝒫_Q(t)=\frac{2\nu }{\pi ^2}_0^{\mathrm{}}d\rho \rho \mathrm{sinh}(\pi \rho )\mathrm{e}^{\frac{u}{4}\left(\rho ^2+1\right)}_0^{\mathrm{}}\frac{\mathrm{d}\tau }{\tau }\mathrm{\Pi }_\rho (\tau )\mathrm{\Pi }_\rho (t/\tau ),$$ (50) Eq. (50) is also exact and applicable at any distances. It is interesting to find this function near its maximum and this can be done easily in the limit $`\|xx^{}\|L_c`$. At very large separation between $`x`$ and $`x^{}`$, typical values of $`t`$ are exponentially small in distributions $`𝒫_Q`$ and $`𝒫_\psi `$, Eqs. (8, 10). A dominant contribution to the integral, Eq. (50), comes from the region $`\tau t1`$. Eq. (48) allows us to write the asymptotics of the function $`\mathrm{\Pi }_\rho (\tau )`$ at $`\tau 1`$ and $`\tau 1`$ $`\mathrm{\Pi }_\rho (\tau )Re\left[\mathrm{\Gamma }(i\rho )/4\tau ^{3/2+i\rho /2}\right],\text{ }\tau 1`$ $$\mathrm{\Pi }_\rho (\tau )\pi \mathrm{exp}\left(4\tau ^{1/2}\right)/\tau ^{1/2},\text{ }\tau 1$$ (51) which ensures convergence of the integral over $`\tau `$ at $`\tau 0`$ in Eq. (50). Performing integration in the resulting expression $`𝒫_Q(\mathrm{ln}t)={\displaystyle \frac{\pi }{32}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}d\rho \rho \mathrm{sinh}(\pi \rho ){\displaystyle \frac{\mathrm{\Gamma }(i\rho )\mathrm{\Gamma }(2+i\rho )}{t^{1/2}}}\mathrm{exp}\left({\displaystyle \frac{u}{4}}{\displaystyle \frac{\mathrm{ln}^2t}{4u}}{\displaystyle \frac{u}{4}}(\rho +{\displaystyle \frac{i\mathrm{ln}t}{u}})^2\right),`$in the saddle point approximation, we obtain $$𝒫_Q(\mathrm{ln}t)\frac{1}{2\pi ^{1/2}u^{1/2}}\mathrm{exp}\left(\frac{1}{4u}(u+\mathrm{ln}t)^2\right).$$ (52) Considering the distribution $`𝒫_\psi `$, Eq. (10), under the same assumptions it is not difficult to check that in the limit $`\|xx^{}\|L_c`$ it reduces to Eq. (52). The saddle-point approximation used for the derivation of Eq. (52) is applicable at large distances $`\|xx^{}\|L_c`$ when $`u`$ is large. Eq. (52) shows that in this limit the distribution functions are log-normal. The result about the log-normal distribution of physical quantities like conductance agrees with the corresponding result obtained for chains and is universal. It shows that the localization length of the wave functions or in, other words, the logarithm of the conductance is a self-averaging quantity. Considering very long samples one has to obtain the localization length $`L_c`$ and fluctuations of this quantity can be neglected. At the same time, the log-normal distribution, Eq. (52), is applicable for not very large values of $`t`$ only. The far tail of the distribution function $`𝒫_\psi `$ decays more slowly. The moments of the density-density correlations $`p_{\mathrm{}}^{\left(n\right)}`$, Eq. (46), cannot be obtained from Eq. (52) because the main contribution comes from the far tails of the distribution function. All these peculiarities are discussed in details in Refs. . What is the physical meaning of the fact that all the moments $`p_{\mathrm{}}^{\left(n\right)}`$, Eq. (46) decay in the same way with a larger length $`4L_c`$? This can be understood if we imagine that there are big splashes of the wave functions even far from the localization center. Although, they can give a considerable contribution to the shape of the wave function, their probability decreases exponentially with $`\|x\|`$ (is proportional to $`\mathrm{exp}\left(\|x\|/4L_c\right)`$). Then, such splashes would give a similar contribution to all the moments. As concerns typical values of the squares of wave functions, they decay faster with the localization length $`L_c`$ and give a smaller contribution to the moments. The situation is opposite if one calculates the average logarithm of the wave functions. For the logarithm, the splashes do not give a specially large contribution, being at the same time very rare. So, their contribution to the average logarithm can be neglected. In other words, the average logarithm of the correlations of wave functions characterizes typical wave functions, whereas the moments describe the rare splashes of the wave functions. This understanding will be important for interpretation of the results in a weak magnetic fields obtained in the next Section. ## VI Effect of Magnetic Field on the Tails of Wave Functions This section is central in our paper. We try now describe electron wave functions in disordered wires in a weak magnetic field, such that $`X1`$. General equations describing this situation are very difficult even when they are written for the “ground state” function $`\mathrm{\Psi }`$, Eqs. (32-36). At the same time, any attempt to consider the weak magnetic field as a perturbation fails in the most interesting limit $`\omega 0`$ (one obtains terms like $`X^2/\stackrel{~}{\omega }`$). This reflects a clear physics: travelling a sufficiently long time the electron ‘feels’ any weak magnetic field. The failure of the perturbation theory signals immediately that the effect of the weak magnetic field is more complicated than a slight changing of the eigenenergies of the effective Schödinger equation and hence the localization length. The first excited state for the orthogonal ensemble has the eigenenergy $`_\mathrm{o}=1/2`$ and this energy definitely changes if the magnetic field is applied. But, as we will try to show below, an additional eigenstate with the eigenenergy $`_\mathrm{u}=1/4`$, the energy of the unitary ensemble, appears (strictly speaking, a continuum of states with the energies $`_\rho =(1+\rho ^2)/4`$, where $`\rho `$ is a continuous variable, should be added). Although at $`X1`$ this state contributes to correlation functions with a small weight, the value $`_\mathrm{u}=1/4`$ does not depend on $`X`$. As concerns the eigenenergy of the ground state, it is exactly zero due to the supersymmetry of the model. The eigenfunction $`\mathrm{\Psi }`$ of the ground state is a complicated function $`Q`$ depending on the ratio $`X^2/\stackrel{~}{\omega }.`$ We did not manage to find it for an arbitrary value of this ratio and can write in the limits $`X^2/\stackrel{~}{\omega }0`$ (orthogonal ensemble) and $`X^2/\stackrel{~}{\omega }\mathrm{}`$ (the limit we are interested in now) only. Although the solution for the orthogonal ensemble is very well known in the standard parametrization , it has not been written before in the magnetic parametrization. To find it in the limit $`\lambda _{1d}`$, $`\lambda _{1c}1`$ we use Eqs. (37, 38) taken at $`X=0`$. The solution for the functions $`\mathrm{\Psi }_i`$ $`\left(i=0,1,2\right)`$ can be sought in the form $$\mathrm{\Psi }_i(\lambda _{1d},\lambda _{1c})=\lambda _{1d}R_i\left(z\right)\text{}z=\lambda _{1d}\lambda _{1c}$$ (53) Then, we obtain the following equation for $`R_i`$ $$\left(_o+\stackrel{~}{\omega }z\right)R_0\left(z\right)=0$$ (54) $$\left(_o+\stackrel{~}{\omega }z\right)R_1\left(z\right)+2\stackrel{~}{\omega }zR_0\left(z\right)=0$$ (55) $$\left(_o+\stackrel{~}{\omega }z\right)R_24\stackrel{~}{\omega }zR_1\left(z\right)=0$$ (56) where $$_o=2z^2\frac{d^2}{dz^2}$$ (57) Eq. (54) is exactly the final equation obtained for the orthogonal ensemble in the standard parametrization. Writing the eigenfunction of the ground state $`\mathrm{\Psi }`$ in the form $$\mathrm{\Psi }=\lambda _{1d}R$$ (58) we find for $`R`$ $$R=R_0\left(z\left(1+2\left(\eta \eta ^{}\kappa \kappa ^{}\right)\right)\right)$$ (59) In the opposite limit $`X^2/\stackrel{~}{\omega }\mathrm{},`$ the solution of Eqs. (37, 38) is completely different. Putting $`\stackrel{~}{\omega }=0`$ in this equations we find a solution $$\mathrm{\Psi }_0\left(\lambda _{1c}\right)=\left(1+X\lambda _{1c}\right)\mathrm{exp}\left(X\lambda _{1c}\right)$$ (60) $`\mathrm{\Psi }_1=\mathrm{\Psi }_2=0`$Formally, one could write non-zero solutions for $`\mathrm{\Psi }_1`$ and $`\mathrm{\Psi }_2`$ proportional to $`\mathrm{\Psi }_0`$. However, we understand that in the magnetic parametrization the ground state function $`\mathrm{\Psi }`$ should not contain at $`\omega =0`$ Grassmann variables because the effective Hamiltonian does not contain them. So, in the limit $`\omega 0`$ the form of the ground state function $`\mathrm{\Psi }`$ changes discontinuously from Eqs. (58) to Eq. (60) as a magnetic field is applied. Physically, it may correspond to lifting degeneracies of some states of the initial electron problem by the magnetic field. The ground state solution, Eq. (60), is written for $`\lambda _{1d}`$, $`\lambda _{1c}1`$. We have checked that this state survives and does not considerably change its form solving Eqs. (34) in the limit $`\omega =0`$. The limit of the vanishing frequency considerably simplifies the equations because the “diffuson variable” $`\lambda _{1d}`$ decouples from the “cooperon variables” $`\lambda _{1c}`$ and $`\lambda _c`$. Then, one can search for a solution $`\mathrm{\Psi }`$ depending on the cooperon variables only. We used the standard over-relaxation method with the Chebyshev acceleration to found numerically that up to $`X`$ as large as $`0.3`$, Eq. (60), describes the solution rather well. At higher values of $`X`$, no dramatic changes occurred to this state either. It is clear that in the limit $`X\mathrm{}`$ essential values of the cooperon variables $`\lambda _{1c}`$ and $`\lambda _c`$ are close to unity for an arbitrary $`\omega `$ and one obtains the function $`\mathrm{\Psi }`$ for the unitary ensemble. The solution found at $`\omega =0`$ can be used for non-zero frequencies provided the variable $`\lambda _{1d}`$ is not very large. From Eqs. (37,38) we can estimate the region of the applicability of the solution as $$\lambda _{1d}\lambda _{1c}\stackrel{~}{\omega }1$$ (61) In order to find the density-density correlation function $`p_{\mathrm{}}\left(xx^{}\right)`$ and its moments $`p_{\mathrm{}}^{\left(n\right)}\left(xx^{}\right)`$ at large distances we have to find not only the ground state but also at least the first excited state $`\varphi _1`$. After that one should calculate the proper “matrix elements” in Eq. (25) or use Eq. (22). Although we argue that the eigenvalue $`_\mathrm{u}=1/4`$ of the lowest excited state is exact, we are not able, as previously, to write the proper wave function in all regions of the variables exactly. The solution for $`\varphi _1`$ can be sought in the form, Eq. (30), corresponding to the symmetry of the block $`Q^{12}`$ of the supermatrix $`Q`$, which is the usual way of constructing the proper excited eigenstates . Explicit formulae can be written at $`X1`$ provided the variables $`\lambda _{1d}`$, $`\lambda _{1c}`$ are large, $`\lambda _{1d}`$, $`\lambda _{1d}1`$, but limited from above by the inequality (61). In this limit one may keep the function $`f_1`$ only and we come to the following equation $$\left(_X^{\left(0\right)}2\lambda _{1d}\frac{}{\lambda _{1d}}\right)f_{1\rho }=\left(\rho \right)f_{1\rho }$$ (62) $`_X^{\left(0\right)}`$ is given by Eq. (38) and $`\rho `$ is a continuous real variable. The variables $`\lambda _{1d}`$ and $`\lambda _{1c}`$ in Eq. (62) separate from each other and we write the solution $`f_1`$ (see Eq. (30) in the form $$f_{1\rho }=\lambda _{1d}^{1/2}\chi _\rho \left(\mathrm{ln}\lambda _{1d}\right)\mathrm{\Psi }_0\left(\lambda _{1c}\right)$$ (63) with the function $`\mathrm{\Psi }_0`$, given by Eq. (60). The function $`\chi \left(v\right)`$ should be written as $$\chi _\rho \left(v\right)=\left(\frac{2}{\pi }\right)^{1/2}\mathrm{sin}\left[\left(\rho /2\right)v+\delta \left(\rho \right)\right]$$ (64) Eq. (64) corresponds to a picture of “free motion” suggested in Ref. for a description of the region of not very large $`\lambda _1`$. In the case considered here, Eq. (63), is applicable in the region determined by the inequality (61). In calculation of asymptotics of the density-density correlator at large distances small values of $`\rho 1`$ give the main contribution and, in this limit, the phase $`\delta `$ is also small, $`\delta \left(\rho \right)\rho `$ (see Ref. ). From Eq. (63) we obtain immediately $$\left(\rho \right)=\frac{1}{4}\left(1+\rho ^2\right)$$ (65) The form of the solution $`f_1`$, Eq. (63), differs from the ground state by a dependence on the diffuson variables $`\lambda _{1d}`$. This is the same difference as the one for the unitary ensemble. It is not difficult to check that the variable $`\lambda _{1d\text{ }}`$ separates from the cooperon variables even if $`\lambda _{1c}`$ is not large, $`\lambda _{1c}1`$. Then, the solution (63) may still be used in the region specified Eq. (61) provided the function $`\mathrm{\Psi }_0\left(\lambda _{1c}\right)`$ is replaced by the solution $`\mathrm{\Psi }`$ (depending only on the cooperon variables) of Eqs. (33, 34) taken at $`\omega =0`$. Eq. (63) with the function $`\mathrm{\Psi }`$ can be used also for an arbitrary $`X`$. In the limit $`X1`$, the function $`\mathrm{\Psi }`$ decays fast as the variables $`\lambda _{1c}`$ and $`\lambda _c`$ deviate from $`1`$. This means that the cooperon variables are ‘frozen’ and we come to the unitary ensemble. Then, the restriction (61) is not important and all formulae can be written for arbitrary $`\lambda _{1d}\omega `$ . Of course, in order to prove the existence of the states with the eigenvalues $`\left(\rho \right)`$, (63), rigorously one should investigate the exact equations for arbitrary values of the cooperon and diffuson variables and prove that the solutions ‘behave well’ everywhere. It is a difficult task that can apparently be performed only numerically. Assuming nevertheless that this state exists and is given for $`\lambda _{1d}`$, $`\lambda _{1c}1`$ by Eq. (63) (the inequality (61) should also be fulfilled) we can estimate the proper matrix elements, Eq. (25), determining the average density-density correlation function $`p_{\mathrm{}}\left(x\right)`$ and its moments $`p_{\mathrm{}}^{\left(n\right)}\left(x\right)`$, Eq. (1). However, trying to calculate the integrals in Eq. (25) one encounters a difficulty. The integrals over the matrix elements of $`Q`$ can be reduced to integrals over the variables $`\lambda _{1d}`$, $`\lambda _{1c},\lambda _d`$. $`\lambda _c`$ and the Grassmann variables. Proceeding in this way one should use the Berezinian $`J`$ of the transformation to these variables, Eq. (31), which is singular at $`\lambda _{1c}=\lambda _c=1`$. The singularity in the Berezinian leads usually to additional contributions. To avoid explicit calculations of these contributions it is more convenient to calculate not the physical quantities of interest themselves but their derivative over $`X`$. This leads to additional factors $`\lambda _{1c}^2\lambda _c^2`$ in integrands, thus compensating the singularities of the Berezinians. Calculation of the physical quantities using the spectral expansion, Eq. (25), is still very difficult. The alternative way of calculations is solving Eq. (21, and calculating the integral in Eq. (22). A simplified differential form of Eq. (21) can be written for the function $`P_k^{\left(n\right)}`$ determining the correlator $`p_{\mathrm{}}^{\left(n\right)}`$ as $$_X^{\left(0\right)}P_k^{(n)}+i\kappa P_k^{(n)}=\lambda _{1d}^n\lambda _{1c}^n\mathrm{\Psi }_0\left(\lambda _{1c}\right)$$ (66) where $`\kappa =kL_{\mathrm{cu}}`$. In Eq. (66), we do not write Grassmann variables explicitly implying that $`P^{\left(n\right)}`$ is the non-compact central part similar to $`\left(1\lambda _{1c}^2\right)^{1/2}f_1`$ in Eq. (30). The derivative over the field $`\left(T_{nn}\right)_X^{}`$ of the correlator $`T_{11},`$ Eq. (22) can be reduced to the form $$\frac{dT_{nn}\left(k\right)}{dX}\frac{d}{dX}\lambda _{1d}^{n2}\lambda _{1c}^{n4}\stackrel{~}{P}_{k33}^{\left(n\right)}𝑑\lambda _{1d}𝑑\lambda _{1c}$$ (67) Writing Eq. (67) the Berezinian $`J`$, Eq. (31), was used in the limit $`\lambda _{1c}`$, $`\lambda _{1d}1`$ and $`\stackrel{~}{P}_k^{\left(n\right)}=P_k^{\left(n\right)}+P_k^{\left(n\right)}`$. Using the fact that the Hamiltonian $`_X^{\left(0\right)}`$ contains the variables $`\lambda _{1d}`$ and $`\lambda _{1c}`$ separately let us expand the function $`P_k^{\left(n\right)}(\lambda _{1d},\lambda _{1c})`$ in the eigenfunctions of the diffuson part of the Hamiltonian. Writing $$P_k^{\left(n\right)}(\lambda _{1d},\lambda _{1c})=\underset{\rho }{}\lambda _{1d}^{1/2}\chi _\rho \left(\mathrm{ln}\lambda _{1d}\right)P_{k\rho }^{\left(n\right)}\left(\lambda _{1c}\right)$$ (68) we obtain the following equation for $`P_{k\rho }^{\left(n\right)}\left(\lambda _{1c}\right)`$ $$\left(\left(\rho \right)+i\kappa \right)P_{k\rho }^{\left(n\right)}\left(\lambda _{1c}\right)+\stackrel{~}{H}_X^{\left(0\right)}P_{k\rho }^{\left(n\right)}\left(\lambda _{1c}\right)=A\left(\lambda _{1c}\right)$$ (69) where $`A\left(\lambda _{1c}\right)={\displaystyle \lambda _{1d}^{n3/2}\chi _\rho \left(\mathrm{ln}\lambda _{1d}\right)\lambda _{1c}^n\mathrm{\Psi }_0\left(\lambda _{1c}\right)𝑑\lambda _{1d}}`$The operator $`\stackrel{~}{H}_X^{\left(0\right)}`$ is the part of the Hamiltonian $`_X^{\left(0\right)}`$ in Eq. (38) acting on the variables $`\lambda _{1c}`$. A similar expansion was done in Ref. to consider the case of a strong disorder. The main contribution in the integral over $`\lambda _{1d}`$ comes from large $`\lambda _{1d}`$ and the integral formally diverges. When considering the unitary ensemble one had to cutoff the integral by $`\lambda _{1c}1/\stackrel{~}{\omega }`$. At weak magnetic fields considered now, the cutoff can be determined from Eqs. (16, 29). We see that the term with the frequency can be neglected if $`\lambda _{1d}\lambda _{1c}\stackrel{~}{\omega }1`$, which leads to the cutoff $`\lambda _{1d}\left(\lambda _{1c}\stackrel{~}{\omega }\right)^1`$. So, we estimate the function $`A\left(\lambda _{1c}\right)`$ at $`\rho 1`$ as $$A\left(\lambda _{1c}\right)\left(i\stackrel{~}{\omega }\right)^{1/2n}\rho \lambda _{1c}^{1/2}\mathrm{\Psi }_0\left(\lambda _{1c}\right)$$ (70) In a similar way we can reduce Eq. (67) to the form $$\frac{dT_{nn}\left(k\right)}{dX}\left(i\stackrel{~}{\omega }\right)^{1/2n}\frac{d}{dX}P_{k\rho }^{\left(n\right)}\left(\lambda _{1c}\right)\mathrm{\Psi }_0\left(\lambda _{1c}\right)\frac{d\lambda _{1c}}{\lambda _{1c}^{7/2}}$$ (71) Now, we have to solve Eq. (69) and calculate the integral, Eq. (71). In order to solve Eq. (69) we introduce a Green function $`g`$ of this equation $$\left(\left(\rho \right)+i\kappa \right)g_{k\rho }(\lambda _{1c},\lambda _{1c}^{})+\stackrel{~}{H}_X^{\left(0\right)}g_{k\rho }(\lambda _{1c},\lambda _{1c}^{})=\lambda _{1c}^4\delta \left(\lambda _{1c}\lambda _{1c}^{}\right)$$ (72) and write the solution as $$P_{k\rho }^{\left(n\right)}=g_{k\rho }(\lambda _{1c},\lambda _{1c}^{})A\left(\lambda _{1c}^{}\right)\frac{d\lambda _{1c}^{}}{\left(\lambda _{1c}^{}\right)^4}$$ (73) Then, Eq. (71) is reduced to the form $$\frac{dT_{nn}\left(k\right)}{dX}\left(i\stackrel{~}{\omega }\right)^{12n}\frac{d}{dX}g_{k\rho }(\lambda _{1c},\lambda _{1c}^{})\mathrm{\Psi }_0\left(\lambda _{1c}\right)\mathrm{\Psi }_0\left(\lambda _{1c}^{}\right)\frac{\rho ^2d\rho d\lambda _{1c}d\lambda _{1c}^{}}{\left(\lambda _{1c}\lambda _{1c}^{}\right)^{7/2}}$$ (74) The last step to be done is to find the function $`g`$ and calculate the integral in Eq. (74). As usual, the combination $`\left(\rho \right)+i\kappa `$ enters the final formulae and one should perform Fourier transform to get the coordinate dependence. The point $`\kappa =i\left(\rho \right)`$ is a branching point and the integral over $`\kappa `$ can be shifted such that the integration is performed along the edges of the cut $`(i\left(\rho \right),i\mathrm{})`$. The point $`i\left(0\right)=i/4`$ gives the value of the localization length whereas the integration over $`\rho `$ leads to an additional power law prefactor $`\left\|x\right\|^{3/2}`$. Being interested mainly in determining the exponential decay we can immediately understand that all correlators $`T_{nn}`$ decay as $`\mathrm{exp}\left(\left\|x\right\|/4L_{\mathrm{cu}}\right)`$. To calculate the prefactor we may put $`\kappa =i\left(\rho \right)`$ in Eq. (72) and find the proper solution $`g_0(\lambda _{1c},\lambda _{1c}^{})`$, which satisfies the equation $$\stackrel{~}{H}_X^{\left(0\right)}g_0(\lambda _{1c},\lambda _{1c}^{})=\lambda _{1c}^4\delta \left(\lambda _{1c}\lambda _{1c}^{}\right)$$ (75) and has to be substituted into Eq. (74). The solution of Eq. (75) can be found easily because one of the solutions $`\mathrm{\Psi }_0\left(\lambda _{1c}\right)`$, Eq. (60), of the corresponding homogeneous equation is known. We can write immediately the second solution $`\mathrm{\Phi }_0\left(\lambda _{1c}\right)`$ in the form $$\mathrm{\Phi }_0\left(\lambda _{1c}\right)=\mathrm{\Psi }_0\left(\lambda _{1c}\right)_0^{\lambda _{1c}}\frac{\lambda ^2}{\mathrm{\Psi }_0^2\left(\lambda \right)}𝑑\lambda $$ (76) The function $`\mathrm{\Phi }_0\left(\lambda _{1c}\right)`$ grows exponentially at $`\lambda _{1c}\mathrm{}`$ and is proportional to $`\lambda _{1c}^3`$ for small $`\lambda _{1c}`$. Using the functions $`\mathrm{\Psi }_0\left(\lambda _{1c}\right)`$ and $`\mathrm{\Phi }_0\left(\lambda _{1c}\right)`$ we write the Green function $`g_0(\lambda _{1c},\lambda _{1c}^{})`$ as $$g_0(\lambda _{1c},\lambda _{1c}^{})=\{\begin{array}{cc}\mathrm{\Phi }_0\left(\lambda _{1c}\right)\mathrm{\Psi }_0\left(\lambda _{1c}^{}\right),& \lambda _{1c}\lambda _{1c}^{}\\ \mathrm{\Psi }_0\left(\lambda _{1c}\right)\mathrm{\Phi }_0\left(\lambda _{1c}^{}\right),& \lambda _{1c}\lambda _{1c}^{}\end{array}$$ (77) Substituting the function $`g_0`$, Eqs. (77, 60, 76), into Eq. (74) we see that in the region $`1\lambda _{1c},\lambda _{1c}^{}1/X`$ the integral is logarithmic and is safely cut from below and above by the limits of this region. It is important that we remain in the region of large $`\lambda _{1c}`$ where our approximations are valid. As a result of the calculation we obtain for the moments $`p_{\mathrm{}}^{\left(n\right)}\left(x\right)`$ of the density-density correlation function $$\frac{dp_{\mathrm{}}^{\left(n\right)}\left(x\right)}{dX}X\mathrm{ln}\left(1/X\right)\left\|x\right\|^{3/2}\mathrm{exp}\left(\left\|x\right\|/4L_{\mathrm{cu}}\right)$$ (78) It is interesting to note that the same prefactor determines correlations at coinciding points. Writing $`p_{\mathrm{}}^{\left(n\right)}\left(0\right)`$ as $`{\displaystyle \frac{dp_{\mathrm{}}^{\left(n\right)}\left(0\right)}{dX}}\left(i\stackrel{~}{\omega }\right)^{2n1}{\displaystyle \frac{\left(\lambda _{1d}\lambda _{1c}\right)^{2n}d\lambda _{1d}d\lambda _{1c}}{\lambda _{1d}^2\lambda _{1c}^2}}`$and integrating as previously over the region $`0<\lambda _{1d}\left(\lambda _{1c}\stackrel{~}{\omega }\right)^1`$, $`1\lambda _{1c}1/X`$ we come to the result $$\frac{dp_{\mathrm{}}^{\left(n\right)}\left(0\right)}{dX}X\mathrm{ln}\left(1/X\right)$$ (79) Eq. (78) shows that the derivative over the magnetic field of all moments of the density-density correlation function decays with the localization length $`L_{\mathrm{cu}}`$ of the unitary ensemble. The moments themselves can be obtained by integrating Eq. (78) over the magnetic field. The result can be written as $$p_{\mathrm{}}^{(n)}(x)\left(\frac{1}{x}\right)^{3/2}\left[\mathrm{exp}\left(\frac{\left\|x\right\|}{4L_{\mathrm{co}}}\right)+C_nX^2\mathrm{ln}\left(\frac{1}{X}\right)\mathrm{exp}\left(\frac{\left\|x\right\|}{4L_{\mathrm{cu}}}\right)\right].$$ (80) where $`C_n`$ is a coefficient depending on $`n`$. The first term in Eq. (80) is added assuming that at zero magnetic field one should have the standard result for the orthogonal ensemble. (In our previous work we wrote $`\mathrm{ln}^2\left(1/X\right)`$ instead of $`\mathrm{ln}\left(1/X\right)`$ which was a result of a not sufficiently accurate estimate). The second term in Eq. (80) becomes larger than the first one at distances exceeding $$x_HL_c\mathrm{ln}\left(1/X\right)$$ (81) and decays in the same way for all the moments. As we have discussed in the previous Section, this should mean that this term is determined by rare splashes of the wave functions. These splashes can be due to a hybridization of states with eigenenergies very close to each other and it is not surprising that they are sensitive to the magnetic field. A more delicate question concerns the distribution of the logarithm of the wave functions, governing the behavior of “typical” wave functions. This quantity is necessary for comparison with a recent numerical study, where the logarithm of the transmittance was shown to behave smoothly between the orthogonal and unitary ensembles. Our new state with the eigenenergy $`=\left(1+\rho ^2\right)/4`$ contributes to the distribution of the logarithms also. Unfortunately, the procedure developed for calculation of the moments cannot be extended to a calculation of the distribution function in the most interesting relation of the variables $`t\mathrm{exp}\left(\left\|x\right\|/L_{\mathrm{cu}}\right)`$ which would give the distribution of the localization lengths. The problem is that one needs to know the analytical properties of the Green function $`g_{k\rho }`$, Eq. (72), for arbitrary complex $`\rho `$. One can understand this on the example of the pure ensembles. In Eq. (52), computation of the distribution function required shifting the contour of the integration over $`\rho `$ to the saddle point $`\rho =i\mathrm{ln}t/ui`$. One could perform this procedure there because the integrand was known. In the crossover regime, one can rely only on the picture of the “free motion” applicable at small $`\rho `$, which make calculation of the distribution function not feasible. We can say only that the distribution function has an additional weight at $`t\mathrm{exp}\left(\left\|x\right\|/L_{\mathrm{cu}}\right)`$ so that the entire function is broader that for the pure ensembles. At the same time, the main body of the function is not known, which does not allow us to say much about “typical” wave functions. So, the results that we really have do not seem to contradict to the results of the numerical study of the average logarithms , where the second length was not seen. They suggest also that the information the Borland conjecture refers to, the behavior of the averaged logarithm of the wave function, may be not sufficient to characterize the crossover. Instead, the entire distribution function is of interest since it can differ from the standard logarithmically normal even for $`xL_c`$. As concerns a numerical observation of the second length $`L_{\mathrm{cu}}`$ in the moments of transmittances or density-density correlations, it may be very difficult because one needs to make computation for a very large number of configurations to be able to make a reliable statistics of the rare events. In addition, a smearing due to the presence of leads can make the observation even more difficult. ## VII Discussion of the Results We see from Eq. (80) that for weak magnetic fields the first term dominates at small distances. However, whatever weak the magnetic field is, the second term always prevails at sufficiently large distances. This is due to the fact that, although the small crossover parameter $`X`$ enters the pre-exponential factor, the exponential function containing the localization length $`L_{\mathrm{cu}}=2L_{\mathrm{co}}`$ decays much slower. Comparing these two terms, we obtained the characteristic distance $`x_H`$, Eq. (81) where the both terms are of the same order of magnitude To explain Eq. (80) qualitatively one can follow argument suggested by Mott for describing the AC conductivity with minor modifications necessary in our case. Actually, the arguments presented below apply for localized states in any dimensionality and allow us to obtain the scale $`x_H`$, Eq. (81), without any calculation. Naturally, the localization length $`L_c`$ in Eq. (81) depends on the dimensionality of the space. However, to be precise, we will speak in terms of the localization in wires. In quasi-one-dimensional samples all states are localized (in other words, $`_{\mathrm{}}^+\mathrm{}p_{\mathrm{}}^{(1)}(x)𝑑x=1`$). At the same time, for frequencies exceeding the energy $`E_c=D_0/L_c^2`$ the system behaves like a metal. Therefore, eigenenergies of states localized in a region of the wire of size $`L_c`$ are separated by the mean level spacing of the order of $`E_c`$. This picture allows us to estimate immediately the magnetic field $`H_c`$ at which the crossover between the orthogonal and unitary ensembles occurs. One should simply consider the region of the size $`L_c`$ as a closed grain and recall that the crossover between the orthogonal and unitary ensembles in a metallic grain occurs at flux $`\varphi `$ through the grain given by the relation $$E_T\left(\varphi /\varphi _0\right)^2\mathrm{\Delta }_G,$$ (82) where $`\varphi _0`$ is the flux quantum, $`\mathrm{\Delta }_G=\left(\nu V\right)^1`$ is the mean level spacing in the grain ($`V`$ is the volume), having the same order of magnitude as the Thouless energy $`E_T=D_0/L^2`$. For a quasi-one-dimensional grain of length $`L=L_c`$ one obtains the characteristic value $`X1`$, which means that the magnetic flux through the segment of the length $`L_c`$ is equal to the flux quantum $`\varphi _0`$ . At such fields the localization length of the main body of the wave functions changes considerably and this was clearly seen in the numerical work. At the same time, we know that the low frequency asymptotics of the level-level correlation function in an isolated grain changes from $`\omega `$ to $`\omega ^2`$ as soon as an arbitrarily weak magnetic field is applied. This corresponds to a finite probability of having two different levels at any small distance $`\omega `$. These levels hybridize by the magnetic field and this leads to the low frequency asymptotics $`\omega ^2`$ typical for the unitary ensemble. The two-scale localization considered in the previous sections is of a similar origin, namely, one can always find two localized states that are arbitrarily close to each other in energy, although they may be separated by a large distance in space. Using this fact Mott calculated the AC conductivity at finite frequencies $`\omega `$. We do not study the conductivity itself but rather different kinetic quantities like e.g. the density-density correlation function. In our case, the role of an external perturbation is taken over by the magnetic field. Following Mott’s arguments, let us try to find states whose energies slightly differ from the energy of a chosen state. As we have mentioned above, the typical separation in energy of states localized within a region of the length $`L_c`$ is of order $`E_c`$. In neighboring regions of the lengths $`L_c`$ one cannot find such states either: Even if two states were occasionally closely located in two initially isolated systems, they would hybridize and split after bringing these two segments of the wire in contact. The value of the splitting would be $`E_c`$ again. However, if two states are separated by a large distance $`xL_c`$, the splitting energy $`\mathrm{\Delta }_x`$ decays exponentially $$\mathrm{\Delta }_xE_c\mathrm{exp}(x/L_c)$$ (83) due to exponentially decaying overlap between the localized states. (Mott considered a tight-binding model; hence, the energy characterizing the splitting in his consideration was $`\mathrm{\Delta }_xI_0\mathrm{exp}(x/L_c)`$, where $`I_0`$ is of the order of band width.) Thus, one can find states $`a`$ and $`b`$ with an exponentially small energy difference considering states localized sufficiently far from each other. Now, turning on an external perturbation (magnetic field in our case) we want to mix these states, which, like in isolated granules, would lead to a behavior typical for the unitary ensemble (with the doubled localization length $`L_{\mathrm{cu}}=2L_{\mathrm{co}}`$). However, in order to mix the states we need a sufficiently large matrix element $`A_{ab}`$ of the vector potential between the states. If the states did not hybridize at all this matrix element would be exponentially small and this is similar the situation one encounters when calculating the AC conductivity at low frequencies. Mott suggested to take into account the states when they start hybridizing and we use this idea for our estimates. In this case the matrix element $`A_{ab}`$ is not exponentially small because the hybridized wave function is concentrated equally in the both localization centers. Using the modelling of the localization center in terms of an isolated granule of the length $`L_c`$ and applying the magnetic field $`H`$ we obtain easily a characteristic energy window $`E_H`$ of the order of $$E_HX^2E_c.$$ (84) This estimate follows from Eq. (16) by using the fact that in wires the energies $`E_c`$, $`E_T`$ and $`\mathrm{\Delta }_G`$ are of the same order. In an isolated grain states within this window are mixed and their correlations obey the statistics of the unitary ensemble. So, we come to the following picture: If two localized states $`a`$ and $`b`$ are separated by a distance $`x`$ such that the corresponding overlap integral $`\mathrm{\Delta }_x`$ exceeds $`E_H`$, the states are not considerably affected by the magnetic field and its influence can be considered perturbatively. If two localized states $`a`$ and $`b`$ are separated by a very large distance $`x`$ such that the overlap $`\mathrm{\Delta }_x`$ is much smaller than $`E_H`$, then the matrix $`A_{ab}`$ is exponentially small and the effect of the magnetic field on the states can be neglected. The magnetic field is important if the distance $`x`$ between the localized states $`a`$ and $`b`$ is such that $`\mathrm{\Delta }_xE_H`$. For such states the matrix element $`A_{ab}`$ is not small but, at the same time, the levels are close to each other. The magnetic field $`H`$ mixes the states $`a`$ and $`b`$ and one can expect for the hybridized states the statistics of the unitary ensemble. Comparing the energies $`\mathrm{\Delta }_x`$ and $`E_H`$ $$\mathrm{exp}(x/L_c)X^2,$$ (85) we obtain the characteristic distance $`x_H`$, Eq. (81). Assuming that the mixing of the levels by the magnetic field leads to the correlations given by the unitary ensemble we come to the conclusion that the exponential decay of the density-density correlation function and other correlation functions characterizing wave functions should be determined at distances $`xx_H`$ by the localization length $`L_{\mathrm{cu}}`$ of the unitary ensemble. These simple arguments explain the two-scale behavior, Eq. (80), and give the correct second scale $`x_H`$, although the doubling of the localization length can be obtained by the exact calculation only. At the same time, the probability to obtain states at a desired distance with eigenenergies very close to each other is small. This small probability may be compensated by a large contribution to moments of the density and therefore can be essential when calculating these quantities. As concerns calculation of the logarithm the large splashes are not crucial and their contribution is small due to the small probability, which is in agreement with the results of the numerical investigation of the transmittance. (actually, the authors of this work suggested as one of possible explanation of the difference between the averaged logarithm and the moments existence of “anomalously” localized states). An additional suppression of the splashes can be a consequence of a smearing due to inelastic scattering and presence of leads. To observe the splashes one has to be at distances from the ends of the wire exceeding $`x_H`$. At distances smaller that $`x_H`$ the weak magnetic field cannot have any considerable effect of the wave functions because the validity of the perturbation theory in $`H`$ is determined by the value of the smearing and not by the energy separation of the relevant states. But the transmittance probes the wave functions at the end of the sample. In this region, the magnetic field can become important only if the relevant magnetic energy $`E_H`$ becomes of the order of the smearing energy and the latter is of the order of $`E_c`$ near the ends. Recently, localization in disordered wires in a magnetic field was studied analyzing the autocorrelation function of spectral determinants (ASD) and no indication of the two-scale localization has been found. In our opinion, this cannot be a great surprise because the ASD is a quantity that can give only a rough picture of what is going in reality. By definition, the ASD probes only correlations of the energy levels and not the spatial structure of the wave functions. Any localized states contribute to the ASD independent of the distance between the centers of localization. This makes hardly possible to say anything about the overlap of the localized states, which is so important for the existence of the second scale. We think that the fact that the second scale was not identified in Ref. is not in contradiction with our finding but rather demonstrates that the ASD is not a proper quantity to describe the effect involved. Let us now discuss how one can try to check experimentally our predictions. Considering the transmittance of single wire one should speak about the logarithm of the transmittance that is a self-averaging quantity. In this case, the rare splashes we discussed previously are not important and one should see only one localization length which is at weak magnetic fields close to the length $`L_{\mathrm{co}}`$ of the orthogonal ensemble. However, one can measure the conductance of a system of a large number of wires connected in parallel. This is exactly the set up of the experiment of Ref. . The large number of wires in this case leads to averaging of the transmittance itself rather than of its logarithm. In such a situation one might see, in principle, the second scale $`L_{\mathrm{cu}}`$. Physically this means that although the splashes are rare, a wire, where such a splash occurs, gives the main contribution to the conductivity. Unfortunately, in the experiment , the localization length was not measured directly but was extracted from the dependence of the hopping conductivity on temperature. It is known that the conductance at very low temperatures is possible due to activations through Mott resonances, which leads to the picture of the variable range hopping. The localization length $`L_c`$ enters directly the activation energy in this conductivity $$\sigma =\sigma _0\mathrm{exp}\left[\left(T_0/T\right)^{1/(1+d)}\right]\text{}T_0\left(\nu L_c^d\right)^1.$$ (86) where $`d`$ is the dimensionality of the system. However, the localization length entering this formula is the length for a typical wave function and the splashes are not very important. As a matter of fact, Eq. (86) does not apply to 1D samples. The conductivity of 1D chains and wires should obey a simple Arrhenius-type law instead $$\sigma =\sigma _0\mathrm{exp}\left(T_0/T\right)\text{}T_0\left(\nu SL_c\right)^1.$$ (87) Remarkably, even accounting for electron-electron interactions preserves this kind of dependence. It is this law that was discovered in recent experiments in the regime of strong localization and activation energy $`T_0`$ was used to extract the localization length $`L_c`$. Which localization length enters the conductivity in the $`1D`$ case is not as clear as for higher dimensions and there could be a chance that the second scale can be seen here. Performing such experimental analysis at different temperatures one could try to observe the doubling of the localization length considered, because decreasing the temperature should result in a crossover from the activation behavior with the activation energy $`T_0`$ in Eq. (87), to the activation energy $`T_0/2`$. Indeed, at very low temperatures, $`T<\left(\nu Sx_H\right)^1`$, the electron hopping due to the overlap of far tails plays the dominant role, so that the larger localization length $`L_{\mathrm{cu}}`$ should be used in Eq. (87). For temperatures $`\left(\nu Sx_H\right)^1<T<\left(\nu SL_c\right)^1`$, hopping between the main body of localized states is essential and the localization length $`L_{\mathrm{co}}`$ should be substituted for $`L_c`$ in Eq. (87). Needless to say that the second scale may be observed only if the wires are sufficiently long such that the main contribution to the resistivity comes from the variable-range hopping inside the sample and the wire can be considered effectively as infinite. In this limit the resistivity should be proportional to the length of the wire. Measuring the conductance of a finite wire at ultralow temperatures where the resistivity grows exponentially with the sample length cannot help extracting the two-scale behavior as it has been discussed previously. ## VIII Conclusions Studying localization properties of infinitely long disordered wires within supersymmetry technique combined with the transfer-matrix technique, we managed to calculate correlation functions of interest. This procedure is based on working in the coordinate representation, Eq. (15), and generating any quantity of interest from the two-point correlation function, Eq. (7). Using this approach we obtained all moments and distribution functions of the density-density correlator and Landauer-like conductivity for the orthogonal and unitary ensembles, Eqs. (47,50,45). These expressions are exact and they are valid for arbitrary distances. Studying the crossover between the unitary and orthogonal ensembles for weak magnetic fields we have found moments of the density-density correlations at large distances $`xL_c`$. We have demonstrated that the far tail of these correlation functions characterizing the wave functions decays at arbitrary weak magnetic field with the localization length which is double as large as that of the main part. As all the moments decay with the same length, we argue that this behavior is due to rare splashes of the wave functions that are very sensitive to the magnetic field. In contrast to this, the localization length characterizing a typical wave function might exhibit a smooth crossover. We emphasize, however, that in the crossover regime, not only study of the typical localization length, but also of the whole distribution of the logarithm is of interest. It should differ from the standard log-normal one even for $`xL_\mathrm{c}`$ so long as $`xx_H`$, Eq. (81), being asymmetric and broad. At present, we can not predict how the entire distribution function of the averaged logarithm of the wave function will develop in the limit of infinitely long distances. Slightly revising the Mott arguments for the AC conductance, we have shown how one can extract the characteristic scale of the problem $`x_H`$, without performing any calculations. These arguments are quite general, which implies that our results may be valid for a disordered system of any dimension. From the exact treatment as well as from the qualitative arguments of the previous section we understand that the effect we have found is very sensitive to different kinds of the level smearing. Therefore, it may be observed at low temperatures only. As concerns a possibility of an experimental observation of the effect it might be seen at low temperatures in the regime of the hopping conductivity in a system of wires connected in parallel. The wires should be sufficiently long, such that the system would obey Ohm’s law. Then, if the tails of the wave functions determine the hopping conductivity, one should expect the Arrhenius law but with different activation energies in the limit of very low temperatures. We hope that it can become possible to check our finding both numerically and experimentally, although this can be a difficult task due to a small probability of the wave functions with the high sensitivity to the magnetic field. ## IX Acknowledgments This work was supported by SFB 237 “Unordnung und Große Fluktuationen”. A.V.K. acknowledges hospitality of the Instituut-Lorentz at the Leiden University. We are indebted to I.L. Aleiner, B.L. Altshuler, C.W.J. Beenakker, M. Janssen, and H. Schomerus for useful discussions.
warning/0005/astro-ph0005288.html	ar5iv	text	# 1 Introduction ## 1 Introduction Planetary Nebulae (PNe) and proto-Planetary Nebulae (PPNe) are believed to be the penultimate evolutionary stage of low and intermediate mass stars (zero age main sequence mass range $`0.8M_{}<M_{zams}5M_{}`$). PNe appear on the sky as expanding ionized plasma clouds surrounding a hot central star. As the resolution of ground based telescopes increased, the “typical” shape of a PN went from spherical (Osterbrock 1972) to elliptical and/or bipolar (Balick 1987). More recently, deeper and higher resolution studies (particularly those with the HST) have shown many PNe with narrow collimated features that are better described as jets than bipolar lobes. In some cases the jets appear as extremely well collimated bipolar outflows: M2-9 (Balick 2000): Hen-401 (Sahai 2000), AFGL 2688 (The Egg) (Sahai et al.1998). In other cases they appear as FLIERS (Fast Low-Ionization Emission Regions), a single knot expanding away from a (usually) elliptical nebula. In addition to the presence of jets it also appears that many PNe show evidence for multi-shell structures or multi-polarity. In these objects two nested pairs of bipolar lobes appear. In multi-shell structures the lobes are aligned along the same axis of symmetry (Hb 12: Welch et al.1999) while multi-polar cases show different orientations for the inner and outer lobes (M2-46: Manchado et al.2000). While considerable progress has been made in understanding the hydrodynamics or magneto-hydrodynamics of shaping bipolar lobes (see Frank 1999 for a review), the origin of jets, point symmetry and multi-polar outflows in PNe poses fundamental challenges for theory. The basic hydrodynamics of bipolar outflows can be described by the classic Generalized Interacting Stellar Winds (GISW) model (Kwok et al.1978, Kahn & West 1986, Balick 1987, Icke 1988). In this scenario a slow ($`10`$ km s<sup>-1</sup>), dense ($`10^4`$ M yr<sup>-1</sup>) toroidal wind expelled during the AGB is followed by a fast ($`1000`$ km s<sup>-1</sup>), tenuous ($`10^7`$ M yr<sup>-1</sup>) wind driven off the contracting proto-white dwarf during the PNe phase. The GISW model can explain bipolar morphologies and even jets. Point-symmetry and multi-polar bubbles however, do not fall easily into line with the concept of a large-scale collimating torus. Models invoking a toroidal magnetic field embedded in a normal radiation driven stellar wind have recently shown some promise. This so-called Magnetized Wind Bubble (MWB) model was first proposed by (Chevalier & Luo 1994) and has been studied numerically by (Rozyczka & Franco 1996) and (García-Segura et al.1999). In these models the field at the star is dipolar but assumes a toroidal topology due to rapid stellar rotation. Collimation is not activated until the wind passes through the inner shock. Then hoop stresses associated with the toroidal field dominate over isotropic gas pressure forces and material is drawn towards the axis producing a collimated flow. This mechanism has been shown capable of producing a wide variety of outflow morphologies including well collimated jets. When precession of the magnetic axis is included in fully 3-D simulations, the MWB model is capable of recovering point-symmetric morphologies as well (García-Segura 1997). It is noteworthy that the wind in these models differs substantially from standard MHD driven winds in disk and stellar wind contexts, and the fundamental issue of collimation before the wind-wind interaction remains to be addressed. In addition it is not clear that such models can be applied to multi-polar bubbles due to the difference in collimation orientations. Finally neither the GISW or MWB models fully address the origin of the wind. This becomes a critical issue when one is dealing with Proto-Planetary Nebulae (PPNe). There is growing evidence suggesting that the PN shaping process begins in the PPN phase via well collimated jets (Sahai & Trauger 1998). The creation of fast ($`V>100`$ km/s) high mass loss rate ($`\dot{M}>10^6`$ M yr<sup>-1</sup>) winds during these early phases when the star is an F or G type is not easily explained via classic line driven wind theory as the available power is insufficient (Alcolea et al. 2000). Thus for PPN, theorists face a dilemma similar to the problem of YSO outflows in that they must explain both wind acceleration and collimation (Pudritz 1991, Königl & Ruden 1993, Shu et al.1994). The above considerations motivate further study of magnetized wind models for the launching and collimation of outflows from PN. In particular, the basic question of the available MHD wind power is important to address. Magnetically driven outflows are a leading paradigm for a ubiquity of source classes in nature. This includes such contexts as protostellar jets (c.f. Smith 1998; Frank 1999 for reviews), the Solar wind (c.f. Parker 1979), jets of Galactic accretors (c.f. Mirabel and Rodriguez 1999), active galactic nuclei (c.f. Ferrari 1998; Blandford 2000 for reviews), and gamma-ray bursts (e.g. Usov 1992; Duncan & Thomspon 1992; Blackman et al. 1996; Meszaros & Rees 1997; Ruderman et al. 2000 ). While the parameter regimes are varied and the extent and nature of collimation, particle acceleration, and radiation process are different in the different settings, an underlying principle remains: magnetic fields can act as a drive belt between gravity and energy deposition at large radii, extracting the rotational energy of the rotator into an outflow. We note that the potential for accretion disks to exist in PPN systems has been raised (Morris 1987; Soker & Livio 1994; Reyes-Ruiz & Lopez 1999). In the case of PN, as in other star-disk or compact object-disk systems, the magnetically driven outflows can in principle emanate from both the central object or the disk. Because the relative collimation and power evolution of the disk and stellar winds can be different, it may not be unreasonable to see outflows of different character coexisting in the same object. In order to investigate this properly however, it is necessary to allow that the disk and star are coupled and this requires a model for the disk evolution. A careful investigation of the relevant disk formation is presented in Reyes-Ruiz & Lopez (1999). The disk in this study forms when a binary system undergoes common envelope evolution. The ejection of the envelope involves transfer of angular momentum from the secondary to the envelope after which the secondary loses enough angular momentum to fall to a separation such that it can fill its Roche lobe and form the disk. Reyes-Ruiz & Lopez (1999) discuss a number of important constraints on the constituents and properties of binaries which lead to disk formation. The most likely system turns out to be an evolved AGB star with mass $`2.5<M_{agb}/M_{}<5`$ for the primary, with a secondary of mass $`\text{ }<0.08M_{}`$ and with an initial binary separation of $`>50100R_{}`$ (Iben 1991). The AGB star will shed 80% of its mass during the common envelope ejection, leaving a post AGB stellar core of Mass $`M_c`$ surrounded by a thin residual convective shell. The time scale for the disk to move to its inner inner radius is of order or shorter than the viscous time scale. Even for systems whose initial outer radii are of order $`100R_{}`$, the relevant time scale would be no more than a few years. Since the disk would form only after the period of common envelope ejection, $`<`$ 1yr, we would expect the disk and the stellar shell to form concurrently within a time of order a few years. Because the relevant stellar field will be exposed once the envelope ejection occurs, and because any stellar dynamo growth time in the disk (discussed later) is less than 1 year, we suggest the rotating stellar shell and disk could both support co-operative jets at a time coincident within about a year after the disk forms. This represents the initial time appearing in our calculations. In section 2, we determine the MHD wind power from PPN stars and disks. This requires a calculation of the expected field strength in the disk and a model of how the accretion affects the stellar field, the stellar spin, and the disk inner radius. We study the cases for which the stellar shell is initially rapidly or slowly rotating, and calculate the maximum disk and stellar magnetic wind powers for each case. In section 3 we discuss the observational implications of the results of the power calculation and speculate on the magnetic shaping. We summarize and conclude in section 4. ## 2 Winds from Disks and Stars There are many models of magnetic jet production and collimation in the literature. The extent and explicit details of collimation are perhaps less agreed upon than the principles of launching. Magnetic launching mechanisms can be divided into two basic classes, “spring” mechanisms (Uchida & Shibata 1985; Contopoulos 1995 Lynden-Bell 1996) and “fling” mechanisms (Blandford & Payne 1982; Lovelace et al. 1987). In the former class of models, the initial driving force is a magnetic spring, i.e. the magnetic field energy density is of order the kinetic energy density in the disk and the outflow is driven by toroidal field pressure. In the fling mechanism, the initial driving is centrifugal along poloidal field lines. Higher up in the corona the wind becomes magnetically driven. In this study we compute the total available magnetic luminosity $`L_m`$ from the disk and central star. The magnetic luminosity provides an upper limit to the associated wind kinetic luminosity. The $`L_m`$ is the integral of the Poynting flux over the surface area of the rotator. We thus have $`L_w\dot{M}_wV_w^2ϵL_m=ϵ{\displaystyle (𝐄\times 𝐁)𝑑𝐒}ϵ{\displaystyle _{R_i}^{R_o}}(\mathrm{\Omega }R)B_pB_\varphi R𝑑R,`$ (1) where $`ϵ`$ is an efficiency, and $`B_\varphi `$ and $`B_p`$ are the toroidal and poloidal field respectively, $`\mathrm{\Omega }`$ and $`R`$ are the rotational frequency and radius of the magnetized wind source, $`\dot{M}_m`$ is the wind outflow rate, and $`V_w`$ is the wind outflow speed. In what follows we discuss this equation in the context of PPN disks and post AGB stars. For all cases of interest, for disks the integrand falls off fast enough such that the inner radius matters most. Thus for disks we have $`L_{dw}ϵ_dL_{dm}ϵ_dB_p(R_i)B_\varphi (R_i)\mathrm{\Omega }_d(R_i)R_i^3,`$ (2) where $`\mathrm{\Omega }_d`$ is the disk angular speed $`L_{dm}`$ is the available magnetic luminosity associated with the disk, and $`ϵ_d`$ is the unknown efficiency factor. For stars, we have $`L_{sw}ϵ_sL_{sm}B_p(R_{})B_\varphi (R_{})\mathrm{\Omega }_{}(R_{})R_{}^3,`$ (3) where $`\mathrm{\Omega }_{}`$ is the angular speed of the stellar surface in which the field lines are anchored, $`L_{sm}`$ is the available magnetic luminosity associated with the star, and $`ϵ_s`$ is the stellar wind efficiency factor. The above disk and stellar wind luminosities are not independent. The poloidal and toroidal fields are coupled through differential rotation in both the disk and star, and the angular velocity of the star may depend on the accretion. In addition, the inner radius of the disk depends on the stellar magnetic field. Thus (2) and (3) are very much interdependent. We now proceed to solve for the coupling relations and consider the stellar and disk MHD wind luminosities for two cases: * The star is initially a strongly magnetized rotator independent of the accretion. * The star is initially a slow rotator but is spun up by the disk accretion. As we will see, the evolution in both cases depends on estimates of the disk and star magnetic fields. We must also solve for the disk inner radius, and the stellar surface angular speed evolution. ### 2.1 Magnetic Field of Disk If the jet/wind is due to a magnetic field in the disk, this field must be produced in situ for the PN case. This is because the companion star, the progenitor of the disk, would have to be a brown dwarf or a Jupiter type planet (Reyes-Ruiz & Lopez 1999) which typically would only have fields of order or less than 10G. Even a field several orders of magnitude larger would be far too small to have any influence on the disk, which forms from shredding this mass, if the disk field were only that resulting from flux freezing. We estimate the field strength roughly from dynamo or turbulent amplification in the disk. An alternative dynamo approach to ours is given by Reyes-Ruiz & Stepinski (1995). For an accretion disk whose angular momentum is transported by a magneto-shearing instability (c.f. Balbus-Hawley 1991, 1998), the time scale for growth of unstable modes and the time scale for the largest eddy turnover time are both approximately equal to the rotation time scale. Very roughly, the eddy turnover time can be written as $`L/v_T`$ where $`v_T`$ is the dominant turbulent velocity and $`L`$ is the dominant correlation scale. We can estimate the strength of the mean magnetic field by first estimating the total magnetic energy. Using the Shakura-Sunyaev (1973) viscosity prescription $$\nu =\alpha _{ss}c_sHv_T^2/\mathrm{\Omega },$$ (4) where $`\alpha _{ss},c_s,H`$ and $`L/v_T`$ are the disk viscosity parameter, the sound speed, the scale height and the dominant eddy turnover time respectively. Using the fact that turbulent stretching leads to $`v_Av_T`$, where $`v_A`$ is the Alfvén speed (ignoring a possible factor of $`\sqrt{2}`$ on $`v_T`$, e.g. Balbus & Hawley 1998) , and $`\mathrm{\Omega }R=c_sR/H`$ for a thin accretion disk, we then have straight away $$v_A^2=\alpha _{ss}c_s^2.$$ (5) When modeled as a mean-field $`\alpha _d\mathrm{\Omega }`$ dynamo (c.f. Parker 1979; Reyes-Ruiz & Stepinski 1995; Blackman 2000), it can be shown that due to the differential shear, the mean toroidal field exceeds the saturated mean poloidal field strength in the disk by a factor of $`H/L`$. From (4) and (5) we have $`L\alpha _{ss}^{1/2}H`$ and thus $`B_p\alpha _{ss}^{1/2}B_\varphi .`$ (6) Whether this relation holds true in the base of the disk corona where the jet would launch is unclear, but since the field in the corona is no greater than the field in the disk, the overall disk mean field provides a good upper limit at least to the total mean field energy density. Note that the total mean stress $`B_\varphi B_p=\overline{B}_\varphi \overline{B}_p+b_\varphi b_p`$, where the second term is a correlation of fluctuating fields. Actually, both of these terms can drive a wind. Exploring this point is outside the scope of the present paper, and here we simply stick with the standard approach (e.g. Blandford & Payne 1982) where the large scale fields drive the wind. Thus we ignore the $`b_\varphi b_p`$ term. The upper limit of the mean field is of order the random field, so from (5) we have, for the dominant disk field component, $`B_\varphi ^24\pi \rho _d\alpha _{ss}c_s^2.`$ (7) From the mass continuity equation of the disk we have, $`\rho _d=\dot{M}_a/(4v_R\pi HR),`$ (8) where $`v_R`$ is the radial infall speed. For (thin or ADAF) disks $`v_R\alpha _{ss}c_sH/R\alpha _{ss}v_k(H/R)^2,`$ (9) where $`v_k^2=GM/R`$, the Keplerian speed. Then using (7) and (8) and (9), we have $`B_\varphi B_p=\alpha _{ss}^{1/2}B_\varphi ^2={\displaystyle \frac{\alpha _{ss}^{1/2}(R/H)v_k(R)\dot{M}_a}{R^2}}.`$ (10) This expression can be used to estimate either the magnetic field strength or Maxwell stress avialable in the disk for launching a wind. ### 2.2 Magnetic field of the stellar shell The relevant stellar field is the field at the anchoring radius of the post AGB star from which the wind emanates. Observational interpretation suggests that the common envelope ejection, which would precede disk formation, removes $`80`$% of an initially $`3M_{}`$ star (Schönberner 1993). Convective stellar models for a 3$`M_{}`$ star (Kawalar, private communication 2000) then tell us that outer layer will have a convective shell containing a mass of order $`0.01M_{}`$. We allow for a distinction between the core of the star $`M_c`$ and the mass in which the field is anchored, $`M_{}`$. The latter is likely to be equal to the core mass if the field is generated by an AGB interface dynamo (Blackman et al. 2000), but if the field were somehow anchored only in the thin shell, $`M_{}`$ might be $`<<M_c`$. Since the convective shell is likely to be differentially rotating with respect to the core we assume that the shear rate is of order the anchoring material’s angular speed, $`\mathrm{\Omega }_{}`$. The shear will stretch a dynamo produced poloidal field linearly in time. The linear streching could operate coherently over a vertical diffusion time $`\tau _DR_{}^2/\eta _T`$, where $`\eta _TLv_T`$ is the turbulent diffision coefficient for turbulent velocity $`v_T`$ and scale $`L`$. (This would need to be calculated for specific interface dynamo models, c.f. Parker 1993; Markiel & Thomas 1999; Blackman et al 2000). Using numbers from Kawaler (private communication 2000), the diffusion time could be of order a few times the convective overturn time for the largest eddies in the post-AGB star, about $`0.05`$yr. In any case, for linear growth of the toroidal field, the toroidal and poloidal fields are then related by $`B_\varphi B_p\mathrm{\Omega }_{}\tau _D,`$ (11) which follows from the magnetic induction equation. Although the poloidal field $`B_p`$ will likely also depend on $`\mathrm{\Omega }_{}`$, it will also depend on other properties of the turbulence and the mechanism of field origin (c.f. Parker 1979, 1993; Pascoli 1997; Soker 1998; Markiel & Thomas 1999; Blackman et al. 2000). For an initial treatment of the problem, we simply consider $`B_p`$ as an input condition. ### 2.3 Disk inner radius Now that we have the stellar magnetic field, we can calculate the disk inner radius. Related calculations were performed by Ghosh & Lamb (1978) and Ostriker & Shu (1995). We estimate the inner radius by balancing the infall ram pressure from the disk with the magnetic pressure of the star. For the infall velocity we use the free fall velocity, as that corresponds to the field strength at which the disk must absolutely truncate. Thus our calculation represents a lower limit on the inner radius. Note for example that when the Alfvén speed in the disk, using the stellar field strength, equals the sound speed there, magneto-shearing instabilities will be shut off. At this field strength the nature of the angular momentum transport must change (to e.g. global modes, B.Chandran, E.Ostriker, 2000 personal communication). We assume that a suitable change occurs and thus still use the larger Keplerian speed for the balance below. Using the $`r^3`$ dependence for a dipole field and considering only the values at the equator $`r=R`$, we have $`{\displaystyle \frac{1}{4\pi }}B_{}^2\left({\displaystyle \frac{R_{}}{R_i}}\right)^6\rho _{d,i}v_{k,i}^2,`$ (12) where $`v_{k,i}`$ is the Keplerian ($``$ free-fall) speed at $`R_i`$. Note that when $`\mathrm{\Omega }_{}\tau _D\text{ }>1`$, $`B_{}^2B_\varphi ^2`$, otherwise $`B_{}^2B_p^2`$. This is important because the angular speed couples in the former case but not in the latter. The next step is to use the mass continuity equation (8) for $`R_i`$. The result is $`R_i=B_p^{4/7}(1+\mathrm{\Omega }_{}^2\tau _D^2)^{2/7}\alpha _{ss}^{2/7}(H_i/R_i)^{6/7}R_{}^{12/7}(GM_c)^{1/7}\dot{M}_a^{2/7}.`$ (13) The inner radius will be time dependent because of the time dependent accretion rate and angular speed of the anchoring material. We now compute the evolution of this angular speed. ### 2.4 Angular Speed of the Star The disk and stellar wind luminosities depend on the angular velocity $`\mathrm{\Omega }_{}`$ of the stellar material where the field is anchored. As shown above, $`\mathrm{\Omega }_{}`$ is one factor which determines the magnetic pressure of the star, and thus $`R_i`$. For small $`\tau _D`$, $`\mathrm{\Omega }_{}`$ decouples from $`R_i`$. If material is accreting onto the star from the disk at rates based on the disk formation models of Reyes-Ruiz and Lopez (1999), then $`\dot{M}_a6\times 10^{22}\left({\displaystyle \frac{t}{1\mathrm{y}\mathrm{r}}}\right)^{5/4}\mathrm{g}/\mathrm{s},`$ (14) and there may be some contribution to the rotation from the accreted material. There are many complications in this accretion process. Here we simply assume that the material imparts its angular momentum to the shell (e.g. via the magnetic field) as it accretes. We also assume that the accreted material does not change the thickness of the star. The rotation speed of the relevant stellar envelope is then determined by balancing the angular momentum gained from the accretion with that lost from the MHD stellar wind. We have the following approximate equation $`M_{}R_{}^2{\displaystyle \frac{d\mathrm{\Omega }_{}}{dt}}=\dot{M}_a\mathrm{\Omega }_{}R_{}^2+\dot{M}_a\mathrm{\Omega }_{k,i}R_i^2B_p^2\mathrm{\Omega }_{}R_{}^3\tau _D.`$ (15) where $`\mathrm{\Omega }_{k,i}`$ is the Keplerian angular speed at the disk inner edge, and we have assumed a constant $`R_{}`$. The last term on the right hand side is due to the magnetic torque applied by the wind ($`B_pB_\varphi `$). This equation applies when $`R_i`$ corresponds to a radius at which the disk is spinning faster than field lines of the star at that radius. Thus it is assumed that the disk accretion can spin up the stellar convective envelope. There exists a regime in which the star spins fast enough to disrupt the disk, and angular momentum of the star is lost to the disk. We do not consider that regime. Note further that if $`R_i`$ were greater than the radius at which the co-rotation speed exceeded the Keplerian speed of the star, there would be no accretion (Armitage & Clarke 1996). Re-arranging, and expanding the Keplerian term, we have $`{\displaystyle \frac{d\mathrm{\Omega }_{}}{dt}}+\mathrm{\Omega }_{}(\dot{M}_a/M_{}+B_p^2R_{}\tau _D/M_{}){\displaystyle \frac{(GM_c)^{1/2}\dot{M}_aR_i^{1/2}}{R_{}^2M_{}}}.`$ (16) In the simple limit of no disk, only the 1st and 3rd terms of (16) would contribute and the solution would be $`\mathrm{\Omega }_{}=\mathrm{\Omega }_0e^{t/t_m},`$ (17) where we define the magnetic spin-down time scale $`t_m`$ $`t_m{\displaystyle \frac{M_{}}{B_p^2R_{}\tau _D}}.`$ (18) This quantity is the spin-down time from MHD powered rotational energy loss. For typical parameters, $`t_m50\left({\displaystyle \frac{M_{}}{M_{}}}\right)\left({\displaystyle \frac{B_p}{10\mathrm{k}\mathrm{G}}}\right)^2\left({\displaystyle \frac{R_{}}{10^{11}\mathrm{cm}}}\right)^1\left({\displaystyle \frac{\tau _D}{10^5\mathrm{sec}}}\right)^1\mathrm{yr}`$ (19) As we shall see there are two regimes of interest to pursue in analytic approximation depending on whether $`t/t_m`$ is large or small. For the kind of disks considered here, the second term in (16) can be ignored compared with the third. Also, for the maximum accretion rates we consider $`(i.e.10^3M_{}\mathrm{yr})(t/1\mathrm{y}\mathrm{r})^{5/4}`$, a mass shell of 0.01 $`M_{}`$ will at most have a 30% gain in mass over the lifetime, because of the form of the time dependence, so we can assume $`M_{}`$ is a constant after the 1 year initial system formation period. The ODE can then be solved by standard integrating factor techniques, and is in fact analogous to an RL circuit. We find $`\mathrm{\Omega }_{}=CExp\left[{\displaystyle ^t}{\displaystyle \frac{dt^{}}{t_m}}\right]+Exp\left[{\displaystyle ^t}{\displaystyle \frac{dt^{}}{t_m}}\right]\times {\displaystyle ^t}Exp\left[{\displaystyle ^t}{\displaystyle \frac{dt}{t_m}}\right]q(t^{})𝑑t^{}`$ (20) where we define $`q(t){\displaystyle \frac{(GM_c)^{1/2}\dot{M}_a(t)R_i^{1/2}(t)}{R_{}^2M_{}}}.`$ (21) Note again that we do not consider the solution before $`t=1`$yr when the accretion disk is still forming. For the approximate early time solution, $`t/1\mathrm{y}\mathrm{r}<t/t_m<1`$, we set the exponentials equal to 1. The result (after the 1 yr disk formation grace period) is then $`\mathrm{\Omega }_{}(t)\mathrm{\Omega }_0+\omega _0\left(1\left(\frac{t}{1\mathrm{y}\mathrm{r}}\right)^{1/14}\right),`$ (22) where $`\omega _o={\displaystyle \frac{3\times 10^8(GM_c)^{1/2}\dot{M}_{a0}R_{i0}^{1/2}}{R_{}^2M_{}}},`$ (23) and where the subscript $`0`$ indicates the quantity at its initial time $`(t=1\mathrm{y}\mathrm{r})`$. We have ignored the time dependence in $`R_i`$ from $`\mathrm{\Omega }_{}`$, while including the implicit time dependence from $`\dot{M}_a`$. This is a good approximation for this regime because $`R_i^{1/2}`$ depends on $`\mathrm{\Omega }_{}^{2/7}\dot{M}_a^{1/7}`$. Note that expression (22) allows us to define the difference between initially fast ($`\mathrm{\Omega }_0>\omega _0`$) and slow ($`\mathrm{\Omega }_0<\omega _0`$) stellar rotation. When $`t/t_m>1`$ the exponential in (20) begins to evolve quickly compared to the power law decay of the $`\dot{M}_aR_i^{1/2}`$ factor inside the last integral. Factoring this out and noting that the product of integrals in this last term then cancel, we have $`\mathrm{\Omega }_{}(t)q(t)t_m+\left(Kq(t)t_m\right)e^{(t/t_m)},`$ (24) where $`K=e\mathrm{\Omega }_{}(t_m)+q(t_m)t_m(1e)`$. For large times $`(t/t_m)>1`$, the first term on the right of (24) dominates and we have $`\mathrm{\Omega }_{}(t)q(t)t_m={\displaystyle \frac{(GM_c)^{1/2}\dot{M}_aR_i^{1/2}}{B_p^2R_{}^3\tau _D}}=\alpha _{ss}^{1/7}{\displaystyle \frac{(GM_c)^{3/7}\dot{M}_a^{6/7}(H/R)^{3/7}}{B_p^{12/7}R_{}^{15/7}\tau _D}}(\mathrm{for}\mathrm{\Omega }_{}\tau _D<1)`$ $`\alpha _{ss}^{1/5}{\displaystyle \frac{(GM_c)^{3/5}\dot{M}_a^{6/5}(H/R)^{3/5}}{B_p^{12/5}R_{}^3\tau _D}}(\mathrm{for}\mathrm{\Omega }_{}\tau _D>1),`$ (25) where the latter represents the case for which $`R_i`$ depends on $`\mathrm{\Omega }_{}`$ from (13). Note that the equations of (25) should be used only when it produces a value below the corresponding escape speed. ### 2.5 Disk Wind Luminosities Plugging (10) into (2), we have $`L_{dw}=\dot{M}_{dw}V_{dw}^2=ϵ_d\alpha _{ss}^{1/2}{\displaystyle B_\varphi ^2\mathrm{\Omega }_kR^2𝑑R}ϵ_d\alpha _{ss}^{1/2}\dot{M}_a(R_i/H_i)GM_c/R_i,`$ (26) where $`\dot{M}_{dw}`$ indicates the outflow rate from the disk wind and $`V_{dw}`$ is the disk wind outflow speed. We now consider several cases which involve using the appropriate value of $`R_i`$ from section 2.3 in (26), then using the appropriate solution for $`\mathrm{\Omega }_{}`$ from section 2.4 in (13) for $`R_i`$. We will also use (14) for $`\dot{M}_a`$. We first consider the case of slow stellar rotation ($`\mathrm{\Omega }_0<\omega _0`$). The maximally luminous case will have the smallest disk inner radius. For a given poloidal field, the toroidal stellar field will be weakest for slowly rotating stars as per our discussion above, so in this case the disk inner radius will be the smallest. When $`R_i`$ as calculated in (13) satisfies $`R_iR_{}`$, we use $`R_i=R_{}`$ and the disk luminosity is given by $`L_{dw}{\displaystyle \frac{ϵ_d\alpha _{ss}^{1/2}GM_c\dot{M}_a}{R_{}}}\left({\displaystyle \frac{R_{}}{H_i}}\right)10^{37}{\displaystyle \frac{ϵ_d\alpha _{ss,2}^{1/2}\dot{M}_{a0,22}M_{c,1}}{R_{,11}}}\left({\displaystyle \frac{t}{1\mathrm{y}\mathrm{r}}}\right)^{5/4}\left({\displaystyle \frac{R_{}/H_i}{10}}\right)\mathrm{erg}/\mathrm{s}.`$ (27) In this, and the expressions which follow, we scale $`\alpha _{ss}`$ to $`0.01`$, magnetic fields to $`10^4`$ G, radii to $`10^{11}`$ cm, diffusion times to $`10^5`$ yr, stellar core masses to $`1M_{}`$, stellar shell masses to $`10^2M_{}`$, angular speeds to $`10^4Hz`$ and mass loss rates to $`10^{22}\mathrm{g}\mathrm{s}^1`$. These scalings are indicated by the subscripts. We now consider the case of an initially rapidly spinning star. The disk and star are strongly coupled through $`R_i`$ and we must consider the different temporal regimes relative to the ratios $`t/t_m`$ and $`\mathrm{\Omega }_{}\tau _D`$. When $`t<t_m`$ the first term on the right hand side of (22) dominates. Then using (13) for $`\mathrm{\Omega }_{}\tau _D>1`$, we have $`L_{dw}{\displaystyle \frac{ϵ_d\alpha _{ss}^{\frac{3}{14}}(GM_c)^{\frac{8}{7}}\dot{M}_a^{\frac{9}{7}}}{B_p^{\frac{4}{7}}\tau _D^{4/7}\mathrm{\Omega }_{}^{\frac{4}{7}}R_{}^{\frac{12}{7}}}}\left({\displaystyle \frac{R_i}{H_i}}\right)^{13/7}7\times 10^{37}{\displaystyle \frac{ϵ_d\alpha _{ss,2}^{\frac{3}{14}}M_{c,1}^{\frac{8}{7}}\dot{M}_{a0,22}^{\frac{9}{7}}}{\left(B_{p,4}\tau _{D,5}\mathrm{\Omega }_{0,4}\right)^{\frac{4}{7}}R_{,11}^{\frac{12}{7}}}}\left({\displaystyle \frac{t}{1\mathrm{y}\mathrm{r}}}\right)^{\frac{45}{28}}\left({\displaystyle \frac{R_i/H_i}{10}}\right)^{\frac{13}{7}}{\displaystyle \frac{\mathrm{erg}}{\mathrm{s}}}.`$ (28) For $`t>t_m`$, using (13) (under the assumption that $`B_\varphi \text{ }>B_p`$, or equivalently $`\mathrm{\Omega }_{}\tau _D>1`$), and (25) in combination with (26), we have $`L_{dw}ϵ_d\alpha _{ss}^{\frac{1}{70}}(GM_c)^{\frac{4}{5}}\dot{M}_a^{\frac{3}{5}}B_p^{\frac{4}{5}}\left({\displaystyle \frac{R_i}{H_i}}\right)^{\frac{11}{5}}10^{38}ϵ_d\alpha _{ss,2}^{\frac{1}{70}}M_{c,1}^{\frac{4}{5}}\dot{M}_{a0,22}^{\frac{3}{5}}B_{p,4}^{\frac{4}{5}}\left({\displaystyle \frac{t}{50\mathrm{y}\mathrm{r}}}\right)^{\frac{3}{4}}\left({\displaystyle \frac{R_i/H_i}{10}}\right)^{\frac{11}{5}}{\displaystyle \frac{\mathrm{erg}}{\mathrm{s}}}`$ (29) In the limit that $`\mathrm{\Omega }_{}\tau _D<1`$ or $`B_p>B_\varphi `$, from (13) we obtain instead, $`L_{dw}{\displaystyle \frac{ϵ_d\alpha _{ss}^{\frac{1}{14}}(GM_c)^{\frac{8}{7}}\dot{M}_a^{\frac{9}{7}}}{B_p^{\frac{4}{7}}R_{}^{\frac{12}{7}}}}\left({\displaystyle \frac{R_i}{H_i}}\right)^{\frac{13}{7}}7\times 10^{35}{\displaystyle \frac{ϵ_d\alpha _{ss,2}^{\frac{1}{14}}M_{c,1}^{\frac{8}{7}}\dot{M}_{a0,22}^{\frac{9}{7}}}{B_{p,4}^{\frac{4}{7}}R_{,11}^{\frac{12}{7}}}}\left({\displaystyle \frac{t}{50\mathrm{y}\mathrm{r}}}\right)^{\frac{45}{28}}\left({\displaystyle \frac{R_i/H_i}{10}}\right)^{\frac{13}{7}}{\displaystyle \frac{\mathrm{erg}}{\mathrm{s}}}.`$ (30) For the characteristic numbers used, (29) does not apply because a check of $`\mathrm{\Omega }_{}\tau _D`$ from (25) reveals that $`\mathrm{\Omega }_{}\tau _D<1`$ for $`t>t_m`$ for our parameter regime. The above formalism is valid when the magnetic luminosity is $`\text{ }<`$ the accretion luminosity $`G\dot{M}_{}M_c/R_i`$. If the magnetic luminosity exceeds the disk luminosity, then the disk structure would be disrupted and the formalism would have to be revisited. ### 2.6 Stellar wind luminosities The stellar wind luminosity is given by (3). We note again the hidden dependence of the field on the rotational velocity. If a dynamo generates the stellar field, then both $`B_p`$ and $`B_\varphi `$ can depend on $`\mathrm{\Omega }_{}`$. For example the rotation would be required to generate the pseudoscalar helicity that sustains a steady field in dynamo models. Also, the toroidal field may induce poloidal field by springing outward. More subtle dependences should be considered in future work, but the important point here is the recognition of some dependence on the rotation, and thus on the supply of angular momentum from the disk. Here we simply assume that the poloidal field of the star is kept at a constant value either by a dynamo or by buoyancy of toroidal field, while the toroidal field depends linearly on the rotational speed as implied by growth from shear. In this case, equation (3) gives $`L_{sw}ϵ_sB_pB_\varphi \mathrm{\Omega }_{}R_{}^3B_p^2\mathrm{\Omega }_{}^2\tau _DR_{}^3.`$ (31) We now consider the same temporal regimes for the stellar wind power that were examined for the disk wind. For a rapidly rotating star in the $`t<t_m`$ regime, the first term on the right hand side of (22) is dominant and from (31) we then have, $`L_{sw}ϵ_sB_p^2\mathrm{\Omega }_0^2\tau _DR_{}^3=10^{38}ϵ_sB_{p,4}^2\mathrm{\Omega }_{0,4}^2\tau _{D,5}R_{,11}^3\mathrm{erg}/\mathrm{s}.`$ (32) For a very slowly rotating star in the $`t<t_m`$ regime, the second term on the right of (22) is dominant. The spin evolves to within a factor of 1/10 of $`\omega _o`$ after a few years. In this case $`\mathrm{\Omega }_{}`$ is such that $`R_i`$ from (13) is less than $`R_{}`$ so we take $`R_i=R_{}`$. The stellar rotation rate and wind power thus approach $`\mathrm{\Omega }_{}(3\times 10^7)\sqrt{{\displaystyle \frac{GM_c\dot{M}_{a0}^2}{R_{}^3M_{}^2}}},`$ (33) and $`L_{sw}3\times 10^{31}ϵ_s{\displaystyle \frac{B_{p,4}^2\tau _{D,5}M_{c,1}\dot{M}_{a0,22}^2}{M_{,1}^2}}\mathrm{erg}/\mathrm{s}.`$ (34) respectively. The slow time evolution in (22), which results in the approximation that the angular speed and luminosity are constant in (33) and (34), reflects the fact that during this period, the angular momentum gain from accreted material and the loss of angular momentum from the wind nearly balance. For $`t>t_m`$, but for $`\mathrm{\Omega }_{}\tau _D>1`$, using (25) and (13) with (31), we have $`L_{sw}{\displaystyle \frac{ϵ_s\alpha _{ss}^{\frac{2}{5}}(GM_c)^{\frac{6}{5}}\dot{M}_a^{\frac{12}{5}}}{B_p^{14/5}R_{}^3\tau _D}}\left({\displaystyle \frac{H_i}{R_i}}\right)^{\frac{6}{5}}=10^{28}{\displaystyle \frac{ϵ_s\alpha _{ss,2}^{\frac{2}{5}}M_{c,1}^{\frac{6}{5}}\dot{M}_{a0,22}^{\frac{12}{5}}}{B_{p,4}^{\frac{14}{5}}R_{,11}^3\tau _{D,5}}}\left({\displaystyle \frac{t}{50\mathrm{y}\mathrm{r}}}\right)^3\left({\displaystyle \frac{R_i/H_i}{10}}\right)^{\frac{6}{5}}{\displaystyle \frac{\mathrm{erg}}{\mathrm{s}}}.`$ (35) For $`t>t_m`$, when $`\mathrm{\Omega }_{}\tau _D<1`$, we have instead $`L_{sw}{\displaystyle \frac{ϵ_s\alpha _{ss}^{\frac{2}{7}}(GM_c)^{\frac{6}{7}}\dot{M}_a^{\frac{12}{7}}}{B_p^{\frac{10}{7}}R_{}^{\frac{9}{7}}\tau _D}}\left({\displaystyle \frac{H_i}{R_i}}\right)^{\frac{6}{7}}=10^{30}{\displaystyle \frac{ϵ_s\alpha _{ss,2}^{2/7}M_{c,1}^{\frac{6}{7}}\dot{M}_{a0,22}^{\frac{12}{7}}}{B_{p,4}^{\frac{10}{7}}R_{,11}^{\frac{9}{7}}\tau _{D,5}}}\left({\displaystyle \frac{t}{50\mathrm{y}\mathrm{r}}}\right)^{\frac{15}{7}}\left({\displaystyle \frac{R_i/H_i}{10}}\right)^{\frac{6}{5}}{\displaystyle \frac{\mathrm{erg}}{\mathrm{s}}},`$ (36) which is larger in comparison to (35) because a smaller field means a slower decay of $`\mathrm{\Omega }_{}`$. For the characteristic numbers used, $`\mathrm{\Omega }_{}\tau _D`$ from (25) satisfies $`\mathrm{\Omega }_{}\tau _D<1`$ for $`t>t_m`$, as also mentioned above, so here only (36) would apply for our parameter regime and not (35). ### 2.7 Slow and Rapid Initial Stellar Rotation Cases In fig 1 we show plots of the disk to stellar wind luminosity ratios using $`ϵ_s=ϵ_d`$ for two different stellar rotation regimes. Slow Rotation: In fig 1a and 1b plots of $`L_{dw}/L_{sw}`$ for the slow rotator case ($`\mathrm{\Omega }_0<\omega _0`$) are given. Fig 1a shows essentially the ratio of (27) to (34) and fig 1b shows the ratio (27) to (36). For both the $`t<<t_m`$ and $`t>>t_m`$ regimes, the disk wind luminosity strongly dominates the stellar wind luminosity by orders of magnitude. This dominance decreases with time in the former regime and increases in the latter regime. However the main feature is that for the entire range, the disk wind dominates the stellar wind. The star is therefore slaved to the disk throughout the evolution in this case. The stellar wind and spin axes are also slaved to those of the disk because the toroidal field is stretched perpendicularly to the axis of rotation and the wind emanates perpendicularly to the toroidal field. The disk and stellar winds would thus be coaxial, but the disk wind would strongly dominate. Rapid Stellar Rotation: In fig 1c and 1d we present plots of $`L_{dw}/L_{sw}`$ for the rapid rotator case ($`\mathrm{\Omega }_0>\omega _0`$). For $`t<t_m`$ the plot shows basically the ratio of (28) to (32) and for $`t>t_m`$ the plot shows the ratio of (30) to (36). For $`t<t_m`$, the stellar wind is comparable to, and even slightly dominates the disk wind luminosity for the parameters used, while for $`t>t_m`$ the disk wind dominates strongly. As in the slow rotator case, the disk wind domination decreases slowly for $`t<t_m`$ and increases for $`t>t_m`$. Thus for the large $`\mathrm{\Omega }_0`$ case, with $`\mathrm{\Omega }_0`$ near maximal speeds, the stellar wind is independent of the disk until $`t=t_m`$. After this time its MHD wind luminosity is negligible compared to that of the disk. (If there were no disk at all, one would see only the rapid fall of of the stellar wind as in equation (17)). The initial independence of the star from the disk means that the axes of the two winds can be misaligned. The origin of the stellar spin and magnetic axes are decoupled since the spin axis of the disk is independently determined by the initial plane of the binary orbit. Thus the disk and stellar wind system can produce a multi-polar wind system of up to $`10^{38}`$erg/s, depending on the efficiency of conversion of rotational energy. ## 3 Discussion Our results show that powerful magneto-centrifugal winds can be driven in PNe systems. In the context of accretion disks formed from the break-up of a secondary star in a binary system, these winds are time dependent and have a finite lifetime. Thus our results would predict that transient, magnetized outflows can be driven during the transition from AGB to PNe. Magneto-centrifugally driven winds potentially solve two problems which have emerged from PNe studies (Frank 2000). First, the large powers computed above can explain how strong winds can be driven independently of radiation line driving processes as required by observations in the PPNe phase (e.g. Alcolea et al 2000). Second, they may explain the high degree of collimation seen in some sources (e.g. Sahai & Trauger 1998), and multi-polar outflows when the disk and stellar winds are contemporaneous. Below we discuss these points further. ### 3.1 Accounting for the power of PN or PPN winds: The calculations above show that magnetically driven winds can account for the large powers observed in PPNe. Alcolea et al. (2000) make the case that observations of PPNe are the best place to infer the required powers and shaping mechanisms of PNe. None of the 3 objects for which Alcolea et al. (2000) have measured kinetic ages (OH 231.8; M 1-92; M2-56) can be radiation driven at the PPNe stage. The required outflow powers range from $`10^{35}`$ erg/s to $`10^{37}`$erg/s. The kinetic ages are fairly similar, ranging from 750 yr to 1800 yr. The radiation driving falls short of being able to provide the power by at least two orders of magnitude. The 9 objects for which Alcolea et. al (2000) have measured wind energies but not kinetic ages look reasonably similar in terms of radiation power and energetics. Since the kinetic ages for the few measured objects are not widely varying, it is likely that others in the population of 9 also probably require high outflow powers. The powers computed in section 2 for MHD wind driving can be as high as $`10^{38}`$ erg/s for the parameters used. Such powers are therefore sufficient for PNe and PPNe. There is likely to be some efficiency factor with which the driving can take place, so to explain a $`10^{37}`$ erg/s wind, we would require an efficiency factor of $`ϵ10\%`$. More study is needed to determine the $`ϵ`$ factors dynamically. For $`t>t_m`$, the disk wind power falls by two orders of magnitude or so, while the stellar MHD wind is of negligible power. The MHD disk wind could still supply e.g. $`10^{36}`$ ergs/s after $`50`$ yr, continuing to fall as $`t^{1.6}`$ from (30). Specific examples of kinetic powers required in this approximate range include Hb5 ($`V_w\text{ }<400`$ km/s and $`\dot{M}6\times 10^5\dot{M}_{}`$/yr; \[Corradi & Schwarz 1995; Pishmish et al. 2000\]) thus with a kinetic luminosity $`L_w\text{ }<3.6\times 10^{36}`$erg/s; IC4406 ($`V_w65`$ km/s \[Corradi al. 1997\] and $`\dot{M}10^4\dot{M}_{}`$/yr \[Sahai et al. 1991; Cox et al. 1991\]) with $`L_w2\times 10^{35}`$erg/s; M2-9 ($`V_w46`$ km/s \[Solf 2000\] and $`\dot{M}4\times 10^6\dot{M}_{}`$/yr \[Zweigle et al. 1997\]) with $`L_w5\times 10^{33}`$erg/s. Observations of NGC 3242, NGC 6828, NGC 7009, and NGC 7662 PNe (Balick et al. 1998; Dwarkadas & Balick 1998) are also consistent with $`10^{33}\mathrm{erg}/\mathrm{s}L_w10^{35}\mathrm{erg}/\mathrm{s}`$. In all cases, a dynamo in either a disk or the star or both may be responsible for the field driving the wind. Magnetic flaring and the production of hard coronal X-rays might then be expected in PPN engines. Observations with Chandra and other X-ray telescopes are therefore be highly desirable. ### 3.2 Accounting for the shapes of PN or PPN winds ?: Given that the fields can power jets, they may also play a role in shaping them. One possible approach to MHD collimation in these sources is simulated in García-Segura et al. (1999). Since we have not yet dynamically calculated or simulated collimation mechanisms as applied to planetary nebulae the discussion that follows is speculative. The fact that the power of the disk and stellar winds can be comparable motivates us to proceed. It is known that in the context of disks (e.g. Königl & Pudritz 2000 for review) and even in the context of stars (Tsinganos & Bogovalov 2000) that magnetic self collimation can be an important effect. The literature on collimation of disk winds is extensive. It is noteworthy that in emphasizing that fast rotating stellar magnetospheres can also collimate winds, Tsinganos & Bogovalov (2000) point out that such collimation from an underlying rotator will likely be effective when the quantity $`Q=0.12\left({\displaystyle \frac{\psi }{\psi _{}}}\right)\left({\displaystyle \frac{\mathrm{\Omega }}{\mathrm{\Omega }_{}}}\right)\left({\displaystyle \frac{\dot{M}_w}{\dot{M}_{}}}\right)^{1/2}\left({\displaystyle \frac{V_w}{V_{}}}\right)^{3/2}`$ (37) is greater than 1. Here $`\psi `$ is the magnetic flux anchored on the rotator spinning at angular speed $`\mathrm{\Omega }`$, $`\dot{M}_w`$ is the outflow mass loss rate, $`V_w`$ is the outflow speed, and these values are scaled to the respective values for the Sun, namely $`V_{}=400`$km/s, $`\dot{M}_{}=1.6\times 10^{12}`$g/s, $`\mathrm{\Omega }_{}=3\times 10^6`$/s and $`\psi _{}=10^{22}`$G-cm<sup>2</sup>. We then have $`Q/Q_{}4.2\left({\displaystyle \frac{B_{,4}R_{,11}^2}{10^{26}\mathrm{G}\mathrm{cm}^2}}\right)\left({\displaystyle \frac{\mathrm{\Omega }_{}}{10^4/\mathrm{s}}}\right)\left({\displaystyle \frac{\dot{M}_{sw}}{6\times 10^{21}\mathrm{g}/\mathrm{s}}}\right)^{1/2}\left({\displaystyle \frac{V_{sw}}{400\mathrm{k}\mathrm{m}/\mathrm{s}}}\right)^{3/2},`$ (38) where we have scaled to our fiducial post-AGB star parameters, taking the upper limit of the magnetic flux, and assuming a kinetic luminosity of the wind $`L_{kin}\dot{M}_{sw}V_{sw}^210^{37}`$erg/s. We can see that $`Q/Q_{}`$ can be greater than 1 for the chosen parameters, and thus collimation in excess of that for the solar wind can occur. If we instead use representative outflow parameters for later stages of PNe (e.g. Balick et al. 1998), such as $`\dot{M}=5\times 10^7M_{}\mathrm{yr}^1`$ and $`V=200\mathrm{km}\mathrm{s}^1`$), we find $`Q/Q_{}\text{ }>100`$ which implies significant collimation. In general, these estimates suggest that self-collimation of magnetized winds in PPNe and PNe might be possible without requiring shocks to aid in the process. The shock-aided approach has been studied in García-Segura et al. (1999). More study is warranted. The relative self-collimation of the disk wind and the stellar wind may be different for different systems. This needs to be studied. Assuming collimation can ensue, the shape of any multi-polar magnetically driven outflow in our picture will also depend on the stellar spin down time scale $`t_m`$, and the time scale for the star to become a PN, $`t_{PN}`$. The ratio of $`t_m/t_{PN}`$ will determine how much shaping will be controlled by the a combined star+disk wind. When $`t_m/t_{PN}<<1`$, the stellar wind is likely to provide micro-structures such as knots, bullets, or ansae in mature PNe. When $`t_m/t_{PN}\text{ }>1`$ then structures such as multi-polar bubbles, which develop during the period when $`L_{dw}L_{sw}`$, may leave a large enough imprint to survive into the PN phase. Another important parameter is the ratio of the wind momentum densities, $`p_{dw}/p_{sw}=\rho _{dw}V_{dw}/(\rho _{sw}V_{sw})`$ where the ratio to be used here is that computed for each type of wind independent of the actual presence of the other. The momentum density depends on collimation and is thus a function of angle from a wind’s axis. The terminal wind speed of material at a given radius from any wind axis depends on the details of the wind launching process but should be of order several times the the escape speed at the radius at which the field lines threading that radius are anchored (e.g. Königl & Pudritz 2000 for a review). The ratio $`p_{dw}/p_{sw}`$ can in part determine the shape of the resulting observed outflow combination by determining which wind-driven “lobe” would extend farther when the two winds are superimposed. The difference in initial launch times of the two winds is also very important in determining the observed shapes. For two concurrent winds satisfying $`p_{dw}/p_{sw}>1`$ on average, and for which the star is rotating slowly ($`\mathrm{\Omega }_o<\omega _o`$), we might expect two nested aligned bipolar bubbles (as shown in Fig 2a) when the terminal speed of the stellar wind, calculated independently from the disk wind, is faster over most of the outflow cross-section. The stellar wind can sweep up disk material, but because of the lower overall momentum density, it cannot proceed as far in distance. An example of a PN with this shape is Hubble 12 (Welch et al.1999). When the stellar and disk magnetic and rotational axes are misaligned, a nested multi-polar bubble may emerge as shown in Fig 2c. A possible example of such an outflow is M2-46 (Manchado 2000) Please note again: Fig 2 represents cartoon speculations, not simulations. When $`p_{dw}/p_{sw}<1`$ and the star is rapidly rotating ($`\mathrm{\Omega }_o>\omega _o`$), the stellar wind could push ahead of the disk wind and may produce the features seen in e.g. NGC 7007 (Fig 2b, Balick et al.1998). In the case where the star is initially rapidly rotating and the stellar and disk axes are not aligned, then a second form of the nested multi-polar bubble may emerge as is shown in Fig 2d. An example of such an outflow may be He2-155 (Sahai 2000). In the rapidly rotating case, the stellar wind remains more or less constant over a spin-down period while the disk wind falls. Thus we might expect the stellar wind to appear more or less continuous, whereas the disk wind could produce structures only of comparable luminosity to the stellar wind for the first several years. In this initial phase, antipodal knots amidst a continuous stellar wind might appear where the disk and stellar winds interact, or where the disk interacts with the ambient medium. Such processes might have been occurring in the Egg nebula (Jura et al. 2000). Given that typical sound speeds for observed outflows are of order 20-30 km/s, while the wind velocities can be $`>100`$km/s, we can estimate the size of outflow structures compared with the distance from the star. For a hybrid disk-star outflow with strongest power phase lasting $``$ few years, but with an age or order hundreds of years, one should see knots (or shells depending on collimation) with a width of order $`10^{16}`$ cm. If the outflow velocity were 100 km/s, then the knot/shell distance from the central core would be $`5`$ times the shell thickness. Such features are seen in the Cat’s Eye nebula (Miranda & Solf 1992) Note that two interacting winds can lead to more than 2 multiple symmetry axis due to edge features. This and the features described above could also result from time dependent effects of separate interacting stellar winds or separate interacting disk winds where the star or disk is allowed to precess (e.g. Manchado et al. 1996). Precession would enable successive outbursts to be emitted on multiple axes. If there is precession, then one might see an additional point symmetry. These ideas may be applied to other coupled disk+star systems (Soker & Livio 1994, Mastrodemos & Morris 1998). ## 4 Summary and Conclusions We have shown that MHD winds from either disks or post-AGB stars can comfortably power PNe outflows using reasonable estimates of stellar and disk field strengths and rotation parameters. The available MHD wind power turns out to be many orders of magnitude larger than that available for radiation/line driven outflows inferred from observations of PPNe (e.g Alcolea et al. 2000). Our study also highlights the importance of understanding the interplay between disk winds and stellar winds in contributing to the diversity of observed multi-polar PNe morphologies. We derived maximal disk and stellar MHD wind powers for a coupled disk-star system at times $`t<t_m`$, $`t>t_m`$, (where $`t_m`$ is the stellar spin down time) and for fast and slow stellar rotators, $`\mathrm{\Omega }_{}/\omega _0<\mathrm{or}>1`$, and values of $`\mathrm{\Omega }_{}\tau _D<\mathrm{or}>1`$. Our calculations take the initial time to be the onset of disk formation in a binary system from common envelope expulsion. The initial onset of the star and disk wind should be nearly contemporaneous. The layers of the star with the strong magnetic field that power the wind are only exposed after the common envelope is ejected, and the accretion disk forms within a year from the common envelope ejection. The resulting luminosities for both winds depend only on the disk accretion rate and properties of the star. This is because the inner regions of the disk are the main contributor to the disk wind and the disk properties can be constrained using properties of magneto-shearing disks and the interaction with the star. We now summarize some numbers for the maximum wind powers: * For $`t<<t_m`$ and our choice of fiducial scalings, the maximum MHD luminosity of a stellar wind from an initially rapidly rotating post AGB star with field strengths consistent with those estimated for AGB dynamos of $`10^4`$ Gauss satisfies $`L_{sw}\text{ }<10^{38}`$ erg/s, and is approximately constant in time. The disk wind satisfies $`L_{dw}\text{ }<7\times 10^{37}(t/1\mathrm{y}\mathrm{r})^{1.6}`$ as the accretion rate falls. * After $`t_m>>50\mathrm{y}\mathrm{r}`$, $`L_{sw}\text{ }<10^{30}(t/50\mathrm{y}\mathrm{r})^{2.1}`$, and the disk wind power goes as $`L_{dw}\text{ }<7\times 10^{35}(t/50\mathrm{y}\mathrm{r})^{1.6}`$. * Whilst the dominance of disk vs. stellar wind depends on $`\mathrm{\Omega }_0`$ for $`t<<t_m`$, for all $`t>>t_m`$, the disk wind always dominates, regardless of $`\mathrm{\Omega }_0`$. Since the MHD disk winds dominate the stellar winds for $`t>t_m`$ regardless of the initial stellar spin rate, the star is slaved to the disk in that regime. For $`t<t_m`$, the stellar wind power can equal or exceed that of the disk, and in this regime the stellar wind can emanate along a different symmetry axis than the disk wind. Multi-polar outflows can then be produced. For a rapidly rotating star, the stellar MHD wind power is steady for a time of order $`t=t_m`$, while the disk wind power falls rapidly. For $`t>t_m`$, the disk wind power falls less rapidly than the stellar wind power. Generally, we would suggest that if a sustained MHD driven bipolar outflow extending to the core lasts for $`>500`$ yr, the disk is dominant. The choice of scalings we use in the text were mainly to provide a general framework. The spin down time $`t_m`$ could be longer than our scale of 50 yr depending on radius and stellar convective diffusion time, (c.f. eqn (18)). Disk winds imply initial formation from a binary system. An initially rapidly rotating star would slow down rapidly by this later time, and would only be able to contribute significantly if its wind were radiatively driven at this stage. If a stellar wind operates powerfully during $`t<t_m`$ it means that the central star was initially rotating rapidly. There is no guarantee that a binary system would form, or that the star would be initially rapidly rotating for every post AGB star. Thus for all systems with bipolar outflows, observations that can reveal the rotation rate of the central star, its magnetic field, or the presence of binary companions are ultimately essential for testing MHD disk/stellar wind paradigms. Further studies of the disk wind-stellar wind interplay are most certainly needed. Given the ubiquity of the accretion-central object-outflow connection throughout the universe the same kind of investigation may ultimately be important for other stellar and compact object accretion systems as well. Acknowledgements: E.B. was supported in part by NSF grant PHY94-07194 at the ITP, and acknowledges support from a DOE Plasma Physics Junior Faculty Development Award. E.B. also thanks R. Krasnolpolsky, R. Blandford, B. Chandran, and E. Ostriker for discussions. Many thanks to G. García-Segura and M. Reyes-Ruiz for their careful readings which led to important changes. Figure Captions Figure 1: Plot of the MHD wind power ratio for the disk and star for four regimes with $`ϵ_s=ϵ_d`$: a) $`t<t_m`$ for the initially slowly rotating star b) $`t>t_m`$ for the initially slowly spinning star c) $`t<t_m`$ regime for initially fast rotating star. d) $`t>t_m`$ regime for initially fast rotating star. See text. Figure 2: Cartoon representing our speculation on a possible correspondence between between HST PN images and multi-polar wind structures from the hybrid stellar-disk wind paradigm, for different initial stellar spins and thrust ratios $`\mathrm{\Pi }_{dw}/\mathrm{\Pi }_{sw}`$ (see text). (Note that these are not simulations, but speculations.) a) top left: $`\mathrm{\Omega }_0<\omega _0`$ and $`\mathrm{\Pi }_{dw}/\mathrm{\Pi }_{sw}>1`$ so star is slaved to disk and the only possibility is nested winds. 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warning/0005/quant-ph0005029.html	ar5iv	text	# Quantum decoherence and the Glauber dynamics from the Stochastic limit ## 1 Introduction In the present paper we investigate a general model of quantum system interacting with a bosonic reservoir via an Hamiltonian of the form $$H=H_0+\lambda H_I$$ where $`H_0`$ is called the free Hamiltonian and $`H_I`$ the interaction Hamiltonian. The stochastic golden rules, which arise in the stochastic limit of quantum theory as natural generalizations of Fermi golden rule , , provide a natural tool to associate a stochastic flow, driven by a white noise equation (stochastic Schrödinger) equation, to any discrete system interacting with a quantum field. This white noise Hamiltonian equation which, when put in normal order becomes equivalent to a quantum stochastic differential equation. The Langevin (stochastic Heisenberg) and master equations are deduced from this white noise equation by means of standard procedures which are described in . We use these equations to investigate the decoherence in quantum systems. In the work , extending previous results obtained with perturbative techniques by , it was shown on the example of the spin–boson Hamiltonian that the decoherence in quantum systems can be controlled by the following constants (cf. section 2 for the definition of the quantities involved) $$\text{Re}(g\|g)_\omega ^+=𝑑k\|g(k)\|^22\pi \delta (\omega (k)\omega )N(k),$$ For the simplest case of the equalibrium state of the reservoir with the temperature $`\beta ^1=kT`$ this constant will be equal to $`\frac{kT}{\mathrm{}}`$ (actually this is true for large temperatures and for the dispersion function $`\omega (k)=\|k\|`$). In this paper we extend the approach of from $`2`$–level systems to arbitrary quantum systems with discrete spectrum. Our resuts show that the stocastic limit technique gives us an effective method to control quantum decoherence. We find that, under the above mentioned interaction, all the off–diagonal matrix elements, of the density matrix of a generic discrete quantum system, will decay exponentially if $`\text{Re}(g\|g)`$ are nonzero. In other words we obtain the asymptotic diagonalization of the density matrix. Moreover, we show that for generic quantum system the off–diagonal elements of the density matrix decay exponentially as $`\mathrm{exp}(N\text{Re}(g\|g)t)`$, with the exponent proportional to the number $`N`$ of particles in the system. Therefore for generic macroscopic (large $`N`$) systems the quantum state will collapse into the classical state very quickly. This effect was built in by hands in several phenomenological models of the quantum measurement process. In the stochastic limit approach it is deduced from the Hamiltonian model. This observation contributes to the clarification of one of the old problems of quantum theory: Why macroscopic systems usually behave classically? i.e. why do we observe classical states although the evolution of the system is a unitary operator described by the Shrödinger equation? Moreover, this result allows to distinguish between macroscopic systems (that behave classically) and microscopic systems (where quantum effects are important). Quantum effects (or effects of quantum interference) are connected with the off–diagonal elements of the density matrix. Therefore the following notion is natural: the macroscopic system is a system where off–diagonal elements of the density matrix decay quickly (faster than the minimal time of observation). Using that off–diagonal elements decay as $`\mathrm{exp}(N\text{Re}(g\|g)t)`$, we get the following definition of the macroscopic system: $`N\text{Re}(g\|g)>>1`$. The quantum Markov semigroup we obtain lives invariant the algebra of the spectral projections of the system Hamiltonian and the associated master equation, when restricted to the diagonal part of the density matrix, takes the form of a standard classical kinetic equation, describing the convergence to equilibrium (Gibbs state) of the system, coupled with the given reservoir (quantum field). Summing up: the convergence to equilibrium is a result of quantum decoherence. If we can control the interaction so that some of the constants $`\text{Re}(g\|g)`$ are zero, then the corresponding matrix elements will not decay in the stochastic approximation, i.e. in a time scale which is extremely long with respect to the slow clock of the discrete system. In this sense the stochastic limit approach provides a method for controlling quantum coherence. The general idea of the stochastic limit (see ) is to make the time rescaling $`tt/\lambda ^2`$ in the solution of the Schrödinger (or Heisenberg) equation in interaction picture $`U_t^{(\lambda )}=e^{itH_0}e^{itH}`$, associated to the Hamiltonian $`H`$, i.e. $$\frac{}{t}U_t^{(\lambda )}=i\lambda H_I(t)U_t^{(\lambda )}$$ with $`H_I(t)=e^{itH_0}H_Ie^{itH_0}`$. This gives the rescaled equation $$\frac{}{t}U_{t/\lambda ^2}^{(\lambda )}=\frac{i}{\lambda }H_I(t/\lambda ^2)U_{t/\lambda ^2}^{(\lambda )}$$ (1) and one wants to study the limits, in a topology to be specified, $$\underset{\lambda 0}{lim}U_{t/\lambda ^2}^{(\lambda )}=U_t;\underset{\lambda 0}{lim}\frac{1}{\lambda }H_I\left(\frac{t}{\lambda ^2}\right)=H_t$$ (2) The limit $`\lambda 0`$ after the rescaling $`tt/\lambda ^2`$ is equivalent to the simultaneous limit $`\lambda 0`$, $`t\mathrm{}`$ under the condition that $`\lambda ^2t`$ tends to a constant (interpreted as a new slow time scale). This limit captures the dominating contributions to the dynamics, in a regime of long times and small coupling, arising from the cumulative effects, on a large time scale, of small interactions $`(\lambda 0)`$. The physical idea is that, looked from the slow time scale of the atom, the field looks like a very chaotic object: a quantum white noise, i.e. a $`\delta `$–correlated (in time) quantum field $`b^{}(t,k)`$, $`b(t,k)`$ also called a master field. The structure of the present paper is as follows. In section 2 we introduce the model and consider its stochastic limit. In section 3 we derive the Langevin equation. In section 4 we derive the master equation for the density matrix and show that for non–zero decoherence the master equation describes the collapse of the density matrix to the classical Gibbs distribution and discuss the connection of this fact with the procedure of quantum measurement. In section 5, using the characterization of quantum decoherence obtained in section 4 and generalizing arguments of , we find that our general model exhibits macroscopic quantum effects (in particular, conservation of quantum coherence). These effects are controllable by the state of the reservoir (that can be controlled by filtering). In section 6 we apply our general scheme to the model of a quantum system of spins interacting with bosonic field and derive a quantum extension of the Glauber dynamics. Thus our stochastic limit approach provides a microscopic interpretation, in terms of fundamental Hamiltonian models, to the dynamics of quantum spin systems. Moreover we deduce the full stochastic equation and not only the master equation. This is new even in the case of classical spin systems. ## 2 The model and it’s stochastic limit In the present paper we consider a general model, describing the interaction of a system $`S`$ with a reservoir, represented by a bosonic quantum field. Particular cases of this general model were investigated in , , . The total Hamiltonian is $$H=H_0+\lambda H_I=H_S+H_R+\lambda H_I$$ where $`H_R`$ is the free Hamiltonian of a bosonic reservoir $`R`$: $$H_R=\omega (k)a^{}(k)a(k)𝑑k$$ acting in the representation space $``$ corresponding to the state $``$ of bosonic reservoir generated by the density matrix $`𝐍`$ that we take in the algebra of spectral projections of the reservoir Hamiltonian. The reference state $``$ of the field is a mean zero gauge invariant Gaussian state, characterized by the second order correlation function equal to $$a(k)a^{}(k^{})=(N(k)+1)\delta (kk^{})$$ $$a^{}(k)a(k^{})=N(k)\delta (kk^{})$$ where the function $`N(k)`$ describes the density of bosons with frequency $`k`$. One of the examples is the (gaussian) bosonic equilibrium state at temperature $`\beta ^1`$. The system Hamiltonian has the following spectral decomposition $$H_S=\underset{r}{}\epsilon _rP_{\epsilon _r}$$ where the index $`r`$ labels the spectral projections of $`H_S`$. For example, for a non–degenerate eigenvalue $`\epsilon _r`$ of $`H_S`$ the corresponding spectral projection is $$P_{\epsilon _r}=\|\epsilon _r\epsilon _r\|$$ where $`\|\epsilon _r`$ is the corresponding eigenvector. The interaction Hamiltonian $`H_I`$ (acting in $`_S`$) has the form $$H_I=\underset{j}{}\left(D_j^{}A(g_j)+D_jA^{}(g_j)\right),A(g)=𝑑k\overline{g}(k)a(k),$$ where $`A(g)`$ is a smeared quantum field with cutoff function (form factor) $`g(k)`$. To perform the construction of the stochastic limit one needs to calculate the free evolution of the interaction Hamiltonian: $`H_I(t)=e^{itH_0}H_Ie^{itH_0}`$. Using the identity $$1=\underset{r}{}P_{\epsilon _r}$$ we write the interaction Hamiltonian in the form $$H_I=\underset{j}{}\underset{rr^{}}{}P_{\epsilon _r}D_j^{}P_{\epsilon _r^{}}𝑑k\overline{g}_j(k)a(k)+h.c.$$ (3) Let us introduce the set of energy differences (Bohr frequencies) $$F=\{\omega =\epsilon _r\epsilon _r^{}:\epsilon _r,\epsilon _r^{}\text{Spec}H_S\}$$ and the set of all energies of the form $$F_\omega =\{\epsilon _r:\epsilon _r^{}(\epsilon _r,\epsilon _r^{}\text{Spec}H_S)\text{ such that }\epsilon _r\epsilon _r^{}=\omega \}$$ With these notations we rewrite the interaction Hamiltonian (3) in the form $$H_I=\underset{j}{}\underset{\omega F}{}\underset{\epsilon _rF_\omega }{}P_{\epsilon _r}D_j^{}P_{\epsilon _r\omega }𝑑k\overline{g}_j(k)a(k)+h.c.=$$ $$=\underset{j}{}\underset{\omega F}{}E_\omega ^{}\left(D_j\right)𝑑k\overline{g}_j(k)a(k)+h.c.$$ (4) where $$E_\omega (X):=\underset{\epsilon _rF_\omega }{}P_{\epsilon _r\omega }XP_{\epsilon _r}$$ (5) It is easy to see that the free volution of $`E_\omega (X)`$ is $$e^{itH_S}E_\omega (X)e^{itH_S}=e^{it\omega }E_\omega (X)$$ Using the formula for the free evolution of bosonic fields $$e^{itH_R}a(k)e^{itH_R}=e^{it\omega (k)}a(k)$$ we get for the free evolution of the interaction Hamiltonian: $$H_I(t)=\underset{j}{}\underset{\omega F}{}E_\omega ^{}\left(D_j\right)𝑑k\overline{g}_j(k)e^{it(\omega (k)\omega )}a(k)+h.c.$$ (6) In the stochastic limit the field $`H_I(t)`$ gives rise to a family of quantum white noises, or master fields. To investigate these noises, let us suppose the following: 1) $`\omega (k)0`$, $`k`$; 2) The $`d1`$–dimensional Lebesgue measure of the surface $`\{k:\omega (k)=0\}`$ is equal to zero (so that $`\delta (\omega (k))=0`$) (for example $`\omega (k)=k^2+m`$ with $`m0`$). Now let us investigate the limit of $`H_I(t/\lambda ^2)`$ using one of the basic formulae of the stochastic limit: $$\underset{\lambda 0}{lim}\frac{1}{\lambda ^2}\mathrm{exp}\left(\frac{it}{\lambda ^2}f(k)\right)=2\pi \delta (t)\delta (f(k))$$ (7) which shows that the term $`\delta (f(k))`$ in (7) is not identically equal to zero only if $`f(k)=0`$ for some $`k`$ in a set of nonzero $`d1`$–dimensional Lebesgue measure. This explains condition (2) above. The rescaled interaction Hamiltonian is expressed in terms of the rescaled creation and annihilation operators $$a_{\lambda ,\omega }(t,k)=\frac{1}{\lambda }e^{i\frac{t}{\lambda ^2}(\omega (k)\omega )}a(k),\omega F$$ After the stochastic limit every rescaled annihilation operator corresponding to any transition from $`\epsilon _r^{}`$ to $`\epsilon _r`$ with the frequency $`\omega =\epsilon _r\epsilon _r^{}`$ generates one non–trivial quantum white noise $$b_\omega (t,k)=\underset{\lambda 0}{lim}a_{\lambda ,\omega }(t,k)=\underset{\lambda 0}{lim}\frac{1}{\lambda }e^{i\frac{t}{\lambda ^2}(\omega (k)\omega )}a(k)$$ with the relations $$[b_\omega (t,k),b_\omega ^{}(t^{},k^{})]=\underset{\lambda 0}{lim}[a_{\lambda ,\omega }(t,k),a_{\lambda ,\omega }^{}(t^{},k^{})]=$$ $$=\underset{\lambda 0}{lim}\frac{1}{\lambda ^2}e^{i\frac{tt^{}}{\lambda ^2}(\omega (k)\omega )}\delta (kk^{})=2\pi \delta (tt^{})\delta (\omega (k)\omega )\delta (kk^{})$$ (8) $$[b_\omega (t,k),b_\omega ^{}^{}(t^{},k^{})]=0$$ (cf. 7). This shows, in particular that quantum white noises, corresponding to different Bohr frequencies, are mutually independent. The stochastic limit of the interaction Hamiltonian is therefore equal to $$h(t)=\underset{j}{}\underset{\omega F}{}E_\omega ^{}\left(D_j\right)𝑑k\overline{g}_j(k)b_\omega (t,k)+h.c.$$ (9) The state of the master field (white noise) $`b_\omega (t,k)`$, corresponding to our choice of the initial state of the field, is the mean zero gauge invariant Gaussian state with correlations: $$b_\omega ^{}(t,k)b_\omega (t^{},k^{})=2\pi \delta (tt^{})\delta (\omega (k)\omega )\delta (kk^{})N(k)$$ $$b_\omega (t,k)b_\omega ^{}(t^{},k^{})=2\pi \delta (tt^{})\delta (\omega (k)\omega )\delta (kk^{})(N(k)+1)$$ and vanishes for noises corresponding to different Bohr frequences. Now let us investigate the evolution equation in interaction picture for our model. According to the general scheme of the stochastic limit, we get the (singular) white noise equation $$\frac{d}{dt}U_t=ih(t)U_t$$ (10) whose normally ordered form is the quantum stochastic differential equation $$dU_t=(idH(t)Gdt)U_t$$ (11) where $`h(t)`$ is the white noise Hamiltonian (9) given by the stochastic limit of the interaction Hamiltonian and $$dH(t)=\underset{j}{}\underset{\omega F}{}\left(E_\omega ^{}\left(D_j\right)dB_{j\omega }(t)+E_\omega \left(D_j\right)dB_{j\omega }^{}(t)\right)$$ (12) $$dB_{j\omega }(t)=𝑑k\overline{g}_j(k)_t^{t+dt}b_\omega (\tau ,k)𝑑\tau $$ (13) According to the stochastic golden rule (11) the limit dynamical equation is obtained as follows: the first term in (11) is just the limit of the iterated series solution for (1) $$\underset{\lambda 0}{lim}\frac{1}{\lambda }_t^{t+dt}H_I\left(\frac{\tau }{\lambda ^2}\right)𝑑\tau $$ The second term $`Gdt`$, called the drift, is equal to the limit of the expectation value in the reservoir state of the second term in the iterated series solution for (1) $$\underset{\lambda 0}{lim}\frac{1}{\lambda ^2}_t^{t+dt}𝑑t_1_t^{t_1}𝑑t_2H_I\left(\frac{t_1}{\lambda ^2}\right)H_I\left(\frac{t_2}{\lambda ^2}\right)$$ Making in this formula the change of variables $`\tau =t_2t_1`$ we get $$\underset{\lambda 0}{lim}\frac{1}{\lambda ^2}_t^{t+dt}𝑑t_1_{tt_1}^0𝑑\tau H_I\left(\frac{t_1}{\lambda ^2}\right)H_I\left(\frac{t_1}{\lambda ^2}+\frac{\tau }{\lambda ^2}\right)$$ (14) Computing the expectation value and using the fact that the limits of oscillating factors of the form $`\underset{\lambda 0}{lim}e^{\frac{ict_1}{\lambda ^2}}`$ vanish unless the constant $`c`$ is equal to zero, we see that we can have non–zero limit only when all oscillating factors of a kind $`e^{\frac{ict_1}{\lambda ^2}}`$ (with $`t_1`$) in (14) cancel. In conclusion we get $$G=\underset{ij}{}\underset{\omega F}{}_{\mathrm{}}^0d\tau (dk\overline{g_i(k)}g_j(k)e^{i\tau (\omega (k)\omega )}(N(k)+1)E_\omega ^{}\left(D_i\right)E_\omega \left(D_j\right)+$$ $$+dkg_i(k)\overline{g_j(k)}e^{i\tau (\omega (k)\omega )}N(k)E_\omega \left(D_i\right)E_\omega ^{}\left(D_j\right))$$ and therefore, from the formula $$_{\mathrm{}}^0e^{it\omega }𝑑t=\frac{i}{\omega i0}=\pi \delta (\omega )i\text{P.P.}\frac{1}{\omega }$$ (15) we get the following expression for the drift $`G`$: $$\underset{ij}{}\underset{\omega F}{}(dk\overline{g_i(k)}g_j(k)\frac{i(N(k)+1)}{\omega (k)\omega i0}E_\omega ^{}\left(D_i\right)E_\omega \left(D_j\right)+$$ $$+dkg_i(k)\overline{g_j(k)}\frac{iN(k)}{\omega (k)\omega +i0}E_\omega \left(D_i\right)E_\omega ^{}\left(D_j\right))=$$ $$=\underset{ij}{}\underset{\omega F}{}\left((g_i\|g_j)_\omega ^{}E_\omega ^{}\left(D_i\right)E_\omega \left(D_j\right)+\overline{(g_i\|g_j)}_\omega ^+E_\omega \left(D_i\right)E_\omega ^{}\left(D_j\right)\right)$$ (16) Let us note that for (16) we have the following Cheshire Cat effect found in : even if the frequency $`\omega `$ is negative and therefore does not generate a quantum white noise the corresponding values $`(g\|g)_\omega ^\pm `$ in (16) will be non–zero. In other terms: negative Bohr frequencies contribute to an energy shift in the system, but not to its damping. Remark If $`F`$ is any subset of $`\text{Spec}H_S`$ and $`X_r`$ are arbitrary bounded operators on $`_S`$ then for any $`tR`$ $$e^{itH_S}\underset{\epsilon _rF}{}P_{\epsilon _r}X_rP_{\epsilon _r}=\underset{\epsilon _rF}{}e^{it\epsilon _r}P_{\epsilon _r}X_rP_{\epsilon _r}=\underset{\epsilon _rF}{}P_{\epsilon _r}X_rP_{\epsilon _r}e^{itH_S}$$ In other words: $`_{\epsilon _rF}P_{\epsilon _r}X_rP_{\epsilon _r}`$ belongs to the commutant $`L^{\mathrm{}}\left(H_S\right)^{}`$ of the abelian algebra $`L^{\mathrm{}}\left(H_S\right)`$, generated by the spectral projections of $`H_S`$. A corollary of this remark is that, for each $`\omega F`$, for any bounded operator $`XL^{\mathrm{}}\left(H_S\right)^{}`$ and for each pair of indices $`(i,j)`$ the operators $$E_\omega \left(D_i\right)XE_\omega ^{}\left(D_j\right),E_\omega ^{}\left(D_i\right)XE_\omega \left(D_j\right)$$ (17) belong to the commutant $`L^{\mathrm{}}\left(H_S\right)^{}`$ of $`L^{\mathrm{}}\left(H_S\right)`$. In particular, if $`H_S`$ has non–degenerate spectrum so that $`L^{\mathrm{}}\left(H_S\right)`$ is a maximal abelian subalgebra of $`B\left(_S\right)`$, the operators (17) also belong to $`L^{\mathrm{}}\left(H_S\right)`$. ## 3 The Langevin equation Now we will find the Langevin equation, which is the limit of the Heisenberg evolution, in interaction representation. Let $`X`$ be an observable. The Langevin equation is the equation satisfied by the stochastic flow $`j_t`$, defined by: $$j_t(X)=U_t^{}XU_t$$ where $`U_t`$ satisfies equation (11) in the previous section, i.e. $$dU_t=(idH(t)Gdt)U_t$$ (18) To derive the Langevin equation we consider $$dj_t(X)=j_{t+dt}(X)j_t(X)=dU_t^{}XU_t+U_t^{}XdU_t+dU_t^{}XdU_t$$ (19) The only nonvanishing products in the quantum stochastic differentials are $$dB_{i\omega }(t)dB_{j\omega }^{}(t)=2\text{Re}(g_i\|g_j)_\omega ^{}dt,dB_{i\omega }^{}(t)dB_{j\omega }(t)=2\text{Re}(g_i\|g_j)_\omega ^+dt$$ (20) Combining the terms in (19) and using (18), (12), (16) and (20) we get the Langevin equation $$dj_t(X)=\underset{\alpha }{}j_t\theta _\alpha (X)dM^\alpha (t)=\underset{n=1,1;j\omega }{}j_t\theta _{nj\omega }(X)dM^{nj\omega }(t)+j_t\theta _0(X)dt$$ (21) where $$dM^{1,j\omega }(t)=dB_{j\omega }(t),\theta _{1,j\omega }(X)=i[X,E_\omega ^{}\left(D_j\right)]$$ (22) $$dM^{1,j\omega }(t)=dB_{j\omega }^{}(t),\theta _{1,j\omega }(X)=i[X,E_\omega \left(D_j\right)]$$ (23) and $$\theta _0(X)=\underset{ij}{}\underset{\omega F}{}(i\text{ Im }(g_i\|g_j)_\omega ^{}[X,E_\omega ^{}\left(D_i\right)E_\omega \left(D_j\right)]+i\text{ Im }(g_i\|g_j)_\omega ^+[X,E_\omega \left(D_i\right)E_\omega ^{}\left(D_j\right)]+$$ $$+2\text{Re}(g_i\|g_j)_\omega ^{}\left(E_\omega ^{}\left(D_i\right)XE_\omega \left(D_j\right)\frac{1}{2}\{X,E_\omega ^{}\left(D_i\right)E_\omega \left(D_j\right)\}\right)+$$ $$+2\text{Re}(g_i\|g_j)_\omega ^+(E_\omega \left(D_i\right)XE_\omega ^{}\left(D_j\right)\frac{1}{2}\{X,E_\omega \left(D_i\right)E_\omega ^{}\left(D_j\right)\}))$$ (24) is a quantum Markovian generator. The structure map $`\theta _0(X)`$ has the standard form of the generator of a master equation $$\theta _0(X)=\mathrm{\Psi }(X)\frac{1}{2}\{\mathrm{\Psi }(1),X\}+i[H,X]$$ where $`\mathrm{\Psi }`$ is a completely positive map and $`H`$ is selfadjoint. In our case $`\mathrm{\Psi }(X)`$ is a linear combination of terms of the type $$E_\omega ^{}\left(D_i\right)XE_\omega \left(D_j\right)$$ Remark A corollary of the remark at the end of section 2 is that the Markovian generator $`\theta _0`$ maps $`L^{\mathrm{}}(H_S)^{}`$ into itself. Moreover, if $`X`$ in (3) belongs to the $`L^{\mathrm{}}(H_S)`$ then the Hamiltonian part of $`\theta _0(X)`$ vanishes and only the dissipative part remains. In particular, if $`H_S`$ has non–degenerate spectrum then $`\theta _0(X)`$ maps $`L^{\mathrm{}}(H_S)`$ and has the form $$\theta _0(X)=\underset{ij}{}\underset{\omega F}{}(2\text{Re}(g_i\|g_j)_\omega ^{}(E_\omega ^{}\left(D_i\right)XE_\omega \left(D_j\right)XE_\omega ^{}\left(D_i\right)E_\omega \left(D_j\right))+$$ $$+2\text{Re}(g_i\|g_j)_\omega ^+(E_\omega \left(D_i\right)XE_\omega ^{}\left(D_j\right)XE_\omega \left(D_i\right)E_\omega ^{}\left(D_j\right)))$$ for any $`XL^{\mathrm{}}(H_S)`$. The structure maps $`\theta _\alpha `$ in (21) satisfy the following stochastic Leibnitz rule, see the paper . Theorem. For any pair of operators in the system algebra $`X`$, $`Y`$, the structure maps in the Langevin equation (21) satisfy the equation $$\theta _\alpha (XY)=\theta _\alpha (X)Y+X\theta _\alpha (Y)+\underset{\beta ,\gamma }{}c_\alpha ^{\beta \gamma }\theta _\beta (X)\theta _\gamma (Y)$$ where the structure constants $`c_\alpha ^{\beta \gamma }`$ is given by the Ito table $$dM^\beta (t)dM^\gamma (t)=\underset{\alpha }{}c_\alpha ^{\beta \gamma }dM^\alpha (t)$$ The conjugation rules of $`dM^\alpha (t)`$ and $`\theta _\alpha `$ are connected in such a way that formula (21) defines a $``$–flow ($`j_t=j_t`$). ### 3.1 Evolution for the density matrix Let us now investigate the master equation for the density matrix $`\rho `$. We will show that if the reservoir is in the equilibrium state at temperature $`\beta ^1`$ then for the generic system with decoherence the solution of the master equation $`\rho (t)`$ with $`t\mathrm{}`$ tends to the classical Gibbs state with the same temperature $`\beta ^1`$. This phenomenon realizes the quantum measurement procedure — the quantum state (density matrix) collapses into the classical state. To show this we use the control of quantum decoherence that was found in the stochastic approximation of quantum theory, see and discussion below. Let us consider the evolution of the state (positive normed linear functional on system observables) given by the density matrix $`\rho `$, $`\rho (X)=\text{tr}\widehat{\rho }X`$. The evolution of the state is defined as follows $$\rho _t=j_t^{}(\rho )=\rho j_t$$ Therefore from (21) we get the evolution equation $$d\rho _t(X)=\rho dj_t(X)=\rho \underset{\alpha }{}j_t\theta _\alpha (X)dM^\alpha (t)=\underset{\alpha }{}\rho _t\left(\theta _\alpha (X)dM^\alpha (t)\right)$$ Only the stochastic differential $`dt`$ in this formula will survive and we get the master equation $$\frac{d}{dt}\rho _t(X)=\rho _t\theta _0(X)\theta _0^{}(\rho _t)(X)$$ (25) Let us consider the density matrix $`\widehat{\rho }=\widehat{\rho }_S\widehat{\rho }_R`$, $$\widehat{\rho }_{S,t}=\underset{\mu ,\nu }{}\rho (\mu ,\nu ,t)\|\mu \nu \|$$ where $`\|\mu `$, $`\|\nu `$ are eigenvectors of the system Hamiltonian $`H_S`$. Using the form (3) of $`\theta _0`$ and the identities $$\text{tr }Y[X,A]=\text{tr }[Y,A]X$$ $$\text{tr }Y\left(AXB\frac{1}{2}\{X,AB\}\right)=\text{tr }\left(BYA\frac{1}{2}\{Y,AB\}\right)X$$ the master equation (25) will take the form $$\underset{\mu ,\nu }{}\frac{d}{dt}\rho (\mu ,\nu ,t)\|\mu \nu \|=\underset{\mu ,\nu }{}\rho (\mu ,\nu ,t)$$ $$\underset{ij}{}\underset{\omega F}{}(i\text{ Im }(g_i\|g_j)_\omega ^{}[\|\mu \nu \|,E_\omega ^{}\left(D_i\right)E_\omega \left(D_j\right)]i\text{ Im }(g_i\|g_j)_\omega ^+[\|\mu \nu \|,E_\omega \left(D_i\right)E_\omega ^{}\left(D_j\right)]+$$ $$+2\text{Re}(g_i\|g_j)_\omega ^{}\left(E_\omega \left(D_j\right)\|\mu \nu \|E_\omega ^{}\left(D_i\right)\frac{1}{2}\{\|\mu \nu \|,E_\omega ^{}\left(D_i\right)E_\omega \left(D_j\right)\}\right)+$$ $$+2\text{Re}(g_i\|g_j)_\omega ^+(E_\omega ^{}\left(D_j\right)\|\mu \nu \|E_\omega \left(D_i\right)\frac{1}{2}\{\|\mu \nu \|,E_\omega \left(D_i\right)E_\omega ^{}\left(D_j\right)\}))=$$ $$=\underset{\mu ,\nu }{}\rho (\mu ,\nu ,t)\underset{ij}{}\underset{\omega F}{}(i\text{ Im }(g_i\|g_j)_\omega ^{}(\|\mu \nu \|\chi _\omega (\epsilon _\nu )D_i^{}P_{\epsilon _\nu \omega }D_jP_{\epsilon _\nu }\chi _\omega (\epsilon _\mu )P_{\epsilon _\mu }D_i^{}P_{\epsilon _\mu \omega }D_j\|\mu \nu \|)+$$ $$+2\text{Re}(g_i\|g_j)_\omega ^{}(\chi _\omega (\epsilon _\mu )\chi _\omega (\epsilon _\nu )P_{\epsilon _\mu \omega }D_j\|\mu \nu \|D_i^{}P_{\epsilon _\nu \omega }$$ $$\frac{1}{2}(\|\mu \nu \|\chi _\omega (\epsilon _\nu )D_i^{}P_{\epsilon _\nu \omega }D_jP_{\epsilon _\nu }+\chi _\omega (\epsilon _\mu )P_{\epsilon _\mu }D_i^{}P_{\epsilon _\mu \omega }D_j\|\mu \nu \|))$$ $$i\text{ Im }(g_i\|g_j)_\omega ^+\left(\|\mu \nu \|\chi _\omega (\epsilon _\nu +\omega )D_iP_{\epsilon _\nu +\omega }D_j^{}P_{\epsilon _\nu }\chi _\omega (\epsilon _\mu +\omega )P_{\epsilon _\mu }D_iP_{\epsilon _\mu +\omega }D_j^{}\|\mu \nu \|\right)+$$ $$+2\text{Re}(g_i\|g_j)_\omega ^+(\chi _\omega (\epsilon _\mu +\omega )\chi _\omega (\epsilon _\nu +\omega )P_{\epsilon _\mu +\omega }D_j^{}\|\mu \nu \|D_iP_{\epsilon _\nu +\omega }$$ $$\frac{1}{2}(\|\mu \nu \|\chi _\omega (\epsilon _\nu +\omega )D_iP_{\epsilon _\nu +\omega }D_j^{}P_{\epsilon _\nu }+\chi _\omega (\epsilon _\mu +\omega )P_{\epsilon _\mu }D_iP_{\epsilon _\mu +\omega }D_j^{}\|\mu \nu \|)))$$ (26) where $`\chi _\omega (\epsilon _\mu )=1`$ if $`\epsilon _\mu F_\omega `$ and equals to 0 otherwise. ## 4 Dynamics for generic systems Let us investigate the behavior of a system with dynamics defined by (3.1). This dynamics will depend on the Hamiltonian of the system. We will call the Hamiltonian $`H_S`$ generic, if: 1) The spectrum $`\text{ Spec }H_S`$ of the Hamiltonian is non degenerate. 2) For any Bohr frequency $`\omega `$ there exists a unique pair of energy levels $`\epsilon `$, $`\epsilon ^{}\text{ Spec }H_S`$ such that: $$\omega =\epsilon \epsilon ^{}$$ We investigate (3.1) for generic Hamiltonian. We also consider the case of one test function $`g_i(k)=g(k)`$, although this is not important. In this case $$E_\omega (X)=\|\sigma ^{}\sigma ^{}\|X\|\sigma \sigma \|=\|\sigma ^{}\sigma \|\sigma ^{}\|X\|\sigma $$ where $`\omega =\epsilon _\sigma \epsilon _\sigma ^{}`$. The Markovian generator $`\theta _0^{}`$ in (3.1) takes the form $$\theta _0^{}(X)=\underset{\sigma ,\sigma ^{}}{}\|\sigma ^{}\|D\|\sigma \|^2(i\text{ Im }(g\|g)_{\sigma \sigma ^{}}^{}[X,\|\sigma \sigma \|]+$$ $$+2\text{Re}(g\|g)_{\sigma \sigma ^{}}^{}\left(\|\sigma ^{}\sigma ^{}\|\sigma \|X\|\sigma \frac{1}{2}\{X,\|\sigma \sigma \|\}\right)$$ $$i\text{ Im }(g\|g)_{\sigma \sigma ^{}}^+[X,\|\sigma ^{}\sigma ^{}\|]+2\text{Re}(g\|g)_{\sigma \sigma ^{}}^+(\|\sigma \sigma \|\sigma ^{}\|X\|\sigma ^{}\frac{1}{2}\{X,\|\sigma ^{}\sigma ^{}\|\}))$$ (27) We use here the notion $$(g\|g)_{\mu \sigma }=(g\|g)_{\epsilon _\mu \epsilon _\sigma }$$ Notice that the factors $`\text{Re}(g\|g)_{\sigma \sigma ^{}}^\pm `$ are $`>0`$ only for $`\epsilon _\sigma >\epsilon _\sigma ^{}`$ and vanish for the opposite case. It is easy to see that the terms in (4) of the form $$\|\sigma \sigma \|\sigma ^{}\|X\|\sigma ^{}$$ for off–diagonal elements of the density matrix $`X=\|\mu \nu \|`$ are equal to zero. We will show that in such case the equation (3.1) will predict fast damping of the states of the kind $`\|\mu \nu \|`$. In the non–generic case one can expect the fast damping of the state $`\|\mu \nu \|`$ with different energies $`\epsilon _\mu `$ and $`\epsilon _\nu `$. With the given assumptions the action of $`\theta _0^{}`$ on the off–diagonal matrix unit $`\|\mu \nu \|`$, $`\epsilon _\mu \epsilon _\nu `$ is equal to $`A_{\mu \nu }\|\mu \nu \|`$ where the number $`A_{\mu \nu }`$ is given by the following $$A_{\mu \nu }=\underset{\sigma }{}(i\text{Im}(g\|g)_{\mu \sigma }^{}\|\sigma \|D\|\mu \|^2i\text{Im}(g\|g)_{\nu \sigma }^{}\|\sigma \|D\|\nu \|^2i\text{Im}(g\|g)_{\sigma \mu }^+\|\mu \|D\|\sigma \|^2+$$ $$+i\text{Im}(g\|g)_{\sigma \nu }^+\|\nu \|D\|\sigma \|^2\text{Re}(g\|g)_{\mu \sigma }^{}\|\sigma \|D\|\mu \|^2$$ $$\text{Re}(g\|g)_{\nu \sigma }^{}\|\sigma \|D\|\nu \|^2\text{Re}(g\|g)_{\sigma \mu }^+\|\mu \|D\|\sigma \|^2\text{Re}(g\|g)_{\sigma \nu }^+\|\nu \|D\|\sigma \|^2)$$ (28) The map $`\theta _0^{}`$ multiplies off–diagonal matrix elements of the density matrix $`\widehat{\rho }_S`$ by a number $`A_{\mu \nu }`$. Let us note that $$\text{Re}A_{\mu \nu }0$$ Moreover, for generic Hamiltonian the map $`\theta _0^{}`$ mixes diagonal elements of the density matrix but does not mix diagonal and off–diagonal elements (the action of $`\theta _0^{}`$ on diagonal element is equal to the linear combination of diagonal elements). The equation (3.1) for the generic case takes the form $$\underset{\mu ,\nu }{}\frac{d}{dt}\rho (\mu ,\nu ,t)\|\mu \nu \|=\underset{\mu \nu }{}A_{\mu \nu }\rho (\mu ,\nu ,t)\|\mu \nu \|+$$ $$+\underset{\sigma }{}\|\sigma \sigma \|\underset{\sigma ^{}}{}(\rho (\sigma ^{},t)(2\text{Re }(g\|g)_{\sigma ^{}\sigma }^{}\|\sigma \|D\|\sigma ^{}\|^2+2\text{Re }(g\|g)_{\sigma \sigma ^{}}^+\|\sigma ^{}\|D\|\sigma \|^2)$$ $$\rho (\sigma ,t)(2\text{Re }(g\|g)_{\sigma ^{}\sigma }^+\|\sigma \|D\|\sigma ^{}\|^2+2\text{Re }(g\|g)_{\sigma \sigma ^{}}^{}\|\sigma ^{}\|D\|\sigma \|^2))$$ (29) with $`A_{\mu \nu }`$ given by (4) and $`\rho (\sigma ,t)=\rho (\sigma ,\sigma ,t)`$. For instance we get $$j_t^{}(\|\mu \nu \|)=\mathrm{exp}(A_{\mu \nu }t)\|\mu \nu \|$$ We see that if any of $`\text{Re}(g\|g)\|\beta \|D\|\alpha \|^2`$ in (4) is non–zero then the corresponding off–diagonal matrix element of the density matrix decays. We obtain an effect of the diagonalization of the density matrix. This gives an effective criterium for quantum decoherence in the stochastic approximation: the system will exhibit decoherence if the constants $`\text{Re}(g\|g)^\pm `$ are non–zero. Now we estimate the velocity of decay of the density matrix $`\|\mu \nu \|`$ for a quantum system with $`N`$ particles. The eigenstate $`\|\mu `$ of the Hamiltonian of such a system can be considered as a tensor product over degrees of freedom of the system of some substates. Let us estimate from below the number of degrees of freedom of the system by the number of particles that belong to the system (for each particle we have few degrees of freedom). To get the estimate from below for the velocity of decay we assume that $`\|\sigma \|D\|\mu \|^2`$ in (4) is non–zero only if the state $`\sigma `$ differs from the state $`\mu `$ only for one degree of freedom. Then the summation over $`\omega `$ (or equivalently over $`\sigma `$) in (4) can be estimated by the summation over the degrees of freedom, or over particles belonging to the system. If we have total decoherence, i.e. all $`\text{Re}(g\|g)`$ are non–zero, then, taking all corresponding $`\|\sigma \|D\|\mu \|^2=1`$, we can estimate (4) as $`N\text{Re}(g\|g)`$, where $`N`$ is the number of particles in the system, or $$j_t^{}(\|\mu \nu \|)=\mathrm{exp}(N\text{Re}(g\|g)t)\|\mu \nu \|$$ (30) The off–diagonal element of the density matrix decays exponentially, with the exponent proportional to the number of particles in the system. Therefore for macroscopic (large $`N`$) systems with decoherence the quantum state will collapse into the classical state very quickly. This observation clarifies, why macroscopic quantum systems usually behave classically. The equation (30) describes such type of behavior, predicting that the quantum state damps at least as quickly as $`\mathrm{exp}(N\text{Re}(g\|g)t)`$. Therefore a macroscopic system (large $`N`$) will become classical in a time of order $`(N\text{Re}(g\|g))^1`$. Let us estimate the constant $`\text{Re}(g\|g)`$ for the equilibrium state of the reservoir with the temperature $`\beta ^1=kT`$. In this case $$\text{Re}(g\|g)_\omega ^{}=\frac{1}{\mathrm{}}\frac{\pi }{1e^{\beta \omega }}𝑑k\|g(k)\|^2\delta (\omega (k)\omega )$$ Taking $`g(k)=1`$ and using that the dispersion function $`\omega (k)`$ depends only on $`\|k\|`$ we get $$𝑑k\|g(k)\|^2\delta (\omega (k)\omega )=𝑑\mathrm{\Omega }_0^{\mathrm{}}𝑑\rho \delta (\omega (\rho )\omega )=4\pi _0^{\mathrm{}}𝑑\rho \delta (\omega (\rho )\omega )$$ where $`𝑑\mathrm{\Omega }`$ is the integration over angles. If we take the dispersion function $`\omega (k)=\|k\|`$, then for this integral we get $`4\pi \omega `$. Therefore for this choice of dispersion function we get $$\text{Re}(g\|g)_\omega ^{}=\frac{1}{\mathrm{}}\frac{4\pi ^2\omega }{1e^{\beta \omega }}$$ and analogously $$\text{Re}(g\|g)_\omega ^+=\frac{1}{\mathrm{}}\frac{4\pi ^2\omega }{e^{\beta \omega }1}$$ In the limit of small $`\beta `$ (or high temperatures) $`\beta ^1=kT>>\omega `$ both these integrals tend to $$4\pi ^2\frac{kT}{\mathrm{}}$$ Summing up, we get that the constant $`\text{Re}(g\|g)_\omega ^\pm `$ for the case of high temperature will be equal to $`\frac{kT}{\mathrm{}}`$ up to multiplication by a dimensionless constant depending on the model. This means that every degree of freeedom that energetically admissible ($`kT>>\omega `$) and not forbidden by the model ($`\|\nu \|D\|\mu \|^20`$) gives the term of order $`\frac{kT}{\mathrm{}}`$ in the exponent for dumping of off–diagonal matrix elements. The off–diagonal matrix element will dump as $`\mathrm{exp}(tN\frac{kT}{\mathrm{}})`$, where $`N`$ is the number of degrees of freedom (that for a generic system can be taken proportional to the number of particles). Off–diagonal matrix elements describe the quantum interference. Our result for the dumping of off–diagonal matrix elements (30) gives us a possibility to distinguish between microscopic systems (where quantum effects such as quantum interference are important) and macroscopic system which can be described by classical mechanics. The macroscopic system is a system satisfying $$N\frac{kT}{\mathrm{}}>>1$$ (31) Actually the value $`N\frac{kT}{\mathrm{}}`$ is of dimension of $`t^1`$, and (31) means that $`\left(N\frac{kT}{\mathrm{}}\right)^1`$ is much less than the time of observation. In the last section of the present paper we will illustrate the collapse phenomenon (30) using the quantum extension of the Glauber dynamics for a system of spins. We see that the stochastic limit predicts the collapse of a quantum state into a classical state and, moreover, allows us to estimate the velocity of the collapse (30). One can consider (30) as a more detailed formulation of the Fermi golden rule: the Fermi golden rule predicts exponential decay of quantum states; formula (30) also relates the speed of the decay to the dimensions (number of particles) of the sstem. Consider now the system density matrix $`\widehat{\rho }_S𝒞`$, where $`𝒞`$ is the algebra generated by the spectral projections of the system Hamiltonian $`H_S`$, and consider the master equation (4) (we consider the generic case). We will find that the evolution defined by this master equation will conserve the algebra $`𝒞`$ and therefore will be a classical evolution. We will show that this classical evolution in fact describes quantum phenomena. For $`\widehat{\rho }_{S,t}𝒞`$ we define the evolved density matrix of the system $$\widehat{\rho }_{S,t}=\underset{\sigma }{}\rho (\sigma ,t)\|\sigma \sigma \|$$ For this density matrix the master equation (25) takes the form $$\frac{d}{dt}\rho (\sigma ,t)=\underset{\sigma ^{}}{}(\rho (\sigma ^{},t)(2\text{Re }(g\|g)_{\sigma ^{}\sigma }^{}\|\sigma \|D\|\sigma ^{}\|^2+2\text{Re }(g\|g)_{\sigma \sigma ^{}}^+\|\sigma ^{}\|D\|\sigma \|^2)$$ $$\rho (\sigma ,t)(2\text{Re }(g\|g)_{\sigma ^{}\sigma }^+\|\sigma \|D\|\sigma ^{}\|^2+2\text{Re }(g\|g)_{\sigma \sigma ^{}}^{}\|\sigma ^{}\|D\|\sigma \|^2))$$ (32) Let us note that if $`\rho (\sigma ,t)`$ satisfies the detailed balance condition $$\rho (\sigma ,t)2\text{Re }(g\|g)_{\sigma \sigma ^{}}^{}=\rho (\sigma ^{},t)2\text{Re }(g\|g)_{\sigma \sigma ^{}}^+$$ (33) then $`\rho (\sigma ,t)`$ is the stationary solution for (4). Let us investigate (4), (33) for the equilibrium state of the field. In this case $$2\text{Re }(g\|g)_{\sigma \sigma ^{}}^{}=2\pi 𝑑k\|g(k)\|^2\delta (\omega (k)+\epsilon _\sigma ^{}\epsilon _\sigma )\frac{1}{1e^{\beta \omega (k)}}=$$ $$=2\pi 𝑑k\|g(k)\|^2\delta (\omega (k)+\epsilon _\sigma ^{}\epsilon _\sigma )\frac{1}{1e^{\beta (\epsilon _\sigma \epsilon _\sigma ^{})}}=\frac{C_{\sigma \sigma ^{}}}{1e^{\beta (\epsilon _\sigma \epsilon _\sigma ^{})}}$$ $$2\text{Re }(g\|g)_{\sigma \sigma ^{}}^+=\frac{C_{\sigma \sigma ^{}}}{e^{\beta (\epsilon _\sigma \epsilon _\sigma ^{})}1}$$ The equation (4) takes the form $$\frac{d}{dt}\rho (\sigma ,t)e^{\beta \epsilon _\sigma }=\underset{\sigma ^{}}{}\frac{C_{\sigma \sigma ^{}}\|\sigma ^{}\|D\|\sigma \|^2C_{\sigma ^{}\sigma }\|\sigma \|D\|\sigma ^{}\|^2}{1e^{\beta (\epsilon _\sigma \epsilon _\sigma ^{})}}\left(\rho (\sigma ^{},t)e^{\beta \epsilon _\sigma ^{}}\rho (\sigma ,t)e^{\beta \epsilon _\sigma }\right)$$ (34) Let us note that $`C_{\sigma \sigma ^{}}`$ are non–zero (and therefore positive) only if denominators in (34) are positive and $`C_{\sigma ^{}\sigma }`$ are non–zero only if the corresponding denominators are negative. If the system possesses decoherence then $`C_{\sigma \sigma ^{}}`$, $`C_{\sigma ^{}\sigma }`$ are non–zero and the solution of equation (34) for $`t\mathrm{}`$ tends to the stationary solution given by the detailed balance condition (33) $$\frac{\rho (\sigma ,t)}{1e^{\beta (\epsilon _\sigma \epsilon _\sigma ^{})}}=\frac{\rho (\sigma ^{},t)}{e^{\beta (\epsilon _\sigma \epsilon _\sigma ^{})}1}$$ or $$\rho (\sigma ,t)e^{\beta \epsilon _\sigma }=\rho (\sigma ^{},t)e^{\beta \epsilon _\sigma ^{}}$$ This means that the stationary solution (33) of (4) describes the equilibrium state of the system $$\rho (\sigma ,t)=\frac{e^{\beta \epsilon _\sigma }}{_\sigma ^{}e^{\beta \epsilon _\sigma ^{}}}$$ For a system with decoherence the density matrix will tend, as $`t\mathrm{}`$, to the stationary solution (33) of (4). In particular, as $`t\mathrm{}`$, the density matrix collapses to the classical Gibbs distribution. The phenomenon of a collapse of a quantum state into a classical state is connected with the quantum measurement procedure. The quantum uncertainty will be concentrated at the degrees of freedom of the quantum field and vanishes after the averaging procedure. One can speculate that the collapse of the wave function is a property of open quantum systems: we can observe the collapse of the wave function of the system averaging over the degrees of freedom of the reservoir interacting with the system. Usually the collapse of a wave function is interpreted as a projection onto a classical state (the von Neumann interpretation). The picture emerging from our considearations is more general: the collapse is a result of the unitary quantum evolution and conditional expectation (averaging over the degrees of freedom of quantum field). This is a generalization of the projection: it is easy to see that every projection $`P`$ generates a (non identity preserving) conditional expectation $`E_P(X)=PXP`$, more generally a set of projections $`P_i`$ generates the conditional expectation $$\underset{i}{}\alpha _iE_{P_i},\alpha _i0$$ but not every conditional expectation could be given in this way. We have found the effect of the collapse of density matrix for $`\rho (t)=U_t\rho U_t^{}`$, where $`U_t=lim_{\lambda 0}e^{itH_0}e^{itH}`$ is the stochastic limit of interacting evolution. The same effect of collapse will be valid for the limit of the full evolution $`e^{itH}`$, because the full evolution is the composition of interacting and free evolution. The free evolution leaves invariant the elements of diagonal subalgebra and multiplies the considered above nondiagonal element $`\|\sigma ^{}\sigma \|`$ by the oscillating factor $`e^{it(\epsilon _\sigma ^{}\epsilon _\sigma )}`$. Therefore for the full evolution we get the additional oscillating factor, and the collapse phenomenon will survive. ## 5 Control of coherence In this section we generalize the approach of and investigate different regimes of qualitative behavior for the considered model. The master equation (4) at first sight looks completely classical. In the present paper we derived this equation using quantum arguments. Now we will show that (4) in fact describes a quantum behavior. To show this we consider the following example. Let us rewrite (4) using the particular form (16) of $`(g\|g)^\pm `$. Using (15), (16) we get $$\frac{d}{dt}\rho (\sigma ,t)=\underset{\sigma ^{}}{}2\pi dk\|g(k)\|^2((N(k)+1)$$ $$\left(\rho (\sigma ^{},t)\delta (\omega (k)+\epsilon _\sigma \epsilon _\sigma ^{})\|\sigma \|D\|\sigma ^{}\|^2\rho (\sigma ,t)\delta (\omega (k)+\epsilon _\sigma ^{}\epsilon _\sigma )\|\sigma ^{}\|D\|\sigma \|^2\right)+$$ $$+N(k)(\rho (\sigma ^{},t)\delta (\omega (k)+\epsilon _\sigma ^{}\epsilon _\sigma )\|\sigma ^{}\|D\|\sigma \|^2\rho (\sigma ,t)\delta (\omega (k)+\epsilon _\sigma \epsilon _\sigma ^{})\|\sigma \|D\|\sigma ^{}\|^2))$$ (35) The first term (integrated with $`N(k)+1`$) on the RHS of this equation describes the emission of bosons and the second term (integrated with $`N(k)`$) describes the absorption of bosons. For the emission term the part with $`N(k)`$ describes the induced emission and the part with $`1`$ the spontaneous emission of bosons. Let us note that the Einstein relation for probabilities of emission and absorption of bosons with quantum number $`k`$ $$\frac{\text{probability of emission}}{\text{probability of absorption}}=\frac{N(k)+1}{N(k)}$$ is satisfied in the stochastic approximation. The formula (5) describes a macroscopic quantum effect. To show this let us take the spectrum of a system Hamiltonian (the set of system states $`\mathrm{\Sigma }=\{\sigma \}`$) as follows: let $`\mathrm{\Sigma }`$ contain two groups $`\mathrm{\Sigma }_1`$ and $`\mathrm{\Sigma }_2`$ of states with the energy gap between these groups (or, for simplicity, two states $`\sigma _1`$ and $`\sigma _2`$ with $`\epsilon _{\sigma _2}>\epsilon _{\sigma _1}`$). This type of Hamiltonian was considered in different models of quantum optics, see for review (for the case of two states we get the spin–boson Hamiltonian investigated in using the stochastic limit). Let the state $``$ of the bosonic field be taken in such a way that the density $`N(k)`$, of quanta of the bosonic field, has support in a set of momentum variables $`k`$ such that $$0<\omega (k)<\omega _0<\|\epsilon _{\sigma _1}\epsilon _{\sigma _2}\|,k\text{supp}N(k)$$ (36) This means that high–energetic bosons are absent. It is natural to consider the state $``$ as a sum of equilibrium state at temperature $`\beta ^1`$ and non–equilibrium part. Therefore the density $`N(k)`$ will be non–zero for small $`k`$ because the equilibrium state satisfies this property. Under the considered assumption (36) the integral of $`\delta `$–function $`\delta (\omega (k)+\epsilon _{\sigma _1}\epsilon _{\sigma _2})`$ with $`N(k)`$ in (5) equals identically to zero. Therefore the RHS of (5) will be equal to $$\underset{\sigma ^{}}{}2\pi dk\|g(k)\|^2(\rho (\sigma ^{},t)\delta (\omega (k)+\epsilon _\sigma \epsilon _\sigma ^{})\|\sigma \|D\|\sigma ^{}\|^2$$ $$\rho (\sigma ,t)\delta (\omega (k)+\epsilon _\sigma ^{}\epsilon _\sigma )\|\sigma ^{}\|D\|\sigma \|^2)$$ It is natural to consider this value (corresponding to the spontaneous emission of bosons by the system) as small with respect to the induced emission (for $`N(k)>>1`$). In this case the density matrix $`\rho (\sigma ,t)`$ will be almost constant in time. This is an effect of conservation of quantum coherence: in the absence of bosons with the energy $`\omega (k)`$ equal to $`\epsilon _{\sigma _1}\epsilon _{\sigma _2}`$ the system cannot jump between the states $`\sigma _1`$ and $`\sigma _2`$ (or, at least, this transition is very slow), because in the stochastic limit such jump corresponds to quantum white noise that must be on a mass shell. At the same time, the transitions between states inside the groups $`\mathrm{\Sigma }_1`$ and $`\mathrm{\Sigma }_2`$ are not forbidden by (36), because these transitions are connected with the soft bosons (with small $`k`$) that are present in the equilibrium part of $``$. In the above assumptions equation (5) describes the transition of the system to intermediate equilibrium, where the transitions between groups of states $`\mathrm{\Sigma }_1`$ and $`\mathrm{\Sigma }_2`$ are forbidden. If the state $``$ does not satisfy the property (36), then the system undergoes fast transitions between states $`\sigma _1`$ and $`\sigma _2`$. We can switch on such a transition by switching on the bosons with the frequency $`\omega (k)=\epsilon _{\sigma _2}\epsilon _{\sigma _1}`$. In conclusion: equation (5) describes a macroscopic quantum effect controlled by the distribution of bosons $`N(k)`$ which can be physically controlled for example by filtering. ## 6 The Glauber dynamics In the present section we apply the master equation (4) to the derivation of the quantum extension of the classical Glauber dynamics. The Glauber dynamics is a dynamics for a spin lattice with nearest neighbor interaction, see , . We will prove that the Glauber dynamics can be considered as a dynamics generated by the master equation of the type (4) derived from a stochastic limit for a quantum spin system interacting with a bosonic quantum field. We take the bosonic reservoir space $``$ corresponding to the bosonic equilibrium state at temperature $`\beta ^1`$. Thus the reservoir state is Gaussian with mean zero and correlations given by $$a^{}(k)a(k^{})=\frac{1}{e^{\beta \omega (k)}1}\delta (kk^{})$$ For simplicity we only consider the case of a one dimensional spin lattice, but our considerations extend without any change to multi–dimensional spin lattices. The spin variables are labeled by integer numbers $`Z`$, and, for each finite subset $`\mathrm{\Lambda }Z`$ with cardinality $`\|\mathrm{\Lambda }\|`$, the system Hilbert space is $$_S=_\mathrm{\Lambda }=_{r\mathrm{\Lambda }}C^2$$ and the system Hamiltonian has the form $$H_S=H_\mathrm{\Lambda }=\frac{1}{2}\underset{r,s\mathrm{\Lambda }}{}J_{rs}\sigma _r^z\sigma _s^z$$ where $`\sigma _r^x`$, $`\sigma _r^y`$, $`\sigma _r^z`$ are Pauli matrices $`(r\mathrm{\Lambda })`$ at the $`r`$-th site in the tensor product $$\sigma _r^i=1\mathrm{}1\sigma ^i1\mathrm{}1$$ For any $`r`$, $`s\mathrm{\Lambda }`$ $$J_{rs}=J_{sr}R,J_{rr}=0$$ We consider for simplicity the system Hamiltonian that describes the interaction of spin with the nearest neighbors (Ising model): $$J_{rs}=J_{r,r+1}$$ The interaction Hamiltonian $`H_I`$ (acting in $`_S`$) has the form $$H_I=\underset{r\mathrm{\Lambda }}{}\sigma _r^x\psi (g_r),\psi (g)=A(g)+A^{}(g),A(g)=𝑑k\overline{g}(k)a(k),$$ where $`\psi `$ is a field operator, $`A(g)`$ is a smeared quantum field with cutoff function (form factor) $`g(k)`$. The eigenvectors $`\|\sigma `$ of the system Hamiltonian $`H_S`$ can be labeled by spin configurations $`\sigma `$ (sequences of $`\pm 1`$), which label the natural basis in $`_S`$ consisting of tensor products of eigenvectors of $`\sigma _r^z`$ (spin up and spin down vectors $`\|\epsilon _r`$, corresponding to eigenvalues $`\epsilon _r=\pm 1`$) $$\|\sigma =_{r\mathrm{\Lambda }}\|\epsilon _r$$ In the present section we denote $`\epsilon _r`$ the energy of the spin at site $`r`$, and denote as $`E(\sigma )`$ the energy of the spin configuration $`\sigma `$ $$E(\sigma )=\frac{1}{2}\underset{r,s\mathrm{\Lambda }}{}J_{rs}\epsilon _r\epsilon _s$$ The action of the operator $`\sigma _r^x`$ on the spin configuration $`\sigma `$ is defined using the action of $`\sigma _r^x`$ on the corresponding eigenvector $`\|\sigma `$: so the operator $`\sigma _r^x`$ flips the spin at the $`r`$–th site in the sequence $`\sigma `$ (i.e. it maps the vector $`\|\epsilon _r`$ in the tensor product into the vector $`\|\epsilon _r`$). From the form of $`H_S`$ and $`H_I`$ it follows that, in (4), the matrix element $`\sigma \|D\|\sigma ^{}`$ of any two eigenvectors, corresponding to the spin configurations $`\sigma `$, $`\sigma ^{}`$, will be non–zero only if the configurations $`\sigma `$, $`\sigma ^{}`$ differ exactly at one site. If the configurations $`\sigma `$, $`\sigma ^{}`$ differ exactly at one site then $`\sigma \|D\|\sigma ^{}=1`$. The (classical) Glauber dynamics will be given by the master equation for the density matrix laying in the algebra of spectral projections of the system Hamiltonian (4) $$\frac{d}{dt}\rho (\sigma ,t)=\underset{r\mathrm{\Lambda }}{}(\rho (\sigma _r^x\sigma ,t)(2\text{Re }(g\|g)_{\sigma _r^x\sigma ,\sigma }^{}+2\text{Re }(g\|g)_{\sigma ,\sigma _r^x\sigma }^+)$$ $$\rho (\sigma ,t)(2\text{Re }(g\|g)_{\sigma _r^x\sigma ,\sigma }^++2\text{Re }(g\|g)_{\sigma ,\sigma _r^x\sigma }^{}))$$ (37) that gives the Glauber dynamics of a system of spins, see , . Here $$2\text{Re }(g\|g)_{\sigma ,\sigma _r^x\sigma }^{}=2\pi 𝑑k\|g(k)\|^2\delta (\omega (k)J_{r1,r}\epsilon _{r1}J_{r,r+1}\epsilon _{r+1})\frac{1}{1e^{\beta \omega (k)}}$$ (38) and analogously all the other $`(g\|g)^\pm `$. Up to now we have investigated the dynamics for the diagonal part of the density matrix. The master equation for the off–diagonal part of the density matrix (25) will give the quantum extension of the Glauber dynamics. We consider now this off–diagonal part: $$\underset{\mu \nu }{}\rho (\mu ,\nu ,t)\|\mu \nu \|$$ From (25), (4) we obtain the equation for the off–diagonal elements of the density matrix $$\frac{d}{dt}\rho (\mu ,\nu ,t)=A_{\mu \nu }\rho (\mu ,\nu ,t)$$ (39) $$A_{\mu \nu }=\underset{r\mathrm{\Lambda }}{}(i\text{Im}(g\|g)_{\mu ,\sigma _r^x\mu }^{}i\text{Im}(g\|g)_{\nu ,\sigma _r^x\nu }^{}i\text{Im}(g\|g)_{\sigma _r^x\mu ,\mu }^++i\text{Im}(g\|g)_{\sigma _r^x\nu ,\nu }^+$$ $$\text{Re}(g\|g)_{\mu ,\sigma _r^x\mu }^{}\text{Re}(g\|g)_{\nu ,\sigma _r^x\nu }^{}\text{Re}(g\|g)_{\sigma _r^x\mu ,\mu }^+\text{Re}(g\|g)_{\sigma _r^x\nu ,\nu }^+)$$ (40) Equations (6), (39), (6) describe the quantum extension of the classical Glauber dynamics (6). As it was already noted in section 4, the coefficient $`A_{\mu \nu }`$ in (6) is proportional to $`\|\mathrm{\Lambda }\|`$ (the number of particles in the system). Due to the summation on $`r\mathrm{\Lambda }`$ the coefficient $`A_{\mu \nu }`$ will diverge for large $`\|\mathrm{\Lambda }\|`$ (the real part of $`A_{\mu \nu }`$ will tend to $`\mathrm{}`$). Therefore the density matrix will collapse to the diagonal subalgebra (the classical distribution function) very quickly. Let us consider now the particular case of one dimensional system with translationally invariant Hamiltonian: $$J_{rs}=J_{r,r+1}=J>0$$ The translationally invariant Hamiltonian does not satisfy the generic non degeneracy conditions on the system spectrum that we have used in the derivation of equations (6), (39) and therefore we cannot apply these equations to describe the dynamics for this Hamiltonian. However in the translation invariant one–dimensional case we can investigate these equations by direct methods. In this case the $`(g\|g)^\pm `$, given by (38), are non–zero only if $`\epsilon _{r1}=\epsilon _{r+1}=1`$, and we get for (38) $$2\text{Re }(g\|g)_{\sigma ,\sigma _r^x\sigma }^{}=2\pi 𝑑k\|g(k)\|^2\delta (\omega (k)2J)\frac{1}{1e^{2\beta J}}=\frac{C}{1e^{2\beta J}}$$ (41) Therefore for one–dimensional translation invariant Hamiltonians we get for (6), compare with , $$\frac{d}{dt}\rho (\sigma ,t)=\frac{C}{1e^{2\beta J}}(\underset{r\mathrm{\Lambda };E(\sigma )>E(\sigma _r^x\sigma )}{}(e^{2\beta J}\rho (\sigma _r^x\sigma ,t)\rho (\sigma ,t))+$$ $$+\underset{r\mathrm{\Lambda };E(\sigma )<E(\sigma _r^x\sigma )}{}(\rho (\sigma _r^x\sigma ,t)e^{2\beta J}\rho (\sigma ,t)))$$ (42) The detailed balance stationary solution of (6) satisfy the following: for two spin configurations $`\sigma `$, $`\sigma _r^x\sigma `$ that differ by the flip of spin at site $`r`$ the energy of corresponding configurations differ by $`2J`$. The expectation $`\rho (\mu )`$, $`\mu =\sigma ,\sigma _r^x\sigma `$ of configuration with the higher energy will be $`e^{2\beta J}`$ times less. For the off–diagonal part of the density matrix for the case of one–dimensional translation invariant Hamiltonian the terms in the imaginary part of (6) cancel and using (41) we get for (6) $$A_{\mu \nu }=\underset{r\mathrm{\Lambda }}{}\left(\frac{2C}{1e^{2\beta J}}+\frac{2C}{e^{2\beta J}1}\right)=2C\underset{r\mathrm{\Lambda }}{}\frac{1+e^{2\beta J}}{1e^{2\beta J}}$$ This sum, over $`r`$, of equal terms diverges with $`\|\mathrm{\Lambda }\|\mathrm{}`$. Therefore the off–diagonal elements of the density matrix that satisfy (39) will decay very quickly and for sufficiently large $`t`$, $`\|\mathrm{\Lambda }\|`$ the dynamics of the system will be given by the classical Glauber dynamics. For the master equation considered above we used the master equation for generic (non–degenerate) Hamiltonian. This gives us the Glauber dynamics. But the translation invariant Hamiltonian is degenerate. Therefore in the translation invariant case we will get some generalization of the Glauber dynamics. To derive this generalization let us consider the general form (3.1) of the master equation. For the considered spin system this gives $$\underset{\mu ,\nu }{}\frac{d}{dt}\rho (\mu ,\nu ,t)\|\mu \nu \|=\underset{\mu ,\nu }{}\rho (\mu ,\nu ,t)\underset{a,b\mathrm{\Lambda }}{}(i(\|\mu P_{E(\nu )}\sigma _a^xP_\nu \sigma _b^x\nu \|\|P_{E(\mu )}\sigma _a^xP_\mu \sigma _b^x\mu \nu \|)+$$ $$+\frac{C}{1e^{2\beta J}}\left(\|P_{E(\mu )J}\sigma _a^x\mu P_{E(\nu )J}\sigma _b^x\nu \|\frac{1}{2}\left(\|\mu P_{E(\nu )}\sigma _a^xP_{E(\nu )J}\sigma _b^x\nu \|+\|P_{E(\mu )}\sigma _a^xP_{E(\mu )J}\sigma _b^x\mu \nu \|\right)\right)$$ $$i\left(\|\mu P_{E(\nu )}\sigma _a^xP_{\nu +}\sigma _b^x\nu \|\|P_{E(\mu )}\sigma _a^xP_{\mu +}\sigma _b^x\mu \nu \|\right)+$$ $$+\frac{C}{e^{2\beta J}1}(\|P_{E(\mu )+J}\sigma _a^x\mu P_{E(\nu )+J}\sigma _b^x\nu \|\frac{1}{2}(\|\mu P_{E(\nu )}\sigma _a^xP_{E(\nu )+J}\sigma _b^x\nu \|+\|P_{E(\mu )}\sigma _a^xP_{E(\mu )+J}\sigma _b^x\mu \nu \|)))$$ (43) Here $`C`$ is given by (41), operator $`P_{E(\mu )}`$ is a projector onto the states with the energy $`E(\mu )`$, operator $`P_\nu `$ is given by $$P_\nu =\text{ Im }(g\|g)_1^{}P_{E(\nu )J}+\text{ Im }(g\|g)_0^{}P_{E(\nu )}+\text{ Im }(g\|g)_1^{}P_{E(\nu )+J}$$ $$\text{ Im }(g\|g)_a^{}=\text{P.P.}𝑑k\|g(k)\|^2\frac{1}{\omega (k)+aJ}\frac{1}{1e^{\beta \omega (k)}},a=1,0,1$$ For the operator $`P_{\nu +}`$ we get the analogous expression $$P_{\nu +}=\text{ Im }(g\|g)_1^+P_{E(\nu )J}+\text{ Im }(g\|g)_0^+P_{E(\nu )}+\text{ Im }(g\|g)_1^+P_{E(\nu )+J}$$ with the coefficients $`(g\|g)_a^+`$: $$(g\|g)_a^+=\text{P.P.}𝑑k\|g(k)\|^2\frac{1}{\omega (k)aJ}\frac{1}{e^{\beta \omega (k)}1},a=1,0,1$$ The equation (6) gives the quantum generalization of the Glauber dynamics. The matrix elements $`\rho (\mu ,\nu ,t)`$ of the density matrix corresponding to the states $`\mu `$, $`\nu `$ with different energies will decay quickly. But for the translation invariant Hamiltonian there exist different $`\mu `$, $`\nu `$ with equal energies. Corresponding matrix element will decay with the same speed as the diagonal elements of the density matrix. Moreover one can expect non–ergodic behavior for this model. Therefore the generalization (6) of the Glauber dynamics is non–trivial. ### 6.1 Evolution for subalgebra of local operators In this section to compare with the results of we consider the dynamics of spin systems, described in the previous section, for Hamiltonian with non necessarily finite set of spins $`\mathrm{\Lambda }`$ but for local observable $`X`$. The observable $`X`$ is local if it belongs to the local algebra, that is UHF–algebra (uniformly hyperfinite algebra) $$𝒜=\underset{\mathrm{\Lambda }\text{ is finite }}{}𝒜_\mathrm{\Lambda }$$ where $`𝒜_\mathrm{\Lambda }`$ is the $``$–algebra generated by the elements $$_iX_i,X_i=1\text{ for }i\mathrm{\Lambda }$$ Consider now the action of $`\theta _0`$ on local $`X`$ $$\theta _0(X)=\underset{ij}{}\underset{\omega F}{}(i\text{ Im }(g_i\|g_j)_\omega ^{}[X,E_\omega ^{}\left(D_i\right)E_\omega \left(D_j\right)]+i\text{ Im }(g_i\|g_j)_\omega ^+[X,E_\omega \left(D_i\right)E_\omega ^{}\left(D_j\right)]+$$ $$+2\text{Re}(g_i\|g_j)_\omega ^{}\left(E_\omega ^{}\left(D_i\right)XE_\omega \left(D_j\right)\frac{1}{2}\{X,E_\omega ^{}\left(D_i\right)E_\omega \left(D_j\right)\}\right)+$$ $$+2\text{Re}(g_i\|g_j)_\omega ^+(E_\omega \left(D_i\right)XE_\omega ^{}\left(D_j\right)\frac{1}{2}\{X,E_\omega \left(D_i\right)E_\omega ^{}\left(D_j\right)\}))$$ (44) For $`E_\omega \left(D_i\right)`$ we get $$E_\omega \left(D_i\right)=\underset{E(r)F_\omega }{}P_{E(r)\omega }D_iP_{E(r)}=$$ $$=1\|\epsilon _{i1}\epsilon _{i1}\|\|\epsilon _i\epsilon _i\|\|\epsilon _{i+1}\epsilon _{i+1}\|1+1\|\epsilon _{i1}\epsilon _{i1}\|\|\epsilon _i\epsilon _i\|\|\epsilon _{i+1}\epsilon _{i+1}\|1;$$ (45) with the frequency $`\omega `$ of the following form $$\omega =J_{i1,i}\epsilon _{i1}+J_{i,i+1}\epsilon _{i+1}$$ Therefore the operator $`E_\omega `$ given by (6.1) is local and moreover, corresponding map $`\theta _0`$ given by (6.1) maps $`𝒜`$ into itself. The formula (6.1) explains the physical meaning of the operator $`E_\omega \left(D_i\right)`$. For positive $`\omega `$ this it flips the spin at site $`i`$ along the direction of the mean field of its neighbors (for negative $`\omega `$ it flips the same spin into the opposite direction). Acknowledgments Sergei Kozyrev is grateful to Luigi Accardi and Centro Vito Volterra where this work was done for kind hospitality. The authors are grateful to I.V.Volovich for discussions. This work was partially supported by INTAS 99–00545 grant. Sergei Kozyrev was partially supported by RFFI 990100866 grant.
warning/0005/hep-th0005135.html	ar5iv	text	# 1 Introduction ## 1 Introduction The irreducible representations of the four-dimensional superconformal group $`SU(2,2\|N)`$ have been much studied in the literature . The theory of such representations has assumed some significance recently in the light of the Maldacena conjecture relating IIB supergravity and string theory on backgrounds with an $`AdS_5`$ factor to four-dimensional superconformal field theories on the boundary. A particularly important class of operators that can arise consists of those operators which correspond to short representations of the superconformal group since these are expected to be protected from quantum corrections and thus not acquire anomalous dimensions. (See, for example, ). One way of constructing representations explicitly is the method of parabolic induction. In the context of superconformal groups this was discussed in some detail in , although a direct comparison with the more algebraic group-theoretic results of was not made at the time. However, it was shown how certain ultra-short representations, which correspond to on-shell massless supermultiplets and which can be described by constrained superfields in ordinary $`N`$-extended Minkowski superspace , can be described by “generalised chiral” or “Grassmann analytic” (G-analytic) superfields in harmonic superspaces . A simple example of a composite operator - the $`N=4`$ supercurrent - was also constructed as the trace of the square of the $`N=4`$ Yang-Mills field-strength superfield. In it was shown that there is in fact a series of such G-analytic superfields obtained by simply taking the trace of powers of the field strength, and it was conjectured that this circumstance might lead to restrictions on the correlations functions of such operators, an idea that has been exploited using both $`N=4`$ and $`N=2`$ superspaces in refs . These composites all have maximum spin 2 and it was subsequently realised in that the superconformal representations which are realised this way are dual to the bulk Kaluza-Klein representations which arise in the compactification of IIB supergravity on $`AdS_5\times S^5`$ . In the case of an Abelian field strength this construction can easily be generalised to the case of the $`D=6,(2,0)`$ tensor multiplet and the $`D=3,N=8`$ scalar multiplet, these being the multiplets which are relevant to M-theory compactified on $`AdS_4\times S^7`$ and $`AdS_7\times S^4`$ respectively (for other discussions of $`D=6,(2,0)`$ representations see ). In a recent series of papers short representations in $`D=4`$ , $`D=6`$ and $`D=3`$ have been discussed in the context of harmonic superspaces which have the form of super Minkowski space times an internal coset space with isotropy group equal to the maximal torus in the internal symmetry group. As well as reproducing the KK multiplets the authors of these papers were able to obtain many other short representations by taking products of the underlying massless multiplets and then imposing harmonic analyticity on the outcome. In the current paper we adopt a more general approach which makes use of all of the harmonic superspaces which are compatible with superconformal symmetry. We first discuss the short representations directly using the method of parabolic induction. To do this it is necessary to complexify spacetime and the superconformal group since the superspaces we are interested in are not cosets of the real superconformal group. One advantage of this approach is that superconformal symmetry is manifest throughout. The short representations act on superfields which depend on fewer odd coordinates than are present in $`N`$-extended super Minkowski space; such G-analytic superfields obey constraints which generalise the constraints satisfied by chiral superfields. We show that each short representation generically admits realisations on different homogeneous superspaces which incorporate different numbers of G-analyticity constraints. Amongst these realisations there are those which are “maximally efficient” in that the maximum number of G-analyticity constraints compatible with the representation is explicitly incorporated. It will be seen that the number of G-analyticity constraints of a given realisation places constraints on the geometry of the internal manifold of the associated harmonic superspace. The fewer G-analyticity constraints there are (equivalently, the larger the number of odd coordinates of the superspace concerned) the larger the number of representations that will be compatible with it. Thus, G-analyticity of type (1,1) is compatible with all representations of series C while G-analyticity of type (1,0) (or (0,1)) is compatible with all short representations (see, for example, for discussions of these representations). After a brief review of homogeneous superspaces of the superconformal group and parabolic induction, we discuss the short representations in detail in section 3. In section 4, we discuss on-shell massless supermultiplets from this viewpoint; these are special superconformal fields in that they carry ultra-short representations. Each of these massless supermultiplets can be realised in many different ways in general and this allows us to construct general representations out of them by taking products on appropriate superspaces. This approach has the virtue of not requiring the projection step of \- the product superfields are naturally superconformal fields on the appropriate superspace. In particular, in non-Abelian $`N=4`$ supersymmetric Yang-Mills theory the single, double and triple trace gauge-invariant composite operators can be constructed straightforwardly as G-analytic and harmonic analytic superfields by making use of the three different realisations of the underlying field strength multiplet. Aspects of multiple trace operators, which are believed to correspond to multiparticle states in the AdS/CFT context, have been studied in and in harmonic superspace in . The paper concludes in section 5 with an explicit discussion of superconformal transformations acting on shortened supermultiplets defined as analytic superfields on real harmonic superspaces. ## 2 Preliminaries ### 2.1 Harmonic superspace Harmonic superspaces are usually taken to be superspaces of the form $`M\times F`$ where $`M`$ is super Minkowski space and $`F=H\backslash G`$ is a coset space of the internal symmetry group $`G`$ with isotropy group $`H`$ . Usually $`F`$ is a compact, complex manifold; in $`N`$-extended supersymmetry in four dimensions it can be taken to be of the form $`\text{𝔽}_𝐤(N)=S\left(U(k_1)\times U(k_2k_1)\mathrm{}U(Nk_{\mathrm{}})\right)\backslash SU(N)`$ where $`𝐤:=(k_1,\mathrm{}k_{\mathrm{}})`$ is a sequence of positive integers with $`k_1<k_2\mathrm{}k_{\mathrm{}}<N`$ . The space $`\text{𝔽}_𝐤(N)`$ is a flag manifold: a point in $`\text{𝔽}_𝐤(N)`$ is a flag in $`\text{}^N`$, i.e. a sequence of subspaces $`\{V_i\}`$, $`dimV_i=k_i`$, such that $`V_1V_2\mathrm{}V_{\mathrm{}}\text{}^N`$. Special cases of flag manifolds include projective spaces and Grassmannians. A complete flag is one with $`𝐤=(1,2,3,\mathrm{}N1)`$. The space of complete flags, which we shall denote $`\text{𝔽}_B`$, is $`U(1)^{N1}\backslash SU(N)`$; it can also be thought of as the coset space $`B\backslash SL(N,\text{})`$ where $`B`$ is the Borel subgroup of $`SL(N)`$. In the GIKOS formalism one works with fields defined on $`M\times G`$ which are equivariant with respect to $`H`$. That is, one has fields $`A(z,u),zM,uG`$ which take their values in some representation space of $`H`$ and which satisfy $`A(z,hu)=R(h)A(z,u)`$, $`R`$ being the representation of $`H`$ in question. Such fields are equivalent to sections of vector bundles over harmonic superspace. In index notation we write $`u_I^i`$ where the upper case index indicates a representation of $`H`$ and the lower case index the defining representation of $`G`$, and where both indices run from $`1`$ to $`N`$. The coordinates of Minkowski superspace are denoted by $`(x^{\alpha \dot{\alpha }},\theta ^{\alpha i},\overline{\theta }_i^{\dot{\alpha }})`$ where $`\alpha `$ and $`\dot{\alpha }`$ are two-component spinor indices. The covariant derivatives $`(_{\alpha \dot{\alpha }},D_{\alpha i},\overline{D}_{\dot{\alpha }}^i)`$ satisfy $$[D_{\alpha i},\overline{D}_{\dot{\alpha }}^j]=i\delta _i{}_{}{}^{j}_{\alpha \dot{\alpha }}^{}$$ (1) with all other (anti)-commutators being zero. With the aid of $`u`$ we can convert $`G`$ indices to $`H`$ indices; this means that we can single out subsets of $`D`$’s and $`\overline{D}`$’s which anticommute with each other and which can therefore consistently be taken to annihilate superfields. Let us suppose that the index $`I`$ splits as follows under $`H`$: $`I=(r,r^{\prime \prime },r^{})`$, where $`r=1,\mathrm{}p,r^{\prime \prime }=p+1,\mathrm{}Nq,r^{}=Nq+1,\mathrm{}N`$, and where $`(p+q)N`$, then we can have superfields $`A`$ which obey the constraints $$D_{\alpha r}A=\overline{D}_{\dot{\alpha }}^r^{}A=0$$ (2) where $$D_{\alpha r}:=u_r{}_{}{}^{i}D_{\alpha i}^{};\overline{D}_{\dot{\alpha }}^r^{}=\overline{D}_{\dot{\alpha }}^i(u^1)_i^r^{}$$ (3) Constraints such as these are called Grassmann analyticity (G-analyticity) conditions following ; they generalise the notion of chirality which was first introduced in the context of $`N=1`$ supersymmetry in . We shall refer to the G-analyticity conditions of equation (2) as being of type $`(p,q)`$. The space of G-analytic structures of type $`(p,q)`$ on $`N`$-extended Minkowski superspace (in four dimensions) is the flag manifold $`\text{𝔽}_{p,Nq}(N)=S(U(p)\times U(N(p+q))\times U(q))\backslash SU(N)`$, and the harmonic superspace with this internal coset is referred to as $`(N,p,q)`$ harmonic superspace. Later on we shall make a refinement of this definition. The derivatives on $`G`$ can be taken to be the right-invariant vector fields, $`D_I{}_{}{}^{J},D_I{}_{}{}^{I}=0`$. They act as follows: $$D_I{}_{}{}^{J}u_{K}^{}{}_{}{}^{k}=\delta _K{}_{}{}^{J}u_{I}^{}{}_{}{}^{k}\frac{1}{N}\delta _I{}_{}{}^{J}u_{K}^{}^k$$ (4) and they obey the reality condition $$\overline{D}^J{}_{I}{}^{}=D_I^J$$ (5) These derivatives divide into three sets: the set $`\{D_r{}_{}{}^{s^{}},D_r{}_{}{}^{s^{\prime \prime }},D_{r^{\prime \prime }}{}_{}{}^{s^{}}\}`$, corresponding to the $`\overline{}`$ operator on the complex coset space (i.e. holomorphic functions are annihilated by $`\overline{}`$), the complex conjugate set $`\{D_r^{}{}_{}{}^{s},D_{r^{\prime \prime }}{}_{}{}^{s},D_r^{}{}_{}{}^{s^{\prime \prime }}\}`$ which corresponds to the $``$ operator, and the set $`\{D_r{}_{}{}^{s},D_r^{}{}_{}{}^{s^{}},D_{r^{\prime \prime }}{}_{}{}^{s^{\prime \prime }}\}`$ which corresponds to the isotropy algebra. As well as imposing G-analyticity we can impose harmonic analyticity (H-analyticity) by demanding that a superfield be annihilated by the first of these sets of derivatives. Indeed, since the algebra of all of these operators and the spinorial derivatives defining G-analyticity closes, it is consistent to impose both G-analyticity and H-analyticity on the same superfield. Such superfields will be referred to as CR-analytic, or simply analytic. It was observed in that analyticity conditions such as these define a generalisation of a complex structure, known as a CR structure, on harmonic superspace. In this language, a G-analytic structure can be considered to be a CR-structure which is purely odd. ### 2.2 Harmonic conjugation It was noted in the first paper on harmonic superspace that it is necessary to use a modified conjugation in order to handle real representations of $`SU(2)`$ in a straightforward manner. This conjugation combines the antipodal map on the internal manifold $`\text{}P^1`$ with ordinary complex conjugation. This was generalised to $`(N,p,q)=(3,1,1)`$ in and to any $`(N,p,p)`$ space in . In a general $`(N,p,q)`$ superspace, with internal manifold $`\text{𝔽}_{p,Nq}(N)`$ the analogue of the antipodal map becomes a map from $`\text{𝔽}_{p,Nq}(N)`$ to $`\text{𝔽}_{q,Np}(N)`$; if we combine this transformation with complex conjugation we obtain a map which maps the CR-structure on $`(N,p,q)`$ harmonic superspace to the CR structure on $`(N,q,p)`$ superspace. We shall refer to this as harmonic conjugation. Working on $`M\times SU(N)`$ we write $`u_I^i`$ as before as $`(u_r{}_{}{}^{i},u_{r^{\prime \prime }}{}_{}{}^{i},u_r^{}{}_{}{}^{i})`$ with $`r=1,\mathrm{}p`$ and $`r^{}=Nq+1,\mathrm{}N`$. This is the appropriate splitting of $`I`$ for $`(N,p,q)`$ superspace. We denote the splitting relevant to $`(N,q,p)`$ by $`I=(R,R^{\prime \prime },R^{})`$ with $`R=1,\mathrm{}q,R^{}=Np+1,\mathrm{}N`$. The mapping we start with is given by $$uku$$ (6) where $`k`$ is the matrix $$k=\left(\begin{array}{ccc}0& 0& 1_q\\ 0& 1_r& 0\\ 1_p& 0& 0\end{array}\right)$$ (7) It is easy to check that $`khk^1S(U(q)\times U(N(p+q))\times U(p))`$ provided that $`hS(U(p)\times U(N(p+q))\times U(q))`$, so that this does indeed define a map from one coset space to the other, i.e. from $`\text{𝔽}_{p,Nq}(N)`$ to $`\text{𝔽}_{q,Np}(N)`$. If we combine this with complex conjugation we find $`u_r^i`$ $``$ $`(u^1)_i^R^{}`$ $`u_{r^{\prime \prime }}^i`$ $``$ $`(u^1)_i^{R^{\prime \prime }}`$ $`u_r^{}^i`$ $``$ $`(u^1)_i^R`$ (8) For the inverses we have $`(u^1)_i^r`$ $``$ $`u_R^{}^i`$ $`(u^1)_i^{r^{\prime \prime }}`$ $``$ $`u_{R^{\prime \prime }}^i`$ $`(u^1)_i^r^{}`$ $``$ $`u_R^i`$ (9) We then find $`D_{\alpha r}`$ $``$ $`\overline{D}_{\dot{\alpha }}^R^{}`$ $`\overline{D}_{\dot{\alpha }}^r^{}`$ $``$ $`D_{\alpha R}`$ (10) while $`D_r^s^{}`$ $``$ $`D_S^R^{}`$ $`D_r^{s^{\prime \prime }}`$ $``$ $`D_{S^{\prime \prime }}^R^{}`$ $`D_{r^{\prime \prime }}^s^{}`$ $``$ $`D_S^{R^{\prime \prime }}`$ (11) so that the derivatives of the $`(N,p,q)`$ CR-structure do indeed transform into those of the $`(N,q,p)`$ CR-structure, and therefore analytic fields on $`(N,p,q)`$ superspace transform into analytic fields on $`(N,q,p)`$ superspace. ### 2.3 Complex superspaces It is the analytic superfields defined on harmonic superspaces of the above type, or related ones, that we shall be interested in for the description of short representations of the superconformal group. However, from a group-theoretical point of view, it is convenient to take a slightly different approach to harmonic superspaces, namely as coset spaces of the superconformal group. A slight technical problem presents itself here, because the spaces we are ultimately interested in, i.e. the so-called analytic superspaces where the G-analyticity constraints are automatically solved, are not coset spaces of the real superconformal group $`SU(2,2\|N)`$; we can circumvent this problem by complexifying both the group and spacetime. The fields we are interested in will then be holomorphic fields defined on coset spaces of the complexified superconformal group $`SL(4\|N)`$ (we omit the which is always understood). The isotropy groups will be parabolic subgroups, which means that the superspaces themselves will be flag supermanifolds, that is, spaces whose points correspond to nested sequences of sub-supervector spaces of $`\text{}^{4\|N}`$. The internal parts of these spaces will, however, be the same as in the real case above, the only difference being that we treat them exclusively as complex manifolds. When we restrict spacetime and the odd coordinates to be real we recover analytic fields of the type we have discussed above on $`(N,p,q)`$ harmonic superspaces. The virtue of this approach is that superconformal covariance is manifest. We begin by considering complexified Minkowski space which can be viewed as an open subset of the coset space $`P\backslash SL(4)`$, $`SL(4)`$ being the complexified conformal group and $`P`$ the subgroup of matrices of the following shape: $$\left(\begin{array}{cccc}& & & \\ & & & \\ & & & \\ & & & \end{array}\right)$$ (12) where the bullets denote elements which do not have to be zero. The blank region can be thought of as corresponding to spacetime. Indeed, we can choose a coset representative of the form $$Mxs(x)=\left(\begin{array}{cc}1_2& x\\ 0_2& 1_2\end{array}\right)$$ (13) where each entry is a two-by-two matrix. From this one can easily work out that the transformation of $`x`$ under the conformal group is the usual one by using standard homogeneous space techniques. Complexified super Minkowski space has the form $`P\backslash SL(4\|N)`$ where $`P`$ consists of matrices of the form $$\left(\begin{array}{ccccccc}& & & & & & \\ & & & & & & \\ & & & & & .& \\ & & & & & .& \\ & & & & & & \\ & & & & & .& \\ & & & & .& .& .\\ & & & & & .& \end{array}\right)$$ (14) The standard coset representative is $$Mzs(z)=\left(\begin{array}{ccccc}1_2& x& & \theta & \\ 0& 1_2& & 0& \\ & & & & \\ & & & & \\ 0& \phi & & 1_N& \end{array}\right)$$ (15) where $`\phi `$ denotes the $`N`$ dotted spinorial coordinates which become the complex conjugates of the $`\theta `$’s in the real case. Complexified $`(N,p,q)`$ harmonic superspace is given by the following subgroup (16) Locally, this space has the form of complex super Minkowski space times the internal flag manifold $`\text{𝔽}_{p,Nq}(N)`$. The related $`(N,p,q)`$ analytic superspace has the same body but fewer odd coordinates; the relevant subgroup is (17) A CR-analytic field defined on real $`(N,p,q)`$ harmonic superspace, when continued to complexified superspace, will be a holomorphic field defined on the above $`(N,p,q)`$ analytic superspace. ### 2.4 Parabolic induction If $`G`$ is a Lie group and $`H`$ a subgroup, representations of $`G`$ can be induced from representations of $`H`$ by considering $`H`$-vector bundles over $`H\backslash G`$. Sections of such bundles are, as we mentioned previously, equivalent to equivariant maps $`F:GV`$, where $`V`$ is the representation space of $`H`$ with representation $`R`$, say. Thus we have (for $`uG`$) $$F(hu)=R(h)F(u)$$ (18) The induced representation itself is given by $`FgF,gG`$ where $$(gF)(u)=F(ug)$$ (19) It is easy to see that this does define a (left) representation of $`G`$ carried by the space of sections of the vector bundle in question. Parabolic induction refers to the case where the subgroup is parabolic. In this case, if the representation $`R`$ is irreducible then the representation of $`G`$ constructed in this manner is also irreducible. (For a detailed discussion of parabolic subgroups and flag manifolds in the context of twistor theory see .) Let G be a complex, simple Lie group and let g be its Lie algebra. A Borel subalgebra b is a maximal solvable subalgebra of g and a parabolic subalgebra is one which contains a Borel subalgebra. For the case of $`\text{s}\text{l}(N)`$ we may take the Borel subalgebra to be the algebra of all lower triangular matrices (with non-zero entries on the diagonal allowed), and a parabolic subalgebra p is one which is block lower triangular. Thus it consists of matrices of the form: (20) where, if the squares have sides $`k_1,k_2,\mathrm{}k_{\mathrm{}}`$, the parabolic specifies the flag manifold $`\text{𝔽}_𝐤(N)`$, where $`𝐤=(k_1,k_2,\mathrm{}k_{\mathrm{}})`$. The parabolic p can be represented by placing a cross on each of the nodes $`k_1,k_2,\mathrm{}k_{\mathrm{}}`$ of the Dynkin diagram for $`\text{s}\text{l}(N)`$. For example, the flag manifold $`\text{𝔽}_{1,3}(4)`$, as well as the corresponding parabolic subalgebra, is represented by the $`\text{s}\text{l}(4)`$ Dynkin diagram . Thus we have two constructions of flag manifolds - as homogeneous spaces of $`SU(N)`$ with isotropy groups of the type discussed previously, which makes it clear that they are compact, and as homogeneous spaces of $`SL(N)`$ defined by parabolic isotropy groups, which makes it clear that they are complex manifolds. As we mentioned previously, the Borel subgroup defines the space of flags of type $`𝐤=(1,2,3,\mathrm{}(N1))`$; it is the same space as the homogeneous space of $`SU(N)`$ which has the maximal torus as its isotropy group. The sub-algebra of p consisting of the blocks along the diagonal is called the Levi sub-algebra l and we have $`\text{p}=\text{l}\text{u}`$ where u consists of the remaining matrices; clearly u is nilpotent, $`\text{u}_k=0`$ for some $`k`$ where $`\text{u}_k`$ is defined iteratively by $`\text{u}_k=[\text{u}_{k1},\text{u}_{k1}];\text{u}_0=\text{u}`$. For a general flag manifold $`\text{𝔽}_𝐤(N)`$, the Levi subalgebra is $$\text{l}=\text{s}\text{l}(k_1)\text{s}\text{l}(k_2k_1)\mathrm{}\text{s}\text{l}(Nk_{\mathrm{}})\text{}^{\mathrm{}}$$ (21) The finite-dimensional irreducible representations of $`\text{s}\text{l}(N)`$ are specified by highest weights $`\lambda `$. These representations can alternatively be constructed by parabolic induction. In fact they act on holomorphic sections of homogeneous vector bundles over the flag manifolds $`P\backslash SL(N)`$. This procedure can be described explicitly. Choose a representation with highest weight $`\lambda `$ and fix a parabolic p. Any weight for $`\text{s}\text{l}(N)`$ is also a weight for any p and so highest weights determine irreducible representations of p in some vector space $`V_\lambda `$. In fact, since u is nilpotent, it acts trivially on such a representation, so that we actually have a representation of the Levi subalgebra of p on $`V_\lambda `$. We then construct the vector bundle $`𝒱_\lambda `$ over $`P\backslash SL(N)`$, i.e. the bundle whose standard fibre is the $`P`$-representation space $`V_\lambda `$. A simple application of the Bott-Borel-Weil theorem gives the desired result: the space of holomorphic sections of $`𝒱_\lambda `$ is isomorphic to the representation space of $`\text{s}\text{l}(N)`$ determined by the highest weight $`\lambda `$. Note that a given representation can be presented in this fashion in as many different ways as there are parabolic subalgebras of $`\text{s}\text{l}(N)`$. ( A discussion of this theorem in the context of twistor theory is given in ). This discussion can be made more concrete as follows: a holomorphic section, $`f`$, of the vector bundle $`𝒱_\lambda `$ is the same thing as a holomorphic map, $`F`$, from $`SL(N)`$ to the representation space $`V`$ which is equivariant with respect to the subgroup $`P`$: $$F(pu)=R(p)F(u)$$ (22) where $`R`$ denotes the representation in question (and is only non-trivial for the Levi subgroup $`L`$), and where $`pP,uSL(N)`$. If the representation in question has Dynkin labels $`(a_1,\mathrm{},a_{N1})`$, it can be represented by a Young tableau with $`m_l`$ boxes in the $`l`$th row, $`l=1,\mathrm{},N1`$ with $$m_l=\underset{k=l}{\overset{N1}{}}a_k$$ (23) Denote an element of $`SL(N)`$ by $`u_I^i`$. The index $`I`$ splits under the Levi sub-algebra as $`I=(r_1,r_2,\mathrm{}r_{\mathrm{}+1})`$, where $`r_i`$ runs from $`k_{i1}`$ to $`k_i`$, with $`k_0=1`$ and $`k_{\mathrm{}+1}=N`$. The desired holomorphic section corresponding to the representation with Dynkin labels is schematically of the form $$F(u)(u_1)^{m_1+m_2+\mathrm{}m_{k_1}}(u_2)^{m_{k_1+1}+\mathrm{}m_{k_2}}\mathrm{}(u_{\mathrm{}+1})^{m_k_{\mathrm{}}\mathrm{}m_{N1}}\widehat{F}$$ (24) where $`u_1`$ denotes the $`k_1\times N`$ matrix $`u_{r_1}^i`$, etc. and where $`\widehat{F}`$ on the right-hand side denotes the SL(N) tensor with a total of $`m=_{k=1}^{N1}ka_k`$ (subscript) indices. ### 2.5 Flag supermanifolds A super-flag in a supervector space $`\text{}^{m\|n}`$ is a nested sequence of sub-supervector spaces $`V_1V_2\mathrm{}V_{\mathrm{}}\text{}^{m\|n}`$ where the super-dimension of $`V_i`$ is $`K_i=(k_i\|\kappa _i)`$, $`k_i(\kappa _i)`$ being the dimensions of the even (odd) subspaces of $`V_i`$ respectively, with $`K_i<K_j`$ for $`i<j`$. This notation means that $`k_ik_j,\kappa _i\kappa _j`$, but with equality for both even and odd indices prohibited. The space of such super-flags will be denoted $`\text{𝔽}_𝐊(m\|n)`$, where $`𝐊=(K_1,\mathrm{}K_{\mathrm{}})`$ . In the present paper we are interested in super-flags in $`\text{}^{4\|N}`$, and they are determined by various parabolic subgroups of $`SL(4\|N)`$. In general, there are three main classes of such spaces according to whether the body is spacetime times an internal flag manifold, a twistor space (in the usual sense), or a combination of both, i.e. a twistor space times an internal flag manifold. We shall be interested in the first of these types of space since they correspond to (complexified) harmonic superspaces. Moreover, we shall be interested in superspaces which correspond to G-analyticity of type $`(p,q)`$. In the complex setting this means that we want to restrict our attention to spaces which have $`(Np)`$ two-component undotted spinor coordinates and $`(Nq)`$ two-component dotted spinor coordinates. The relevant flag supermanifolds are thus given by $`𝐊=(2\|p,2\|p+k_1,\mathrm{}2\|p+k_{\mathrm{}},2\|(Nq))`$. The body of such a supermanifold is locally the product of spacetime with the internal flag manifold $`\text{𝔽}_{p,p+k_1,\mathrm{},p+k_{\mathrm{}},Nq}(N)`$. Such supermanifolds are generalisations of the $`(N,p,q)`$ analytic superspaces introduced in . The latter are the flag supermanifolds $`\text{𝔽}_{2\|p,2\|Nq}(4\|N)`$; they are the spaces with the smallest internal manifolds for a given choice of $`(p,q)`$. The $`(N,p,q)`$ harmonic superspaces have internal flag manifolds $`\text{𝔽}_{p,Nq}(N)`$ with associated Levi sub-algebras $`\text{s}\text{l}(p)\text{s}\text{l}(q)\text{s}\text{l}(r)\text{}^2`$, where $`r=N(p+q)`$. The generalised $`(N,p,q)`$ spaces differ in that a non-trivial internal space is allowed corresponding to the central $`r\times r`$ sector of the internal symmetry group. An important consequence of these definitions is that not all possible internal flag manifolds are compatible with a given $`(p,q)`$ G-analyticity and superconformal symmetry. For example, if either $`p`$ or $`q`$, or both, is bigger than 1, then the internal parabolic cannot be taken to be the Borel subalgebra. <sup>1</sup><sup>1</sup>1This would seeem to contradict refs . However, the harmonic superfields for the underlying massless supermultiplets used by these authors satisfy further harmonic anti-analytic constraints which means that they actually live on $`(N,p,q)`$ spaces; in fact, they are the same fields as those of . This may be clearly seen from the following pictorial representation of the parabolic subalgebra of $`\text{s}\text{l}(4\|N)`$ for $`(N,p,q)`$ analytic superspace: (25) Multiplication of odd elements of the algebra (i.e. the bottom left times the top right) clearly generates the $`p\times p`$ and $`q\times q`$ parts of the internal algebra which must therefore be filled with bullets. The remaining central part of the algebra can be completed in many different ways so that an internal parabolic subalgebra is obtained; we have depicted the simplest example here corresponding to the flag space $`\text{𝔽}_{p,Nq}(N)`$. The group $`SL(4\|N)`$ acts to the left on $`\text{}^{4\|N}`$ which may be considered as $`N`$-extended non-projective super twistor space . To exhibit the flag nature of the above superspaces more clearly it is convenient to make a change of basis of $`\text{}^{4\|N}`$. If a supertwistor $`Z`$ is written as $$Z_A=\left(\begin{array}{c}z^\alpha \\ z_{\dot{\alpha }}\\ \zeta _i\end{array}\right)$$ (26) then the required change of basis swaps $`z_{\dot{\alpha }}`$ and $`\zeta _i`$. $$Z_A\left(\begin{array}{c}z^\alpha \\ \zeta _i\\ z_{\dot{\alpha }}\end{array}\right)$$ (27) In this basis the above parabolic takes the following form: (28) This brings it to block lower-triangular form as in the bosonic case. <sup>2</sup><sup>2</sup>2In this basis the Borel subalgebra consists of the lower triangular matrices; note that it is possible to have inequivalent Borel subalgebras in the super case. The Levi subalgebra for this parabolic is $$\text{l}=\text{s}(\text{g}\text{l}(2\|p)\text{g}\text{l}(2\|q)\text{g}\text{l}(r))$$ (29) provided that $`N4`$. For non-exceptional cases, i.e. when neither $`p`$ nor $`q`$ equals $`2`$ and when $`r0`$, this algebra is isomorphic to $`\text{s}\text{l}(2\|p)\text{s}\text{l}(2\|q)\text{s}\text{l}(r)\text{}^2`$. We shall discuss the exceptional cases in the next section in more detail. In the case of generalised $`(N,p,q)`$ spaces, the central $`\text{s}\text{l}(r)`$ summand is replaced by an appropriate parabolic subalgebra. ## 3 Short representations We can apply the method of parabolic induction straightforwardly in the supersymmetric case to obtain representations of $`SL(4\|N)`$ as holomorphic sections of homogeneous vector bundles over flag supermanifolds defined by parabolic subgroups of the type we have discussed above. The short representations are characterised by being short multiplets and thus having lower component spins than unconstrained superfields on Minkowski superspace. Such representations act naturally on superfields defined on analytic superspaces since they have fewer odd coordinates than Minkowski superspace. The superfields carrying the representations should transform under irreducible representations of Levi subalgebras of the form $`\text{l}=\text{s}(\text{g}\text{l}(2\|p)\text{g}\text{l}(2\|q)\text{g}\text{l}(r))`$; however, in order to ensure that the representations are indeed short these superfields must not carry any spacetime indices. They must therefore transform trivially under any supergroup factors of the Levi subgroup. In the generic case this means that they transform only under $`\text{s}\text{l}(r)\text{}^2`$, where we recall that $`r=N(p+q)`$. Another point to be borne in mind is that the same representation can be realised on different superspaces corresponding to different parabolic subgroups; in particular, for a given representation there will be a class of ”maximally efficient” realisations, by which we mean that the shortening is the largest possible, i.e. $`p`$ and $`q`$ take their largest possible values for the given representation. There may be several representations of this type due to the fact that one is free to choose the central factor of the internal parabolic in different ways. In order to keep matters as simple as possible, we shall concentrate for the time being on $`(N,p,q)`$ superspaces which have internal flag manifolds $`\text{𝔽}_{p,Nq}(N)`$. The representations to be studied can then be represented by modified Dynkin diagrams of the following type: (30) For the reasons discussed above the first $`(p1)`$ and the last $`(q1)`$ Dynkin labels must vanish, leaving $`(r1)`$ labels to specify the representation of the central $`\text{s}\text{l}(r)`$ and two further labels which specify the charges. The Young tableau for this representation is (31) From a purely algebraic point of view representations of superconformal groups in four dimensions are determined by the following set of labels: $`(J_1,J_2,L,R;a_1,a_2,\mathrm{}a_{N1})`$, where $`J_1,J_2`$ are spins, $`L,R`$ are charges corresponding to dilations and R-symmetry respectively while the $`a^{}s`$ are $`\text{s}\text{l}(N)`$ Dynkin labels (see, e.g., for reviews). For representations of the above type $`J_1=J_2=0`$, as are the Dynkin labels $`a_1,\mathrm{}a_{p1};a_{Nq+1}\mathrm{}a_{N1}`$, while the $`L`$ and $`R`$ charges are, except when either $`p`$ or $`q=0,N4`$ determined by these labels. The representations that occur belong to series C with $`J_1=J_2=0`$ for which one has $$L=m_1;R=\frac{2m}{N}m_1$$ (32) where $$m:=\underset{k=1}{\overset{N1}{}}ka_k;m_1:=\underset{k=1}{\overset{N1}{}}a_k$$ (33) Note that $`m`$ is the total number of boxes in the $`\text{s}\text{l}(N)`$ Young tableau, while $`m_1`$ is the number of boxes in the first row. When $`q=0,N4`$, one has the weaker condition $$LR=2m_1\frac{2m}{N}$$ (34) while if $`p=0,N4`$, one finds $$L+R=\frac{2m}{N}$$ (35) Both of these conditions are compatible with either series C or series B. In the latter case one has, for zero spins, $`Lm_1+1`$, together with either (34) or (35). In these cases one more label ($`L`$ or $`R`$) has to be specified in order to fix the representation. For $`N=4`$ the same formulae hold, but now $`R=0`$ as well. In particular, this means that the representations are fully determined by the Dynkin labels even if $`p`$ or $`q=0`$. It should be noted that there is a second possibility for $`N=4`$ when the group is $`SL(4\|4)`$ and not $`PSL(4\|4)`$; this case is the same as general $`N`$. To demonstrate the above explicitly, we first note that, in the basis in which we are working, the matrices representing $`L`$ and $`R`$ are $$L=\frac{1}{2}\left(\begin{array}{ccccc}1_2& & & & \\ & & & & \\ & & & & \\ & & & & \\ & & & & \\ & & & & \\ & & & & 1_2\end{array}\right)$$ (36) and $$R=\frac{1}{2}\left(\begin{array}{ccccc}1_2& & & & \\ & & & & \\ & & & & \\ & & \frac{4}{N}1_N& & \\ & & & & \\ & & & & \\ & & & & 1_2\end{array}\right)$$ (37) As we have mentioned, for most of these representations, fixing $`(p,q)`$ and the Dynkin labels determines the values of $`L`$ and $`R`$ as well. The key point is that, as we remarked above, factors such as $`\text{s}\text{l}(2\|p)`$ must be represented trivially in order for the fields we are interested in to have $`J_1=J_2=0`$. To put it another way, the sub-algebra $`\text{g}\text{l}(2\|p)`$ is represented by its supertrace. This means that any element with vanishing supertrace becomes zero when acting on fields of interest. In particular, $$\left(\begin{array}{ccccc}\frac{1}{2}1_2& & & & \\ & & & & \\ & \frac{1}{p}1_p& & & \\ & & & & \\ & & & & \\ & & & & \end{array}\right)0$$ (38) and, similarly $$\left(\begin{array}{ccccc}& & & & \\ & & & & \\ & & & & \\ & & & & \\ & & & \frac{1}{q}1_q& \\ & & & & \\ & & & & \frac{1}{2}1_2\end{array}\right)0$$ (39) The operator $`1_p`$ on the representation in question counts the total number of boxes in the first $`p`$ rows, of the Young tableau (31), while the operator $`1_N`$ counts the total number of boxes in the diagram. If we denote the field by $`A`$, then $`1_pA=pm_1A`$, $`1_NA=mA`$; since the tableau (31) has boxes only in the first $`Nq`$ rows it follows that $`1_qA=0`$. It might be thought that when either $`p`$ or $`q`$ equals $`2`$ there could be an extra charge due to the fact that $`\text{s}\text{l}(n\|n)`$ is not simple since the unit matrix has vanishing supertrace. However, for vector representation spaces which do not carry any indices for this supergroup, one can show that the unit matrix must be represented by zero (because the anticommutator of two odd generators includes it) and this means that the above rules for evaluating $`L`$ and $`R`$ apply in this case, too. In $`N=4`$ the superconformal group can be taken to be either $`SL(4\|4)`$ or $`PSL(4\|4)`$. The latter group is simple and is the one which is relevant for SYM, while the former is needed for certain representations such as the spin $`3/2`$ supermultiplet. In the second case case it is still easier to work with $`SL(4\|4)`$ and impose the constraint $`R=0`$ to get to $`PSL(4\|4)`$. Note that in either case R-symmetry does not act on the coordinates. We shall now go through the various cases that can arise, explicitly evaluating $`L`$ and $`R`$ for the representations in question: 1. The generic case: $`p\mathbf{,}q\mathbf{,}r\mathbf{}\mathrm{𝟎}`$ In this case, the residual symmetry algebra acting on the representation space of the isotropy group $`V`$ is $`\text{s}\text{l}(r)\text{}^2`$. Using the above rules for evaluating the unit matrices $`1_N,1_p,1_q`$, we find, as expected, $$L=m_1;R=\frac{2m}{N}m_1$$ (40) The representation belongs to series C. 2. $`p\mathbf{,}q\mathbf{}\mathrm{𝟎}\mathbf{;}r\mathbf{=}\mathrm{𝟎}`$ In this case the residual symmetry algebra is simply and the Dynkin diagram reduces to (41) so that all the labels are zero except for $`a_p`$. $`L`$ and $`R`$ are given by the same formulae with $`m=pa_p;m_1=a_p`$. The Young tableau is (42) 3. $`p\mathbf{}\mathrm{𝟎}\mathbf{,}q\mathbf{=}\mathrm{𝟎}`$ The Dynkin diagram now has the form: (43) In this case the Levi subalgebra is $`s(\text{g}\text{l}(2\|p)\text{g}\text{l}(Np)\text{g}\text{l}(2))`$, where the last term arises from the fact that $`q=0`$. In this case the constraint (39) is lost, and one therefore only finds that $$LR=2m_1\frac{2m}{N}$$ (44) This is compatible with both series B and C representations. 4. $`p\mathbf{=}\mathrm{𝟎}\mathbf{,}q\mathbf{}\mathrm{𝟎}`$ The Dynkin diagram now has the form: (45) In this case the Levi subalgebra is $`s(\text{g}\text{l}(2\|q)\text{g}\text{l}(Nq)\text{g}\text{l}(2))`$, where the last term now arises because $`p=0`$. Again the constraint (38) is lost, and one therefore only finds that $$L+R=\frac{2m}{N}$$ (46) This is also compatible with series B and C. For both $`(p,q)=(p,0)`$ and $`(p,q)=(0,q)`$ there is the possibility of having superfields with all Dynkin labels equal to zero. These fields just behave as holomorphic functions with regard to the internal space and therefore do not depend on the coordinates of the internal space by compactness. G-analyticity then implies that they are chiral or anti-chiral scalar fields (for $`(0,q)`$ or $`(p,0)`$ respectively). For $`L=1`$ such superfields are not irreducible and satisfy an additional constraint, as we shall discuss in the next section, but for other (higher) values of $`L`$ they are irreducible superconformal fields with $`L=\pm R`$. The remaining examples all refer to $`N=4`$ with the group $`PSL(4\|4)`$. 5. $`\mathbf{(}p\mathbf{,}q\mathbf{)}\mathbf{=}\mathbf{(}\mathrm{𝟏}\mathbf{,}\mathrm{𝟑}\mathbf{)}`$ This is a particular example of $`(N,p,Np)`$ analytic superspace. In principle, one would expect to have Dynkin diagrams of the form (47) but the imposition of $`R=0`$ implies that $`a_1=0`$, so that there are no short representations on this space. 6. $`\mathbf{(}p\mathbf{,}q\mathbf{)}\mathbf{=}\mathbf{(}\mathrm{𝟐}\mathbf{,}\mathrm{𝟐}\mathbf{)}`$ The Dynkin diagram is: (48) and this is compatible with $`R=0`$. 7. $`\mathbf{(}p\mathbf{,}q\mathbf{)}\mathbf{=}\mathbf{(}\mathrm{𝟏}\mathbf{,}\mathrm{𝟐}\mathbf{)}`$ We expect to have (49) but $`R=0`$ implies that $`a_1=0`$; hence there are no new $`(1,2)`$ representations, since the ones that are allowed can be formulated on $`(2,2)`$ superspace. 8. $`\mathbf{(}p\mathbf{,}q\mathbf{)}\mathbf{=}\mathbf{(}\mathrm{𝟏}\mathbf{,}\mathrm{𝟏}\mathbf{)}`$ In this case we could expect to have arbitrary representations of $`\text{s}\text{l}(4)`$, (50) However, $`R=0`$ implies that $`a_1=a_3`$. 9. $`\mathbf{(}p\mathbf{,}q\mathbf{)}\mathbf{=}\mathbf{(}p\mathbf{,}\mathrm{𝟎}\mathbf{)}`$ The possible diagrams are (51) for $`p=1,2,3`$ respectively. In all cases we find that $$L=2m_1\frac{m}{2}$$ (52) because $`R=0`$ and $`q=0`$. This can be rewritten as $$L=m_1+\frac{1}{2}(a_1a_3)$$ (53) Since the superfields under consideration are bosonic it follows that $`L`$ must be integral, and hence $`a_1a_3=2b`$, so that $$L=m_1+bm_1+1\mathrm{for}b1$$ (54) There are therefore no $`p=3`$ representations, while for $`p=2`$ only $`a_3=0`$ is allowed. Finally, for $`p=1`$ we have $`a_1a_3=2b`$ so that, for $`b0`$ we have series B with $`L=m_1+b`$, while if $`b=0`$ we have series C with $`a_1=a_3`$. 10. $`\mathbf{(}p\mathbf{,}q\mathbf{)}\mathbf{=}\mathbf{(}\mathrm{𝟎}\mathbf{,}q\mathbf{)}`$ The diagrams are the same as in the previous case, but now they correspond to $`q=3,2,1`$ respectively. In all cases we have $$L=\frac{m}{2}=m_1+\frac{1}{2}(a_3a_1)$$ (55) There are no representations for $`q=3`$, for $`q=2`$ we can only have $`a_20`$ and for $`q=1`$ we can have series B with $`L=m_1+b`$ and $`a_3a_1=2b2`$, or series C with $`L=m_1`$, $`a_1=a_3`$. We conclude this survey of the various short representations with a brief discussion of the general case. For the generic case, we can have representations of the following type: (56) This diagram is read as follows: the crosses determine the internal parabolic subalgebra of $`\text{s}\text{l}(N)`$, the corresponding internal flag manifold being $`\text{𝔽}_{p,p+k_1,\mathrm{}p+k_{\mathrm{}},Nq}(N)`$; the two outer crosses, at the $`p`$th and $`(Nq)`$th nodes, determine the explicitly realised G-analyticity to be of type $`(p,q)`$, and the non-zero Dynkin labels are $`a_s\mathrm{}a_t`$ where $`sp`$ and $`tNq`$. If $`s=p`$, $`t=Nq`$ the representation is maximally efficient, but otherwise this is not the case. This holds if neither $`p`$ nor $`q`$ is zero. The representation of the superconformal group is fixed by the diagram since $`L`$ and $`R`$ can be determined and the spins are zero. If either of $`p`$ or $`q`$, say $`q`$, equals zero, then the only restriction is that the left-most cross must be placed at node $`p`$, with $`sp`$. As we have seen, in this case the representation is only fully determined when an extra label is specified unless $`N=4`$ and the group is $`PSL(4\|4)`$ where the fact that $`R=0`$ implies that no extra label is required. Finally, we note that any short representation can be realised on $`(1,0)`$ superspace, while any representation from series $`C`$ can be realised on $`(1,1)`$ superspace, although these realisations will not be maximally efficient in general. ## 4 Massless supermultiplets On-shell massless supermultiplets with maximal helicity $`s`$, where $`[\frac{N}{2}]2s<N`$, ($`[\frac{N}{2}]`$ denotes the nearest integer greater than or equal to $`\frac{N}{2}`$), are described in (real) super Minkowski space, $`M`$, by superfields $`W`$ which have $`p=2s`$ totally antisymmetric internal indices and which satisfy $`\overline{D}_{\dot{\alpha }}^iW_{j_1\mathrm{}j_p}`$ $`=`$ $`{\displaystyle \frac{p(1)^{p1}}{Np+1}}\delta _{[j_1}^i\overline{D}_{\dot{\alpha }}^kW_{j_2\mathrm{}j_p]k}`$ $`D_{\alpha i}W_{j_1\mathrm{}j_p}`$ $`=`$ $`D_{\alpha [i}W_{j_1\mathrm{}j_p]}`$ (57) For each such superfield there is a conjugate superfield $`\stackrel{~}{W}_{i_1\mathrm{}i_{Np}}`$ defined by $$\stackrel{~}{W}_{i_1\mathrm{}i_{Np}}=\frac{1}{p!}\epsilon _{i_1\mathrm{}i_{Np}j_{Np+1}\mathrm{}j_N}\overline{W}^{j_{Np+1}\mathrm{}j_N}$$ (58) The conjugate superfield obeys similar constraints: $`\overline{D}_{\dot{\alpha }}^i\stackrel{~}{W}_{j_1\mathrm{}j_{Np}}`$ $`=`$ $`{\displaystyle \frac{(Np)(1)^{Np1}}{p+1}}\delta _{[j_1}^i\overline{D}_{\dot{\alpha }}^k\stackrel{~}{W}_{j_2\mathrm{}j_{Np}]k}`$ $`D_{\alpha i}\stackrel{~}{W}_{j_1\mathrm{}j_{Np}}`$ $`=`$ $`D_{\alpha [i}\stackrel{~}{W}_{j_1\mathrm{}j_{Np}]}`$ (59) When $`s=\frac{1}{4}N`$, the multiplet is self-conjugate: $$\overline{W}^{i_1\mathrm{}i_p}=\frac{1}{p!}\epsilon ^{i_1\mathrm{}i_pj_1\mathrm{}j_p}W_{j_1\mathrm{}j_p}.$$ (60) For each $`N`$ there is therefore a range of scalar superfields $`W_{i_1\mathrm{}i_p}`$ antisymmetric on $`p`$ internal indices with $`p`$ ranging from 1 to $`N1`$ all of which obey the constraints (57). We can extend this to $`N`$, since for such a superfield the constraints (57) imply that it is anti-chiral. Its conjugate has no indices and is chiral. Such a chiral field describes an on-shell massless supermultiplet (with maximum spin $`N/2`$) if it satisfies the additional constraint $$D_{\alpha i}D_j^\alpha W=0$$ (61) Note that this equation is only superconformal if the dilation weight $`L=1`$; for other values of $`L`$ a chiral field is superconformally irreducible and (61) should not be imposed. The superfields $`W_{i_i\mathrm{}i_p}`$ can be naturally described as CR-analytic fields on $`(N,p,Np)`$ harmonic superspaces . On $`M\times SU(N)`$ define the superfield $`W^{(p)}`$ by $$W^{(p)}=\frac{1}{p!}\epsilon ^{r_1\mathrm{}r_p}u_{r_1}{}_{}{}^{i_1}\mathrm{}u_{r_p}{}_{}{}^{i_p}W_{i_1\mathrm{}i_p}^{},$$ (62) then equation (57) implies that $`W^{(p)}`$ satisfies the G-analyticity constraints $`D_{\alpha r}W^{(p)}`$ $`=`$ $`0`$ (63) $`\overline{D}_{\dot{\alpha }}^r^{}W^{(p)}`$ $`=`$ $`0,`$ (64) and the H-analyticity constraint, $$D_r{}_{}{}^{s^{}}W_{}^{(p)}=0$$ (65) In addition $`W^{(p)}`$ is clearly equivariant with respect to the isotropy group $`S(U(p)\times U(Np))`$ of the internal flag manifold $`\text{𝔽}_{p,Np}(N)`$, which in this case is the Grassmannian of $`p`$-planes in $`\text{}^N`$, $`\text{𝔾}_p(N)`$. Furthermore, an equivariant field with the same charge as $`W^{(p)}`$ satisfying the above analyticity constraints is equivalent to a superfield on super Minkowski superspace satisfying (57). We can apply this construction to the entire family of scalar superfields with $`p`$ varying from $`0`$ to $`N`$, bearing in mind that the chiral fields satisfy the additional constraint (61). Under harmonic conjugation the superfield $`W^{(p)}\stackrel{~}{(W^{(p)})}=W^{(Np)}`$. In complexified superspace, there are corresponding fields $`W^{(p)}`$ and $`\stackrel{~}{W}^{(p)}`$ defined on $`(N,p,Np)`$ and $`(N,Np,p)`$ analytic superspace respectively. Furthermore, it is possible to extend harmonic conjugation to this case straightforwardly. The Dynkin diagram for $`W^{(p)}`$ is (66) where the single non-zero label is above the $`p`$th node, and thus corresponds to the representation of $`\text{s}\text{l}(N)`$ on $`^P(\text{}^N)`$. The diagram for $`\stackrel{~}{W}^{(p)}`$ is similar, but with $`p`$ replaced by $`Np`$. Although $`W^{(p)}`$ provides the most efficient realisation of this representation, the maximum spin for $`W^{(p)}`$, when $`pN/2`$, is in fact only $`(0,\frac{p}{2})`$, and not $`(\frac{Np}{2},\frac{p}{2})`$; these representations are therefore ultra-short. For $`\stackrel{~}{W}^{(p)}`$ the maximum spin is $`(\frac{p}{2},0)`$. In the self-conjugate case, $`W^{(p)}`$ contains both $`(0,\frac{p}{2})`$ and $`(\frac{p}{2},0)`$. There are many other ways of representing such multiplets which are “less efficient” in that the superspaces have more odd coordinates. Following the discussion at the end of the previous section we can simply place crosses where we like on the above Dynkin diagram with the restrictions that the cross furthest to the left must have node number $`k`$ less or equal to $`p`$ and the cross furthest to the right must have node number $`l`$ at least equal to $`p`$. The corresponding superspace will then have G-analyticity of type $`(k,l)`$. For example, any such field (excluding $`p=0,N`$) can be realised with $`(1,1)`$ G-analyticity on an internal flag defined by the Borel subgroup (i.e. crosses at all nodes). To illustrate this procedure we consider the $`N=4`$ Maxwell supermultiplet which is represented on $`N=4`$ super Minkowski space by the familiar self-conjugate Sohnius superfield $`W_{ij}`$ . The possible realisations of this superfield have G-analyticities $`(2,2),(1,1)`$ and $`(1,0)`$. The $`(2,2)`$ superfield is the one introduced in . The (1,1) superfield can be written $`W_{1r},r\{2,3\}`$; in $`(4,1,1)`$ harmonic superspace it obeys the constraints $$D_{\alpha 1}W_{1r}=\overline{D}_{\dot{\alpha }}^4W_{1r}=0$$ (67) Under harmonic conjugation we have $$W_{1r}\stackrel{~}{W}^{4r}=ϵ^{rs}W_{1s}$$ (68) The $`(1,0)`$ version is a superfield $`W_{1r}`$, $`r\{2,3,4\}`$ satisfying $$D_{\alpha 1}W_{1r}=0$$ (69) while the $`(0,1)`$ version is a superfield $`W_{RS},R1,2,3`$ obeying $$\overline{D}_{\dot{\alpha }}^4W_{RS}=0$$ (70) Under harmonic conjugation we find $$W_{1r}\stackrel{~}{W}^{4R}=\frac{1}{2}ϵ^{RST}W_{ST}$$ (71) Due to self-conjugacy, the constraint (69) together with harmonic analyticity is sufficient to reproduce the standard constraints obeyed by the Sohnius superfield. The above superfields have been defined on $`(4,1,1)`$ and $`(4,1,0)`$ or $`(4,0,1)`$ superspaces which have the smallest possible internal flags. It is possible to relax this; for example we could use the maximal flag space determined by the Borel subalgebra. In the $`(1,1)`$ case we would then split the indices as $`I=1,2,3,4`$; this enables us to define the field $`W_{12}`$ which is $`(1,1)`$ G-analytic and H-analytic on $`\text{𝔽}_B`$. Note that this field satisfies further first-order differential constraints; it is actually $`(2,2)`$ G-analytic and is also annihilated by $`D_2^1`$ and $`D_4^3`$ ## 5 Tensor products General short multiplets can be obtained as tensor products of the massless supermultiplets in a straightforward fashion. If one wishes to find such a short multiplet with $`(p,q)`$ analyticity the optimal way to construct it from massless multiplets is to use realisations of the latter with the same analyticity. One then multiplies an appropriate number of massless multiplets together to get the desired representation. Because the massless multiplets will in general transform under non-trivial representations of the internal symmetry algebra ($`\text{s}\text{l}(r)`$ in the case of $`(N,p,q)`$ superspace), it will be necessary to project out an irreducible representation of this algebra in order to obtain the desired irreducible short representation. ### 5.1 Generating all short representations There is a very simple way of generating all possible short representations from the set of available massless supermultiplets for a given $`N`$. Consider first those belonging to series C. We realise each of the superfields $`W^{(p)}`$ on $`(N,1,1)`$ superspace with $`\text{𝔽}_B`$ as the internal flag space. Thus we have superfields $`W_1,W_{12},W_{123},\mathrm{}W_{123\mathrm{}N1}`$ all of which are G-analytic with respect to $`D_{\alpha 1}`$ and $`\overline{D}_{\dot{\alpha }}^N`$. Let $`(a_1,\mathrm{}a_{N1})`$ be the Dynkin labels of the representation we are interested in and simply form the product $$A=\underset{k=1}{\overset{N1}{}}\left(W_{12\mathrm{}k}\right)^{a_k}$$ (72) The leading component of this superfield transforms according to the irreducible representation of $`SL(N)`$ specified by the Dynkin labels. The dilation weight is just equal to the sum of the Dynkin labels, that is $`m_1`$, since each underlying field has weight 1. This construction of a general representation also illustrates what happens when the Dynkin labels are restricted. For example, suppose that the first $`p1`$ and the last $`q1`$ labels are zero, then the only superfields present in the above product will be $`W_{12\mathrm{}p},\mathrm{}W_{12\mathrm{}Nq}`$. Each of these superfields is G-analytic with respect to $`D_{\alpha 1}\mathrm{}D_{\alpha p}`$ and $`\overline{D}_{\dot{\alpha }}^{Nq+1}\mathrm{}\overline{D}_{\dot{\alpha }}^N`$; furthermore, as well as being H-analytic on $`\text{𝔽}_B`$ it satisfies further harmonic anti-analyticity conditions and can therefore be thought of as a superfield on $`\text{𝔽}_{p,p+1,p+2,\mathrm{}Nq1,Nq}(N)`$ superspace. For series B representations we can work on $`(N,0,1)`$ superspace again with $`\text{𝔽}_B`$ as the internal flag manifold. We use the same family of superfields as before augmented with $`W_o`$ which satisfies $`\overline{D}_{\dot{\alpha }}^NW_o=0`$ together with the additional condition (61). We can then form the product superfield $$A^{}=(W_o)^b\underset{k=1}{\overset{N1}{}}\left(W_{12\mathrm{}k}\right)^{a_k}$$ (73) where $`b`$ is a positive integer. The $`L`$ weight of this field is $`b`$ plus the dilation weight of the previous example, so $`L=m_1+b`$. Likewise the R-symmetry weight differs from the previous example by $`b`$, so $`R=2m/N(m_1+b)`$ and therefore $`L+R=2m/N`$ as required. In general this superfield will only be annihilated by $`\overline{D}_{\dot{\alpha }}^N`$, but if the last $`q1`$ Dynkin labels vanish, the superfield will be defined on $`(N,0,q)`$ superspace. If the field $`A`$ of (72) is multiplied by products of the anti-chiral field with $`L=1,R=1`$ we get the B series of representations on $`(N,1,0)`$ superspace with $`LR=2m_12m/N`$. ### 5.2 Products of a given massless multiplet In order to find the short representation composites that can be constructed using a given underlying massless supermultiplet one first of all computes the different Young tableau that can arise by multiplying together the tableaux for $`k`$ factors of $`W`$ and $`l`$ factors of $`\stackrel{~}{W}`$. For series C, these tableaux need to be subjected to the following restrictions: firstly, there should be $`k+l`$ boxes in the first row, secondly, there should be a total of $`pk+(Np)l`$ boxes with at most $`N1`$ in the first column and thirdly, a number of representaions will be absent due to symmetrisation. In the case of series B, one can either have columns with $`N`$ boxes, in which case after these have been discarded the remaining tableaux will clearly have fewer than $`pk+(Np)l`$ boxes, or one is still restricted to at most $`N1`$ boxes in a column, but with the constraint that there should be $`k+l`$ boxes in the first row dropped. The problem of restrictions due to symmetrisation can be overcome if one is allowed to use different copies of the same underlying multiplet. As a simple example consider the $`N=3`$ Maxwell multiplet which is described by an on-shell superfield $`W_{ij}`$ of the above type. The corresponding $`W`$ is CR-analytic on $`(3,2,1)`$ harmonic superspace, and its conjugate $`\stackrel{~}{W}`$ is CR-analytic on $`(3,1,2)`$ harmonic superspace. Both of them can be realised as CR-analytic superfields on $`(3,1,1)`$ superspace, with $$W_{12}=u_1{}_{}{}^{i}u_{2}^{}{}_{}{}^{j}W_{ij}^{};\stackrel{~}{W}_1=u_1{}_{}{}^{i}\stackrel{~}{W}_{i}^{}$$ (74) On $`(3,2,1)`$ superspace, the products of $`W`$ with itself have maximum spins $`\frac{1}{2},1`$, and corresponding tableaux with 2 rows of the same length. The simplest of these supermultiplets, $`S`$, say, has its leading component in the 6-dimensional representation of $`SU(3)`$. In order to construct the energy-momentum tensor multiplet one has to take the product of $`W`$ and $`\stackrel{~}{W}`$ on $`(3,1,1)`$ superspace. This multiplet, $`T`$, say, has maximum spin $`(1,1)`$ and its leading component is in the $`8`$-dimensional representation of $`SU(3)`$. The multiplets $`T`$ and $`S`$ are the $`N=3`$ “components” of the $`N=4`$ supercurrent. (See for more details of $`N=3`$ mulitplets.) An interesting set of multiplets can be found for supergravity theories with $`N5`$. The basic scalar superfield with maximum spin 2 has four indices and is self-conjugate in $`N=8`$. The analogue of the energy-momentum tensor in some senses for this multiplet is the Bel-Robinson tensor ($`C_{\alpha \beta \gamma \delta }\overline{C}_{\dot{\alpha }\dot{\beta }\dot{\gamma }\dot{\delta }}+\mathrm{}`$ in two-component spinor notation, where $`C`$ is the Weyl spinor). This spin 4 field is the highest component of the composite multiplet obtained by multiplying $`W`$ and its conjugate. This superfield lives naturally on $`(N,N4,N4)`$ superspace. If we square it, we obtain the three-loop supergravity counterterm supermultiplet . In $`N=8`$ this is simply $`W^4`$ since $`W`$ is self-conjugate . ### 5.3 N=4 composites We conclude this discussion of composite operators with some comments on $`N=4`$ SYM since this is of most interest in the AdS/CFT context. The superconformal group in this case is $`PSU(2,2\|4)`$ and these fields were discussed from a different perspective in . We begin with Maxwell multiplets which we denote by $`A,B,C`$, etc. If we regard these as superfields on $`(4,2,2)`$ superspace, the only composites we can build are of the form $`A^k`$ with Dynkin labels $`(0,k,0)`$ and highest spin $`(1,1)`$. On (1,1) superspace we can generate all the series C representations using two fields $`A_{1r},B_{1r}`$, where $`r\{2,3\}`$. The product $`A^kB^l`$, where we may assume $`kl`$, gives rise to a series of representations with Dynkin labels $`(n,k+l2n,n)`$ where $`n=0,1,\mathrm{}l`$. The point here is that products of $`A`$ with itself must be symmetric and can therefore have tableaux with only two rows. To get three rows it is necessary to introduce a second field $`B`$. Finally, to get representations from series B one has to use a third field $`C`$ and work on $`(1,0)`$ superspace. In this manner one can generate tableaux with four rows, although one then discards the columns which do have four rows. In the non-Abelian theory the field strength superfield $`W`$ takes its values in the Lie algebra of the gauge group which we shall take to be $`SU(N_c)`$; it is not itself a CR-analytic superfield on $`(4,2,2)`$ harmonic superspace, rather it is H-analytic and covariantly G-analytic . However, gauge-invariant products of powers of $`W`$ are CR-analytic and hence can be used to construct short representations of the $`N=4`$ superconformal group. The simplest series is the KK series $`A_k:=\mathrm{Tr}(W^k)`$; this series of operators on $`(4,2,2)`$ space has maximum spin $`(1,1)`$ and includes the supercurrent for $`k=2`$ . To obtain other short representations from series C we consider the field strength tensor as a $`(4,1,1)`$ superfield $`W_{1r},r\{2,3\}`$. We can then form the double trace operators $`\mathrm{Tr}(W_{1r}^k)\mathrm{Tr}(W_{1s}^l)`$, which of course decompose into irreducible representations under $`\text{s}\text{l}(2)`$. Finally, we can obtain series B representations by taking triple traces of a similar type using the basic superfield $`W_{1r},r\{2,3,4\}`$ which is covariantly G-analytic on $`(4,1,0)`$ superspace. Each factor in these multi-trace operators, if we assume the gauge-group trace to be symmetrised, is one of the operators in the $`A_k`$ series, so that another way of obtaining these operators is as products of these. For series C has to write these operators in $`(4,1,1)`$ superspace, thus $`A_k=\mathrm{Tr}(W^k)`$ is replaced by $`A_{r_1\mathrm{}r_k}:=\mathrm{Tr}(W_{1r_1}\mathrm{}W_{1r_k})`$. One can then form the products $`A_{r_1\mathrm{}r_k}A_{s_1\mathrm{}s_l}`$ which can subsequently be decomposed into irreducible representations under $`SU(2)`$. Similarly the triple trace operators can be viewed as products of sets of three $`A_k`$ operators written as superfields on $`(4,1,0)`$ superspace. The multi-trace operators are also discussed in . It would seem that further gauge-invariant operators could be constructed using the structure constants, ; for example, on $`(1,0)`$ superspace, there is an operator $`\epsilon ^{rst}\mathrm{tr}(W_{1r}W_{1s}W_{1t})`$ where the trace is no longer symmetrised. ## 6 Superconformal transformations on harmonic superspace The short representations obtained by parabolic induction in complex superspaces transform under the superconformal group in the standard way. These transformations can then be adapted to real spacetime. However, if one is interested in real harmonic superspaces in the GIKOS formalism it is perhaps just as easy to derive the explicit transformation rules by other means. This is what we shall do in this section. Most of the results given here were obtained previously in , but the derivation is perhaps slightly simpler. We start from superconformal transformations on super Minkowski space, extend these in a simple way to $`M\times SU(N)`$ and then make a final modification in order to preserve analyticity. We recall that a superconformal transformation on super Minkowski space is a diffeomorphism which preserves the odd tangent space (spanned by the odd covariant derivatives). Dually, such a transformation can be thought of as one which preserves the even cotangent bundle spanned by the differential forms $`E^{\alpha \dot{\alpha }}:=dx^{\alpha \dot{\alpha }}+\frac{i}{2}d\theta ^{\alpha i}\overline{\theta }_i^{\dot{\alpha }}+\frac{i}{2}d\overline{\theta }_i^{\dot{\alpha }}\theta ^{\alpha i}`$. An infinitesimal such transformation is generated by a superconformal Killing vector $`V`$ where $$V=F^{\alpha \dot{\alpha }}_{\alpha \dot{\alpha }}+f^{\alpha i}D_{\alpha i}\overline{f}_i^{\dot{\alpha }}\overline{D}_{\dot{\alpha }}^i$$ (75) with $$D_{\alpha i}F_{\beta \dot{\beta }}=i\epsilon _{\alpha \beta }\overline{f}_{\dot{\beta }i}$$ (76) so that $$f^{\alpha i}=\frac{i}{2}\overline{D}_{\dot{\alpha }}^iF^{\alpha \dot{\alpha }}$$ (77) Thus such a vector field is determined by the function $`F^{\alpha \dot{\alpha }}`$ obeying the above constraint. The Lie algebra can therefore be presented in terms of $`F`$, $$[V_F,V_F^{}]=V_{[F,F^{}]}$$ (78) where $$[F,F^{}]^{\alpha \dot{\alpha }}:=F^{{}_{}{}^{\prime \prime }\alpha \dot{\alpha }}=(FF^{\alpha \dot{\alpha }}if^{\alpha i}\overline{f}_i^{{}_{}{}^{}\dot{\alpha }})(FF^{})$$ (79) It is not difficult to verify that the superfield $`F`$ subject to (76) does indeed specify a superconformal transformation. In $`N=4`$ there is an additional constraint, namely $$Df+\overline{D}\overline{f};=D_{\alpha i}f^{\alpha i}+\overline{D}_{\dot{\alpha }}^i\overline{f}_i^{\dot{\alpha }}=0$$ (80) For general $`N`$ one has the identity $$F\frac{2}{N}\left(Df\overline{D}\overline{f}\right)=0$$ (81) and we identify the dilation and R-symmetry parameters $`L`$ and $`R`$ by $`F`$ $`=`$ $`4L+\mathrm{}`$ $`Df+\overline{D}\overline{f}`$ $`=`$ $`2i(N4)R,\mathrm{N}4`$ (82) The $`N=4`$ constraint (80) has the geometrical consequence that the $`N=4`$ chiral measure is superconformally invariant, since it implies that $`FDf=0`$. The above transformation can be extended to a transformation of $`M\times SU(N)`$. The corresponding Killing vector in this case has the form $$𝒱=V_o+V_u$$ (83) where $`V_o`$ is a superconformal Killing vector on super Minkowski space and $$V_u:=f_I{}_{}{}^{J}D_{J}^{}^I$$ (84) with $$f_I{}_{}{}^{J}=\frac{1}{2}(D_{\alpha I}f^{\alpha J}\frac{1}{N}\delta _I{}_{}{}^{J}D_{\alpha K}^{}f^{\alpha K})$$ (85) The parameter $`f_I^J`$ is the parameter superfield $`f_i^j`$ with indices converted to capital ones by means of $`u`$’s. The leading component of $`f_i^j`$ is the $`\text{s}\text{u}(N)`$ parameter. It is straightforward to verify that $`𝒱`$ gives a representation of the superconformal algebra. If $`A`$ is an equivariant superfield on $`M\times SU(N)`$ then it would be natural to look for a transformation rule of the form $`\delta A=𝒱A`$. This preserves H-analyticy, but not preserve G-analyticity. This deficiency can be overcome by including an extra term in the transformation rule. This term is constructed using the chiral function $`K`$ defined by $$K=\frac{1}{4}F+\frac{1}{2(N4)}\left(Df+\overline{D}\overline{f}\right)$$ (86) which is normalised so that its leading component is $`SiR`$. It is also proportional to the transformation factor of the chiral measure. A term of the form $`(aK+b\overline{K})A`$ can be added to the transformation rule for $`A`$ because $$V_FK_F^{}V_F^{}K_F=K_{[F,F^{}]}$$ (87) Moreover, it is the only available function which preserves H-analyticity. We consider a field $`A`$ which has $`(p,q)`$ G-analyticity and which is H-analytic in the sense of $`(N,p,q)`$ harmonic superspace. As before, we set $`I=(r,r^{\prime \prime },r^{})`$. The field $`A`$ satisfies the G-analyticity conditions $$D_{\alpha r}A=\overline{D}_{\dot{\alpha }}^r^{}A=0$$ (88) as well as the H-analyticity conditions $$D_r{}_{}{}^{s^{}}A=D_r{}_{}{}^{s^{\prime \prime }}A=D_{r^{\prime \prime }}{}_{}{}^{s^{}}A=0$$ (89) $`A`$ cannot transform under $`SU(p)`$ (acting on $`r`$) or $`SU(q)`$ (acting on $`r^{}`$) because this is incompatible with G-analyticity. Thus we find $`V_uA`$ $`=`$ $`f_r{}_{}{}^{s^{}}D_{s^{}}^{}{}_{}{}^{r}A+f_r{}_{}{}^{s^{\prime \prime }}D_{s^{\prime \prime }}^{}{}_{}{}^{r}A+f_{r^{\prime \prime }}{}_{}{}^{s^{}}D_{s^{}}^{}{}_{}{}^{r^{\prime \prime }}A+f_{r^{\prime \prime }}{}_{}{}^{s^{\prime \prime }}\widehat{D}_{s^{\prime \prime }}^{}{}_{}{}^{r^{\prime \prime }}A`$ (90) $`+f_oD_oA+f_o^{}D_o^{}A+f_o^{\prime \prime }D_o^{\prime \prime }A`$ The first three terms involve the coset derivatives of the internal flag manifold and the absence of their conjugates reflects the H-analyticity of the field $`A`$. The fourth term takes care of the transformation properties of $`A`$ under $`SU(r)`$ ($`\widehat{D}_{r^{\prime \prime }}{}_{}{}^{r^{\prime \prime }}=0`$). The next two terms are left over from $`f_r{}_{}{}^{s}D_{s}^{}{}_{}{}^{r}A`$ and $`+f_r^{}{}_{}{}^{s^{}}D_{s^{}}^{}{}_{}{}^{r^{}}A`$ which can only be non-trivial for the trace parts, and the final term arises from removing the trace from $`D_{r^{\prime \prime }}^{s^{\prime \prime }}`$. In the above we have used the definition $`D_o={\displaystyle \frac{1}{p}}D_r{}_{}{}^{r};f_o=f_r^r`$ (91) with similar definitions for primed and double-primed indices. The fact that we are working with $`SU(N)`$ then implies $`pD_o+qD_o^{}+rD_o^{\prime \prime }`$ $`=`$ $`0`$ $`f_o+f_o^{}+f_o^{\prime \prime }`$ $`=`$ $`0`$ (92) A short calculation shows that $`D_{\alpha r}(𝒱A)`$ $`=`$ $`{\displaystyle \frac{i}{2}}_{\alpha \dot{\alpha }}\overline{f}_r^{\dot{\alpha }}D_oA`$ $`\overline{D}_{\dot{\alpha }}^r^{}(𝒱A)`$ $`=`$ $`{\displaystyle \frac{i}{2}}_{\alpha \dot{\alpha }}f^{\alpha r^{}}D_o^{}A`$ (93) Since $`D_{\alpha i}K`$ $`=`$ $`{\displaystyle \frac{i}{2}}_{\alpha \dot{\alpha }}\overline{f}_i^{\dot{\alpha }}`$ $`\overline{D}_{\dot{\alpha }}{}_{}{}^{i}\overline{K}`$ $`=`$ $`{\displaystyle \frac{i}{2}}_{\alpha \dot{\alpha }}f^{\alpha i}`$ (94) we see that the desired transformation rule for $`A`$ is $$\delta A=𝒱A+(cKc^{}\overline{K})A$$ (95) where $$D_oA=cA;D_o^{}A=c^{}A$$ (96) This transformation rule can be rewritten in the form $`\delta A`$ $`=`$ $`f_r{}_{}{}^{s^{}}D_{s^{}}^{}{}_{}{}^{r}A+f_r{}_{}{}^{s^{\prime \prime }}D_{s^{\prime \prime }}^{}{}_{}{}^{r}A+f_{r^{\prime \prime }}{}_{}{}^{s^{}}D_{s^{}}^{}{}_{}{}^{r^{\prime \prime }}A+f_{r^{\prime \prime }}{}_{}{}^{s^{\prime \prime }}\widehat{D}_{s^{\prime \prime }}^{}{}_{}{}^{r^{\prime \prime }}A`$ (97) $`+c(K+f_o{\displaystyle \frac{p}{r}}f_o^{\prime \prime })Ac^{}(\overline{K}f_o^{}+{\displaystyle \frac{q}{r}}f_o^{\prime \prime })A`$ The parameters $`f_o\frac{p}{r}f_o^{\prime \prime }`$ and $`f_o^{}\frac{q}{r}f_o^{\prime \prime }`$ are independent of the dilation and R-symmetry parameters, so that the weights for these two transformations are determined by the parameters $`c`$ and $`c^{}`$. Using the fact that $`D_ou_r`$ $`=`$ $`{\displaystyle \frac{Np}{pN}}u_r`$ $`D_ou_r^{}`$ $`=`$ $`{\displaystyle \frac{1}{N}}u_r^{}`$ $`D_ou_{r^{\prime \prime }}`$ $`=`$ $`{\displaystyle \frac{1}{N}}u_{r^{\prime \prime }}`$ (98) together with similar formulae for $`D_o^{}`$ and $`D_o^{\prime \prime }`$, it is straightforward to show that $$c=m_1\frac{m}{N},c^{}=\frac{m}{N}$$ (99) We therefore find $$L=m_1,R=\frac{2m}{N}m_1$$ (100) which is the correct series C result. It is straightforward to adapt the above calculation to the special cases $`p=0`$, $`q=0`$. For $`N=4`$ we have two possibilities, according to whether the superconformal group is chosen to be $`PSU(2,2\|4)`$ or $`SU(2,2\|4)`$. In both cases R-symmetry does not act on the coordinates, and in the former case it is absent. The parameter $`Df+\overline{D}\overline{f}`$ does act on the coordinates and is therefore not the R-parameter; it is the parameter of $`U(1)_Y`$ transformations in the terminology of . Since the corresponding generator does not belong to the superconformal algebra in either case the parameter $`Df+\overline{D}\overline{f}`$ must be set to zero. For $`PSU(2,2\|4)`$ the results derived above can be straightforwardly taken over provided one sets $`R=0`$. In the $`SU(2,2\|4)`$ case, things are a little more subtle. It is possible to infer the existence of a real superfield parameter $`\widehat{R}`$, say, whose leading component is the $`R`$ parameter and which satisfies $$D_{\alpha i}\widehat{R}=\frac{1}{4}_{\alpha \dot{\alpha }}\overline{f}_i^{\dot{\alpha }}$$ (101) One can then define the chiral superfield $$K=\frac{1}{4}Fi\widehat{R}$$ (102) and proceed as before. The only difference is that the R-transformation so defined does not act on the coordinates, although it does act on the superfield itself multiplicatively. Acknowledgement This research was supported in part by PPARC SPG grant 613.
warning/0005/astro-ph0005211.html	ar5iv	text	# ON YOUNG NEUTRON STARS AS PROPELLERS AND ACCRETORS WITH CONVENTIONAL MAGNETIC FIELDS ## 1 Introduction Anomalous x-ray pulsars (AXPs) are characterized by periods in the range 6-12 secs (see Mereghetti 1999 for a review and an extensive list of references therein for properties of specific AXPs). RXJ0720.4-3125, one of the dim nearby ROSAT point sources thought to be thermally emitting neutron stars (Walter, Wolk & Neuhauser 1996, Walter & Matthews 1997) has a measured period (Haberl et al. 1997) of 8.4 s, while a recently discovered member of this class, RXJ0420.0-5022, is reported to have a period of 22.7 s (Haberl, Pietsch & Motch 1999). Among the four confirmed soft gamma ray repeaters (SGRs- Mazets et al. 1979, 1981, Cline et al. 1980, Woods et al. 1999) with quiescent X-ray sources, there are two with measured pulse periods. Both periods are in the AXP period range, with P = 7.47s for SGR 1806-20 (Kouveliotou et al. 1998) and P = 5.16 s for SGR 1900+14 (Hurley et al.1999, Kouveliotou et al. 1999). For the AXPs it was pointed out (van Paradijs, Taam & van der Heuvel 1995, Ghosh, Angelini & White 1995) that these periods would obtain as equilibrium periods for the accretion rates $`\dot{M}10^{15}gm/s`$ corresponding to the observed luminosities and for 10<sup>11</sup>-10<sup>12</sup> G magnetic fields. Values of $`\dot{P}`$ measured from the AXPs and SGRs, together with the long rotation periods in the 5-12 s range have led to the suggestion that these sources are magnetars (Kouveliotou et al. 1998, 1999, Duncan & Thompson 1992, Thompson & Duncan 1993, 1995), neutron stars with very strong magnetic fields B$``$ 10<sup>14</sup>-10<sup>15</sup> G spinning down through magnetic dipole radiation. The short cooling age estimated for RXJ0720.4-3125 has also been suggested (Heyl & Hernquist 1998) as evidence for a magnetar if the source is an isolated pulsar born with a short rotation period. We propose here an alternative explanation, extending the accretion hypothesis to include propeller phases, and linking these classes of neutron stars as stars with conventional $``$ 10<sup>12</sup> G magnetic fields. Mass inflow at rates $`\dot{M}`$ 5 $`\times `$ 10<sup>15</sup> \- 4 $`\times `$ 10<sup>16</sup> gm/s from a disk (van Paradijs, Taam & van den Heuvel 1995, Ghosh, Angelini & White 1997) around the neutron star will produce torques on the neutron star of the order of the spindown torques observed. The AXPs and SGRs are accreting at least part of the incoming mass flow. The similarity of the periods simply reflects similar conditions. All these systems are asymptotically approaching equilibrium periods in a common range, defined by the common ranges of magnetic fields and mass transfer rates. As the approach to rotational equilibrium is asymptotic, $`P/\dot{P}`$ is not the true age of these systems. Sources of different ages and different circumstances with similar conventional magnetic fields have similar periods in this asymptotic regime under a fairly large range of mass inflow rates. We propose that the DTNs, AXPs, SGRs and RQNSs constitute the alternative subclasses of young neutron stars produced by supernovae, complementary to the subclass of neutron stars that show up as radio pulsars (Kaspi 1999). In view of the plausible associations of several AXPs and SGRs with supernova remnants, tight upper limits on the mass of a binary companion (Mereghetti, Israel & Stella 1998, Wilson et al. 1998), and statistics and age considerations, these sources are likely to be young neutron stars under mass inflow from a remnant disk formed by fallback material from the supernova explosion (Chevalier, 1989, Lin, Woosley & Bodenheimer 1991, Mineshige, Nomoto & Shigeyama 1993). Most ideas discussed in this paper and their links to related subjects were first presented in an earlier preprint (Alpar 1999). The DTNs are the propeller counterparts to the accreting AXPs. Their luminosity is due to energy dissipation in the neutron star. They may well account for a large fraction of neutron stars formed in supernovae. For the AXPs the scenario of accretion from a remnant disk, preceded with a propeller phase, has been explored also by Chatterjee, Hernquist & Narayan (1999), who have employed specific thin disk evolution models for the time dependence of the mass inflow from the remnant disk, while Marsden et al. (1999) have evaluated the supernova associations of AXPs and SGRs to argue that what unites these sources is their atypical SNR environments rather than an intrinsic property of the neutron stars like superstrong magnetic fields. Under the working hypothesis of the present work, that all supernovae form neutron stars with conventional $``$ 10<sup>12</sup> G magnetic fields, the differences between the DTNs, AXPs and SGRs are due to their different mass inflow environments and histories. At the highest mass inflow rates we have the RQNSs. These are propellers whose pulse periods are not observable because of accumulated circumstellar material that is optically thick to electron scattering. The subclasses unified in this scenario represent the entire range of mass inflow rates, extending from $`\dot{M}0`$ for the radio pulsars to near Eddington rates inferred for the RQNSs. Table 1 displays observed parameters and model estimates for the AXPs, for the two SGRs with $`\dot{P}`$ measurements, the two DTNs with measured periods, and for a typical RQNS. Section 2 discusses the energy dissipation in a neutron star under an external torque. Section 3 introduces propeller torques and applies these ideas to DTNs. Section 4 presents a simple model for asymptotic spindown. Section 5 relates the different classes of neutron stars, including the RQNSs, as supernova products under mass inflow. Section 6 reviews the expected signature of a remnant disk around the neutron star, and Section 7 presents the conclusions. ## 2 Luminosity of a Non-accreting Neutron Star from Energy Dissipation The DTNs yield fits to blackbody spectra with temperatures of 57 eV for RXJ185635-3754 and 79 eV for RXJ0720.4-3125 and luminosities in the $`L_x10^{3132}`$ erg s<sup>-1</sup> range (Walter, Wolk & Neuhauser 1996, Walter & Matthews 1997, Haberl et al. 1997, Motch & Haberl 1998, Kulkarni & van Kerkwijk 1998). Several further candidates for this class (Stocke et al. 1995, Haberl, Motch & Pietsch 1998, Schwope et al. 1999, Motch et al. 1999, Haberl, Pietsch & Motch 1999) also have similar blackbody temperatures, flux values and limits on the ratio of X-ray flux to optical flux. Accretion from the interstellar medium would require unlikely ambient interstellar medium densities and low velocities (for a discussion of accretion from the ISM for the DTNs see Treves et al. 1999 and references therein). The presence of several dim thermal neutron stars within $``$ 100 pc suggests that there are $``$ 10<sup>4</sup> such sources in the galaxy, requiring ages of $``$ 10<sup>6</sup> years or longer, if the birth rate is 10<sup>-2</sup> yr<sup>-1</sup> or less. The cooling of a young neutron star can typically yield the observed thermal luminosities at ages of the order of 10<sup>5</sup>-10<sup>6</sup> yrs (Ögelman 1995). For a neutron star born with a rotation period of the order of 10 ms, (as typically inferred from the P and $`\dot{P}`$ values of young pulsars), to have spun down to the 8.4 s period of RXJ0720.4-3125 in 10<sup>5</sup> yrs as a rotating magnetic dipole would require a mean spindown rate $`\dot{\mathrm{\Omega }}`$ of the order of 10<sup>-13</sup> rad s<sup>-2</sup>, and a magnetic field of the order of 10<sup>14</sup> G or more. (Heyl & Hernquist 1998). There is an alternative source of the thermal luminosity which takes over at $``$ 10<sup>5</sup>-10<sup>6</sup> yrs, after the initial cooling, and lasts longer than the cooling luminosity: There will be energy dissipation (frictional heating) in a neutron star being spun down by some external torque. The rate of energy dissipation is given by (Alpar et al. 1984, Alpar, Nandkumar & Pines 1985) $$\dot{E}_{diss}=I_p\omega \|\dot{\mathrm{\Omega }}\|$$ (1) where I<sub>p</sub> is the moment of inertia of some component of the neutron star whose rotation rate is faster than that of the observed crust by the amount $`\omega `$. $`\dot{E}_{diss}`$ will supply the thermal luminosity of a non-accreting neutron star at ages greater than $``$ 10<sup>6</sup> years (Umeda et al. 1993), as the cooling luminosity rapidly falls below $`\dot{E}_{diss}`$ after the transition from neutrino cooling to surface photon cooling. Among the radio pulsars with X-ray emission (Ögelman 1995, Becker & Trümper 1997), observations (Alpar et al. 1987, Yancopoulos, Hamilton & Helfand 1994) of the pulsar PSR 1929+10, whose spin-down age is 3.1 $`\times `$ 10<sup>6</sup> yrs, provide an upper limit to the thermal luminosity which yields I$`{}_{p}{}^{}\omega <`$ 10<sup>43</sup> gm cm<sup>2</sup> rad s<sup>-1</sup> (Alpar 1998). A lower limit to $`\dot{E}_{diss}`$ can be obtained from the parameters of large pulsar glitches (Alpar 1998). The consistency of the observed glitch parameters of all pulsars with large glitches ($`\mathrm{\Delta }\mathrm{\Omega }/\mathrm{\Omega }`$ 10<sup>-7</sup>) and measured second derivatives $`\ddot{\mathrm{\Omega }}`$ of the rotation rate (Shemar & Lyne 1996) as well as the statistics of the large glitches (Alpar & Baykal 1994) support the hypothesis that large glitches are a universal feature of pulsar dynamics. The current phenomenological description (Alpar et al. 1993, Alpar 1998) of large glitches entails angular momentum exchange between the crust and an interior component (a pinned superfluid in current models) which rotates faster than the crust by the lag $`\omega `$. The typical glitch related change in relative rotation rate of the crust and interior, $`\delta \omega `$ 10<sup>-2</sup> rad s<sup>-2</sup>, inferred from the common behaviour of all pulsars with large glitches, must be less than the lag $`\omega `$. Using values of I$`{}_{p}{}^{}`$ 10<sup>43</sup> gm cm<sup>2</sup> inferred from the detailed postglitch timing measurements available for the Vela pulsar glitches (Alpar et al. 1993) we obtain the lower bound I$`{}_{p}{}^{}\omega >`$I$`{}_{p}{}^{}\delta \omega `$=10<sup>41</sup> gm cm<sup>2</sup> rad s<sup>-1</sup>. This lower bound is independent of the details of the glitch models and rests only on the assumption that the large glitches involve angular momentum exchange within the neutron star. Further, as the mode of angular momentum transfer inside the star depends only on neutron star structure, the same parameter I$`{}_{p}{}^{}\omega `$ would determine the energy dissipation rates in all neutron stars under external torques, also when the source of the external torque is not magnetic dipole radiation. The upper and lower bounds together imply $$10^{41}\|\dot{\mathrm{\Omega }}\|ergs^1<L=\dot{E}_{diss}<10^{43}\|\dot{\mathrm{\Omega }}\|ergs^1.$$ (2) The luminosities of RXJ0720.4-3125, RXJ185635-3754 and the other DTNs (Treves et al. 2000) give $$\|\dot{\mathrm{\Omega }}\|10^{13}10^{10}rads^2.$$ (3) The measured spindown rates of AXPs and SGRs all fall in this range, providing another similarity, in addition to the periods, between the DTNs, AXPs and SGRs. AXP spindown rates actually suggest that the DTN spindown rate may be closer to $`\|\dot{\mathrm{\Omega }}\|`$ 10<sup>-12</sup> rad s<sup>-2</sup>, which means that $`\dot{E}_{diss}`$ is closer to the upper limit in Eq.(2). With the 8.4s period of RXJ0720.4-3125 these spindown rates imply surface magnetic fields in excess of 10<sup>14</sup> G if magnetic dipole spindown is assumed. Are there other spindown mechanisms that will give high spindown rates, larger than 10<sup>-12</sup> rad s<sup>-2</sup>, with 10<sup>12</sup> G magnetic fields typical of the canonical radio pulsars and of the accreting neutron stars with observed cyclotron lines? ## 3 Propeller Spindown High spindown rates, larger than 10<sup>-12</sup> rad s<sup>-2</sup>, can indeed be expected for neutron stars with conventional 10<sup>12</sup> G fields under the typical spindown torques for certain phases of accreting sources. For the AXPs accretion is a possibility that has already been explored (Mereghetti & Stella 1995, van Paradijs, Taam & van den Heuvel 1995, Ghosh, Angelini & White 1995) and will be pursued below. In connection with the DTNs, we note that a neutron star subject to mass inflow will experience high spindown rates even when the inflowing mass is not accreted, because of the star’s centrifugal barrier (the propeller effect) (Illarionov & Sunyaev 1975). We propose that RXJ 0720.4-3125 and the other DTNs are neutron stars with magnetic fields of the order of 10<sup>12</sup> G, spinning down under propeller torques. There is no accretion yet in the propeller phase. The luminosities of the DTNs are produced by energy dissipation in the neutron star. These luminosities are dimmer, by several orders of magnitude, than the accretion luminosities that would have been produced if the mass inflow were accreted. The propeller torques depend on the magnetic moment of the neutron star and on the rate of mass inflow. Using the spindown rates estimated from the energy dissipation luminosities together with order of magnitude expressions for propeller torques, we show that the neutron star has a dipole magnetic field of the order of 10<sup>12</sup> G. The estimated range of mass inflow rates coincides with the range of mass accretion rates inferred from the X-ray luminosities of the accreting AXPs and SGRs, suggesting that the DTNs are in the propeller stage and the AXPs and SGRs in the accretion phase under similar mass inflow circumstances. Occasional quiescent states of some X-ray transients are likely examples of the propeller phase (Cui et al 1997, Zhang et al 1998, Campana et al 1998, Menou et al. 1999). Sources that are persistently in the propeller phase have not been detected previously. This is understandable since they are not lit up with an accretion luminosity. Thus RXJ 0720.4-3125, RXJ185635-3754 and other DTNs are the first observed examples of propellers, observed only through their dissipation luminosities, only as the nearest members of the DTN class, and only by ROSAT, in view of their surface temperatures in the soft X-ray band. In the propeller phase there is a spindown torque exerted on the neutron star, through the interaction of its magnetosphere with the ambient material. The order of magnitude of this spindown torque is $$N=I\|\dot{\mathrm{\Omega }}\|=\mu ^2/r_{A}^{}{}_{}{}^{3}$$ (4) where $`\mu =BR^3`$ is the magnetic moment of the neutron star with surface dipole magnetic field B and radius R, and $$r_A=9.85\times 10^8cm\mu _{30}^{}{}_{}{}^{4/7}\dot{M}_{15}^{}{}_{}{}^{2/7}m^{1/7}$$ (5) is the Alfven radius (Lipunov 1992, Frank, King & Raine 1992). Here $`\mu _{30}`$ denotes the magnetic moment in units of 10<sup>30</sup> G cm<sup>3</sup>, $`\dot{M}_{15}`$ is the mass inflow rate in units of 10<sup>15</sup> gm s<sup>-1</sup> and m is the neutron star mass in solar mass units. The propeller phase will last until accretion starts when a critical rotation period is reached. This critical period is of the order of, but somewhat smaller than the equilibrium rotation period $$P_{eq}=16.8s\mu _{30}^{}{}_{}{}^{6/7}\dot{M}_{15}^{}{}_{}{}^{3/7}m^{5/7}$$ (6) which obtains when the star’s rotation rate equals the Keplerean rotation rate of ambient matter at the Alfven radius; that is, when the corotation radius r<sub>c</sub> = (GM)$`{}_{}{}^{1/3}\mathrm{\Omega }_{}^{2/3}`$ becomes equal to r<sub>A</sub> . In the propeller phase the source settles to an asymptotic spindown towards the equilibrium. A source is most likely to be observed during its asymptotic phase. Since in the asymptotic regime r<sub>c</sub> is close to r<sub>A</sub> we can estimate the magnetic moment of a propeller from its spindown rate, and rotation rate, without knowing the mass inflow rate, simply by substituting r<sub>c</sub> in Eq.(4). This leads to $$\mu =\frac{(I\|\dot{\mathrm{\Omega }}\|GM)^{1/2}}{\mathrm{\Omega }}.$$ (7) We take the DTN RXJ0720.4-3125, with P = 8.4 s, as our typical example (Haberl et al. 1997). The other DTN with a reported rotation period, RXJ 0420.0-5022, P = 22.7 s (Haberl, Pietsch & Motch 1999), yields similar results. For RXJ0720.4-3125 we obtain $`\mu _{30}(0.353.5)`$ using the bounds on the spindown rate, Eq.(2). Thus if RXJ0720.4-3125 is a propeller it has a conventional magnetic field, $$3.5\times 10^{11}G<B<3.5\times 10^{12}G$$ (8) Using Eq.(6) and setting the observed period to be equal to the equilibrium period, we can now obtain bounds on the mass inflow rates: $$\dot{M}_{15}=(16.8s/P_{eq})^{7/3}(\mu _{30})^2m^{5/3}.$$ (9) This yields $`\dot{M}_{15}(0.6262)`$ for the mass inflow rate for RXJ0720.4-3125. Model parameters obtained from Eqs. (7) and (9) are given in Table 1, assuming $`\|\dot{\mathrm{\Omega }}\|=`$ 2.6 $`\times `$ 10<sup>-12</sup> rad s<sup>-2</sup> from Eq.(2) for RXJ0720.4-3125, and $`\|\dot{\mathrm{\Omega }}\|=`$ 10<sup>-11</sup> rad s<sup>-2</sup> for RXJ 0420.0-5022. For most DTNs and for all RQNSs, which, as we propose below, are also propeller candidates, we do not have observed rotation periods. For these sources, setting the rotation period to be equal to the equilibrium period, $$\dot{M}_{15}\mu _{30}^{}{}_{}{}^{1/3}\|\dot{\mathrm{\Omega }}\|_{12}^{}{}_{}{}^{7/6}m^{1/2}I_{45}^{}{}_{}{}^{7/6}$$ (10) from Eqs.(6) and (7). With the range of $`\|\dot{\mathrm{\Omega }}\|`$ obtained from DTN luminosities (Eqn. (3)), $`\dot{M}_{15}(0.1200)`$ is obtained for these DTNs for $`\mu _{30}`$ = 1. If the source of mass inflow is not depleted, the propeller phase is followed by an accretion phase starting at some critical period. The neutron star will continue to spin down, now as an accreting source, as its period evolves from the critical period towards the equilibrium period. Similarity of the rotation periods of the AXPs with the 8.4 s period of the DTN RXJ0720.4-3125 suggests that the approach to rotational equilibrium is an asymptotic process extending through the propeller phase on to the accretion phase. These sources are asymptotically close to their individual equilibrium rotation periods, which lie in a narrow range determined by their magnetic moments and mass inflow histories. The mass accretion rates inferred from the X-ray luminosities of the AXPs and SGRs all fall in the range $`\dot{M}_{15}1100`$. The corresponding magnetic moments for accreting sources near rotational equilibrium are in the 10<sup>12</sup> G range given the 5-15 s rotation periods: $$\mu _{30}=(P_{eq}/16.8s)^{7/6}\dot{M}_{15}^{}{}_{}{}^{1/2}m^{5/6}.$$ (11) To summarize, (i) the spindown rates inferred for the DTN sources through the interpretation of their luminosities as due to energy dissipation (Eqs.(1)-(3)) agree with the observed spindown rates for the AXPs and SGRs to order of magnitude. (ii) Estimates of the magnetic moment, from the spindown rates for the DTNs as propellers close to rotational equilibrium, coincide in the same range as the magnetic moments inferred from the mass accretion rates implied by the X-ray luminosities of AXPs and SGRs as accretors asymptotically close to rotational equilibrium. This is the conventional 10<sup>12</sup> G range of most young neutron stars observed as rotation powered pulsars and in high mass X-ray binaries. (iii) Furthermore, the mass inflow rates inferred for the DTNs from the assumption that they are close to rotational equilibrium, agree with the accretion rates of the AXPs and SGRs. These observations strongly support the unifying hypothesis that the sources are related instances of asymptotic spindown near similar equilibrium periods. ## 4 The Asymptotic Spindown Regime We now turn to a simple model for the asymptotic spindown. A neutron star in the presence of inflowing matter experiences both spindown and spin-up torques. The long term evolution is determined by the balance between these, described by a function $`n(\mathrm{\Omega }/\mathrm{\Omega }_{eq}\omega _c)`$ (Ghosh & Lamb 1991) which goes through zero when $`\mathrm{\Omega }=\omega _c\mathrm{\Omega }_{eq}`$ where $`\omega _c`$ is of order one, and $`\mathrm{\Omega }_{eq}=2\pi /P_{eq}`$. In accretion from a disk, the relative specific angular momentum brought in by the accreting material to spin the neutron star up is $`[(GMr_A)^{1/2}\mathrm{\Omega }r_{A}^{}{}_{}{}^{2}]`$. Since the dimensional torque is $`\mu ^2/r_{A}^{}{}_{}{}^{3}\dot{M}(GMr_A)^{1/2}`$, spindown near equilibrium can be modelled as $$\dot{\mathrm{\Omega }}=(\mu ^2/Ir_{A}^{}{}_{}{}^{3})(1\mathrm{\Omega }/\mathrm{\Omega }_{eq})=(\mathrm{\Omega }_{eq}\mathrm{\Omega })/t_0$$ (12) where t<sub>0</sub> = $`\mathrm{\Omega }_{eq}Ir_{A}^{}{}_{}{}^{3}/\mu ^2`$. Here the zero of the torque (the end of the spindown era) is taken to be at $`\mathrm{\Omega }=\mathrm{\Omega }_{eq}`$ rather than $`\omega _c\mathrm{\Omega }_{eq}`$, for simplicity. The AXPs and SGRs must then be sources that have emerged from propeller spindown and started to accrete while still under a net spindown torque as they have not yet reached rotational equilibrium. Another simplifying assumption we make is that the mass inflow rate is constant in each source. In a more realistic model the mass inflow rate should be time dependent, especially if its source is a remnant disk with a finite lifetime, rather than a binary companion. A decaying $`\dot{M}`$ leads to a natural explanation for the prevalence of spindown in an accreting source: the equilibrium period increases as $`\dot{M}`$ decreases, and the source will be spinning down as its period tracks the equilibrium period (Ghosh, Angelini & White 1997, Chatterjee, Hernquist & Narayan 1999), answering the criticism of accretion models raised by Xi (1999). Chatterjee, Hernquist & Narayan (1999) have employed an $`\dot{M}`$ decaying as a power law in time, corresponding to the viscous evolution of a thin disk (Mineshige, Nomoto & Shigeyama 1993, Canizzo et al 1997). The actual situation warrants taking into account propeller boundary conditions. The propeller activity could probably support a thick disk or corona. It would also lead to a different characteristic time dependence of $`\dot{M}`$. Here we take a constant $`\dot{M}`$ to represent an average of the actual $`\dot{M}`$(t). Even under constant mass inflow rate, the torque history must be more complicated. In X-ray binaries, the occurrence of spindown and spinup torques of comparable magnitudes, and torque reversals on all timescales signals the importance of electromagnetic torques. The real torques could lead to initial power law decays of $`\mathrm{\Omega }`$ as a function of time. Whatever the form of the initial spindown may be, once the rotation rate is close to the equilibrium value, the factor $`(\mathrm{\Omega }_{eq}\mathrm{\Omega })`$ dominates the asymptotic evolution, which becomes an exponential decay. Here we employ the simple asymptotic model of Eq.(12) assuming that the magnetic moment is constant, and that the mass inflow rate is also a constant representing the long term average $`\dot{M}`$. The results are given in Table 1 and Figure 1. With constant $`\mu `$ and $`r_A`$ , the spindown leads $`\mathrm{\Omega }`$ towards $`\mathrm{\Omega }_{eq}`$: $$\mathrm{\Omega }(t)=[\mathrm{\Omega }(0)\mathrm{\Omega }_{eq}]exp(t/t_0)+\mathrm{\Omega }_{eq}.$$ (13) This solution is to describe the spindown through both the propeller and the accretion phases, with mass accretion starting at some $`\mathrm{\Omega }_{acc}<\mathrm{\Omega }_{eq}`$. From Eqs.(5),(6) and (12) we obtain $$t_0=1.3\times 10^{12}sI_{45}\dot{M}_{15}^{}{}_{}{}^{1}m^{2/3}\mathrm{\Omega }_{eq}^{}{}_{}{}^{4/3}.$$ (14) Substituting the expression for t<sub>0</sub> in the spindown equation, Eq.(12), $`\mathrm{\Omega }_{eq}`$ can be obtained for each source with known $`\mathrm{\Omega }`$ and $`\dot{\mathrm{\Omega }}`$ by solving $$\mathrm{\Omega }_{eq}=\mathrm{\Omega }+1.3\times 10^{12}s\dot{\mathrm{\Omega }}I_{45}\dot{M}_{15}^{}{}_{}{}^{1}m^{2/3}\mathrm{\Omega }_{eq}^{}{}_{}{}^{4/3}.$$ (15) Once the equilibrium rotation rate is estimated, the magnetic moment $`\mu `$ can be obtained from Eq.(11). The time since the beginning of the propeller phase can be obtained from Eq.(13) for a given initial rotation period, which can be taken to be P(0) = 0.01 s for a newborn neutron star. The age t estimated from Eq.(13) is: $$t=t_0log[(\mathrm{\Omega }(0)\mathrm{\Omega }_{eq})/(\mathrm{\Omega }\mathrm{\Omega }_{eq})].$$ (16) Sample solutions are given in Table 1 for the AXPs, the SGRs with measured P and $`\dot{P}`$ and for the DTNs with measured periods. The AXPs show changes in spindown rate, by up to an order of magnitude, on timescales of several years. Long term average values of the spindown rate are used, as appropriate for the model to describe the long term average trend. Mass inflow rates are inferred from the observed luminosities taken as accretion luminosities onto neutron stars, $`\dot{M_{acc}}`$ = RL/GM . The luminosity has been observed to change by as much as a factor of 15 in AXPs (Torii et al. 1998). Furthermore, the mass inflow rate from a remnant disk will be decreasing as a function of time, so that the constant representative $`\dot{M}`$ appropriate for our simple model should be larger than the present time $`\dot{M_{acc}}`$ inferred from the acretion luminosity. For 1E1048.1-5937, the current $`\dot{M_{acc}}`$ leads to P<sub>eq</sub> = 64 s., while for 1E2259+586 the current $`\dot{M_{acc}}`$ gives a relaxation time t<sub>0</sub> that may be larger than the age of the associated SNR. For these two sources, taking examples of mass inflow rates $`\dot{M}>\dot{M_{acc}}`$, one can obtain solutions in agreement with the asymptotic model. These solutions are given in Table 1. Three out of six AXPs provide direct evidence of youth through their likely supernova associations. These AXPs are situated close to the centers of the SNRs. The associations with SNRs are not as certain for the SGRs: The SGRs are at edges of the SNRs, requiring velocities $``$ 1000 Km s<sup>-1</sup> if they were born at the center (for a recent discussion of the SNR associations of the AXPs and SGRs, and references to the individual associations, see Gaensler 1999). For both AXPs and SGRs, the associated SNRs are not typical; but rather of a particularly dense morphology. This point was realized and elaborated by Marsden et al (1999) in their discussion of propeller and accretion scenarios, as evidence that these classes of sources are related through similarities in their environments rather than a common intrinsic property like very strong magnetic fields. Here we adopt the possibility that a remnant disk around the neutron star, formed from the debris of the core collapse in a supernova explosion (Chevalier 1989, Mineshige, Nomoto & Shigeyama 1993), provides the mass inflow. The propeller and accretion phases would terminate when the supply of $`\dot{M}`$ is depleted. ## 5 The Fallback Mass Inflow and the Signature of the Neutron Star Born in a Supernova Does the narrow range of observed rotation periods indicate a very restricted range of $`\dot{M}`$ ? The range of estimates for the sources in Table 1 extends from 10<sup>15</sup> gm s<sup>-1</sup> to 10<sup>17</sup> gm s<sup>-1</sup>. To estimate the parameter range corresponding to the narrow range of equilibrium periods, we note that the magnetic moment and mass inflow rate or equivalently, the magnetic moment and equilibrium rotation rate are independent parameters. For magnetic moments in the range 5 $`\times `$ 10<sup>11</sup> G - 5 $`\times `$ 10<sup>12</sup> G, the range of mass inflow or accretion rates which would lead to equilibrium periods of 5-15 s extends over three orders of magnitude, up to about half the Eddington rate: $$3.2\times 10^{14}gms^1<\dot{M}<4.2\times 10^{17}gms^1.$$ (17) This is not a very restricted range. Fig.1 shows the mass inflow rate as a function of the equilibrium period for 0.5 $`<\mu _{30}<`$ 5. Also shown are the model solutions for equilibrium periods and mass inflow rates for all the sources presented in Table 1. The narrow range of observed periods lies between the vertical lines. We note that the narrow range of observed and equilibrium periods corresponds to a rather wide range of $`\dot{M}`$, extending over more than two orders of magnitude. Sample solutions for different sources are scattered in $`\dot{M}`$ and $`\mu `$, as they should be, since the magnetic moments of neutron stars should not be correlated with the mass inflow from their environments. This scatter indicates that the model does not require an artificial correlation between these parameters. The mass inflow rates we obtain are all above $``$ 10<sup>14</sup> gm s<sup>-1</sup>. This means that for lower $`\dot{M}`$ the propeller phase never starts, and the neutron star promptly starts life as a radio pulsar. At the other extreme extended exposure to the highest mass inflow rates may lead to accumulation of enough circumstellar material such that the optical thickness to electron scattering washes away any beaming of the X-rays from the neutron star and the rotation period is not detected (Lamb et al 1985, Ghosh, Angelini & White 1997). If all supernovae left neutron stars that went through an AXP phase under mass inflow from the debris, then we would expect at least 100 such objects for a lifetime $`>`$ 10<sup>4</sup> yrs and galactic supernova rate of 10<sup>-2</sup>, since the AXPs would be observable from all galactic distances. That we only observe a few indicates that such sources rarely reach the accretion phase. The number of DTNs must be much larger, since we see several within a distance of the order of 100 pc. In the Galaxy we would expect about 10<sup>4</sup> DTNs, and a lifetime of the order of 10<sup>6</sup> yrs or longer, for a birth rate of 10<sup>-2</sup> yr<sup>-1</sup> or lower. The much smaller number of the AXPs may be due to the depletion of the supply of mass inflow in most sources by the end of the propeller phase. If the mass inflow is depleted before the end of the propeller phase, subsequent accretion phases never occur. This may be the reason why the numbers of AXPs and SGRs are much less than the number of DTNs. By the same token, mass inflow rates $`<`$ 3.2 $`\times `$ 10<sup>14</sup> gm s<sup>-1</sup> are not inferred: at lower mass inflow rates the duration of the propeller phase exceeds the lifetime of the mass inflow source. The rare AXPs and SGRs are thus the young objects born with conventional magnetic fields but in circumstances of large enough mass inflow towards the neutron star, so that the time t<sub>0</sub>, which scales with $`\dot{M}^1`$, is short compared to the lifetime of the mass supply: These are the rare sources which have gone through the propeller phase in time to start accretion before the matter supply is depleted. The more common source is the DTN, in a propeller phase of duration 10<sup>6</sup> yrs, surviving the disappearance of the supernova remnant. While it is reasonable to expect that the lifetime of the mass supply increases at the lower mass inflow rates, one needs a model of the disk evolution to make the comparison with the propeller spindown quantitative. If DTNs are born at a rate comparable to the total supernova rate, then a significant fraction of SNRs must leave the neutron star under conditions of low enough $`\dot{M}`$. If all AXPs are to be descendants of DTNs, either the propeller phase lasts 1000 times longer than the AXP accretion phase, or only 10<sup>-3</sup> of DTNs turn into AXPs, according to the comparison of total numbers of these two types of sources in the Galaxy. Taking into account the evidence for possible association with supernova remnants, and the agreement between SNR ages and the age estimates the asymptotic model, mass supply from a debris disk around a single neutron star seems to be the likely scenario, rather than mass supply from a binary companion. The SGRs are like the AXPs in their X-ray properties. Both classes of sources are detectable throughout the galaxy, and their comparable numbers suggest the SGRs are also in the same rare or relatively short beginning accretor phase as the AXPs. The claimed associations of most SGRs with SNRs are less certain than the AXP-SNR associations. Unlike the AXPs the SGRs are located at the outskirts of their candidate SNR associations. Nevertheless these SNRs are of the same atypical dense class as the SNRs associated with the AXPs (Marsden et al 1999). For the SGRs the magnetar hypothesis is on a different, stronger footing than is the case for the AXPs, in that it provides the energy store and a detailed dyamical model for the soft gamma ray bursts. While the present work does not address the mechanisms of the gamma ray burst phenomenon, it is nevertheless intriguing that as Wang and Robertson (1985) have noted propellers (and probably their descendants the early accretors) can support relativistic particle luminosities and gamma ray production in the surrounding accumulated matter. The energy source could be sporadic accretion of accumulated circumstellar mass released through instabilities. The remnant disk scenarios were first recalled in connection with the problem of planet formation around pulsars (Lin, Woosley & Bodenheimer 1991), and disk instabilities involving planet-like masses (Mineshige, Nomoto & Shigeyama 1993) are consistent with the luminosities of the SGRs. It may not be unreasonable to speculate that SGRs are also a class of neutron stars with conventional fields undergoing disk instabilities. Why is it that the AXPs are observed as pulsars while the LMXBs are not ? The important single exception is the recently discovered source XTE J1808-369. The explanation for the rarity of LMXB pulsars has been that comptonization by circumstellar material will destroy pulses by washing out the beams emerging from the neutron star if the material is optically thick to electron scattering (Lamb et al. 1985). This was noted in connection with the AXPs by Ghosh, Angelini & White (1997). The present picture is consistent with this: the AXPs are observed as pulsars because they are in the beginning stages of accretion, and the corona around them does not have significant optical thickness $`\tau _{es}>`$ 1 to destroy the beaming. Work in progress to test this hypothesis by modelling the AXP spectra as unsaturated comptonization and estimating the optical thicknesses will be reported separately. The radio quiet neutron stars found in centers of the supernova remnants Cas A, Puppis A, RCW 103 and 296.5+10 (Chakrabarty et al. 2000, Gaensler, Bock & Stappers 2000, Petre, Becker & Winkler 1996, Gotthelf, Petre & Hwang, 1997, Gotthelf, Petre & Vasisht 1999, Mereghetti, Bignami & Caraveo 1996, Vasisht et al. 1997) may be the propeller sources for which the cumulative effect of the mass inflow has indeed set up a corona of $`\tau _{es}>`$ 1 around the neutron star. The source properties are very similar. Variability of the flux (Gotthelf, Petre & Vasisht 1999) favours accretion or propeller driven energy dissipation rather than the cooling luminosity of an isolated neutron star. The low luminosities, L$`{}_{x}{}^{}`$ 10<sup>33</sup> \- 10<sup>34</sup> erg s<sup>-1</sup> suggest that these sources are in the propeller phase. Accretion would imply $`\dot{M}`$ 10<sup>13</sup> \- 10<sup>14</sup> gm/s, which can accrete onto the neutron star only if the rotation period is longer than 50 s (Eq.(6)) for B $``$ 10 <sup>12</sup> G. But then at these low accretion rates there is no reason that we should not see pulses at the rotation period. Blackbody effective areas of 1 km<sup>2</sup> also enhance the expectation of coherent pulsations at the rotation period from these sources while upper limits of 13% and 25% have been obtained for the RQNSs in RCW103 and Cas A respectively (Gotthelf, Petre & Vasisht 1999, Chakrabarty et al. 2000). Taking the RQNSs to be propeller sources with L<sub>x</sub> = $`\dot{E}_{diss}`$, as we did for the DTNs, leads to $`\|\dot{\mathrm{\Omega }}\|`$10<sup>-10</sup> \- 10<sup>-8</sup> rad s<sup>-2</sup>. Since the rotation periods of the RQNSs are not known, we use Eq.(10) with B $``$ 10 $`{}_{}{}^{1}2`$ G, and obtain $$\dot{M}>2\times 10^{17}gms^1.$$ (18) For the RQNSs the mass inflow rates $`\dot{M}`$ are indeed higher than those inferred for the DTNs and AXPs, consistent with the proposal that the radio quiet neutron stars are the high $`\dot{M}`$ end of our spectrum of neutron stars under mass inflow from remnant disks. Model parameters for the RQNS in Cas A are displayed in Table 1. The luminosity is typical of the RQNSs. As a propeller with luminosity arising from energy dissipation, taking $`\|\dot{\mathrm{\Omega }}\|>`$ 3 $`\times `$ 10<sup>-10</sup> rad s<sup>-2</sup> with $`\mu _{30}`$ = 1 and obtain $`\dot{M}>`$ 7.4 $`\times `$ 10<sup>17</sup> gm s<sup>-1</sup>, in agreement with the interpretation here. P$`{}_{eq}{}^{}<`$1 s , corresponding to the high mass inflow rate. The unobserved periods of the RQNSs are smaller than the observed range of periods of the AXPs and DTNs. ## 6 Luminosity of the Disk For both the DTNs and the AXPs there will be another source of energy dissipation in the disk or circumstellar material, due to accretion down to $`r_A`$: $$L(r_A)=3/2GM\dot{M}/r_A3/2(GM)^{2/3}\mathrm{\Omega }^{2/3}\dot{M}1.3\times 10^{33}ergs^1\dot{M}_{15}m^{2/3}P^{2/3}$$ (19) Near rotational equilibrium, L(r<sub>A</sub>) = I$`\mathrm{\Omega }\dot{\mathrm{\Omega }}`$, the power expended on the star by the spindown torque. For the AXPs this luminosity is smaller than the accretion luminosity by a factor $$L(r_A)/L3/2R/r_A(\dot{M}/\dot{M}_{acc})10^2(\dot{M}/\dot{M}_{acc})m^{1/3}P^{2/3}$$ (20) noting that the mass inflow may not be all accreted even in the AXPs. For the DTNs L(r<sub>A</sub>) is actually larger than the observed luminosity $`\dot{E}_{diss}`$ by a factor $$L(r_A)/\dot{E}_{diss}(I\mathrm{\Omega }\|\dot{\mathrm{\Omega }}\|)/(I_p\omega \|\dot{\mathrm{\Omega }}\|)(10^210^4)\mathrm{\Omega }.$$ (21) If this luminosity is dissipated entirely at the boundary region, r $``$ r<sub>A</sub>, the effective temperature is: $$T(r_A)=(L(r_A)/4\pi r_{A}^{}{}_{}{}^{2}\sigma )^{1/4}9.5\times 10^4K\dot{M}_{15}^{}{}_{}{}^{1/4}P^{1/2}$$ (22) If the disk temperature lies in the euv range this would be extremely difficult to detect, but being nearby sources, DTNs with low neutral hydrogen column density might provide a chance of looking for a disk luminosity as a test of the present model. For the thin disk model employed by Chatterjee, Hernquist & Narayan (1999) calculations of the expected disk spectrum and lumionosity are reported by Perna, Hernquist & Narayan (2000). For T(r<sub>A</sub>) in the optical or IR, observational upper limits for disks around the AXPs (Coe and Pightling 1998, Hulleman et al 1999) pose a problem for thin disk models. The upper limits allow thin disks with disk outer radii $``$ 10<sup>10</sup> cm, which requires a very young age, t $``$ 100 yrs, in terms of the viscous evolution of an isolated thin disk. However several considerations suggest that the remnant disks may be consistent with the observational upper bounds. The propeller effect might well lead to boundary conditions that alter the disk evolution. The work done on the inflowing material by the propeller, at the rate I$`\mathrm{\Omega }\dot{\mathrm{\Omega }}`$, may lift and sustain the material at distances much larger than r<sub>A</sub>. The emerging radiation may have a lower effective temperature, from a larger effective area. The remnant disk may be a thick disk, and it may be enshrouded in a comptonizing corona, as we invoked above. Modelling of remnant disks and circumstellar material with propeller boundary conditions, and future observations on the spectra of DTNs, AXPs and radio quiet neutron stars will help to resolve these issues. ## 7 Discussion and Conclusions A unified picture is proposed to account for all neutron stars formed in supernovae, and to include RQNSs, DTNs and AXPs and perhaps SGRs. The salient features of this picture are: (i) The signature of a young neutron star depends on the presence and nature of the mass inflow of fallback material from the supernova explosion. The related classes of sources represent different pathways under different mass inflow rates and histories. (ii) All neutron stars are born with 10<sup>12</sup> G fields. (iii) Radio pulsars are formed if there is no mass inflow or not enough to protrude the light cylinder. Higher mass inflow rates do not allow radio pulsar activity, leading to DTNs, AXPs and RQNSs. (iv) The similarity in the rotation periods of these sources is not a coincidence but rather a consequence of the asymptotic approach to rotational equilibrium under a wide range of mass inflow rates. (v) With low $`\dot{M}`$ the spindown is too slow for the source to go through the propeller stage and reach the accretion stage before the circumstellar mass is depleted. These sources, observed as DTNs, are propellers for as long as 10<sup>5</sup>-10<sup>6</sup> yrs, and make up the most numerous class. (vi) The DTNs are the first observed examples of neutron stars in the propeller phase. (vii) A non-accreting neutron star under propeller spindown has a luminosity provided by energy dissipation inside the star. (viii) The observability of the rotation period in the AXPs but not in the LMXBs can be explained qualitatively in terms of comptonization as supported by an interpretation of their spectra. The RQNSs are the young neutron stars for which comptonization washes out beaming so that the rotation period is not observable. (vi) With low $`\dot{M}`$ the spindown is too slow for the source to go through the propeller stage and reach the accretion stage before the circumstellar mass is depleted. These sources, observed as DTNs, are propellers for as long as 10<sup>5</sup>-10<sup>6</sup> yrs, and make up the most numerous class. (vii) The DTNs are the first observed examples of neutron stars in the propeller phase. (viii) A non-accreting neutron star under propeller spindown has a luminosity provided by energy dissipation inside the star. Arguments that AXPs and DTNs are magnetars are based on the grounds that isolated neutron stars with ordinary 10<sup>12</sup> G fields cannot have spun down to $``$ 10 s periods within the estimated ages. The above picture obviates the need to postulate magnetars for AXPs and DTNs. For the SGRs the required energy budgets and dynamical arguments make a case for the magnetar hypothesis. Marsden et al (1999) have noted that AXPs and SGRs share a common morphology of their associated supernova remnants, which they link with the high density of the surrounding interstellar medium. These authors argue that the common properties of these classes of sources must therefore be due to a common set of environmental factors, rather than intrinsic neutron star properties like superstrong magnetic fields. Whether the large spindown rates observed are due to spindown by interaction with ambient matter, as proposed here, as well as by Chatterjee, Hernquist & Narayan (1999) and Marsden et al. (1999) or to spindown by a magnetar can be decided by detailed analysis of the fluctuations (noise) in the spindown process, as the timing noise characteristic of accretion powered neutron stars is quite distinguishable from the timing noise in the typically much quieter isolated rotation powered pulsars. This analysis will require frequently sampled timing observations of the AXPs and SGRs. The recently reported extended quiet spindown phases in 1E2259+586 and 1RXS J170849.0-400910 (Kaspi, Chakrabarty & Steinberger 1999) is not sufficient to conclude that this source is undergoing spindown under magnetic dipole radiation since the quiet spindown phase was preceded by episodes of higher spindown rates with large timing noise strengths (Baykal & Swank 1996). Whether such changes in the torque repeat periodically as foreseen for a precessing magnetar (Melatos 1999) remains to be checked in future timing observations. A well known accreting source, the LMXB 4U1626-67, has exhibited similar quiet episodes as well as intervals of strong timing fluctuations that are typical of accreting sources (Chakrabarty et al. 1997). It is encouraging from the unified point of view of this work that during the quiet spindown epoch that would allow the dtection of glitches, 1E 2259+586 has indeed exhibited a glitch that is very similar to the glitches of the Vela pulsar and other radio pulsars (Kaspi, Lackey & Chakrabarty 2000). If the AXPs are confirmed as accreting sources, the glitch from 1E2259+586 will constitute strong support for our starting hypothesis that all neutron stars have the same internal dynamics and the associated energy dissipation rates. Observation of a spindown rate in the expected range from RXJ0720.4-3125 or periods and spindown rates from other DTNs would constitute strong evidence for the propeller hypothesis. 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warning/0005/cond-mat0005157.html	ar5iv	text	# Observation of Supershell Structure in Alkali Metal Nanowires \[ ## Abstract Nanowires are formed by indenting and subsequently retracting two pieces of sodium metal. Their cross-section gradually reduces upon retraction and the diameters can be obtained from the conductance. In previous work we have demonstrated that when one constructs a histogram of diameters from large numbers of indentation-retraction cycles, such histograms show a periodic pattern of stable nanowire diameters due to shell structure in the conductance modes. Here, we report the observation of a modulation of this periodic pattern, in agreement with predictions of a supershell structure. \] Metallic nanowires clearly exhibit quantum properties in their conductance and structure. At the level of a few atoms in cross-section the conductance through metallic nanowires can be described in terms of a finite number of quantum modes, and it has been shown that this number is determined by the number of valence orbitals of the metal atoms involved . For monovalent free-electron metals such as gold and, in particular, the alkali metals the conductance in the smallest contacts evolves roughly through the successive opening of distinct conductance channels as the contact size increases . This quantum character of electrical transport was already inferred from histograms of conductance values for atomic-size point contacts, which show that the contacts have a preference for multiples of the conductance quantum $`G_0=2e^2/h`$ , after correction for a small series resistance. It had been suggested that the formation of quantum modes in nanowires should not only determine the conductance but also the cohesive energy . Recently, we have shown that the stability of nanowires in the range of cross-sections up to about 130 atoms is determined by electronic shell structure for quantum modes . The shell structure was observed in conductance histograms, which were measured by many times indenting and retracting two metal electrodes by means of a mechanically controllable break junction (MCBJ) . The key evidence for shell structure is a regular spacing of diameters for wires with enhanced stability, where the diameters were obtained from the semiclassical expression for the conductance . The shell structure observed here is a close analogue of the shell structure observed for alkali metal clusters , and it is the same principle that also applies to electrons in atoms and to protons and neutrons in nuclei. For metal clusters produced in vapor jets in vacuum, and analyzed by mass selection, it was observed that clusters with certain ‘magic numbers’ of atoms, 8, 20 ,40 , 58, etc., are more abundant than others. This was explained by their enhanced stability due to the closing of the shells of electronic states, modeled as free electron waves confined to a spherical potential well. The magic numbers can even be obtained from a semiclassical expansion for the oscillating part of the density of single-particle levels , where the stable clusters are determined by those diameters for which a bouncing electron wave traveling along a closed classical path inscribed inside the spherical cluster walls obeys the Bohr-Sommerfeld quantization condition with the bulk metal Fermi wave vector, $`k_F`$. The possible trajectories are illustrated in Fig. 1. It was shown that for spherical clusters the triangular and square orbits, with indices $`(3,1)`$ and $`(4,1)`$, dominantly determine the magic number series. Since the two dominant orbits, triangle and square, lead to slightly different series of stable diameters, the interference between the two series gives rise to a beat pattern, known as supershell structure . For alkali metal clusters a single beat due to this effect has indeed been reported . Here, we report the observation of supershell structure in sodium nanowires. We show that, in contrast to clusters but in agreement with the theory for nanowires, the diametric orbit (labeled $`(2,1)`$ in Fig. 1) has a strong contribution. Due to the larger separation between the periods for diametric and higher orbits, several beat minima of the supershell structure can be observed. The experiment is performed using an MCBJ, modified to accommodate the very reactive alkali metals . A rectangular piece of sodium metal is cut, while immersed in paraffin oil for protection against oxidation, and fixed upon a phosphor bronze bending beam with four 1 mm screws. A notch is cut in the center of the sodium sample, and the assembly is taken out of the paraffin and quickly mounted in a sample holder in a three-point bending configuration. Current and voltage leads are connected to both sides of the notch and the sample holder is evacuated and cooled down to liquid helium temperature. By applying force to the bending beam the sodium sample is broken at the notch, by which two clean fracture surfaces are exposed. Atomic-size contacts can then be established by relaxing the force, using a piezoelectric element for fine control. A heater and thermometer permit controlling the temperature from 4.2 K to above 100 K, while the vacuum can remains immersed in liquid helium. The conductance of the contacts is measured with a four-point dc-voltage bias circuit. The signal goes via a current-to-voltage converter through a 16 bit, $`10^5`$ samples/s analog-to-digital converter to a pc-based controller. The software also drives the two halves of the sample into and out of contact by controlling the piezovoltage. Conductance values are automatically accumulated into histograms for typically over $`10^4`$ contact breaking cycles, with a resolution of about 10 bins per $`G_0`$. In order to reduce digitization noise a three bin wide smoothing function is applied to the histograms. At low temperatures and low voltage bias only four pronounced peaks at low conductance are seen in the histogram, near 1, 3, 5 and 6 $`G_0`$, which have been attributed to the successive occupation of distinct quantum modes . At higher conductance the histogram shows a number of rather wide hills which grow into sharp peaks as we increase the temperature . Fig. 2a shows an example of a histogram for sodium at a sample temperature of 90 K. The histogram is similar to the one presented in Ref. , but the modulation of the peak intensities is more pronounced here . The radius $`R`$ of the nanowires can be obtained from the semiclassical expression for the conductance , $$\frac{G}{G_0}=\left(\frac{k_FR}{2}\right)^2\left(1\frac{2}{k_FR}\right).$$ (1) The inset in Fig. 2a shows the radii, in units $`k_F^1`$, corresponding to the first 6 maxima against their sequential number. Indeed, for shell structure we expect regularly spaced peaks as a function of the radius. The modulation of the peak intensities and their periodicity is more clearly illustrated by subtracting a smooth background and plotting the histogram as a function of the radius, see Fig. 2b. Two beat minima are clearly visible at about $`k_FR7`$ and 11.5, and a third can be seen on the expanded scale in the inset. Similar modulations have been observed in histograms for potassium and lithium, which will be discussed elsewhere. The components in the periodic structure are separated by making a Fourier transform of the curve in Fig. 2b, which is shown in Fig. 3a. We observe two frequency components, one at 0.6–0.65 and the other at 0.8–0.9 $`(k_FR)^1`$. We will argue that these are the components of the supershell structure. The quantum modes in a nanowire give rise to an oscillating contribution in the density of electronic states as a function of electron energy at constant radius or as a function of the radius of the wire at the Fermi energy. The modes each form a one-dimensional band, with quantum numbers determined by the confinement in the two transverse dimensions. This leads to an oscillating structure in the thermodynamic potential . The latter determines the stability of a particular structure, which leads to an oscillating probability for nanowire diameters. Also the Sharvin conductance of the nanowires is modified by an oscillating contribution, but this contribution appears to be of secondary importance to the experiment . The leading terms in a semiclassical expansion of the oscillating part of thermodynamic potential are , $`\mathrm{\Omega }^{osc}={\displaystyle \frac{2\epsilon _Fk_F}{\sqrt{2\pi r^{\prime \prime }}r}}{\displaystyle \underset{M=2}{\overset{\mathrm{}}{}}}{\displaystyle \underset{Q=1}{\overset{M/2}{}}}{\displaystyle \frac{1}{M^{7/2}}}\left[\mathrm{sin}\left({\displaystyle \frac{\pi Q}{M}}\right)\right]^{3/2}\times `$ (2) $`\mathrm{cos}\left[2Mr\mathrm{sin}\left({\displaystyle \frac{\pi Q}{M}}\right)+{\displaystyle \frac{\pi }{2}}\left(M+{\displaystyle \frac{1}{2}}\right)\right],`$ (3) where $`r=k_FR`$, $`r^{\prime \prime }=^2r/z^2`$ is the curvature along the direction of the wire axis, and $`r`$ and $`r^{\prime \prime }`$ are taken at the narrowest cross-section, $`z=0`$. The number $`M`$ corresponds to the number of vertices in the semiclassical orbit (see Fig 1) and $`Q`$ is the winding number. In order to evaluate this expression we need to make some assumption about a smooth evolution of the curvature of the wire, but this will only affect the prefactor and not the oscillating structure itself. Following Ref. we assume a parabolic wire shape, $`r=r_0+4(Ar_o)(z/L)^2`$. This shape grows smoothly from an initial cylindrical wire of length $`L_0`$ and radius $`A`$ while maintaining a constant volume, so that the radius at the narrowest point can be expressed in terms of its length $`L`$ as $`r_0=(A/4)(\sqrt{30L_0/L5}1)`$. Substituting these expressions into Eq. 3 we calculate the Fourier transform of $`\mathrm{\Omega }^{osc}(r)`$ in the same range of $`r`$ as for the experimental data. The result is given in Fig. 3b, where we have included terms up to $`M=7`$ for $`Q=1`$ (the higher $`Q`$ give rise to lower amplitude harmonics and we concentrate on the fundamental components). The dominant peak is due to the diametric orbit $`(M,Q)=(2,1)`$ and the peaks for higher $`M`$ rapidly decrease in amplitude and their frequency converges at 1, marking the end of the first series ($`Q=1`$). In the Fourier transform we can only resolve two peaks in addition to the one for $`M=2`$, corresponding to the triangular ($`M=3`$) and square orbits ($`M=4`$). The experiment clearly shows the same groups of frequencies, one corresponding to the diametric orbit and one peak at the position of the triangular and square orbits, which cannot be individually resolved. Some of the higher frequency maxima with small intensities are probably partly due to harmonics of the main bands. The relative intensities of three calculated peaks cannot be directly compared to the theoretical spectrum since the latter results from the thermodynamic potential, which we do not measure directly. The experimental spectrum is derived from the conductance histogram and is expected to reflect the same components as the thermodynamic potential, but the intensities depend, amongst others, on kinetic factors for surface diffusion of atoms and the time available for reaching the proper minima in the potential. Further differences may arise from deviations in cylindrical symmetry of the wire and scattering of the electrons on residual surface roughness. In particular the later mechanism would favor the higher order orbits, since specular reflection increases for smaller angle of incidence with the surface. In contrast to cylindrical systems, for spherical systems the diametric orbit ($`M=2`$) is negligible as it has a significantly smaller degeneracy compared to the triangular ($`M=3`$) and squared ($`M=4`$) ones : For a given position of one vertex the triangle and square have an additional degree of freedom being the rotation around the normal to this point. The contribution of the higher order orbits decays as a high power of the index $`M`$ of the orbit. The closeness of the frequencies for the triangular and square orbits leads to a long beating period. The first node of the shell structure amplitude should be observed after approximately 12 periods of principal oscillations. The natural limit on the number of chemical elements makes the number of shells in the periodic table too small for supershell structure to be observable. The same applies to shell structure in atomic nuclei. Although the amplitude of such a high number of oscillations is greatly diminished, the first node of the beating pattern in cluster abundance spectra was observed in Refs. . For nanowires, the degeneracy of the diametric, triangle and square orbits is of the same order and the larger separation between the $`M=2`$ frequency and the higher ones make the beating pattern more readily visible in the nanowire shell structure as compared to the metal cluster experiments. Thus, for the first time a direct Fourier transform of the experimental data yields proof for the existence of semiclassical orbits responsible for the oscillations in the thermodynamic potential of the system as a function of the radius.
warning/0005/cond-mat0005156.html	ar5iv	text	# Statistical properties of the 2D attached Rouse chain ## Abstract We study various dynamical properties (winding angles, areas) of a set of harmonically bound Brownian particles (monomers), one endpoint of this chain being kept fixed at the origin 0. In particular, we show that, for long times $`t`$, the areas $`\{A_i\}`$ enclosed by the monomers scale like $`t^{1/2}`$, with correlated gaussian distributions. This is at variance with the winding angles $`\{\theta _i\}`$ around fixed points that scale like $`t`$ and are distributed according to independent Cauchy laws. Laboratoire de Physique Théorique et Modèles Statistiques. Université Paris-Sud, Bât. 100, F-91405 Orsay Cedex, France. Laboratoire de Physique Théorique des Liquides, UPMC, 4 place Jussieu 75252 Paris Cedex 05 In this paper, we will study the planar motion of a chain of $`n`$ harmonically bound brownian particles. This model is usually refered in the litterature as the Rouse chain and has shown to be historically very important in polymer science . We will consider such a chain attached at the origin 0 and examine some of its properties from the Brownian motion viewpoint. Representing a given configuration of the chain by a complex $`n`$-vector $`z`$ (the components $`z_i,`$ $`i=1,\mathrm{},n`$, are the complex coordinates of the particles), we consider the set of all the closed trajectories of length $`t`$, i.e. $`z(t)=z(0)`$, and this for all the starting configurations $`z(0)`$. Practically, we will not weight the starting configurations with any thermodynamical factor. We are aware that this approach is quite different from the one taken in polymer physics where, at $`t=0`$, the chain is supposed to be in equilibrium with the environment at some finite temperature $`T`$. $`A_j`$ and $`\theta _j`$ being the area enclosed by the $`j^{th}`$ particle and its winding angle around 0, our goal is to compute the joint probability distributions $`P(\{A_i\})`$ and $`P(\{\theta _i\})`$ for such trajectories. In order to make comparisons, we now recall some of the results concerning the planar Brownian motion. We first quote the area and winding angle distributions, respectively $`P(A)`$ (Lévy’s law ) and $`P(\theta )`$ (Spitzer’s law ) for a particle allowed to wander everywhere in the plane: $$P(A)=\frac{\pi }{2t}\frac{1}{\mathrm{cosh}^2\frac{\pi A}{t}}$$ (1) $$P(\theta )=\frac{2}{\pi \mathrm{ln}t}\frac{1}{1+\left(\frac{2\theta }{\mathrm{ln}t}\right)^2}$$ (2) (the last one holds, in the limit $`t\mathrm{}`$, for open curves, the final point being left unspecified). Those two laws were obtained more than 40 years ago and since that time many refinements have been brought. For instance, in , the authors pointed out the importance of the small windings occuring when the particle is close to 0. Excluding an arbitrary small zone around 0, they showed that the variance $`\theta ^2`$ becomes finite in contrast with the Spitzer’s result, eq.(2). On the other hand, for Brownian motion on bounded domains , the scaling variables in the limit $`t\mathrm{}`$, become, resp., $`A/\sqrt{t}`$ and $`\theta /t`$ with still an infinite variance $`\theta ^2`$. We close here this brief recall and start our chain study with the following set of coupled Langevin equations : $`\dot{z}_1`$ $`=`$ $`k(z_22z_1)+\eta _1`$ $`\dot{z}_l`$ $`=`$ $`k(z_{l+1}+z_{l1}2z_l)+\eta _l,2ln1`$ (3) $`\dot{z}_n`$ $`=`$ $`k(z_{n1}z_n)+\eta _n`$ where $`k`$ is the spring constant and $`\eta _m`$ ($`\eta _{mx}+i\eta _{my}`$) a gaussian white noise: $`\eta _m(t)`$ $`=`$ $`0`$ $`\eta _m(t)\eta _m^{}(t^{})`$ $`=`$ $`2\delta _{mm^{}}\delta (tt^{})`$ (4) Introducing the complex $`n`$-vector $`\eta `$, eq. (3) can be written in a matrix form: $$\dot{z}=k𝐌z+\eta $$ (5) where $`𝐌`$ is the tridiagonal $`(n\times n)`$ matrix: $$𝐌=\left(\begin{array}{ccccc}2& 1& 0& \mathrm{}& 0\\ 1& 2& 1& \mathrm{}& 0\\ 0& 1& 2& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& 0& 0& \mathrm{}& 1\end{array}\right)$$ with an inverse given by: $$𝐌^1=\left(\begin{array}{ccccc}1& 1& 1& \mathrm{}& 1\\ 1& 2& 2& \mathrm{}& 2\\ 1& 2& 3& \mathrm{}& 3\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ 1& 2& 3& \mathrm{}& n\end{array}\right)$$ The eigenvalues of $`𝐌`$ are: $$\omega _j=\mathrm{\hspace{0.17em}2}\left(1\mathrm{cos}\frac{\pi (2j1)}{2n+1}\right),1jn$$ (6) With the matrix $`\omega =\mathrm{diag}(\omega _i)`$, we can write: $`\omega `$ $`=`$ $`𝐑^1𝐌𝐑`$ (7) $`z`$ $`=`$ $`𝐑Z`$ (8) where $`𝐑`$ is an orthogonal matrix and the components of $`Z`$ are the normal coordinates. Let us call $`𝒫(z,z_0,t)`$ the probability for the chain to go from $`z_0`$ at $`t=0`$ to $`z`$ at time $`t`$. $`𝒫`$ satisfies a Fokker-Planck equation : $$_t𝒫=\left({}_{}{}^{t}_{z}^{}k𝐌z+^t_{\overline{z}}k𝐌\overline{z}+2^t_{\overline{z}}_z\right)𝒫$$ (9) where $`_z`$ (resp. $`_{\overline{z}}`$) is a $`n`$-vector of components $`_{z_i}`$ (resp. $`_{\overline{z}_i}`$) and $`{}_{}{}^{t}_{z}^{}`$ (resp. $`{}_{}{}^{t}_{\overline{z}}^{}`$) is the transpose of $`_z`$ (resp. $`_{\overline{z}}`$). The solution can be written in terms of a path integral ($`𝒟z𝒟\overline{z}=_{i=1}^n𝒟z_i𝒟\overline{z_i}`$): $`𝒫(z,z_0,t)`$ $`=`$ $`det\left(e^{tk𝐌}\right){\displaystyle _{z(0)=z_0}^{z(t)=z}}𝒟z𝒟\overline{z}\mathrm{exp}({\displaystyle \frac{1}{2}}{\displaystyle _0^t}{}_{}{}^{t}(\dot{\overline{z}}+k𝐌\overline{z})(\dot{z}+k𝐌z)d\tau )`$ $``$ $`F(z,z_0,t).G(z,z_0,t)`$ with $`F(z,z_0,t)`$ $`=`$ $`det\left(e^{tk𝐌}\right)e^{\frac{1}{2}\left({}_{}{}^{t}\overline{z}k𝐌z^t\overline{z}_0k𝐌z_0\right)}`$ $`G(z,z_0,t)`$ $`=`$ $`{\displaystyle _{z(0)=z_0}^{z(t)=z}}𝒟z𝒟\overline{z}\mathrm{exp}\left({\displaystyle \frac{1}{2}}{\displaystyle _0^t}\left({}_{}{}^{t}\dot{\overline{z}}\dot{z}+k^2{}_{}{}^{t}\overline{z}𝐌^2z\right)𝑑\tau \right)=z\|e^{tH_0}\|z_0`$ (11) $`=`$ $`det\left({\displaystyle \frac{𝐒}{2\pi }}\right)\mathrm{exp}\left({\displaystyle \frac{1}{2}}\left({}_{}{}^{t}\overline{z}𝐂z+^t\overline{z}_0𝐂z_0^t\overline{z}𝐒z_0^t\overline{z}_0𝐒z\right)\right)`$ (12) $`H_0`$ $`=`$ $`2^t_{\overline{z}}_z+{\displaystyle \frac{1}{2}}k^2{}_{}{}^{t}\overline{z}𝐌^2z`$ (13) The matrices $`𝐒`$ and $`𝐂`$ appearing in (12) are defined as: $$𝐒=k𝐌\left(\mathrm{sinh}(tk𝐌)\right)^1,𝐂=k𝐌\mathrm{coth}(tk𝐌)$$ (14) In fact, $`𝒫`$, eq. ( Statistical properties of the 2D attached Rouse chain), can be easily deduced from the gaussian distribution of $`\eta `$ (use (5); det($`e^{tk𝐌}`$) is simply the functional Jacobian for the change of variable $`\eta z`$ ). (12) is a generalization of the harmonic oscillator propagator . It is obtained by using the normal coordinates. Furthermore, as can be easily checked, $`𝒫`$ is properly normalized: $`dzd\overline{z}𝒫(z,z_0,t)=1`$. Remark that an effective measure can be built for a distinguished monomer of the chain : this can be done by integrating the Wiener measure ( Statistical properties of the 2D attached Rouse chain) over all the paths of the other monomers. The result is a complicated expression that contains, in particular, a non local part (in time) exhibiting the non-Markovian character of the process for this monomer. Nevertheless, we will show, in the sequel, that, despite this complication, we can compute some joint laws for several monomers (and a fortiori for one monomer). So, let us turn to the computation of the area distribution $`P(\{A_i\})`$ for closed trajectories. Inserting the constraint $$\underset{j=1}{\overset{n}{}}\delta \left(A_j\frac{1}{4i}_0^t(z_j\dot{\overline{z}_j}\overline{z}_j\dot{z_j})𝑑\tau \right)$$ (15) in the Wiener measure and using $`\delta (x)=\frac{1}{2\pi }e^{iBx}dB`$, we get the lagrangian for $`n`$ particles submitted to uniform magnetic fields orthogonal to the motion plane (in addition to the harmonic interactions). Remark that, in principle, the magnetic fields are not the same for all the particles. Introducing the $`(n\times n)`$ diagonal matrix $`𝐁`$ ($`𝐁_{ij}=B_i\delta _{ij}`$), we obtain $`P(\{A_i\})`$ $`=`$ $`{\displaystyle \left(\underset{j=1}{\overset{n}{}}\frac{dB_j}{2\pi }e^{iB_jA_j}\right)\left(\frac{Z_𝐁(t)}{Z_0(t)}\right)}`$ (16) $`Z_𝐁(t)`$ $`=`$ $`\text{Tr}e^{tH_𝐁}`$ $`H_𝐁`$ $`=`$ $`H_0+V`$ $`V`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(^tz𝐁_z+^t\overline{z}𝐁_{\overline{z}}\right)+{\displaystyle \frac{1}{8}}^t\overline{z}𝐁^2z`$ (17) In general, the matrices $`𝐁`$ and $`𝐌`$ do not commute and it is a difficult task to get the partition function $`Z_𝐁(t)`$. On the other hand, the distribution of the total area $`A=_{i=1}^nA_i`$ is obtained by taking $`B_j=B`$ for all $`j`$. In this case, $`𝐁`$ and $`𝐌`$ commute. Using normal coordinates and known results about the partition function of the “2D harmonic oscillator + uniform magnetic field” problem , we get the characteristic function of $`A`$ ($`𝐈_𝐧`$ is the $`(n\times n)`$ unit matrix): $`{\displaystyle \frac{Z_𝐁(t)}{Z_0(t)}}`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{n}{}}}\left({\displaystyle \frac{\mathrm{cosh}(tk\omega _j)1}{\mathrm{cosh}(t\sqrt{(k\omega _j)^2+(\frac{B}{2})^2})\mathrm{cosh}(t\frac{B}{2})}}\right)`$ (18) $`=`$ $`{\displaystyle \frac{det\left(\mathrm{cosh}(tk𝐌)𝐈_𝐧\right)}{det\left(\mathrm{cosh}(t\sqrt{(k𝐌)^2+(\frac{𝐁}{2})^2})\mathrm{cosh}(t\frac{𝐁}{2})\right)}}`$ (19) (the $`\omega _j`$’s are defined in (6)). By Fourier transformation, the above $`j^{th}`$ factor gives: $`P_{\omega _j}(A)`$ $`=`$ $`4{\displaystyle \frac{\alpha _j}{\pi }}\mathrm{sinh}^2\left({\displaystyle \frac{\alpha _jt}{2}}\right){\displaystyle \underset{m,r=0}{\overset{\mathrm{}}{}}}\left(\sqrt{{\displaystyle \frac{\gamma _{jr}}{\beta _{jm}}}}K_1(2\sqrt{\beta _{jm}\gamma _{jr}})+\sqrt{{\displaystyle \frac{\beta _{jm}}{\gamma _{jr}}}}K_1(2\sqrt{\beta _{jm}\gamma _{jr}})\right)`$ (20) $`\gamma _{jr}`$ $`=`$ $`\alpha _jt(r+1/2)+i\alpha _jA`$ $`\beta _{jm}`$ $`=`$ $`\alpha _jt(m+1/2)i\alpha _jA`$ $`\alpha _j`$ $`=`$ $`k\omega _j`$ where the $`K_{\pm 1}`$ are modified Bessel functions. Thus, in the general case, $`P(A)`$ can be obtained by a $`n`$-convolution product of the $`P_{\omega _j}`$’s. However, we are afraid that the final result could be awkward! Nevertheless, if we consider the limit $`t\mathrm{}`$, (18) leads to: $$\frac{Z_𝐁(t)}{Z_0(t)}\mathrm{exp}\left(\frac{tB^2}{8k}\underset{i=1}{\overset{n}{}}\frac{1}{\omega _i}\right)=\mathrm{exp}\left(\frac{tB^2n(n+1)}{16k}\right)$$ (21) Then, Fourier transformation shows that, in the large $`t`$ limit, $`A`$ is gaussian and scales like $`\sqrt{t}`$. Such a scaling is expected for all the areas $`A_i`$. This is what we will demonstrate by perturbation theory. When $`t\mathrm{}`$, we have $$Z_𝐁(t)e^{tE_0(𝐁)}$$ (22) where $`E_0(𝐁)`$ is the ground state energy. Moreover, due to the large oscillations of the factor $`e^{iB_jA_j}`$ in (16) when $`A_j\mathrm{}`$, only small values of $`B_j`$ will contribute. So, it is enough to compute $`E_0(𝐁)`$ at lowest order in $`𝐁`$. We will use the normal coordinates $`Z_i`$. The eigenstates of $`H_0`$ are given by $$\mathrm{\Psi }_{\{m_j\},\{n_j\}}(\{Z_j\})=\underset{j=1}{\overset{n}{}}\left(\sqrt{\frac{\omega _jn_j!}{\pi (n_j+\|m_j\|)!}}e^{im_j\theta _j}(\omega _j\|Z_j\|^2)^{\|m_j\|/2}L_{n_j}^{\|m_j\|}(\omega _j\|Z_j\|^2)e^{\frac{1}{2}\omega _j\|Z_j\|^2}\right)$$ (23) $$E_{\{m_j\},\{n_j\}}=\underset{j=1}{\overset{n}{}}(2n_j+\|m_j\|+1)\omega _j$$ (24) where $`L_{n_j}^{\|m_j\|}`$ is a Laguerre polynomial and the ground state is $`\mathrm{\Psi }_{\{0\},\{0\}}`$. The perturbation $`V`$, (17), writes: $$V=\frac{1}{2}\left(^tZ𝐑^1\mathrm{𝐁𝐑}_Z+^t\overline{Z}𝐑^1\mathrm{𝐁𝐑}_{\overline{Z}}\right)+\frac{1}{8}^t\overline{Z}𝐑^1𝐁^2𝐑Z$$ (25) At first order in $`V`$, we get: $$\mathrm{\Delta }E_0^{(1)}(𝐁)=\mathrm{\Psi }_{\{0\},\{0\}}^{}V\mathrm{\Psi }_{\{0\},\{0\}}=\frac{1}{8k}\text{Tr}(𝐁^2𝐌^1)=\frac{1}{8k}\underset{m=1}{\overset{n}{}}mB_m^2$$ (26) Quadratic terms in $`𝐁`$ will also be produced at second order in $`V`$. The non-vanishing contributions will come out from the transitions from the ground-state to the states $`\{m_j=\pm 1,m_l=1,m_i=0\text{if}ij,l\},\{n_k=0\}`$. The computation gives: $`\mathrm{\Delta }E_0^{(2)}(𝐁)`$ $`=`$ $`{\displaystyle \frac{1}{16k}}{\displaystyle \underset{m,m^{}=1}{\overset{n}{}}}B_mB_m^{}{\displaystyle \underset{jl}{}}R_{ml}R_{mj}R_{m^{}l}R_{m^{}j}\left({\displaystyle \frac{\frac{\omega _l}{\omega _j}+\frac{\omega _j}{\omega _l}}{\omega _l+\omega _j}}\right)`$ (27) $`=`$ $`{\displaystyle \frac{1}{8k}}{\displaystyle \underset{m=1}{\overset{n}{}}}mB_m^2+{\displaystyle \frac{1}{2k}}{\displaystyle \underset{m,m^{}=1}{\overset{n}{}}}B_mD_{m,m^{}}B_m^{}`$ (28) with: $$D_{m,m^{}}=\frac{1}{4}_0^{\mathrm{}}\left[\left(e^{\tau 𝐌}\right)_{m,m^{}}\right]^2𝑑\tau +\frac{1}{8}\underset{l=1}{\overset{n}{}}\frac{R_{ml}^2R_{m^{}l}^2}{\omega _l}$$ (29) So, to lowest order in $`𝐁`$, we get: $$\frac{Z_𝐁(t)}{Z_0(t)}\mathrm{exp}\left(\frac{t}{2k}\underset{m,m^{}=1}{\overset{n}{}}B_mD_{m,m^{}}B_m^{}\right)$$ (30) As can be easily checked, (21) is recovered if we set $`B_m=B,m`$. With (16), we arrive at the probability distribution: $$P(\{A_i\})=\left(\frac{k}{2\pi t}\right)^{n/2}\frac{1}{\sqrt{detD}}\mathrm{exp}\left(\frac{k}{2t}\underset{m,m^{}=1}{\overset{n}{}}A_m(D^1)_{m,m^{}}A_m^{}\right)$$ (31) Thus, we observe that the areas $`A_i`$ are correlated gaussian variables and that they scale like $`t^{1/2}`$ as expected. For the special case $`n=2`$, we have: $$P(A_1,A_2)=\sqrt{\frac{5}{3}}\frac{2k}{\pi t}\mathrm{exp}\left(\frac{2k}{9t}\left(23A_1^214A_1A_2+8A_2^2\right)\right)$$ (32) The width of $`A_2`$ is larger than the one of $`A_1`$: this is related to the fact that the second particle is, in average, farther from 0 than the first one. So, it sweeps larger areas. Now, going to the winding angles $`\{\theta _i\}`$ around 0, we proceed as before and insert the constraint $$\underset{j=1}{\overset{n}{}}\delta \left(\theta _j\frac{1}{2i}_0^t\left(\frac{z_j\dot{\overline{z}_j}\overline{z}_j\dot{z_j}}{z_j\overline{z}_j}\right)𝑑\tau \right)$$ (33) in the Wiener measure. We are now faced to the problem of $`n`$ harmonically bound particles submitted to the magnetic fields of point-like vortices located at the origin. The corresponding hamiltonian is: $`H_\lambda `$ $`=`$ $`H_0+W`$ (34) $`W`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}\lambda _i\left({\displaystyle \frac{1}{z_i}}_{\overline{z}_i}{\displaystyle \frac{1}{\overline{z}_i}}_{z_i}\right)+{\displaystyle \underset{i=1}{\overset{n}{}}}{\displaystyle \frac{\lambda _i^2}{2z_i\overline{z}_i}}`$ (35) and the distribution $`P(\{\theta _i\})`$ is given by: $$P(\{\theta _i\})=\left(\underset{j=1}{\overset{n}{}}\frac{d\lambda _j}{2\pi }e^{i\lambda _j\theta _j}\right)\left(\frac{Z_\lambda (t)}{Z_0(t)}\right)$$ (36) Studying the limit $`t\mathrm{}`$, we cannot develop directly as before a perturbation theory with $`W`$: this is because of the last term in $`W`$ that leads to a singular perturbation . Due to this term, all the eigenfunctions of $`H_\lambda `$ must vanish in 0 at least as $`_{i=1}^n(z_i\overline{z}_i)^{\|\lambda _i\|/2}(U)`$ . So we redefine those eigenfunctions : $$\mathrm{\Psi }=U\stackrel{~}{\mathrm{\Psi }}$$ (37) The new hamiltonian acting on $`\stackrel{~}{\mathrm{\Psi }}`$ is: $`\stackrel{~}{H}`$ $`=`$ $`H_0+\stackrel{~}{W}`$ (38) $`\stackrel{~}{W}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}\left((\lambda _i\|\lambda _i\|){\displaystyle \frac{1}{z_i}}_{\overline{z}_i}(\lambda _i+\|\lambda _i\|){\displaystyle \frac{1}{\overline{z}_i}}_{z_i}\right)`$ (39) That time, we can compute $`\mathrm{\Delta }E_0(\lambda )`$ perturbatively and it will appear that only first order is necessary. Integrals of the form $$e^{\frac{1}{2}^t\overline{z}k𝐌z}\frac{1}{\overline{z}_i}_{z_i}e^{\frac{1}{2}^t\overline{z}k𝐌z}dzd\overline{z}$$ (40) are involved. Integrating by parts and using $`_{z_i}\left(\frac{1}{\overline{z}_i}\right)=\pi \delta \left(z_i\right)`$, we get, after some algebra: $$\mathrm{\Delta }E_0(\lambda )=k\underset{j=1}{\overset{n}{}}\frac{\|\lambda _j\|}{(𝐌^1)_{jj}}=k\underset{j=1}{\overset{n}{}}\frac{\|\lambda _j\|}{j}\underset{j=1}{\overset{n}{}}\mu _j\|\lambda _j\|$$ (41) So, for the winding angle distribution, we obtain: $`P(\{\theta _i\})`$ $`=`$ $`{\displaystyle \left(\underset{j=1}{\overset{n}{}}\frac{d\lambda _j}{2\pi }e^{i\lambda _j\theta _j}\right)e^{t_{j=1}^n\mu _j\|\lambda _j\|}}`$ (42) $`=`$ $`{\displaystyle \underset{j=1}{\overset{n}{}}}\left({\displaystyle \frac{1}{\pi \mu _jt}}{\displaystyle \frac{1}{1+(\frac{\theta _j}{\mu _jt})^2}}\right)`$ (43) At large times, the winding angles are uncorrelated, they scale like $`t`$ and are distributed according to Cauchy laws. The variance $`\theta _j^2`$ is infinite: this is, of course, due to the “small windings” occuring in the vicinity of the origin as will be seen explicitly at the end of this paper. Moreover, we observe that $`\theta _j`$ scales like $`\mu _j`$, i.e. like $`1/j`$. This is reasonable because, when $`j`$ increases, the considered particle is, in average, farther from 0 and, consequently, its winding angle must decrease. What is somewhat unexpected is such a simple dependence of $`\theta _j`$ on $`j`$. We also addressed the problem of winding angles around $`n`$ different points of complex coordinates $`b_l`$, $`l=1,\mathrm{}n`$. $`\theta _j^{}`$ being the angle wound by the particle $`j`$ around the point $`b_j`$, we obtained for the set of variables $`\{\theta _j^{}\}`$ the same joint law as (43) except for the change of $`\mu _j`$ into $`\mu _j^{}`$: $$\mu _j^{}=\mu _je^{\mu _j\|b_j\|^2}$$ (44) Owing to the rotationnal symmetry breaking when $`b_j0`$, the winding angles $`\theta _j^{}`$ are statistically reduced by the factor $`e^{\mu _j\|b_j\|^2}`$. Nevertheless, even for large $`\|b_j\|`$’s, the variance $`(\theta _j^{})^2`$ is infinite. Setting all the $`b_j`$’s to zero, we recover (43). This is what we will consider now and assume that we count the winding angles $`\theta _j`$ only when $`\|z_j\|>r_0`$ (i.e. the so-called “big windings” ). Still when $`t\mathrm{}`$, the perturbation $`W`$, eq.(35), can now be used because $`\lambda _j=0`$ when $`\|z_j\|<r_0`$. At first order in $`W`$, the linear contributions in $`\lambda _j`$ will cancel. In the limit of a small, but finite $`r_0`$, we get, for the remaining contribution: $$\mathrm{\Delta }E_0^{(1)}(\lambda )k\|\mathrm{ln}r_0\|\underset{j=1}{\overset{n}{}}\frac{\lambda _j^2}{j}$$ (45) The quadratic contributions in the $`\lambda _j`$’s coming out from the second order in $`W`$ will be finite (thus subleading) when $`r_00^+`$. Finally, we get for the big winding angles asymptotic distribution: $$P(\{\theta _j\})=\underset{j=1}{\overset{n}{}}\sqrt{\frac{j}{4\pi tk\|\mathrm{ln}r_0\|}}\mathrm{exp}\left(\frac{j}{4tk\|\mathrm{ln}r_0\|}\theta _j^2\right)$$ (46) In this limit,the variables $`\theta _j`$ are uncorrelated: the correlations get smaller and smaller when $`r_0`$ decreases. They are now gaussian and scale like $`\sqrt{t\|\mathrm{ln}r_0\|/j}`$. Their variance grows to infinity when $`r_0`$ goes to 0, showing the increasing contribution of the small windings around 0. To summarize, we have computed explicitly the asymptotic joint laws for the areas (that scale like $`\sqrt{t}`$) and for the winding angles (that scale like $`t`$ when no critical region is excluded). The scaling variables we have got compare well with those involved in the Brownian motion on finite domains: this is not so surprising since the chain is bound to a fixed point. Moreover, we have shown that physical interactions between particles (harmonic interactions here) can lead to statistical correlations (case of the areas) or not (case of the winding angles): it depends on the quantity we consider. In a forthcoming paper , we will study the statistical properties of the free Rouse chain. We will especially show that the areas and winding angles distributions are very different from those presented in this work. This is essentially due to the translation invariance that holds when the chain is free. One of us (O. B.) acknowledges Dr. G. Oshanin for drawing his attention to this problem. e-mail: benichou@lptl.jussieu.fr desbois@ipno.in2p3.fr
warning/0005/hep-ex0005002.html	ar5iv	text	# 1 Introduction ## 1 Introduction We describe measurements and searches using a data sample of photonic events with large missing energy collected in 1998 with the OPAL detector at LEP. The events result from $`\mathrm{e}^+\mathrm{e}^{}`$ collisions at a centre-of-mass energy of 188.6 GeV with an integrated luminosity of 177.3 pb<sup>-1</sup>. The present paper builds on publications from earlier data samples at lower centre-of-mass energies . This data-set at 189 GeV gives discovery potential in a new kinematic regime with about a four-fold increase in integrated luminosity. Measurements of photonic event production have also been made by the other LEP collaborations at centre-of-mass energies above the W pair threshold , including new results from L3 and DELPHI at $`\sqrt{s}=`$ 189 GeV . The single-photon and acoplanar-photons search topologies presented here are designed to select events with one or more photons and significant missing transverse energy, indicating the presence of at least one neutrino-like invisible particle which interacts only weakly with matter. The event selections for these search topologies are similar to those used in our recent publication. The single-photon search topology is sensitive to events in which there are one or two photons and missing energy which, within the Standard Model, are expected from the $`\mathrm{e}^+\mathrm{e}^{}\nu \overline{\nu }\gamma (\gamma )`$ process<sup>1</sup><sup>1</sup>1The photon in parentheses denotes that the presence of this photon is allowed but not required.. The acoplanar-photons search topology is designed to select events with two or more photons and significant missing transverse energy which, within the Standard Model, are expected from the $`\mathrm{e}^+\mathrm{e}^{}\nu \overline{\nu }\gamma \gamma (\gamma )`$ process. The single photon topology provides a direct measurement of the invisible width of the $`\mathrm{Z}^0`$ and can probe charged and neutral triple gauge couplings. The acoplanar-photons topology can probe WW$`\gamma \gamma `$ quartic couplings in the $`\mathrm{e}^+\mathrm{e}^{}\nu _\mathrm{e}\overline{\nu _\mathrm{e}}\gamma \gamma `$ process. The neutral and quartic gauge coupling measurements will be described in forthcoming papers based on the event selections described herein. These photonic final-state topologies are sensitive to several different new physics scenarios. A generic classification is $`\mathrm{e}^+\mathrm{e}^{}\mathrm{XY}`$ or $`\mathrm{e}^+\mathrm{e}^{}\mathrm{XX}`$ where $`\mathrm{X}`$ is neutral and can decay radiatively ($`\mathrm{X}\mathrm{Y}\gamma `$) and $`\mathrm{Y}`$ is stable and only weakly interacting. For the general case of massive $`\mathrm{X}`$ and $`\mathrm{Y}`$ this includes conventional supersymmetric processes $`(\mathrm{X}=\stackrel{~}{\chi }_{}^{0}{}_{2}{}^{},\mathrm{Y}=\stackrel{~}{\chi }_{}^{0}{}_{1}{}^{})`$. These topologies also have particularly good sensitivity for the special case of $`M_\mathrm{Y}0`$. This applies both to the production of excited neutrinos $`(\mathrm{X}=\nu ^{},\mathrm{Y}=\nu )`$ and to supersymmetric models in which the lightest supersymmetric particle (LSP) is a light gravitino and $`\stackrel{~}{\chi }_{}^{0}{}_{1}{}^{}`$ is the next-to-lightest supersymmetric particle (NLSP) which decays to a gravitino and a photon ($`\mathrm{X}=\stackrel{~}{\chi }_{}^{0}{}_{1}{}^{},\mathrm{Y}=\stackrel{~}{\mathrm{G}}`$). The neutralino life-time in such models is a free parameter and so we also address the possibility of neutralino-pair production with macroscopic decay lengths. One type of new physics which could be seen in the single-photon topology is the production of an invisible particle in association with a photon. An example of this is graviton-photon production, $`\mathrm{e}^+\mathrm{e}^{}\mathrm{G}\gamma `$ . This can occur within string theory models which allow gravitons to propagate in a higher-dimensional space but restrict Standard Model particles to the usual four space-time dimensions . Another type of new physics is the production of invisible particles tagged by initial-state radiation. One example<sup>2</sup><sup>2</sup>2The initial-state radiation diagram is one of many that contribute to this final state. is production of a pair of gravitinos, $`\mathrm{e}^+\mathrm{e}^{}\stackrel{~}{\mathrm{G}}\stackrel{~}{\mathrm{G}}\gamma `$, as in the superlight gravitino model . The acoplanar-photons search topology also has sensitivity to the production of two particles, one invisible, or with an invisible decay mode, and the other decaying into two photons. This paper will first describe the OPAL detector and the Monte Carlo samples used. A brief summary of the event selections will then be given, followed by cross-section measurements for $`\mathrm{e}^+\mathrm{e}^{}\nu \overline{\nu }\gamma (\gamma )`$ and $`\mathrm{e}^+\mathrm{e}^{}\nu \overline{\nu }\gamma \gamma (\gamma )`$ and comparisons with Standard Model expectations. The new physics search results will then be discussed. ## 2 Detector and Monte Carlo Samples The OPAL detector, which is described in detail in , contains a silicon micro-vertex detector surrounded by a pressurized central tracking system operating inside a solenoid with a magnetic field of 0.435 T. The region outside the solenoid (barrel) and the pressure bell (endcap) is instrumented with scintillation counters, presamplers and the lead-glass electromagnetic calorimeter (ECAL). The magnet return yoke is instrumented for hadron calorimetry and is surrounded by external muon chambers. Electromagnetic calorimeters close to the beam axis measure luminosity and complete the acceptance. The measurements presented here are mainly based on the observation of clusters of energy deposited in the lead-glass electromagnetic calorimeter. This consists of an array of 9,440 lead-glass blocks in the barrel ($`\|\mathrm{cos}\theta \|<0.82`$) with a quasi-pointing geometry and two dome-shaped endcap arrays, each of 1,132 lead-glass blocks, covering the polar angle<sup>3</sup><sup>3</sup>3In the OPAL coordinate system, $`\theta `$ is the polar angle defined with respect to the electron beam direction and $`\varphi `$ is the azimuthal angle. range ($`0.81<\|\mathrm{cos}\theta \|<0.984`$). Fully hermetic electromagnetic calorimeter coverage is achieved beyond the end of the ECAL down to small polar angles with the use of the the gamma-catcher calorimeter, the forward calorimeter (FD) and the silicon-tungsten calorimeter (SW). Scintillators in the barrel and endcap regions are used to reject backgrounds from cosmic ray interactions by providing time measurements for the large fraction ($``$ 80%) of photons which convert in the material in front of the ECAL. The barrel time-of-flight (TOF) scintillator bars are located outside the solenoid in front of the barrel ECAL and match its geometrical acceptance ($`\|\mathrm{cos}\theta \|<0.82`$). Tile endcap (TE) scintillator arrays are located at $`0.81<\|\mathrm{cos}\theta \|<0.955`$ behind the pressure bell and in front of the endcap ECAL. The integrated luminosities of the data samples are determined to better than 1% from small-angle Bhabha scattering events in the SW calorimeter. Triggers based on electromagnetic energy deposits in either the barrel or endcap electromagnetic calorimeters lead to full trigger efficiency for photonic events passing the event selection criteria described in the following section. The KORALZ and NUNUGPV98 Monte Carlo generators were used to simulate the expected Standard Model signal process, $`\mathrm{e}^+\mathrm{e}^{}\nu \overline{\nu }`$ \+ photon(s). For other expected Standard Model processes, a number of different generators were used: RADCOR for $`\mathrm{e}^+\mathrm{e}^{}\gamma \gamma (\gamma )`$; BHWIDE for $`\mathrm{e}^+\mathrm{e}^{}\mathrm{e}^+\mathrm{e}^{}`$; TEEGG for $`\mathrm{e}^+\mathrm{e}^{}\mathrm{e}^+\mathrm{e}^{}\gamma `$; KORALW using grc4f matrix elements for $`\mathrm{e}^+\mathrm{e}^{}\mathrm{}^+\mathrm{}^{}\nu \overline{\nu }(\gamma )`$ and $`\mathrm{e}^+\mathrm{e}^{}\nu \overline{\nu }\mathrm{q}\overline{\mathrm{q}}`$, and KORALZ for $`\mathrm{e}^+\mathrm{e}^{}\mu ^+\mu ^{}(\gamma )`$ and $`\mathrm{e}^+\mathrm{e}^{}\tau ^+\tau ^{}(\gamma )`$. The Vermaseren program and grc4f were used for $`\mathrm{e}^+\mathrm{e}^{}\mathrm{e}^+\mathrm{e}^{}\mathrm{}^+\mathrm{}^{}`$. The expected contributions from each of these Standard Model processes were evaluated using a total equivalent integrated luminosity at least five times larger than the integrated luminosity of the data sample. To simulate possible new physics processes of the type $`\mathrm{e}^+\mathrm{e}^{}\mathrm{XY}`$ and $`\mathrm{e}^+\mathrm{e}^{}\mathrm{XX}`$ where $`\mathrm{X}`$ decays to $`\mathrm{Y}\gamma `$ and $`\mathrm{Y}`$ escapes detection, a modified version of the SUSYGEN Monte Carlo generator was used to produce neutralino pair events of the type $`\mathrm{e}^+\mathrm{e}^{}\stackrel{~}{\chi }_{}^{0}{}_{2}{}^{}\stackrel{~}{\chi }_{}^{0}{}_{1}{}^{}`$ and $`\mathrm{e}^+\mathrm{e}^{}\stackrel{~}{\chi }_{}^{0}{}_{2}{}^{}\stackrel{~}{\chi }_{}^{0}{}_{2}{}^{}`$, $`\stackrel{~}{\chi }_{}^{0}{}_{2}{}^{}\stackrel{~}{\chi }_{}^{0}{}_{1}{}^{}\gamma `$, with isotropic angular distributions for the production and decay of $`\stackrel{~}{\chi }_{}^{0}{}_{2}{}^{}`$ and including initial-state radiation. Monte Carlo events were generated at 48 (for $`\mathrm{XY}`$ production) and 42 (for $`\mathrm{XX}`$ production) points in the kinematically accessible region of the ($`M_\mathrm{X}`$, $`M_\mathrm{Y}`$) plane. A Monte Carlo generator was written to simulate the superlight gravitino signature $`\stackrel{~}{\mathrm{G}}\stackrel{~}{\mathrm{G}}\gamma `$, discussed in Section 5.1.4. The KORALZ $`\nu \overline{\nu }\gamma (\gamma )`$ sample was used to determine the efficiency for graviton-photon production in the context of additional space dimensions, by means of an event reweighting (see section 5.1.5). The same procedure was also used to calculate the efficiency for $`\stackrel{~}{\mathrm{G}}\stackrel{~}{\mathrm{G}}\gamma `$ production, and compared with the prediction of the direct Monte Carlo simulation. All the Monte Carlo samples described above were processed through the OPAL detector simulation . Simulation of $`\mathrm{e}^+\mathrm{e}^{}\mathrm{XX}`$, $`\mathrm{X}\mathrm{Y}\gamma `$ signal events with $`M_\mathrm{Y}0`$ where $`M_\mathrm{X}`$ has a finite non-zero lifetime $`\tau _X`$ was implemented in the full simulation of the OPAL detector. In particular, the massive quasi-stable neutral particle $`\mathrm{X}`$ was assigned properties similar to a heavy neutrino for the purpose of propagation, and then propagated within the GEANT framework, according to its initial kinematics and $`\tau _X`$ before forcing the $`\mathrm{X}\mathrm{Y}\gamma `$ decay. ## 3 Photonic Event Selection This section summarizes the criteria for selecting single-photon and acoplanar-photons events. The kinematic acceptance of each selection is defined in terms of the photon energy, $`E_\gamma `$, and the photon polar angle, $`\theta `$. In addition, the scaled energy, $`x_\gamma `$, is defined as $`E_\gamma /E_{\mathrm{beam}}`$, and the scaled transverse energy, $`x_T`$, as $`x_\gamma \mathrm{sin}\theta `$. One or two photons accompanied by invisible particle(s): * At least one photon with $`x_T>0.05`$ and with $`15^{}<\theta <165^{}`$ ($`\|\mathrm{cos}\theta \|<0.966)`$. Two or more photons accompanied by invisible particle(s): * At least two photons, each with $`x_\gamma >0.05`$ and $`15^{}<\theta <165^{}`$, or one photon with $`E_\gamma >1.75`$ GeV and $`\|\mathrm{cos}\theta \|<0.8`$ and a second photon with $`E_\gamma >1.75`$ GeV and $`15^{}<\theta <165^{}`$. * The transverse momentum $`p_T^{\gamma \gamma }`$ of the two-photon system consisting of the two highest energy photons must satisfy $`p_T^{\gamma \gamma }/E_{\mathrm{beam}}>0.05`$. In each of the two cases we retain acceptance for events with additional photons in which the resulting photonic system is still consistent with the presence of significant missing energy. This reduces the sensitivity of each measurement to the modelling of higher-order contributions. Consequently, a large fraction of the kinematic acceptance of the acoplanar-photons selection is also contained in the kinematic acceptance of the single-photon selection. ### 3.1 Single-Photon Event Selection The single-photon selection criteria are explained in detail in a previous publication. After defining the kinematic acceptance, additional cuts on cluster quality, forward energy, muon chamber and hadron calorimeter information, and multiphoton kinematics are used to remove cosmic ray backgrounds, beam related backgrounds and standard model physics backgrounds that have no missing energy. Then events are classed as either having or not having photon conversion candidates based on tracking chamber information. Each of these two classes has somewhat different additional selection criteria based on charged track activity, time-of-flight information, and other background suppression cuts. Several improvements have been made to the single-photon selection used in this analysis compared with the previous analysis. They are described in detail in Appendix A. ### 3.2 Acoplanar-Photons Event Selection The acoplanar-photons selection has classes of selection requirements similar to those of the single-photon selection. The details of the acoplanar-photons selection are described in our previous analysis; changes with respect to the previous analysis are minimal and are given in Appendix A. The requirement of a second photon in an event reduces many of the backgrounds which are otherwise a problem for the single-photon analysis. For this reason, the acoplanar photons selection can tolerate a lower energy threshold for the most energetic photon, as well as looser but more inclusive acceptance requirements for photon conversions. On the other hand, the single-photon selection has more acceptance for events having no time-of-flight information for the photons. In order to obtain the best overall acceptance for acoplanar-photons, we have added to the acoplanar-photons selection that part of the single-photon selection which contains two photons within the kinematic acceptance of the acoplanar-photons selection. This addition results in a relative increase in efficiency of 9.6% for Standard Model $`\mathrm{e}^+\mathrm{e}^{}\nu \overline{\nu }\gamma \gamma (\gamma )`$ events. ## 4 Selection Results The results of the single-photon and acoplanar-photons selections and the corresponding cross-section measurements and other measured event quantities are given below in sections 4.1 and 4.2. Results from both selections are summarized in Table 1. ### 4.1 Single-Photon After applying the single-photon selection criteria to the data sample, 643 events are selected. The expected contribution from cosmic ray and beam-related backgrounds is 4.6 $`\mathbf{\pm }`$ 1.5 events. These backgrounds have been estimated from events having out-of-time TOF or TE information, but passing all other selection criteria, and from events selected with looser criteria that have been visually scanned. Of the expected physics backgrounds from plausible sources, $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝐞^\mathbf{+}𝐞^{\mathbf{}}𝜸`$, $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathbf{}^\mathbf{+}\mathbf{}^{\mathbf{}}𝝂\overline{𝝂}\mathbf{(}𝜸\mathbf{)}`$, $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}^\mathbf{+}\mathbf{}^{\mathbf{}}`$, $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂\overline{𝝂}𝐪\overline{𝐪}`$, $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝁^\mathbf{+}𝝁^{\mathbf{}}𝜸`$ and $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝉^\mathbf{+}𝝉^{\mathbf{}}𝜸`$ have non-negligible contributions<sup>4</sup><sup>4</sup>4 The expected number of events from the Standard Model process $`\mathrm{e}^+\mathrm{e}^{}\nu \overline{\nu }\nu \overline{\nu }\gamma (\gamma )`$ is negligible and is neglected in both the efficiency and background estimates. The total number of expected physics background events is 4.6 $`\mathbf{\pm }`$ 0.5 and the contributing sources are summarized in Table 2. The number of events expected from the Standard Model process $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂\overline{𝝂}𝜸\mathbf{(}𝜸\mathbf{)}`$ as predicted by KORALZ is 679 $`\mathbf{\pm }`$ 5 (stat) $`\mathbf{\pm }`$ 14 (sys) where the systematic error is from experimental sources as discussed later. The number of events observed is consistent with the number expected from $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂\overline{𝝂}𝜸\mathbf{(}𝜸\mathbf{)}`$ plus the background. The efficiency for selecting $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂\overline{𝝂}𝜸\mathbf{(}𝜸\mathbf{)}`$ events within the kinematic acceptance of the single-photon selection is $`\mathbf{(}\mathbf{82.1}\mathbf{\pm }\mathbf{1.7}\mathbf{)}`$%. For both the single-photon and acoplanar-photons selections, efficiency losses due to vetoes on random detector occupancy range from about (3-5)%. Quoted efficiencies include these losses. The single-photon cross-section for $`\sqrt{𝒔}\mathbf{=}\mathrm{𝟏𝟖𝟗}`$ GeV, accounting for detector and selection efficiencies and subtracting the estimated background, is 4.35 $`\mathbf{\pm }`$ 0.17 (stat) $`\mathbf{\pm }`$ 0.09 (sys) pb. The systematic error on the cross-section measurement from knowledge of the efficiency and normalisation is estimated to be 2.1%. The contributing uncertainties are summarized in Table 3. The dominant systematics are the selection efficiency uncertainty (1.5%) and the uncertainty on the detector occupancy estimate (1%). The event selection efficiency is controlled using a data sample of around 1200 $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝜸𝜸`$ events. In particular, the efficiency and time response of the TB and TE scintillators are measured from these data and small correction factors applied to the efficiency. The selection efficiency systematic includes the uncertainties on such corrections and an estimate of the residual uncertainty on the efficiency. The occupancy, estimated from random beam-crossing triggers, is typically 4% and we assign an error of 1%. We have compared the estimated efficiency using the NUNUGPV98 and KORALZ event generators. We find that the relative efficiency difference is ($`\mathbf{0.5}\mathbf{\pm }\mathbf{0.4}`$)% and assign an error of 0.5% to take into account the sensitivity to the modelling of the photon energy, angle and multiplicity distributions. This source of error is considerably reduced from our previous publications, largely because both generators that are now used are designed to model accurately $`𝝂\overline{𝝂}𝜸𝜸`$ events. Additional systematic errors on the efficiency arise from uncertainties in modelling the material close to or inside the beam-pipe which accounts for photons which convert early (20% relative uncertainty) and from systematic effects in the track reconstruction and parameter estimation which affect the association of tracks to ECAL clusters. Uncertainties in the photon energy scale, resolution and angular measurement give small additional systematic contributions. The cross-section as a function of centre-of-mass energy is plotted in Figure 1. Cross-section results from the current analysis and from our earlier publications are plotted. The curve shows the predicted cross-section from the KORALZ event generator for the Standard Model process $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂\overline{𝝂}𝜸\mathbf{(}𝜸\mathbf{)}`$. The data are consistent with the prediction. Figure 2a shows the recoil mass distribution, where the recoil mass $`𝑴_{\mathrm{𝐫𝐞𝐜𝐨𝐢𝐥}}`$ is defined as the mass recoiling against the photon (or against the two-photon system). The peak in the distribution at $`𝑴_𝐙`$ is due to a large contribution from the decay $`𝐙^\mathrm{𝟎}\mathbf{}𝝂\overline{𝝂}`$. Figure 2b shows the polar angle distribution along with the $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂\overline{𝝂}𝜸\mathbf{(}𝜸\mathbf{)}`$ Monte Carlo expectation. In both distributions, there is consistency between data and Monte Carlo, although we note a deficit in the radiative return peak and a slight excess in the low energy, high recoil mass region. The single-photon selection is designed to allow for the presence of a second photon in order to accept events from the $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂\overline{𝝂}𝜸𝜸`$ process. Thirty-six observed events are considered to be two-photon events (i.e. have a second photon with deposited energy exceeding 300 MeV and with $`\mathrm{𝟏𝟓}^{\mathbf{}}\mathbf{<}𝜽\mathbf{<}\mathrm{𝟏𝟔𝟓}^{\mathbf{}}`$), consistent with the expectation of 33.6 $`\mathbf{\pm }`$ 1.5 events from the KORALZ Monte Carlo. Of these, 20 fall within the kinematic acceptance of the acoplanar photon selection, as compared to the KORALZ expectation of $`\mathbf{23.5}\mathbf{\pm }\mathbf{1.0}`$. ### 4.2 Acoplanar-Photons The acoplanar-photons selection applied to the data sample yields 24 events, in good agreement with the KORALZ prediction of $`\mathbf{26.9}\mathbf{\pm }\mathbf{1.2}`$ events for the Standard Model $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂\overline{𝝂}𝜸𝜸\mathbf{(}𝜸\mathbf{)}`$ contribution. The expected contribution from other Standard Model processes and from cosmic ray and beam related backgrounds is $`\mathbf{0.11}\mathbf{\pm }\mathbf{0.04}`$ events. The selection efficiency for $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂\overline{𝝂}𝜸𝜸\mathbf{(}𝜸\mathbf{)}`$ events within the kinematic acceptance of the selection is ($`\mathbf{66.4}\mathbf{\pm }\mathbf{2.9}`$)%. The corresponding cross-section is $`\mathbf{0.204}\mathbf{\pm }\mathbf{0.043}`$ pb, compared to a KORALZ prediction of $`\mathbf{0.228}\mathbf{\pm }\mathbf{0.002}`$ pb. The OPAL measurements of the cross-sections at $`\sqrt{𝒔}\mathbf{=}`$ 130, 136, 161, 172, 183 and 189 GeV are summarized in Table 4. Results for $`\sqrt{𝒔}\mathbf{<}\mathrm{𝟏𝟖𝟗}`$ GeV have been taken from our previous publications. The dominant source of systematic uncertainties is modelling of the reconstruction efficiency, especially the simulation of the detector material and consequent photon conversion probabilities. Other sources arise from uncertainties on the electromagnetic calorimeter energy scale and resolution, on the integrated luminosity measurement, on detector occupancy estimates and from comparisons of different Monte Carlo event generators for the process $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂\overline{𝝂}𝜸𝜸\mathbf{(}𝜸\mathbf{)}`$. The relative systematic error from all sources is 4.3%. The kinematic properties of the selected events are displayed in Figure 3 where they are compared with the predicted distributions for $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂\overline{𝝂}𝜸𝜸\mathbf{(}𝜸\mathbf{)}`$ obtained using the KORALZ generator normalized to the integrated luminosity of the data. Plot (a) shows the recoil mass distribution of the selected acoplanar-photon pairs (or of the two most-energetic photons in the case of events with three or more photons). The distribution is peaked near the mass of the $`𝐙^\mathrm{𝟎}`$ as is expected for contributions from $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂\overline{𝝂}𝜸𝜸\mathbf{(}𝜸\mathbf{)}`$. The resolution of the recoil mass is typically 3-5 GeV for $`𝑴_{\mathrm{𝐫𝐞𝐜𝐨𝐢𝐥}}\mathbf{}𝑴_𝐙`$. Plot (b) shows the distribution of the scaled energy of the second most energetic photon. Plot (c) shows the $`𝜸𝜸`$ invariant-mass distribution for which the mass resolution is typically 0.7-1.9 GeV. Plot (d) shows the distribution in scaled transverse momentum of the selected two-photon system. There is 1 selected event having a third photon with deposited energy above 300 MeV and within the polar-angle acceptance of the selection. The corresponding expectation from KORALZ is $`\mathbf{1.20}\mathbf{\pm }\mathbf{0.08}`$ events. ## 5 Data Interpretation The results of the single-photon and acoplanar-photons selections are used to test the Standard Model and search for new physics contributions. For the XY and XX searches, we test the range of the following product branching ratios which are consistent with the data: $`𝝈\mathbf{(}𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐘}\mathbf{)}\mathbf{}\mathrm{𝐁𝐑}\mathbf{(}𝐗\mathbf{}𝐘𝜸\mathbf{)}`$ and $`𝝈\mathbf{(}𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐗}\mathbf{)}\mathbf{}\mathrm{𝐁𝐑}^\mathrm{𝟐}\mathbf{(}𝐗\mathbf{}𝐘𝜸\mathbf{)}`$. Given lack of evidence for signal, we then set 95% CL upper limits on these quantities. This is done both for the general case of massive $`𝐗`$ and $`𝐘`$, and also separately for the special case of $`𝑴_𝐘\mathbf{}\mathrm{𝟎}`$. All efficiencies are first evaluated under the assumption that the decay length of $`𝐗`$ is zero. For the XX search in the special case of $`𝑴_𝐘\mathbf{}\mathrm{𝟎}`$, we evaluate the efficiency as a function of the non-zero lifetime of the X particle, thus quantifying the sensitivity of the search for the general case of non-prompt decay. For both the XY and XX searches, Monte Carlo samples were generated for a variety of mass points in the kinematically accessible region of the $`\mathbf{(}𝑴_𝐗\mathbf{,}𝑴_𝐘\mathbf{)}`$ plane. To set limits for arbitrary $`𝑴_𝐗`$ and $`𝑴_𝐘`$, the efficiency over the entire $`\mathbf{(}𝑴_𝐗\mathbf{,}𝑴_𝐘\mathbf{)}`$ plane is parameterized using the efficiencies calculated at the generated mass points. As justified previously, we restrict the searches to $`𝑴_𝐗\mathbf{+}𝑴_𝐘\mathbf{>}𝑴_𝐙`$ for the single-photon topology and to $`𝑴_𝐗`$ values larger than about $`𝑴_𝐙`$/2 for the acoplanar-photons topology. ### 5.1 Single-Photon With the single-photon topology we measure the number of light neutrino species and observe a rise in the cross-section at low photon energies which is consistent with the additional cross-section expected from charged current contributions to $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂_𝒆\overline{𝝂_𝒆}𝜸\mathbf{(}𝜸\mathbf{)}`$. We also give results of searches for XY production, superlight gravitino pair production, and graviton-photon production in the context of models with extra space dimensions. #### 5.1.1 Neutrino Counting Single-photon events are expected within the Standard Model from the process $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂\overline{𝝂}𝜸\mathbf{(}𝜸\mathbf{)}`$. At the tree level, the cross-section for muon-type and tau-type neutrinos is attributable solely to $`𝒔`$-channel Z production with initial-state radiation. For electron-type neutrinos, the single photon cross-section arises from Feynman diagrams corresponding to $`𝒔`$-channel Z production with initial-state radiation and $`𝒕`$-channel W exchange with radiation from the initial-state or the exchanged W. Higher order electroweak processes such as WW boxes are expected to give a negligible contribution to the cross-section. The results of the single-photon selection are used to measure the size of the $`𝒔`$-channel Z production contributions and the W-related contributions (W amplitude squared plus the W-Z interference), which we parameterize in terms of $`𝐍_𝝂`$ and $`𝒇_𝑾`$. The measurement of the pure $`𝒔`$-channel Z production contributions is a direct measurement of the Z invisible width, which is related to the effective number of light neutrino generations, $`𝐍_𝝂`$, defined as the ratio of the Z invisible width to the expected width in the Standard Model for one neutrino generation. The W-related contributions are parameterized by a multiplicative scale factor, $`𝒇_𝑾`$, defined to be 1 for the Standard Model expectation. In order to measure $`𝐍_𝝂`$ and/or $`𝒇_𝑾`$, we perform a binned fit to both the overall event rate and to the shape of the photon energy distribution by minimising the negative log-likelihood: $$\mathbf{}\mathrm{𝐥𝐨𝐠}𝑳\mathbf{=}\mathbf{}\mathrm{𝐥𝐨𝐠}𝓟\mathbf{(}𝒏_{𝒆𝒙𝒑}\mathbf{(}𝑬_𝜸\mathbf{,}𝐍_𝝂\mathbf{,}𝒇_𝑾\mathbf{)}\mathbf{}𝒏_{𝒐𝒃𝒔}\mathbf{(}𝑬_𝜸\mathbf{)}\mathbf{)}$$ (1) where $`𝒏_{𝒆𝒙𝒑}`$ is given by $$𝒏_{𝒆𝒙𝒑}\mathbf{(}𝑬_𝜸\mathbf{,}𝐍_𝝂\mathbf{,}𝒇_𝑾\mathbf{)}\mathbf{=}𝒇_𝑾𝒏_{𝒆𝒙𝒑}^𝑾\mathbf{(}𝑬_𝜸\mathbf{)}\mathbf{+}𝐍_𝝂𝒏_{𝒆𝒙𝒑}^𝒁\mathbf{(}𝑬_𝜸\mathbf{)}\mathbf{.}$$ (2) Here $`𝒏_{𝒐𝒃𝒔}`$ is the observed number of events, $`𝒏_{𝒆𝒙𝒑}`$ is the expected number of events as a function of $`𝐍_𝝂`$ and $`𝒇_𝑾`$ derived by reweighting fully simulated $`𝝂\overline{𝝂}𝜸\mathbf{(}𝜸\mathbf{)}`$ events, and $`𝓟`$ represents the Poisson probability for observing $`𝒏_{𝒐𝒃𝒔}`$ events given an expectation of $`𝒏_{𝒆𝒙𝒑}`$. $`𝒏_{𝒆𝒙𝒑}^𝑾`$ is the number of events for Standard Model W contributions<sup>5</sup><sup>5</sup>5Both the pure W $`t`$-channel, which is dominant, and its interference with the $`\mathrm{Z}^0`$ $`s`$-channel are taken into account assuming Standard Model couplings of electrons and electron-type neutrinos. and $`𝒏_{𝒆𝒙𝒑}^𝒁`$ is the number of $`𝒔`$-channel events expected per neutrino generation; both $`𝒏_{𝒆𝒙𝒑}^𝑾`$ and $`𝒏_{𝒆𝒙𝒑}^𝒁`$ were evaluated using the NUNUGPV98 generator<sup>6</sup><sup>6</sup>6 Using the NUNUGPV98 feature which allows $`\nu _\mathrm{e}\overline{\nu }_\mathrm{e}\gamma (\gamma )`$ and $`\nu _\mu \overline{\nu }_\mu \gamma (\gamma )`$ events to be generated separately . The dominant systematic errors arise from the uncertainty on the selection efficiency and from the theoretical uncertainty on the expected number of events, with a minor contribution from the limited Monte Carlo statistics. The relative uncertainty on the selection efficiency is 2.1% and was discussed in Section 4. We assign a value of 2% to the theoretical uncertainty based on comparisons and estimated precisions of the NUNUGPV98, KORALZ and grc$`𝝂𝝂𝜸`$ event generators. In addition, an uncertainty related to the modelling of the photon energy spectrum is assigned, again estimated by comparing the NUNUGPV98 and KORALZ event generators. Other sources of systematic error, such as the uncertainty on the centre-of-mass energy and the expected background, result in negligible contributions. In a first step, we check the consistency with the Standard Model predictions by fitting for both $`𝐍_𝝂`$ and $`𝒇_𝑾`$. The results are $`𝒇_𝑾\mathbf{=}\mathbf{1.12}\mathbf{\pm }\mathbf{0.13}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{)}\mathbf{\pm }\mathbf{0.12}\mathbf{(}\mathrm{𝐬𝐲𝐬}\mathbf{)}`$ and $`𝑵_𝝂\mathbf{=}\mathbf{2.63}\mathbf{\pm }\mathbf{0.15}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{)}\mathbf{\pm }\mathbf{0.11}\mathbf{(}\mathrm{𝐬𝐲𝐬}\mathbf{)}`$ with a correlation coefficient of $`\mathbf{}\mathrm{𝟒𝟏}`$%. Figure 4 shows the 70%, 95% and 99% probability contours in the space of the two parameters, while Figure 5 shows the photon energy distribution for the events selected at 189 GeV, compared to the expectation with $`𝒇_𝑾\mathbf{=}\mathbf{1.12}`$, $`𝐍_𝝂\mathbf{=}\mathbf{2.63}`$. The data are seen to be in fair agreement with the Standard Model prediction for $`𝑵_𝝂\mathbf{=}\mathrm{𝟑}`$ and $`𝒇_𝑾\mathbf{=}\mathrm{𝟏}`$. In particular, the W contributions are observed with a high degree of significance, even when allowing $`𝑵_𝝂`$ to be unconstrained. In a second step, we assume the W contributions to be as predicted by the Standard Model and fit for $`𝑵_𝝂`$. The result is: $$𝑵_𝝂\mathbf{=}\mathbf{2.69}\mathbf{\pm }\mathbf{0.13}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{)}\mathbf{\pm }\mathbf{0.11}\mathbf{(}\mathrm{𝐬𝐲𝐬}\mathbf{)}\mathbf{.}$$ (3) Alternatively, one can assume $`𝑵_𝝂\mathbf{=}\mathrm{𝟑}`$ as indicated by the precise, but less direct, measurements of the Z lineshape , and fit for the relative size of the W-contributions. The result is $`𝒇_𝑾\mathbf{=}\mathbf{0.99}\mathbf{\pm }\mathbf{0.11}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{)}\mathbf{\pm }\mathbf{0.12}\mathbf{(}\mathrm{𝐬𝐲𝐬}\mathbf{)}`$, establishing that the W-contributions are observed and are consistent with the expectations from the Standard Model. #### 5.1.2 Search for $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐘}`$, $`𝐗\mathbf{}𝐘𝜸`$ ; General case: $`𝑴_𝐘\mathbf{}\mathrm{𝟎}`$ We search for evidence of new physics processes of type $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐘}`$$`𝐗\mathbf{}𝐘𝜸`$. To select candidate new physics events, kinematic-consistency and degraded-resolution cuts are applied to events in the single-photon event sample. The kinematic-consistency cuts require that the energy of the most energetic photon be within the range kinematically consistent with mass values $`\mathbf{(}𝑴_𝐗\mathbf{,}𝑴_𝐘\mathbf{)}`$ after accounting for energy resolution effects. The degraded-resolution cuts reject events in which the most energetic photon is in the angular regions $`\mathbf{0.78}\mathbf{<}\mathbf{\|}\mathrm{𝐜𝐨𝐬}𝜽\mathbf{\|}\mathbf{<}\mathbf{0.82}`$ or $`\mathbf{\|}\mathrm{𝐜𝐨𝐬}𝜽\mathbf{\|}\mathbf{>}\mathbf{0.95}`$; energy resolution is significantly degraded in these regions and the Monte Carlo simulation is less reliable. These cuts are approximately 94% efficient in selecting signal events within the kinematic acceptance of the single-photon selection, assuming uniform energy and $`\mathrm{𝐜𝐨𝐬}𝜽`$ signal distributions. Two methods are used to search for, and place upper limits on, contributions from $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐘}`$$`𝐗\mathbf{}𝐘𝜸`$. The first is an event-counting method, based simply on the total number of events selected as candidates. The kinematic cuts, described above, are sufficiently loose that the results are relatively insensitive to the shapes of signal distributions in photon energy or $`\mathrm{𝐜𝐨𝐬}𝜽`$ that result from the specifics of a particular model. Thus this method gives the more generally applicable cross-section limits for XY production. In many models, however, the distributions for the production and decay angles in the process $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐘}`$$`𝐗\mathbf{}𝐘𝜸`$ are approximately isotropic. We therefore also perform a likelihood-based analysis to search for new physics contributions under the assumption of isotropically distributed production and decay angles. This results in a substantial increase in experimental sensitivity, since energy and angular distributions can now be used to differentiate signal from background. In both methods, Standard Model background is assumed to be from $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂\overline{𝝂}𝜸\mathbf{(}𝜸\mathbf{)}`$ only. Background from other sources, including the estimated cosmic ray and beam-related background, is small and is neglected in the limit calculations. Since uncertainties due to energy scale and resolution effects in the low $`𝑴_𝐗\mathbf{}𝑴_𝐘`$ region lead to large relative uncertainties in the estimated efficiencies, we restrict our search to the region in which $`𝑴_𝐗\mathbf{}𝑴_𝐘\mathbf{>}\mathrm{𝟓}`$ GeV. ##### Event-Counting Method: <br> Upper limits are calculated based on the total numbers of observed and expected background events using the method described in reference. This procedure is similar to the method used in our previous publication, with the exception that only the kinematic-consistency and degraded-resolution cuts are applied in order to improve the generality of the results. The total signal efficiency within the single-photon kinematic acceptance is approximately 85%, varying by less than 1% over most of the $`\mathbf{(}𝑴_𝐗\mathbf{,}𝑴_𝐘\mathbf{)}`$ plane. However, the efficiency decreases significantly when the mass difference $`𝑴_𝐗\mathbf{}𝑴_𝐘`$ becomes small, due to low photon energies resulting from the $`𝐗\mathbf{}𝐘𝜸`$ decay. Contributions to the systematic error on the efficiency for selecting events from potential new physics sources are similar in nature and size to those discussed in our previous publication. The total relative systematic error varies from 2.5 to 5.5%, depending on $`𝑴_𝐗`$ and $`𝑴_𝐘`$. These systematic errors have been treated according to the method in reference; the effect on the upper limits is small. Uncertainties in the $`𝝂\overline{𝝂}𝜸\mathbf{(}𝜸\mathbf{)}`$ background estimate have a more significant effect on the upper limits, and are treated by means of a convolution within the limit calculations. A relative 3% uncertainty has been assigned to the expected background contribution, based on factors considered in the cross-section measurement and assigning a 2% theoretical uncertainty. There are 552 events selected after the degraded-resolution cuts are applied; the KORALZ expectation is $`\mathrm{𝟔𝟎𝟏}\mathbf{\pm }\mathrm{𝟏𝟒}`$. The photon energy distribution for these events is shown in Figure 6. The number of selected events after further application of the kinematic-consistency cuts depends on the values for $`𝑴_𝐗`$ and $`𝑴_𝐘`$; ranges for observed and expected events are displayed as contour plots in Figure 7. The consistency between the number of observed events ($`𝒏_{𝒐𝒃𝒔}`$) and the number of expected background events ($`𝒏_𝒃`$) is given in Figure 8. Plotted is $`𝑷_{\mathrm{𝐟𝐥𝐮𝐜𝐭}}`$, the probability for observing at least $`𝒏_{𝒐𝒃𝒔}`$ events given a background expectation of $`𝒏_𝒃\mathbf{\pm }𝝈_𝒃`$ events, defined by: $$𝑷_{\mathrm{𝐟𝐥𝐮𝐜𝐭}}\mathbf{=}\underset{𝒏\mathbf{=}𝒏_{𝒐𝒃𝒔}}{\overset{\mathbf{}}{\mathbf{}}}\mathbf{\left[}\mathbf{}_{\mathbf{}\mathbf{}}^{\mathbf{}}𝒅𝝁\mathbf{\left(}\frac{𝒆^{\mathbf{(}𝒏_𝒃\mathbf{}𝝁\mathbf{)}^\mathrm{𝟐}\mathbf{/}\mathrm{𝟐}𝝈_𝒃^\mathrm{𝟐}}}{\sqrt{\mathrm{𝟐}𝝅𝝈_𝒃^\mathrm{𝟐}}}\mathbf{\right)}\mathbf{}\mathbf{\left(}𝒆^\mathbf{}𝝁\frac{𝝁^𝒏}{𝒏\mathbf{!}}\mathbf{\right)}\mathbf{\right]}\mathbf{.}$$ (4) This figure is intended to highlight regions in the $`\mathbf{(}𝑴_𝐗\mathbf{,}𝑴_𝐘\mathbf{)}`$ plane with an excess of events that could potentially be indicative of new physics. Such an excess would appear as a very low probability $`\mathbf{(}𝑷_{\mathrm{𝐟𝐥𝐮𝐜𝐭}}\mathbf{}\mathrm{𝟏}\mathbf{\%}`$). However, a deficit of events is actually observed for most of the $`\mathbf{(}𝑴_𝐗\mathbf{,}𝑴_𝐘\mathbf{)}`$ plane. This leads to high values for $`𝑷_{\mathrm{𝐟𝐥𝐮𝐜𝐭}}`$; for example, for $`𝒏_{𝒐𝒃𝒔}\mathbf{}\mathrm{𝟏}`$, a value $`𝑷_{\mathrm{𝐟𝐥𝐮𝐜𝐭}}\mathbf{=}\mathbf{0.97}`$ indicates a deficit consistent at approximately the 3% level. The calculated values of $`𝑷_{\mathrm{𝐟𝐥𝐮𝐜𝐭}}`$ range from 0.71% to 99.94%. Values with rather low and high probability occur partly because we consider several thousand points in the $`\mathbf{(}𝑴_𝐗\mathbf{,}𝑴_𝐘\mathbf{)}`$ plane; each of these has a different set of kinematic-consistency cuts which select different parts of the photon energy spectrum. The net effect is that almost every possible energy range is selected for some $`\mathbf{(}𝑴_𝐗\mathbf{,}𝑴_𝐘\mathbf{)}`$ value. For example, as seen in Figure 6, there is a deficit in the energy range from approximately 26 to 38 GeV; this range corresponds roughly to $`𝑴_𝐗\mathbf{=}\mathrm{𝟏𝟏𝟑}`$ GeV and $`𝑴_𝐘\mathbf{=}\mathrm{𝟕𝟒}`$ GeV. The general structure of Figure 8 can also be understood from the photon energy spectrum. There is an overall deficit of approximately 49 events ($`\mathbf{}\mathbf{1.7}𝝈`$). As seen in Figure 6, most of the deficit lies in the region of the radiative return peak, while there is an excess in the low-energy region. For regions in the $`\mathbf{(}𝑴_𝐗\mathbf{,}𝑴_𝐘\mathbf{)}`$ plane for which low-energy photons are not kinematically allowed, events populating the low-energy region of the energy spectrum are rejected and the significance of the remaining deficit is enhanced. In particular, this occurs when $`𝑴_𝐗\mathbf{}𝑴_𝐘`$ is relatively large, and accounts for the large $`𝑷_{\mathrm{𝐟𝐥𝐮𝐜𝐭}}`$ values in the low-$`𝑴_𝐘`$ region of Figure 8. Conversely, regions with a small mass difference $`𝑴_𝐗\mathbf{}𝑴_𝐘`$ correspond to low photon energies, which in turn correspond to the region of the photon energy spectrum in which there is an excess. This accounts for the small $`𝑷_{\mathrm{𝐟𝐥𝐮𝐜𝐭}}`$ values in the low $`𝑴_𝐗\mathbf{}𝑴_𝐘`$ region of Figure 8. Cross-section times branching-ratio upper limits calculated using the event-counting method are shown in Figure 9 and range from 43 to 409 fb. Results for one $`𝑴_𝐗\mathbf{,}𝑴_𝐘`$ pair ($`𝑴_𝐗\mathbf{=}\mathrm{𝟏𝟐𝟓}\mathbf{,}𝑴_𝐘\mathbf{=}\mathrm{𝟔𝟑}`$ GeV) are shown as an addition to the expected $`𝝂\overline{𝝂}𝜸\mathbf{(}𝜸\mathbf{)}`$ background in Figure 6. The added signal contribution is that which would be expected from a cross-section equal to the 95% CL upper limit. ##### Likelihood Method: <br> The likelihood-based analysis is a straightforward extension of an extended maximum likelihood fit. Upper limits are calculated using information from the photon energy and angular distributions as well as the total number of observed events. The number of observed events, number of expected background events, and signal efficiencies are the same as in the event-counting method. The likelihood function is given by $$𝑳\mathbf{(}𝝈_𝒔\mathbf{)}\mathbf{=}𝓟\mathbf{(}\mathbf{(}𝒏_𝒃\mathbf{+}𝒏_𝒔\mathbf{)}\mathbf{}𝒏_{𝒐𝒃𝒔}\mathbf{)}\mathbf{}\underset{𝒊}{\mathbf{}}\mathbf{[}𝒇_𝒃𝑷_𝒃\mathbf{(}𝑬_𝒊\mathbf{,}𝜽_𝒊\mathbf{)}\mathbf{+}𝒇_𝒔𝑷_𝒔\mathbf{(}𝑬_𝒊\mathbf{,}𝜽_𝒊\mathbf{)}\mathbf{]}$$ (5) where the product is over all selected events in the data, and * $`𝒏_𝒃`$ = expected number of background events, * $`𝒏_𝒔\mathbf{=}\mathit{ϵ}_𝒔𝝈_𝒔𝓛\mathbf{=}`$ expected number of signal events, with $`\mathit{ϵ}_𝒔`$ the efficiency for observing signal events, $`𝝈_𝒔`$ the signal cross-section times branching ratio, and $`𝓛`$ the integrated luminosity, * $`𝒏_{𝒐𝒃𝒔}`$ = number of candidate events observed in the data, * $`𝓟\mathbf{(}\mathbf{(}𝒏_𝒃\mathbf{+}𝒏_𝒔\mathbf{)}\mathbf{}𝒏_{𝒐𝒃𝒔}\mathbf{)}\mathbf{=}𝒆^{\mathbf{}\mathbf{(}𝒏_𝒃\mathbf{+}𝒏_𝒔\mathbf{)}}\mathbf{(}𝒏_𝒃\mathbf{+}𝒏_𝒔\mathbf{)}^{𝒏_{𝒐𝒃𝒔}}\mathbf{/}𝒏_{𝒐𝒃𝒔}\mathbf{!}`$ = the Poisson probability to observe $`𝒏_{𝒐𝒃𝒔}`$ events given an expectation of $`\mathbf{(}𝒏_𝒃\mathbf{+}𝒏_𝒔\mathbf{)}`$, * $`𝑷_{𝒃\mathbf{,}𝒔}\mathbf{(}𝑬_𝒊\mathbf{,}𝜽_𝒊\mathbf{)}`$ = probability density (normalized to one) for a photon $`𝒊`$ resulting from a background ($`𝒃`$) or signal ($`𝒔`$) process to have an energy $`𝑬_𝒊`$ and polar angle $`𝜽_𝒊`$, and * $`𝒇_{𝒃\mathbf{,}𝒔}\mathbf{=}𝒏_{𝒃\mathbf{,}𝒔}\mathbf{/}\mathbf{(}𝒏_𝒃\mathbf{+}𝒏_𝒔\mathbf{)}`$ = relative fraction of background or signal events. The expected number of background events, $`𝒏_𝒃`$, is determined from a Monte Carlo sample of events generated by KORALZ. The signal efficiency, $`\mathit{ϵ}_𝒔`$, is calculated by integrating an energy and angular distribution function over the region of kinematic acceptance (including kinematic-consistency and degraded-resolution cuts), and scaling the result to account for additional efficiency losses due to other cuts in the single-photon selection. This energy and angular distribution function serves as the the signal probability density function $`𝑷_𝒔\mathbf{(}𝑬\mathbf{,}𝜽\mathbf{)}`$. The background probability density function $`𝑷_𝒃\mathbf{(}𝑬\mathbf{,}𝜽\mathbf{)}`$ is obtained from a parameterization of the distributions formed by simulated $`𝝂\overline{𝝂}𝜸\mathbf{(}𝜸\mathbf{)}`$ events generated by KORALZ and NUNUGPV98. In the region $`\mathbf{\|}\mathrm{𝐜𝐨𝐬}𝜽\mathbf{\|}\mathbf{<}\mathbf{0.72}`$, the photon energy and angular distributions are independent, with the angular distribution given by 1/$`\mathrm{𝐬𝐢𝐧}^\mathrm{𝟐}𝜽`$ and the energy distribution determined with a parameterization. In the region $`\mathbf{\|}\mathrm{𝐜𝐨𝐬}𝜽\mathbf{\|}\mathbf{>}\mathbf{0.72}`$, energy resolution is dependent on $`\mathrm{𝐜𝐨𝐬}𝜽`$. In this region, energy distributions are parameterized separately for slices of width 0.1 in $`\mathbf{\|}\mathrm{𝐜𝐨𝐬}𝜽\mathbf{\|}`$; when properly normalized, these parameterizations together form the 2-dimensional probability distribution. The value for $`𝝈_𝒔`$ is not fixed; it is instead treated as a free parameter in the likelihood function. When properly normalized, the likelihood function can be thought of as a probability density function for a hypothesized $`𝝈_𝒔`$ consistent with the observed data. Restricting $`𝝈_𝒔`$ to be non-negative, the 95% CL upper limit $`𝝈^{\mathrm{𝟗𝟓}}`$ is therefore determined from the following equation: $$\mathbf{0.95}\mathbf{=}\frac{{\displaystyle \mathbf{}_\mathrm{𝟎}^{𝝈^{\mathrm{𝟗𝟓}}}}𝑳\mathbf{(}𝝈_𝒔\mathbf{)}𝒅𝝈_𝒔}{{\displaystyle \mathbf{}_\mathrm{𝟎}^{\mathbf{}}}𝑳\mathbf{(}𝝈_𝒔\mathbf{)}𝒅𝝈_𝒔}\mathbf{.}$$ (6) The validity of the likelihood-based method was tested with Monte Carlo simulations. Two types of test were performed. The first test relies on the following definition of a 95% CL limit: if the true signal cross-section happened to be equal to the 95% CL limit, then at least 95% of a large number of identical experiments would result in data that are more signal-like (as defined by a likelihood comparison) than the experiment from which the 95% CL likelihood was derived. This type of test was performed for various combinations of large and small numbers of expected signal and background events, using events generated randomly according to distributions derived from the expected XY signal and $`𝝂\overline{𝝂}𝜸\mathbf{(}𝜸\mathbf{)}`$ background. The second type of test involved the calculation of a large number of 95% CL upper limits from samples generated with an identical signal cross-section. In this case, at least 95% of the calculated upper limits should be larger than the actual signal cross-section. Both types of test yielded agreement with expectations to within statistical errors (typically less than 0.5%). In addition, tests with fully simulated Monte Carlo events were used to test for any biases due to the signal and background parameterizations. No significant biases were found. Sources of systematic errors are the same as for the event-counting method, and are treated in the same manner. However, when calculating limits using the likelihood method, an additional 4% relative uncertainty in experimental sensitivity has been added to account for effects due to uncertainties in the background and signal parameterizations. The 4% estimate is the result of tests using different parameterizations; the uncertainty is again treated according to the method given in reference. The 95% CL upper limits resulting from the likelihood method are shown as a function of $`𝑴_𝐗`$ and $`𝑴_𝐘`$ in Figure 10. The values range from 23 fb to 371 fb. #### 5.1.3 Search for $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐘}`$, $`𝐗\mathbf{}𝐘𝜸`$ ; Special case: $`𝑴_𝐘\mathbf{}\mathrm{𝟎}`$ The case $`𝑴_𝐘\mathbf{}\mathrm{𝟎}`$ is applicable to excited neutrino models and to some supersymmetric models mentioned earlier. The results presented above include this case and no separate analysis is performed, although the results are highlighted here. The upper limits on $`𝝈\mathbf{(}𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐘}\mathbf{)}\mathbf{}𝐁𝑹\mathbf{(}𝐗\mathbf{}𝐘𝜸\mathbf{)}`$ for $`𝑴_𝐘\mathbf{}\mathrm{𝟎}`$ as a function of $`𝑴_𝐗`$ range from 23 to 107 fb using the likelihood method, and from 43 to 170 fb using the event-counting method. The limits are shown in Figure 11 for both methods. Also shown is the average limit expected in the absence of signal for the event-counting method. The actual limit is significantly lower than the expected limit (see discussion in Section 5.1.2). Table 5 gives more details of the search results. #### 5.1.4 Gravitino Pair Production A supersymmetric model has been proposed in which the gravitino is very light. This model predicts a new source of single-photon events from the process $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\stackrel{\mathbf{~}}{𝐆}\stackrel{\mathbf{~}}{𝐆}𝜸`$. We use the single-photon topology to place constraints on the cross-section and therefore the gravitino mass in this model. The differential cross-section is given by $$\frac{𝒅^\mathrm{𝟐}𝝈}{𝒅𝒙_𝜸𝒅\mathrm{𝐜𝐨𝐬}𝜽}\mathbf{=}\mathbf{\left(}\frac{𝜶𝑮_{𝑵}^{}{}_{}{}^{\mathrm{𝟐}}}{\mathrm{𝟒𝟓}}\mathbf{\right)}\frac{𝒔^\mathrm{𝟑}}{𝒎_{\stackrel{\mathbf{~}}{𝐆}}^{}{}_{}{}^{\mathrm{𝟒}}}𝒇_{\stackrel{\mathbf{~}}{𝐆}\stackrel{\mathbf{~}}{𝐆}𝜸}\mathbf{(}𝒙_𝜸\mathbf{,}\mathrm{𝐜𝐨𝐬}𝜽\mathbf{)}$$ (7) where $`𝜶`$ is the fine structure constant, $`𝑮_𝑵`$ is the gravitational constant, $`𝒎_{\stackrel{\mathbf{~}}{𝐆}}`$ is the gravitino mass, $`𝒙_𝜸`$ is the photon scaled energy ($`𝑬_𝜸\mathbf{/}𝑬_{\mathrm{𝐛𝐞𝐚𝐦}}`$) and $`𝜽`$ is the polar angle, with $$𝒇_{\stackrel{\mathbf{~}}{𝐆}\stackrel{\mathbf{~}}{𝐆}𝜸}\mathbf{(}𝒙\mathbf{,}\mathrm{𝐜𝐨𝐬}𝜽\mathbf{)}\mathbf{=}\mathrm{𝟐}\mathbf{(}\mathrm{𝟏}\mathbf{}𝒙\mathbf{)}^\mathrm{𝟐}\mathbf{\left[}\frac{\mathbf{(}\mathrm{𝟏}\mathbf{}𝒙\mathbf{)}\mathbf{(}\mathrm{𝟐}\mathbf{}\mathrm{𝟐}𝒙\mathbf{+}𝒙^\mathrm{𝟐}\mathbf{)}}{𝒙\mathrm{𝐬𝐢𝐧}^\mathrm{𝟐}𝜽}\mathbf{+}\frac{𝒙\mathbf{(}\mathbf{}\mathrm{𝟔}\mathbf{+}\mathrm{𝟔}𝒙\mathbf{+}𝒙^\mathrm{𝟐}\mathbf{)}}{\mathrm{𝟏𝟔}}\mathbf{}\frac{𝒙^\mathrm{𝟑}\mathrm{𝐬𝐢𝐧}^\mathrm{𝟐}𝜽}{\mathrm{𝟑𝟐}}\mathbf{\right]}\mathbf{,}$$ (8) leading to a soft photon energy spectrum. Two methods were used to determine the efficiency for observing events from $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\stackrel{\mathbf{~}}{𝐆}\stackrel{\mathbf{~}}{𝐆}𝜸`$. The first is by reweighting simulated $`𝝂\overline{𝝂}𝜸\mathbf{(}𝜸\mathbf{)}`$ events generated by KORALZ. Each simulated event is weighted by $`𝒇_{\stackrel{\mathbf{~}}{𝐆}\stackrel{\mathbf{~}}{𝐆}𝜸}\mathbf{(}𝒙\mathbf{,}\mathrm{𝐜𝐨𝐬}𝜽\mathbf{)}\mathbf{/}𝒇_{𝑲𝒁}\mathbf{(}𝒙\mathbf{,}\mathrm{𝐜𝐨𝐬}𝜽\mathbf{)}`$, where $`𝒇_{𝑲𝒁}`$ is the 2-dimensional energy and angular distribution for the most energetic photon in events generated by KORALZ. We use a parameterization $`𝒇_{𝑲𝒁}`$ which is valid up to photon energies of about 60 GeV, and thus we establish a kinematic acceptance region for $`\stackrel{\mathbf{~}}{𝐆}\stackrel{\mathbf{~}}{𝐆}𝜸`$ production: $`𝒙_𝑻\mathbf{>}\mathbf{0.05}`$, $`\mathrm{𝟏𝟓}^{\mathbf{}}\mathbf{<}𝜽\mathbf{<}\mathrm{𝟏𝟔𝟓}^{\mathbf{}}`$, and $`𝑬_𝜸\mathbf{<}\mathrm{𝟔𝟎}`$ GeV. We note that the cross-section for $`\stackrel{\mathbf{~}}{𝐆}\stackrel{\mathbf{~}}{𝐆}𝜸`$ production with photon energies above 60 GeV is negligible. In addition, a Monte Carlo generator was written to generate events according to the distribution given in Equation 8. Initial state radiation was treated in a manner identical to that used by the EXOTIC generator. Efficiencies were calculated and found to be identical to those using the reweighting method to within statistical errors. To achieve maximum sensitivity<sup>7</sup><sup>7</sup>7 The optimization condition chosen was that the expected upper limit on the cross-section for contributions from the $`\mathrm{e}^+\mathrm{e}^{}\stackrel{~}{\mathrm{G}}\stackrel{~}{\mathrm{G}}\gamma `$ signal process be minimised, where the expected upper limit is defined as the average limit one would expect to set in the absence of signal. , we place a maximum energy requirement on the observed photon energy: $`𝑬_𝜸\mathbf{<}\mathrm{𝟑𝟎}`$ GeV. With this requirement, we observe 195 candidates with an expected $`𝝂\overline{𝝂}𝜸`$ background of $`\mathbf{179.6}\mathbf{\pm }\mathbf{5.4}`$. The efficiency within the above kinematic acceptance region is estimated by the Monte Carlo to be 82.8%, and was assigned a relative 5% uncertainty. Using the event-counting method described earlier, we place a 95% CL cross-section upper limit of 293 fb, giving a lower limit<sup>8</sup><sup>8</sup>8Evaluated with $`\alpha =\frac{1}{128}`$ on the gravitino mass of 8.7 $`𝝁`$eV. #### 5.1.5 Graviton-Photon: Search for Extra Space Dimensions There has been recent interest in string theory models which postulate the existence of additional compactified space dimensions; these models allow gravitons to propagate freely in the higher-dimensional space while restricting Standard Model particles to a 3+1 dimensional hypersurface. The fundamental mass scale in this class of theories, $`𝑴_𝑫`$, governs the rate of graviton production; it is possible that direct graviton-photon production could occur at significant rates at LEP2 energies with an experimental signature of a single photon with missing energy. The single-photon search topology can therefore be used to place constraints on the fundamental mass scale $`𝑴_𝑫`$ (or, equivalently, on the size of the extra dimensions) by placing limits on the graviton-photon cross-section. The differential cross-section is given by $$\frac{𝒅^\mathrm{𝟐}𝝈}{𝒅𝒙_𝜸𝒅\mathrm{𝐜𝐨𝐬}𝜽}\mathbf{=}\frac{𝜶𝑺_{𝜹\mathbf{}\mathrm{𝟏}}}{\mathrm{𝟔𝟒}𝑴_{𝑫}^{}{}_{}{}^{\mathrm{𝟐}}}\mathbf{\left(}\frac{\sqrt{𝒔}}{𝑴_𝑫}\mathbf{\right)}^𝜹𝒇_{𝑮𝜸}\mathbf{(}𝒙_𝜸\mathbf{,}\mathrm{𝐜𝐨𝐬}𝜽\mathbf{)}$$ (9) where $`𝜹`$ is the number of extra dimensions and $`𝑺_{𝜹\mathbf{}\mathrm{𝟏}}`$ is the surface area of a $`𝜹`$-dimensional sphere of unit radius, with $$𝒇_{𝑮𝜸}\mathbf{(}𝒙\mathbf{,}\mathrm{𝐜𝐨𝐬}𝜽\mathbf{)}\mathbf{=}\frac{\mathrm{𝟐}\mathbf{(}\mathrm{𝟏}\mathbf{}𝒙\mathbf{)}^{\frac{𝜹}{\mathrm{𝟐}}\mathbf{}\mathrm{𝟏}}}{𝒙\mathbf{(}\mathrm{𝟏}\mathbf{}\mathrm{𝐜𝐨𝐬}^\mathrm{𝟐}𝜽\mathbf{)}}\mathbf{\left[}\mathbf{(}\mathrm{𝟐}\mathbf{}𝒙\mathbf{)}^\mathrm{𝟐}\mathbf{(}\mathrm{𝟏}\mathbf{}𝒙\mathbf{+}𝒙^\mathrm{𝟐}\mathbf{)}\mathbf{}\mathrm{𝟑}𝒙^\mathrm{𝟐}\mathrm{𝐜𝐨𝐬}^\mathrm{𝟐}𝜽\mathbf{(}\mathrm{𝟏}\mathbf{}𝒙\mathbf{)}\mathbf{}𝒙^\mathrm{𝟒}\mathrm{𝐜𝐨𝐬}^\mathrm{𝟒}𝜽\mathbf{\right]}\mathbf{.}$$ (10) We use the event-counting method to place limits on graviton-photon production in the cases $`\mathrm{𝟐}\mathbf{}𝜹\mathbf{}\mathrm{𝟕}`$. We use the same kinematic acceptance region as described in Section 5.1.4: $`𝒙_𝑻\mathbf{>}\mathbf{0.05}`$, $`𝑬_𝜸\mathbf{<}\mathrm{𝟔𝟎}`$ GeV, and $`\mathrm{𝟏𝟓}^{\mathbf{}}\mathbf{<}𝜽\mathbf{<}\mathrm{𝟏𝟔𝟓}^{\mathbf{}}`$. The expected photon energy spectrum is soft, so in order to improve sensitivity, we require observed photon energies to be less than 34 GeV for all values of $`𝜹`$. The efficiency for observing events from $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝑮𝜸`$ is determined by reweighting simulated $`𝝂\overline{𝝂}𝜸\mathbf{(}𝜸\mathbf{)}`$ events generated by KORALZ in a manner equivalent to that described in Section 5.1.4: simulated events are weighted by $`𝒇_{𝑮𝜸}\mathbf{(}𝒙\mathbf{,}\mathrm{𝐜𝐨𝐬}𝜽\mathbf{)}\mathbf{/}𝒇_{𝑲𝒁}\mathbf{(}𝒙\mathbf{,}\mathrm{𝐜𝐨𝐬}𝜽\mathbf{)}`$, where $`𝒇_{𝑲𝒁}`$ is the same KORALZ parameterization used earlier. The distributions $`𝒇_{𝑲𝒁}`$ and $`𝒇_{𝑮𝜸}`$ have roughly similar shapes, and the calculated efficiency is insensitive to small differences in the KORALZ parameterization. We observe 208 candidates with an expected $`𝝂\overline{𝝂}𝜸`$ contribution of $`\mathbf{196.0}\mathbf{\pm }\mathbf{5.9}`$ events. Assuming a 5% relative uncertainty in the signal efficiency, upper limits on the cross-section and corresponding lower limits on $`𝑴_𝑫`$ for $`\mathrm{𝟐}\mathbf{}𝜹\mathbf{}\mathrm{𝟕}`$ have been calculated<sup>9</sup><sup>9</sup>9Evaluated also with $`\alpha =\frac{1}{128}`$. The results are given in Table 6 based on the convention<sup>10</sup><sup>10</sup>10The convention of reference differs. of Equation 2 in reference. ### 5.2 Acoplanar-Photons #### 5.2.1 Search for $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐗}`$, $`𝐗\mathbf{}𝐘𝜸`$ ; General case: $`𝑴_𝐘\mathbf{}\mathrm{𝟎}`$ The searches for $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐗}`$, $`𝐗\mathbf{}𝐘𝜸`$, both for the general case discussed here and the special case of $`𝑴_𝐘\mathbf{}\mathrm{𝟎}`$ discussed in 5.2.2, use the methods described in our previous publication. Selected events are classified as consistent with a given value of $`𝑴_𝐗`$ and $`𝑴_𝐘`$ if the energy of each of the photons falls within the region kinematically accessible to photons from the process $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐗}`$, $`𝐗\mathbf{}𝐘𝜸`$, including resolution effects. The selection efficiencies at each generated grid point for the $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐗}`$, $`𝐗\mathbf{}𝐘𝜸`$ Monte Carlo events are shown in Table 7. These values include the efficiency of the kinematic consistency requirement which is higher than 95% at each generated point in the $`\mathbf{(}𝑴_𝐗\mathbf{,}𝑴_𝐘\mathbf{)}`$ plane for which $`𝑴_𝐗\mathbf{}𝑴_𝐘\mathbf{>}\mathrm{𝟓}`$ GeV. Events from $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂\overline{𝝂}𝜸𝜸\mathbf{(}𝜸\mathbf{)}`$ are typically characterized by a high-energy photon from the radiative return to the $`𝐙^\mathrm{𝟎}`$ and a second lower energy photon. The kinematic consistency requirement is such that the two photons must have energies within the same (kinematically accessible) region. As $`𝑴_𝐗`$ and $`𝑴_𝐘`$ increase, the allowed range of energy for the photons narrows, and fewer $`𝝂\overline{𝝂}𝜸𝜸\mathbf{(}𝜸\mathbf{)}`$ events will be accepted. For the 24 selected events, the distribution of the number of events consistent with a given mass point ($`𝑴_𝐗`$,$`𝑴_𝐘`$) is consistent with the expectation from $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂\overline{𝝂}𝜸𝜸\mathbf{(}𝜸\mathbf{)}`$ Monte Carlo, over the full ($`𝑴_𝐗`$,$`𝑴_𝐘`$) plane. In our previous publication, because of concerns about the modelling of the Standard Model process $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂\overline{𝝂}𝜸𝜸\mathbf{(}𝜸\mathbf{)}`$, all limits derived from the acoplanar photons analysis were obtained without accounting for the expected background. The theoretical situation is now greatly improved, with two event generators agreeing to better than 1% for the total cross-section for this process within the kinematic acceptance of this analysis. For that reason, in this paper, all limits derived from this selection have been calculated taking the (KORALZ) background estimate into account. Figure 12 shows the 95% CL exclusion regions for $`𝝈\mathbf{(}𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐗}\mathbf{)}\mathbf{}𝐁𝑹^\mathrm{𝟐}\mathbf{(}𝐗\mathbf{}𝐘𝜸\mathbf{)}`$. The limits vary from 33 fb to 103 fb for $`𝑴_𝐗\mathbf{>}\mathrm{𝟒𝟓}`$ GeV and $`𝑴_𝐗\mathbf{}𝑴_𝐘\mathbf{>}\mathrm{𝟓}`$ GeV. In the region 2.5 GeV $`\mathbf{}𝑴_𝐗\mathbf{}𝑴_𝐘\mathbf{<}\mathbf{5.0}`$ GeV, the efficiency falls off rapidly (see Table 7). As this rapid fall increases the associated uncertainty in the efficiency, no limits have been set in this region. Systematic errors are due primarily to limited Monte Carlo statistics at the generated ($`𝑴_𝐗\mathbf{,}𝑴_𝐘`$) points and the uncertainty on the efficiency parameterization across the ($`𝑴_𝐗\mathbf{,}𝑴_𝐘`$) plane. The combined relative uncertainty on the efficiency varies from about (3-6)% across the plane (for $`𝑴_𝐗\mathbf{}𝑴_𝐘\mathbf{>}\mathrm{𝟓}`$ GeV). All systematic uncertainties are accounted for in the manner advocated in reference. This also applies to the limits for the $`𝑴_𝐘\mathbf{}\mathrm{𝟎}`$ case, presented in the next section. #### 5.2.2 Search for $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐗}`$, $`𝐗\mathbf{}𝐘𝜸`$ ; Special case: $`𝑴_𝐘\mathbf{}\mathrm{𝟎}`$ For the special case of $`𝑴_𝐘\mathbf{}\mathrm{𝟎}`$ the kinematic consistency requirements differ from those used for the general case. One can calculate the maximum mass, $`𝑴_𝐗^{\mathrm{𝐦𝐚𝐱}}`$, which is consistent with the measured three-momenta of the two photons, assuming a massless $`𝐘`$. A cut on $`𝑴_𝐗^{\mathrm{𝐦𝐚𝐱}}`$ provides further suppression of the $`𝝂\overline{𝝂}𝜸𝜸\mathbf{(}𝜸\mathbf{)}`$ background while retaining high efficiency for the signal hypothesis. This is discussed in more detail in reference. We require that the maximum kinematically allowed mass be greater than $`𝑴_𝐗\mathbf{}\mathrm{𝟓}`$ GeV, which retains $`\mathbf{(}\mathbf{95.5}_{\mathbf{}\mathbf{1.0}}^{\mathbf{+}\mathbf{2.0}}\mathbf{)}`$% relative efficiency for signal at all values of $`𝑴_𝐗`$ while suppressing much of the remaining $`𝝂\overline{𝝂}𝜸𝜸\mathbf{(}𝜸\mathbf{)}`$ background. In our previous publication we also applied a recoil-mass cut at 80 GeV. In the case where background is not accounted for in the limit calculation, such a cut improves the expected sensitivity of the analysis. With background subtraction this is no longer the case. Therefore that cut has been removed. Figure 13 shows the expected $`𝑴_𝐗^{\mathrm{𝐦𝐚𝐱}}`$ distribution for signal Monte Carlo events with $`𝑴_𝐗\mathbf{=}\mathrm{𝟗𝟎}`$ GeV and for $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂\overline{𝝂}𝜸𝜸\mathbf{(}𝜸\mathbf{)}`$ Monte Carlo events. Also shown is the distribution of the selected data events. For the $`𝑴_𝐘\mathbf{}\mathrm{𝟎}`$ case, the efficiencies calculated from Monte Carlo events are shown in Table 8 after application of the event selection criteria and then after the cut on $`𝑴_𝐗^{\mathrm{𝐦𝐚𝐱}}`$. Also shown in Table 8 are the number of selected events consistent with each value of $`𝑴_𝐗`$ as well as the expected number of events from $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂\overline{𝝂}𝜸𝜸\mathbf{(}𝜸\mathbf{)}`$. The number of selected events consistent with a given value of $`𝑴_𝐗`$ varies from 14, for $`𝑴_𝐗\mathbf{}`$ 45 GeV, to 3 events at the kinematic limit. The expected number of events decreases from $`\mathbf{15.8}\mathbf{\pm }\mathbf{0.7}`$ at $`𝑴_𝐗\mathbf{}\mathrm{𝟒𝟓}`$ GeV to $`\mathbf{1.34}\mathbf{\pm }\mathbf{0.07}`$ consistent with $`𝑴_𝐗\mathbf{}\mathrm{𝟗𝟒}`$ GeV. Based on the efficiencies and the number of selected events, we calculate a 95% CL upper limit on $`𝝈\mathbf{(}𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐗}\mathbf{)}\mathbf{}𝐁𝑹^\mathrm{𝟐}\mathbf{(}𝐗\mathbf{}𝐘𝜸\mathbf{)}`$ for $`𝑴_𝐘\mathbf{}\mathrm{𝟎}`$ as a function of $`𝑴_𝐗`$. This is shown as the solid line in Figure 14. The limit is between 50 and 80 fb for $`𝑴_𝐗`$ values from 45 GeV up to the kinematic limit. Also shown as a dashed line is the expected limit, defined as the average limit one would expect to set in the absence of signal. The limits can be used to set model-dependent limits on the mass of the lightest neutralino in supersymmetric models in which the NLSP is the lightest neutralino and the LSP is a light gravitino ($`𝐗\mathbf{=}\stackrel{\mathbf{~}}{𝝌}_{}^{\mathrm{𝟎}}{}_{\mathrm{𝟏}}{}^{}\mathbf{,}𝐘\mathbf{=}\stackrel{\mathbf{~}}{𝐆}`$). Shown in Figure 14 as a dotted line is the (Born-level) cross-section prediction from a specific light gravitino LSP model in which the neutralino composition is purely bino, with $`𝒎_{\stackrel{\mathbf{~}}{𝒆}_𝑹}\mathbf{=}\mathbf{1.35}𝒎_{\stackrel{\mathbf{~}}{𝝌}_{}^{\mathrm{𝟎}}{}_{\mathrm{𝟏}}{}^{}}`$ and $`𝒎_{\stackrel{\mathbf{~}}{𝒆}_𝑳}\mathbf{=}\mathbf{2.7}𝒎_{\stackrel{\mathbf{~}}{𝝌}_{}^{\mathrm{𝟎}}{}_{\mathrm{𝟏}}{}^{}}`$. Within the framework of this model, $`\stackrel{\mathbf{~}}{𝝌}_{}^{\mathrm{𝟎}}{}_{\mathrm{𝟏}}{}^{}`$ masses between 45 and 88.3 GeV are excluded at 95% CL. As described in section 2, the efficiencies over the full angular range have been calculated using isotropic angular distributions for production and decay of $`𝐗`$. The validity of this model has been examined based on the angular distributions calculated for photino pair production in reference. For models proposed in reference, the production angular distributions are more central and so this procedure is conservative. For a $`\mathrm{𝟏}\mathbf{+}\mathrm{𝐜𝐨𝐬}^\mathrm{𝟐}𝜽`$ production angular distribution expected for t-channel exchange of a very heavy particle according to reference, the relative efficiency reduction would be less than 2% at all points in the ($`𝑴_𝐗\mathbf{,}𝑴_𝐘`$) plane. #### 5.2.3 Search for $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐗}`$, $`𝐗\mathbf{}𝐘𝜸`$ ; Special case: $`𝑴_𝐘\mathbf{}\mathrm{𝟎}`$ and macroscopic decay length As an extension to the special case of $`𝑴_𝐘\mathbf{}\mathrm{𝟎}`$, we consider the sensitivity of the acoplanar-photons search to $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐗}`$, $`𝐗\mathbf{}𝐘𝜸`$ when $`𝐗`$ has a macroscopic decay length. This extension evaluates the selection efficiency using signal Monte Carlo samples with various lifetime values $`𝝉_𝑿`$. The $`𝐗\mathbf{}𝐘𝜸`$ decay is treated by a modified OPAL detector simulation package designed to handle delayed decays (as described in section 2). The range of lifetimes considered extends from the near-zero lifetime value of $`𝝉_𝑿\mathbf{=}\mathrm{𝟏𝟎}^{\mathbf{}\mathrm{𝟏𝟓}}`$ s to $`𝝉_𝑿\mathbf{=}\mathrm{𝟏𝟎}^\mathbf{}\mathrm{𝟕}`$ s at which point c$`𝝉`$ is 30 m. Such long lifetimes lead to decays outside the sensitive volume of the OPAL detector (determined by the outer radius of the ECAL), which implies a natural cutoff in sensitivity. As with the zero lifetime case, the ($`𝑴_𝐗^{\mathrm{𝐦𝐚𝐱}}\mathbf{>}𝑴_𝐗\mathbf{}\mathrm{𝟓}`$ GeV) requirement is applied. The relative efficiency of this cut decreases with increasing lifetime, but exceeds 90% for all $`\mathbf{(}𝑴_𝐗\mathbf{,}𝝉_𝑿\mathbf{)}`$ combinations with high (greater than 50%) selection efficiencies. For $`\mathbf{(}𝑴_𝐗\mathbf{,}𝝉_𝑿\mathbf{)}`$ combinations with lower selection efficiencies, the relative efficiency falls more rapidly with increasing $`𝝉_𝑿`$, but does not drop below 60% for the $`\mathbf{(}𝑴_𝐗\mathbf{,}𝝉_𝑿\mathbf{)}`$ combinations considered. The loss in relative efficiency of the $`𝑴_𝐗^{\mathrm{𝐦𝐚𝐱}}`$ cut is expected due to the definition of $`𝑴_𝐗^{\mathrm{𝐦𝐚𝐱}}`$, but the cut is maintained to allow comparison between the prompt decay and macroscopic decay length selection efficiencies. The resulting selection efficiencies for macroscopic decay lengths are given in Table 9. Also listed in Table 9 are the number of events observed in the data and the expected $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂\overline{𝝂}𝜸𝜸\mathbf{(}𝜸\mathbf{)}`$ background. The efficiencies indicate that sensitivity to macroscopic decay lengths is maintained up to lifetimes of $`\mathrm{𝟏𝟎}^\mathbf{}\mathrm{𝟗}`$ s but falls rapidly for longer lifetimes. This is as expected, since for $`𝝉_𝑿\mathbf{=}\mathrm{𝟏𝟎}^\mathbf{}\mathrm{𝟗}`$ s the decay length ranges from 30 to 3 cm for the values of $`𝑴_𝐗`$ considered, so that the $`𝐗\mathbf{}𝐘𝜸`$ decays typically occur well within the confines of the detector. For lifetimes greater than $`\mathrm{𝟏𝟎}^\mathbf{}\mathrm{𝟗}`$ s, the decrease in efficiency is two-fold. The primary loss is due to the increased decay length, implying fewer $`𝐗\mathbf{}𝐘𝜸`$ decays within the sensitive detector volume. Selection efficiency is also lost with increasing $`𝝉_𝑿`$ since larger lifetimes result in a delayed arrival time of the photon at the ECAL, which, if sufficiently late, causes the event to be vetoed due to the timing cuts imposed by the analysis. The effect of the event kinematics is also seen in the increase in selection efficiency for a given lifetime as $`𝑴_𝐗`$ approaches its threshold value, corresponding to a drop in $`𝜷`$ from 0.84 at $`𝑴_𝐗`$=50 GeV to 0.10 at $`𝑴_𝐗`$=94 GeV. Contributions to the systematic error in the selection efficiency are (in order of significance) the limited Monte Carlo statistics, the efficiency parameterisation in the $`\mathbf{(}𝑴_𝐗\mathbf{,}𝝉_𝑿\mathbf{)}`$ plane and the modelling of the timing cuts. The total contributions are given in Table 9. Given the efficiencies and the number of selected events, the 95% CL upper limit on $`𝝈\mathbf{(}𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐗}\mathbf{)}\mathbf{}𝐁𝑹^\mathrm{𝟐}\mathbf{(}𝐗\mathbf{}𝐘𝜸\mathbf{)}`$ for $`𝑴_𝐘\mathbf{}\mathrm{𝟎}`$ as a function of $`𝑴_𝐗`$ and $`𝝉_𝑿`$ is calculated, and is shown in figure 15. The background-subtracted limits from the prompt decay analysis are used to define the profile as a function of $`𝑴_𝐗`$ for $`𝝉_𝑿\mathbf{=}\mathrm{𝟏𝟎}^{\mathbf{}\mathrm{𝟏𝟓}}`$ s. The interpolation of the limit to larger lifetimes is done by factoring the selection efficiencies of Table 9 into the 95% confidence level limits on $`𝝈\mathbf{(}𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐗}\mathbf{)}\mathbf{}𝐁𝑹^\mathrm{𝟐}\mathbf{(}𝐗\mathbf{}𝐘𝜸\mathbf{)}`$ for $`𝑴_𝐘\mathbf{}\mathrm{𝟎}`$. As in Section 5.2.2 the Born-level cross-section of the neutralino NLSP Gravitino LSP model of is used to give an exclusion region in the (neutralino mass, lifetime) plane. This exclusion region is superimposed on Figure 15 and gives the excluded domain in the ($`𝑴_𝐗`$,$`𝝉_𝑿`$) plane within the same specific light gravitino LSP model discussed above. ## 6 Conclusions We have searched for photonic events with large missing energy in two topologies in data taken with the OPAL detector at LEP, at a centre-of-mass energy of 189 GeV. In the single-photon selection, which requires at least one photon with $`𝒙_𝑻\mathbf{>}\mathbf{0.05}`$ in the region $`\mathrm{𝟏𝟓}^{\mathbf{}}\mathbf{<}𝜽\mathbf{<}\mathrm{𝟏𝟔𝟓}^{\mathbf{}}`$ ($`\mathbf{\|}\mathrm{𝐜𝐨𝐬}𝜽\mathbf{\|}\mathbf{<}\mathbf{0.966}\mathbf{)}`$, 643 events are observed in the data. The background-subtracted cross-section measurement of 4.35 $`\mathbf{\pm }`$ 0.17 (stat) $`\mathbf{\pm }`$ 0.09 (sys) pb is consistent with the KORALZ prediction of 4.66 pb from the Standard Model $`𝝂\overline{𝝂}𝜸\mathbf{(}𝜸\mathbf{)}`$ process. Interpreting the results as a measurement of the effective number of light neutrino species, we measure $`𝐍_𝝂\mathbf{=}\mathbf{2.69}\mathbf{\pm }\mathbf{0.13}\mathbf{(}\mathrm{𝐬𝐭𝐚𝐭}\mathbf{)}\mathbf{\pm }\mathbf{0.11}\mathbf{(}\mathrm{𝐬𝐲𝐬}\mathbf{)}`$. We also observe significant W contributions to the cross-section with a rate consistent with the Standard Model expectation. We calculate upper limits on the cross-section times branching ratio for the process $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐘}`$, $`𝐗\mathbf{}𝐘𝜸`$ using two methods: an event-counting method, which is insensitive to energy or angular distribution shapes and is therefore relatively model-independent, and a likelihood method, which assumes isotropic production and decay angular distributions and has greater sensitivity. In the region of interest in the $`\mathbf{(}𝑴_𝐗\mathbf{,}𝑴_𝐘\mathbf{)}`$ plane, the limits vary from 43 to 409 fb using the event-counting method and from 23 to 371 fb using the likelihood method. These limits include the special case of $`𝑴_𝐘\mathbf{}\mathrm{𝟎}`$, where the limit varies between 43 and 170 fb using the event-counting method and from 23 to 107 fb using the likelihood method. We note that some of these limits are much more stringent than would be expected on average in the absence of signal contributions. We set a 95% CL cross-section upper limit on $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\stackrel{\mathbf{~}}{𝐆}\stackrel{\mathbf{~}}{𝐆}𝜸`$ production of 293 fb implying a lower limit on the gravitino mass of 8.7 $`𝝁`$eV in the superlight gravitino model of reference. In the context of string theory models with extra space dimensions, upper limits on the cross-section for graviton-photon production of between 309 and 271 fb at 95% CL are set, giving lower limits on the fundamental mass scale varying between 1086 and 470 GeV for between 2 and 7 additional dimensions respectively. The acoplanar-photons selection requires at least two photons with scaled energy $`𝒙_𝜸\mathbf{>}\mathbf{0.05}`$ within the polar angle region $`\mathrm{𝟏𝟓}^{\mathbf{}}\mathbf{<}𝜽\mathbf{<}\mathrm{𝟏𝟔𝟓}^{\mathbf{}}`$ or at least two photons with energy $`𝑬_𝜸\mathbf{>}\mathbf{1.75}`$ GeV with one satisfying $`\mathbf{\|}\mathrm{𝐜𝐨𝐬}𝜽\mathbf{\|}\mathbf{<}\mathbf{0.8}`$ and the other satisfying $`\mathrm{𝟏𝟓}^{\mathbf{}}\mathbf{<}𝜽\mathbf{<}\mathrm{𝟏𝟔𝟓}^{\mathbf{}}`$. In each case, the requirement $`𝒑_𝑻^{𝜸𝜸}\mathbf{/}𝑬_{\mathrm{𝐛𝐞𝐚𝐦}}\mathbf{>}\mathbf{0.05}`$ is also applied. There are 24 events selected. The KORALZ prediction for the contribution from $`𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}𝝂\overline{𝝂}𝜸𝜸\mathbf{(}𝜸\mathbf{)}`$ is $`\mathbf{26.9}\mathbf{\pm }\mathbf{1.2}`$ events; contribution from other sources is 0.11 events. The number of events observed in the data and their kinematic distributions are consistent with Standard Model expectations. We derive 95% CL upper limits on $`𝝈\mathbf{(}𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐗}\mathbf{)}\mathbf{}𝐁𝑹^\mathrm{𝟐}\mathbf{(}𝐗\mathbf{}𝐘𝜸\mathbf{)}`$ ranging from 33 to 103 fb for the general case of massive $`𝐗`$ and $`𝐘`$. For the special case of $`𝑴_𝐘\mathbf{}\mathrm{𝟎}`$, the 95% CL upper limits on $`𝝈\mathbf{(}𝐞^\mathbf{+}𝐞^{\mathbf{}}\mathbf{}\mathrm{𝐗𝐗}\mathbf{)}\mathbf{}𝐁𝑹^\mathrm{𝟐}\mathbf{(}𝐗\mathbf{}𝐘𝜸\mathbf{)}`$ range from 50 to 80 fb. For the case of macroscopic decay lengths, these values range from 50 fb at short lifetimes to 8.8 pb for lifetimes as long as $`\mathrm{𝟏𝟎}^\mathbf{}\mathrm{𝟕}`$ s. The results of the acoplanar-photons search are used to place model-dependent lower limits on the $`\stackrel{\mathbf{~}}{𝝌}_{}^{\mathrm{𝟎}}{}_{\mathrm{𝟏}}{}^{}`$ mass in a specific light gravitino LSP model. Masses between 45 and 88.3 GeV are excluded at 95% CL for promptly decaying neutralinos, while for a $`\stackrel{\mathbf{~}}{𝝌}_{}^{\mathrm{𝟎}}{}_{\mathrm{𝟏}}{}^{}`$ lifetime of $`\mathrm{𝟏𝟎}^\mathbf{}\mathrm{𝟖}`$ s ($`𝒄𝝉\mathbf{=}\mathrm{𝟑}`$m), masses between 45 and 81 GeV are excluded. Appendix ## Appendix A Improvements to Event Selections Detailed descriptions of the event selection criteria are given in a previous publication. This appendix describes changes and updates to the selections used in the present publication. ### A.1 Single-Photon Selection There have been several improvements to the single-photon selection: * Non-conversion candidates in the endcap region ($`\mathbf{\|}\mathrm{𝐜𝐨𝐬}𝜽\mathbf{\|}\mathbf{>}\mathbf{0.82}`$) are no longer required to have a good in-time associated TE hit, although candidates with associated out-of-time TE hits are still rejected. This change reduces sensitivity to the modelling of the material between the interaction point and the TE scintillators, and results in a relative increase in efficiency of about 19% in the endcap region with an associated increase in cosmic and beam-related background of about 1.3 events. The special background vetos continue to be applied to all non-conversion candidates in the endcap regardless of the presence of TE timing. * The cluster extent cut has been modified for events in the endcap region. This cut had been fixed at lower centre-of-mass energies and was no longer fully efficient at $`\sqrt{𝒔}\mathbf{=}\mathrm{𝟏𝟖𝟗}`$ GeV, particularly for beam energy photons. Events with the primary photon candidate in the region $`\mathbf{\|}\mathrm{𝐜𝐨𝐬}𝜽\mathbf{\|}\mathbf{>}\mathbf{0.82}`$ are now rejected if $`𝚫\mathit{\varphi }\mathbf{}\mathrm{𝐬𝐢𝐧}^{\mathbf{1.46}}𝜽\mathbf{/}\mathrm{𝐥𝐧}𝑬_𝜸\mathbf{>}\mathbf{0.05}`$ radians<sup>11</sup><sup>11</sup>11 The $`\mathrm{sin}\theta `$ dependence was parameterized for beam energy photons using samples of real and simulated $`\mathrm{e}^+\mathrm{e}^{}\gamma \gamma `$ events., where $`𝚫\mathit{\varphi }`$ is the $`\mathit{\varphi }`$ extent of the cluster and $`𝑬_𝜸`$ is the cluster energy in GeV. The result is a relative increase in efficiency of about 13% in the endcap region with no significant change in expected cosmic and beam-related background. * Small changes have been made to the timing requirements. The timing cuts for the photon candidates with TOF associated hits ($`\mathbf{\|}\mathrm{𝐜𝐨𝐬}𝜽\mathbf{\|}\mathbf{<}\mathbf{0.82}`$) have been relaxed for certain periods of data-taking, resulting in an efficiency increase of 2.3% for events with a photon in $`\mathbf{\|}\mathrm{𝐜𝐨𝐬}𝜽\mathbf{\|}\mathbf{<}\mathbf{0.82}`$ with a TOF associated hit. The expected cosmic ray background increases by about 10%. In addition, the majority of events rejected by poorly measured TOF times are recovered. In cases where the ECAL cluster and the TOF z-measurement from time difference differ by more than 40 cm, modified arrival times based on the ECAL cluster position and the time measurement at each end of the scintillator are constructed. Events are retained if either modified arrival time is within 5 ns of the expected time for a photon originating from the interaction point. As a result, the probability for a signal event to be rejected due to bad TOF timing falls from 0.6% to less than 0.1% with no measurable change in the cosmic and beam-related background. * The efficiency of the photon conversion consistency criteria has been improved substantially. In particular, all reconstructed tracks in the event (usually two for signal events) are now used to test association of a charged track to an ECAL cluster, and the association test also considers the disfavoured solution of the jet-chamber left-right ambiguity to account for cases where the incorrect solution has been chosen. * Improved redundancy in the rejection of beam-related backgrounds for conversion candidates with $`𝒙_𝑻\mathbf{<}\mathbf{0.1}`$ and $`\mathbf{\|}\mathrm{𝐜𝐨𝐬}𝜽\mathbf{\|}\mathbf{>}\mathbf{0.90}`$ is introduced, based on the measured z-difference of the track’s point of closest approach to the interaction point. * An additional beam-halo veto was implemented; events are rejected if ECAL clusters exceeding 100 MeV in deposited energy are found within 20 cm in the $`𝒓\mathbf{}\mathit{\varphi }`$ plane but on opposite ends of the ECAL endcap. This veto provides better rejection of beam-halo background and has negligible effect on signal efficiencies. * A minor addition was made to the special background vetoes used for non-conversion candidates. Events are now explicitly rejected if there is a large amount of hadronic energy in an HCAL cluster associated with the ECAL cluster of the primary photon candidate. This provides better redundancy with existing cuts in terms of background rejection, with negligible inefficiency even for beam energy photons. ### A.2 Acoplanar-Photons Selection The acoplanar-photons selection is almost identical to that described in our previous publication. The only difference involves the timing requirements for candidates with associated TOF hits. These were relaxed for certain periods of the data-taking in the same manner as described above for the single-photon selection. ## Acknowledgements The authors wish to thank G. Giudice, M. Peskin and F. Zwirner for helpful clarifications. We particularly wish to thank the SL Division for the efficient operation of the LEP accelerator at all energies and for their continuing close cooperation with our experimental group. We thank our colleagues from CEA, DAPNIA/SPP, CE-Saclay for their efforts over the years on the time-of-flight and trigger systems which we continue to use. In addition to the support staff at our own institutions we are pleased to acknowledge the Department of Energy, USA, National Science Foundation, USA, Particle Physics and Astronomy Research Council, UK, Natural Sciences and Engineering Research Council, Canada, Israel Science Foundation, administered by the Israel Academy of Science and Humanities, Minerva Gesellschaft, Benoziyo Center for High Energy Physics, Japanese Ministry of Education, Science and Culture (the Monbusho) and a grant under the Monbusho International Science Research Program, Japanese Society for the Promotion of Science (JSPS), German Israeli Bi-national Science Foundation (GIF), Bundesministerium für Bildung, Wissenschaft, Forschung und Technologie, Germany, National Research Council of Canada, Research Corporation, USA, Hungarian Foundation for Scientific Research, OTKA T-029328, T023793 and OTKA F-023259.
warning/0005/hep-th0005176.html	ar5iv	text	# A Note on the Path Integral Representation of the Boundary State of D-brane ## 1 Introduction It is well known that D-branes in play very important roles in the study of nonperturbative string theory dynamics. However, we do not know how to construct an exact theory of D-brane in the general background until now. So we have to limit ourselves within a very small circle, namely in some special backgrounds. This kind of study may have some heuristic implications in seeking a good theory describing D-brane. On the other hand many recent researches show the boundary state description indeed seizes some basic features of D-brane. In fact the concept of the boundary state is introduced into string theory before that of D-brane. The BPS object can be represented by boundary state, which satisfies the constraint $`(Q+\stackrel{~}{Q})\|B=0`$. We can read off the tension and R-R charge from the tree diagram of the boundary state. These are the most important proofs supporting the conjecture that D-brane is the classical soliton solution with R-R charge, namely p-brane. Many relevant literatures have been included in . The concept of the boundary state is fundamental. Firstly, the tadpole can be represented with the boundary state which is the source of closed string. Naively, it seems reasonable that the boundary state should be seen as quantum state of an independent dynamic object arising from string vacuum. D-brane is just defined as that object where open string end points end. Secondly, the boundary state can be defined according to the idea that one loop amplitude of open string can be regarded as tree amplitude of closed string. From this physical picture, the holographic principle can be argued because non-abelian gauge field theory appear from the low energy limits of open string theory and the low energy limit of closed string theory can lead to effective theory of gravity. Thirdly in Cardy generalizes the boundary state concept so as to get an exact description of the boundary conformal field theory. Now the boundary conformal field theory of Gepner model, which is important in the heterotic string theory, has been proposed in . This generalizes D-brane concept, since the original D-brane only comes from the system of Type I string, Type IIA string and Type IIB string. Fourthly, in , the path integral representation of boundary state provides a simple proof supporting the fundamental argument in Matrix theory that physics of $`D_p`$ brane is equivalent to that of infinitely many $`D_{(p2r)}`$ branes. Fifthly, the path integral formalism has been incorporated into recent development of Noncommutative geometry in string theory in . Finally, in the classical $`p`$-brane solution can be constructed from one point amplitude of boundary state with closed string states. This is also an important proof that D-brane is the classical soliton solution with R-R charge. Here we only give some examples of the path integral representation of the conformal invariant boundary state <sup>1</sup><sup>1</sup>1 The Feynman path integral is one way to represent a quantum theory, and it is a very natural method for describing interactions in string theory. Since we have had a systematic study of string theory using the Polyakov path integral, we may try to explore the path integral representation of D-brane.. Since the boundary conformal field theory is a candidate for D-brane theory, we must find solutions of the boundary conformal field theory. In the case of path integral representation, algebraic constraint equation can be transformed into differential equation which is easily solved to obtain the solution represented by string modes. Although the conformal invariant boundary state satisfying Cardy condition can be constructed with Ishibashi state, this formalism is not often used in the practical calculation. For example in those review articles the mode representation of the boundary state is used calculating various amplitudes <sup>2</sup><sup>2</sup>2 It is necessary to notice the path integral representation of boundary state is easily connected with divergence. Although the path integral representation of the boundary state in non-constant $`U(1)`$ gauge field background in can be obtained, those boundary states contain bad divergences in and so need regularization and renormalization. Another example is the path integral representation of the boundary state in the linear dilaton background whose procedure of regularization and renormalization will be discussed in Appendix B. . The organization of this paper is as follows. In section 2, only from conformal invariance we deduce the path integral representation of boundary state in constant $`U(1)`$ gauge field background <sup>3</sup><sup>3</sup>3 Although the path integral representation of boundary state in constant $`U(1)`$ gauge field background in had been argued so as to determine the normalization factor, here the familiar form of path integral representation in can be deduced directly from the conformal invariance.. and show some applications of this formalism in T-duality and Matrix theory. In section 3, we construct the conformal invariant boundary state in the linear dilaton background with the path integral representation, and at the same time construct the path integral representation of the boundary state in open string tachyon background. Section 4 is devoted to discussions. ## 2 The path integral representation of the boundary state in the $`U(1)`$ gauge field background The boundary state must be conformal invariant, which is the requirement of the reparametrization invariance on the world sheet. So the conformal invariance can be directly used to construct the path integral representation of boundary state in the general background. In this section, firstly this procedure will be given in detail. Secondly T-duality of D-brane to construct $`D_p`$ brane from $`D_1`$ brane can be carried out in the path integral formalism of the boundary state. Then the simplest solution of those differential equations is found right to be the path integral representation of boundary state in constant $`U(1)`$ gauge field background. ### 2.1 The path integral representation of the boundary state In this paper only consider the bosonic string case. In fact the generalization to supersymmetrical case is straightforward. The holomorphic stress tensor reads $$T=\frac{1}{2}:(X)^2:,$$ (1) and a similar anti-holomorphic counterpart. It is necessary to work with mode expansions $`X(\tau ,\sigma )`$ $`=`$ $`q+(\alpha _0+\stackrel{~}{\alpha }_0)\tau +i{\displaystyle \underset{n0}{}}{\displaystyle \frac{1}{n}}\left(\alpha _ne^{in(\tau \sigma )}+\stackrel{~}{\alpha }_ne^{in(\tau +\sigma )}\right),`$ (2) $`X(z,\overline{z})`$ $`=`$ $`qi(\alpha _0\text{ln}z+\stackrel{~}{\alpha }_0\text{ln}\overline{z})i{\displaystyle \underset{n0}{}}{\displaystyle \frac{1}{n}}\left(\alpha _nz^n+\stackrel{~}{\alpha }_n\overline{z}^n\right),`$ (3) $`L_n`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{m}{}}:\alpha _{m+n}\alpha _m:`$ (4) in which we have taken the Regge slope $`\alpha ^{}=2`$. The commutators are $`[\alpha _m,\alpha _n]=m\delta _{m+n,0}`$, and similarly for the right-moving modes. Let $`K_n=L_n\stackrel{~}{L}_n`$. The boundary condition is entirely encoded in the boundary state $`\|B`$, and the conformal invariant condition is $`K_n\|B=0`$. To solve equations $`K_n\|B,p=0`$, it is essential to adopt the coherent state technique introduced in . Introduce the following coherent state which satisfies $$(\alpha _n\stackrel{~}{\alpha }_nx_n)\|x,p=0,$$ (5) where $`n`$ can be either positive or negative. Further the conjugate coherent state should satisfy $$x,p\|(\alpha _n\stackrel{~}{\alpha }_nx_n)=0.$$ (6) If requiring $`x,p\|=(\|x,p)^{}`$, then $`x_n=(x_n)^{}`$. Solve this equation (5), the correspondent solution is $$\|x,p=\mathrm{exp}\left(\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n}[\frac{1}{2}x_nx_n+\alpha _n\stackrel{~}{\alpha }_n+x_n\alpha _nx_n\stackrel{~}{\alpha }_n]\right)\|p.$$ (7) This set of states form a complete orthogonal basis as can be checked. The formulas $`e^{A+B}=e^Ae^Be^{\frac{1}{2}[A,B]}`$, and $`e^Ae^B=e^Be^Ae^{[A,B]}`$ when $`[[A,B],A]=[[A,B],B]=0`$ are very helpful, and will be often used in the involved calculation in this paper. Observing this coherent state, we can find each oscillation mode can be traced in this state! Basing on this point, our goal may be realized that transforming the algebraic equations into the differential equations. So the detailed operation is in essence a kind of substitution $`\alpha _n\|x,p=\left(n{\displaystyle \frac{}{x_n}}+{\displaystyle \frac{1}{2}}x_n\right)\|x,p,`$ $`\stackrel{~}{\alpha }_n\|x,p=\left(n{\displaystyle \frac{}{x_n}}{\displaystyle \frac{1}{2}}x_n\right)\|x,p,`$ with the requirement of $`n0`$. We postulate that the boundary state with the boundary interaction is of the form $$\|B,p=[dx]\|x,p\mathrm{exp}\left(S(x)\right),$$ (9) in which $`S(x)`$ represents the special boundary interaction. In fact the general form of the boundary interaction should admit string zero mode $`\widehat{q}`$ or differential operator$`\frac{}{x}`$. To solve $`K_n\|B,p=0`$, one first computes $`K_n\|x,p`$ $`=`$ $`\left(px_n+{\displaystyle \underset{m0,n}{\overset{\mathrm{}}{}}}mx_{m+n}{\displaystyle \frac{}{x_m}}\right)\|x,p=\overline{p},`$ (10) in which $`n0`$ is necessary. Substituting the above relations into $`K_n\|B,p=0`$ and integrating by parts, the conformal invariance condition for $`S(x)`$ is transformed into $$\left(px_n+\underset{m0,n}{\overset{\mathrm{}}{}}mx_{m+n}\frac{}{x_m}\right)S(x)=0.$$ (11) However for $`n=0`$ case the thing is not similar. Applying the above procedure to $$K_0\|B,p=0,$$ (12) we obtain the equation $$\underset{m0}{\overset{\mathrm{}}{}}mx_m\frac{}{x_m}S(x)=0.$$ (13) To our special ansatz of the boundary interaction and the special choice of vacuum state, the solution of the equation (13) and the equation (11) corresponds to a special conformal invariant boundary state. In order to obtain one boundary state, we may take various ansatz of vacuum state and boundary interaction. For example, if we take vacuum state $`\|p=\overline{p}`$, then in the $`n0`$ case the equation is the form, $$\left(\underset{m0}{\overset{\mathrm{}}{}}mx_{m+n}\frac{}{x_m}\right)S(x)=0.$$ (14) ### 2.2 T-duality of D-brane Since T-duality interchanges Neumann and Dirichlet boundary conditions, a further T-duality in a direction tanget to a Dp-brane reduces it to a $`D_{(p1)}`$ brane, while a T-duality in an orthogonal direction turns it into a $`D_{(p+1)}`$ brane. In the path integral representation of the boundary state of D-brane, T-duality of $`D_p`$ brane can be carried out. The coherent state which satisfies $$(\alpha _n\stackrel{~}{\alpha }_nx_n)\|x,p=\overline{p}=0$$ (15) has to satisfy the following constraint $$(\alpha _n+\stackrel{~}{\alpha }_n)\|x,p=\overline{p}=2n\frac{}{x_n}\|x,p=\overline{p}.$$ (16) Therefore the boundary state $$\|B_N=[dx]\|x,p=\overline{p}$$ (17) satisfies Neumann condition $`_\tau X(\sigma )\|B_N=0`$. By carrying out the path integral, its detailed formalism is $$\|B_N=A\mathrm{exp}\left(\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n}[\alpha _n\stackrel{~}{\alpha }_n]\right)\|p=\overline{p},$$ (18) where A is a proportional constant. Further, another coherent state has to be also defined in order to realize our goal in this subsection. Firstly, some formulas should be given by $`X(\sigma )`$ $`=`$ $`q+i\sqrt{{\displaystyle \frac{\alpha ^{}}{2}}}{\displaystyle \underset{n0}{}}{\displaystyle \frac{1}{n}}\left(\alpha _n\stackrel{~}{\alpha }_n\right)e^{in\sigma },`$ (19) $`_\tau X(\sigma )`$ $`=`$ $`\alpha ^{}P+\sqrt{{\displaystyle \frac{\alpha ^{}}{2}}}{\displaystyle \underset{n0}{}}\left(\alpha _n+\stackrel{~}{\alpha }_n\right)e^{in\sigma },`$ (20) where $`P`$ is the total string momentum. Their commutator is $`[X(\sigma ),P(\sigma ^{})]=i\delta (\sigma \sigma ^{})`$, in which the string momentum $`P(\sigma )`$ at the point labeled by the $`\sigma `$ is defined as $`\frac{1}{2\pi \alpha ^{}}_\tau X(\sigma )`$. The Fourier expansion of $`\delta (\sigma \sigma ^{})`$ takes the form $`\frac{1}{2\pi }_m\mathrm{exp}(im(\sigma \sigma ^{}))`$. Now define a new coherent state $$X(\sigma )\|x=x(\sigma )\|x.$$ (21) The difference between the coherent state (5) and (21) only lies in that the latter is the position eigenstate $`\widehat{q}\|x=q\|x`$ and the former is the momentum eigenstate $`\widehat{P}\|x,p=P\|x,p`$. Naturally, they may be connected by the transform $$Dx(\sigma )\|x=[dx]\|x,p=\overline{p}.$$ (22) And the coherent state (21) can be constructed with the boundary state of instanton namely $`D_1`$ brane, $$\|x=\mathrm{exp}\left(i𝑑\sigma P(\sigma )x(\sigma )\right)\|B_1,$$ (23) where the boundary state of instanton satisfies the constraint $`X(\sigma )\|B_1=0`$. We conclude this subsection by constructing $`D_p`$ brane boundary state with $`D_1`$ brane boundary state $$\|B_p=\underset{i=0}{\overset{p}{}}Dx_i(\sigma )\mathrm{exp}\left(i𝑑\sigma P^i(\sigma )x_i(\sigma )\right)\|B_1$$ (24) with $`i`$ is limited to $`i=0,1,\mathrm{},p`$ <sup>4</sup><sup>4</sup>4 The idea used in where the path integral representation was used to support open string T-duality has been generalized. Here as the result of a T-duality in an orthogonal direction the boundary state of $`D_p`$ brane can be constructed from the boundary state of $`D_1`$ brane. In (24) a further T-duality in a direction tanget to a Dp-brane is equivalent to changing the measure $`_{i=0}^pDx_i(\sigma )`$ to $`_{i=0}^{p1}Dx_i(\sigma )`$. . It can be checked such $`D_p`$ brane boundary state will satisfy the boundary conditions, $`X^\mu (\sigma )\|B_p=0,\mu =p+1,\mathrm{},D1;`$ (25) $`P^i(\sigma )\|B_p=0,i=0,1,\mathrm{},p.`$ (26) On the other hand only from the above expression of (24), T-duality of boundary state seems to be Fourier transform in the configuration space involved with T-duality. ### 2.3 The path integral representation of the boundary state in the $`U(1)`$ gauge field background We now turn to solve the differential equation (13) and the differential equation (14). It is easy to find one simple solution $$S(x)=\frac{1}{4}F_{\mu \nu }\underset{m0}{}\frac{x_{m}^{}{}_{}{}^{\mu }x_{m}^{}{}_{}{}^{\nu }}{m}$$ (27) where $`x_0=0`$ and $`F_{\mu \nu }=F_{\nu \mu }`$ are necessary and the factor $`\frac{1}{4}`$ has been fixed due to the following calculation of the normalization factor. The correspondent boundary state is the form $$\|B,p=0=[dx]\mathrm{exp}\left(\frac{1}{4}F_{\mu \nu }\underset{m0}{}\frac{x_{m}^{}{}_{}{}^{\mu }x_{m}^{}{}_{}{}^{\nu }}{m}\right)\|x,p=\overline{p}=0.$$ (28) Firstly this above solution <sup>5</sup><sup>5</sup>5 It is nontrivial to point out that in the presence of a constant $`U(1)`$ gauge field background Virasoro generators are not modified. can be integrated out to show the final result is right the boundary state in the constant $`U(1)`$ gauge field background in . The final result of (28) in the Euclidean spacetime is $$\|B,p=0=N(F)\mathrm{exp}\left(\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n}\alpha _n\left(\frac{1F}{1+F}\right)\stackrel{~}{\alpha }_n\right)\|p=\overline{p}=0,$$ (29) with the normalization factor $`N(F)=\left(det(1+F)\right)^{\frac{1}{2}}`$. In bosonic string case the normalization factor is just the effective action of gauge potential in <sup>6</sup><sup>6</sup>6We would like to thank Professor A.A.Tseytlin for bringing our attention to the paper. In fact, our purpose which is different from of is to give an interesting example to show potential value of path integral representation. The familiar boundary state in $`U(1)`$ constant gauge background can be only determined by conformal invariance. The boundary state can be used to determine boundary condition from which the action can be constructed. Naturally this example also shows the effective action or the normalization factor can be determined directly from the conformal invariance.. The purpose of arguing the path integral representation in is to obtain the normalization factor. Here shows this approach to get normalization factor is rather natural in contrast with the other methods. The method in is to compare the one point amplitude of the boundary state with that from Dirac-Born-Infeld action. In addition the detailed form of the normalization factor in is in the requirement of gauge invariance of D-brane source term contained in the action of closed string field theory. With the form (19) and the following $$_\sigma X(\sigma )=\underset{n0}{}\left(\alpha _n\stackrel{~}{\alpha }_n\right)e^{in\sigma },$$ (30) the path integral representation of (28) may be written with another form $$\|B,p=0=[dx]\mathrm{exp}\left(\frac{i}{8\pi }F_{\mu \nu }_{0}^{}{}_{}{}^{2\pi }𝑑\sigma x^\mu (\sigma )_\sigma x^\nu \right)\|x,p=\overline{p}=0$$ (31) which shows the boundary interaction of $`U(1)`$ constant gauge field background. With the help of (23), the above boundary state can be changed into the form $$\|B,p=0=Dx(\sigma )\mathrm{exp}\left(\frac{i}{8\pi }F_{\mu \nu }_{0}^{}{}_{}{}^{2\pi }𝑑\sigma x^\mu (\sigma )_\sigma x^\nu i_{0}^{}{}_{}{}^{2\pi }𝑑\sigma P_\mu x^\mu \right)\|B_1$$ (32) from which the boundary condition determining such above boundary state may be found. Indeed, one can show that the following identity holds: $`0`$ $`=`$ $`{\displaystyle Dx(\sigma )\frac{\delta }{\delta x^\mu (\sigma )}\mathrm{exp}\left(\frac{i}{8\pi }F_{\mu \nu }_{0}^{}{}_{}{}^{2\pi }𝑑\sigma x^\mu (\sigma )_\sigma x^\nu i_{0}^{}{}_{}{}^{2\pi }𝑑\sigma P_\mu x^\mu \right)\|B_1}`$ (33) $`=`$ $`[{\displaystyle \frac{i}{4\pi }}F_{\mu \nu }_\sigma X^\nu iP_\mu (\sigma )]\|B,p=0,`$ where the boundary condition $$\left(_\tau X_\mu F_{\mu \nu }_\sigma X^\nu \right)\|B,p=0=0$$ (34) can be extracted. Here it is both the boundary condition and the conformal invariance which require $$p=\overline{p}=0.$$ (35) Such the boundary condition just corresponds to the open string theory $$S=𝑑\tau _0^{2\pi }𝑑\sigma \left\{\frac{1}{4\pi \alpha ^{}}[(_\tau X)^2(_\sigma X)^2][\delta (\sigma )\delta (\sigma 2\pi )]\dot{X}^\mu A_\mu (X)\right\}.$$ (36) The following is a simple review to where the path integral representation of the boundary state in the $`U(1)`$ constant gauge field background can be also argued from the point of Matrix theory. Notice that their conventions are a little different with ours since the definition of $`X(\tau ,\sigma )`$ and two order antisymmetry tensor $`F_{\mu \nu }`$ are not completely fixed. The configuration of infinitely many D-instantons can be expressed by $`\mathrm{}\times \mathrm{}`$ hermitian matrices $`X^M(M=0,\mathrm{},D1)`$. The one we consider is $`X^i`$ $`=`$ $`\widehat{Q}^i,(i=0,\mathrm{},p)`$ $`X^M`$ $`=`$ $`0(M=p+1,\mathrm{},D1),`$ (37) where $`\widehat{Q}^i(i=0\mathrm{},p)`$ satisfy $$[\widehat{Q}^i,\widehat{Q}^j]=i\theta ^{ij}.$$ (38) Here $`(p+1)\times (p+1)`$ matrix $`\theta `$ is assumed to be invertible. In Matrix theory this configuration of D-instantons is equivalent to a D$`p`$-brane. A quick way to see the equivalence is to look at the boundary states. The boundary state $`\|B`$ corresponding to the configuration eq. (37) can be written as follows: $$\|B=\text{TrX}e^{i_0^{2\pi }𝑑\sigma P_i\widehat{X}^i}\|B_1.$$ (39) $`\|B_1`$ includes also ghost part which is not relevant to the discussion here. The factor in front of $`\|B_1`$ is an analogue of Wilson loop and corresponds to the background eq. (37). Eq. (39) can be rewritten with the path integral as $$\|B,p=0=Dx(\sigma )\mathrm{exp}\left(\frac{i}{2}𝑑\sigma x^i_\sigma x^jF_{ij}i𝑑\sigma P_ix^i\right)\|B_1,$$ (40) where $`F_{ij}=(\theta ^1)_{ij}`$. It is straightforward to perform the path integral in eq. (40). In conclusion $`\|B`$ is equivalent to the boundary state for a D$`p`$-brane with $`U(1)`$ gauge field strength $`F_{ij}`$ on the worldvolume. ### 2.4 Some notes of the path integral representation of the boundary state To construct a good theory of $`D_p`$ brane in the general background, we can argue $$\|B_p^N=\underset{i=0}{\overset{p}{}}Dx_i(\sigma )\mathrm{exp}\left(S(\widehat{q},x_n,\frac{}{x_n},\mathrm{})\right)\|x(\sigma ),$$ (41) where $`\|B_p^N`$ is the Neumann part of the boundary state $`\|B_p`$ which is defined as $`\|B_p`$ $`=`$ $`\|B_p^N\|B_p^D`$ (42) with $`\|B_p^D=_{\mu =p+1}^{D1}\|B_1^\mu `$. In essence (41) may be seen as Fourier transform between the boundary state $`\|B_p^N`$ with boundary interaction and the coherent state $`\|x(\sigma )`$ <sup>7</sup><sup>7</sup>7From the above expression, it can be argued that physics of the world volume of $`D_p`$ brane, which is represented by $`\|B_p^N`$, may be equivalent to the complete physics of $`(p+1)`$ particles in space time, which is represented by quantum state $`_{i=0}^p\|x(\sigma )^i`$. The involved motion of these quantum states is with the weight $`\left(S(\widehat{q},x_n,\frac{}{x_n},\mathrm{})\right)`$ which is determined by the respective paths.. For the path integral representation of the boundary state $$\|B,p=[dx]\|x,p=\overline{p}\mathrm{exp}\left(S(x)\right)$$ (43) we can check those equations $`\left({\displaystyle \underset{m0}{\overset{\mathrm{}}{}}}mx_{m+n}{\displaystyle \frac{}{x_m}}\right)S(x)`$ $`=`$ $`0,n0,`$ $`{\displaystyle \underset{m0}{\overset{\mathrm{}}{}}}mx_m{\displaystyle \frac{}{x_m}}S(x)`$ $`=`$ $`0,`$ $`x_0`$ $`=`$ $`0`$ (44) from the conformal invariance are consistent with the following equations $`2n{\displaystyle \frac{}{x_n^\mu }}S(x)`$ $`=`$ $`F_{\mu \nu }x_n^\nu ,n0,`$ $`P`$ $`=`$ $`0`$ (45) from the mixed boundary condition determining the boundary state in the constant $`U(1)`$ gauge field background. Therefore when we take the ansatz $`S(x)`$, we will know such path integral representation shows the physics of $`D`$ brane in constant $`U(1)`$ gauge field background before the detailed path integral calculation in Appendix A. In fact it is important to firstly estimate the form of $`S`$ with some physical consideration before determining the detailed form of $`S`$ by the conformal invariance. In this paper some ansatz can be changed. Define the coherent state $$(\alpha _n+\stackrel{~}{\alpha }_np_n)\|pp=0$$ (46) then the boundary state with the boundary interaction can be expanded into the sum of the above coherent state, $$\|B,P=[dp]\mathrm{exp}\left(S(\widehat{q},p_n,\frac{}{p_n},\mathrm{})\right)\|pp=0.$$ (47) So all the calculation in this paper can be carried out from the new starting point. ## 3 The path integral representation of the boundary state in other special backgrounds Besides the $`U(1)`$ gauge field background, the linear dilaton background and the tachyon background are important in recent research. Especially, it seems attracting how to construct the boundary state of D-brane in Type 0 string under the closed string tachyon background and make it consistent with the argument that D-brane is the classical soliton solution with R-R charge. ### 3.1 The path integral representation of the boundary state in the linear dilaton background We will construct the conformal invariant boundary state in the linear dilaton background with path integral representation <sup>8</sup><sup>8</sup>8 Professor Miao Li introduced me his paper and advised me to read this paper carefully. In the conformal invariance was directly used to construct the path integral representation of boundary state in the linear dilaton background. However the path integral representation of the boundary state in can be found not to be exactly conformal invariant. In the following the conformal invariant boundary state will be constructed.. Now the starting point is still the stress tensor $$T=\frac{1}{2}(\varphi )^2+Q^2\varphi $$ (48) and a similar anti-holomorphic counterpart. The central charge of this free scalar is $`c=1+12Q^2`$, and $`Q=\sqrt{2}`$ in two dimensional string theory. Consider a unit disk, the conformal invariance condition on the boundary means no net energy-momentum flow out of the boundary. It is convenient to work with mode expansions $`\varphi `$ $`=`$ $`\phi _0i(p\text{ln}z+\overline{p}\text{ln}\overline{z})i{\displaystyle \underset{n0}{}}{\displaystyle \frac{1}{n}}\left(\alpha _nz^n+\stackrel{~}{\alpha }_n\overline{z}^n\right),`$ (49) $`L_n`$ $`=`$ $`[\alpha _0+iQ(n+1)]\alpha _n+{\displaystyle \frac{1}{2}}{\displaystyle \underset{m0,n}{}}\alpha _{m+n}\alpha _m,n0,`$ (50) $`L_0`$ $`=`$ $`[{\displaystyle \frac{\alpha _0}{2}}+iQ]\alpha _0+{\displaystyle \frac{1}{2}}{\displaystyle \underset{m0}{}}\alpha _m\alpha _m.`$ (51) The similar formula for $`\stackrel{~}{L}_n`$. The commutators are $`[\alpha _m,\alpha _n]=m\delta _{m+n,0}`$, and similarly for the right-moving modes. Let $`K_n=L_n\stackrel{~}{L}_n`$. The boundary condition is entirely encoded in the boundary state $`\|B`$, and the conformal invariance condition is $`K_n\|B=0`$. The usual Neumann boundary condition is given by $`_r\varphi =0`$ on the boundary of the unit disk. In terms of the boundary state, it states that $$P\|B_N=(\alpha _n+\stackrel{~}{\alpha }_n)\|B_N=0.$$ Due to the existence of the background charge $`Q`$, one has to modify the boundary condition a bit: $`p=iQ`$. So there must be a net momentum flow out of the boundary (in view of spacetime $`\varphi `$). One way to see this is to consider the commutators $$[K_m,\alpha _n+\stackrel{~}{\alpha }_n]=2iQm\delta _{m+n,0}n\left(\alpha _{m+n}+\stackrel{~}{\alpha }_{mn}\right).$$ (52) When the case $`n+m=0`$ occurs, the above commutator can be showed in the clearer form $$[K_m,\alpha _n+\stackrel{~}{\alpha }_n]=m(2iQ+\alpha _0+\stackrel{~}{\alpha }_0).$$ (53) Actually, we have taken the following conventions, $$2p=\alpha _0+\stackrel{~}{\alpha }_0,\alpha _0=\stackrel{~}{\alpha }_0=p.$$ (54) So when $`p=iQ`$, the center term disappears, and it is possible to impose both the conformal invariance condition and Neumann boundary condition. To be as close to the ordinary Dirichlet condition as possible, one requires that a net momentum transfer is possible if one scatters string states against the object described by the boundary state. So $`\|B,p`$ is an eigen-state of $`p`$ with arbitrary number $`p`$. To solve equations $`K_n\|B,p=0`$, we may construct the path integral representation of the conformal invariant boundary state. However the ansatz $`S(x,\widehat{q},\widehat{\stackrel{~}{q}})`$ is one little special form with the contribution of zero mode, $$S(x,\widehat{q},\widehat{\stackrel{~}{q}})=V_{coup}^n\mathrm{exp}\left(\frac{\widehat{q}+\widehat{\stackrel{~}{q}}}{2Q}\right)\mathrm{\Phi }(x),$$ (55) where the operator $`\widehat{q}`$ and $`\widehat{\stackrel{~}{q}}`$ satisfy the commutative relation $`[\widehat{q},\widehat{\stackrel{~}{q}}]=0`$, $`[\alpha _0,\widehat{q}]=i`$ and $`[\stackrel{~}{\alpha }_0,\widehat{\stackrel{~}{q}}]=i`$. And $`V_{coup}^n`$ is the coupling constant before renormalization procedure. The boundary state has been changed into the form $$\|B=[dx]\mathrm{exp}\left(V_{coup}^n\mathrm{exp}\left(\frac{\widehat{q}+\widehat{\stackrel{~}{q}}}{2Q}\right)\mathrm{\Phi }(x)\right)\|x,p=\overline{p}=iQ.$$ (56) So the equation $`K_n\|B,p=0`$ now can be changed into the following two equations, $`{\displaystyle \frac{i}{2Q}}x_n\mathrm{\Phi }(x)+2iQn^2{\displaystyle \frac{}{x_n}}\mathrm{\Phi }(x)`$ $`=`$ $`{\displaystyle \underset{m0,n}{\overset{\mathrm{}}{}}}mx_{m+n}{\displaystyle \frac{}{x_m}}\mathrm{\Phi }(x),`$ (57) $`{\displaystyle \underset{m0}{\overset{\mathrm{}}{}}}mx_m{\displaystyle \frac{}{x_m}}\mathrm{\Phi }(x)`$ $`=`$ $`0.`$ (58) The special solution of the two equations could be found as follow, $$\mathrm{\Phi }(x)=\frac{dz}{z}\mathrm{exp}\left(\underset{m0}{}\frac{ix_m}{2Qmz^m}\right).$$ (59) As we have claimed, such above solution contains the divergence. Naturally the renormalized boundary state which will be given in Appendix B is conformal invariant only in some specific cases. Actually the conformal invariant boundary state in the linear dilaton background can be constructed without considering the path integral representation. Define the screening charge $$QQ=𝑑z:\mathrm{exp}\left(ikX(z)\right):$$ (60) with $`k`$ must satisfy the following equation $$k^2+2iQk=2.$$ (61) Since the commutative relation $`[L_n,QQ]=0`$, the conformal invariant boundary state <sup>9</sup><sup>9</sup>9This conformal invariant boundary state is useful in Appendix B. is $$\|B=Fun(QQ)\|B,p=\overline{p}=iQ_N$$ (62) in which $`Fun`$ is function of $`QQ`$. The boundary state $`\|B,p=\overline{p}=iQ_N`$ is the Neumanm boundary state with the momemtum which satisfies $`(\alpha _n+\stackrel{~}{\alpha }_n)\|B,p=\overline{p}=iQ_N`$ $`=0,n0,`$ $`(\alpha _0+\stackrel{~}{\alpha }_0)\|B,p=\overline{p}=iQ_N`$ $`=2p\|B,p=\overline{p}=iQ_N,`$ $`(L_n\stackrel{~}{L}_n)\|B,p=\overline{p}=iQ_N`$ $`=0.`$ (63) Here the boundary state (62)is still rather specifically. In principle the number of conformal invariant boundary state in the linear dilaton background is infinitely many. ### 3.2 The path integral representation of the boundary state in the open string tachyon background The path integral representation of the boundary state in the open string tachyon background will be constructed, which provides some clues to construct the conformal invariant boundary state of D-brane from Type $`0`$ string in the tachyon background <sup>10</sup><sup>10</sup>10 It is necessary to point out that tachyon appearing in Type 0 string theory is closed string state but here tachyon is open string state. Since the boundary state is the closed string source, it is possible to construct the boundary state of D-brane in Type 0 string. On the other hand in order to support the argument the path integral representation of boundary state can provide an exact description of D-brane in general background, we have to give a nontrivial example. So the work to ensure the boundary state representing the boundary state of D-brane in Type 0 string is just nontrivial. . The conformal invariant boundary state in the open string tachyon background <sup>11</sup><sup>11</sup>11 The boundary state in the open string tachyon background has been researched in . is the form $$\|B=\mathrm{exp}(dz:\mathrm{exp}(ikX(z)):)\|B,p=\overline{p}_N,$$ (64) with the convention $`k^2=2`$. So the path integral representation of the boundary state in the open string tachyon background will have one form $`\|B={\displaystyle [dx]\|x,p=\overline{p}}+`$ (65) $`{\displaystyle [dx]\underset{l=1}{\overset{\mathrm{}}{}}\mathrm{\Phi }_l\|x,p=\overline{p}}.`$ (66) The differential equations from the conformal invariance constraint are $$\mathrm{\Phi }_l\left(nl\frac{}{x_n}+\frac{1}{2}lx_n\right)+\mathrm{\Phi }_l\left(\underset{m0}{}mx_{m+n}\frac{}{x_m}\right)=0,n0,$$ (67) and for $`n=0`$ case the differential equation is $$\mathrm{\Phi }_l\left(\frac{1}{2}l^2k^2+lkp\right)+\mathrm{\Phi }_l\left(\underset{m0}{}mx_m\frac{}{x_m}\right)=0.$$ (68) Finally the ansatz of $`\mathrm{\Phi }_l`$, namely $`\mathrm{\Phi }_l(\widehat{q},p,x_n,x_n,\frac{}{x_n})`$ is solved as $`\mathrm{\Phi }_l={\displaystyle \underset{i=1}{\overset{l}{}}}{\displaystyle 𝑑z_i\mathrm{exp}^{ik\widehat{q}}z_i^{kp}z_i^{(i1)k^2}}`$ $`\mathrm{exp}\left({\displaystyle \underset{n=1}{}}{\displaystyle \frac{kx_n}{2nz_i^n}}\right)\mathrm{exp}\left({\displaystyle \underset{n0}{}}kz_i^n{\displaystyle \frac{}{x_n}}\right)\mathrm{exp}\left({\displaystyle \underset{n=1}{}}{\displaystyle \frac{kx_nz_i^n}{2n}}\right),`$ (69) with the definition of $`_{i=1}^l[dz_i][dz_1][dz_2]\mathrm{}[dz_l]`$. In essence this solution is relevant with the form of (64). The factor $`z_i^{(i1)k^2}`$ appearing in the (69) is from the noncommutative relation between zero modes, for example $`l=2`$ case, $`{\displaystyle 𝑑z_2\mathrm{exp}^{ik\widehat{q}}z_{2}^{}{}_{}{}^{k\widehat{p}}𝑑z_1\mathrm{exp}^{ik\widehat{q}}z_{1}^{}{}_{}{}^{kp}\|x,p=\overline{p}}`$ (70) $`={\displaystyle 𝑑z_1\mathrm{exp}^{ik\widehat{q}}z_{1}^{}{}_{}{}^{kp}𝑑z_2\mathrm{exp}^{ik\widehat{q}}z_{2}^{}{}_{}{}^{kp}z_{2}^{}{}_{}{}^{k^2}\|x,p=\overline{p}}.`$ (71) Although such above calculation is not easy, it shows the path integral representation of the boundary state is very useful if we can take some tricks in the given background. ## 4 Discussions In this paper some examples of the path integral representation of the boundary state are given in some special backgrounds such as the $`U(1)`$ gauge field background, the linear dilaton background and the open string tachyon background. The path integral representation of the boundary state in the $`U(1)`$ gauge field background contains much essential information, especially the normalization factor which is the effective action of gauge field in the bosonic string theory. The application in Matrix theory and Noncommutative geometry hints that Path integral representation of the boundary state could be rather fundamental in describing D-brane in the general background. The initial purpose of this paper is to construct a general solution of the boundary conformal field theory with analytical approach, mainly for the constraint equations $`(L_n\stackrel{~}{L}_n)\|B=0`$ are difficult to be solved to obtain the solution represented by string modes from the pure algebraic approach. However in the path integral representation it is easy transforming those algebraic equations into the differential formalism which can be solved. In addition, Cardy condition which ensures the existence of open string theory in the boundary conformal field theory is vital to this view point that the boundary conformal field theory is an exact description of D-brane in the general background. We will consider Cardy condition realization in the path integral representation in the future work. Finally, it is also our wish that the path integral representation of boundary state should be used to support or interpret the known ansatz in <sup>12</sup><sup>12</sup>12 It has been recently argued that holographic principle should be deduced from noncommutative geometry, such as in . And it seems that two fundamental principles are intrinsic consistent in . In fact they can be both argued from the point of string/M theory. So it is possible to deduce AdS/CFT duality from string/M theory. Naturally we wish that this procedure may be depend on the path integral representation of boundary state. From string amplitude with D brane, which may be represented by saturating string state with boundary state in, much useful information about supergravity can be obtained in. However those cases are in flat spacetime, but now we have to face curved spacetime such as AdS. We wish we could continue this research because this will provide a truely nontrivial example of path integral representation of boundary state of D-brane. , $$\mathrm{exp}_{S^d}\varphi _0O_{CFT}=Z_S(\varphi _0).$$ (72) It is possible to supersymmetrize all the results in this paper. There are still a lot of work to be done. The path integral representation of constant gauge field configuration should be generalized to nonconstant background in the conformal invariant requirement. The path integral representation of the boundary state in the linear dilaton background might be useful to construct the boundary state under $`AdS_3`$ background, especially in the light cone gauge in . Most important, the construction of the boundary state in tachyon background is potentially useful for the generalization of $`AdS/CFT`$ in non-super symmetrical case. The recent research on D-brane in Type 0 string shows we must face the tachyon problem in arriving at the purpose. However, it seems that the usual boundary state is not reasonable which is constructed from the general procedure without considering the tachyon background effect, because the classical p-brane from the usual boundary state is not the solution in . This can be the result of the introduction of the tachyon background, which makes the equations corresponding to p-brane solution become nonlinear. But the usual boundary state is only the linear combination of closed string states. Since nonsupersymmetrical generalization of $`AdS/CFT`$ in Type $`0`$ string case seems more reasonable than in that approach of supersymmetrical breaking case used in , it is possible to construct the boundary state with the tachyon background in Type $`0`$ string case. ###### Acknowledgments. We would like to thank Miao Li for helpful comments on the manuscript and Yi-Hong Gao for a helpful discussion. ## Appendix A The calculation of the normalization factor $`N(F)`$ In the following the calculation of the normalization factor $`N(F)`$ is given. The special solution of the differential equation (13) and the differential equation (14) is $`S(x)`$ $`=`$ $`{\displaystyle \frac{1}{4}}F_{\mu \nu }{\displaystyle \underset{m0}{}}{\displaystyle \frac{x_{m}^{}{}_{}{}^{\mu }x_{m}^{}{}_{}{}^{\nu }}{m}},`$ (73) $`=`$ $`{\displaystyle \frac{1}{2}}F_{\mu \nu }{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{x_{m}^{}{}_{}{}^{\mu }x_{m}^{}{}_{}{}^{\nu }}{m}}.`$ The entire calculation of the path integral of (28) has to be put in the Euclidean spacetime. The calculation of the (28) is as follow, $`\|B,p=0`$ $`=`$ $`{\displaystyle [dx]\mathrm{exp}\left(\frac{1}{4}F_{\mu \nu }\underset{m0}{}\frac{x_{m}^{}{}_{}{}^{\mu }x_{m}^{}{}_{}{}^{\nu }}{m}\right)\|x,p=\overline{p}=0},`$ (74) $`=`$ $`{\displaystyle [dx]\mathrm{exp}\left(\frac{1}{2}\underset{m=1}{\overset{\mathrm{}}{}}\frac{1}{m}(x_m2\alpha _m\frac{1}{1+F})(1+F)(x_m+2\frac{1}{1+F}\stackrel{~}{\alpha }_m)\right)}`$ $`\mathrm{exp}\left({\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m}}\alpha _m\left({\displaystyle \frac{1F}{1+F}}\right)\stackrel{~}{\alpha }_m\right)\|p=\overline{p}=0.`$ The integral measure in the above formula has to be chosen as $`{\displaystyle [dx]}`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}\frac{dx_ndx_n}{n}},`$ (75) $`{\displaystyle \frac{dx_ndx_n}{n}}`$ $`{\displaystyle \frac{dadb}{(2\pi )^D}},`$ (76) with the definition of $`x_n=\sqrt{n}(a+ib)`$ and $`x_n=\sqrt{n}(aib)`$. The choice of (76) aims at using the following formula in our calculation, $$\frac{dadb}{(2\pi )^D}\mathrm{exp}\left(\frac{1}{2}(aib)A(a+ib)\right)=\frac{1}{A^D}.$$ (77) So the final result can be given, $`\|B,p=0`$ $`=`$ $`{\displaystyle [dx]\mathrm{exp}\left(\frac{1}{4}F_{\mu \nu }\underset{m0}{}\frac{x_{m}^{}{}_{}{}^{\mu }x_{m}^{}{}_{}{}^{\nu }}{m}\right)\|x,p=\overline{p}=0},`$ (78) $`=`$ $`\left(Det(1+F)\right)^1\mathrm{exp}\left({\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{m}}\alpha _m\left({\displaystyle \frac{1F}{1+F}}\right)\stackrel{~}{\alpha }_m\right)`$ $`\|p=\overline{p}=0.`$ In conclusion, the normalization factor should take the form $$N(F)=\left(Det(1+F)\right)^1.$$ (79) Since $`Det(1+F)={\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}det(1+F),`$ $`{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}1=\zeta (0)={\displaystyle \frac{1}{2}},`$ (80) we obtain $$N(F)=\left(det(1+F)\right)^{\frac{1}{2}}.$$ (81) ## Appendix B The conformal path integral representation of the boundary state in the linear dilaton background The path integral representation of the boundary state in the linear dilaton background is the form, $$\|B=[dx]\mathrm{exp}\left(V_{coup}^n\mathrm{exp}\left(\frac{\widehat{q}+\widehat{\stackrel{~}{q}}}{2Q}\right)\mathrm{\Phi }(x)\right)\|x,p=\overline{p}=iQ.$$ (82) The conformal invariant constraint of the boundary state, $`(L_n\stackrel{~}{L}_n)\|B=0`$, can be expressed by $`(L_n\stackrel{~}{L}_n)\|B`$ $`=`$ $`{\displaystyle [dx][\alpha _0\alpha _n\stackrel{~}{\alpha }_0\stackrel{~}{\alpha }_n,\mathrm{exp}S(x,\widehat{q},\widehat{\stackrel{~}{q}})]\|x,p=\overline{p}=iQ}+`$ (83) $`{\displaystyle [dx]\mathrm{exp}S(x,\widehat{q},\widehat{\stackrel{~}{q}})\left(2iQn^2\frac{}{x_n}+\underset{m0,n}{\overset{\mathrm{}}{}}mx_{m+n}\frac{}{x_m}\right)}`$ $`\|x,p=\overline{p}=iQ=0.`$ So the equation $`K_n\|B,p=0`$ now can be changed into the following two equations, $`2iQn^2{\displaystyle \frac{}{x_n}}\mathrm{\Phi }(x)`$ $`=`$ $`{\displaystyle \frac{i}{2Q}}x_n\mathrm{\Phi }(x)+{\displaystyle \underset{m0,n}{\overset{\mathrm{}}{}}}mx_{m+n}{\displaystyle \frac{}{x_m}}\mathrm{\Phi }(x),n0,`$ (84) $`{\displaystyle \underset{m0}{\overset{\mathrm{}}{}}}mx_m{\displaystyle \frac{}{x_m}}\mathrm{\Phi }(x)`$ $`=`$ $`0,`$ (85) with the commutative relation $`[\alpha _0\alpha _n\stackrel{~}{\alpha }_0\stackrel{~}{\alpha }_n,\mathrm{exp}S(x,\widehat{q},\widehat{\stackrel{~}{q}})]`$ $`={\displaystyle \frac{i}{2Q}}(\alpha _n\stackrel{~}{\alpha }_n)S(x,\widehat{q},\widehat{\stackrel{~}{q}})\mathrm{exp}S(x,\widehat{q},\widehat{\stackrel{~}{q}}).`$ (86) The special solution of the two equations could be found, $$\mathrm{\Phi }(x)=\frac{dz}{z}\mathrm{exp}\left(\underset{m0}{}\frac{ix_m}{2Qmz^m}\right).$$ (87) Now we come to check the solution. The right part of the (85) is $`{\displaystyle \frac{i}{2Q}}x_n\mathrm{\Phi }(x)+{\displaystyle \underset{m0,n}{\overset{\mathrm{}}{}}}mx_{m+n}{\displaystyle \frac{}{x_m}}\mathrm{\Phi }(x)`$ $`={\displaystyle 𝑑z(z^n)\frac{d}{dz}\mathrm{exp}\left(\underset{m0}{}\frac{ix_m}{2Qmz^m}\right)};`$ (88) and the left part of the (85) is $$2iQn^2\frac{}{x_n}\mathrm{\Phi }(x)=𝑑znz^{n1}\mathrm{exp}\left(\underset{m0}{}\frac{ix_m}{2Qmz^m}\right);$$ (89) so the difference between the above expressions is zero integral <sup>13</sup><sup>13</sup>13 if $`x_00`$ then the conclusion of the zero integral is not correct.. With $`(\alpha _n\stackrel{~}{\alpha }_nx_n)\|x,p=\overline{p}=iQ=0`$, and $`\widehat{q}\|x,p=\overline{p}=iQ=\widehat{\stackrel{~}{q}}\|x,p=\overline{p}=iQ`$, the boundary state is transformed into $`\|B`$ $`=`$ $`{\displaystyle [dx]\mathrm{exp}\left(V_{coup}^n\mathrm{exp}\left(\frac{\widehat{q}+\widehat{\stackrel{~}{q}}}{2Q}\right)\mathrm{\Phi }(x)\right)\|x,p=\overline{p}=iQ}`$ (90) $`=`$ $`\mathrm{exp}\left(V_{coup}^n\mathrm{exp}\left({\displaystyle \frac{\widehat{q}+\widehat{\stackrel{~}{q}}}{2Q}}\right){\displaystyle \frac{dz}{z}\mathrm{exp}\left(\underset{m0}{}\frac{i(\alpha _m\stackrel{~}{\alpha }_m)}{2Qmz^m}\right)}\right)\|B,p=\overline{p}=iQ_N`$ $`=`$ $`\mathrm{exp}\left(V_{coup}^n\mathrm{exp}({\displaystyle \frac{\widehat{q}}{Q}}){\displaystyle \frac{dz}{z}\mathrm{exp}\left(\underset{m0}{}\frac{i\alpha _m}{Qmz^m}\right)}\right)\|B,p=\overline{p}=iQ_N`$ which contains the divergence, so needs normalizing and renormalization. With the commutator $`[\alpha _n,\alpha _m]=n(1ϵ)^{\|n\|}\delta _{n+m,0}`$, the expression $`\mathrm{exp}\left(_{m0}\frac{i\alpha _m}{Qmz^m}\right)`$ may be represented by $$\mathrm{exp}\left(\underset{m=1}{\overset{\mathrm{}}{}}\frac{i\alpha _mz^m}{Qm}\right)\mathrm{exp}\left(\underset{m=1}{\overset{\mathrm{}}{}}\frac{i\alpha _m}{Qmz^m}\right)(ϵ^{\frac{1}{2}Q^2}).$$ (91) With the renormalized coupling constant $`V_{coup}^r`$, the boundary state in the path integral representation has the final form $$\|B=\mathrm{exp}\left(V_{coup}^r\mathrm{exp}(\frac{\widehat{q}}{Q})\frac{dz}{z}\mathrm{exp}\left(\underset{m=1}{\overset{\mathrm{}}{}}\frac{\alpha _mz^m}{iQm}\right)\mathrm{exp}\left(\underset{m=1}{\overset{\mathrm{}}{}}\frac{\alpha _m}{iQmz^m}\right)\right)\|B,p=\overline{p}=iQ_N.$$ (92) Since the boundary state $`\|B,p=\overline{p}=iQ_N`$ is conformal invariant, we have to verify that the former part of $`\|B,p=\overline{p}=iQ_N`$ should commute with Virasoro generator $`L_n`$ to ensure the conformal invariance of $`\|B`$ . In fact that part is similar to the screening charge (60), $`QQ`$ $`=`$ $`{\displaystyle 𝑑z}:\mathrm{exp}\left(ikX(z)\right):`$ (93) $`=`$ $`{\displaystyle 𝑑z\mathrm{exp}(ik\widehat{q})z^{k\widehat{p}}\mathrm{exp}\left(\underset{m=1}{\overset{\mathrm{}}{}}\frac{k\alpha _mz^m}{m}\right)\mathrm{exp}\left(\underset{m=1}{\overset{\mathrm{}}{}}\frac{k\alpha _m}{mz^m}\right)}`$ $`\|B,p=\overline{p}=iQ_N,`$ with $`k`$ must satisfy the following equation, $$k^2+2iQk=2.$$ (94) With $`k=\frac{1}{iQ}`$, the screening charge is just the former factor of the first order term of coupling constant of the boundary state. In addition, the solutions of equation (94) may be labeled by $`\frac{1}{a_+}`$ or $`\frac{1}{a_{}}`$. The $`a_+`$ and $`a_{}`$ are respectively $`\frac{1}{2}(Q+\sqrt{Q^22})`$ and $`\frac{1}{2}(Q\sqrt{Q^22})`$. So the solution $`k=1/(iQ)`$ is good in the large Q limit with the corrections suppressed by $`1/Q^2`$ etc. to $`k`$. Such a limit will change the boundary state$`\|B`$ into $`\|B,p=\overline{p}=iQ_N`$ which is conformal invariant. Therefore it is reasonable that the renormalized boundary state should be conformal invariant in some special backgrounds.
warning/0005/math0005243.html	ar5iv	text	# On representations of a q-analogue of the ∗-algebra 𝑃⁢𝑜⁢𝑙⁢(𝑀⁢𝑎⁢𝑡_{2,2}) ## Abstract Bounded Hilbert space $``$-representations are studied for a $`q`$-analogue of the $``$-algebra $`Pol(Mat_{2,2})`$ of polynomials on the space $`Mat_{2,2}`$ of complex $`2\times 2`$ matrices. <sup>0</sup><sup>0</sup>footnotetext: 2000 Mathematics Subject Classification: Primary 17B37; Secondary 20G42 1. Introduction The study of q-analogues of the Cartan domains (irreducible bounded symmetric domains) was initiated by S. Sinel’shchikov and L. Vaksman in \[SV\]. In particular, for each Cartan domain they defined the $``$-algebra $`Pol(𝔤_1)_q`$, a q-analogue of the polynomial algebra on the prehomogenous vector space $`𝔤_1`$, and set a problem on investigation of their representations. The theory of representations of the $``$-algebras corresponding to domains of rank 1 is well-understood. In this paper our purpose is to study such representations for one of the popular Cartan domains of rank 2, the matrix ball in the space $`Mat_{2,2}`$ of complex $`2\times 2`$ matrices. Following \[SSV\] we will denote this $``$-algebra by $`Pol(Mat_{2,2})_q`$. A description of $`Pol(Mat_{m,n})_q`$, $`m`$, $`n`$, in terms of generators and relations is given in \[SSV\]. In the paper we classify all irreducible representations of $`Pol(Mat_{2,2})_q`$ by bounded operators on a Hilbert space. The method which we use here is based on the study of some dynamical system arising on a spectrum of a commutative $``$-subalgebra of $`Pol(Mat_{2,2})_q`$ (see \[OS\]). Note that the $``$-algebra has also unbounded $``$-representation. One can easily define a “well-behaved” class of such unbounded representations and classify them up to unitary equivalence using the same technique. In the paper we use the following standard notations: $``$ is the set of real numbers, $`^+`$ is the set of nonnegative real numbers, $``$ denotes the set of integers, $`^+=\{0,1,2,\mathrm{}\}`$. 1. The $``$-algebra $`Pol(Mat_{2,2})_q`$ and its $``$-representations Let $`q(0,1)`$. The $``$-algebra $`Pol(Mat_{2,2})_q`$, a $`q`$-analogues of polynomials on the space $`Mat_{2,2}`$ of complex $`2\times 2`$ matrices, is given by its generators $`\{z_a^\alpha \}_{a=1,2;\alpha =1,2}`$ and the following commutation relations: $$\begin{array}{cccccc}\hfill z_1^1z_2^1& =& qz_2^1z_1^1,\hfill & \hfill z_2^1z_1^2& =& z_1^2z_2^1,\hfill \\ \hfill z_1^1z_1^2& =& qz_1^2z_1^1,\hfill & \hfill z_2^1z_2^2& =& qz_2^2z_2^1,\hfill \\ \hfill z_1^1z_2^2z_2^2z_1^1& =& (qq^1)z_1^2z_2^1,\hfill & \hfill z_1^2z_2^2& =& qz_2^2z_1^2,\hfill \end{array}$$ (1) $$\begin{array}{ccc}\hfill (z_1^1)^{}z_1^1& =& q^2z_1^1(z_1^1)^{}(1q^2)(z_2^1(z_2^1)^{}+z_1^2(z_1^2)^{})+\hfill \\ & & +q^2(1q^2)^2z_2^2(z_2^2)^{}+1q^2,\hfill \\ \hfill (z_2^1)^{}z_2^1& =& q^2z_2^1(z_2^1)^{}(1q^2)z_2^2(z_2^2)^{}+1q^2,\hfill \\ \hfill (z_1^2)^{}z_1^2& =& q^2z_1^2(z_1^2)^{}(1q^2)z_2^2(z_2^2)^{}+1q^2,\hfill \\ \hfill (z_2^2)^{}z_2^2& =& q^2z_2^2(z_2^2)^{}+1q^2,\hfill \end{array}$$ (2) and $`\begin{array}{cccccc}\hfill (z_1^1)^{}z_2^1qz_2^1(z_1^1)^{}& =& (qq^1)z_2^2(z_1^2)^{},\hfill & \hfill (z_2^2)^{}z_2^1& =& qz_2^1(z_2^2)^{},\hfill \\ \hfill (z_1^1)^{}z_1^2qz_1^2(z_1^1)^{}& =& (qq^1)z_2^2(z_2^1)^{},\hfill & \hfill (z_2^2)^{}z_1^2& =& qz_1^2(z_2^2)^{},\hfill \\ \hfill (z_1^1)^{}z_2^2& =& z_2^2(z_1^1)^{},\hfill & \hfill (z_2^1)^{}z_1^2& =& z_1^2(z_2^1)^{}.\hfill \end{array}`$ (6) Consider a representation $`\pi `$ of $`Pol(Mat_{2,2})_q`$ on a separable Hilbert space $`H`$ by bounded operators. The theorem below gives the complete classification of such irreducible representations up to a unitary equivalence. ###### Theorem 1 Any irreducible representation $`\pi `$ is unitarily equivalent to one from the following 6 series: 1) one-dimensional representations $`\xi _{\phi _1,\phi _2}`$ $$\xi _{\phi _1,\phi _2}(z_1^1)=q^1e^{i\phi _1},\xi _{\phi _1,\phi _2}(z_2^1)=\xi _{\phi _1,\phi _2}(z_1^2)=0,\xi _{\phi _1,\phi _2}(z_2^2)=e^{i\phi _2},$$ (7) $`\phi _i[0,2\pi )`$; 2) infinite-dimensional representations $`\pi _\phi `$ on $`H=l_2(^+)`$ $`\begin{array}{ccc}\hfill \pi _\phi (z_1^1)e_k& =& q^1\sqrt{1q^{2(k+1)}}e_{k+1},\hfill \\ \hfill \pi _\phi (z_2^2)e_k& =& e^{i\phi }e_k,\hfill \end{array}`$ (10) $`\pi _\phi (z_2^1)=\pi _\phi (z_1^2)=0,`$ $`\phi [0,2\pi )`$; 3) infinite-dimensional representations $`\rho _{\phi _1,\phi _2}`$ on $`H=l_2(^+)`$ $`\begin{array}{ccc}\hfill \rho _{\phi _1,\phi _2}(z_1^1)e_k& =& e^{i(\phi _1+\phi _2)}q^1\sqrt{1q^{2k}}e_{k1},\hfill \\ \hfill \rho _{\phi _1,\phi _2}(z_2^1)e_k& =& e^{i\phi _1}q^ke_k,\hfill \\ \hfill \rho _{\phi _1,\phi _2}(z_1^2)e_k& =& e^{i\phi _2}q^ke_k,\hfill \\ \hfill \rho _{\phi _1,\phi _2}(z_2^2)e_k& =& \sqrt{1q^{2(k+1)}}e_{k+1},\hfill \end{array}`$ (15) $`\phi _i[0,2\pi )`$; 4a) infinite-dimensional representations $`\rho _\phi ^1`$ on $`H=l_2(^+\times ^+)`$ $`\begin{array}{ccc}\hfill \rho _\phi ^1(z_1^1)e_{m,k}& =& e^{i\phi }q^1\sqrt{1q^{2(m+1)}}\sqrt{1q^{2k}}e_{m+1,k1},\hfill \\ \hfill \rho _\phi ^1(z_2^1)e_{m,k}& =& q^k\sqrt{1q^{2(m+1)}}e_{m+1,k},\hfill \\ \hfill \rho _\phi ^1(z_1^2)e_{m,k}& =& e^{i\phi }q^ke_{m,k},\hfill \\ \hfill \rho _\phi ^1(z_2^2)e_{m,k}& =& \sqrt{1q^{2(k+1)}}e_{m,k+1},\hfill \end{array}`$ (20) $`\phi [0,2\pi )`$; 4b) infinite-dimensional representations $`\rho _\phi ^2`$ on $`H=l_2(^+\times ^+)`$ $`\begin{array}{ccc}\hfill \rho _\phi ^2(z_1^1)e_{m,k}& =& e^{i\phi }q^1\sqrt{1q^{2(m+1)}}\sqrt{1q^{2k}}e_{m+1,k1},\hfill \\ \hfill \rho _\phi ^2(z_2^1)e_{m,k}& =& e^{i\phi }q^ke_{m,k},\hfill \\ \hfill \rho _\phi ^2(z_1^2)e_{m,k}& =& q^k\sqrt{1q^{2(m+1)}}e_{m+1,k},\hfill \\ \hfill \rho _\phi ^2(z_2^2)e_{m,k}& =& \sqrt{1q^{2(k+1)}}e_{m,k+1},\hfill \end{array}`$ (25) $`\phi [0,2\pi )`$; 5) infinite-dimensional representations $`\widehat{\rho }_\phi `$ on $`H=l_2(^+\times ^+\times ^+)`$ $`\begin{array}{ccc}\hfill \rho (z_1^1)e_{m,l,k}& =& e^{i\phi }q^{m+l}e_{m,l,k}\hfill \\ & & q^1\sqrt{(1q^{2(l+1)})(1q^{2(m+1)})(1q^{2k})}e_{m+1,l+1,k1},\hfill \\ \hfill \rho (z_2^1)e_{m,l,k}& =& q^k\sqrt{1q^{2(m+1)}}e_{m+1,l,k},\hfill \\ \hfill \rho (z_1^2)e_{m,l,k}& =& q^k\sqrt{1q^{2(l+1)}}e_{m,l+1,k},\hfill \\ \hfill \rho (z_2^2)e_{m,l,k}& =& \sqrt{1q^{2(k+1)}}e_{m,l,k+1},\hfill \end{array}`$ (31) $`\phi [0,2\pi )`$; 6) infinite-dimensional representation $`\rho `$ on $`H=l_2(^+\times ^+\times ^+\times ^+)`$ $`\begin{array}{ccc}\hfill \rho (z_1^1)e_{s,m,l,k}& =& q^{m+l}\sqrt{1q^{2(s+1)}}e_{s+1,m,l,k}\hfill \\ & & q^1\sqrt{(1q^{2(l+1))}(1q^{2(m+1)})(1q^{2k})}e_{s,m+1,l+1,k1},\hfill \\ \hfill \rho (z_2^1)e_{s,m,l,k}& =& q^k\sqrt{1q^{2(m+1)}}e_{s,m+1,l,k},\hfill \\ \hfill \rho (z_1^2)e_{s,m,l,k}& =& q^k\sqrt{1q^{2(l+1)}}e_{s,m,l+1,k},\hfill \\ \hfill \rho (z_2^2)e_{s,m,l,k}& =& \sqrt{1q^{2(k+1)}}e_{s,m,l,k+1}.\hfill \end{array}`$ (37) Proof. Let us consider a $``$-subalgebra $``$ of $`Pol(Mat_{2,2})_q`$ which is generated by $`z_2^1`$, $`z_1^2`$, $`z_2^2`$ and $`(z_2^1)^{}`$, $`(z_1^2)^{}`$, $`(z_2^2)^{}`$. Direct computation shows that $`z_2^1(z_2^1)^{}`$, $`z_1^2(z_1^2)^{}`$, $`z_2^2(z_2^2)^{}`$ generate a commutative $``$-subalgebra of $``$ and satisfy the following relations: $$(z_a^\alpha (z_a^\alpha )^{})z_b^\beta =z_b^\beta F_{ba}^{\beta \alpha }(z_2^1(z_2^1)^{},z_1^2(z_1^2)^{},z_2^2(z_2^2)^{})$$ (38) where $$\begin{array}{c}𝔽_{21}(x_1,x_2,x_3)=(F_{22}^{11}(x_1,x_2,x_3),F_{21}^{12}(x_1,x_2,x_3),F_{22}^{12}(x_1,x_2,x_3))=\hfill \\ \hfill =(q^2x_1(1q^2)(x_31),x_2,x_3),\\ \hfill 𝔽_{12}(x_1,x_2,x_3)=(F_{12}^{21}(x_1,x_2,x_3),F_{11}^{22}(x_1,x_2,x_3),F_{12}^{22}(x_1,x_2,x_3))=\\ \hfill =(x_1,q^2x_2(1q^2)(x_31),x_3),\\ \hfill 𝔽_{22}(x_1,x_2,x_3)=(F_{22}^{21}(x_1,x_2,x_3),F_{21}^{22}(x_1,x_2,x_3),F_{22}^{22}(x_1,x_2,x_3))=\\ \hfill =(q^2x_1,q^2x_2,q^2(x_31)+1).\end{array}$$ The functions $`𝔽_{21}`$, $`𝔽_{12}`$, $`𝔽_{22}:^3^3`$ define an action of $`^3`$ on $`^3`$ with orbits $$\begin{array}{c}\mathrm{\Omega }_{x_1,x_2,x_3}=\{𝔽_{21}^{(m)}(𝔽_{12}^{(l)}(𝔽_{22}^{(k)}(x_1,x_2,x_3)))=\hfill \\ \hfill =(q^{2k}(q^{2m}x_1(1q^{2m})(x_31)),q^{2k}(q^{2l}x_2(1q^{2l})(x_31)),\\ \hfill q^{2k}(x_31)+1),m,l,k\}.\end{array}$$ Here and in the sequel we denote by $`𝔽_{a\alpha }^{(m)}`$ the $`m`$-th iteration of $`𝔽_{a\alpha }`$ and $`(𝔽_{a\alpha }^{(m)})_i`$, $`i=1,2,3`$, the $`i`$-th coordinate of $`𝔽_{a\alpha }^{(m)}`$. Let $`\pi `$ be a $``$-representation of $`Pol(Mat_{2,2})_q`$ on a Hilbert space $`H`$ by bounded operators, let $`𝔼()`$ be the resolution of the identity for the commutative family $`𝔸_\pi `$ of the positive operators $`\pi (z_2^1)\pi (z_2^1)^{}`$, $`\pi (z_1^2)\pi (z_1^2)^{}`$, $`\pi (z_2^2)\pi (z_2^2)^{}`$ and let $`\sigma _\pi `$ be the joint spectrum of the family $`𝔸_\pi `$. Next step is to show that any irreducible representation is concentrated on an orbit of this dynamical system. ###### Lemma 1 If $`\pi `$ is an irreducible representation of $`Pol(Mat_{2,2})_q)`$ then the spectral measure $`𝔼()`$ is ergodic with respect to the action of the dynamical system generated by $`𝔽_{21}`$, $`𝔽_{12}`$, $`𝔽_{22}`$ and there exists an orbit $`\mathrm{\Omega }_{x_1,x_2,x_3}`$ such that $`𝔼(\mathrm{\Omega }_{x_1,x_2,x_3})=I`$. Proof. From (38) and the spectral theorem it follows that $`\begin{array}{ccc}\hfill 𝔼(\mathrm{\Delta })\pi (z_b^\beta )& =& \pi (z_b^\beta )𝔼(𝔽_{b\beta }^{(1)}(\mathrm{\Delta })),\hfill \\ \hfill 𝔼(\mathrm{\Delta })\pi (z_b^\beta )^{}& =& \pi (z_b^\beta )^{}𝔼(𝔽_{b\beta }(\mathrm{\Delta })),\hfill \end{array}`$ for any $`\mathrm{\Delta }𝔅(^3)`$. Hence any subset $`\mathrm{\Delta }`$ such that $`𝔽_{b\beta }^{(1)}(\mathrm{\Delta })\mathrm{\Delta }`$, $`𝔽_{b\beta }(\mathrm{\Delta })\mathrm{\Delta }`$, $`(b,\beta )=(2,1)`$, $`(1,2)`$, $`(2,2)`$ defines a subspace $`𝔼(\mathrm{\Delta })H`$ which is invariant with respect to the operators $`\pi (z_b^\beta )`$, $`\pi (z_b^\beta )^{}`$ for any $`(b,\beta )`$ as above. Moreover, such subspace is invariant with respect to any operator of the representation $`\pi `$. In fact, the following relations hold in $`Pol(Mat_{2,2})_q`$ $`z_a^\alpha (z_a^\alpha )^{}z_1^1=z_1^1z_a^\alpha (z_a^\alpha )^{}(1)^{a+\alpha }(qq^1)z_2^1z_1^2(z_2^2)^{}`$ (40) $`(a,\alpha )=(2,1)`$, $`(1,2)`$, $`(2,2)`$, which gives $`𝔼(^3\mathrm{\Delta })\pi (z_a^\alpha (z_a^\alpha )^{})\pi (z_1^1)𝔼(\mathrm{\Delta })=𝔼(^3\mathrm{\Delta })\pi (z_1^1)\pi (z_a^\alpha (z_a^\alpha )^{})𝔼(\mathrm{\Delta })`$ $`(1)^{a+\alpha }(qq^1)𝔼(^3\mathrm{\Delta })\pi (z_2^1)\pi (z_1^2)\pi (z_2^2)^{}𝔼(\mathrm{\Delta })`$ Therefore if $`\mathrm{\Delta }𝔅(^3)`$ is invariant with respect to all $`𝔽_{b\beta }^{(1)}`$ and $`𝔽_{b\beta }`$ we obtain $$\pi (z_a^\alpha (z_a^\alpha )^{})𝔼(^3\mathrm{\Delta })\pi (z_1^1)𝔼(\mathrm{\Delta })=𝔼(^3\mathrm{\Delta })\pi (z_1^1)𝔼(\mathrm{\Delta })\pi (z_a^\alpha (z_a^\alpha )^{})$$ and hence $$𝔼(\mathrm{\Delta }^{})𝔼(^3\mathrm{\Delta })\pi (z_1^1)𝔼(\mathrm{\Delta })=𝔼(^3\mathrm{\Delta })\pi (z_1^1)𝔼(\mathrm{\Delta })𝔼(\mathrm{\Delta }^{})$$ for any $`\mathrm{\Delta }^{}𝔅(^3)`$. Taking $`\mathrm{\Delta }^{}=\mathrm{\Delta }`$ gives $`𝔼(^3\mathrm{\Delta })\pi (z_1^1)𝔼(\mathrm{\Delta })=0`$, i.e. $`\pi (z_1^1)𝔼(\mathrm{\Delta })H𝔼(\mathrm{\Delta })H`$. Similarly, $`\pi (z_1^1)^{}𝔼(\mathrm{\Delta })H𝔼(\mathrm{\Delta })H`$. The ergodicity of the measure $`𝔼()`$ follows immediately, i.e., $`𝔼(\mathrm{\Delta })=I`$ or $`0`$ for any Borel $`\mathrm{\Delta }`$ which is invariant with respect to $`𝔽_{b\beta }`$, $`𝔽_{b\beta }^{(1)}`$. The simplest invariant sets are the orbits of the dynamical system. The next step is to show that only atomic measures concentrated on an orbit give rise to irreducible representation of the $``$-algebra. It is easily seen that the dynamical system generated by $`𝔽_{b\beta }`$ is one-to-one and possesses a measurable section, i.e., a set $`\tau 𝔅(^3)`$ which intersects any orbit in a single point. This implies that any ergodic measure is concentrated on a single orbit of the dynamical system and therefore $`𝔼(\mathrm{\Omega }_{x_1,x_2,x_3})=I`$ for some orbit $`\mathrm{\Omega }_{x_1,x_2,x_3}`$. We now clarify which orbits $`\mathrm{\Omega }_{x_1,x_2,x_3}`$ give rise to bounded irreducible representation $`\pi `$, i.e., $`\sigma _\pi \mathrm{\Omega }_{x_1,x_2,x_3}`$, and classify all such representations up to unitary equivalence. We claim first that there is no bounded representations $`\pi `$ with $`\sigma _\pi \mathrm{\Omega }_{x_1,x_2,x_3}`$ if $`x_3>1`$. From (38) we have $$\pi (z_b^\beta )H_xH_{𝔽_{b\beta }(x)},\pi (z_b^\beta )^{}H_xH_{𝔽_{b\beta }^{(1)}(x)},$$ (41) where $`H_x`$ is the eigenspace for $`𝔸_\pi `$ corresponding to the eigenvalue $`x^3`$. Since $`y=(y_1,y_2,y_3)\mathrm{\Omega }_{x_1,x_2,x_3}`$, where $`x_3>1`$, implies $`y_3>1`$ we conclude that $`\pi (z_2^2)\pi (z_2^2)^{}1`$ and $`\mathrm{ker}\pi (z_2^2)=\mathrm{ker}\pi (z_2^2)^{}=\{0\}`$. This clearly forces $`𝔽_{22}^{(k)}(y)\sigma _\pi `$ for any $`k`$. However, the set $`\{𝔽_{22}^{(k)}(y),k\}`$ is unbounded which contradicts the boundness of the representation $`\pi `$. Similar arguments show that there is no bounded representation $`\pi `$ with $`\sigma _\pi \mathrm{\Omega }_{x_1,x_2,1}`$, $`x_10`$ or $`x_20`$. In this case $`\mathrm{\Omega }_{x_1,x_2,1}=\{(q^{2(k+m)}x_1,q^{2(k+l)}x_2,1),k,l,m\}`$. The only possibility is $`\sigma _\pi =\mathrm{\Omega }_{0,0,1}=\{(0,0,1)\}`$ and in this case we obtain $`\pi (z_2^1)=\pi (z_1^2)=0`$, $`\pi (z_2^2)\pi (z_2^2)^{}=I`$. It follows now from (1)–(6) that $`\pi (z_2^2)`$, $`\pi (z_1^1)`$ satisfy the relations $`\pi (z_1^1)^{}\pi (z_1^1)=q^2\pi (z_1^1)\pi (z_1^1)^{}+(q^21),`$ $`[\pi (z_1^1),\pi (z_1^2)]=0,[\pi (z_1^1)^{},\pi (z_1^2)]=0,`$ (42) $`\pi (z_2^2)^{}\pi (z_2^2)=\pi (z_2^2)\pi (z_2^2)^{}=I.`$ This implies that $`\pi (z_2^2)`$ commutes with all images of the generators in the algebra under th e representation $`\pi `$ and therefore $`\pi (z_2^2)`$ is a multiple of the identity operator if $`\pi `$ is irreducible. By (On representations of a q-analogue of the $``$-algebra $`Pol(Mat_{2,2})`$) we have also $`\pi (z_2^2)=e^{i\phi _2}I`$, $`\phi _2[0,2\pi )`$. Irreducible representations of the relation $`(z_1^1)^{}z_1^1=q^2z_1^1(z_1^1)^{}+(q^21)`$ are well-known and can be easily calculated using the method of dynamical systems (see \[OS, Chapter 2\]). Any such representation is either one-dimensional: $`\xi _{\phi _1}(z_1^1)=q^1e^{i\phi _1}`$, $`\phi _1[0,2\pi )`$, or infinite-dimensional which is unitary equivalent to the following one $`\pi _\phi (z_1^1)e_k=q^1\sqrt{1q^{2(k+1)}}e_{k+1}`$. The corresponding irreducible representations of $`Pol(Mat_{2,2})_q`$ are $`\xi _{\phi _1,\phi _2}`$ and $`\pi _\phi `$. Since $`\sigma _\pi (^+)^3`$ and $`(𝔽_{22}^{(k)})_3(x_1,x_2,x_3)=q^{2k}(x_31)+1\mathrm{}`$ as $`k\mathrm{}`$, it follows from (41) that $`\mathrm{ker}\pi (z_2^2)^{}\{0\}`$, $`\mathrm{ker}\pi (z_2^2)\pi (z_2^2)^{}\{0\}`$ and the corresponding orbit contains a point $`(x_1,x_2,0)`$. We have $`\mathrm{\Omega }_{x_1,x_2,0}=\{(q^{2k}(q^{2m}(x_11)+1),q^{2k}(q^{2l}(x_21)+1),1q^{2k}),m,l,k\}`$. Similar arguments show that $`\sigma _\pi \mathrm{\Omega }_{x_1,x_2,0}`$, where $`x_1>1`$ or $`x_2>1`$, is impossible if the representation $`\pi `$ is bounded. From the positiveness of $`\sigma _\pi `$ we obtain also that the only orbits corresponding to irreducible representation of the $``$-algebra are $`\mathrm{\Omega }_{1,1,0}`$, $`\mathrm{\Omega }_{1,0,0}`$, $`\mathrm{\Omega }_{0,1,0}`$, $`\mathrm{\Omega }_{0,0,0}`$ and $`\mathrm{\Omega }_{0,0,1}`$ The last one was treated above. We consider now the case $`\sigma _\pi \mathrm{\Omega }_{x_1,x_2,x_3}`$, $`x_3=0`$. Let $`P_y`$, $`y=(y_1,y_2,y_3)`$ be the projection onto the eigenspace corresponding to the eigenvalue $`y`$. Using (40) we get $$(z_ky_k)P_z\pi (z_1^1)P_y=\pm (qq^1)P_z\pi (z_2^1)\pi (z_1^2)\pi (z_2^2)^{}P_y$$ (“+” for $`k=1,2`$ and “-” for $`k=3`$) $`z,y^3`$. By (41) we have $`\pi (z_2^1)\pi (z_1^2)\pi (z_2^2)^{}H_yH_{𝔽_{21}(𝔽_{12}(𝔽_{22}^{(1)}(y)))}`$ and $$\pi (z_2^1)\pi (z_1^2)\pi (z_2^2)^{}P_y=P_{𝔽_{21}(𝔽_{12}(𝔽_{22}^{(1)}(y)))}\pi (z_2^1)\pi (z_1^2)\pi (z_2^2)^{}P_y.$$ Setting $`P_{m,l,k}`$ the projection onto an eigenspace which corresponds to the eigenvalue $`𝔽_{21}^{(m)}(𝔽_{12}^{(l)}(𝔽_{22}^{(m)}(x_1,x_2,0)))`$ we obtain $$\pi (z_1^1)P_{m,l,k}=P_{m,l,k}\pi (z_1^1)P_{m,l,k}+P_{m+1,l+1,k1}\pi (z_1^1)P_{m,l,k},$$ i.e., $$\pi (z_1^1)H_{m,l,k}H_{m,l,k}H_{m+1,l+1,k1}.$$ Moreover, $`P_{m+1,l+1,k1}\pi (z_1^1)P_{m,l,k}=q^{12k}\pi (z_2^1)\pi (z_1^2)\pi (z_2^2)^{}P_{m,l,k}`$. The operator $`\pi (z_1^1)`$ can be written as a sum of its diagonal part $`\pi (z_1^1)_0=_{m,l,k}P_{m,l,k}\pi (z_1^1)P_{m,l,k}`$, and the operator $`_{m,l,k}q^{12k}\pi (z_2^1)\pi (z_1^2)\pi (z_2^2)^{}P_{m,l,k}=q\pi (z_2^1)\pi (z_1^2)\pi (z_2^2)^{}(1\pi (z_2^2(z_2^2)^{}))^1`$. Let now $`\sigma _\pi \mathrm{\Omega }_{x_1,x_2,0}`$, where $`x_10`$ or $`x_20`$. It follows from (1)-(6) by direct computation that $$\pi (z_1^1)_0^{}\pi (z_1^1)_0=q^2\pi (z_1^1)_0\pi (z_1^1)_0^{}.$$ The only bounded operator $`\pi (z_1^1)_0`$ satisfying this relation is the zero-operator. Therefore $$\pi (z_1^1)=q\pi (z_2^1)\pi (z_1^2)\pi (z_2^2)^{}(1\pi (z_2^2(z_2^2)^{}))^1$$ and $`\pi `$ is irreducible iff so is the family $`(\pi (z_2^1),\pi (z_1^2),\pi (z_2^2),\pi (z_2^1)^{},\pi (z_1^2)^{},\pi (z_2^2)^{})`$. Let $`\pi (z_a^\alpha )=U_a^\alpha \sqrt{\pi (z_a^\alpha )^{}\pi (z_a^\alpha )}`$ be the polar decomposition of $`\pi (z_a^\alpha )`$. Using easy arguments one can show that $`[U_a^\alpha ,U_b^\beta ]=[(U_a^\alpha )^{},U_b^\beta ]=0`$, $`(a,\alpha )(b,\beta )`$ and $$(z_a^\alpha (z_a^\alpha )^{})(U_b^\beta )=(U_b^\beta )F_{ba}^{\beta \alpha }(z_2^1(z_2^1)^{},z_1^2(z_1^2)^{},z_2^2(z_2^2)^{}).$$ Here $`(a,\alpha )`$, $`(b,\beta )\{(1,2),(2,1),(2,2)\}`$. Moreover, if $`\sigma _\pi \mathrm{\Omega }_{1,x_2,0}`$ ($`\sigma _\pi \mathrm{\Omega }_{x_1,1,0}`$) we have $`U_2^1`$ ($`U_1^2`$ respectively) commutes with any operators from the family $`𝔸_\pi `$ and therefore with any operator of the representation. This clearly forces $`U_2^1=e^{i\phi _1}I`$, $`\phi _1[0,2\pi )`$ ($`U_1^2=e^{i\phi _2}I`$, $`\phi _2[0,2\pi )`$ respectively). Let $`\sigma _\pi \mathrm{\Omega }_{0,1,0}`$. Consider $`e_{k,l}=(U_2^1)^k(U_2^2)^le`$, $`e\mathrm{ker}\pi (z_2)\pi (z_2^2)^{}\mathrm{ker}\pi (z_2^1)\pi (z_2^1)^{}`$, $`k,l^+`$. Then $`\{e_{k,l},k,l^+\}`$ is an orthonormal system which defines an invariant subspace. The corresponding irreducible representation is $`\rho _\phi ^2`$. Analogously $`(U_1^2)^k(U_2^2)^le=e_{k,l}`$, $`e\mathrm{ker}\pi (z_2)\pi (z_2^2)^{}\mathrm{ker}\pi (z_1^2)\pi (z_1^2)^{}`$, $`k,l^+`$, build an orthonormal basis of an irreducible representation space if $`\sigma _\pi \mathrm{\Omega }_{1,0,0}`$, the corresponding action is given by formulae (25). If $`\sigma _\pi \mathrm{\Omega }_{1,1,0}`$ we have that $`l.s.\{(U_2^2)^ke=e_k,k^+\}`$, $`e\mathrm{ker}\pi (z_2^2)\pi (z_2^2)^{}`$ is invariant with the corresponding action given by (15). We now turn to the case $`\sigma _\pi \mathrm{\Omega }_{0,0,0}`$. From (1)–(6) we have $`\pi (z_1^1)_0\pi (z_1^2)=q\pi (z_1^2)\pi (z_1^1)_0\pi (z_1^1)_0^{}\pi (z_1^2)=q\pi (z_1^2)\pi (z_1^1)_0^{},`$ $`\pi (z_1^1)_0\pi (z_2^1)=q\pi (z_2^1)\pi (z_1^1)_0\pi (z_1^1)_0^{}\pi (z_2^1)=q\pi (z_2^1)\pi (z_1^1)_0^{},`$ $`\pi (z_1^1)_0\pi (z_2^2)=\pi (z_2^2)\pi (z_1^1)_0\pi (z_1^1)_0^{}\pi (z_2^2)=\pi (z_2^2)\pi (z_1^1)_0^{},`$ $$\pi (z_1^1)_0^{}\pi (z_1^1)_0P_{m,l,k}=q^2\pi (z_1^1)_0\pi (z_1^1)_0^{}P_{m,l,k}+(1q^2)q^{2(m+l)}P_{m,l,k}.$$ Note that $`\pi (z_1^1)_0P_{m,l,k}HP_{m,l,k}H`$, $`\pi (z_1^1)_0^{}P_{m,l,k}HP_{m,l,k}H`$. Moreover, it follows from the above relation that if $`\pi `$ is irreducible then the family $`(\pi (z_1^1)_0,\pi (z_1^1)_0^{})`$ restricted to the subspace $`P_{m,l,k}H`$ is irreducible for any $`m,l,k^+`$. We have $$a^{}a=q^2aa^{}+(1q^2),$$ where $`a=\pi (z_1^1)_0P_{0,0,0}`$ Any irreducible family ($`a`$, $`a^{}`$) is either one-dimensional and given by $`a=e^{i\phi }`$, $`\phi [0,2\pi )`$, or infinite dimensional defined on $`l_2(^+)`$ by $`ae_s=\sqrt{1q^{2(s+1)}}e_{s+1}`$. These representations give rise to irreducible representations of the $``$-algebra $`Pol(Mat_{2,2})_q`$. Namely, in the first case we have that $`e_{m,l,k}=(U_2^1)^m(U_1^2)^l(U_2^2)^ke`$, where $`eP_{0,0,0}H=\mathrm{ker}\pi (z_2^2)\pi (z_2^2)^{}\mathrm{ker}\pi (z_1^2)\pi (z_1^2)^{}\mathrm{ker}\pi (z_2^1)\pi (z_2^1)^{}`$, $`m,l,k^+`$, define an orthonormal basis of the space where the irreducible representation $`\widehat{\rho }_\phi `$ acts, and for the second irreducible family we have that $`e_{s,m,l,k}=(U_2^1)^m(U_1^2)^l(U_2^2)^ke_s`$, $`s,m,l,k^+`$, define an orthonormal basis of the space where the irreducible representation $`\rho `$ acts. This finishes the proof. Comments. It follows from the proof that for any representation $`\pi `$ on a Hilbert space $`H_\pi `$ the family of self-adjoint operators $`\pi (z_2^2(z_2^2)^{})`$, $`\pi (z_2^1(z_2^1)^{})`$, $`\pi (z_1^2(z_1^2)^{})`$, $`\pi (z_1^1)_0\pi (z_1^1)_0^{}`$, where $$\pi (z_1^1)_0=\pi (z_1^1)\{\begin{array}{cc}0,\hfill & \pi (z_2^2(z_2^2)^{})=I\hfill \\ q\pi (z_2^1)\pi (z_1^2)\pi (z_2^2)^{}(1\pi (z_2^2(z_2^2)^{})^1,\hfill & \pi (z_2^2(z_2^2)^{})I\hfill \end{array}$$ generates a commutative $``$-subalgebra $`𝔸`$ in $`B(H_\pi )`$, the bounded operators on $`H_\pi `$. Moreover, any irreducible representation of $`Pol(Mat_{2,2})_q`$ is a weight representation with respect to this algebra, i.e., $`𝔸`$ can be diagonalized, and the spectrum of $`𝔸`$ is simple. A question which arise here is how to generalise the method to higher dimension matrix balls and classify $``$-representations of the corresponding $``$-algebras. In principle, just analysing the commutation relations between the generators in the $``$-algebra one can find a commutative $``$-subalgebra of $`Pol(Mat_{m,n})_q`$ or some its localisation and show that any irreducible representation $`\pi `$ is a weight representation with respect to this commutative $``$-algebra having a simple spectrum in this representation. However, the computations can be extremely difficult in general. ###### Remark 1 The polynomial algebra on the vector space $`Mat_{2,2}`$ can be supplied with a Poisson structure. Writing $`q=e^h`$ we have that $`Pol(Mat_{2,2})_{exp(h)}`$ is an associative algebra over the ring of formal series $`[[h]]`$ and $$Pol(Mat_{2,2})Pol(Mat_{2,2})_{exp(h)}/hPol(Mat_{2,2})_{exp(h)}.$$ The Poisson bracket now is given by $$\{amodh,bmodh\}=ih^1(abba)modh$$ for any $`a`$, $`bPol(Mat_{2,2})_{exp(h)}`$. The problem now is to define the symplectic leaves of this Poisson structure. Any primitive ideal $`\mathrm{ker}\pi `$, where $`\pi `$ is an irreducible representation of $`Pol(Mat_{2,2})_q`$, defines a maximal Poisson ideal $`I_\pi =\mathrm{ker}\pi modh`$ of the algebra $`Pol(Mat_{2,2})`$ ordered by inclusion and hence the closure of a symplectic leaf which is given by $`\{xMat_{2,2}f(x)=0,fI_\pi \}`$. As in the case of $`(SU(n))_q`$ (see \[SoV\]) one can expect that there is a one-to-one correspondence between irreducible representations (bounded irreducible representations) of $`Pol(Mat_{2,2})_q`$ and symplectic leaves (bounded symplectic leaves) in $`Mat_{2,2}`$. Acknowledgement. I am sincerely grateful to Prof. L. Vaksman for having communicated the problem to me and for a number helpful discussions during the preparation of the paper.
warning/0005/hep-th0005069.html	ar5iv	text	# The Renormalized Thermal Mass with Non-Zero Charge Density ## I Introduction It has been known for some time that conventional perturbation theory is inadequate for describing high-temperature field theory , with the perturbation expansion breaking down at some order in the coupling constant in the parameter regime where the temperature is very much greater than the bare mass. There have been several attempts to circumvent this problem (for example Refs. use resummed perturbative expansions and systematically include all relevant diagrams). A temperature-dependent renormalization group approach has been used in Ref. . The method we will employ in this paper is the linear $`\delta `$ expansion (see, for example ), which contains elements of both of the above methods. It involves a mass shift which is determined order-by-order in the expansion by a non-perturbative criterion, but prior to the (crucial) optimization stage it merely uses low-order perturbation theory with modified propagators and vertices. Pinto and Ramos have successfully applied the method to finite-temperature symmetry breaking in a scalar field theory with a $`\phi ^4`$ interaction up to second order in $`\delta `$, which in particular involves evaluating the non-trivial “setting-sun” diagram. In the present paper we extend this approach to tackle the problem of a non-zero chemical potential; in so doing we consider a complex scalar field rather than a real scalar field. For $`\mu =0`$, the effect is simply to alter the symmetry factors of the one- and two-loop diagrams, but for $`\mu 0`$ there is a shift in the energy component of the Matsubara propagator, which makes the diagrams considerably more difficult to evaluate but presents no problem of principle. To our knowledge there have been only a few attempts at this problem. Benson et al. performed a one-loop calculation, while Funakubo and Sakamoto used a renormalization group approach in the large-$`N`$ limit of an $`\mathrm{O}(N)`$ theory (in which the setting-sun diagram does not appear). In the standard lattice Monte-Carlo approach it is impossible to incorporate a non-zero chemical potential, because the Euclidean action then becomes complex and cannot be used as a real, positive statistical weight. We outline the linear $`\delta `$ expansion technique below and explain how it can generate nonperturbative results, while in Sec. II we present briefly the construction of the theory in the imaginary time formalism and include the chemical potential associated with a conserved charge. We then go on to calculate all diagrams contributing to the self-energy, up to second order in $`\delta `$. Sec. III includes results for the critical temperature and its behaviour as one allows $`\mu `$ to take non-zero values. The appendix contains a detailed calculation of the two-loop setting sun diagram in the high temperature expansion. ### A The Linear Delta Expansion The linear $`\delta `$ expansion (LDE) is a technique that allows the use of an analytic approach to probe the nonperturbative sector of the field theory to which it is applied. It has been employed with success in a wide variety of areas , with convergence of the expansion rigorously demonstrated in some simple zero- and one-dimensional models (see, for example, ). The method involves the introduction of an artificial expansion parameter, $`\delta `$, and the modification of the action of the theory under consideration via $$SS^\delta =(1\delta )S_0(\{\eta _i\})+\delta S,$$ (1) where $`\{\eta _i\}`$ is some set of variational parameters and $`S_0`$ represents the action of some soluble theory. This trial action is not determined by the method; however, an appropriate choice is usually suggested by the form of the theory under consideration. The $`\delta `$-modified action can then be used to evaluate some desired physical quantity as a power series in $`\delta `$ (truncated at some finite order, $`N`$), with $`\delta `$ then set equal to 1 at the end of the calculation. We label this quantity $`P_N`$, noting that it will, in general, have a residual dependence on the $`\{\eta _i\}`$. The variational element is introduced by fixing these parameters according to some specified criterion. We shall use the principle of minimal sensitivity (PMS) whereby the $`\{\eta _i\}`$ are chosen at a stationary point of the quantity $`P_N`$, namely: $$\frac{P_N}{\eta _i}\|_{\eta _i=\overline{\eta }_i}=0.$$ (2) It is this order by order fixing that allows for non-perturbative behaviour to emerge (the optimized variational parameters becoming functions of the order of truncation, and so being more correctly labelled by $`\{\eta _i(N)\}`$), and can also provide for convergence of the expansion, two particularly desirable features. ## II The Charged Scalar Field At Finite Temperature In Minkowski space, the Lagrangian density of a massive, complex scalar field $`\phi (x)`$ with a $`𝒱(\phi ^{}\phi )`$ interaction term is given by $$=(^\mu \phi ^{})(_\mu \phi )m^2\phi ^{}\phi 𝒱(\phi ^{}\phi ).$$ (3) This Lagrangian density has a well-known global $`U(1)`$ gauge symmetry which leads to a conserved charge $`𝒬`$ and an associated chemical potential $`\mu `$, so that the grand partition function is $$Z(\beta ,\mu )=\mathrm{T}re^{\beta (\widehat{H}\mu \widehat{Q})}.$$ (4) Proceeding through the Hamiltonian form of the path integral for $`Z`$, in the imaginary-time formalism of finite-temperature field theory , we arrive at $$Z(\beta ,\mu )=𝒟(\phi ^{},\phi )e^{S_E(\beta ,\mu )},$$ (5) where the Euclideanized action is (dropping the suffix for future notational convenience) $$S(\beta ,\mu )=S_F(\beta ,\mu )+S_I(\beta ,\mu )+S_C(\beta ,\mu ),$$ (6) with $`S_F(\beta ,\mu )`$ $`=`$ $`{\displaystyle _T}d^4x\phi ^{}\left[{\displaystyle \frac{^2}{\tau ^2}}+2\mu {\displaystyle \frac{}{\tau }}^2+m^2\mu ^2\right]\phi ,`$ (7) $`S_I(\beta ,\mu )`$ $`=`$ $`{\displaystyle _T}d^4x𝒱(\phi ^{}\phi ),`$ (8) $`S_C(\beta ,\mu )`$ $`=`$ $`{\displaystyle _T}d^4x\phi ^{}\left[A\left({\displaystyle \frac{^2}{\tau ^2}}2\mu {\displaystyle \frac{}{\tau }}+^2+\mu ^2\right)+Bm^2\right]\phi +{\displaystyle _T}d^4xC𝒱(\phi \phi ),`$ (9) where $`\tau =it`$, $`x=(\tau ,𝐱)`$ and $`_Td^4x=_0^\beta 𝑑\tau d^3𝐱.`$ Here $`S_C(\beta ,\mu )`$ represents those counterterms required to render the model finite. A sensible choice for our delta-modified action, in the presence of a $`(\phi ^{}\phi )^2`$ self-interaction, $`𝒱(\phi ^{}\phi )=(\lambda /4)(\phi ^{}\phi )^2`$, would then be $$S^\delta (\beta ,\mu )=S_F^\delta (\beta ,\mu )+S_I^\delta (\beta ,\mu )+S_C^\delta (\beta ,\mu ),$$ (10) with $`S_F^\delta (\beta ,\mu )`$ $`=`$ $`{\displaystyle _T}d^4x\phi ^{}\left[{\displaystyle \frac{^2}{\tau ^2}}+2\mu {\displaystyle \frac{}{\tau }}^2+\mathrm{\Omega }^2\mu ^2\right]\phi ,`$ (11) $`S_I^\delta (\beta ,\mu )`$ $`=`$ $`{\displaystyle _T}d^4x\left[{\displaystyle \frac{\delta \lambda }{4}}(\phi ^{}\phi )^2\delta \eta ^2\phi ^{}\phi \right],`$ (12) $`S_C^\delta (\beta ,\mu )`$ $`=`$ $`{\displaystyle _T}d^4x\phi ^{}\left[A^\delta \left({\displaystyle \frac{^2}{\tau ^2}}2\mu {\displaystyle \frac{}{\tau }}+^2+\mu ^2\right)+B^\delta \mathrm{\Omega }^2\right]\phi +{\displaystyle _T}d^4x\left[C^\delta {\displaystyle \frac{\delta \lambda }{4}}(\phi ^{}\phi )^2B^\delta \delta \eta ^2\phi ^{}\phi \right],`$ (13) introducing the variational parameter $`\mathrm{\Omega }`$, or equivalently $`\eta `$, given by $`\eta ^2=\mathrm{\Omega }^2m^2`$. In subsequent sections we will evaluate the thermal mass up to $`𝒪(\delta ^2)`$, using dimensional regularization to renormalize the theory. It is an important point, raised in , that the arbitrary variational parameter $`\eta `$ becomes a function of the bare parameters (which is how non-perturbative behaviour arises in the LDE), so we must eliminate any divergences before we apply the PMS optimization procedure. The renormalization procedure we use is identical to that of and is based on ; we forgo any detailed discussion of the cancelling of temperature-dependent divergences and any systematic calculation of the coefficients of the counterterms $`A^\delta `$, $`B^\delta `$ and $`C^\delta `$, the details being essentially identical to those discussed in . ### A Calculating the thermal mass We recall that the Feynman rules in frequency space in the presence of an overall charge are as follows: 1. To every line of the diagram assign a factor $`\mathrm{\Delta }_F(i\omega _n,𝐤)[(i\omega _n+\mu )^2+\omega _k^2]^1`$, where $`\omega _n=2\pi n/\beta `$ and $`\omega _k=\sqrt{(𝐤^2+\mathrm{\Omega }^2)}`$, and an arrow in the direction of momentum flow 2. Assign a factor $`\delta \lambda `$ to each vertex 3. Assign a factor $`\delta \eta ^2`$ to each insertion 4. Integrate over every internal line with the measure $`T_nd^3𝐤/(2\pi )^3`$ 5. There must be conservation of charge at each vertex, i. e. the number of arrows entering a vertex must be the same as the number leaving We are now in a position to estimate corrections to the thermal mass (which we do up to $`𝒪(\delta ^2)`$), defined by $$m_{T,\mu }^2=\mathrm{\Omega }^2+\mathrm{\Pi }(i\omega _n,𝐩),$$ (14) evaluated on-shell, i.e. $`i\omega _n=\mathrm{\Omega }\mu ,𝐩=\mathrm{𝟎}`$, where $`\mathrm{\Pi }`$ is the thermal self-energy. This choice of four-momentum is discussed in the Appendix, where the setting-sun diagram is calculated, the only momentum-dependent contribution to $`\mathrm{\Pi }`$ up to $`𝒪(\delta ^2)`$. ### B The thermal mass at first order To lowest order, the relevant contributions are $`\mathrm{\Pi }_{T,\mu }^\delta `$ $`=`$ $`\mathrm{\Pi }_1^\delta +\mathrm{\Pi }_2^\delta `$ (15) $`=`$ $`\delta \eta ^2+\delta \lambda T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^3𝐤}{(2\pi )^3}\mathrm{\Delta }_F(i\omega _n,𝐤)}.`$ (16) The frequency sum in (15) can be treated in a particularly concise and efficient manner through the mixed representation $`\mathrm{\Delta }_F(\tau ,𝐤)`$ of $`\mathrm{\Delta }_F`$: $`\mathrm{\Delta }_F(\tau ,𝐤)`$ $`=`$ $`T{\displaystyle \underset{n}{}}e^{i\omega _n\tau }\mathrm{\Delta }_F(i\omega _n,𝐤)`$ (17) $`=`$ $`{\displaystyle \frac{1}{2\omega _k}}\left[(1+n_k^{})e^{(\omega _k\mu )\tau }+n_k^+e^{(\omega _k+\mu )\tau }\right],`$ (18) where $$n_k^\pm =\frac{1}{e^{\beta (\omega _k\pm \mu )}1}$$ (19) is the Bose-Einstein distribution function in the presence of a non-zero chemical potential. The Matsubara propagator is recovered by Fourier transforming the mixed propagator with respect to the $`\tau `$ variable: $$\mathrm{\Delta }_F(i\omega _n,𝐤)=_0^\beta 𝑑\tau e^{i\omega _n\tau }\mathrm{\Delta }_F(\tau ,𝐤).$$ (20) The contribution of $`\mathrm{\Pi }_2^\delta `$ in (15) can then be calculated trivially to give $`\mathrm{\Pi }_2^\delta `$ $`=`$ $`\delta \lambda T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^3𝐤}{(2\pi )^3}_0^\beta 𝑑\tau e^{i\omega _n\tau }\mathrm{\Delta }_F(\tau ,𝐤)}`$ (21) $`=`$ $`\delta \lambda T{\displaystyle \frac{d^3𝐤}{(2\pi )^3}\frac{1}{2\omega _k}(1+n_k^{}+n_k^+)}.`$ (22) Using dimensional regularization one finds that the $`𝒪(\delta )`$ contribution to the thermal mass is $$m_{T,\mu }^2=\mathrm{\Omega }^2\mu ^2\delta \eta ^2+\delta \frac{\lambda }{16\pi ^2}\mathrm{\Omega }^2\left[\frac{1}{ϵ}+\mathrm{ln}\left(\frac{\mathrm{\Omega }^2}{4\pi M^2}\right)+\gamma 1\right]+\delta \frac{\lambda T^2}{\pi ^2}h_3^e(y,r),$$ (23) where $`M`$ is a mass scale introduced by dimensional regularization and, in the notation of , $`h_3^e(y,r)`$ $`=`$ $`\frac{1}{2}[h_3(y,r)+h_3(y,r)],`$ (24) $`h_3(y,r)`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Gamma }(3)}}{\displaystyle _0^{\mathrm{}}}𝑑x{\displaystyle \frac{x^2}{\sqrt{x^2+y^2}}}\left[{\displaystyle \frac{1}{\mathrm{exp}(\sqrt{x^2+y^2}ry)1}}\right],`$ (25) where $`x=\beta k`$, $`y=\beta \mathrm{\Omega }`$ and $`r=\mu /\mathrm{\Omega }`$. Using the minimal subtraction scheme (MS) we can eliminate the divergent term arising from the temperature-independent part of the self energy using the $`𝒪(\delta )`$ mass counterterm $$\mathrm{\Pi }_{\mathrm{ct}}^\delta =\left(\delta \frac{\lambda }{16\pi ^2ϵ}\right)\mathrm{\Omega }^2.$$ (26) We now analyze (24) in the limit $`y0`$ (the high temperature expansion) keeping $`r`$ fixed : $$h_3^e(y,r)=\frac{\pi ^2}{12}\frac{\pi y}{4}\sqrt{1r^2}\frac{y^2}{8}\left[\mathrm{ln}\left(\frac{y}{4\pi }\right)+\gamma \frac{1}{2}+r^2\right]+\mathrm{},$$ (27) leading to the following expression for the renormalized thermal mass at first order: $$m_{T,\mu }^2=\mathrm{\Omega }^2\delta \eta ^2+\delta \frac{\lambda T^2}{12}\delta \frac{\lambda T\mathrm{\Omega }}{4\pi }\sqrt{1r^2}+\delta \frac{\lambda \mathrm{\Omega }^2}{16\pi ^2}\left[\mathrm{ln}\left(\frac{4\pi T^2}{M^2}\right)\gamma 2r^2\right].$$ (28) In this and all subsequent calculations we neglect terms of $`𝒪(1/T)`$. ### C The thermal mass at second order At $`𝒪(\delta ^2)`$ there are five diagrams (shown in Fig. 2) which provide corrections to the self-energy. We begin by evaluating the first diagram, which in the high temperature limit is $$\mathrm{\Pi }_1^{\delta ^2}=\delta ^2\frac{\lambda T\eta ^2}{8\pi \mathrm{\Omega }}\frac{1}{\sqrt{1r^2}}\delta ^2\frac{\lambda \eta ^2}{16\pi ^2}\left[\frac{1}{ϵ}+\mathrm{ln}\left(\frac{4\pi T^2}{M^2}\right)\gamma \right].$$ (29) The next two diagrams are constructed using mass and vertex counterterms to give $`𝒪(\delta ^2)`$ diagrams that eliminate the temperature-dependent divergences arising from $`\mathrm{\Pi }_4^{\delta ^2}`$. Explicitly, we obtain in the high temperature limit $`\mathrm{\Pi }_2^{\delta ^2}`$ $`=`$ $`\delta ^2{\displaystyle \frac{\lambda ^2\mathrm{\Omega }^2}{(16\pi ^2)^2ϵ^2}}+\delta ^2{\displaystyle \frac{\lambda ^2}{16\pi ^2ϵ}}\left\{{\displaystyle \frac{T\mathrm{\Omega }}{8\pi }}{\displaystyle \frac{1}{\sqrt{1r^2}}}+{\displaystyle \frac{\mathrm{\Omega }^2}{16\pi ^2}}\left[\mathrm{ln}\left({\displaystyle \frac{4\pi T^2}{M^2}}\right)\gamma \right]\right\}`$ (31) $`\delta ^2{\displaystyle \frac{\lambda ^2\mathrm{\Omega }^2}{2(16\pi ^2)^2}}\left\{\left[\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Omega }^2}{4\pi M^2}}\right)+\gamma \right]^2+{\displaystyle \frac{\pi ^2}{6}}\right\},`$ $`\mathrm{\Pi }_3^{\delta ^2}`$ $`=`$ $`\delta ^2{\displaystyle \frac{5\lambda ^2\mathrm{\Omega }^2}{2(16\pi ^2)^2ϵ^2}}+\delta ^2{\displaystyle \frac{5\lambda ^2}{32\pi ^2ϵ}}\left\{{\displaystyle \frac{T^2}{12}}{\displaystyle \frac{T\mathrm{\Omega }}{4\pi }}\sqrt{1r^2}+{\displaystyle \frac{\mathrm{\Omega }^2}{16\pi ^2}}\left[\mathrm{ln}\left({\displaystyle \frac{4\pi T^2}{M^2}}\right)\gamma 2r^2\right]\right\}`$ (33) $`\delta ^2{\displaystyle \frac{5\lambda ^2\mathrm{\Omega }^2}{(32\pi ^2)^2}}\left\{\left[\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Omega }^2}{4\pi M^2}}\right)+\gamma 1\right]^2+1+{\displaystyle \frac{\pi ^2}{6}}\right\}.`$ The momentum-independent two-loop diagram is given by $`\mathrm{\Pi }_4^{\delta ^2}`$ $`=`$ $`\delta ^2{\displaystyle \frac{\lambda ^2\mathrm{\Omega }^2}{(16\pi ^2)^2}}{\displaystyle \frac{1}{ϵ^2}}\delta ^2{\displaystyle \frac{\lambda ^2}{16\pi ^2}}{\displaystyle \frac{1}{ϵ}}\left\{{\displaystyle \frac{T^2}{12}}{\displaystyle \frac{T\mathrm{\Omega }}{4\pi }}\sqrt{1r^2}{\displaystyle \frac{T\mathrm{\Omega }}{8\pi }}{\displaystyle \frac{1}{\sqrt{1r^2}}}+{\displaystyle \frac{\mathrm{\Omega }^2}{8\pi ^2}}\left[\mathrm{ln}\left({\displaystyle \frac{4\pi T^2}{M^2}}\right)\gamma r^2\right]\right\}`$ (36) $`\delta ^2{\displaystyle \frac{\lambda ^2T^3}{96\pi \mathrm{\Omega }}}{\displaystyle \frac{1}{\sqrt{1r^2}}}+\delta ^2{\displaystyle \frac{\lambda ^2T^2}{32\pi ^2}}+\delta ^2{\displaystyle \frac{\lambda ^2}{(8\pi )^2}}\left\{{\displaystyle \frac{T^3}{3}}{\displaystyle \frac{T\mathrm{\Omega }}{\pi }}\sqrt{1r^2}{\displaystyle \frac{T\mathrm{\Omega }}{2\pi }}{\displaystyle \frac{1}{\sqrt{1r^2}}}+{\displaystyle \frac{\mathrm{\Omega }^2}{4\pi ^2}}\left[\mathrm{ln}\left({\displaystyle \frac{4\pi T^2}{M^2}}\right)\gamma 2r^2\right]\right\}`$ $`\times \left[\mathrm{ln}\left({\displaystyle \frac{4\pi T^2}{M^2}}\right)\gamma \right]+\delta ^2{\displaystyle \frac{\lambda ^2T\mathrm{\Omega }}{64\pi ^3}}{\displaystyle \frac{r^2}{\sqrt{1r^2}}}+\delta ^2{\displaystyle \frac{\lambda ^2\mathrm{\Omega }^2}{(16\pi ^2)^2}}\left[\mathrm{ln}^2\left({\displaystyle \frac{\mathrm{\Omega }^2}{4\pi M^2}}\right)+(2\gamma 1)\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Omega }^2}{4\pi M^2}}\right)+2.4\right]`$ Finally, we have the momentum-dependent setting sun diagram, which we evaluate with the Euclidean self-energy on shell (the details of the calculation are provided in Appendix A): $`\mathrm{\Pi }_5^{\delta ^2}`$ $`=`$ $`\delta ^2{\displaystyle \frac{3\lambda ^2\mathrm{\Omega }^2}{(32\pi ^2)^2}}{\displaystyle \frac{1}{ϵ^2}}+\delta ^2{\displaystyle \frac{3\lambda ^2\mathrm{\Omega }^2}{(32\pi ^2)^2}}{\displaystyle \frac{1}{ϵ}}+\delta ^2{\displaystyle \frac{\lambda ^2p^2}{(32\pi ^2)^2}}{\displaystyle \frac{1}{2ϵ}}\delta ^2{\displaystyle \frac{3\lambda ^2}{16\pi ^2ϵ}}\left\{{\displaystyle \frac{T^2}{24}}{\displaystyle \frac{T\mathrm{\Omega }}{8\pi }}\sqrt{1r^2}+{\displaystyle \frac{\mathrm{\Omega }^2}{32\pi ^2}}\left[\mathrm{ln}\left({\displaystyle \frac{4\pi T^2}{M^2}}\right)\gamma 2r^2\right]\right\}`$ (40) $`+\delta ^2{\displaystyle \frac{3\lambda ^2\mathrm{\Omega }^2}{2(16\pi ^2)^2}}\left[\mathrm{ln}^2\left({\displaystyle \frac{\mathrm{\Omega }^2}{4\pi M^2}}\right)+(2\gamma \frac{17}{6})\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Omega }^2}{4\pi M^2}}\right)+1.9785\right]\delta ^2{\displaystyle \frac{3\lambda ^2}{16\pi ^2}}\left[\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Omega }^2}{4\pi M^2}}\right)+2\gamma \right]`$ $`\times \left\{{\displaystyle \frac{T^2}{24}}{\displaystyle \frac{T\mathrm{\Omega }}{8\pi }}\sqrt{1r^2}{\displaystyle \frac{\mathrm{\Omega }^2}{16\pi ^2}}\left[\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Omega }}{4\pi T}}\right)+\gamma \frac{1}{2}+r^2\right]\right\}`$ $`+\delta ^2{\displaystyle \frac{\lambda ^2T^2}{128\pi ^2}}\left[\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Omega }^2}{T^2}}\right)+5.0669+\frac{3}{2}\mathrm{ln}(1r^2){\displaystyle \frac{2r^2}{\pi ^2}}(2\mathrm{ln}21){\displaystyle \frac{1}{\pi ^2}}\mathrm{ln}^2\left({\displaystyle \frac{1+r}{1r}}\right)\mathrm{ln}2\right].`$ The divergent parts of these diagrams can all be eliminated by a suitable choice of the counterterm in Eq. (13). As mentioned above, we use the minimal subtraction prescription. ## III Results Having obtained an expression for the renormalized thermal mass, we can set $`\delta =1`$ and obtain numerical results for this mass, $`m_{T,\mu }^2(\overline{\eta })`$, where $`\overline{\eta }`$ is determined via the PMS condition: $$\frac{m_{T,\mu }^2(\eta )}{\eta }\|_{\eta =\overline{\eta }}=0.$$ (41) At $`𝒪(\delta )`$, $`\overline{\eta }`$ does not depend on the coupling and so does not generate non-perturbative information, but non-perturbative behaviour emerges at $`𝒪(\delta ^2)`$. Fig. 3 is a typical plot, showing a clear maximum in $`m_{T,\mu }^2(\eta )`$. It will be of particular interest to study the $`\mu `$ dependence of the critical temperature $`T_c`$, the signal for the phase transition being taken to be $`m_{T_c,\mu }^2=\mu ^2`$, which leads to an implicit equation for $`T_c=T_c(\mu )`$ (with dependence on the renormalized bare mass and coupling of the theory suppressed in the notation). Fig. 4 shows a plot of the contours of $`m_{T,\mu }^2`$ for $`m^2=M^2,\lambda =0.5`$ in a region of the $`(T,\mu )`$ plane, with the thicker line indicating the critical temperature. Below this line the LDE breaks down - the thermal mass being a monotonically decreasing function of $`\eta `$ and thus lacking extrema. In Fig. 5 we present a comparison of $`T_c(\mu )`$ as calculated by the LDE with the first-order estimate $$m^2+\frac{\lambda T_c^2}{12}=\mu ^2$$ (42) provided by perturbation theory in the high temperature limit . The two curves are quite similar in shape, and converge at large $`T_c`$. The reason for this behaviour is that in the high $`T`$ limit, it so happens that $`\overline{\eta }^2\mu ^2m^2`$, which is equivalent to $`r^21`$. Examining those contributions from $`\mathrm{\Pi }_4`$ in (36) we see that if this is the case, we require the coefficients of the $`1/\sqrt{(1r^2)}`$ terms to disappear, which is exactly the same condition as (42). Fig. 6 illustrates the high-$`T`$ behaviour of the second-order approximation to the thermal mass, with $`\mu /T_c`$ approaching a constant as $`T_c`$ increases; the constant being determined from (42) to be $`\sqrt{(\lambda /12)}`$. ## IV Conclusions In this paper we have shown how the problem of finite chemical potential for a charged scalar theory can be formulated in the context of the linear delta expansion. The graphs to be evaluated are essentially those of ordinary perturbation theory, with modified mass and coupling parameters. The chemical potential appears explicitly in the Lagrangian, and also in the Bose-Einstein factors occurring in the momentum integrals. Renormalization has been implemented on the lines of Ref. , where the importance of renormalizing before applying the variational aspect of the method was emphasized. We obtain unambiguous PMS points, as illustrated in Fig. 3. The results have been plotted in Fig. 4 as a contour plot of the thermal mass in the $`(T,\mu )`$ plane, and in Fig. 5 as $`T_c`$ versus $`\mu `$ for given values of $`m^2,\lambda `$. This latter curve approaches the result of resummed perturbation theory at high temperature, but differs significantly from it at lower temperatures. The reason for the convergence at higher temperatures can be understood in terms of the properties of the stationary points in the variational parameter $`\eta `$. As emphasized in the introduction, the problem of a non-zero chemical potential is not amenable to treatment by the usual Monte-Carlo lattice method. A lattice version of the present calculation is in progress. An extension of the present work which we intend to pursue in the near future is to free it from the dependence on the high temperature expansion, thereby enabling us to consider a system with non-zero chemical potential at low temperature. ## acknowledgments We are very grateful to Dr. T. S. Evans for useful discussions and the benefit of his expertise in thermal field theory. One of us (PP) wishes to thank the Particle Physics and Astronomy Research Council of the UK for financial support. ## A In this appendix we will derive in some detail the expression for the two-loop setting sun diagram. We begin by mimicking the approach of Parwani to split up the diagram into three parts; containing zero, one and two Bose-Einstein factors respectively. Explicitly, $$\mathrm{\Pi }_5^{\delta ^2}=\frac{\delta ^2\lambda ^2}{2}M^{4ϵ}T^2\underset{l,m}{}\frac{d^{d1}𝐤}{(2\pi )^{d1}}\frac{d^{d1}𝐪}{(2\pi )^{d1}}\mathrm{\Delta }_F(i\omega _l,𝐤)\mathrm{\Delta }_F(i\omega _m,𝐪)\mathrm{\Delta }_F(i\omega _r,𝐫),$$ (A1) where $`d=42ϵ`$ is the space-time dimension and $`𝐫=𝐩𝐤𝐪`$. If we now use the mixed representation of the Matsubara propagator defined in (20), the frequency sums become trivial and eventually we have $$\mathrm{\Pi }_5^{\delta ^2}(i\omega _n,𝐩)=\frac{\delta ^2\lambda ^2}{2}\left(G_0(i\omega _n,𝐩)+G_1(i\omega _n,𝐩)+G_2(i\omega _n,𝐩)\right),$$ (A2) where $`G_0(i\omega _n,𝐩)`$ $`=`$ $`{\displaystyle d[𝐤,𝐪]S_\mu (\omega _k,\omega _q,\omega _r)},`$ (A3) $`G_1(i\omega _n,𝐩)`$ $`=`$ $`{\displaystyle }d[𝐤,𝐪][n_k^+(S_\mu (\omega _k,\omega _q,\omega _r)+S_\mu (\omega _k,\omega _q,\omega _r)+S_\mu ^+(\omega _k,\omega _q,\omega _r)+S_\mu ^{}(\omega _k,\omega _q,\omega _r))`$ (A5) $`+n_k^{}(S_\mu (\omega _k,\omega _q,\omega _r)+S_\mu (\omega _k,\omega _q,\omega _r)+S_\mu ^+(\omega _k,\omega _q,\omega _r)+S_\mu ^{}(\omega _k,\omega _q,\omega _r))],`$ $`G_2(i\omega _n,𝐩)`$ $`=`$ $`{\displaystyle }d[𝐤,𝐪][n_k^+n_q^+(S_\mu ^+(\omega _k,\omega _q,\omega _r)+S_\mu (\omega _k,\omega _q,\omega _r)+S_\mu (\omega _k,\omega _q,\omega _r)S_\mu ^+(\omega _k,\omega _q,\omega _r))`$ (A9) $`+n_k^+n_q^{}\left(S_\mu (\omega _k,\omega _q,\omega _r)+S_\mu ^{}(\omega _k,\omega _q,\omega _r)+S_\mu ^+(\omega _k,\omega _q,\omega _r)S_\mu (\omega _k,\omega _q,\omega _r)\right)`$ $`+n_k^{}n_q^+\left(S_\mu (\omega _k,\omega _q,\omega _r)+S_\mu ^+(\omega _k,\omega _q,\omega _r)+S_\mu ^{}(\omega _k,\omega _q,\omega _r)S_\mu (\omega _k,\omega _q,\omega _r)\right)`$ $`+n_k^{}n_q^{}(S_\mu ^{}(\omega _k,\omega _q,\omega _r)+S_\mu (\omega _k,\omega _q,\omega _r)+S_\mu (\omega _k,\omega _q,\omega _r)S_\mu ^{}(\omega _k,\omega _q,\omega _r))],`$ with the definitions $`S_\mu ^\pm (\omega _k,\omega _q,\omega _r)`$ $`=`$ $`{\displaystyle \frac{1}{\pm (i\omega _n+\mu )+\omega _k+\omega _q+\omega _r}},`$ (A10) $`S_\mu (\omega _k,\omega _q,\omega _r)`$ $`=`$ $`S_\mu ^+(\omega _k,\omega _q,\omega _r)+S_\mu ^{}(\omega _k,\omega _q,\omega _r),`$ (A11) $`d[𝐤,𝐪]`$ $`=`$ $`M^{4ϵ}{\displaystyle \frac{d^{d1}𝐤}{(2\pi )^{d1}}}{\displaystyle \frac{d^{d1}𝐪}{(2\pi )^{d1}}}{\displaystyle \frac{1}{8\omega _k\omega _q\omega _r}}.`$ (A12) We choose to evaluate the self-energy on shell ($`𝐩=\mathrm{𝟎},i\omega _n=\mathrm{\Omega }\mu `$). The reason we evaluate at this point rather than at $`𝐩=\mathrm{𝟎},i\omega _n=\mathrm{\Omega }`$ is best explained as follows. In setting up the theory, we are dealing with an effective Hamiltonian, $`\widehat{H}_{\mathrm{eff}}=\widehat{H}\mu \widehat{Q}`$ rather than the real Hamiltonian, $`\widehat{H}`$. We would naturally choose to evaluate the self-energy at $`i\omega _n=\mathrm{\Omega }`$ for the real Hamiltonian, which means evaluating at $`i\omega _n=\mathrm{\Omega }\mu `$ in the effective theory. Following this prescription, we find that $`\mathrm{}[G_0(\mathrm{\Omega }\mu ,\mathrm{𝟎})]`$ $`=`$ $`{\displaystyle \frac{3\mathrm{\Omega }^2}{2(16\pi ^2)^2}}\left[{\displaystyle \frac{1}{ϵ^2}}+{\displaystyle \frac{32\gamma }{ϵ}}{\displaystyle \frac{2}{ϵ}}\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Omega }^2}{4\pi M^2}}\right)\right]{\displaystyle \frac{p^2}{(32\pi ^2)^2}}{\displaystyle \frac{1}{ϵ}}`$ (A14) $`{\displaystyle \frac{3\mathrm{\Omega }^2}{(16\pi ^2)^2}}\left[\mathrm{ln}^2\left({\displaystyle \frac{\mathrm{\Omega }^2}{4\pi M^2}}\right)+(2\gamma \frac{17}{6})\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Omega }^2}{4\pi M^2}}\right)+1.9785\right].`$ We now calculate the contribution from $`G_1`$. Writing $$G_1(i\omega _n,𝐩)=G_1^+(i\omega _n,𝐩)+G_1^{}(i\omega _n,𝐩),$$ (A15) where $`G_1^\pm `$ represents that part of $`G_1`$ with an associated BE factor of $`n_k^\pm `$ respectively, and decomposing this factor into a UV divergent part and a UV finite part in the manner of the $`\mu =0`$ case one obtains $$\mathrm{}[G_1^\pm (\mathrm{\Omega }\mu ,\mathrm{𝟎})]=F_0^\pm +F_1^\pm +F_2^\pm (\mathrm{\Omega }^2).$$ (A16) Here $`F_0^\pm `$ $`=`$ $`{\displaystyle \frac{3}{32\pi ^2}}{\displaystyle \frac{T^2}{\pi ^2}}h_3(y,\pm r){\displaystyle \frac{1}{ϵ}},`$ (A17) $`F_1^\pm `$ $`=`$ $`{\displaystyle \frac{3}{32\pi ^2}}{\displaystyle \frac{T^2}{\pi ^2}}h_3(y,\pm r)\left[\mathrm{ln}\left({\displaystyle \frac{4\pi M^2}{\mathrm{\Omega }^2}}\right)+2\gamma \right],`$ (A18) and $$F_2^\pm (\mathrm{\Omega }^2)=\frac{1}{4(2\pi )^4}_0^{\mathrm{}}𝑑k\frac{kn_k^\pm }{\omega _k}_0^{\mathrm{}}\frac{dq}{\omega _q}\left[q\mathrm{ln}\left\|\frac{X_+^\pm }{X_{}^\pm }\right\|6k\right],$$ (A19) with $$X_\pm ^\pm =[\mathrm{\Omega }^2(\omega _k+\omega _q+\omega _{k\pm q})^2][\mathrm{\Omega }^2(\omega _k+\omega _q+\omega _{k\pm q})^2][(\mathrm{\Omega }\pm \omega _k)^2(\omega _q+\omega _{k\pm q})^2],$$ (A20) where the subscript $`\pm `$ refers to the $`\omega _{k\pm q}`$ term. Following a similar procedure to that of Ref. we find that in the high-$`T`$ limit $$F_2^\pm (\mathrm{\Omega }^2)\frac{T^2}{128\pi ^2}\left[\mathrm{ln}\left(\frac{\mathrm{\Omega }}{T}\right)0.54597\right]$$ (A21) Thus, collecting all the terms together, $$\mathrm{}[G_1(\mathrm{\Omega }\mu ,\mathrm{𝟎})]=F_0+F_1+F_2(\mathrm{\Omega }^2),$$ (A22) with $`F_0`$ $`=`$ $`{\displaystyle \frac{3T^2}{16\pi ^4}}h_3^e(y,r){\displaystyle \frac{1}{ϵ}},`$ (A23) $`F_1`$ $`=`$ $`{\displaystyle \frac{3T^2}{16\pi ^4}}h_3^e(y,r)\left[\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Omega }^2}{4\pi M^2}}\right)+2\gamma \right],`$ (A24) and $$F_2(\mathrm{\Omega }^2)\frac{T^2}{64\pi ^2}\left[\mathrm{ln}\left(\frac{\mathrm{\Omega }}{T}\right)0.54597\right].$$ (A25) Finally, we consider the contribution from $`G_2`$, which contains a BE factor for each loop and so is UV finite. We write $$\mathrm{}[G_2(\mathrm{\Omega }\mu ,\mathrm{𝟎})]=H^{++}(\mathrm{\Omega }^2)+H^+(\mathrm{\Omega }^2)+H^+(\mathrm{\Omega }^2)+H^{}(\mathrm{\Omega }^2)$$ (A26) with $$H^{\pm \pm }(\mathrm{\Omega }^2)=\frac{1}{4(2\pi )^4}_0^{\mathrm{}}𝑑k\frac{kn_k^\pm }{\omega _k}_0^{\mathrm{}}𝑑q\frac{qn_q^\pm }{\omega _q}\mathrm{ln}\left\|\frac{Y_+^{\pm \pm }}{Y_{}^{\pm \pm }}\right\|,$$ (A27) where $`Y_\pm ^{++}`$ $`=`$ $`[\mathrm{\Omega }^2(\omega _k+\omega _q+\omega _{k\pm q})^2][\mathrm{\Omega }^2(\omega _k\omega _q+\omega _{k\pm q})^2][(\mathrm{\Omega }+\omega _k+\omega _q)^2\omega _{k\pm q}^2],`$ (A28) $`Y_\pm ^+`$ $`=`$ $`[\mathrm{\Omega }^2(\omega _k+\omega _q+\omega _{k\pm q})^2][\mathrm{\Omega }^2(\omega _k+\omega _q\omega _{k\pm q})^2][(\mathrm{\Omega }+\omega _k\omega _q)^2\omega _{k\pm q}^2],`$ (A29) $`Y_\pm ^+`$ $`=`$ $`[\mathrm{\Omega }^2(\omega _k+\omega _q+\omega _{k\pm q})^2][\mathrm{\Omega }^2(\omega _k+\omega _q\omega _{k\pm q})^2][(\mathrm{\Omega }\omega _k+\omega _q)^2\omega _{k\pm q}^2],`$ (A30) $`Y_\pm ^{}`$ $`=`$ $`[\mathrm{\Omega }^2(\omega _k+\omega _q+\omega _{k\pm q})^2][\mathrm{\Omega }^2(\omega _k\omega _q+\omega _{k\pm q})^2][(\mathrm{\Omega }\omega _k\omega _q)^2\omega _{k\pm q}^2].`$ (A31) Each of the $`H^{\pm \pm }`$ pieces has a logarithmic IR divergence as $`\mathrm{\Omega }0`$. To deal with this, we extract the leading order behaviour of $`\mathrm{ln}\|Y_+^{\pm \pm }/Y_{}^{\pm \pm }\|`$ in this limit, separating into three terms as follows: $$\mathrm{ln}\left\|\frac{Y_+^{+\pm }}{Y_{}^{+\pm }}\right\|\mathrm{ln}\left\|\frac{Y_+^{}}{Y_{}^{}}\right\|\pm \mathrm{ln}2\pm \mathrm{ln}\left\|\frac{kq}{\mathrm{\Omega }\sqrt{k^2q^2}}\right\|+\frac{3}{2}\mathrm{ln}\left\|\frac{k+q}{kq}\right\|.$$ (A32) Collecting those terms involving the $`\mathrm{ln}2`$ piece alone, which we can take outside the integral, we calculate the coefficient of this term in the high-temperature limit $`a\beta \mathrm{\Omega }1`$, which we call $`C_{\mathrm{ln}2}`$, to be $$C_{\mathrm{ln}2}=\frac{T^2}{(8\pi ^2)^2}\mathrm{ln}^2\left(\frac{1+r}{1r}\right).$$ (A33) The remainder of $`\mathrm{}[G_2(\mathrm{\Omega }\mu ,\mathrm{𝟎})]`$ splits into two pieces which we define by $`H^1(\mathrm{\Omega }^2)`$ $`=`$ $`{\displaystyle \frac{1}{4(2\pi )^4}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{kdk}{\omega _k}}(n_k^+n_k^{}){\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{qdq}{\omega _q}}(n_q^+n_q^{})\mathrm{ln}\left\|{\displaystyle \frac{kq}{\mathrm{\Omega }\sqrt{k^2q^2}}}\right\|,`$ (A34) $`H^2(\mathrm{\Omega }^2)`$ $`=`$ $`{\displaystyle \frac{3}{8(2\pi )^4}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{kdk}{\omega _k}}(n_k^++n_k^{}){\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{qdq}{\omega _q}}(n_q^++n_q^{})\mathrm{ln}\left\|{\displaystyle \frac{k+q}{kq}}\right\|.`$ (A35) Extracting the high-temperature limit of $`H^1(\mathrm{\Omega }^2)`$ gives, after using the symmetry of the integrand under $`kq`$ to restrict the range of the $`q`$ integration and making the changes of variable $`\mathrm{\Omega }\xi =\omega _k,\mathrm{\Omega }\eta =\omega _q`$, $$H^1(\mathrm{\Omega }^2)=\frac{\mathrm{\Omega }^2}{4(2\pi )^4}_1^{\mathrm{}}𝑑\xi (n_\xi ^+n_\xi ^{})_1^\xi 𝑑\eta (n_\eta ^+n_\eta ^{})\mathrm{ln}\left\|\frac{(\xi ^21)(\eta ^21)}{\xi ^2\eta ^2}\right\|,$$ (A36) where $`n_\xi ^\pm =1/(e^{a(\xi \pm r)}1).`$ After making the successive approximationsThe second approximation in (A37) is strictly valid only for $`\|a\xi \|1`$; however, its validity is assured by the convergence of the resulting integral and has been checked numerically. $$n_\xi ^+n_\xi ^{}2ra\frac{e^{a\xi }}{(e^{a\xi }1)^2}\frac{2r}{a\xi ^2}$$ (A37) in the integrand, we have, for $`a1`$, $$H^1(\mathrm{\Omega }^2)=\frac{T^2r^2}{(2\pi )^4}(\mathrm{ln}2\frac{1}{2}).$$ (A38) The calculation of $`H^2(\mathrm{\Omega }^2)`$ parallels the case for $`\mu =0`$, leading to $$H^2(\mathrm{\Omega }^2)=\frac{3T^2}{64\pi ^2}\left[\frac{1}{2}\mathrm{ln}\left(\frac{\mathrm{\Omega }^2\mu ^2}{T^2}\right)1.50699\right].$$ (A39)
warning/0005/hep-ph0005320.html	ar5iv	text	# Δ⁢𝐺⁢(𝑥,𝜇²) from jet and prompt photon production at Rhic ## I Introduction One of the most remarkable features of high-energy hadron collisions is certainly their complexity. Particle multiplicities in single events are frequently of the order $`𝒪(10^2`$-$`10^3)`$ and will continue to rise with the advent of new colliders like Bnl-Rhic and Cern-Lhc. On the other hand, perturbative QCD (pQCD) can only deal with a limited number of final-state partons so that we have to ask, how these are related to experimentally measured distributions of hadrons and leptons. This is where event generators (EGs) come to the rescue. Heavy use of EGs is made during the whole life cycle of a high energy physics experiment for purposes such as estimating absolute production rates, studying the background-to-signal ratio, studying the detector acceptance, testing the reliability of event reconstruction during experimental analysis, etc. For collisions of polarized hadrons this list must be extended to also include studying the size of the hadronic double spin asymmetries of the signal and the background, determining the sensitivity to the polarized parton distribution which is to be measured, calculating the polarized acceptance, etc. The last point is especially important for fragmented detector systems with a small geometrical coverage like Phenix at Rhic (see, e.g., ). EGs for the simulation of collisions of longitudinally polarized particles must therefore include a correct treatment of particle helicity. In order to see how this goal can be achieved let us summarize the principle of current EGs. The generation process is subdivided into several steps where the first one deals with the hard partonic scattering. For nucleon-nucleon scattering, parton flavours and kinematics are selected according to the factorization formula for the unpolarized cross section, $$\mathrm{d}\sigma ^{(AB)}(P_A;P_B)=\underset{a,b}{}dx_Adx_Bf_a^A(x_A,\mu ^2)f_b^B(x_B,\mu ^2)d\widehat{\sigma }_{a,b}(x_AP_A,x_BP_B,\mu ).$$ (1) A, B (a, b) are the involved hadrons (partons), $`\mu `$ denotes the factorization and renormalization scale and the partonic cross section $`\mathrm{d}\widehat{\sigma }_{a,b}`$ is evaluated to leading order (LO) in the strong coupling $`\alpha _s(\mu )`$. Please note, that the momenta of the colliding partons are so far parallel to the momenta of the colliding beams. Of course this picture is very simplified and is significantly altered by radiative corrections which are specific to every hard partonic reaction. However, certain parts of these radiative corrections, consisting of the collinear poles of each parton splitting, are universal and can be used to formulate a parton shower (PS) algorithm which can be applied to every partonic subprocess . In this way an approximation of the complete perturbative series is obtained which becomes exact in the limit of vanishing decay angles and fails for large angles. The parton showers can be somewhat arbitrarily subdivided into an initial state shower (ISS) and a final state shower (FSS). The former generates all particles which directly ‘connect’ the two initial colliding partons with the ones that finally participate in the hard partonic reaction. It mimics the DGLAP evolution of parton densities in normal perturbative calculations but there is also an important difference: each parton branching also involves transverse momentum which in general leads to a Lorentz boost and rotation of the partonic scattering system. Thus, the ISS changes the kinematics of the partons (for details see ). Finally, the numerous partons produced by the PS fragment into stable and unstable hadrons which are then allowed to decay. Thus, EGs provide complete events which are in principle as detailed as real events that could be measured with a perfect detector. In an ideal case, they would also show the same fluctuation properties. An approximate treatment of polarization can be achieved relatively easy and has so far been used by the members of the Rhic Spin Program to prepare the experiments : * Generate unpolarized events with the unpolarized EG Pythia . Extract the kinematics ($`x_A`$, $`x_B`$, $`\widehat{s}`$, $`\widehat{t}`$, $`\widehat{u}`$, $`\mu `$) and the flavours of the hard partonic interaction. The partonic Mandelstam variables follow the usual definitions $$\widehat{s}=(x_AP_A+x_BP_B)^2,\widehat{t}=(x_AP_Ak_1)^2,\mathrm{and}\widehat{u}=(x_AP_Ak_2)^2,$$ (2) where $`k_{1,2}^\mu `$ are the momenta of the outgoing partons. * For each observable plot two histograms. The first histogram with the weight set to 1 gives the unpolarized cross section. For the second one the weight is taken to be the hadronic double spin asymmetry $$A_{LL}^{(AB)}(x_A,x_B,\widehat{s},\widehat{t},\widehat{u},\mu ^2)=\frac{\mathrm{\Delta }f_a^A(x_A,\mu )}{f_a^A(x_A,\mu )}\frac{\mathrm{\Delta }f_b^B(x_B,\mu )}{f_b^B(x_B,\mu )}\frac{\mathrm{\Delta }\widehat{\sigma }_{a,b}(\widehat{s},\widehat{t},\widehat{u},\mu ^2)}{\widehat{\sigma }_{a,b}(\widehat{s},\widehat{t},\widehat{u},\mu ^2)}.$$ (3) The formulae for the partonic double spin asymmetries $`\mathrm{\Delta }\widehat{\sigma }_{a,b}/\widehat{\sigma }_{a,b}`$ for the processes relevant for this work can be found in Table I. * After the event generation is completed, divide the second histogram by the first one to obtain the hadronic double spin asymmetry. Finally, normalize the first histogram appropriately. We shall call this prescription the method of asymmetry weights (MAW) in the remainder of this paper. It is clearly only an approximation to the ’true answer’ within the EG model because it does not describe how parton showers depend on the parton helicities. A simple way to overcome this deficit was presented by the authors of the EG Sphinx which is a polarized version of Pythia. It uses hadronic and partonic cross sections which depend on the helicities of the incoming hadrons $`H_{A,B}`$ and partons $`h_{a,b}`$ and which are summed over the polarization of the outgoing particles. The kinematics of the partonic reaction is therefore generated according to $`\mathrm{d}\sigma ^{(AB)}(P_A,H_A;P_B,H_B)`$ $`=`$ $`{\displaystyle \underset{a,h_a,b,h_b}{}}{\displaystyle dx_Adx_Bf_{a,h_a}^{A,H_A}(x_A,\mu ^2)f_{b,h_b}^{B,H_B}(x_B,\mu ^2)d\widehat{\sigma }_{ah_a,bh_b}(x_AP_A,x_BP_B,\mu ^2)},`$ (4) $`f_{a,h_a}^{A,H_A}(x_A,\mu ^2)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[f_a^A(x_A,\mu ^2)+\delta _{ϵ(H_A),ϵ(h_a)}\mathrm{\Delta }f_a^A(x_A,\mu ^2)\delta _{ϵ(H_A),ϵ(h_a)}\mathrm{\Delta }f_a^A(x_A,\mu ^2)\right],`$ (5) where $`ϵ(x)`$ is the sign function . Similarly, the helicity of all partons participating in the ISS is chosen according to the helicity dependent DGLAP splitting functions $$P_{bh_b,ah_a}(z)=\frac{1}{2}\left[P_{ba}(z)+\delta _{ϵ(h_b),ϵ(h_a)}\mathrm{\Delta }P_{ba}(z)\delta _{ϵ(h_b),ϵ(h_a)}\mathrm{\Delta }P_{ba}(z)\right]$$ (6) instead of using the usual unpolarized $`P_{ba}(z)`$ as in Pythia. Following this procedure one directly gets the cross sections for different helicity configurations of the colliding hadrons. They can be combined according to $$A_{LL}^{(AB)}=\frac{\mathrm{d}\sigma ^{(AB)}(P_A,+;P_B,+)\mathrm{d}\sigma ^{(AB)}(P_A,+;P_B,)}{\mathrm{d}\sigma ^{(AB)}(P_A,;P_B,+)+\mathrm{d}\sigma ^{(AB)}(P_A,+;P_B,)}$$ (7) to obtain the hadronic double spin asymmetry. It is clear, that one needs events as generated by Sphinx to test the experimental analysis software. Sphinx has, however, a severe disadvantage in comparison to the MAW. Longitudinal double spin asymmetries in $`\stackrel{}{p}\stackrel{}{p}`$-collisions are frequently of the order of 1% which means that $`10^6`$ accepted events per bin are necessary to achieve a relative statistical Monte Carlo (MC) error of 10%. The MAW needs much less computation time since it calculates the asymmetry directly and not by computing the relative difference between two similarly large numbers. Furthermore, it allows for the simultaneous usage of multiple parameterizations of polarized parton distributions because they are irrelevant for the generation of the event itself. Depending on the size of the asymmetry computation times easily differ by a factor of the order $`𝒪(10^2)`$. However, this speed advantage is useless if both methods yield largely different results for the hadronic double spin asymmetries. So far, no systematic comparison using realistic experimental cuts for the relevant observables has been performed for the Rhic Spin Project or any other experiment according to the authors’ knowledge. Since this question is so important we have done the necessary simulations and present our results in this paper. We would also like to stress that rates and asymmetries calculated with Sphinx should not be regarded as the ultimate truth. EGs possess some serious shortcomings which their users should always be aware of. The foremost problem stems from the fact that LO expressions for the cross sections and splitting functions are used throughout. The resulting large scale dependences prohibit reliable predictions of absolute cross sections or rates. Changes in some of the numerous necessary parameters of the PS model also may have quite strong effects on counting rates. Being aware of this issue we should try to clarify if the polarized EGs nevertheless yield a good description of nature. Unfortunately, this is currently impossible because so far no $`\stackrel{}{p}\stackrel{}{p}`$-collisions have been studied. We can only compare with the unpolarized results for the wide range of available $`pp`$\- and $`p\overline{p}`$-data but even this task is non-trivial without access to the analysis software of the various collaborations. Since we are more interested in the polarized sector anyway, we decided to rather study the uncertainty of theoretical predictions by comparing the EG results with next-to-leading order (NLO) QCD calculations. In some sense, both methods are complementary because they take different parts of the perturbative series into account. Furthermore, NLO predictions depend on a minimal set of parameters which includes only the intrinsic QCD scale $`\mathrm{\Lambda }_{\mathrm{QCD}}`$, the renormalization, factorization and in some cases the fragmentation scales. Usually this dependence is much smaller than in LO calculations effectively giving NLO QCD predictive power with regards to absolute rates. In some cases, agreement with experiment has been quite spectacular. E.g., the NLO calculation of the transverse energy spectrum of inclusive 1-jet production at Tevatron at $`\sqrt{S}=1.8`$ GeV agrees very well with data even though the unpolarized cross section varies over a range of six orders in magnitude . Unfortunately, NLO calculations tend to get unrealiable if highly restrictive experimental acceptances are involved. We assume that NLO QCD will also do a good job for polarized jet production at Rhic which was shown to be very sensitive to the polarized gluon distribution $`\mathrm{\Delta }g(x,\mu ^2)`$ , the prime target of most measurements at Rhic and take the differences when compared to EG results as estimate for the theoretical systematic error. Originally, prompt photon production has been promoted as the most useful process for the measurement of $`\mathrm{\Delta }g(x,\mu ^2)`$ but currently the situation in the unpolarized sector is far from satisfactory. Aside from the fact that NLO QCD cannot fit all available data sets with one set of scales in the fixed-target energy range, there also seem to be inconsistencies between some data sets . Under very lucky circumstances these problems may happen to cancel when taking ratios of polarized and unpolarized cross sections but of course we should not count on it. Whatever the solution to the current problems will be, NLO QCD will always be an integral part of it. For this reason, we also decided to not only compare jet observables but to also undertake a systematic comparison of prompt photon observables calculated with EGs and NLO QCD. For this purpose we use the parton generators of which use MC techniques to integrate NLO cross sections and permit the implementation of simple cuts on the final state. After discussing the necessary technical preliminaries we will present our results for jet observables in section 2 whereas section 3 will deal with prompt photon production only. The final section summarizes our findings and contains the conclusions. ## II Production of QCD-jets The authors of Ref. have already studied 1-jet and 2-jet observables in longitudinally polarized $`\stackrel{}{p}\stackrel{}{p}`$-collisions at the maximal Rhic cm-energy of $`\sqrt{S}=500`$ GeV using their NLO QCD parton generator. Star is the only available detector system at Rhic which is able to perform such measurements. After installation of the endcap its acceptance region will cover $`1<\eta <2`$ in pseudorapidity and $`0<\varphi <2\pi `$ in azimuthal angle. The main results of are that the perturbative series of the polarized cross section is under control, that parameterizations with different $`\mathrm{\Delta }g(x,\mu ^2)`$ yield very distinctive asymmetries and that the statistical experimental error will be sufficiently small to allow for a precise measurement from which the polarized gluon distribution can be determined. Therefore, is an ideal basis for our studies and we can use the same observables without any worries about perturbative stability. Before we give a detailed description of them let us make a few remarks on the choice of the parton distributions and the QCD running coupling constant. Nowadays, polarized inclusive lepton-nucleon scattering still has the by far largest impact on the determination of the polarized parton distributions. Unfortunately, the current data are not precise enough and cover a too small $`x`$-$`Q^2`$-range to give any useful restriction on $`\mathrm{\Delta }g(x,\mu ^2)`$. As a consequence, LO and NLO fits based on the same assumptions and restrictions may give widely different results for the polarized gluon distribution, i.e. the ratio $$K_{\mathrm{\Delta }g}(x,\mu ^2)=\frac{\mathrm{\Delta }g^{\mathrm{NLO}}(x,\mu ^2)}{\mathrm{\Delta }g^{\mathrm{LO}}(x,\mu ^2)}$$ (8) can deviate strongly from 1. According to Fig. 1 this is indeed the case for two of the three sets of polarized parton distributions we use in our studies. Although the EGs are only based on LO QCD we would therefore be ill advised to use them in combination with LO parton distributions. Any agreement with NLO QCD results would be purely accidental, especially in view of the sensitivity of jet production towards $`\mathrm{\Delta }g(x,\mu ^2)`$. Accordingly, we also use the NLO expression for the strong coupling $`\alpha _s(\mu ^2)`$ throughout to ensure that the EG without PS and the lowest order perturbative calculation will give the same results. The unpolarized GRV as well as the polarized GRSV and DSS parameterizations we use are already quite old. Nevertheless, they still fulfill all our needs since we are interested in comparing different methods of calculation rather than making precise predictions. It is more important here to ensure that for each polarized set we use the correct unpolarized set of distributions which in all cases is GRV. Otherwise Sphinx will not give correct results because it makes use of the unpolarized and polarized set simultaneously, so that only one value can be given for $`\mathrm{\Lambda }_{\mathrm{QCD}}`$. The three selected sets of $`\mathrm{\Delta }f(x,\mu ^2)`$ pretty much cover the whole range of possible values for $`\mathrm{\Delta }g(x,\mu ^2)`$ with GRSV $`g_{\mathrm{max}}`$ giving the largest and DSS3 the smallest values. In all cases the scales were set equal to the average transverse momentum of the partons leaving the hard partonic interaction: $$\mu =\frac{1}{n}\underset{i=1}{\overset{n}{}}\|k_{T,i}\|.$$ (9) ### A The observables We already mentioned that the number of final-state partons is much larger in the EG model than in perturbative calculations even though both intend to describe the same physics. For our studies we should therefore pick a jet-clustering algorithm which depends only weakly on particle multiplicities. The Ellis-Soper $`k_T`$-algorithm with the preferred choice of $`R=1`$ suits our needs. For each event it gives the number of reconstructed jets as well as their transverse energies $`E_T`$, pseudorapidities $`\eta `$ and azimuthal direction $`\varphi `$. From these we construct four observables<sup></sup><sup></sup>* We actually studied more observables but they did not yield any new insights. which are defined as follows: * the $`E_T`$-dependent 1-jet inclusive cross section $`\mathrm{d}(\mathrm{\Delta })\sigma /\mathrm{d}E_T`$, 14 GeV$`<E_T<50`$ GeV; * the $`\eta `$-dependent 1-jet inclusive cross section $`\mathrm{d}(\mathrm{\Delta })\sigma /\mathrm{d}\eta `$, $`0<\eta <2`$, $`E_T>15`$ GeV. Since the cross section is symmetric in $`\eta `$ we restrict ourselves to the region of positive pseudorapidity. For all 2-jet inclusive observables we denote the hardest and second hardest jet by $`J_1`$ and $`J_2`$ and demand $$E_{T,J_1}>15\mathrm{GeV},E_{T,J_2}>10\mathrm{GeV},\|\eta _{J_1}\|<1,\|\eta _{J_2}\|<1.$$ (10) The asymmetric cut on the transverse energies reduces contributions from configurations in which the two leading jets are exactly back-to-back and which cannot be calculated with finite order perturbation theory . All events fulfilling the requirement (10) are considered for * the distribution of the rapidity difference $`\mathrm{d}(\mathrm{\Delta })\sigma /\mathrm{d}\mathrm{\Delta }\eta `$, $`0<\mathrm{\Delta }\eta <2`$, $`\mathrm{\Delta }\eta =\eta _{J_1}\eta _{J_2}`$; * the distribution of the invariant jet mass $`M_{JJ}`$ defined by $`\mathrm{d}(\mathrm{\Delta })\sigma /\mathrm{d}M_{JJ}`$, $`M_{JJ}<100`$ GeV, $`M_{JJ}^2=(p_{J_1}+p_{J_2})^2`$. ### B Comparison of the EG methods Event generators are very complex programs which often consist of several ten thousand lines of code and are equipped with a large number of options. In order to avoid errors we should try to reuse as much program code as possible for our comparisons of EG methods and NLO calculations. This can be easily achieved, because Sphinx can also be operated in an unpolarized mode, in which it corresponds to Pythia 5.6 and only needs to be extended with a function that calculates the asymmetry weights of Table I. Furthermore, the EG and the NLO parton generators use the same code to evaluate the parton distributions and the strong coupling constant. The absolute statistical MC error of the asymmetries calculated with Sphinx is only determined by the number of accepted events in a certain bin and is therefore independent of the actual size of the asymmetries. In order to maximize the precision and reduce the required computer time we should use parton densities which yield large asymmetries, i.e. GRSV $`g_{\mathrm{max}}`$ is the preferred set. The artifical set $`\mathrm{\Delta }f(x,\mu ^2)=f(x,\mu ^2)`$ generally yields even larger asymmetries but cannot be used within the PS model since it is inconsistent with DGLAP evolution. By running Sphinx and Pythia with PS switched off, we made sure that the implementation of the asymmetry weights is indeed correct. The relative deviations between asymmetries calculated with Sphinx and the MAW did not exceed 2% and were within the statistical MC accuracy. We want to stress, that during all EG runs the hadronization and intrinsic partonic transverse momentum were switched off so that the results can be better compared to NLO calculations. Fig. 2 shows the results for the (un)polarized cross sections, asymmetries and the relative deviations $$R_\sigma =\frac{\mathrm{d}\sigma _{\mathrm{MAW}}\mathrm{d}\sigma _{\mathrm{Sphinx}}}{\mathrm{d}\sigma _{\mathrm{MAW}}},R_A=\frac{A_{LL,\mathrm{MAW}}A_{LL,\mathrm{Sphinx}}}{A_{LL,\mathrm{MAW}}}$$ (11) between the MAW ($`810^6`$ generated events) and Sphinx ($`410^7`$ generated events) with the PS switched on. Let us concentrate on the right column, in which the full black line represents $`R_\sigma `$. In all cases it is either compatible with or very close to zero so that we proved once more that the spin averaging of Sphinx indeed works. The dashed crosses represent $`R_A`$ calculated with polarized treatment switched on in the ISS of Sphinx. For all four observables systematic deviations show up. The asymmetries for the two 1-jet observables calculated with the MAW are generally about 5% larger than the ones which are calculated with Sphinx. The differences are especially strong in the large $`\eta `$-region, where they reach 20%. However, these findings cannot be generalized to all observables. For the $`\mathrm{\Delta }\eta `$-distributions and the region of small invariant mass the situation is reversed. Within the limited accumulated MC statistics, both methods seem to yield the same asymmetries for large $`M_{JJ}`$ but even here we can note a tendency for systematic differences towards positive $`R_A`$. In order to pin down the origin of the large deviations we also carried out an analysis of events generated by Sphinx for which the polarized treatment in the ISS was switched off (dotted crosses in Fig. 2). From the fact, that $`R_A(\eta )`$ is now compatible with zero for large pseudorapidity we conclude, that the largest part of the deviations is indeed caused by the missing treatment of parton helicity within the MAW approach. However this cannot be the full story because $`R_A`$ is not compatible with zero for all other observables. But keeping in mind that the ISS changes the kinematics of the hard partonic interaction, that within the MAW the ISS is applied to unpolarized events whereas within Sphinx it is applied to polarized events with different configurations of proton helicities, it is clear that the ISS can slightly modify the asymmetry even though polarized treatment is switched off. Unfortunately we were unable to check if the same results also hold true for the two other sets of polarized parton distributions. As we will see in Fig. 4, the set GRSV std. (DSS3) yields asymmetries which are smaller by a factor of five (twenty) so that we would have had to generate $`110^9`$ ($`1.610^{10}`$) events with Sphinx to obtain the same relative statistical accuracy. This was simply beyond the means of our computer equipment. ### C Comparison of the MAW with NLO QCD predictions In view of the fact that parton helicity plays a significant role in the ISS we would have liked to compare Sphinx with NLO QCD. Unfortunately, this was impossible for the same exact reasons as were given at the end of the previous paragraph. We resort to comparing the results of the MAW with NLO QCD predictions so that several different parameterizations of polarized parton densities can be used. In order to find a common basis we again start out by comparing Born-results, i.e. observables calculated with LO partonic cross sections, NLO parton distribution, and a NLO expression for the running strong coupling. For this purpose, Fig. 3 shows the relative deviations of the unpolarized cross sections and the asymmetries defined by $$R_\sigma \frac{\mathrm{d}\sigma _{(\mathrm{N})\mathrm{LO}\mathrm{QCD}}\mathrm{d}\sigma _{\mathrm{MAW}}}{\mathrm{d}\sigma _{(\mathrm{N})\mathrm{LO}\mathrm{QCD}}},R_A\frac{A_{LL,(\mathrm{N})\mathrm{LO}\mathrm{QCD}}A_{LL,\mathrm{MAW}}}{A_{LL,(\mathrm{N})\mathrm{LO}\mathrm{QCD}}}.$$ (12) Only the error bars of $`R_\sigma `$ are plotted for the sake of clarity, as the error bars of $`R_A`$ can be estimated from the statistical fluctuation between two neighbouring bins. The unpolarized cross sections and the asymmetries for both GRSV sets mainly show a satisfactory, but not perfect agreement for the bins with high statistics. This is unfortunately not the case for the DSS3 set – even for bins with good statistics deviations of 5% show up although both programs use the same routines to evaluate the parton distributions, as was already mentioned before. The reason for these problems is still under investigation but we assume, that it is most likely caused by the different treatment of colour interference terms (c.f. ) and the fact, that the relatively small asymmetries of the large number of contributing partonic processes undergo a delicate cancellation process. We will soon see that PS or NLO corrections to the asymmetries will be by far stronger so that we can effectively neglect the just mentioned problems. In an ideal case, the results of NLO QCD and the MAW would just agree. This is of course very unlikely. Just by considering the scale dependence of the observables, which is much larger for the EG, we see that agreement can only be achieved for one special scale or even not at all, because the results might not only differ in their normalization but also in their shape. So it is not surprising, that the cross sections calculated with both methods deviate by up to 40% when using the standard choice of scales and the standard PS parameters (see Fig 4). However, the objects we are mainly interested in, namely the asymmetries, show a much better agreement. $`R_A`$ seldomly exceeds 20% for any of the three polarized sets and the agreement is about twice as ‘good’ for $`A_{LL}(M_{JJ})`$. The PS needs a number of parameters which can be freely varied within certain limits and have to be determined by fitting to data. We picked three of of them, listed in Table II, and tried to refit them using the NLO QCD calculation as ‘data set’. Due to the long necessary computation time an automated fitting procedure is unpractical and we resort to manually testing a few values. The $`E_T`$\- and $`\eta `$-dependent observables are the most inclusive ones so that their normalization should be fixed first, aside from improving the agreement of the asymmetries. Accordingly, the absolute rates predicted by the EG have to be reduced which can be achieved by reducing the activity of the ISS. Table II shows the default as well as the fitted values of the three parameters. Since the fitting was done manually with a very limited number of different values it is unlikely that the optimal set has been found. Nevertheless, a look at Fig. 5 shows that it completely fulfills its purpose: the unpolarized cross sections for the 1-jet inclusive observables as well as the asymmetries for the two GRSV sets now agree at a 10% level whereas the situation for the set DSS3 is not quite as good. On the other hand, the agreement for the 2-jet inclusive rates got much worse. In order to get a feeling for the size of the relative deviations we study the scale dependence of the NLO predictions. In Fig. 6 the standard choice of the scales (9) is varied by a factor of 0.5 and 2, resp., which results in the bands of cross section displayed in the left column of the figure. By dividing the minimal (maximal) polarized cross section by the maximal (minimal) unpolarized cross section we obtain the limits on the bands of asymmetries, which are shown in the right column. According to the figure, the optimized EG asymmetries are always within the bands. Aside from the normalization of the 2-jet inclusive cross sections, the EG results are therefore compatible with NLO QCD. ## III Prompt photon production Let us now turn to the production of prompt photons in $`\stackrel{}{p}\stackrel{}{p}`$-collisions. This process receives a lot of interest mainly due to fact that the quark-gluon fusion $`qgq\gamma `$ strongly dominates at LO for medium sized and large $`\mathrm{\Delta }g(x,\mu ^2)`$. Before the trouble with the fixed-target data emerged, the analysis therefore seemed to be especially easy and clean. At the large cm-energies accessible at Rhic another problem arises: the majority of photons are actually not produced in the hard partonic interaction but through pion-decay into two photons. These fragmentation photons can luckily be completely filtered out by an isolation prescription defined in Ref. . Unless the usual cone-isolation prescription, which limits the hadronic energy contained in a cone drawn around the photon direction in $`\eta `$-$`\varphi `$-space with opening angle $`\delta _0`$, Frixione’s prescription allows less hadronic energy the closer the hadrons are to the photon and no hadronic energy for exactly parallel hadrons. Since fragmentation is an exactly collinear process in QCD, this prescription effectively eliminates all fragmentation photons while at the same time the cancellation of soft and collinear divergences stays untouched. For our studies we choose $`\delta _0=0.7`$. ### A The observables Common experience with polarized prompt photon calculations shows that asymmetries tend to become smaller when the cm-energy of the experiment is increased because smaller $`x`$-values are probed. For our event generator studies we therefore choose $`\sqrt{S}=200`$ GeV instead of the maximal Rhic cm-energy which is beneficial for the precision of all results calculated with Sphinx. Our choice of observables follows closely the NLO QCD study of Ref. for the same reasons that were given in section II. The cuts were inspired by the geometrical coverage of Star. Frixione’s isolation criterion in combination with the ES jet algorithm yields the transverse momentum, rapidity and azimuthal angle of the photon and a certain number of jets for all events which pass the isolation cut. These are then used to calculate two inclusive observables: * the $`p_{T,\gamma }`$-dependent cross section $`\mathrm{d}(\mathrm{\Delta })\sigma /\mathrm{d}p_{T,\gamma }`$, 10 GeV$`<p_{T,\gamma }<`$50 GeV, $`\|\eta _\gamma \|<0.5`$; * the $`\eta _\gamma `$-dependent cross section $`\mathrm{d}(\mathrm{\Delta })\sigma /\mathrm{d}\eta _\gamma `$, $`0<\eta _\gamma <2`$, $`p_{T,\gamma }>10`$ GeV. Due to symmetry we again restrict ourselves to the positive $`\eta `$-region. Inclusive prompt photon observables have the advantage that counting rates are quite large. On the other hand, the kinematics of the hard partonic reaction can of course not be reconstructed. Thus, there is no one-to-one correspondence between the measured transverse momentum of the photon and the probed gluonic $`x`$ even in a LO analysis. This problem can be fixed if only events with one reconstructed jet are considered. Additionally, we require $`p_{T,\gamma }>10\mathrm{GeV},E_{T,J}>11\mathrm{GeV},1<\eta _{\gamma ,J}<2,\|\mathrm{\Delta }\varphi \|=\|\varphi _J\varphi _\gamma \|>{\displaystyle \frac{\pi }{2}}`$ (13) which effectively means that the jet must be located in the hemisphere opposite to the photon. Events with multiple jets are not considered here. Their simulation would require a NLO parton generator with so far unknown matrix elements of order $`𝒪(\alpha \alpha _s^3)`$ or higher. All events that pass the criteria (13) are considered for two $`\gamma J`$-observables: * the distribution of the rapidity difference $`\mathrm{d}(\mathrm{\Delta })\sigma /\mathrm{d}\mathrm{\Delta }\eta `$, $`0<\mathrm{\Delta }\eta <2.4`$, $`\mathrm{\Delta }\eta =\eta _J\eta _\gamma `$; * the distribution of the invariant mass of the $`\gamma J`$-pair $`\mathrm{d}(\mathrm{\Delta })\sigma /\mathrm{d}M_{\gamma J}`$, $`M_{\gamma J}<100`$ GeV, $`M_{\gamma J}^2=(p_\gamma +p_J)^2`$. ### B Comparison of the EG methods For the comparison of the EG methods we unfortunately had to restrict ourselves again to the set GRSV $`g_{\mathrm{max}}`$. The mandatory test of the implementation of the MAW did not yield any surprise. The asymmetries calculated from Born-cross sections with both methods agreed within the relative statistical accuracy of a few per cent. Before discussing the full results including parton showers we should be careful to check if Frixione’s isolation method also works for the string fragmentation model used in Pythia. Hadronization is not an exactly collinear process here, so that some fragmentation photons might actually fulfill the isolation criterion. This is indeed the case, but only for very few events so that we do not need to worry about them and can safely keep fragmentation switched off. Unlike for NLO QCD calculations, where all perturbatively produced photons originate directly in the hard interaction, we need to make a further distinction for perturbatively produced photons in the EG model. Aside from ‘direct’ photons, which are produced in the LO hard processes $`qgq\gamma `$ and $`q\overline{q}g\gamma `$, there is also the possibility that the photon is produced by emission from a quark line in the PS ‘before’ or ‘after’ one of the purely strong interacting reactions $`qqqq`$, $`qgqg`$, $`gggg`$, etc. occurs. We denote this second source of photons as ‘bremsstrahlung’ photons. Depending on their transverse momentum they make up 20-40% of all isolated photons at $`\sqrt{S}=200`$ GeV according to Fig. 7. The simulation of the bremsstrahlung contribution is very tedious because under normal circumstances only every thousandth event contains an isolated photon. This very low efficiency can be somewhat improved by multiplying the value of $`\alpha `$ by a small factor (e.g. 3) so that the ratio of QED- to QCD-emissions becomes larger without changing the event characteristics. But even then huge numbers of events are needed for decent statistical accuracies, especially when using Sphinx. We therefore only present the comparison of ’direct’ photons in Fig. 8, for which we produced $`1.710^8`$ events. We proceed along the same lines as for the jet studies in the last section. At first, the correct treatment of particle helicity in the ISS of Sphinx is switched on resulting in the dashed crosses in the right column. Then the ISS is operated in unpolarized mode resulting in the dotted crosses for $`R_A`$. The overall picture is very similar to Fig. 2 but less pronounced. Especially the $`p_{T,\gamma }`$\- and $`\mathrm{\Delta }\eta `$-dependent asymmetries show relative deviations of up to 10% which are again due to the missing treatment of particle helicity in the ISS within the MAW. Keeping in mind that we do not know if the inclusion of bremsstrahlung photons would change the overall picture, we can state that the usage of the MAW for prompt photon simulations appears to be better justified than for jet physics. We generated an additional $`2.810^8`$ events to study the bremsstrahlung contribution but unfortunately the insufficient statistical accuracy did not allow us to draw any conclusions. ### C Comparison of the MAW with NLO QCD predictions Before we proceed to the comparison with NLO QCD we should again try to establish a common basis for both methods by checking the Born-level results. Our tests showed that cross sections and asymmetries calculated with the MAW and the parton generator agreed on a 1% level so that everything is in perfect order. This is a strong indication for the correctness of our explanation for the deviations seen in Fig. 3 because colour interference does not exist in LO prompt photon production (see also Table I). In Fig. 9 the full EG results, i.e. including PS and including QED-bremsstrahlung, are compared with the NLO QCD predictions. Please note, that an overall $`K`$-factor of $`K=1.5`$ has been applied to all cross sections calculated with the MAW to bring its rates better in line with the NLO predictions. Of course this correction is far from being ideal but it can be justified by the large scale dependence of the EG results. Indeed, it is common practice to more or less trust the shape of cross sections predicted by an EG and to adjust the normalization within certain limits. A slightly larger $`K=1.67`$ would have led to a very good agreement for $`\mathrm{d}\sigma /\mathrm{d}\mathrm{\Delta }\eta `$ and $`\mathrm{d}\sigma /\mathrm{d}M_{\gamma J}`$ but for the price of larger discrepancies between the inclusive cross sections whose slopes are already way off. With the exception of the critical set DSS3, the asymmetries surprisingly deviate only by 5-10% on average! Due to the small gluon parameterization of DSS3, resulting in the small asymmetries of Fig. 9, the quark-gluon fusion process does not dominate as strongly and other processes become increasingly relevant. The situation becomes a lot more complicated with the effect that the discrepancies between asymmetries reach up to 25%. Unfortunately, the very long computation times prohibited any attempt to improve the agreement by repeating the fit of the shower parameters discussed in section II. This leaves us with the study of the scale dependence of the NLO predictions presented in Fig. 10. It turns out to be much larger than for jet production so that the different bands of asymmetries partly overlap. With the exception of two bins even the sizable deviations of DSS3 are within the limits. The asymmetries for the two GRSV parameterizations are located near the center of the corresponding bands. At least in this restricted sense, the asymmetries calculated with the MAW and NLO QCD are consistent with each other. ## IV Resume Our detailed comparisons between Sphinx and the method of asymmetry weights show, that the correct description of particle polarization in the initial state shower generally has to be implemented into any polarized event generator if one requires a precision of order 10%. For some of the studied observables we find even differences as large as 20% between the MAW and Sphinx. The most drastic deviations arise for the region of large pseudorapidity ($`\eta 2`$) of 1-jet inclusive observables. Luckily this region of phase-space will be irrelevant for the experiments at Rhic. The geometrical coverage of Star only reaches out to $`\eta =2`$ so that under normal circumstances the maximal pseudorapidity for the jet-axis will be $`\eta 1`$. So we are left with maximal deviations of approximately 10%, which are however still substantial. In view of the tremendous advantage in speed of the method of asymmetry weights we therefore recommend the following compromise: in a first step, all studies should be performed using the method of asymmetry weights. Once the optimal energies, experimental cuts, etc. have been found, Sphinx should be used to check the results. The comparisons between event generator and NLO QCD predictions are encouraging despite of the substantial discrepancies in the unpolarized sector. The longitudinal double spin asymmetries are the objects we are mainly interested in and the agreement is much better here. With very few exceptions, the deviations between the method of asymmetry weights and NLO QCD are smaller than the uncertainties of the NLO QCD predictions due to scale dependence. Both methods yield consistent results in this restricted sense. During the starting phase of the Rhic spin project these theoretical uncertainties on the asymmetries of approximately 10% are probably tolerable. However, in the long run more precise event generators will be needed. Their development can in our opinion not be achieved without substantial conceptual progress in event generator techniques. Most likely a fusion of NLO matrix elements with the parton shower methods is needed. The development of such a NLO EG has been a hot topic of scientific debate for several years already due to its urgency. So far, no working solution is available. In view of the substantial costs of the Rhic spin program and the scientific importance of a polarized NLO event generator it would seem advisable to us to also invest the necessary manpower and funds into the development of such a code, which, however, does not seem to happen. Thus, we might soon be in the situation that wonderfully rich data is available but cannot be appropriately analyzed. ## Acknowledgments This work was supported by the BMBF and the “Deutsche Forschungsgemeinschaft”. We thank the authors of Refs. for making their parton generator codes available to us. O. Martin is grateful to N. Saito for the invitation to the RIKEN-BNL research center workshop “Event Generator for Spin Physics” and thanks A. Kirchner, L. Niedermeier, S. Schäfer, E. Stein and M. Stratmann for lively discussions during our weekly meetings. ## Figure captions * The ratio $`K_{(\mathrm{\Delta })g}(x,\mu ^2)`$ defined in Eq. (8) for the unpolarized and polarized parton densities of our analysis in the relevant $`x`$-range. The line style coding of this figure will be used throughout the paper. * Comparison of jet production calculated with the MAW and Sphinx: (un)polarized cross sections and asymmetries calculated with GRSV $`g_{\mathrm{max}}`$ and their relative deviations $`R_{\sigma ,A}`$ defined in Eq. (11). Left and center columns: full lines represent the MAW results, dashed lines represent the Sphinx results. Right column: full lines represent $`R_\sigma `$, dashed (dotted) crosses represent $`R_A`$ calculated with (without) polarized treatment in the ISS of Sphinx. * The relative deviations $`R_{\sigma /A}`$ defined in Eq. (12) for the Born results calculated with the MAW and the parton generator. Line styles are chosen according to Fig. 1. The figure is based on $`210^7`$ generated events. * Comparison of jet production calculated with NLO QCD and the MAW with switched-on PS and standard PS parameters. Left column: MAW (NLO QCD) results are represented by full (dashed) lines. Center column: the line styles are chosen according to Fig. 1 with one line representing the NLO QCD and the other one the MAW result for each parametrization. They can be identified from the third column, where the line styles are again chosen according to Fig. 1, i.e. full crosses represent $`R_\sigma `$ and all other lines represent $`R_A`$ for the various parameterizations of polarized parton distributions. * same as Fig. 4 but with tuned PS parameters which were chosen to improve the agreement between both methods (see text). * Comparison of the optimized EG results (markers) with the bands of NLO QCD results which are obtained, when the scales are varied by a factor of 2 or 1/2, resp. Left column: full lines represent the band of possible unpolarized NLO QCD cross sections, dashed and dotted lines represent the bands of polarized NLO QCD cross sections and follow the same pattern as in Fig. 1. Right column: corresponding minimal and maximal NLO QCD predictions for the asymmetries. Line Styles are again chosen according to Fig. 1. The large scale dependence for 20 GeV$`<M_{JJ}<`$30 GeV is due to the fact that this bin receives no contribution to $`𝒪(\alpha _s^2)`$ so that all values are effectively only LO predictions. * Relative contribution of ‘direct’ and ‘bremsstrahlung’ photons to the total production rate of isolated photons. * Same as Fig. 2 but for the production of ‘direct’ photons (without ‘bremsstrahlung’ contribution). * Same as Fig. 4 but for the production of isolated photons (including the ‘bremsstrahlung contribution’). A $`K`$-factor of $`K=1.5`$ is applied to all EG cross sections. The figure is based on $`310^8`$ generated events. * Same as Fig. 6 but for the production of isolated photons. A $`K`$-factor of $`K=1.5`$ is applied to all EG cross sections.
warning/0005/nlin0005038.html	ar5iv	text	# Renormalization of Quantum Anosov Maps: Reduction to Fixed Boundary Conditions ## Abstract A renormalization scheme is introduced to study quantum Anosov maps (QAMs) on a torus for general boundary conditions (BCs), whose number ($`k`$) is always finite. It is shown that the quasienergy eigenvalue problem of a QAM for all $`k`$ BCs is exactly equivalent to that of the renormalized QAM (with Planck’s constant $`\mathrm{}^{}=\mathrm{}/k`$) at some fixed BCs that can be of four types. The quantum cat maps are, up to time reversal, fixed points of the renormalization transformation. Several results at fixed BCs, in particular the existence of a complete basis of “crystalline” eigenstates in a classical limit, can then be derived and understood in a simple and transparent way in the general-BCs framework. Nonintegrable systems whose dynamics can be reduced to a 2D torus in phase space have attracted much attention in the quantum-chaos literature. When quantizing such a system on the torus, the admissible quantum states must satisfy proper boundary conditions (BCs), i.e., they have to be periodic on the torus up to constant phase factors specified by a Bloch wave vector $`𝐰`$. If the Hamiltonian of the system is periodic in phase space, e.g., the kicked Harper model , its classical dynamics can be reduced to the toral phase space of one unit cell of periodicity, and all Bloch wave vectors $`𝐰`$ in some Brillouin zone (BZ) are allowed. The sensitivity of the eigenstates to continuous variation of $`𝐰`$ in the BZ may be characterized topologically by a Chern index and appears to be strong for eigenstates spread over the chaotic region and weak for eigenstates localized on stability islands. In general, however, the Hamiltonian of a system whose dynamics can be reduced to a torus is not periodic in phase space. An extreme case is that of the purely chaotic, Anosov “cat maps” , whose Hamiltonians are quadratic in the phase-space variables . When quantizing these systems, it turns out that only a finite set of $`𝐰`$’s in the BZ is allowed , see Eq. (2) below, and this set increases with increasing chaotic instability. However, due to the lack of 2D continuity in this set, the sensitivity to variations in $`𝐰`$ cannot be characterized by a Chern index . A large class of cat maps was first quantized by Hannay and Berry for $`𝐰=\mathrm{𝟎}`$ (strict periodicity on the torus) and general BCs were considered in subsequent works . An important recent work has shown the existence of eigenstates whose Husimi density is periodic on a lattice in the torus and whose Husimi zeros form a crystal, for all Planck’s constants in a sequence tending to zero. These “crystalline” eigenstates are the first example of exactly equidistributed eigenstates of quantized chaotic systems in a classical limit. An atypical feature of quantum cat maps is their high spectral degeneracy, which increases in the classical limit . Typical spectral properties, fitting generic eigenvalue statistics, are already found by quantizing torus maps that are very slight perturbations of the cat maps . According to Anosov’s theorem , these maps have essentially the same classical dynamics, in particular they are purely chaotic, as the unperturbed cat maps, which are structurally stable. This ceases to be the case for larger perturbations that cause bifurcations generating elliptic islands and thus a mixed phase space . However, the quantum BCs for a perturbed cat map are the same as those of the unperturbed cat map, independently of the size of the perturbation , see Eq. (2) below. Because of this reason and to simplify terms, general perturbed cat maps are referred to as Anosov maps in this Letter. The importance of these maps is in that they may be viewed as generic torus maps on the basis of a general expression for a smooth torus map derived recently , see below. In this Letter, we introduce a renormalization scheme to study quantum Anosov maps (QAMs) for general BCs. We show that the eigenvalue problem of a QAM for all BCs, determining the full quantum dynamics, can be exactly reduced to a problem at fixed BCs. More specifically, consider a QAM given by the evolution operator $`\widehat{U}`$ quantizing a classical Anosov map. Quantization on a torus requires a Planck’s constant $`\mathrm{}`$ to satisfy $`2\pi /\mathrm{}=p`$, an integer. The finite number of BCs is denoted by $`k`$, which depends only on the classical unperturbed cat map. We define a renormalization transformation $``$ of $`\widehat{U}`$ generating a QAM $`\widehat{U}^{}=(\widehat{U})`$ on the same torus. The number of BCs for $`\widehat{U}^{}`$ is also $`k`$ and $`\widehat{U}^{}`$ is associated with a renormalized Planck’s constant $`\mathrm{}^{}=\mathrm{}/k`$. We then show that the quasienergy eigenvalue problem for $`\widehat{U}`$ for all $`k`$ BCs is equivalent, by a unitary transformation accompanied by a scaling of variables, to that for $`\widehat{U}^{}`$ at some fixed BCs. The latter can be of four types, corresponding to the four “symmetry” points in the renormalized BZ (see Table 1 below). The quantum cat maps are always fixed points of $``$ in the reflection-hyperbolic case. In the ordinary-hyperbolic case, they are fixed points of $``$ only if some conditions are satisfied. Otherwise, they are fixed points of $``$ accompanied by time reversal. Thus, for general $`\widehat{U}`$, the QAMs $`\widehat{U}^{(n)}=^n(\widehat{U})`$ represent perturbations of a given quantum cat map, or its time inverse, in a classical limit, $`\mathrm{}^{(n)}=\mathrm{}/k^n0`$. The spectrum and eigenstates of $`\widehat{U}`$, $`\widehat{U}^{}`$,…, $`\widehat{U}^{(n1)}`$ for all $`k`$ BCs and $`n>1`$ can be fully reproduced from those of $`\widehat{U}^{(n)}`$ at fixed BCs corresponding to either strict periodicity for $`kp`$ even or antiperiodicity for $`kp`$ odd. As a first application of all this scheme, we show that some results at fixed BCs can be derived and understood in a simple and transparent way in the general-BCs framework. In particular, it is easy to show in this framework that the eigenstates of $`\widehat{U}^{(n)}`$ at the fixed BCs are all “crystalline”. Unlike general crystalline eigenstates of the quantum cat maps , the crystal features persist for arbitrarily large perturbations. We denote by $`𝐳=(u,v)`$ the phase-space variables, $`[\widehat{u},\widehat{v}]=i\mathrm{}`$, and we assume that the classical dynamics can be reduced to a $`2\pi \times 2\pi `$ torus $`T^2`$, where it is described by an Anosov map $`M`$. In general, a smooth map $`M`$ on $`T^2`$ can be written as $`𝐳M(𝐳)=\overline{M}(𝐳)`$ mod$`2\pi `$, where the “lifted” map $`\overline{M}`$, defined on the entire phase plane $`(u,v)`$, can be expressed uniquely as the composition of two maps, $`\overline{M}=M_AM_1`$ . Here $`M_A`$ is a linear map, $`M_A(𝐳)=A𝐳`$, where $`A`$ is a $`2\times 2`$ integer matrix with $`det(A)=1`$; by “Anosov” we just mean that $`\|`$Tr$`(A)\|>2`$, a condition generically satisfied by $`A`$. The map $`M_1`$ is defined by $`M_1(𝐳)=𝐳+𝐅(𝐳)`$, where $`𝐅(𝐳)`$ is a $`2\pi `$-periodic vector function of $`𝐳`$. For $`𝐅(𝐳)=\mathrm{𝟎}`$, $`M`$ is a “cat map”. The quantization of $`\overline{M}`$ is the unitary operator $`\widehat{U}=\widehat{U}_A\widehat{U}_1`$, where $`\widehat{U}_A`$ and $`\widehat{U}_1`$ are the quantizations of $`M_A`$ and $`M_1`$, respectively . In the $`u`$ representation, $$u_2\|\widehat{U}_A\|u_1_{\mathrm{}}=\left(\frac{1}{2\pi i\mathrm{}A_{1,2}}\right)^{1/2}\mathrm{exp}\left[\frac{i}{2\mathrm{}A_{1,2}}\left(A_{1,1}u_1^22u_1u_2+A_{2,2}u_2^2\right)\right].$$ (1) We shall assume that $`M_1`$ is the map for a Hamiltonian which is periodic in phase space with unit cell $`T^2`$. As shown in Ref. , this is the case if and only if $`_{T^2}𝐅(𝐳)𝑑𝐳=\mathrm{𝟎}`$. Then $`\widehat{U}_1`$ is the one-step evolution operator for the Weyl quantization of this Hamiltonian and is a periodic operator function $`\widehat{U}_1(\widehat{𝐳};\mathrm{})`$, representable by a well defined Fourier expansion. Quantization on a torus requires now that $`\mathrm{}=2\pi /p`$, $`p`$ integer. The QAM quantizing $`M`$ is $`\widehat{U}`$ restricted to act on the simultaneous eigenstates of the commuting phase-space translations on $`T^2`$, $`\widehat{D}_1=\mathrm{exp}(ip\widehat{u})`$ and $`\widehat{D}_2=\mathrm{exp}(ip\widehat{v})`$; the corresponding eigenvalues are $`\mathrm{exp}(ipw_1)`$ and $`\mathrm{exp}(ipw_2)`$, where $`(w_1,w_2)=𝐰`$ is the Bloch wave vector specifying the toral BCs . An eigenstate $`\mathrm{\Psi }_𝐰`$ of $`\widehat{D}_1`$ and $`\widehat{D}_2`$ can be an eigenstate of $`\widehat{U}`$ only for those values of $`𝐰`$ in the Brillouin zone (BZ: $`0w_1,w_2<2\pi /p`$ ) satisfying the equation $$A𝐰=𝐰+\pi 𝐲\text{mod }2\pi /p,$$ (2) where $`𝐲(A_{1,1}A_{1,2},A_{2,1}A_{2,2})`$. We write the general solution of Eq. (2) as follows: $$𝐰=(2\pi /p)B(𝐫+pE^1𝐲/2)\text{ mod }2\pi /p,$$ (3) where $`B=(AI)^1E`$, $`I`$ is the identity matrix, $`E`$ is an arbitrary integer matrix with $`det(E)=\pm 1`$, and $`𝐫`$ is an integer vector labeling the solutions. There are precisely $`k=\|det(B^1)\|=\|2`$Tr$`(A)\|`$ distinct vectors (3), as the number of fixed points of $`M_A`$ mod $`2\pi `$ , forming a lattice in the BZ. We denote by $`𝒮`$ the space of states $`\mathrm{\Psi }_𝐰`$ for all these $`k`$ values of $`𝐰`$. The subspace $`𝒮_𝐰`$ of $`𝒮`$ with a fixed value of $`𝐰`$ is $`p`$-dimensional, i.e., it is spanned by a basis of $`p`$ independent states , whose general expression in the $`u`$ representation is $$\mathrm{\Psi }_{b,𝐰}(u)=\underset{m=0}{\overset{p1}{}}\varphi _b(m;𝐰)\underset{l=\mathrm{}}{\overset{\mathrm{}}{}}e^{ilpw_2}\delta (uw_12\pi m/p2\pi l),$$ (4) where $`b=1,\mathrm{},p`$. Such a basis is formed, naturally, by the $`p`$ eigenstates of $`\widehat{U}`$ at fixed $`𝐰`$. We now introduce the torus $`T_B^2`$, defined by the vectors $`𝐑_j=2\pi k(B_{1,j},B_{2,j})`$, $`j=1,2`$; $`kB`$ has integer entries and $`T_B^2`$ contains precisely $`k`$ tori $`T^2`$. Since $`B^1AB=E^1AE`$ is an integer matrix, the superlattice with unit cell $`T_B^2`$ is invariant under $`A`$, so that the map $`\overline{M}`$ modulo $`T_B^2`$, denoted by $`M^{(B)}`$, is well defined. To continue, we shall first work out in detail the reflection-hyperbolic case of Tr$`(A)<2`$, choosing $`E=I`$, so that $`[A,B]=0`$ and $`det(B)>0`$. We shall then specify the changes to be made in the ordinary-hyperbolic case of Tr$`(A)>2`$. Let us perform the linear transformation of variables $$𝐳=kB𝐳^{}=\sqrt{k}C𝐳^{},$$ (5) where $`C=\sqrt{k}B`$ and $`𝐳^{}=(u^{},v^{})`$. Eq. (5) is the combination of a linear canonical transformation \[$`det(C)=+1`$, since $`det(B)>0`$\] with a scaling by a factor $`\sqrt{k}`$. Using $`[A,B]=0`$, it is easy to check that the map $`M^{(B)}`$ above is transformed by (5) into a map $`M^{}`$ on $`T^2`$ in the $`𝐳^{}`$ variables, $`M^{}(𝐳^{})=`$ $`\overline{M}^{}(𝐳^{})`$ mod $`2\pi `$, with $`\overline{M}^{}=M_A^{}M_1^{}`$ . Here $`M_A^{}(𝐳^{})=A𝐳^{}`$ and $`M_1^{}(𝐳^{})=(kB)^1M_1(𝐳=kB𝐳^{})`$. The renormalization transformation $`_c`$ in the classical case is then defined by $`M^{}(𝐳^{})=_c[M(𝐳)]`$ . Clearly, the cat maps, with $`𝐅(𝐳)=\mathrm{𝟎}`$, are fixed points of $`_c`$, i.e., $`M^{}(𝐳^{}=𝐳)=M(𝐳)`$. The quantum version of (5) implies that $`[\widehat{u}^{},\widehat{v}^{}]=i\mathrm{}^{}`$, where $`\mathrm{}^{}=\mathrm{}/k=2\pi /p^{}`$, $`p^{}kp`$. The quantization $`\widehat{U}^{}`$ of $`M^{}(𝐳^{})`$ is simply $`\widehat{U}`$ expressed in terms of $`(\widehat{𝐳}^{},\mathrm{}^{})`$ and acting on the space of the simultaneous eigenstates $`\mathrm{\Psi }_𝐰^{}^{}`$ of the phase-space translations on $`T^2`$ in the $`𝐳^{}`$ variables, $`\widehat{D}_1^{}=\mathrm{exp}(ip^{}\widehat{u}^{})`$ and $`\widehat{D}_2^{}=\mathrm{exp}(ip^{}\widehat{v}^{})`$. It is easy to show that $$\widehat{D}_{j+1}^{}=\widehat{D}(𝐑_j)=(1)^{pk^2B_{1,j}B_{2,j}}\widehat{D}_1^{kB_{2,j}}\widehat{D}_2^{kB_{1,j}}$$ (6) ($`j=1,2`$, $`\widehat{D}_3^{}\widehat{D}_1^{}`$), where $`\widehat{D}(𝐑_j)`$ are precisely the Weyl-Heisenberg phase-space translations on $`T_B^2`$. By expressing $`\widehat{U}=\widehat{U}_A\widehat{U}_1`$ in terms of $`(\widehat{𝐳}^{},\mathrm{}^{})`$, using also the theory of linear quantum canonical transformations , we obtain the expected result $`\widehat{U}^{}=\widehat{U}_A^{}\widehat{U}_1^{}`$. Here the $`u^{}`$ representation of $`\widehat{U}_A^{}`$ is given by (1) with $`u`$ and $`\mathrm{}`$ replaced by $`u^{}`$ and $`\mathrm{}^{}`$, respectively, and $`\widehat{U}_1^{}`$ is the operator function $`\widehat{U}_1^{}(\widehat{𝐳}^{};\mathrm{}^{})=\widehat{U}_1(\widehat{𝐳}=kB\widehat{𝐳}^{};\mathrm{}=k\mathrm{}^{})`$ \[the function $`\widehat{U}_1(\widehat{𝐳};\mathrm{})`$ was defined above\]. The renormalization transformation $``$ is then defined by $`\widehat{U}^{}=(\widehat{U})`$. The quantum cat maps ($`\widehat{U}=\widehat{U}_A`$) are fixed points of $``$, i.e., $`u_2^{}=u_2\|\widehat{U}^{}\|u_1^{}=u_1_\mathrm{}^{}=\mathrm{}=u_2\|\widehat{U}\|u_1_{\mathrm{}}`$. The space of states $`\mathrm{\Psi }_𝐰^{}^{}`$ for all the $`k`$ allowed values of $`𝐰^{}`$ will be denoted by $`𝒮^{}`$. We now show that $`𝒮`$ coincides with the subspace $`𝒮_{𝐰_0^{}}^{}`$ of $`𝒮^{}`$ associated with a particular value $`𝐰_0^{}`$. Thus, $`𝒮^{}`$ is $`k`$ times larger than $`𝒮`$. To show this, let us apply $`\widehat{D}_j^{}`$, $`j=1,2`$, on a state $`\mathrm{\Psi }_𝐰`$ of $`𝒮`$. Using (3), (6), and the fact that $`\widehat{D}_j\mathrm{\Psi }_𝐰=`$ $`\mathrm{exp}[i(1)^{j+1}pw_j]\mathrm{\Psi }_𝐰`$, $`j=1,2`$, we obtain $$\widehat{D}_j^{}\mathrm{\Psi }_𝐰=(1)^{pA_{j,j+1}}\mathrm{\Psi }_𝐰$$ (7) ($`A_{2,3}A_{2,1}`$). Rel. (7) means that all $`\mathrm{\Psi }_𝐰`$ in $`𝒮`$ are eigenstates of $`\widehat{D}_j^{}`$, $`j=1,2`$, associated with the same renormalized Bloch wave vector $`𝐰_0^{}`$. The latter can assume only four values, depending on the matrix $`A`$, see Table 1. $`k`$ even $`k`$ odd $`p`$ even $`𝐰_0^{}=\mathrm{𝟎}`$ $`𝐰_0^{}=\mathrm{𝟎}`$ $`p`$ odd $`𝐰_0^{}=(\frac{A_{1,2}\pi }{p^{}},\frac{A_{2,1}\pi }{p^{}})`$ mod $`\frac{2\pi }{p^{}}`$ $`𝐰_0^{}=(\frac{\pi }{p^{}},\frac{\pi }{p^{}})`$ Table 1. It is easy to show that $`𝐰_0^{}`$ is indeed an allowed value of $`𝐰^{}`$ in all four cases. Thus, $`𝒮_{𝐰_0^{}}^{}`$ includes $`𝒮`$, but since both $`𝒮_{𝐰_0^{}}^{}`$ and $`𝒮`$ are $`kp`$-dimensional, they coincide. This completes the proof. Now, by the definition above of $`\widehat{U}^{}`$, the restriction of $`\widehat{U}^{}`$ to $`𝒮_{𝐰_0^{}}^{}=𝒮`$ is just $`\widehat{U}`$. The $`kp`$ eigenstates of $`\widehat{U}`$ for all $`k`$ BCs are then precisely the $`p^{}`$ eigenstates of $`\widehat{U}^{}`$ associated with the value of $`𝐰_0^{}`$ in Table 1. When referred to $`\widehat{U}^{}`$, however, these eigenstates should be expressed in a representation based on the operator $`\widehat{𝐳}^{}`$. If the $`kp`$ eigenstates of $`\widehat{U}`$ are $`\mathrm{\Psi }_{b,𝐰}(u)`$ in the $`u`$ representation, see (4), their $`u^{}`$ representation will be obtained by applying to $`\mathrm{\Psi }_{b,𝐰}(u)`$ the unitary transformation corresponding to a linear canonical transformation with matrix $`C`$, after scaling $`u^{}`$ by a factor $`\sqrt{k}`$. The eigenstates of $`\widehat{U}^{}`$ for $`𝐰^{}=𝐰_0^{}`$ are thus given by $$\mathrm{\Psi }_{b^{},𝐰_0^{}}^{}(u^{})=\left(\frac{p}{4\pi ^2B_{1,2}}\right)^{1/2}\underset{\mathrm{}}{\overset{\mathrm{}}{}}du\mathrm{exp}\left[\frac{ip}{4\pi B_{1,2}}(kB_{1,1}u^{}{}_{}{}^{2}2u^{}u+B_{2,2}u^2)\right]\mathrm{\Psi }_{b,𝐰}(u),$$ (8) where $`b^{}=b^{}(b,𝐰)`$ takes precisely all its $`p^{}`$ values when $`b`$ and $`𝐰`$ take all their $`p`$ and $`k`$ values, respectively; conversely, $`\mathrm{\Psi }_{b,𝐰}(u)`$ can be fully reproduced from $`\mathrm{\Psi }_{b^{},𝐰_0^{}}^{}(u^{})`$ by inverting Rel. (8) and determining $`𝐰`$ by applying $`\widehat{D}_1`$ and $`\widehat{D}_2`$ on $`\mathrm{\Psi }_{b,𝐰}(u)`$. If the quasienergies of $`\widehat{U}`$ are $`\omega _b(𝐰)`$, those of $`\widehat{U}^{}`$ for $`𝐰^{}=𝐰_0^{}`$ are $`\omega _{b^{}(b,𝐰)}^{}(𝐰_0^{})=\omega _b(𝐰)`$. The latter relation and Rel. (8) show the equivalence between the quasienergy eigenvalue problem for $`\widehat{U}`$ for all $`k`$ BCs and that for $`\widehat{U}^{}`$ at the fixed BCs given by $`𝐰^{}=𝐰_0^{}`$. The case of Tr$`(A)>2`$ can be treated similarly only if the integer matrix $`E`$ in $`B=(AI)^1E`$ can be chosen so that $`[A,E]=0`$ and $`det(E)=1`$. Then one has again $`[A,B]=0`$ and $`det(C)=+1`$ in (5), leading to the same results as above. If $`A=K^{2l}`$, where $`K`$ is any integer matrix with $`det(K)=1`$ and $`l`$ is an integer, one can choose $`E=K`$. In general, an integer matrix $`E`$ having the properties above does not exist, and we then make the simple choice $`E_{1,1}=E_{2,2}=1`$, $`E_{1,2}=E_{2,1}=0`$, corresponding to time reversal. As a result, in all the expressions and equations involving the renormalized quantities, including in Table 1, $`A`$ is replaced by $`A^{}=E^1AE`$. Thus, the quantum cat maps are now fixed points of $``$ accompanied by time reversal: Given the eigenstates $`\mathrm{\Psi }_{b,𝐰}(u)`$ and quasienergies $`\omega _b(𝐰)`$ of $`\widehat{U}_A`$ for all BCs, its eigenstates and quasienergies for $`\mathrm{}^{}=2\pi /p^{}`$ and $`𝐰^{}=𝐰_0^{}`$ are, respectively, $`\mathrm{\Psi }_{b^{},𝐰_0^{}}^{}(u^{})`$ and $`\omega _{b^{}(b,𝐰)}^{}(𝐰_0^{})=\omega _b(𝐰)`$, where $`\mathrm{\Psi }_{b^{},𝐰_0^{}}^{}(u^{})`$ is given by (8). Time-reversal symmetry ($`A^{}=A^1`$) allows to choose $`\mathrm{\Psi }_{b^{},𝐰_0^{}}^{}(u^{})`$ as real, $`\mathrm{\Psi }_{b^{},𝐰_0^{}}^{}(u^{})=\mathrm{\Psi }_{b^{},𝐰_0^{}}^{}(u^{})`$. By iterating $``$, one obtains a sequence of QAMs $`\widehat{U}^{(n)}=^n(\widehat{U})`$ on $`T^2`$, associated with the Planck’s constants $`\mathrm{}^{(n)}=2\pi /p^{(n)}`$, $`p^{(n)}k^np`$. General $`\widehat{U}^{(n)}=\widehat{U}_A^{(n)}\widehat{U}_1^{(n)}`$ represent perturbations of the quantum cat map $`\widehat{U}_A`$, or its time inverse $`\widehat{U}_A^{}`$ (only for $`n`$ odd), in a classical limit, $`n\mathrm{}`$. The perturbation $`\widehat{U}_1^{(n)}(\widehat{𝐳}^{(n)};\mathrm{}^{(n)})`$ is periodic in $`\widehat{𝐳}^{(n)}`$ with a unit cell $`T_n^2=(kB)^nT^2`$, which is $`k^n`$ times smaller than $`T^2`$. In the generalization of Table 1 to $`n>1`$, $`p^{(n1)}=k^{n1}p`$ is always even when $`k`$ is even, so that $`𝐰_0^{(n)}`$ can take only two values: $`𝐰_0^{(n)}=\mathrm{𝟎}`$ (for $`kp`$ even) and $`𝐰_0^{(n)}=(\pi /p^{(n)},\pi /p^{(n)})`$ (for $`kp`$ odd), corresponding to strictly periodic and antiperiodic BCs, respectively. The eigenstates of $`\widehat{U}^{(n)}`$ for $`𝐰^{(n)}=𝐰_0^{(n)}`$ are connected with those of $`\widehat{U}^{(n1)}`$ for all $`k`$ BCs by a relation analogous to Rel. (8). The quasienergies are related by $`\omega _{b^{(n)}}^{(n)}(𝐰_0^{(n)})=\omega _{b^{(n1)}}^{(n1)}(𝐰^{(n1)})`$. Thus, the spectrum and eigenstates of $`\widehat{U},\mathrm{},\widehat{U}^{(n1)}`$ for all BCs can be fully reproduced from those of $`\widehat{U}^{(n)}`$ for $`𝐰^{(n)}=𝐰_0^{(n)}`$. We now show that all the eigenstates of $`\widehat{U}^{(n)}`$ for $`𝐰^{(n)}=𝐰_0^{(n)}`$ are “crystalline”. By construction, (8) are eigenstates of $`\widehat{D}_j=\widehat{D}^{}(𝐑_j^{})`$, $`j=1,2`$, where $`\widehat{D}^{}(𝐑_j^{})`$ are Weyl-Heisenberg translations in the $`(u^{},v^{})`$ phase space by vectors $`𝐑_j^{}`$ defining precisely the unit cell $`T_1^2=(kB)^1T^2`$, which is $`k`$ times smaller than $`T^2`$. Thus, the Husimi density of all the eigenstates of $`\widehat{U}^{}`$ for $`𝐰^{}=𝐰_0^{}`$ is exactly periodic on a lattice with unit cell $`T_1^2`$. The eigenstates are then crystalline : each cell contains $`p`$ of the $`kp`$ Husimi zeros in $`T^2`$. For $`n>1`$ and $`𝐰^{(n)}=𝐰_0^{(n)}`$, the $`k^np`$ eigenstates of $`\widehat{U}^{(n)}`$ can be grouped into $`n`$ sets: the $`l`$th set, $`l=1,\mathrm{},n`$, consists of $`k^{l1}p(k+\delta _{l,1}1)`$ crystalline eigenstates with unit cell $`T_{nl+1}^2`$ and $`k^{l1}p`$ Husimi zeros in each cell. Unlike general crystalline eigenstates of the quantum cat maps , the crystal features of all this complete basis of eigenstates persist under an arbitrarily large perturbation $`\widehat{U}_1^{(n)}`$ in $`\widehat{U}^{(n)}`$. Another interesting result at fixed BCs concerns the expansion (4) for the eigenstates (8) and can be easily derived in the case when the vectors (3) form a square lattice in the BZ, $`𝐰=2\pi (r_1,r_2)/(gp)`$, $`r_1,r_2=0,\mathrm{},g1`$, $`g`$ integer (see the conditions for this in note ). Since $`B=I/g`$ , the transformation (5) is simply $`𝐳=g𝐳^{}`$, and the eigenstates $`\mathrm{\Psi }_{b^{},𝐰_0^{}}^{}(u^{})`$ of $`\widehat{U}^{}`$ can be easily determined, without using (8), by just substituting $`u=gu^{}`$ in (4). After rearranging terms, we find that $`\mathrm{\Psi }_{b^{},𝐰_0^{}}^{}(u^{})`$ is given by the expression in Eq. (4) with all the quantities replaced by their primed counterparts and $`𝐰^{}=𝐰_0^{}=\mathrm{𝟎}`$. For given $`b`$ and $`𝐰`$, an expansion coefficient $`\varphi _{b^{}(b,𝐰)}(m^{};𝐰_0^{})`$, $`m^{}=0,\mathrm{},p^{}1`$, is nonzero only if there exists an integer pair $`(m,l)`$, $`m=0,\mathrm{},p1`$, $`l=0,\mathrm{},g1`$, solving the Diophantine equation $`pgl+gm+r_1=m^{}`$. The solution is then unique and $`\varphi _b^{}(m^{};𝐰_0^{})=\mathrm{exp}(2\pi ilr_2/g)\varphi _b(m;𝐰)`$, associated with a sparse expansion. In particular, for $`p=1`$, i.e., $`\mathrm{}^{}=2\pi /g^2`$, we can choose $`\varphi _b(m;𝐰)=1`$, and the only nonzero coefficients are $`\varphi _b^{}(m^{};𝐰_0^{})=\mathrm{exp}(2\pi ilr_2/g)`$ with $`m^{}=gl+r_1`$. This result can be obtained directly by applying the methods in Refs. , where the general case of $`2\pi /\mathrm{}=`$ perfect square was studied in detail. In conclusion, we have shown that the quasienergy problem of general QAMs for all BCs can be exactly reduced to a fixed-BCs problem. This reduction is possible due to a distinctive feature of QAMs: For given $`\mathrm{}=2\pi /p`$, the $`k`$ allowed Bloch wave vectors $`𝐰`$ form a lattice in the BZ, completely determined by the associated matrix $`A`$. Thus, this reduction is not possible for the special torus maps with Tr$`(A)=2`$, e.g., the kicked rotor and the kicked Harper maps, since the number of BCs is infinite for them. The reduction is implemented by a renormalization transformation $``$ exhibiting nontrivial features: For Tr$`(A)<2`$, the quantum cat maps are always fixed points of $``$, while for Tr$`(A)>2`$ they are generally fixed points only of $``$ accompanied by time reversal. These features reflect the scaling invariance of the quantum kernel (1). As a first application, we have shown that some results at fixed BCs can be derived and understood in a simple and transparent way in the general-BCs framework. In particular, one can easily establish the existence of a complete fixed-BCs basis of crystalline eigenstates in a classical limit for arbitrarily large perturbations. The results in this Letter may be applied to study several aspects of the general-BCs problem for QAMs using known results and methods at fixed BCs. An important and unexplored issue is the sensitivity of eigenstates to variations in the BCs and its relation to classical phase-space structures, especially in a regime of mixed phase space (large perturbations). Acknowledgments The author would like to thank S. Nonnenmacher for helpful comments and correspondence. Comments from J.P. Keating and Z. Rudnick are gratefully acknowledged. This work was partially supported by the Israel Science Foundation administered by the Israel Academy of Sciences and Humanities.
warning/0005/astro-ph0005079.html	ar5iv	text	# Dynamical processes, element mixing and chemodynamical cycles in dwarf galaxies ## 1. Introduction Observations of dwarf galaxies show no significant metal abundance gradients throughout the galaxies (Kobulnicky, Kennicutt, & Pizagno 1999; van Zee et al. 1998). A possible mechanism suggested by Tenorio-Tagle (1996) is the mixing of metal-enriched blow-out gas with fresh matter from the environment of the galaxy. After Supernovae (SNe) have caused an outflow of matter (depending on the mass of the galaxy; see Mac Low & Ferrara, 1998), the hot gas cools, condenses and falls back into the galactic body. The moderate-to-low stellar metallicities in dwarf elliptical galaxies (dEs) and the related dwarf spheroidals (dSphs) suggest that extensive gas loss has occured during their evolution by means of SNe-driven galactic winds (Larson 1974; Dekel & Silk 1986). Current starburst dwarf galaxies (SBDGs) are characterized by superwinds (Marlowe et al. 1995) or by large expanding X-ray plumes which are often confined by swept-up H$`\alpha `$ shells. Yet many dSph galaxies show not only a significant intermediate-age stellar population (Hodge 1989; Grebel 1997), but also more recent star formation (SF) events (Smecker at al. 1994; Han et al. 1997) indicating that gas was partly kept in the systems. On the other hand, gas infall might also cause a new SF episode as in NGC 205 (Welch, Sage, & Mitchell 1998). In several SBDGs large HI reservoirs envelope the luminous galactic body (e.g. NGC 4449: Hunter et al. 1998; I Zw 18: van Zee et al. 1998; NGC 1705: Meurer, Staveley-Smith, & Killeen 1998) and obtrude that the starburst is fueled by enhanced infall. Since the infall rate cannot yet be evaluated observationally the task for evolutionary models is to investigate the efficiency of SF and its according gas consumption. ## 2. The Model ### 2.1. The chemodynamical description Because of their low gravitational energy dwarf galaxies are strongly exposed to the energetic impact from processes like stellar winds, supernovae or even stellar radiation. When investigating gas mixing by large-scale dynamics and by small-scale exchanges it is necessary to distinguish between the dynamically separated gas components: the dense cloudy medium (CM) and the hot dilute intercloud medium (ICM). To account for differences in energy and timescale of the stellar yield an evolutionary code also has to divide stars into different stellar mass components: high-mass stars (HMS), intermediate-mass stars (IMS) and the not yet evolved low-mass stars (LMS). The evolution calculation presented here was performed using our chemodynamical evolution code CoDEx, a two dimensional grid-based axisymmetrical evolution code including matter, momentum and energy equations of all relevant processes in a self-consistant way. A detailed description can by found in Samland, Hensler & Theis (1997). Differing from the specifications given there recent HMS yields by Woosley & Weaver (1995) are used which provide secondary nitrogen production (and a small amount of primary nitrogen production). Also, for consistency, the most recent IMS models by van der Hoek & Groenewegen (1997) have been applied. Contrary to simple chemical evolutionary models that do not distinguish between CM and ICM and therefore have to parameterize separation and mixing of elements and possible selective outflow from the galaxy, the chemodynamical description allows to trace the flow and mixing of metal enriched matter in a self-consistent way. ### 2.2. Initial conditions We assume that the protogalactic gas cloud given by a Plummer-Kuzmin model (Satoh 1980) has been formed in a cosmological CDM scenario but was prevented from cooling by the metagalactic UV radiation field (Kepner, Babul, & Spergel, 1997). Not before this background radiation drops, the protogalactic gas cloud can cool due to recombination and subsequently collapes. The total numerical grid size is $`20kpc`$ x $`20kpc`$, the spatial resolution in the central parts amounts to almost $`100pc`$. The evolved numerical model starts with a baryonic mass of about $`10^9M_{}`$. A static dark matter halo according to Burkert (1995) with a mass of $`10^{10}M_{}`$ is added. This model is aimed to represent a dwarf irregular galaxy (dIrr). ## 3. Results ### 3.1. Dynamical phases of the evolution For analytical purposes we devide the model galaxy into four equidistant zones of radial extent of $`0.5kpc`$ along the galactic plane up to $`2.0kpc`$ and with a z-height of $`1.0kpc`$. The four zones represent tori due to the axisymmetrical grid. For consistency we use this scheme throughout the whole evolution, even when the model galaxy changes its shape (e.g. by the settling of a disk-like structure in the galactic plane). In Figure 1 the SF history for these four zones is plotted. The four curves are stacked, i.e. the uppermost curve is the sum of all four zones. An analysis of the model’s kinematics (Rieschick & Hensler, in prep.) implies five distinct dynamical phases of evolution that are indicated in Fig. 1 by vertical lines. Collapse phase ($`0`$ \- $`0.3Gyr`$): The protogalactic gas distribution is dynamically stable but cools and collapses. The net gas infall rate amounts to $`3.210^1M_{}/yr`$ leading to a central density increase by a factor of 100. Thus the SF rate rises steeply according to its quadratic dependence on the CM-density. The succeeding supernova type II (SNII) explosions starting with a time delay of about one $`Myr`$ smooth the further collapse of the CM. The peak in CM density is reached after about $`300Myr`$ with a total SF rate of $`1.810^2M_{}/yr`$. Post-collapse phase ($`0.3`$ \- $`0.8Gyr`$): Evaporation and gas outflow triggered by SNeII lead to further reduction of CM density down to $`1/4`$ of the maximum value. The net gas loss rate averaged over this phase amounts to $`6.610^2M_{}/yr`$. A second much smaller bounce occurs until an equilibrium between gravitation and gas pressure is reached. Transitional phase ($`0.8`$ \- $`2.0Gyr`$): In this phase the dynamical structure of the model galaxy changes completely. While during the collapse and the post-collapse phase the inward or outward motions, respectively, show almost coherence, now they differ locally and the SF splits into several small regions. This process is accompanied by both decreasing total gas density and total SF. Turbulent phase ($`2.0`$ \- $`6.2Gyr`$): Two different zones have formed: the central part with a radius of about $`0.5kpc`$ and a thick disk in the range between $`0.5`$ and $`2kpc`$ distance from the galactic center. These two parts follow separate paths in SF history. While the central SF rate stays almost constant, the disk SF rate increases nearly by a factor of ten in three steps. Thus the outer zones’ SF rate reaches more than $`{}_{}{}^{2}/_{3}^{}`$ of the total value while it has been only about $`20\%`$ before. Regularly this zone splits into a few separate SF regions (Hensler & Rieschick 1999) while SF is wandering outwards in the disk to larger radial distances because a growing number of gas clouds is accreted there due to infall. Irregular phase ($`>6.2Gyr`$): A global quasi-stability seems to be reached now, even though the SF in the outer parts of the galaxy happens in shortly living patches. Due to the smaller gas density the exchange with the surrounding is larger than in the central part leading to a non-existing abundance gradient. Because of the incoherence and inhomogenities of dynamics and local regions we denote this phase as ”irregular”. Figure 2 shows the SF rate at $`8.0Gyr`$ as an example. While the highest SF rate per volume ($`1.610^6M_{}/[Myrpc^3`$\]) exists in the center of the galaxy, three destinct SF regions with SF rates of about $`{}_{}{}^{1}/_{2}^{}`$ of the previous value are visible in the disk. The SF decreases slightly due to less infalling CM gas from the depleted HI reservoir. ### 3.2. The gas flow cycle Especially when studying the chemical evolution of galaxies the mixing and exchange processes between gas phases become important. Matter is expelled by SNe as hot dilute gas (ICM) and can return as clouds of cold gas (CM) if it is gravitationally bound. The gas changes its state from ICM to CM by condensation on the surface of gas clouds or vice versa by evaporation. In Figure 3 we have analyzed the mass flow treated in chemodynamics between the different components averaged between $`6.210.0Gyr`$ over a cylinder with $`r=2kpc`$ and $`z=1kpc`$. From the mass flow rates one can distinguish between an outer and an inner cycle. The outer one is produced by infall of CM. From this about 24 % is consumed by means of SF where 10 % is almost instantaneously rejected by HMS. The inner cycle represents SF and stellar evolution and, therefore, contains different timescales. The production of a minor fraction of hot ICM by SNe II leads to evaporation of remaining CM and escapes from the galactic body as outflow as a consequence of its high energy content. Significantly, almost 80 % of the infall is immediately converted into outflow by only 3.0 % of gas that has gone through the stellar cycle and is puffed up by stellar energy release. The full evaporation rate can exceed the infall rate because shell sweep-up leads to fragmentation and cloud formation in addition to a small amount of condensation. We emphasize that even though a large amount of the outflowing gas is gravitationally unbound and leaves the galactic body, the metals produced in HMS are for the most part kept in the outer gas flow cycle by mixing ICM with continuously infalling clouds. As a result of this mechanism only a few percent of the metals leave the gravitational field of the galaxy with the outflowing ICM. While this mixing itself happens on a timescale of about $`20Myr`$, the complete cycle of metal enrichment takes almost $`1Gyr`$ because of the low infall velocity at later evolutionary stages. In contrary, the inner cycle leads to an efficient self-enrichment of SF regions within $`10Myr`$ (see also Skillman, Dohm-Palmer, & Kobulnicky, 1998). The outer enrichment timescale would be reduced significantly, however, in a scenario of a rapidly infalling intergalactic gas cloud. This interaction scheme is running through all dynamical phases with different amounts but nearly the same mass flux ratios. During the collapse phase the infall rate is larger leading to intensive accumulation of matter in the CM and stellar components. In the post-collapse phase the CM density is even decreasing. ## 4. Conclusions In the previous section we have shown that a constant infall of metal-poor matter from the enveloping ”HI reservoir” drives a permanent gas mixing cycle keeping SNe-produced metals effectively in the gravitational field of the galaxy. Since recent deep HI observations of dIrrs and SBDGs have discovered an increasing number of large gas envelopes, i.e., of enormous gas reservoirs circling or accumulating around the visible body of the BCDGs, the relevance of infall episodes for the DG evolution is obvious and serves as the most promissing explanation for their observed abundances (see Hensler, Rieschick, & Köppen, 1999). Additionally we demonstrate by chemodynamical models that the SF cycle (Fig. 3) is triggered by only a small fraction of infalling matter, but produces sufficient energy to cause major evaporation and to drive more than $`3/4`$ of the infalling gas, incorporated into the ICM, as a galactic outflow back into an outer long-term cycle. Chemodynamical models can provide a fundamental insight into strong interactions between dynamical and energetical processes that happen in these sensitively balanced systems of low gravitational energy. #### Acknowledgments. We gratefully acknowledge cooperation and discussions with J. Köppen and Ch. Theis. A.R. is supported by the Deutsche Forschungsgemeinschaft under grant no. He 1487/23-1. The numerical calculations have been performed at the computer center of the University of Kiel and the NIC in Jülich. ## References Burkert, A. 1995, ApJ, 447, L25 Dekel, A., & Silk, J. 1986, ApJ, 303, 39 Grebel, E. 1997, Rev. Mod. 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warning/0005/math0005139.html	ar5iv	text	# Rational points near curves and small nonzero \|𝑥³-𝑦²\| via lattice reduction ## 1 Introduction One intriguing class of Diophantine problem concerns small values of homogeneous polynomials. In the simplest nontrivial case of a polynomial in three variables defining a projective plane curve $`C:P(X,Y,Z)=0`$, the problem can be reformulated thus: given a plane curve $`C`$, describe for each positive $`N,\delta `$ the rational points of height at most $`N`$ in $`𝐏^2`$ which are at distance at most $`\delta `$ of $`C`$. With present-day methods, hardly any nontrivial results can be proved on the number or existence of such points. But one can still seek numerical evidence, and efficient algorithms for obtaining this evidence. The direct approach is to try all $`x,y`$ with $`\|x\|,\|y\|N`$, and for each pair to solve $`P(x,y,z)=0`$ for $`z[N,N]`$, recording those cases in which $`z`$ is sufficiently close to an integer. This requires space $`O(\mathrm{log}(N))`$ but time $`(N^2+\delta N^3)\mathrm{log}^{O(1)}N`$, which is inefficient once $`\delta `$ is much smaller than $`N^1`$ since for general $`C,N,\delta N^3`$ the number of solutions should be proportional to $`\delta N^3`$. We give a new algorithm, also requiring only $`O(\mathrm{log}(N))`$ space, but with heuristic running time $`(N+\delta N^3)\mathrm{log}^{O(1)}N`$. Thus as long as $`\delta N^2`$ we expect to find all the points of height $`N`$ and distance $`\delta `$ in time only $`\mathrm{log}^{O(1)}N`$ per point. Moreover, our method readily parallelizes, since it divides the computation into many independent subproblems. We describe this algorithm, give the heuristic estimate for its run time, and briefly discuss the problem, which seems quite difficult, of proving our heuristic time estimates. We prove (Thm.2.1) that an alternative description of those points can always be computed in the heuristically expected time. We then discuss natural generalizations to other valuations and higher dimensions. An algorithm for finding rational points near a variety can in particular find rational points on the variety; applying our methods to embeddings of the variety in projective spaces of high dimension we obtain a new approach to this fundamental problem in computational number theory which improves on existing methods in several important cases. This approach also works for non-algebraic varieties, and even yields a theoretical result (Thm.2.4) on the paucity of rational points on non-algebraic analytic curves. We next describe experimental results of the implementation of our algorithm to various curves of interest, notably the Fermat curves of degree $`n>2`$, where some of our experimental findings led us to new polynomial families of small values of $`\|z^ny^nx^n\|`$ (Thm.3.1). We devote a separate section to the case of the cubic Fermat curve, corresponding to small values of $`\|z^3y^3x^3\|`$, a problem for which there is already some literature and the heuristics are subtler. In particular, we found for several integers $`d<10^3`$ the first representation of $`d`$ as a sum of three integer cubes; D.J.Bernstein has since extended the search up to $`N=210^9`$ and beyond, and found many new solutions, including one for $`d=30`$ which was a long-standing open problem. Finally we show how to modify our algorithm to efficiently search for small nonzero values of $`\|x^3y^2\|`$. This is the topic of Hall’s conjecture, which is part of a web of important Diophantine problems surrounding the ABC conjecture of Masser and Oesterlé. The conjecture asserts that $`x^3y^2`$ is either zero or $`_ϵx^{1/2ϵ}`$ for all $`x,y𝐙`$. We are able to find all solutions of $`0<\|x^3y^2\|x^{1/2}`$ with $`xX`$ in time $`O(X^{1/2}\mathrm{log}^{O(1)}X)`$, again using only $`O(\mathrm{log}X)`$ space. Using this improvement on the obvious $`X\mathrm{log}^{O(1)}X`$ method of trying all $`xX`$, we computed all cases of $`0<\|x^3y^2\|<x^{1/2}`$ with $`X10^{18}`$. We found ten new solutions, including most notably $$5853886516781223^3447884928428402042307918^2=1641843$$ with $`x^{1/2}/\|x^3y^2\|=46.600+`$, improving the previous record by a factor of almost $`10`$. In this case the time estimate is not heuristic; its proof not only streamlined the computation but even yields new theorems on the distribution mod $`1`$ of $`(cx)^{3/2}`$ for any positive rational $`c`$. We announce some of these results at the end of the present paper; the full statements and proofs will appear elsewhere. ### 1.1 Acknowledgements Richard K. Guy wrote the book \[G\] that first introduced me to many open problems in number theory including the Diophantine equations $`x^3+y^3+z^3=d`$ \[G, Prob. D5\], and later brought me up to date on recent work on this problem. Dan J. Bernstein efficiently implemented my new algorithm for the problem. Alan Murray told me of the appearance of approximate integer solutions of $`x^{12}+y^{12}=z^{12}`$ on The Simpsons. Frits Beukers and Franz Lemmermeyer filled gaps in my knowledge of earlier work concerning Hall’s conjecture. Barry Mazur suggested that a method for locating points near a variety might also profitably be applied to finding points on the variety; this started me thinking in the direction that led to Theorems 2.2 through 2.4. Alf van der Poorten and Hugh Montgomery directed me to Bombieri and Pila’s work \[BP\] concerning integral points on curves; Peter Sarnak put me in contact with Pila, who noted his more recent paper \[P\]; meanwhile Victor Miller alerted me to results announced by Roger Heath-Brown \[HB2\], who discussed his and Pila’s work with me. Meanwhile, Michel Waldschmidt informed me of relevant results by Weierstrass and others collected in \[M2, Chapter 3\]. I thank them all for these contributions to the present paper. Most of the numerical and symbolic computations reported here were carried out using the gp/pari and macsyma packages. This work was made possible in part by funding from the David and Lucile Packard Foundation. ## 2 The algorithm in theory ### 2.1 Specification and heuristic analysis While we are mainly interested in algebraic plane curves $`C`$, the algorithm does not require so strong a hypothesis: we can find<sup>1</sup><sup>1</sup>1 Our computations indicate that the first example is probably the smallest value of $`\|x^\pi +y^\pi z^\pi \|`$ for positive integers $`x,y,z`$, and at any rate the smallest with $`z10^6`$; and the second is the smallest ratio of $`\|x^7+y^7z^7\|`$ to $`z^7`$, and even to $`z^4`$, for positive integers satisfying $`xy<z10^6`$. See the next section. $`2063^\pi +8093^\pi 8128^\pi =0.019369`$ as well as $`386692^7+411413^7=(11.035\mathrm{}10^{18})441849^7`$. All we need is that $`C`$ is the image of a differentiable map $`\varphi :[0,1]\mathrm{𝐑𝐏}^2`$ with bounded second derivatives. Fix a positive $`\delta 1`$, and assume $`\delta N^2`$ for reasons given in the next paragraph. Partition $`[0,1]`$ into $`O(\delta ^{1/2})`$ intervals $`I_m`$ each of length $`\|I_m\|=O(\delta ^{1/2})`$. On each $`I_m`$, approximate $`\varphi `$ to within $`O(\|I_m\|^2)=O(\delta )`$ by a linear approximation $`\overline{\varphi }`$. Then a point at distance $`\delta `$ from $`\varphi (I_m)`$ remains at distance $`\delta `$ from $`\overline{\varphi }(I_m)`$. We now treat each $`I_m`$ independently. The triples $`(x,y,z)𝐙^3\{\mathrm{𝟎}\}`$ such that $`(x:y:z)𝐏^2`$ has height $`N`$ and is within $`O(\delta )`$ of $`\varphi (I_m)`$ are among the nonzero integer points in a parallelepiped $`P_m`$ of height, length and width proportional to $`N,\delta ^{1/2}N,\delta N`$. Thus we expect that $`\|P_m𝐙^3\|`$ is approximately the volume of $`P_m`$, provided that this volume is $`1`$. This is the case once $`\delta N^2`$. (That is why we insisted that $`\delta N^2`$: choosing smaller $`\delta `$ would only make us work at least as hard to find fewer points.) Listing all the points in $`P_m𝐙^3`$ is a standard application of lattice reduction. Let $`M_m`$ be an invertible $`3\times 3`$ matrix such that $`M_mP_m`$ is the cube $`K=[1,1]^3`$. We are then seeking all $`v𝐙^3`$ such that $`M_mvK`$, or equivalently all vectors in $`KM_m^1𝐙^3`$. We find them by reducing the lattice $`M_m^1𝐙^3`$. This gives us a matrix $`L_m\mathrm{GL}_3(𝐙)`$ such that $`M_mL_m`$ is small. Now $`M_mvK`$ if and only if $`w𝐙^3(M_mL_m)^1K`$ where $`v=L_mw`$. But $`(M_mL_m)^1K`$ is contained in the box centered on the origin whose $`i`$-th side is twice the $`l^1`$ norm of the $`i`$-th row of $`(M_mL_m)^1`$ ($`i=1,2,3`$). For each nonzero integral $`w`$ in this box, calculate $`(x,y,z)=L_mw`$ and test whether $`(x:y:z)`$ in fact has height $`N`$ and lies within $`\delta `$ of $`C`$. Doing this for each $`m`$ yields the full list of such points. As advertised, the algorithm requires only $`O(\mathrm{log}N)`$ space (though much more space is usually needed to store the results of the computation). Also, since each of many intervals $`I_m`$ is treated independently, the computation can be massively parallelized with little loss among processors that interact only by reporting each $`(x:y:z)`$ to headquarters as it is found. How long do we expect the computation to take? We assume that $`\varphi `$ and its derivatives can be calculated to within $`N^{O(1)}`$ in time $`\mathrm{log}^{O(1)}N`$. Such is the case for all curves we consider and for every algebraic plane curve. Then each $`M_m`$ takes only $`\mathrm{log}^{O(1)}N`$ operations to compute. Each lattice reduction can also be done in time polynomial in $`\mathrm{log}N`$, since our lattices are in fixed dimension — and moreover our dimension of $`3`$ is small enough that Minkowski reduction is described explicitly. \[For an overview and further references concerning Minkowski reduction, see \[CS, pp.396–7\].\] So far this amounts to $`\delta ^{1/2}N`$ time up to the usual log factors. Now each $`P_m`$ has volume $`2^3/\|detM_m\|(\delta ^{1/2}N)^3`$. If each $`(M_mL_m)^1`$ had all of its entries $`O(\delta ^{1/2}N)`$ — equivalently, if the shortest nonzero vectors of each lattice $`M_m^1𝐙^3`$ had length $`\delta ^{1/2}/N`$ — then there would only be $`O((\delta ^{1/2}N)^3)`$ choices for $`w`$, which summed over $`m`$ gives $`O(\delta N^3)`$. Thus the total work would indeed be $`\mathrm{log}^{O(1)}N`$ times the expected number of solutions. Unfortunately it is too optimistic to expect that the entries of $`(M_mL_m)^1`$ are all $`\delta ^{1/2}N`$. If the lattices $`M_m^1𝐙^3`$ are randomly distributed in the space of lattices of covolume $`(\delta ^{1/2}/N)^3`$ in $`𝐑^3`$, some of them will have nonzero vectors much shorter than $`\delta ^{1/2}/N`$. However, the average number of lattice vectors in $`K`$ of a random lattice of determinant $`D`$ is still $`O(1/D)`$. Thus we expect — and typically find in practice — that, even accounting for the occasional short lattice vector, we will find all rational points of height $`N`$ that lie within $`\delta `$ of $`C`$, doing on average $`\mathrm{log}^{O(1)}N`$ work per point. ### 2.2 Can the estimates be made rigorous? Our assumption that the lattices $`M_m^1𝐙^3`$ are randomly distributed was not proved; indeed it is false at least for some choices of $`C`$. Most glaringly, if $`C`$ is a rational straight line then there are $`N^2`$ rational points on $`C`$, and a fortiori at least as many at distance $`\delta `$. While we of course will not apply our algorithm to straight lines, we do apply it to the $`n`$-th Fermat curve, which has contact of order $`n`$ with several rational lines such as $`y=z`$; each of those lines contains $`N^{22/n}`$ points at distance $`1/N^2`$ from the curve, exceeding the expected count of $`N\mathrm{log}^{O(1)}N`$ once $`n>2`$. (These are the points we exclude by imposing the inequality $`y<z`$ in $`0<xy<z<N`$.) Assume, then, that $`C`$ has at most finitely many tangent lines which have contact of order $`>2`$ with $`C`$, and for any $`\delta >0`$ let $`C_\delta `$ be the curve consisting of points of $`C`$ at distance $`\delta `$ from each of those higher-order tangent lines. For each point $`P`$ on $`C_\delta `$ we obtain a lattice $`L_\delta (P)𝐑^3`$ whose nonzero short vectors correspond to points near $`P`$ in $`𝐏^2(𝐐)`$, of height $`\delta ^{1/2}`$, lying at distance $`\delta `$ from $`C_\delta `$. This gives a map $`\mathrm{\Lambda }_\delta `$ from $`C_\delta `$ to the moduli space of lattices in $`𝐑^3`$. We would thus like to ask: as $`\delta 0`$, does the image of $`\mathrm{\Lambda }_\delta `$ become uniformly distributed in this moduli space? There are several problems with this formulation of our question. A minor one is that we have not defined $`\mathrm{\Lambda }_\delta `$ precisely enough for the question to make sense, because we have left some $`O`$-constants unspecified. This did not matter for qualitative properties such as whether the lattice has $`O(1)`$ short vectors, but makes it easy to frustrate uniform distribution by simply choosing $`\mathrm{\Lambda }_\delta `$ to avoid a small region in the moduli space. This problem is easy enough to fix for any given $`C`$; for instance, if $`C`$ is given by $`x(x:y(x):1)`$ for some differentiable function $`y:[0,1][1,1]`$ with bounded second derivatives, we may take for $`\mathrm{\Lambda }_\delta (x)`$ the integer span of the columns of $$\left(\begin{array}{ccc}0& 0& \delta \\ 1& 0& x\\ y^{}/\delta & \mathrm{\hspace{0.33em}\hspace{0.17em}1}/\delta & (xy^{}y)/\delta \end{array}\right).$$ (1) But this brings us to a more serious difficulty. The question of whether $`\mathrm{\Lambda }_\delta (C_\delta )`$ is asymptotically uniformly distributed as $`\delta 0`$ is likely to be a very hard problem in analytic number theory. For our purposes we are only concerned with how often and how close does $`\mathrm{\Lambda }_\delta (P)`$ come near the cusp of the moduli space. For instance, we see in the final section that if $`C`$ is a conic then $`\mathrm{\Lambda }_\delta (C_\delta )`$ is restricted to a surface in the moduli space of lattices in $`𝐑^3`$, but within that surface it still approaches the cusp rarely enough that the average number of short vectors in a lattice in $`\mathrm{\Lambda }_\delta (C)`$ is still $`\mathrm{log}(1/\delta )`$. In general, then, what we would like is the following result: as $`\delta 0`$, the average number of vectors of norm $`<1`$ of a lattice in $`\mathrm{\Lambda }_\delta (C_\delta )`$ is $`\mathrm{log}^{O(1)}(1/\delta )`$. This still looks like a very difficult problem. While it remains open, we propose a contingency plan in case the lattices $`L_\delta (P)`$ have many more short vectors than expected. If all the short vectors are multiples of a single vector of small norm, there is no difficulty, because all these multiples yield the same point in $`𝐏^2`$. But there could be two independent short vectors, whose linear combinations yield a line in $`𝐏^2`$ containing many points of small height near $`C`$. We claim that this is in fact the only way that a lattice of covolume $`1`$ could have more than $`O(1)`$ short vectors. This claim is easy enough to check using the description of Minkowski-reduced lattices in $`𝐑^3`$, but we shall later need a generalization to lattices in higher dimension. We thus state and prove the generalization as follows: ###### Lemma 1 For each positive integer $`n`$ and positive real $`t`$ there exists an effective constant $`M_n(t)`$ such that the following bound holds: for any lattice $`\mathrm{\Lambda }𝐑^n`$ whose dual lattice $`\mathrm{\Lambda }^{}`$ has no nonzero vector of length $`<r`$, and for any $`R>0`$, there are at most $`M_n(rR)r^n\|\mathrm{\Lambda }\|^1`$ vectors of length $`R`$ in $`\mathrm{\Lambda }`$. Here $`\|\mathrm{\Lambda }\|`$ is the covolume Vol$`(𝐑^n/\mathrm{\Lambda })`$. The lemma can be obtained as a consequence of the theory of lattice reduction, but it is not easy to extract $`M_n(t)`$ explicitly this way. We thus give the following alternative proof in the spirit of \[C1\] from which explicit (albeit far from optimal) bounds $`M_n(t)`$ may be easily computed if desired. ###### Proof Given $`n`$, choose a positive Schwartz function $`f:𝐑^n𝐑`$ with the following properties: $`f`$ is radial, i.e. $`f(x)`$ depends only on $`\|x\|`$; and the Fourier transform $`\widehat{f}:𝐑^n𝐑`$, defined for $`y𝐑^n`$ by $$\widehat{f}(y):=_{x𝐑^n}f(x)e^{2\pi i(x,y)}𝑑x,$$ (2) satisfies $`\widehat{f}(y)0`$ for all $`y`$ such that $`\|y\|1`$. For instance, we may take $$f(x)=(\|x\|^2+a)e^{\pi c\|x\|^2}$$ (3) where $$0<c<\frac{2\pi }{n},a=\frac{1}{c^2}\frac{n}{2\pi c},$$ (4) because the Fourier transform of a function (3) is $$\widehat{f}(y)=\left(a+\frac{n}{2\pi c}\frac{\|y\|^2}{c^2}\right)e^{\pi \|y\|^2/c}$$ (5) for any $`c>0`$ and $`a𝐑`$. By Poisson summation, $$\underset{x\mathrm{\Lambda }}{}f(rx)=\frac{1}{r^n\|\mathrm{\Lambda }\|}\underset{y\mathrm{\Lambda }^{}}{}\widehat{f}(y/r).$$ (6) Under the hypothesis on $`r`$, the only positive term in the sum over $`y`$ is $`\widehat{f}(0)`$. The sum over $`x`$ is bounded from below by the sum over $`x`$ of length $`R`$, which is at least the number of such vectors times $`\mathrm{min}_{\|x\|R}^{}f(rx)`$. It follows that $`\mathrm{\Lambda }`$ has at most $$\frac{\widehat{f}(0)}{r^n\|\mathrm{\Lambda }\|\mathrm{min}_{\|x\|rR}f(x)}=M_n(rR)r^n\|\mathrm{\Lambda }\|^1$$ (7) vectors of length $`R`$, as claimed. ###### Corollary 1 For each positive integer $`n`$ there exists an effective constant $`A_n`$ such that if a lattice $`\mathrm{\Lambda }𝐑^n`$ has more than $`A_nR^n/\|\mathrm{\Lambda }\|`$ vectors of length $`<R`$ for some $`R>0`$ then all those vectors lie in a hyperplane, which can be computed in polynomial time. ###### Proof Except for the last phrase, this follows from the previous Lemma by taking $`r=1/R`$ and $`A_n=M_n(1)`$, since then $`\mathrm{\Lambda }^{}`$ must have a nonzero $`y`$ of length at most $`r`$, and any vector of $`\mathrm{\Lambda }`$ of length $`<R`$ must be orthogonal to $`y`$. To assure that $`y`$ can be computed in polynomial time, we take $`r=c/R`$ for a positive constant $`c`$ small enough that if $`\mathrm{\Lambda }^{}`$ has a nonzero vector of length at most $`c/R`$ then the LLL algorithm will find a (possibly different) nonzero vector of length at most $`1/R`$. Our Corollary now holds with $`A_n=c^nM_n(c)`$. From the case $`n=3`$ of this Corollary we deduce: ###### Theorem 2.1 Let $`C`$ be the image of a differentiable map $`\varphi :[0,1]\mathrm{𝐑𝐏}^2`$ with bounded second derivatives. Then for each $`N>1`$ and $`\delta N^2`$ one can find $`O(\delta N^3)`$ rational points and $`O(N)`$ rational line segments each of length $`O(1/N)`$ in $`𝐏^2`$ which together include all rational points of height $`N`$ at distance $`\delta `$ from $`C`$. These points and line segments can be computed in time $`\delta N^3\mathrm{log}^{O(1)}N`$. Outside of $`O(\delta N^3\mathrm{log}N)`$ space used only to record each point or segment as it is found, the computation requires space $`\mathrm{log}^{O(1)}N`$. All implied constants depend effectively on $`C`$. Note that here we do not exclude neighborhoods of high-order rational tangents to $`C`$; such tangents will contain some of the lines segments computed by the algorithm. ### 2.3 Variations and generalizations The problem of finding rational points near plane curves is only the first nontrivial example of many analogous problems to which our method can apply. We briefly discuss some of these here. One easy variation is to change the norm: instead of approximating the curve in the real valuation, use a nonarchimedean one, or a combination of several. For instance, one can efficiently seek nontrivial triples of small integers the sum of whose cubes is divisible by a high power of $`2`$ or of $`10`$. Likewise one can replace $`𝐙`$ by $`𝐅_q[T]`$ or similar rings in function fields of positive genus. The lattice-reduction step should then be even easier than in the archimedean case, though in the function-field setting our approach faces strong competition from the method of undetermined coefficients, and it is not clear which is superior. All these comments apply equally to the adaptation of our method to the problem of finding small nonzero values of $`\|x^3y^2\|`$, provided the characteristic is not $`2`$ or $`3`$. For the $`\|x^3y^2\|`$ problem, the work estimates are again rigorous; otherwise, they are still heuristic, but their analysis may be more tractable in the function-field case. Higher dimensions present many new opportunities. The easiest generalization is to a $`𝒞^2`$ hypersurface in $`𝐏^{k1}`$. Here we are seeking small values of a homogeneous function of $`k`$ variables evaluated at an integral point. This time we chop the hypersurface into $`O(\delta ^{(k2)/2})`$ chunks each of diameter $`O(\delta ^{1/2})`$, and replace each chunk by a subset of a hyperplane which approximates it to within $`O(\delta )`$. The points of height $`N`$ that are within $`O(\delta )`$ of this chunk then come from integral points in a parallelepiped in $`𝐑^k`$ whose sides have lengths $`O(N)`$, $`O(N\delta ^{1/2})`$ ($`k2`$ times), and $`O(N\delta )`$. Again most of these this parallelepipeds have volume $`1`$ provided $`\delta N^2`$, and we locate the integral points using lattice reduction in $`𝐑^k`$. So, as long as $`\delta N^2`$, we expect to find on the order of $`\delta N^k`$ points, using $`\mathrm{log}^{O_k(1)}N`$ space and spending $`\mathrm{log}^{O_k(1)}N`$ time per point. For a general hypersurface, this again improves on other approaches to the problem. But the improvement decreases with $`k`$: the direct approach takes time $`N^{k1}\mathrm{log}^{O(1)}N`$, and we lower the exponent by a factor no better than $`(k2)/(k1)`$, which approaches $`1`$ as $`k\mathrm{}`$. Moreover, lattice reduction in $`𝐑^k`$ quickly becomes difficult as $`k`$ grows. Another consideration is that for special surfaces there are known, and simpler, algorithms that take time $`N^{k1}\mathrm{log}^{O(1)}N`$ or less once $`k>3`$. For instance, for Fermat surfaces in $`𝐏^3`$, one readily adapts the method of \[B1\] to find all solutions of $`x^n+y^n=z^n\pm t^n+O(z^{n2})`$ in positive integers with $`tzN`$, in expected time $`N^2\mathrm{log}^{O(1)}N`$, and with no need for lattice reduction in $`𝐑^4`$ or other complicated ingredients. This computation does require space proportional to $`N\mathrm{log}N`$, which however poses no difficulty for practical values of $`N`$. As in the previous paragraph, all that is described in the present paragraph can be done also for a nonarchimedean norm, with similar results except that lattice reduction over a function field is tractable even for large $`k`$. In either case rigorous estimates may become even less accessible as $`k`$ grows. We can generalize further to manifolds $`𝐏^{k1}`$ of codimension $`c>1`$. Here we expect to find on the order of $`\delta ^cN^k`$ rational points of height $`N`$ at distance $`O(\delta )`$ from $``$. We chop $``$ into $`O(\delta ^{(k1c)/2})`$ patches of diameter $`O(\delta ^{1/2})`$, each of which yields a parallelepiped in $`𝐑^k`$ with dimensions of order $`N`$ (once), $`N\delta `$ ($`d`$ times), and $`N\delta ^{1/2}`$ (the remaining $`k1c`$ dimensions). We thus expect to efficiently find all $`\delta ^cN^k`$ points as long as $`\delta N^{2k/(k+c1)}`$. A further possibility emerges if $``$ has bounded derivatives past the second derivatives and has small enough dimension compared with $`k`$: we can then make further headway when $`\delta `$ falls below that threshold. Usually we are only interested in points much closer than $`N^{2k/(k+c1)}`$; but as long as we use only the $`𝒞^2`$ structure we gain nothing by making $`\delta `$ even smaller, so we may as well find all the points at distance $`O(N^{2k/(k+c1)})`$ and locate the best approximations in the resulting list. However, if $``$ is $`𝒞^3`$ and its dimension $`d=k1c`$ is so small that $`k>\left(\genfrac{}{}{0pt}{}{d+2}{2}\right)`$, then a patch of diameter $`ϵ`$ is contained in a box with $`d`$ sides of length $`ϵ`$, a further $`(d^2+d)/2`$ sides of length $`ϵ^2`$, and the remaining $`k\left(\genfrac{}{}{0pt}{}{d+2}{2}\right)`$ sides of length $`ϵ^3`$. This means that we can makes our parallelepipeds thinner in some directions, and thus use wider patches of $``$, covering the entire manifold with fewer of them. This lets us locate the points of height $`N`$ closest to $``$ in time significantly less than it would take to record all the points at distance $`N^{2k/(k+c1)}`$, even though not so efficiently that we only spend $`\mathrm{log}^{O_k(1)}N`$ time per point. More generally if $``$ is a $`𝒞^i`$ manifold we can exploit bounds on the $`i`$-th derivatives once $`k>\left(\genfrac{}{}{0pt}{}{d+i1}{i1}\right)`$. If the ambient projective space is not of high enough dimension, we can still make some use of approximations to $``$ of degree $`i>1`$ by using the $`i`$-th Veronese embedding $`V_i`$ of $`𝐏^{k1}`$ into projective space of dimension $`\left(\genfrac{}{}{0pt}{}{k+i1}{i}\right)1`$. \[The $`i`$-th Veronese embedding takes the point with projective coordinates $`(X_1:\mathrm{}:X_k)`$ to the point whose projective coordinates are all $`\left(\genfrac{}{}{0pt}{}{k+i1}{i}\right)`$ monomials of degree $`i`$ in the $`X_j`$. Thus $`V_i`$ raises all heights to the power $`i`$, and transforms intersections with hypersurfaces of degree $`d`$ in $`𝐏^{k1}`$ into hyperplane sections in a projective space of much higher dimension. For more on Veronese embeddings, see for instance \[FH\], where they arise several times.\] The idea is to surround each patch of $`V_i()`$ by a box containing all points in $`V_i(𝐏^{k1})`$ at distance $`O(\delta )`$ from $`V_i()`$. The resulting asymptotic improvement may be only barely worth it in practice, though. Consider the simplest case of a $`𝒞^3`$ curve $`C𝐏^2`$, embedded in $`𝐏^5`$ by $`V_2`$. Assume for simplicity that the parametrization $`\varphi `$ of $`C`$ has $`\|\varphi ^{\prime \prime }\|`$ bounded away from zero. Then, for $`\delta `$ such that $`ϵ^3\delta ϵ^2`$, the radius-$`\delta `$ neighborhood in $`𝐏^2`$ of an interval of length $`ϵ`$ on $`C`$ maps into a box in $`𝐏^5`$ whose sides are of order $`ϵ,ϵ^2,\delta ,ϵ^3,ϵ^4`$. \[To see this, choose coordinates $`(X_0:X_1:X_2)`$ on $`𝐏^2`$ for which $`\varphi `$ is of the form $`(1:t:t^2+O(t^3))`$ for $`t`$ in a neighborhood of $`0`$, and note that $`V_2`$ takes $`(X_0:X_1:X_2)`$ to $`(X_0^2:X_0X_1:X_0X_2:X_1^2:X_1X_2:X_2^2)`$.\] Thus the points of height at most $`N`$ in that neighborhood map to lattice points in a $`6`$-dimensional parallelepiped of volume $`\delta N^{12}ϵ^{10}`$. (Here $`N`$ occurs to the power $`12`$ rather than $`6`$ because $`V_2`$ squares the height of each rational point.) Thus if we take $`ϵ=(\delta N^{12})^{1/10}`$ we expect to find all points at distance $`\delta `$ from $`C`$, of which there should be about $`\delta N^3`$, in time $`N(\delta N^2)^{1/10}\mathrm{log}^{O(1)}N`$. The condition $`\delta ϵ^3`$ yields $`\delta N^{36/13}`$, so we save a factor of at most $`N^{1/13}`$. We pay not only by missing the points at distance between $`N^{36/13}`$ and $`N^2`$ (which usually do not interest us anyway) but also by reducing lattices of rank $`6`$ rather than $`3`$. This takes more time per lattice, and probably yields parallelepipeds whose average bounding box is larger. Each of these effects amounts to only a constant factor, but these factors may be considerable, and it will be interesting to see how large $`N`$ must be for this use of $`V_2`$ to be practical. ### 2.4 Rational points on varieties In the last paragraph we exploited the fact that points near $``$ map under $`V_i`$ to points that are not only near $`V_i()`$ but exactly on $`V_i(𝐏^{k1})`$. We can go much further when we search for points exactly on $``$. Again we consider the simplest case of a curve. We begin with a curve in one projective space: ###### Theorem 2.2 Let $`C`$ be an algebraic curve in $`M`$-dimensional projective space, defined over $`𝐐`$ and not contained in any hyperplane. Then for any $`N1`$ the rational points of $`C`$ of height at most $`N`$ can be listed in time $`_CN^{2/M}\mathrm{log}^{O_M(1)}N`$. The implied constants depend effectively on $`d`$ and $`C`$. Remarks. As seen above for $`M=2`$, this result applies more generally to a $`𝒞^M`$ curve in $`𝐏^M`$ whose intersection with any hyperplane can be computed in polynomial time. The exponent $`2/M`$ is best possible: a rational normal curve of degree $`M`$ (a.k.a. the image of $`𝐏^1`$ under $`V_M`$) has on the order of $`N^{2/M}`$ rational points of height at most $`N`$, and it takes time $`N^{2/d}\mathrm{log}N`$ just to write them down. The constant implied in $`O_d(1)`$ and/or $`_C`$, while effective, may be unpleasant in practice for large $`M`$, since lattice reduction in dimension $`M+1`$ is involved. ###### Proof A segment of $`C`$ of length $`N^{2/M}`$ is contained in a box whose $`i`$-th side is $`N^{2i/M}`$ ($`i=1,2,\mathrm{},M`$). The rational points of height at most $`N`$ in this box come from points of $`𝐙^{M+1}`$ contained in a box $`B`$ whose $`i`$-th side is $`N^{12i/M}`$ ($`i=0,1,2,\mathrm{},M`$) and thus has volume $`O(1)`$. It takes time $`_C\mathrm{log}^{O_M(1)}N`$ to apply lattice reduction and, by Corollary 1, either list $`𝐙^{M+1}B`$ or find a hyperplane containing $`𝐙^{M+1}B`$. In the former case, we test whether each of the resulting $`O_d(1)`$ points lies in $`C`$. In the latter case, we map this hyperplane to $`𝐏^M`$ and intersect it with $`C`$, finding at most $`\mathrm{deg}(C)=O_C(1)`$ rational points. Thus in either case we find all rational points of height $`N`$ on our segment in time $`_C\mathrm{log}^{O_M(1)}N`$. Since it takes only $`O(N^{2/M})`$ segments to cover $`C`$, we are done. It might seem that this algorithm is superfluous: if $`C`$ has genus $`0`$ then its small rational points may be found directly from a rational parametrization, without any lattice reduction; and if $`C`$ has positive genus then we can find all its points of height $`N`$ in time $`\mathrm{log}^{O(1)}N`$ once we have generators of the Mordell-Weil group of the Jacobian of $`C`$. But the difficulty is that we must first find these generators, and this requires locating rational points on a curve or a higher-dimensional variety. For instance, to find the Mordell-Weil group of an elliptic curve $`E`$ we usually apply a few descents and then search for points on certain principal homogeneous spaces for $`E`$, each of which is a curve $`C`$ of genus $`1`$, usually (in the case of a complete $`2`$-descent) of the form $`y^2=P(x)`$ for some irreducible quartic $`P𝐙[X]`$. One then searches for $`x𝐐`$ of height up to $`H`$ for which $`P(x)𝐐^2`$. There are on the order of $`H^2`$ candidates for $`x`$; one can set up a sieve to efficiently try them all, but this still takes time $`H^2\mathrm{log}^{O(1)}H`$ (and significant space). Instead we can embed $`C`$ in $`𝐏^3`$ as the intersection of two quadrics (by writing $`P(x)`$ as a homogeneous quadric in $`1,x,x^2`$), and use the algorithm of Thm.2.2 with $`N=H^2`$ to find all rational solutions of $`y^2=P(x)`$ with $`x`$ of height $`H`$ in time $`H^{4/3}\mathrm{log}^{O(1)}H`$. For certain $`E`$ one can use Heegner points to locate a rational point on $`C`$ to within $`\delta `$ (see \[E3\]); if $`\delta H^2`$, this is sufficient to identify $`x`$ using continued fractions, a.k.a. lattice reduction in dimension $`2`$. Using the new algorithm, we see that $`\delta H^{4/3}`$ suffices if we use lattice reduction in dimension $`4`$. This saves a constant factor in the computation of $`x`$, since fewer digits and terms are needed in the floating-point computation of Heegner points. When $`C`$ has genus $`>1`$, there are only finitely many rational points by Faltings’ theorem, but they still may be of significant number and/or height. For instance, in \[KK, S\] one finds curves $`C:y^2=P(x)`$ of genus $`2`$ which have hundreds of rational points. In both cases, all points with $`x`$ of height $`10^6`$ were found using the $`H^2\mathrm{log}^{O(1)}H`$ sieve method, a substantial computation. At least in the case considered in \[S\], where the Jacobian of $`C`$ is absolutely simple with large Mordell-Weil rank, it would probably be even more onerous to find all these points by first determining the Mordell-Weil group. But the embedding $`(1:x:x^2:x^3:y)`$ of $`C`$ into $`𝐏^4`$ yields an improvement from $`H^2`$ to $`H^{3/2}`$ with $`5`$-dimensional lattice reduction. We can do even better by mapping the same curve to larger projective spaces. Fix an algebraic curve $`C`$ of genus $`g`$ defined over $`𝐐`$, and a divisor $`D`$ on $`C`$ of degree $`d>0`$. For $`n`$ sufficiently large, the sections of $`nD`$ embed $`C`$ into $`𝐏^{ndg}`$. This embedding sends any rational point on $`C`$ of height (exponential, as usual here) $`H`$ relative to $`D`$ to a point on $`𝐏^{ndg}`$ of height $`H^n`$. By Thm.2.2 again, we can find all such points in time $`H^{2n/(ndg)}\mathrm{log}^{O(1)}H`$. Letting $`n\mathrm{}`$, we conclude: ###### Theorem 2.3 Fix an algebraic curve $`C/𝐐`$ and a divisor $`D`$ on $`C`$ of degree $`d>0`$. For each $`ϵ>0`$ there exists an effectively computable constant $`A_ϵ`$ such that for any $`H1`$ one can find all points of $`C`$ whose height relative to $`D`$ is at most $`H`$ in time $`A_ϵH^{(2/d)+ϵ}`$. For instance, all rational points on $`y^2=P(x)`$ with $`x`$ of height at most $`H`$ can be computed in time $`_ϵH^{1+ϵ}`$. What of varieties $``$ of dimension $`\mathrm{\Delta }>1`$ in $`𝐏^M`$? A chunk of radius $`\delta `$ then yields the intersection of $`𝐙^{M+1}`$ with a box with sides as follows: one of length $`O(N)`$, $`\mathrm{\Delta }`$ sides of length $`O(N\delta )`$, $`\left(\genfrac{}{}{0pt}{}{\mathrm{\Delta }+1}{2}\right)`$ sides of length $`O(N\delta ^2)`$, …, $`\left(\genfrac{}{}{0pt}{}{\mathrm{\Delta }+i1}{i}\right)`$ sides of length $`O(N\delta ^i)`$, … until $`\left(\genfrac{}{}{0pt}{}{\mathrm{\Delta }+j}{j}\right)=_{j=0}^i\left(\genfrac{}{}{0pt}{}{\mathrm{\Delta }+i1}{i}\right)`$ first exceeds $`M`$. As usual we choose $`\delta `$ so that the product of these sides is $`1`$, and apply lattice reduction to each of $`O(\delta ^\mathrm{\Delta })`$ chunks. The difficulty here is that if the lattice is nearly degenerate, the hyperplane found in Corollary 1 meets $``$ not in a finite number of points but in a subvariety of positive dimension $`\mathrm{\Delta }1`$. This suggests an induction on $`\mathrm{\Delta }`$, since we can apply our method to that hyperplane section of $``$. But already for $`\mathrm{\Delta }=2`$ such an argument requires a version of Thm.2.2 with more uniformity in the implied constants than we know how to obtain. However, as with our first nontrivial case of curves in $`𝐏^2`$, we do not expect such degenerate lattices to arise in practice often enough to raise the computational cost above $`O(\delta ^\mathrm{\Delta }\mathrm{log}^{O(1)}N)`$, except for a finite number of proper subvarieties of $``$. If we assume this, we can again obtain better estimates by embedding $``$ in larger projective spaces. Fix an ample divisor $`D`$ on $``$, and ask for all rational points whose height relative to $``$ is at most $`H`$. Using the sections of $`nD`$ to embed $``$ in projective spaces, and letting $`n\mathrm{}`$, we find the following heuristic generalization of Thm.2.3: for each $`ϵ>0`$, there exists a proper subvariety $`_0(ϵ)`$ of $``$ such that all points of $`_0(ϵ)`$ of height at most $`H`$ relative to $`D`$ can be found in time $$O_ϵ(H^{((\mathrm{\Delta }+1)/\|D\|)+ϵ}),$$ (8) where $`\|D\|`$ is the $`\mathrm{\Delta }`$-th root of the intersection number $`D^\mathrm{\Delta }`$. One might even hope that $`_0(ϵ)`$ can be taken independent of $`ϵ`$. For instance, if $``$ is a surface of degree $`d`$ in $`𝐏^3`$ then we expect that, for some union $`_0`$ of curves on $``$, we can find all rational points of height $`N`$ on $`_0`$ in time $`_ϵN^{(3/\sqrt{d})+ϵ}`$. We must admit that this is unlikely to yield a practical improvement over the $`N^2\mathrm{log}^{O(1)}N`$ method we already knew: the first $`V_i`$ that reduces the exponent of $`N`$ below $`2`$ is $`V_3`$, and then (assuming $`d4`$) the exponent drops only to $`24/13`$ — but instead of reducing $`4`$-dimensional lattices we are then faced with lattice reduction in dimension $`20`$. It will probably be a long time before $`N`$ can feasibly be taken large enough that this extra effort is worth the $`N^{2/13}`$ factor gained. Returning to plane curves, we can use this idea to prove an even stronger bound on rational points on a plane curve $`C`$ that is analytic but not algebraic. This is because the homogeneous monomials of degree $`i`$ in the coordinates of $`C`$ are linearly independent for each $`i`$, so $`V_i(C)`$ spans a projective space whose dimension grows quadratically in $`i`$ (whereas for an algebraic curve the growth is always linear). This leads us to the following result: ###### Theorem 2.4 Let $`C`$ be a transcendental analytic arc in $`𝐏^2`$, i.e. $`C=\{f(x):axb\}`$ where $`f`$ is an analytic map from a neighborhood of $`[a,b]`$ to $`𝐏^2`$ whose image is contained in no algebraic curve. Then for each $`ϵ>0`$ there exists a constant $`A_ϵ`$ such that for every $`H1`$ there are fewer than $`A_ϵH^ϵ`$ points of height $`H`$ in $`C𝐏^2(𝐐)`$. ###### Proof For each positive integer $`i`$ consider $`V_i(C)𝐏^{(i^2+3i)/2}`$. Since $`C`$ is transcendental, $`V_i(C)`$ is an analytic arc $`V_if`$ contained in no hyperplane of $`𝐏^{(i^2+3i)/2}`$. Now apply the argument for Thm.2.2 with $`N=H^i`$. As noted in the remarks following the statement of that theorem, the curve need not be algebraic as long as it is $`𝒞^M`$ and its intersection with any hyperplane is of bounded size. (Here we need not compute this intersection numerically, since we are only bounding the number of rational points of small height on $`C`$, not computing them efficiently.) The differentiability is clear since $`V_i(C)`$ is analytic, and the boundedness is proved in the next lemma. We conclude that the number of points of height $`H`$ on $`C`$ is $`_iH^{4/(i+3)}`$. Since $`i`$ can be taken arbitrarily large, our theorem follows. The existence of an upper bound on the size of the intersection of any hyperplane with $`V_i(C)`$ is a special case of the following lemma in complex analysis. Throughout the lemma and its proof we count zeros of an analytic function according to multiplicity, even though in the application to Thm.2.4 a multiple zero is no worse than a simple one. ###### Lemma 2 Let $`E`$ be an open subset of $`𝐂`$ and $`V`$ a finite-dimensional vector space of analytic functions: $`E𝐂`$. Then for any compact subset $`KE`$ there exists an integer $`n`$ such that any nonzero $`fV`$ has at most $`n`$ zeros in $`K`$. ###### Proof Fix $`K`$. We shall say that a compact $`K^{}E`$ is “good” if its boundary $`(K^{})`$ is rectifiable and its interior $`K^{}(K^{})`$ contains $`K`$. Choose a good $`K_1`$, and define a norm on $`V`$ by $`f=sup_{zK^{}}\|f(z)\|`$. Let $`V_1`$ be the unit ball $`\{fV:f=1\}`$. It is sufficient to prove the lemma for $`fV_1`$. For each $`fV_1`$ choose a good $`K_fK^{}`$ such that $`f`$ does not vanish on $`(K_f)`$. Let $`r_f^{}=inf_{z(K_f)}\|f(z)\|`$, and let $`n_f`$ be the number of zeros of $`f`$ in $`K_f`$. By Rouché’s theorem, if $`gV`$ with $`fg<r_f`$ then $`g`$ has at most $`n_f`$ zeros in $`K_f`$, and a fortiori in $`K`$. Now $`V_1`$ is compact and is covered by the open balls $`B_f`$ of radius $`r_f`$ about $`fV_1`$. Thus there is a finite subcover $`\{B_{f_i}\}_{i=1}^M`$. Then $`n:=\mathrm{max}_in_{f_i}`$ is an upper bound for the number of zeros in $`K`$ of any $`fV_1`$, and thus of any nonzero $`FV`$. To recover our result on hyperplane sections of $`V_i(C)`$, take $`K=[a,b]`$, let $`E`$ be a neighborhood of $`K`$ on which $`f`$ is analytic, and choose any analytic functions $`f_0,f_1,f_2`$ on $`E`$ such that $`f=(f_0:f_1:f_2)`$ on $`E`$. Then take for $`V`$ the space of homogeneous polynomials of degree $`i`$ in $`f_0,f_1,f_2`$. If we understand $`f`$ well enough to obtain for each $`i`$ an effective bound $`n`$ in Lemma 2 then the constants $`A_ϵ`$ in Thm.2.4 are effective too. With a little additional work $`𝐐`$ can be replaced by an arbitrary number field $`F`$ embedded in $`𝐂`$, and $`C`$ by $`f(K)`$, where $`K𝐂`$ is any compact subset and $`f`$ is again an analytic map from a neighborhood of $`K`$ to $`𝐏^2`$ whose image is contained in no algebraic curve. A separate approach to bounding the number of rational points on curves was initiated in \[BP\] and pursued further in \[P\] and \[HB2\]. For example, Heath-Brown obtains in \[HB2\] bounds on the number of rational points on an algebraic plane curve that coincide with the time estimates in our Theorems 2.2 and 2.3. Moreover, our Thm.2.4 is contained in \[P, Thm.8\], which asserts that for a number field $`F`$ with $`[F:𝐐]=n`$ the number of $`F`$-rational points of height $`<H`$ on a transcendental analytic arc $`C`$ is at most $`A_{C,n,ϵ}H^ϵ`$. Probably the methods of \[BP, P\] can also prove these results with arcs $`C`$ replaced by compact transcendental curves $`f(K)`$, and our bounds can also be made uniform in $`F`$ given $`[F:𝐐]`$. There is clearly some overlap between the two approaches; for instance the Corollary preceding \[P, Thm.8\] is the same as our Lemma 2, but proved using the determinants of \[BP, P\]. What is not clear, but intriguing, is whether those determinantal methods and our lattice-reduction technique can ultimately be interpreted as facets of the same basic idea. All this also suggests the question of whether a transcendental arc can contain infinitely many rational points, of whatever height. I thank Michel Waldschmidt for pointing out that this question was already asked, and later answered affirmatively, by Weierstrass. See \[M2, Chapter 3\] for this and related results. ## 3 The algorithm in practice In this section we report on the outcome of the application of our algorithm to various plane curves, and on some results suggested by our findings. We suppress details of the explicit constants replacing each $`O(\mathrm{})`$ and $``$; these details are of course crucial in practice, but are straightforward and not enlightening. In each case our curve has some rational points of inflection, and we make sure to truncate our curve enough to avoid the tangents at such points but not so much that we lose approximations near but not on those tangents. In general, for a plane curve given by a homogeneous equation $`P(X,Y,Z)=0`$ of degree $`n`$, we associate to a rational point $`(x:y:z)`$ near but not on the curve the number $$n\mathrm{max}(\|x\|,\|y\|,\|z\|)^{n3}/\|P(x,y,z)\|.,$$ (9) which measures how close the point $`(x:y:z)`$ is to the curve relative to the point’s height. We insert the factor $`n`$ so that we can reasonably compare approximations for curves of different degrees. For instance, for the Fermat curve one expects that as $`x,y`$ vary, the integer $`z^n:=x^n+y^n`$ comes on average within $`\frac{1}{4}nz^{n1}`$ of the nearest $`n`$-th power of an integer, and thus that the smallest value of $`\|z^ny^nx^n\|`$ for $`z[N,2N]`$ is proportional to $`nz^{n3}`$. One could insert further factors to correct for the length and shape of our curve, but these factors are not significant for most of the curves we study. We noted already that the heuristics leading to formulas such as (9) refer to “random” $`(x:y:z)`$ near the curve, not for systematic families of approximations which may attain values of the ratio (9) larger or more often than expected. We again give an example for the Fermat curves, which were the subjects of most of our computations. One usually guesses that for each $`r`$ there will be $`r\mathrm{log}N`$ triples $`(x,y,z)`$ such that the ratio (9) exceeds $`r`$. However, in the identity $$(t+1)^n(t1)^n=2nt^{n1}+O(t^{n3}),$$ (10) we can make $`2nt^{n1}`$ an $`n`$-th power by setting $`t=2nu^n`$; this yields $`N^{1/n}`$ triples with (9) bounded away from zero. We note the special cases $`n=2,3`$ of this identity: for $`n=2`$, the $`O(t^{n3})`$ error vanishes, and we recover a familiar parametrization of Pythagorean triples; for $`n=3`$, the error is constant, and we can scale the identity to obtain the known family of solutions $`(x,y,z)=(6t^2,6t^31,6t^3+1)`$ of $`z^3y^3x^3=2`$. Returning to general $`n`$: in our searches we set the threshold on $`z^{n3}/\|z^ny^nx^n\|`$ low enough to find all the examples coming from (10), as a check on the computation; but we chose a higher threshold for the tabulation of results so that our list is not dominated by this polynomial family. ### 3.1 Fermat curves of degree $`>3`$ We implemented our algorithm to find small values of $`\|z^ny^nx^n\|`$ with $`0<xy<z`$, $`4n20`$, and $`z[10^3,10^6]`$. Since the threshold for “small” depends on the size of $`z`$, we wrote $`[10^3,10^6]`$ as the union of $`10`$ intervals $`[N/2,N]`$ and treated each separately. We also used a direct search for $`z<5000`$, using the overlap region $`[1000,5000]`$ as a check on the computation. We did not attempt to fine-tune the algorithm for efficiency, since we carried it out more as a demonstration project than a major computational undertaking. Thus we programmed the search in gp, using the built-in arithmetic and LLL lattice reduction. We estimate that transcribing the program to C, and replacing LLL by Minkowski reduction in $`𝐑^3`$, would speed the computation by roughly an order of magnitude; of course a machine faster than a Sun Sparcstation Ultra $`1`$ would help too. With a C program and a more powerful machine, it should be feasible to search the range $`n[4,20]`$, $`z<10^9`$ in time on the order of a month. The behavior of the run times and the counts of solutions with $`\|z^ny^nx^n\|z^{n2}`$ seem broadly consistent with our heuristics, though we have not attempted a detailed statistical analysis. We tabulate the most striking examples, those with $$r:=nz^{n3}/(z^ny^nx^n).$$ (11) of absolute value at least $`4`$: $`n`$ $`x`$ $`y`$ $`z`$ $`r`$ 4 167 192 215 $`4.5`$ 4 8191 16253 16509 $`12.9`$ 4 24576 48767 49535 $`64.5`$ 4 49152 97534 99070 $`8.1`$ 4 34231 157972 158059 $`5.2`$ 4 76215 311390 311669 $`14.8`$ 5 13 16 17 $`120.4`$ 5 26 32 34 $`15.1`$ 5 39 48 51 $`4.5`$ 5 42 71 72 $`8.8`$ 5 262 328 347 $`6.2`$ 5 1125 2335 2347 $`5.0`$ 5 5088 16155 16165 $`4.1`$ 5 190512 292329 298900 $`5.5`$ 6 1236 3587 3588 $`12.5`$ 6 6107 8919 9066 $`9.9`$ 7 386692 411413 441849 $`78.4`$ 7 773384 822826 883698 $`9.8`$ $`n`$ $`x`$ $`y`$ $`z`$ $`r`$ 8 209959 629874 629886 $`11.6`$ 8 209945 629826 629838 $`11.6`$ 9 6817 10727 10747 $`5.3`$ 9 21860 25208 25903 $`24.7`$ 10 280 305 316 $`137.1`$ 10 560 610 632 $`17.1`$ 10 840 915 948 $`5.1`$ 10 7533 8834 8999 $`4.4`$ 12 1782 1841 1922 $`6.1`$ 12 3987 4365 4472 $`7.1`$ 12 781769 852723 874456 $`10.3`$ 13 666 806 811 $`8.3`$ 13 5579 8235 8239 $`4.1`$ 15 434437 588129 588544 $`42.9`$ 16 492151 741267 741333 $`4.6`$ 19 79 85 86 $`4.7`$ 19 491 565 567 $`4.9`$ 19 43329 51144 51257 $`5.8`$ 20 4110 4693 4709 $`4.3`$ All decimal values of $`r`$ are rounded to the nearest tenth. If for some integer $`\lambda >1`$ we have $`r>4\lambda ^3`$ then $`(\lambda x,\lambda y,\lambda z)`$ will also appear in the table provided $`\lambda z10^6`$; this happens for $`\lambda =2`$ at $`n=4,5,7,10`$, and for $`\lambda =3`$ at $`n=5,10`$. The first examples for $`n=10`$ and particularly $`n=5`$ (where $`13^5+16^5=17^5+12`$) are small and striking enough that one feels they must have been observed already, but I do not know a reference. On the other hand, the first two examples for $`n=12`$ have been published, and in a most unlikely place: each appeared in a different episode of the popular animated cartoon The Simpsons. Perhaps the third example for $`n=12`$, or an example with $`n=7`$ or $`n=15`$, could be used if the cartoon repeats this theme once more; the relative error $`\|z^ny^nx^n\|/z^n`$ in each case is between $`1`$ and $`2`$ parts in $`10^{18}`$, as compared to $`310^{10}`$ and $`210^{11}`$ for the two four-digit examples… Frivolity aside, one is struck by the pair of examples for $`n=8`$. The values of $`r`$ are far from the largest in the table, but they are almost equal and opposite, and involve nearly equal triples $`(x,y,z)`$ for which $`zy`$ has the same small value of $`12`$. This suggests that we are dealing with a polynomial family $`(x(t),y(t),z(t))`$ specialized at $`t=\pm t_0`$. Indeed we quickly find that these are the cases $`t=\pm 3`$ of $$(32t^9+6t)^8+(32t^8+7)^8=(32t^9+10t)^8+212^{28}t^{40}+O(t^{32}),$$ (12) with $`r=t^5/21+O(t^3)`$. Thus arbitrarily large values of $`r`$ occur, and indeed $`z^8y^8x^8`$ can be as small as $`O(z^{40/9})`$ rather than the expected $`O(z^5)`$. Trying to generalize the identity (12) further, we soon find that there are similar families for any exponent $`n`$ such that $`3n(n2)`$ is a square: ###### Theorem 3.1 Let $`n>1`$ be a positive integer. Then there exist polynomials $`x(t),y(t),z(t)𝐙[t]`$ of the form $$x(t)=Ct^n+D,y(t)=At^{n+1}+Bt,z(t)=At^{n+1}+B^{}t$$ (13) with $`A0`$, $`B^{}B`$ such that $`z^ny^nx^n`$ is a polynomial of degree at most $`n(n3)`$, if and only if $`3n(n2)`$ is a square. In that case, there exist infinitely many integer triples $`(x,y,z)`$ with $`0<x<y<z`$ such that $`z^ny^nx^nz^{(n^23n)/(n+1)}`$. ###### Proof Let $`b,b^{}`$ be the distinct rational numbers $`B/A,B^{}/A`$. Expand $`z^ny^n`$ at infinity: $`z^ny^n`$ $`=`$ $`nA^n(b^{}b)(t^{n^2}+{\displaystyle \frac{n1}{2}}(b^{}+b)t^{n^2n}`$ $`+{\displaystyle \frac{(n1)(n2)}{6}}(b_{}^{}{}_{}{}^{2}+b^{}b+b^2)t^{n^22n}+O(t^{n^23n})).`$ For this to be of the form $`(Ct^n+D)^n+O(t^{n^23n})`$ we must have $$(n1)\left(\frac{n1}{2}(b^{}+b)\right)^2=2n\frac{(n1)(n2)}{6}(b_{}^{}{}_{}{}^{2}+b^{}b+b^2).$$ (15) The discriminant of this quadratic equation in $`b^{}/b`$ is $`3n(n2)`$ times a square; thus (15) has nonzero rational solutions if and only if $`3n(n2)`$ is a square. Explicitly we find that $`b,b^{}`$ are proportional to $`\sqrt{(n^22n)/3}\pm 1`$. Conversely, suppose $`n^22n=3m^2`$ for some integer $`m`$. Let $$z=A(t^{n+1}+c(m+1)t),y=A(t^{n+1}+c(m1)t).$$ (16) Then $$z^ny^n=2cnA^n\left(t^n+\frac{n1}{n}cm\right)^n+O(t^{n^23n}).$$ (17) To make this $`(Ct^n+D)^n+O(t^{n^23n})`$ with $`C,D𝐙`$ we now need only choose nonzero $`c𝐙`$ so that $`2cn`$ is an $`n`$-th power (e.g. take $`c=(2n)^{n1}`$), and then choose $`A`$ so that $`n\|Acm`$. Specializing $`t`$ to sufficiently large integers in the resulting $`(x(t),y(t),z(t))`$ yields infinitely many integer triples $`(x,y,z)`$ with $`0<x<y<z`$ such that $`z^ny^nx^nz^{(n^23n)/(n+1)}`$, as claimed. ∎ The smallest $`n>3`$ such that $`n^22n=3m^2`$ is $`n=8`$. There are infinitely many further examples, starting with $`27`$, $`98`$, $`363`$, …, and parametrized by a Fermat-Pell equation. Dropping the constraint $`n>3`$ yields the further cases $`n=2`$ and $`n=3`$. For $`n=2`$ we again obtain a Pythagorean parametrization, this time with $`x,y,z`$ multiplied by $`t`$; for $`n=3`$ we find $$(9t^3+1)^3+(9t^4)^3(9t^4+3t)^3=1,$$ (18) one of infinitely many polynomial solutions of $`x^3+y^3z^3=1`$. ### 3.2 The Fermat cubic Our algorithm applies to the Fermat cubic as it does to the Fermat curves of higher degree, but we treat it separately both because the heuristic analysis is subtler and because the problem of finding small values of $`\|z^3y^3x^3\|`$ has already attracted some attention. We noted that in general we expect the smallest values of $`\|z^ny^nx^n\|`$ to be comparable with $`z^{n3}`$. For $`n=3`$, we have $`z^{n3}=1`$, and of course (given this case of Fermat’s Last Theorem) $`\|z^3y^3x^3\|`$ can be no smaller than $`1`$ for nonzero integers $`x,y,z`$. Moreover, $`z^3y^3x^3`$ cannot be an arbitrary rational multiple of $`z^{n3}`$: only the discrete values $`\pm 1,\pm 2,\mathrm{}`$ may arise. Thus, instead of a Diophantine inequality $`z^ny^nx^nz^{n3}`$, we have a family of Diophantine equations $`z^3y^3x^3=d`$ ($`d𝐙`$), and new tools can bear on solving them or, failing that, describing their distribution of solutions. These equations have been investigated by various means since the beginning of the computer age; see \[G\] for references to work up to about 1980 (some of which dates back to the 1950’s), and \[B2, CV, HBLR, KTS, PV\] for more recent results. As we shall see, the problem has been approached in several ways, some of which already improve on direct exhaustion over some $`N^2`$ values of $`(x,y)`$. Still, our new linear approximation method is better yet, both in heuristic theory — even though by factors smaller than our accustomed $`N/\mathrm{log}^{O(1)}N`$ — and in practice, as evidenced by the computation of many new solutions. Our discussion here applies with almost no change to other “diagonal” cubics, such as $`x^3+y^3+2z^3`$ which was also singled out in \[G, Prob. D5\]; but we have not yet implemented a search for small values of $`\|x^3+y^3+2z^3\|`$ beyond what has already been reported in the literature. For each nonzero $`d`$, the expected distribution of solutions of $$z^3y^3x^3=d$$ (19) involves not only considerations of size — i.e. of local behavior at the archimedean place of $`𝐐`$ — but also on the behavior of $`z^3y^3x^3`$ at finite primes $`p`$: each $`p`$ contributes a local factor $`f_p(d)`$ that is the ratio of the $`p`$-adic measure of the $`𝐙_p`$-points of (19) to the average of that measure as $`d`$ ranges over $`𝐙_p`$. For instance, if any of those factors $`f_p(d)`$ vanishes, there can be no solutions at all. It is not hard to see that the only such local constraint is $`d\pm 4mod9`$. For such $`d`$, the resulting product over $`p`$ was investigated by Heath-Brown \[HB1\]. He showed that the product does not converge absolutely, but can nevertheless be analyzed and approximated numerically by comparing $`f_p(d)`$ with the factor at $`p`$ of the Euler product for $`\left(\zeta _{𝐐(\sqrt[3]{d})}(s)/\zeta (s)\right)^3`$ at $`s=1`$, which differs from $`f_p(d)`$ by a factor of at most $`1+O(p^{3/2})`$. The product $`_pf_p(d)`$ is then seen to diverge to $`+\mathrm{}`$ if $`d`$ is a cube and to converge to a positive limit when $`d`$ is neither a cube nor congruent to $`\pm 4mod9`$. Heath-Brown thus conjectured in \[HB1\] that all nonzero integers $`d\pm 4mod9`$ occur as $`z^3y^3x^3`$ infinitely often. So far this is only known when $`d`$ is either a cube or twice a cube, thanks to polynomial parametrizations, which the above heuristics do not try to account for. We have already exhibited polynomial solutions for $`d=1,2`$. For many $`d4mod9`$ which are neither cubes nor twice cubes, not a single solution is known for $`z^3y^3x^3=d`$. Heath-Brown observes \[HB1\] that this is not surprising, because for many of these $`d`$ the expected number of solutions with $`z[N,10^6N]`$ is positive but smaller than $`1`$. Guy \[G\] lists the cases with $`d<10^3`$ which were open as of 1980, and while the list is now shorter the question of which integers are the sums of three cubes is not yet settled even in that range. For instance, the case $`d=30`$ was open until 1999, and had been the smallest open case for several decades. We have noted already that a direct search finds all small $`\|z^3x^3y^3\|`$ with $`z<N`$ in time $`N^2\mathrm{log}^{O(1)}N`$. There have been several improvements on this, all obtained by rewriting the equation (19) as $$x^3+d=z^3y^3=(zy)(z^2+yz+y^2).$$ (20) Once $`x^3+d`$ is factored, which takes heuristic time $`N^{o(1)}`$, all solutions of (20) can be found by trying each factor of $`x^3+d`$ for $`zy`$. Given the value of $`d`$, this takes time only $`N^{1+o(1)}`$. In addition to dealing with only one $`d`$ at a time, this method has the disadvantage that the $`N^{o(1)}`$ time required to factor $`x^3+d`$, though subexponential, is still considerable. The advantage of this method is that it finds all solutions with $`xN`$, while $`y,z`$ may be considerably larger, of order up to $`N^{3/2}`$. Many of the new solutions found in \[KTS\] are of this type, with $`y,z`$ large but $`zy`$ very small. Heath-Brown observed that, again given $`d`$, the factorization of $`x^3+d`$ can be simplified by a precomputation in $`𝐙[\sqrt[3]{d}]`$, though the complexity of the precomputation depends unpredictably on $`d`$ via the arithmetic of the number field $`𝐐(\sqrt[3]{d})`$; this approach was implemented in \[HBLR\]. Note that in effect these methods find rational points near the Fermat cubic that are close to the tangents to the curve at its inflection points — the same tangents that demand special care in our algorithm. A further variation which we suggested in 1996 is to use the factorization $$z^3d=x^3+y^3=(x+y)(x^2+xy+y^2)$$ (21) as follows: fix $`x+y`$, solve for $`zmodx+y`$, and try each of the resulting values of $`z`$. Here we only find solutions with $`z`$, not $`x`$, bounded by $`N`$, but the advantage is that factoring costs are greatly diminished. To find all cube roots of $`dmodx+y`$ requires factoring $`x+y`$, a number of size $`N`$ rather than $`N^3`$; and with enough space to set up a sieve the factorization can be avoided entirely. In 1999, Eric Pine, Kim Yarbrough, Wayne Tarrant and Michael Beck, all graduate students at the University of Georgia, took up this suggestion, choosing $`d=30`$, and found the first solution: $$30=2220422932^3283059965^32218888517^3$$ (22) We announced our new algorithm in the same 1996 posting to the NMBRTHRY mailing list, together with results of a search for solutions with $`z<10^7`$ and $`\|d\|<10^3`$. We did our search in gp, making our computation easy to program (since gp already provides multiprecision arithmetic and lattice reduction) but far from optimally efficient. In 1999, unaware of the work of the Georgia group, we asked Dan J. Bernstein for an efficient implementation. He soon wrote a C program that found all solutions with $`z<310^9`$ and $`\|d\|<10^4`$, including (22) and many others. Several values of $`d`$ had not been previously represented as the sum of three cubes. Detailed results and analysis will appear elsewhere. As usual, since we are interested in small $`d`$, not all $`dN`$, the improvement by a factor $`N^{1/13}`$ should apply here as well to find all cases of $`\|z^3y^3x^3\|z^{3/13}`$ with $`z<N`$, but we have not attempted to implement such a computation. ### 3.3 Miscellaneous examples Trinomial units. One sometimes sees in Olympiad-style mathematics contests the question “Is $`z^{1/3}`$ greater or smaller than $`x^{1/3}+y^{1/3}`$?” for some specific positive integers $`x,y,z`$. Of course this is a challenge only when the sign of the difference $`u_3:=z^{1/3}(x^{1/3}+y^{1/3})`$ cannot be determined by inspection. In some cases the question be settled by applying classical inequalities; for instance if $`a>b>0`$ then $`(a+b)^{1/3}+(ab)^{1/3}<2a^{1/3}`$ by convexity of the cube root. The general solution is to compute the norm of $`u_3/z^{1/3}`$, an algebraic number of degree $`9`$ none of whose other conjugates is real unless $`x=y`$. We find that $`u_3`$ has the same sign as $$𝐍(x,y,z):=(zyx)^327xyz.$$ (23) Moreover, given the size of $`x,y,z`$, the smaller $`𝐍(x,y,z)`$ is, the nearer $`u_3`$ will be to $`0`$. In particular, we would like to have $`𝐍(x,y,z)=\pm 1`$, which would make the algebraic integer $`u_3`$ a unit. Thus again we seek rational points close to a plane cubic curve, here $`𝐍(x,y,z)=0`$. This time the curve is rational: by construction, it is parametrized by $`(x:y:z)=(t^3:(1t)^3:1)`$. It is thus not smooth, but its only singularity is the isolated point $`(x:y:z)=(1:1:1)`$ (geometrically a node with complex conjugate tangents), which does not affect our algorithm. The three rational points of inflection at $`xyz=zyx=0`$ do affect our algorithm, but fortunately we are not interested in the points on their tangent lines, since those are the points with $`xyz=0`$. We thus restrict our attention to the portion of the curve with $`x/z,y/z>1/N`$, i.e. with $`tN^{1/3}`$ and $`1tN^{1/3}`$ in the rational parametrization. This takes us far enough from the inflection points that they cause us no difficulty. The situation is now much the same as for $`z^3y^3x^3=d`$. We expect the number of solutions of $`𝐍(x,y,z)=d`$ of height up to $`N`$ to be proportional to $`\mathrm{log}N`$ times a product of local factors $`g_p(d)`$. The only local factor that can vanish is $`g_3(d)`$, which is nonzero if and only if $`9\|d`$ or $`d\pm 1mod9`$. We henceforth assume that $`d`$ is in one of these congruence classes. We can then check whether $`_pg_p`$ converges by comparing it with the $`L`$-series of the projective cubic surface $`𝐍(x,y,z)=dt^3`$. This in turn depends on the Galois structure of the Néron-Severi group of the surface, which can be determined from the action of Galois on the lines on that cubic surface, as explained in \[W1\]. We must be careful here because, unlike $`x^3+y^3+z^3=dt^3`$, the surfaces $`𝐍(x,y,z)=dt^3`$ are not smooth: each has an $`A_2`$ singularity at $`(x:y:z:t)=(1:1:1:0)`$. Thus each has, not $`27`$ lines as usual, but $`15`$, of which $`6`$ go through the singularity; see \[BW\]. Explicitly, these are the preimages under the projection to $`(x:y:z)`$ of the three coordinate axes and the two tangents to the curve at $`(1:1:1)`$. We conclude that, as with (19), $`_pg_p(d)`$ converges unless $`d`$ is a cube. So we expect the number of unparametrized solutions of height $`N`$ to grow as $`\mathrm{log}N`$, except when $`d`$ is a cube, when it should grow faster, albeit still as a power of $`\mathrm{log}N`$ — perhaps $`\mathrm{log}^3N`$, by analogy with Manin’s conjecture for cubic surfaces. Unlike the case of (19), we know of no solutions of $`𝐍(x,y,z)=\pm 1`$ in nonconstant polynomials $`x,y,z𝐙[t]`$, other than the trivial ones with $`xyz=0`$. Nevertheless we can find infinitely many nontrivial integer solutions parametrized by Fermat-Pell equations, and thus show that the number of solutions of height $`N`$ is $`\mathrm{log}N`$. There are several ways to do this. In 1982 we found a somewhat complicated route to such a parametrization, obtaining a family of solutions starting with $`𝐍(16948,31226,186919)=1`$. The details may be found in the pages of \[CM\]. Many years later, we observed that a simpler approach is to factor $`𝐍(x,y,z)=\pm 1`$ as $$27xyz=(zyx)^31=(zyx1)\left[(zyx)^2\pm (zyx)+1\right].$$ (24) For each $`r𝐐^{}`$, we obtain a conic curve $`C_r`$ by setting $`(zyx1)=rx`$ in (24). This can be viewed geometrically as follows: the affine surface $`𝐍(x,y,z)=\pm 1`$ contains the line $`x=(zyx1)=0`$; thus the intersection of the surface with any plane $`(zyx1)=rx`$ containing that line is the union of the line and some residual conic, which is our $`C_r`$. Likewise we could start from the line $`z=(zyx1)=0`$ and intersect it with a variable plane $`(zyx1)=rz`$. For many choices of $`r`$, one of these conics is a hyperbola with infinitely many integral points parametrized by a Fermat-Pell equation. In retrospect this approach to $`𝐍(x,y,z)=\pm 1`$, in which we fiber an affine surface by conics that may be regarded as principal homogeneous spaces for Fermat-Pell equations, seems a remarkable premonition of our later analysis \[E1\] of the projective quartic surface $`A^4+B^4+C^4=D^4`$ via a fibration by genus-$`1`$ curves (principal homogeneous spaces for elliptic curves). In both cases the approach finds infinitely many solutions but does not readily lend itself to efficiently finding all solutions of height $`N`$. Again a later computation found that the solution that was discovered first, because it lies on the first fiber that could contain a solution, is not the one of smallest height. We used our algorithm to find all small values of $`𝐍(x,y,z)`$ with $`0<x,y,z10^6`$. We found that the smallest solution of $`𝐍(x,y,z)=\pm 1`$ is $`(14,84,313)`$ of norm $`+1`$, followed by $`(6818,4996,46879)`$, $`(20388,4881,86830)`$, and $`(2742,32540,96843)`$ each of norm $`1`$, the known $`(16948,31226,186919)`$, and $`(3408,182899,370338)`$ of norm $`+1`$, with no further solutions up to $`10^6`$. We also found several primitive solutions of $`𝐍(x,y,z)=\pm 8`$ and a few sporadic examples with $`d`$ small but not a cube, which could not have been obtained at all using the factorization trick; the smallest of these are $$𝐍(204,115327,162434)=17,𝐍(650,1425,7899)=26.$$ (25) The $`𝐍(x,y,z)=17`$ solution yields a disappointingly large value of $`u_3`$ because the conjugates $`z^{1/3}y^{1/3}e^{\pm 2\pi i/3}x^{1/3}`$ are smaller than usual. An unexpected result — since the identity (10) cannot be used with exponents $`<1`$ — was a polynomial solution of $`𝐍(x,y,z)=108`$, namely $`(4,y(t),y(1t))`$ where $`y(t)=4t^36t+3`$. We can write this symmetrically as $`(8,g(t),g(t))`$ where $`g(t)=8t^312t6t+11`$, a cubic polynomial determined up to scaling by the condition that the Laurent expansion at infinity of $`(g(t))^{1/3}`$ have vanishing $`t^2`$ and $`t^4`$ terms. In this form, $`𝐍(x,y,z)`$ is the larger constant $`864=2^3108`$, but with the bonus that $`x`$ is a cube so $`u_3`$ involves one fewer surd; for instance, taking $`t=7`$ we find that $`\sqrt[3]{3279}`$ is smaller than $`2+5\sqrt[3]{17}`$ by less than $`3.7510^7`$. In this family, as with the first example in (25), $`u_3`$ is of order $`z^2`$, not $`z^{8/3}`$, because two of the conjugates of $`u_3`$ are $`O(1)`$. A similar investigation of $`u_4:=z^{1/4}(x^{1/4}+y^{1/4})`$ was not as productive, perhaps not surprisingly since there are no arithmetic reasons to expect many nonzero small examples. For the record, the smallest $`z^{11/4}\|u_4\|`$ value found for $`z<10^6`$ was $`0.365+`$ for $`(x,y,z)=(241,691,6759)`$, while the smallest $`\|u_4\|`$ in that range was $`(3.23)10^{16}`$ for $`(37792,36109,591093)`$. The $`\pi `$-th Fermat curve. To illustrate our algorithm also for non-algebraic curves, we chose to apply it to the Fermat curve of exponent $`\pi `$. Since $`\pi `$ exceeds $`3`$, but only slightly, we expected that $`\|z^\pi y^\pi x^\pi \|`$ achieves a global minimum over all $`x,y,z`$ with $`0<xy<z`$ but that the minimum might involve numbers of several digits. We were rewarded with the example $$2063^\pi +8093^\pi 8128^\pi =0.019369=8128^{\pi 3}/(184.75+),$$ (26) which seems likely to be the minimum of $`\|z^\pi y^\pi x^\pi \|`$ over all positive integers $`x,y,z`$. At any rate, according to our computations it is the smallest with $`z10^6`$. The ratio $`184.75+`$ is also the largest in that range, though there is also $$1198^\pi +4628^\pi 4649^\pi =(0.04949+)=4649^{\pi 3}/(66.794+).$$ (27) It will probably be a long time before the question of the minimality of (26) is settled; a weaker but still intractable conjecture is that there are only finitely many integer solutions of $`\|z^\pi y^\pi x^\pi \|<1`$. The Klein quartic. All our examples so far were Fermat curves, even though some had unusual exponents $`1/3`$, $`1/4`$, $`\pi `$. Probably the best-known projective plane curve that is not a Fermat curve is the Klein quartic $`K(X,Y,Z)=0`$, where $$K(X,Y,Z):=X^3Y+Y^3Z+Z^3X.$$ (28) We used our algorithm to search for small values of $`K(x,y,z)`$. By symmetry we may assume $`\mathrm{max}(x,y,z)=z`$. We are then seeking rational points near a segment of a plane curve with a single inflection point, at $`x=y=0`$. The tangent $`x=0`$ at this point accounts for the obvious family $`(0,1,z)`$ with $`K(x,y,z)=z`$. Our computation up to height $`10^6`$ quickly revealed a less obvious family, $`K(1,t^2,t^3)=t^2`$, with $`K(x,y,z)`$ growing even more slowly than the height. As usual we also found sporadic examples, though here (as with several other cases we have already seen such as the Fermat quintic) the best ones are small enough that our algorithm was not needed to locate them: $`K(1421,1057,1501)`$ $`=`$ $`49,`$ $`K(7211,8381,11010)`$ $`=`$ $`121,`$ (29) $`K(1550,11817,32615)`$ $`=`$ $`245,`$ with $`z/\|K(x,y,z)\|=30.6`$, $`91.0`$, $`133.1`$ respectively. The largest $`z/\|K(x,y,z)\|`$ found with $`z[10^5,10^6]`$ off the singular cubic $`y^3+x^2z=0`$ was $`6.756+`$, from $`K(7871,175577,829244)=122741`$. ## 4 Hall’s conjecture ### 4.1 Review of Hall’s conjecture By Hall’s conjecture me mean the following assertion: if $`x,y`$ are positive integers such that $$k:=x^3y^2$$ (30) is nonzero (equivalently, such that $`(x,y)(t^2,t^3)`$), then $$\|k\|_ϵx^{1/2ϵ}.$$ (31) (While this accords with current usage, it is not exactly what Hall originally wrote: as F. Beukers points out, Hall \[H\] conjectured $`\|k\|x^{1/2}`$, a stronger statement which is probably false — the usual heuristic suggests that there are at least $`(\delta +o(1))\mathrm{log}X`$ cases of $`0<\|k\|<\delta \sqrt{x}`$ with $`x<X`$ — but unlikely to be soon disproved. See also \[BCHS\] for the early history of this conjecture.) Among several equivalent forms of (31) we note the conjecture that the discriminant of an elliptic curve over $`𝐐`$ in its standard minimal form has absolute value $`_ϵ\|a_4\|^{1/2ϵ}`$. Known lower bounds on $`\|k\|`$ are much weaker than (31). By Siegel’s theorem on the finiteness of integer points on elliptic curves, each nonzero $`k𝐙`$ occurs finitely many times as $`x^3y^2`$, so $`\|k\|\mathrm{}`$ as $`x\mathrm{}`$. Siegel’s proof is ineffective and thus says nothing about how fast $`\|k\|`$ must grow with $`x`$. Starting with Baker’s method, effective bounds have become available, but they are still very weak. For instance, it is not yet possible to prove for any $`\theta >0`$ that $`\|k\|x^\theta `$. Hall’s conjecture is now recognized as an important special case of the Masser-Oesterlé ABC conjecture \[O\] (see also \[L\]). Thus its analogue over function fields is known to be true by Mason’s theorem \[M3\]. In the special case of Hall’s conjecture for polynomials $`x(t),y(t)`$, the fact that $`x^3y^2`$ is either zero or has degree $`>\frac{1}{2}\mathrm{deg}(x)`$ was proved some twenty years earlier by Davenport \[D2\] in response to a question raised in \[BCHS\]. As in \[E2\] it follows that the conjecture cannot be disproved by a polynomial parametrization, and indeed in any polynomial family $`(x(t),y(t)\|t𝐙)`$ we must have $`kx^\theta `$ with $`\theta >1/2`$. One does better with solutions parametrized by Fermat-Pell equations, i.e. $`x,y𝐙[t,\sqrt{at^2+bt+c}]`$ for some $`a,b,c𝐙`$ such that $`u^2=at^2+bt+c`$ has infinitely many solutions. The function field $`𝐐(t,\sqrt{at^2+bt+c})`$ is then still rational, so the Davenport-Mason inequality again holds, but since now there are two places at infinity one can have $`x^3y^2`$ of degree exactly $`\frac{1}{2}\mathrm{deg}(x)`$, and thus attain $`\theta =1/2`$. The existence of a single such family (exhibited below) shows that the exponent in (31) cannot be raised above $`1/2`$. The fact that one cannot reduce $`\theta `$ below $`1/2`$ in this way was again observed in \[E2\] in the more general context of the ABC conjecture. This fact lends some credence to that conjecture, and thus to its special case (31); this contrasts with the situation for $`\|z^ny^nx^n\|`$, where there is no reason why some polynomial or Pell family might not do better than the $`z^{n3ϵ}`$ expected by probabilistic heuristics, and indeed we found such families for some choices of $`n`$. We next digress to say some more on polynomial and Fermat-Pell families that attain the Davenport-Mason bound, both because they are of independent interest and because families of both kinds appear in our numerical results. In either case $`x^3y^2=k`$ is an identity in a genus-zero function field, namely $`𝐐(t)`$ in the polynomial case and $`𝐐(t,\sqrt{at^2+bt+c})`$ in the Fermat-Pell case. Let $`x,y`$ have degrees $`2m,3m`$ respectively, and suppose $`k`$ has the smallest degree possible, i.e. $`m+1`$ in the polynomial case and $`m`$ for Fermat-Pell. Then $`f:=x^3/y^2`$ is a rational function of degree $`6m`$ or $`12m`$ on $`𝐏^1`$ ramified only above $`0,1,\mathrm{}`$. The Riemann existence theorem provides infinitely many such functions $`f=x^3/y^2`$ in $`𝐂(t)`$; this answers the first part of the question raised in \[BCHS, p.68\]. The second part concerns solutions over $`𝐑`$, and can probably be settled by adding data on complex conjugation to the branched covering. But we are most interested in the third part of the question, in which $`f`$ must have rational coefficients. Given any one $`(x(t),y(t))`$, we may trivially obtain others of the form $`(x^{},y^{})=(\lambda ^2x(t^{}),\lambda ^3y(t^{}))`$ where $`t^{}=at+b`$ in the polynomial case, and $`t^{}𝐐[t]`$ with $`\sqrt{at_{}^{}{}_{}{}^{2}+bt^{}+c}/\sqrt{at^2+bt+c}𝐐[t]`$ in the Fermat-Pell case. If we regard such $`(x^{},y^{})`$ and $`(x,y)`$ as equivalent, only a handful of examples over $`𝐐`$ are known, and there may well be no others. We next list representatives of the known examples. In the polynomial case, all known examples have $`m5`$. For $`m=1`$, translation and scaling brings any quadratic $`x(t)`$ to the form $`t^2+2a`$, and then $`y=t^3+3at`$ and $`k=3a^2t^2+8a^3`$. Necessarily $`a0`$, and all such examples are “twists” of each other, becoming isomorphic over $`\overline{𝐐}`$ if not over $`𝐐`$. Note that $`x^3/y^2`$ is a degree-$`3`$ function of $`t^2`$ with a triple zero. This function occurs for instance as the cover of the modular curve $`\mathrm{X}(1)`$ by $`\mathrm{X}_0(2)`$. For $`m=2`$ we again find that the solution is unique up to twist: $`x=t^4+4at`$, $`y=t^6+6at^3+6a^2`$, and $`k=8a^3t^336a^4`$. This time $`x^3/y^2`$ is a degree-$`4`$ function of $`t^3`$, whose ramification identifies it with the modular cover $`\mathrm{X}_0(3)\mathrm{X}(1)`$. Birch found examples of $`(x,y,k)`$ with $`m=3,5`$ and included them in a 29.ix.1961 letter to Chowla; they are reported in \[BCHS\]: $$(36t^6+24t^4+10t^2+1,216t^9+216t^7+126t^5+35t^3+\frac{21}{4}t,\frac{9}{2}t^4+\frac{39}{16}t^2+1),$$ (32) and $$(\frac{t}{9}\left(t^9+6t^6+15t^3+12\right),\frac{t^{15}}{27}+\frac{t^{12}+4t^9+8t^6}{3}+\frac{5t^3+1}{2},\frac{3t^6+14t^3+27}{108}).$$ (33) These yield integer solutions if $`t`$ is a multiple of $`4`$ in (32) or congruent to $`3`$ mod $`6`$ in (33). As noted in \[BCHS\], the second example provides infinitely many integer solutions of $`\|x^3y^2\|x^{3/5}`$; moreover, for this choice of twist, the leading coefficient of $`k(t)`$ is small enough that $`\|x^3y^2\|`$ is even a respectably small multiple of $`x^{1/2}`$ for the first few specializations of $`t`$. The maps $`f=x^3/y^2`$ associated with Birch’s polynomials both have interesting Galois groups. For (32), $`f`$ is a degree-$`9`$ function of $`t^2`$ whose Galois group is $`\mathrm{PSL}_2(𝐅_8)`$ over $`𝐂(t^2)`$ and $`\mathrm{Aut}(\mathrm{PSL}_2(𝐅_8))`$ over $`𝐐(t^2)`$; the Galois closure is the Fricke-Macbeath curve \[F, M1\]. For (33), $`f`$ is a degree-$`10`$ function of $`t^3`$ whose Galois group is $`\mathrm{PSL}_2(𝐅_9)`$. These groups and curves do not arise in connection with classical modular curves, but they can be identified with certain Shimura modular curves, most naturally those associated with with the $`(2,3,7)`$ and $`(2,3,8)`$ arithmetic triangle groups (see for instance \[T, E5\]). Hall \[H, p.185\] gives an example with $`m=4`$: $$x=4(t^8+6t^7+21t^6+50t^5+86t^4+114t^3+109t^2+74t+28);$$ (34) In August 1998 I announced a new example with $`m=5`$ (its computation will be explained elsewhere): $$x=t^{10}+2t^9+33t^8+12t^7+378t^6336t^5+2862t^42652t^3+14397t^29922t+18553.$$ (35) In both cases (as with all the other $`(x,y,k)`$ examples), $`y`$ is obtained by truncating the Laurent expansion at infinity of $`x^{3/2}`$ after the constant term. Neither (34) nor (35) yields an interesting Galois group: the Galois groups of $`x^3/y^2`$ are $`\mathrm{Alt}_{24}`$ and $`\mathrm{Sym}_{30}`$ respectively. While (35), like (33), must yield infinitely many integer solutions of $`\|x^3y^2\|x^{3/5}`$, the leading coefficient of $`k`$ in (35) makes the implied constant much larger, and none of these solutions will appear in our list of small values of $`\|x^3y^2\|`$. The question, raised in \[BCHS\], whether there are any $`x,y,k𝐐[t]`$ of degrees $`2m,3m,m+1`$ with $`m>5`$, remains unsolved. For Fermat-Pell families, the list is even shorter: all known examples are equivalent, and come from the identity $$(t^2+10t+5)^3(t^2+22t+125)(t^2+4t1)^2=1728t.$$ (36) Here $`y`$ is a multiple of $`\sqrt{at^2+bt+c}`$, so $`f`$ factors as a map of degree $`6`$ composed with the double cover of $`𝐐(t)`$ by $`𝐐(t,\sqrt{at^2+bt+c})`$. We noted in \[E4, p.49\] that the resulting degree-$`6`$ map $`f=x^3/y^2:𝐏^1𝐏^1`$ is the cover $`\mathrm{X}_0(5)\mathrm{X}(1)`$ of classical modular curves. Thus the elliptic curves of low discriminant coming from the identity (36) all admit a rational $`5`$-isogeny. Each Fermat-Pell family obtained from (36) by specifying the class of $`t^2+22t+125`$ mod $`𝐐_{}^{}{}_{}{}^{2}`$ yields $`kCx^{1/2}`$ for some nonzero $`C`$. The smallest such $`C`$ is $`5^{5/2}54=.96598\mathrm{}`$, obtained by Danilov \[D1\] by substituting $`125(2t1)`$ for $`t`$ in (36) and dividing by $`20^3`$: $$(5^5t^23000t+719)^3(5^3t^2114t+26)(5^6t^25^3123t+3781)^2=27(2t1).$$ (37) The factor $`5^3t^2114t+26`$ is a square for $`t=5`$, and thus for infinitely many $`t`$. The first case $`t=5`$ of this yields the elliptic curve of discriminant $`11`$ labeled 11-A2(C) in Cremona’s table \[C2\]; it is known that the isogeny class of this curve provides the examples with minimal conductor of a rational $`5`$-isogeny, and indeed of an elliptic curve over $`𝐐`$. ### 4.2 The new algorithm To obtain numerical data with which to compare Hall’s conjecture, we want to find all small nonzero values of $`\|x^3y^2\|`$ with $`x,y𝐙`$ and $`xX`$. So that we can compare our algorithm with other approaches we briefly review previous work on this problem. The most direct approach is to simply compute for each $`xX`$ the integer $`y`$ closest to $`x^{3/2}`$. Since $`x^{3/2}`$ varies smoothly with $`x`$, this can be done quite efficiently, but clearly must take at least time proportional to $`X`$. This is essentially what Hall did in \[H\], with $`X=710^8`$; some three decades later, faster computers make larger $`X`$ feasible, and indeed Frits Beukers reports in an Aug. 1998 e-mail that he performed such a computation for $`X=10^{12}`$. But this is probably close to the practical limit with today’s technology, and at any rate this direct approach is superseded by the $`X^{1/2}\mathrm{log}^{O(1)}X`$ algorithm described below. A fundamentally different approach is taken in \[GPZ\]: for each nonzero $`k[K,K]`$, investigate the arithmetic of the elliptic curve $`E_k:y^2=x^3k`$, and use effective bounds on integral points to find all integer solutions of $`x^3y^2=k`$. In \[GPZ\], Gebel, Pethö, and Zimmer did most of this work for $`K=10^5`$, except for a few values of $`k`$, for which they could not find a generator for $`E_k(𝐐)`$; Wildanger later showed in his doctoral thesis \[W2\] that none of these $`E_k`$ has an integral point, thus completing the computation of integer solutions of $`0<\|x^3y^2\|<10^5`$. It is not clear even heuristically how this method compares with other approaches. It is the only approach used thus far that will provably find all solutions with $`\|k\|K`$. (The recent proof of the modularity conjecture means that Cremona’s algorithms \[C2\] yield another such approach, but to my knowledge it has not been used to solve $`x^3y^2=k`$.) Assuming Hall’s conjecture, $`\|k\|K`$ is equivalent to $`x_ϵK^{2+ϵ}`$, but this begs the question of the constant implied in “$``$”. Neither do we know how to estimate the average work required to find all integer points on a curve $`E_k`$. It may be reasonable to guess that this average work is proportional to $`K^{c+o(1)}`$ for some $`c>0`$. (This estimate certainly holds for Cremona’s algorithms.) The total work would then be $`K^{1+c+o(1)}`$. Under Hall’s conjecture, this is equivalent to $`X^{(1+c)/2+o(1)}`$, so strictly worse (modulo an unknown implied constant) than our $`X^{1/2}\mathrm{log}^{O(1)}X`$ algorithm, though perhaps better than a direct search, depending on whether $`c<1`$. We noted already that the direct search can exploit the smoothness of the function $`x^{3/2}`$. We can try to take further advantage of this by mimicking our approach to rational approximation of curves: surround the segment $`x<X`$ of the semicubical parabola $`y=x^{3/2}`$ by a union of parallelograms each of area $`O(1)`$, and use lattice reduction to quickly find all integer points in each parallelogram. This does give an asymptotic improvement, though a small one: the parallelogram containing a point $`(x,x^{3/2})`$ has length $`x^{1/6}`$, so the computational cost is reduced by at most $`X^{1/6}`$, to $`X^{5/6}\mathrm{log}^{O(1)}X`$. We reduce the exponent of $`X`$ from $`1`$ or $`5/6`$ to $`1/2`$ by a more radical reorganization of the computation that lets us apply lattice reduction more efficiently. More generally, for each positive $`c𝐐`$ we can find all cases of $`0<\|cx^3y^2\|x`$ in time $`O_c(X^{1/2}\mathrm{log}^{O(1)}X)`$. All choices of $`c`$ are essentially equivalent: we get from one to the other by scaling $`x,y`$ and imposing congruence conditions on them. The most convenient choice of $`c`$ turns out to be $`4/3`$. We thus show how to solve $`0<\|4x^33y^2\|x`$; the cases relevant to Hall’s conjecture are those with $`3\|x`$ and $`6\|y`$, when $`(4x^33y^2)/108=(x/3)^3(y/6)^2`$. We begin as in \[H\] by approximating $`x,y`$ by (multiples of) a square and a cube. Any positive integer $`x`$ may be written uniquely as $$x=3\zeta ^2+\eta \mathrm{with}\eta ,\zeta 𝐙,\zeta >0,\eta (3\zeta ,3\zeta ].$$ (38) Then $$(4x^3/3)^{1/2}=6\zeta ^3+3\eta \zeta +\frac{1}{4}\frac{\eta ^2}{\zeta }\frac{1}{72}\left(\frac{\eta }{\zeta }\right)^3+O(1/\zeta ).$$ (39) We thus write $$y=6\zeta ^3+3\eta \zeta +\xi $$ (40) with $`\xi \zeta `$. More precisely, if $$\eta =\beta \zeta $$ (41) Then $`\beta (3,3]`$, and $`\|4x^33y^2\|x`$ if and only if $$\xi =\frac{\eta ^2}{4\zeta }\frac{1}{72}\beta ^3+O(1/\zeta ).$$ (42) At this point, Hall \[H\] imposes the assumption $`\beta \xi ^{1/5}`$. We allow an arbitrary $`\beta (3,3]`$ and approximate it within $`O(X^{1/2})`$ by one of $`O(X^{1/2})`$ evenly spaced points in that interval. Suppose, then, that $`b`$ is one of those points. We approximate (42) by a linear combination of $`\zeta `$, $`\eta b\zeta `$, and $`1`$: $$\xi =\frac{b^2}{4}\zeta +\frac{b}{2}(\eta b\zeta )\frac{b^3}{72}+O(1/\zeta )=\frac{b^2}{4}\zeta +\frac{b}{2}\eta \frac{b^3}{72}+O(1/\zeta ).$$ (43) We now assume that $`\zeta X^{1/2}`$, for instance by requiring that $`x>X/4`$; repeating the computation with $`X`$ replaced by $`X/4`$, $`X/16`$, $`X/64`$,… will then cover the entire range $`xX`$, and if we can cover $`(X/4,X]`$ in time $`O(x^{1/2}\mathrm{log}^{O(1)}X)`$ then the same is true of $`[1,X]`$. Under the assumption $`x(X/4,X]`$, we have the following constraints on $`\xi ,\eta ,\zeta `$: $$\zeta X^{1/2},\eta b\zeta 1,$$ (44) and $$\xi +\frac{b^2}{4}\zeta \frac{b}{2}\eta +\frac{b^3}{72}X^{1/2}.$$ (45) We are thus in a familiar situation: we seek all the integral points in $`O(X^{1/2})`$ parallelepipeds, each of volume $`O(1)`$. The term $`b^3/72`$ in (45) means that the parallelepipeds are no longer centered at the origin, but this causes no difficulty — indeed we already dealt with off-center parallelepipeds in the practical implementation of our algorithm for finding rational points near curves. So again we linearly transform each parallelepiped to a cube and obtain a lattice reduction problem; if these lattices were randomly distributed among three-dimensional lattices, we would almost certainly have only $`O(X^{1/2})`$ points to try, and would thus find all solutions of $`0<\|4x^33y^2\|x`$ with $`xX`$ in time $`O(X^{1/2}\mathrm{log}^{O(1)}X)`$. In fact it turns out that in this case our lattices are not equidistributed: they all lie in a $`2`$-dimensional subspace of the $`5`$-dimensional moduli space of lattices in $`𝐑^3`$. This gives rise to both a minor annoyance and a major advantage. The bad news is that we cannot expect our lattices to have on average $`O(1)`$ vectors of norm $`1`$; but this annoyance is minor because the actual average is proportional to $`\mathrm{log}X`$ and thus can be absorbed into the $`\mathrm{log}^{O(1)}X`$ factor. The good news is that we understand our special lattices well enough to actually prove results that are only heuristic for rational points near curves. The key is that in each case our lattice is a symmetric square of a lattice in $`𝐑^2`$. By this we mean the following. Recall that the symmetric square of a $`2`$-dimensional vector space $`V`$ is the $`3`$-dimensional vector space $`\mathrm{Sym}^2V`$ consisting of symmetric tensors in $`VV`$. Since $`\mathrm{Sym}^2V`$ is defined naturally in terms of $`V`$, any linear transformation of $`V`$ yields a linear transformation of $`\mathrm{Sym}^2V`$. We thus have a homomorphism $`\mathrm{Sym}^2:\mathrm{GL}_2\mathrm{GL}_3`$. To give this map explicitly we choose a basis $`(e_1,e_2)`$ for $`V`$, and use the basis $`(e_1e_1,(e_1e_2+e_2e_1)/2,e_2e_2)`$ for $`\mathrm{Sym}^2V`$. We then calculate that $$\mathrm{Sym}^2\left(\begin{array}{cc}p& q\\ r& s\end{array}\right)=\left(\begin{array}{ccc}p^2& pq& q^2\\ 2pr& ps+qr& 2qs\\ r^2& rs& s^2\end{array}\right).$$ (46) Over any field, $`\mathrm{Sym}^2(\mathrm{SL}_2)`$ is contained in the subgroup of $`\mathrm{SL}_3`$ preserving the discriminant form $`4a_1a_3a_2^2`$ on $`\mathrm{Sym}^2(V)`$; if we worked over an algebraically closed field, that subgroup would coincide with $`\mathrm{Sym}^2(\mathrm{SL}_2)`$. Now (44,45) mean that the column vector $`v=(\xi ,\eta ,\zeta )𝐙^3`$ satisfies $`M_bvu_b1`$ where $`u_b:=(0,0,b^3/72)`$ and $$M_b:=\left(\begin{array}{ccc}0& 0& X^{1/2}\\ 0& 1& b\\ X^{1/2}& X^{1/2}b/2& X^{1/2}b^2/4\end{array}\right)=\mathrm{Sym}^2\left(\begin{array}{cc}0& X^{1/4}\\ X^{1/4}& X^{1/4}b/2\end{array}\right).$$ (47) This is why we went after $`4x^33y^2`$ rather than pursuing $`x^3y^2`$ directly: an analogous approach to $`x^3y^2`$ would yield a matrix that is still a symmetric square but with respect to a different basis, requiring a definition of $`\mathrm{Sym}^2`$ with fractional coefficients and complicating the lattice reduction. Note that the quadratic form $`4\xi \zeta \eta ^2`$ preserved by $`M_b\mathrm{Sym}^2(\mathrm{SL}_2)`$ is already visible in (42). Our algorithm, then, is as follows. For each of our $`O(X^{1/2})`$ choices of $`b`$, calculate the matrix $$N_b:=\left(\begin{array}{cc}0& X^{1/4}\\ X^{1/4}& X^{1/4}b/2\end{array}\right)$$ (48) with $`M_b=\mathrm{Sym}^2=N_b`$. Use lattice reduction to find a matrix $`K_b\mathrm{GL}_2(𝐙)`$ such that $`N_bK_b`$ is as small as possible. Then $$M_b^{}:=\mathrm{Sym}^2(N_bK_b)=\mathrm{Sym}^2N_b\mathrm{Sym}^2K_b=M_b\mathrm{Sym}^2K_b.$$ (49) is small too. Let $`L_b=\mathrm{Sym}^2K_b\mathrm{GL}_3(𝐙)`$. Then $`M_bv=M_b^{}L_b^1v`$. Find a box containing all $`w𝐙^3`$ such that $`M_b^{}wu_b1`$. For each $`w`$ in the box, compute $`v=L_bw`$ and check whether the resulting $`x,y`$ satisfy $`x(X/4,X]`$ and $`0<\|4x^33y^2\|x`$; if they do, output $`x`$ (and check whether $`3\|x`$ and $`6\|y`$ to determine whether this solution also yields a small value of $`x^3y^2`$). This is easier than our usual algorithm because we are reducing a lattice in $`𝐑^2`$ rather than $`𝐑^3`$, which in our case amounts to calculating the continued fraction of $`b/X^{1/2}`$. Moreover, the computational cost of the algorithm can be bounded rigorously: $`M_b^{}`$ will only be large if $`b/X^{1/2}`$ is close to a rational number with numerator and denominator $`X^{1/4}`$, and the effect of such a close rational approximation is easy to determine. Summing over all rationals of height $`X^{1/4}`$ we find that the total number of candidate vectors $`v`$ is $`X^{1/2}\mathrm{log}X`$, and thus that the computation takes time $`O(X^{1/2}\mathrm{log}^{O(1)}X)`$ as claimed. Note that the $`X^{1/2}\mathrm{log}X`$ bound also has the following consequence: there are $`X^{1/2}\mathrm{log}X`$ solutions of $`\|x^3y^2\|x`$ with $`xX`$. Moreover, if $`C`$ is large enough, we can deduce from this analysis that there are $`X^{1/2}`$ solutions of $`0<\|x^3y^2\|<Cx`$ with $`x[X/2,X]`$. More generally, we show that for each positive $`c𝐐`$ there exists $`C`$ such that for each $`r𝐑/𝐙`$ and $`d>1`$ there are at most $`CdX^{1/2}\mathrm{log}X`$ solutions of $`\|(cx^3)^{1/2}(y+r)\|<dX^{1/2}`$ with $`x,y𝐙`$ and $`x<X`$; and, given $`c`$ as above and any $`\theta [0,1)`$, there exists $`C_0`$ such that for any $`r𝐑/𝐙`$ there are $`X^{1/2}`$ solutions of $`\|(cx^3)^{1/2}(y+r)\|<C_0X^{1/2}`$ with $`x,y𝐙`$ and $`x[\theta X,X]`$. The constants $`C,C_0`$ depend effectively on $`c,\theta `$. These results improve considerably on results in this direction available from general exponential-sum techniques for proving uniform distribution mod $`1`$. The detailed proofs of our claims in this paragraph will appear elsewhere. ### 4.3 Numerical results We have implemented our algorithm in a C program using $`64`$-bit integer arithmetic, again replacing each $`O(\mathrm{})`$ and $``$ by explicit bounds, and searched for all solutions of $`0<\|4x^33y^2\|<200x^{1/2}`$ with $`410^6<x<310^{18}`$. The range $`x<10^{10}`$ was covered by a direct search, the overlap $`[410^6,10^{10}]`$ being used as a check on the computation. The code was processed with an optimizing compiler and ran for three weeks during the summer of 1998 on a Sun Sparcstation Ultra $`1`$. As a corollary we obtained all cases of $`0<\|x^3y^2\|<\frac{1}{2}\sqrt{x}`$ with $`x<10^{18}`$. (With currently available hardware the same computation could easily finish in a few days; with parallelization it should be feasible to reach $`10^{23}`$ at least.) The next table lists, for each of the $`25`$ solutions of $`0<\|x^3y^2\|<\frac{1}{2}\sqrt{x}`$, the values of $`k=x^3y^2`$, $`x`$, and $`r=x^{1/2}/\|k\|`$. We need not list $`y`$, which is always the integer nearest to $`x^{3/2}`$. The explanation of the last two columns follows the table. # $`k`$ $`x`$ $`r`$ GPZ? Comments 1 1641843 5853886516781223 46.60 !! 2 30032270 38115991067861271 6.50 ! 3 $``$1090 28187351 4.87 + 4 $``$193234265 810574762403977064 4.66 5 $``$17 5234 4.26 + $`P(3)`$ 6 $``$225 20114 3.77 + 7 $``$24 8158 3.76 + $`P(3)`$ 8 307 939787 3.16 + 9 207 367806 2.93 + 10 $``$28024 3790689201 2.20 + 11 $``$117073 65589428378 2.19 12 $``$4401169 53197086958290 1.66 13 105077952 23415546067124892 1.46 * 14 $`1`$ 2 1.41 15 $``$497218657 471477085999389882 1.38 16 $``$14668 384242766 1.34 + $`P(9)`$ 17 $``$14857 390620082 1.33 + $`P(9)`$ 18 $``$87002345 12813608766102806 1.30 19 2767769 12438517260105 1.27 20 $``$8569 110781386 1.23 + 21 5190544 35495694227489 1.15 22 $``$11492 154319269 1.08 + 23 $``$618 421351 1.05 + 24 548147655 322001299796379844 1.04 D 25 $``$297 93844 1.03 + D The “GPZ” column indicates whether the solution was among the $`13`$ listed in \[GPZ\]. These are the solutions with $`1<\|k\|<10^5`$. Presumably the solution $`2^33^2=1`$ is not on that list because the elliptic curve $`y^2=x^3+1`$ was already known to have rank $`0`$ so Gebel, Pethö and Zimmer were not interested in it. The #1 row is a new record, improving the previous record $`r`$ by a factor of almost $`10`$, whence the notation “!!”. Even row #2, marked “!”, has $`r`$ larger than the old record which is row #3. Either of this suffices to refute Hall’s comment \[H, p.175\], repeated in \[GPZ\], that $`r<5`$ seems to hold in all cases. \*: Obtained from row #1 by scaling $`(x,y,k)`$ to $`(2^2x,2^3y,2^6k)`$. This reduces $`r`$ by a factor of $`32`$, but $`r=46+`$ in row #1 is large enough that even $`r/32`$ still exceeds the threshold of our table. $`P(t)`$: Birch’s polynomial family (33). This has $`r=12/t+O(t^4)`$, so the only values of $`t3mod6`$ that appear on the $`r>1`$ list are $`t=\pm 3`$ and $`\pm 9`$. Already in \[BCHS, p.69\] the specializations $`t=\pm 3`$ are noted as “striking special cases” of (33). D: The first two cases of Danilov’s family (37). The appearance of the larger of these was a welcome check on our computation. Any threshold on $`r`$ is of necessity arbitrary; the next solution has $`r`$ just below our cutoff of $`1`$: $`(x,k,r)=(16544006443618,4090263,0.9944\mathrm{})`$.
warning/0005/nucl-th0005002.html	ar5iv	text	# The Half-Life and the Nuclear Matrix Element of the 2⁢𝜈⁢𝛽⁢𝛽 Decay in ⁷⁶𝐺⁢𝑒 The project supported by National Natural Science Foundation of China and Grant of Academia Sinica. ## (I) Introduction The two-neutrino double beta $`(2\nu \beta \beta )`$ decay is allowed in the standard theory of the electroweak interaction and has been observed in the laboratory<sup></sup>. A number of experiments have reported the direct observation of the $`2\nu \beta \beta `$ modes. For instance, Elliott and collaborators first observed the $`2\nu \beta \beta `$ decay of $`{}_{}{}^{82}Se^{[2]}`$ in 1986. Ejili et al and Alston-Garnjost et al have observed the $`2\nu \beta \beta `$ decay of $`{}_{}{}^{100}Mo^{[3]}`$. And the $`2\nu \beta \beta `$ decay of $`{}_{}{}^{76}Ge`$ have been observed by three groups-Russian<sup></sup>, American<sup></sup> and Heidelberg-Moscow<sup></sup> groups. The future of double beta experiments will be dominated by the use of enriched detectors. Among them, $`{}_{}{}^{76}Ge`$ plays a particular favorable role. Theoretically, one can predict the decay half-lifes with the calculated nuclear matrix elements,since it is parameter-free from the particle physics side. The experiment for the $`2\nu \beta \beta `$ decay will test the theoretical prediction of the nuclear matrix elements. However the situation is different for the neutrinoless double beta $`(0\nu \beta \beta `$ decay, where no neutrino is emitted and lepton number is violated. These processes can only occur, as there is an exchange of Majorana neutrino. Consequently the $`0\nu \beta \beta `$ decay half-life depends besides on the nuclear matrix elements also on unknown parameters of the decay mechanism. Furthermore, both modes are related to each other by the nuclear matrix elements. Thus one can use the same wave functions of nuclei to calculate the nuclear matrix elements for the $`2\nu \beta \beta `$ and $`0\nu \beta \beta `$-decays. The availabe experimental data of the $`2\nu \beta \beta `$ decay will test $`2\nu \beta \beta `$ nuclear matrix elements and tell us how reliable the wave functions are. With a more accurate wave-functions one can extract the neutrino mass information from the experimental half-life limits of the $`0\nu \beta \beta `$ decay. The information about the neutrino mass and neutrino mixing attracts the most attention due to the recent experimental progress<sup></sup>. In this paper, following the shell-model wave-functions of $`{}_{}{}^{76}Ge`$ and $`{}_{}{}^{76}Se`$ which were discussed in the previous publications<sup></sup>, we evaluate the half-life and the nuclear matrix elements of the $`2\nu \beta \beta `$ decay in $`{}_{}{}^{76}Ge`$. The latest Heidelberg-Moscow measurement for the $`2\nu \beta \beta `$ decay of $`{}_{}{}^{76}Ge`$ provides a test to these wave-functions. This comparison confirms our nuclear wave-functions that apply to the $`0\nu \beta \beta `$ decay of $`{}_{}{}^{76}Ge`$. The result shows that the theoretical calculation from our wave-functions is comparable with the experimental data of the $`2\nu \beta \beta `$ decay $`{}_{}{}^{76}Ge`$ and the upper limit of the neutrino mass is less than 0.4eV from the $`0\nu \beta \beta `$ decay of $`{}_{}{}^{76}Ge`$. ## (II) Calculated matrix element of $`2\nu \beta \beta `$ decay $`{}_{}{}^{76}Ge`$ The $`2\nu \beta \beta `$ mode has two electrons and two antineutrinos in the final states and is expected to appear in the standard model in the second order process of weak interaction, such as $`{}_{}{}^{76}Ge^{76}Se+2e^{}+2\stackrel{~}{\nu }_e.`$ (1) The decay amplitude associated with this process takes the form<sup></sup> $`𝒥_{2\nu }`$ $`=`$ $`E_m[{\displaystyle \frac{<e_1e_2\nu _1\nu _2\psi _fH_\beta e_1\nu _1\psi _m><e_1\nu _1\psi _mH_\beta \psi _i>}{E_iE_{e_1}E_{\nu _1}E_m}}`$ (2) $`=`$ $`{\displaystyle \frac{<e_1e_2\nu _1\nu _2\psi _fH_\beta e_2\nu _2\psi _m><e_2\nu _2\psi _mH_\beta \psi _i>}{E_iE_{e_2}E_{\nu _2}E_m}}(e_1e_2)],`$ where $`H_\beta `$ is the hamiltonian of the weak interaction. $`H_\beta ={\displaystyle \frac{G}{\sqrt{2}}}{\displaystyle J_\mu (x)L_\mu (x)d^3x}.`$ (3) The $`\psi _i,\psi _m`$ and $`\psi _f`$ in Eq.(2) are the initial, intermediate and final state wave-functions, respectively, $`E_i`$ and $`E_m`$ are the inital and intermediate state energies. $`G,L_\mu `$ and $`J_\mu `$ in Eq.(3) are Fermi coupling constant, the leptonic and hadronic currents respectively. $`L_\mu `$ and $`J_\mu `$ can be expressed as $`L_\nu =\overline{e}\gamma _\mu (1\gamma _5)\nu `$ (4) and $`J_\mu =\overline{N}\gamma _\mu (F_VF_A\gamma _5)\tau _+N,`$ (5) with $`F_V=1.0`$ and $`F_A=1.25`$. Employing the usual closure approximation to the intermediate states, the half-life $`T_{1/2}^{2\nu }`$ is given by $`T_{1/2}^{2\nu }={\displaystyle \frac{1n2}{f_{GT}M_{GT}^2}}`$ (6) where $`M_{GT}=<\psi _f\mathrm{\Sigma }_{i<j}\tau _+(i)\tau _+(j)\sigma (i)\sigma (j)\psi _i>`$ (7) and $`f_{GT}`$ $`=`$ $`\xi _{2\nu }{\displaystyle \frac{2m_e^{11}G^4}{7!\pi ^7}}\left[{\displaystyle \frac{F^{PR}(Z)}{E_i\overline{E}_mT_0/2m_e}}\right]^2\times `$ (8) $`\times `$ $`\left[1+{\displaystyle \frac{\stackrel{~}{T}_0}{2}}+{\displaystyle \frac{\stackrel{~}{T}_0^2}{9}}+{\displaystyle \frac{\stackrel{~}{T}_0^3}{90}}+{\displaystyle \frac{\stackrel{~}{T}_0^4}{1980}}\right]\times \stackrel{~}{T}_0^7`$ Here $`\stackrel{~}{T}_0=T_0/m_e`$ with $`T_0`$ being the total kinetic energy of the outgoing electrons and $`m_e`$ being the electron mass. For $`{}_{}{}^{76}Ge,T_0=2.045MeV.\overline{E}_m`$ is the average energy of the intermediate nuclear states. According to the statistics study of the $`\beta `$ decay, the average intermediate energy is chosen reasonably and $`\overline{E}_mE_i=7.88MeV`$ in the case of $`{}_{}{}^{76}Ge^{[10]}.F_{PR}(Z)`$ is the nonrelativistic Coulomb correction factor<sup></sup> and its use permits analytic evaluation of the phase space integrals appearing in Eq.(8). $`\xi _{2\nu }`$ in Eq.(8) represents the difference between this approximation and an exact integrals of the phase space. For the $`2\nu \beta \beta `$ decay of $`{}_{}{}^{76}Ge,\xi _{2\nu }=1.65^{[10]}`$. We have calculated the wave-functions of $`{}_{}{}^{76}Ge`$ and $`{}_{}{}^{76}Se`$ by using a simple shell-model structure$`{}_{}{}^{[8,9]}.Ni^{56}`$ is adopted as an inert core for $`{}_{}{}^{76}Ge`$ and $`{}_{}{}^{76}Se`$. The shell-model space for these nuclei includes four single particle particle orbits $`\{1P_{3/2},0f_{5/2},1p_{1/2},0g_{9/2}\}`$. We employed the modified surface delta interaction (MSDI)<sup></sup> as the two body residual interaction. The single-particle energies and first set of parameters of parameters of MSDI in the Table 1 were taken from Faessler et al<sup></sup>. The second set of parameters of MSDI was used to the $`2\nu \beta \beta `$ decay of $`{}_{}{}^{82}Se`$ in our previous paper<sup></sup>. Table 1. The nuclear matrix elements and the half-lifes of the $`2\nu \beta \beta `$ decay in $`{}_{}{}^{76}Ge`$ Corresponding the different parameters of MSDI. \| MSDI \| \| \| \| \| $`\epsilon _i`$(MeV) \| \| \| \| \| \| \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| \| \| \| \| \| \| \| \| \| \| $`M_{GT}`$ \| $`T_{1/2}^{2\nu }(th)`$ \| \| \| $`A_0`$ \| $`A_1`$ \| B \| C \| $`1p_{3/2}`$ \| $`0f_{5/2}`$ \| $`1p_{1/2}`$ \| $`0g_{9/2}`$ \| \| $`\times 10^{21}(years)`$ \| \| 1 \| 0.43 \| 0.35 \| 0.33 \| 0.00 \| 0.00 \| 1.75 \| 2.20 \| 3.39 \| 0.353 \| 5.17 \| \| 2 \| 0.31 \| 0.30 \| 0.30 \| 0.00 \| 0.00 \| 1.75 \| 2.20 \| 3.39 \| 0.479 \| 2.82 \| ## (III) Comparison between the calculated result and the experimental data The recent experimental data for $`2\nu \beta \beta `$-and $`0\nu \beta \beta `$-decay modes of $`{}_{}{}^{76}Ge`$ are given, $`T_{1/2}^{2\nu }(exp)=(1.77_{0.010.11}^{+0.01+0.13})\times 10^{21}years^{[6]}`$ (9) and $`T_{1/2}^{0\nu }(exp)>5.7\times 10^{25}years^{[7]}`$ (10) The calculated result (see the table 1.) is comparable to the observed $`2\nu \beta \beta `$-decay half-life of $`{}_{}{}^{76}Ge`$ and the half-life from the second line is close to the observed value in Eq.(9). Thereore, one can infer that our wave-functions of $`{}_{}{}^{76}Ge`$ and $`{}_{}{}^{76}Se`$ may have more satisfied construction with the reasonable parameter choices. With the same wave-functions one can obtain that the effective neutrino mass is less than 0.4eV by using the latest experimental lower limit (Eq.(10)) for the $`0\nu \beta \beta `$-deday half-life of $`{}_{}{}^{76}Ge`$. As we know, the result of neutrino oscillation experiments on the effective neutrino mass<sup></sup> gives the upper limit may be around $`10^2eV`$. In order to approach this limit from the nuclear $`0\nu \beta \beta `$ decay we suggest to improve the existing experimental limit of the $`0\nu \beta \beta `$-decay half-life by using the enriched $`{}_{}{}^{76}Ge`$. More efforts to improve the theoretical calculation on the nuclear matrix element of $`{}_{}{}^{76}Ge`$ are also necessary.
warning/0005/physics0005031.html	ar5iv	text	# Untitled Document A Nonstandard Derivation for the Special Theory of Relativity<sup></sup><sup></sup>This is an expanded version of the paper presented at the 14 NOV 1992 Meeting of the Mathematical Association of America, Coppin State College, Baltimore MD and as it appears in Herrmann, R. A., Special relativity and a nonstandard substratum, Speculat. Sci. Technol., 17(1994), 2-10. Robert A. Herrmann Mathematics Department U. S. Naval Academy 572C Holloway Rd Annapolis, MD 21402-5002 USA 21 SEP 1993 Abstract: Using properties of the nonstandard physical world, a new fundamental derivation for effects of the Special Theory of Relativity is given. This fundamental derivation removes all the contradictions and logical errors in the original derivation and leads to the fundamental expressions for the Special Theory Lorentz transformation. Necessarily, these are obtained by means of hyperbolic geometry. It is shown that the Special Theory effects are manifestations of the interaction between our natural world and a nonstandard medium, the NSPPM. This derivation eliminates the controversy associated with any physically unexplained absolute time dilation and length contraction. It is shown that there is no such thing as a absolute time dilation and length contraction but, rather, alterations in pure numerical quantities associated with an electromagnetic interaction with an NSP-world NSPPM. 1. The Fundamental Postulates. There are various Principles of Relativity. The most general and least justified is the one as stated by Dingle “There is no meaning in absolute motion. By saying that such motion has no meaning, we assert that there is no observable effect by which we can determine whether an object is absolutely at rest or in motion, or whether it is moving with one velocity or another.”\[1:1\] Then we have Einstein’s statements that “I. The laws of motion are equally valid for all inertial frames of reference. II. The velocity of light is invariant for all inertial systems, being independent of the velocity of its source; more exactly, the measure of this velocity (of light) is constant, $`c`$, for all observers.”\[7:6–7\] I point out that Einstein’s original derivation in his 1905 paper (Ann. der Phys. 17: 891) uses certain well-known processes related to partial differential calculus. In 1981 and 1991 , it was discovered that the intuitive concepts associated with the Newtonian laws of motion were inconsistent with respect to the mathematical theory of infinitesimals when applied to a theory for light propagation. The apparent nonballistic nature for light propagation when transferred to infinitesimal world would also yield a nonballistic behavior. Consequently, there is an absolute contradiction between Einstein’s postulate II and the derivation employed. This contradiction would not have occurred if it had not been assumed that the æther followed the principles of Newtonian physics with respect to electromagnetic propagation. \[Note: On Nov. 14, 1992, when the information in this article was formally presented, I listed various predicates that Einstein used and showed the specific places within the derivations where the predicate’s domain was altered without any additional argument. Thus, I gave specific examples of the model theoretic error of generalization.\] I mention that Lorentz speculated that æther theory need not correspond directly to the mathematical structure but could not show what the correct correspondence would be. Indeed, if one assumes that the Nonstandard Photon-particle Medium, the NSPPM, satisfies the most basic concept associated with an inertial system that a body can be considered in a state of rest or uniform motion unless acted upon by a force, then the expression $`F=ma,`$ among others, may be altered for infinitesimal NS-substratum behavior. Further, the NS-substratum, when light propagation is discussed, does not follow the Galilean rules for velocity composition. The additive rules are followed but no negative real velocities exterior to the Euclidean monads are used since we are only interested in the propagation properties for electromagnetic radiation. The derivation in section 3 removes all contradictions by applying the most simplistic Galilean properties of motion, including the ballistic property, but only to behavior within a Euclidean monad. As discussed in section 3, the use of an NSP-world (i.e. nonstandard physical world) NSPPM allows for the elimination of the well-known Special Theory “interpretation” contradiction that the mathematical model uses the concepts of Newtonian absolute time and space, and, yet, one of the major interpretations is that there is no such thing as absolute time or absolute space. Certain general principles for NSPPM light propagation will be specifically stated in section 3. These principles can be gathered together as follows: (1) There is a portion of the nonstandard photon-particle medium - the NSPPM - that sustains N-world (i.e. natural = physical world) electromagnetic propagation. Such propagation follows the infinitesimally presented laws of Galilean dynamics, when restricted to monadic clusters, and the monadic clusters follow an additive and an actual metric property for linear relative motion when considered collectively. \[The term “nonstandard electromagnetic field” should only be construed as a NSPPM notion, where the propagation of electromagnetic radiation follows slightly different principles than within the natural world.\] (2) The motion of light-clocks within the N-world (natural world) is associated with one single effect. This effect is an alteration in an appropriate light-clock mechanism. \[The light-clock concept will be explicitly defined at the end of section 3.\] It will be shown later that an actual physical cause may be associated with vreified Special Theory physical alterations. Thus the Principle of Relativity, in its general form, and the inconsistent portions of the Einstein principles are eliminated from consideration and, as will be shown, the existence of a special type of medium can be assumed without contradicting experimental evidence. In modern Special Theory interpretations , it is claimed that the effect of “length contraction” has no physical meaning, whereas time dilation does. This is probably true if, indeed, the Special Theory is actually based upon the intrinsic N-world concepts of length and time. What follows will further demonstrate that the Special Theory is a light propagation theory, as has been previously argued by others, and that the so-called “length contraction” and time dilation can both be interpreted as physically real effects when they are described in terms of the NSPPM. The effects are only relative to a theory of light propagation. 2. Pre-derivation Comments. Recently –, nonstandard analysis has proved to be a very significant tool in investigating the mathematical foundations for various physical theories. In 1988 , we discussed how the methods of nonstandard analysis, when applied to the symbols that appear in statements from a physical theory, lead formally to a pregeometry and the entities termed as subparticles. One of the goals of NSP-world research is the re-examination of the foundations for various controversial N-world theories and the eventual elimination of such controversies by viewing such theories as but restrictions of more simplistic NSP-world concepts. This also leads to indirect evidence for the actual existence of the NSP-world. The Special Theory of Relativity still remains a very controversial theory due to its philosophical implications. Prokhovnik produced a derivation that yields all of the appropriate transformation formulas based upon a light propagation theory, but unnecessarily includes an interpretation of the so-called Hubble textural expansion of our universe as an additional ingredient. The new derivation we give in this article shows that properties of a NSPPM also lead to Prokhovnik’s expression (6.3.2) in reference and from which all of the appropriate equations can be derived. However, rather than considering the Hubble expansion as directly related to Special Relativity, it is shown that one only needs to consider simplistic NSP-world behavior for light propagation and the measurement of time by means of N-world light-clocks. This leads to the conclusion that Special Theory effects may be produced by a dense NSPPM within the NSP-world. Such an NSPPM – an æther – yields N-world Special Theory effects. 3. The derivation The major natural system in which we exist locally is a space-time system. “Empty” space-time has only a few characterizations when viewed from an Euclidean perspective. We investigate, from the NSP-world viewpoint, electromagnetic propagation through a Euclidean neighborhood of space-time. Further, we assume that light is such a propagation. One of the basic precepts of infinitesimal modeling is the experimentally verified simplicity for such a local system. For actual time intervals, certain physical processes take on simplistic descriptions. These NSP-world descriptions are represented by the exact same description restricted to infinitesimal intervals. Let $`[a,b],ab,a>0,`$ be an objectively real conceptional time interval and let $`t(a,b).`$ The term “time” as used above is very misunderstood. There are various viewpoints relative to its use within mathematics. Often, it is but a term used in mathematical modeling, especially within the calculus. It is a catalyst so to speak. It is a modeling technique used due to the necessity for infinitesimalizing physical measures. The idealized concept for the “smoothed out” model for distance measure appears acceptable. Such an acceptance comes from the use of the calculus in such areas as quantum electrodynamics where it has great predictive power. In the subatomic region, the assumption that geometric measures have physical meaning, even without the ability to measure by external means, is justified as an appropriate modeling technique. Mathematical procedures applied to regions “smaller than” those dictated by the uncertainty principle are accepted although the reality of the infinitesimals themselves need not be assumed. On the other hand, for this modeling technique to be applied, the rules for ideal infinitesimalizing should be followed. The infinitesimalizing of ideal geometric measures is allowed. But, with respect to the time concept this is not the case. Defining measurements of time as represented by the measurements of some physical periodic process is not the definition upon which the calculus is built. Indeed, such processes cannot be infinitesimalized. To infinitesimalize a physical measurement using physical entities, the entities being observed must be capable of being smoothed out in an ideal sense. This means that only the macroscopic is considered, the atomic or microscopic is ignored. Under this condition, you must be able to subdivide the device into “smaller and smaller” pieces. The behavior of these pieces can then be transferred to the world of the infinitesimals. Newton based the calculus not upon geometric abstractions but upon observable mechanical behavior. It was this mechanical behavior that Newton used to define physical quantities that could be infinitesimalized. This includes the definition of “time.” All of Newton’s ideas are based upon velocities as the defining concept. The notation that uniform (constant) velocity exists for an object when that object is not affected by anything, is the foundation for his mechanical observations. This is an ideal velocity, a universal velocity concept. The modern approach would be to add the term “measured” to this mechanical concept. This will not change the concept, but it will make it more relative to natural world processes and a required theory of measure. This velocity concept is coupled with a smoothed out scale, a ruler, for measurement of distance. Such a ruler can be infinitesimalized. From observation, Newton then infinitesimalized his uniform velocity concept. This produces the theory of fluxions. Where does observer time come into this picture? It is simply a defined quantity based upon the length and velocity concept. Observationally, it is the “thing” we call time that has passed when a test particle with uniform velocity first crosses a point marked on a scale and then crosses a second point marked on the same scale. This is in the absence of any physical process that will alter either the constant velocity or the scale. Again this definition would need to be refined by inserting the word “measured.” Absolute time is the concept that is being measured and cannot be altered as aconcept. Now with Einstein relativity, we are told that measured quantities are effected by various physical processes. All theories must be operational in that the concept of measure must be included. But, the calculus is used. Indeed, used by Einstein in his original derivation. Thus, unless there is an actual physical entity that can be substituted for the Newton’s ideal velocity, then any infinitesimalizing process would contradict the actual rules of application of the calculus to the most basic of physical measures. But, the calculus is used to calculate the measured quantities. Hence, we are in a quandary. Either there is no physical basis for mathematical models based upon the calculus, and hence only selected portions can be realized while other selected portions are simply parameters not related to reality in any manner, or the calculus is the incorrect mathematical structure for the calculations. Fortunately, nature has provided us with the answer as to why the calculus, when properly interpreted, remains such a powerful tool to calculate the measures that describe observed physical behavior. In the 1930s, it was realized that the measured uniform velocity of the to-and-fro velocity of electromagnetic radiation, (i.e. light) is the only known natural entity that will satisfy the Newtonian requirements for an ideal velocity and the concepts of space-time and from which the concept of time itself can be defined. The first to utilize this in relativity theory was Milne. This fact I learned after the first draughts of this paper were written and gives historical verification of this paper’s conclusions. Although, it might be assumed that such a uniform velocity concept as the velocity of light or light paths in vacuo cannot be infinitesimalized, this is not the case. Such infinitesimalizing occurs for light-clocks and from the simple process of “scale changing” for a smoothed out ruler. What this means is that, at its most basic physical level, conceptually absolute or universal Newton time can have operational meaning as a physical foundation for a restricted form of “time” that can be used within the calculus. As H. Dingle states it, “The second point is that the conformability of light to Newton mechanics . . . makes it possible to define corresponding units of space and time in terms of light instead of Newton’s hypothetical ‘uniformly moving body.’ ” \[The Relativity of Time, Nature, 144(1939): 888–890.\] It was Milne who first (1933) attempted, for the Special Theory, to use this definition for a “Kinematic Relativity” \[Kinematic Relativity, Oxford University Press, Oxford, 1948\] but failed to extend it successfully to the space-time environment. In what follows such an operational time concept is being used and infinitesimalized. It will be seen, however, that based upon this absolute time concept another time notion is defined, and this is the actual time notion that must be used to account for the physical changes that seem to occur due to relativistic processes. In practice, the absolute time is eliminated from the calculations and is replaced by defined “Einstein time.” It is shown that Einstein time can be infinitesimalized through the use of the definable “infinitesimal light-clocks” and gives an exact measurement. Our first assumption is based entirely upon the logic of infinitesimal analysis, reasoning, modeling and subparticle theory. (i) “Empty” space within our universe, from the NSP-world viewpoint, is composed of a dense-like nonstandard medium (the NSPPM) that sustains, comprises and yields N-world Special Theory effects. These NSPPM effects are electromagnetic in character. This medium through which the effects appear to propagate comprise the objects that yield these effects. The next assumption is convincingly obtained from a simple and literal translation of the concept of infinitesimal reasoning. (ii) Any N-world position from or through which an electromagnetic effect appears to propagate, when viewed from the NSP-world, is embedded into a disjoint “monadic cluster” of the NSMP, where this monadic cluster mirrors the same unusual order properties, with respect to propagation, as the nonstandard ordering of the nonarchimedian field of hyperreal numbers $`{}_{}{}^{}\mathrm{𝙸𝚁}.`$ A monadic cluster may be a set of NS-substratum subparticles located within a monad of the standard N-world position. The propagation properties within each such monad are identical. In what follows, consider two (local) fundamental pairs of N-world positions $`F_1,F_2`$ that are in nonzero uniform (constant) NSP-world linear and relative motion. Our interest is in what effect such nonzero velocity might have upon such electromagnetic propagation. Within the NSP-world, this uniform and linear motion is measured by the number $`w`$ that is near to a standard number $`\omega `$ and this velocity is measured with respect to conceptional NSP-world time and a stationary subparticle field. \[Note that field expansion can be additionally incorporated.\] The same NSP-world linear ruler is used in both the NSP-world and the N-world. The only difference is that the ruler is restricted to the N-world when such measurements are made. N-world time is measured by only one type of machine – the light-clock. The concept of the light-clock is to be considered as any clock-like apparatus that utilizes either directly or indirectly an equivalent process. As it will be detailed, due to the different propagation effects of electromagnetic radiation within the two “worlds,” measured N-world light-clock time need not be the same as the NSP-world time. Further, the NSP-world ruler is the measure used to define the N-world light-clock. Experiments show that for small time intervals $`[a,b]`$ the Galilean theory of average velocities (velocitys) suffices to give accurate information relative to the compositions of such velocities. Let there be an internal function $`q:^{}[a,b]{}_{}{}^{}\mathrm{𝙸𝚁},`$ where $`q`$ represents in the NSP-world a distance function. Also, let nonnegative and internal $`\mathrm{}:^{}[a,b]{}_{}{}^{}\mathrm{𝙸𝚁}`$ be a function that yields the NSP-world velocity of the electromagnetic propagation at any $`t^{}[a,b].`$ As usual $`\mu (t)`$ denotes the monad of standard $`t,`$ where “$`t`$” is an absolute NSP-world “time” parameter. The general and correct methods of infinitesimal modeling state that, within the internal portion of the NSP-world, two measures $`m_1`$ and $`m_2`$ are indistinguishable for dt (i.e. infinitely close of order one) (notation $`m_1m_2`$) if and only if $`0dt\mu (0),`$ ($`\mu (0)`$ the set of infinitesimals) $$\frac{m_1}{dt}\frac{m_2}{dt}\mu (0).$$ $`(3.1)`$ Intuitively, indistinguishable in this sense means that, although within the NSP-world the two measures are only equivalent and not necessarily equal, the first level (or first-order) effects these measures represent over $`dt`$ are indistinguishable within the N-world (i.e. they appear to be equal.) In the following discussion, we continue to use photon terminology. Within the N-world our photons need not be conceived of as particles in the sense that there is a nonzero finite N-world distance between individual photons. Our photons may be finite combinations of intermediate subparticles that exhibit, when the standard part operator is applied, basic electromagnetic field properties. They need not be discrete objects when viewed from the N-world, but rather they could just as well give the appearance of a dense NS-substratum. Of course, this dense NSPPM portion is not the usual notion of an “æther” (i.e. ether) for it is not a subset of the N-world. This dense-like portion of the NS-substratum containsnonstandard particle medium (NSPPM). Again “photon” can be considered as but a convenient term used to discuss electromagnetic propagation. Now for another of our simplistic physical assumptions. (iii) In an N-world convex space neighborhood $`I`$ traced out over the time interval $`[a,b],`$ the NSPPM disturbances appear to propagate linearly. As we proceed through this derivation, other such assumptions will be identified. The functions $`q,\mathrm{}`$ need to satisfy some simple mathematical characteristic. The best known within nonstandard analysis is the concept of S-continuity . So, where defined, let $`q(x)/x`$ (a velocity type expression) and $`\mathrm{}`$ be S-continuous, and $`\mathrm{}`$ limited (i.e. finite) at each $`p[a,b],(a+\mathrm{at}a,b\mathrm{at}b).`$ From compactness, $`q(x)/x`$ and $`\mathrm{}`$ are S-continuous, and $`\mathrm{}`$ is limited on $`{}_{}{}^{}[a,b].`$ Obviously, both $`q`$ and $`\mathrm{}`$ may have infinitely many totally different NSP-world characteristics of which we could have no knowledge. But the function $`q`$ represents within the NSP-world the distance traveled with linear units by an identifiable NSPPM disturbance. It follows from all of this that for each $`t[a,b]`$ and $`t^{}\mu (t)^{}[a,b],`$ $$\frac{q(t^{})}{t^{}}\frac{q(t)}{t}\mu (0);\mathrm{}(t^{})\mathrm{}(t)\mu (0).$$ $`(3.2)`$ Expressions (3.2) give relations between nonstandard $`t^{}\mu (t)`$ and the standard $`t.`$ Recall that if $`x,y{}_{}{}^{}\mathrm{𝙸𝚁},`$ then $`xy`$ iff $`xy\mu (0).`$ From (3.2), it follows that for each $`dt\mu (0)`$ such that $`t+dt\mu (t)^{}[a,b]`$ $$\frac{q(t+dt)}{t+dt}\frac{q(t)}{t},$$ $`(3.3)`$ $$\mathrm{}(t+dt)+\frac{q(t+dt)}{t+dt}\mathrm{}(t)+\frac{q(t)}{t}.$$ $`(3.4)`$ One important observation is necessary. The fact that the function $`\mathrm{}`$ has been evaluated at $`t+dt`$ is not necessary for (3.4) to hold for it will also hold for any $`t^{}\mu (t)`$ and $`\mathrm{}(t^{})`$ substituted for $`\mathrm{}(t+dt).`$ But since we are free to choice any value $`t^{}\mu (t),`$ selecting particular values will allow our derivation to proceed to an appropriate N-world conclusion. From (3.4), we have that $$\left(\mathrm{}(t+dt)+\frac{q(t+dt)}{t+dt}\right)dt\left(\mathrm{}(t)+\frac{q(t)}{t}\right)dt.$$ $`(3.5)`$ It is now that we begin our application of the concepts of classical Galilean composition of velocities but restrict these ideas to the NSP-world monadic clusters and the notion of indistinguishable effects. You will notice that within the NSP-world the transfer of the classical concept of equality of constant or average quantities is replaced by the idea of indistinguishable. At the moment $`t[a,b]`$ that the standard part operator is applied, an effect is transmitted through the NSPPM as follows: (iv) For each $`dt\mu (0)`$ and $`t[a,b]`$ such that $`t+dt^{}[a,b],`$ the NSP-world distance $`q(t+dt)q(t)`$ (relative to $`dt`$) traveled by the NSPPM effect within a monadic cluster is indistinguishable for $`dt`$ from the distance produced by the Galilean composition of velocities. From (iv), it follows that $$q(t+dt)q(t)\left(\mathrm{}(t+dt)+\frac{q(t+dt)}{t+dt}\right)dt.$$ $`(3.6)`$ And from (3.5), $$q(t+dt)q(t)\left(\mathrm{}(t)+\frac{q(t)}{t}\right)dt.$$ $`(3.7)`$ Expression (3.7) is the basic result that will lead to conclusions relative to the Special Theory of Relativity. In order to find out exactly what standard functions will satisfy (3.7), let arbitrary $`t_1[a,b]`$ be the standard time at which electromagnetic propagation begins from position $`F_1.`$ Next, let $`q=^{}s`$ be an extended standard function and $`s`$ is continuously differentiable on $`[a,b].`$ Applying the definition of $`,`$ yields $$\frac{{}_{}{}^{}s(t+dt)s(t)}{dt}\mathrm{}(t)+\frac{s(t)}{t}.$$ $`(3.8)`$ Note that $`\mathrm{}`$ is microcontinuous on $`{}_{}{}^{}[a,b].`$ For each $`t[a,b],`$ the value $`\mathrm{}(t)`$ is limited. Hence, let $`\mathrm{𝚜𝚝}(\mathrm{}(t))=v(t)\mathrm{𝙸𝚁}.`$ From Theorem 1.1 in or 7.6 in , $`v`$ is continuous on $`[a,b].`$ \[See note 1 part a.\] Now (3.8) may be rewritten as $${}_{}{}^{}\left(\frac{d(s(t)/t)}{dt}\right)=\frac{{}_{}{}^{}v(t)}{t},$$ $`(3.9)`$ where all functions in (3.9) are \-continuous on $`{}_{}{}^{}[a,b].`$ Consequently, we may apply the \-integral to both sides of (3.9). \[See note 1 part b.\] Now (3.9) implies that for $`t[a,b]`$ $$\frac{s(t)}{t}=^{}_{t_1}^t\frac{{}_{}{}^{}v(x)}{x}𝑑x,$$ $`(3.10)`$ where, for $`t_1[a,b],`$ $`s(t_1)`$ has been initialized to be zero. Expression (3.10) is of interest in that it shows that although (iv) is a simplistic requirement for monadic clusters and the requirement that $`q(x)/x`$ be S-continuous is a customary property, they do not lead to a simplistic NSP-world function, even when view at standard NSP-world times. It also shows that the light-clock assumption was necessary in that the time represented by (3.10) is related to the distance traveled and unknown velocity of an identifiable NSPPM disturbance. It is also obvious that for pure NSP-world times the actual path of motion of such propagation effects is highly nonlinear in character, although within a monadic cluster the distance $`{}_{}{}^{}s(t+dt)s(t)`$ is indistinguishable from that produced by the linear-like Galilean composition of velocities. Further, it is the standard function in (3.10) that allows us to cross over to other monadic clusters. Thus, substituting into (3.7) yields, since the propagation behavior in all monadic clusters is identical, $`{}_{}{}^{}s(t+dt)s(t)\left({}_{}{}^{}v(t)+\left({}_{}{}^{}_{t_1}^{t}{}_{}{}^{}v(x)/x𝑑x\right)\right)dt,`$ (3.11)for every $`t[a,b],t+dt\mu (t)^{}[a,b]`$ Consider a second standard position $`F_2`$ at which electromagnetic reflection occurs at $`t_2[a,b],t_2>t_1,t_2+dt\mu (t_2)^{}[a,b].`$ Then (3.11) becomes $`{}_{}{}^{}s(t_2+dt)s(t_2)\left({}_{}{}^{}v(t_2)+\left({}_{}{}^{}_{t_1}^{t_2}{}_{}{}^{}v(x)/x𝑑x\right)\right)dt.`$ (3.12) Our final assumption for monadic cluster behavior is that the classical ballistic property holds with respect to electromagnetic propagation. (v) From the exterior NSP-world viewpoint, at standard time $`t[a,b],`$ the velocity $`{}_{}{}^{}v(t)`$ acquires an additional velocity $`w`$. Applying the classical statement (v), with the indistinguishable concept, means that the distance traveled $`{}_{}{}^{}s(t_2+dt)s(t_2)`$ is indistinguishable from $`(^{}v(t_2)+w)dt.`$ Hence, $$(^{}v(t_2)+w)dt^{}s(t_2+dt)s(t_2)({}_{}{}^{}v(t_2)+\left({}_{}{}^{}_{t_1}^{t_2}\frac{{}_{}{}^{}v(x)}{x}dx\right))dt.$$ $`(3.13)`$ Expression (3.13) implies that $${}_{}{}^{}v(t_2)+w^{}v(t_2)+\left({}_{}{}^{}_{t_1}^{t_2}\frac{{}_{}{}^{}v(x)}{x}𝑑x\right).$$ $`(3.14)`$ Since $`\mathrm{𝚜𝚝}(w)`$ is a standard number, (3.14) becomes after taking the standard part operator, $$\mathrm{𝚜𝚝}(w)=\mathrm{𝚜𝚝}\left({}_{}{}^{}_{t_1}^{t_2}\frac{{}_{}{}^{}v(x)}{x}𝑑x\right).$$ $`(3.15)`$ After reflection, a NSPPM disturbance returns to the first position $`F_1`$ arriving at $`t_3[a,b],t_1<t_2<t_3.`$ Notice that the function $`s`$ does not appear in equation (3.15). Using the nonfavored position concept, a reciprocal argument entails that $$\frac{s_1(t_3)}{t_3}=\mathrm{𝚜𝚝}\left({}_{}{}^{}_{t_2}^{t_3}\frac{{}_{}{}^{}v_{1}^{}(x)}{x}𝑑x\right),$$ $`(3.16)`$ $$\mathrm{𝚜𝚝}(w)=\mathrm{𝚜𝚝}\left({}_{}{}^{}_{t_2}^{t_3}\frac{{}_{}{}^{}v_{1}^{}(x)}{x}𝑑x\right),$$ $`(3.17)`$ where $`s_1(t_2)`$ is initialized to be zero. It is not assumed that $`{}_{}{}^{}v_{1}^{}=^{}v.`$ We now combine (3.10), (3.15), (3.16), (3.17) and obtain an interesting nonmonadic view of the relationship between distance traveled by an NSPPM disturbance and relative velocity. $$s_1(t_3)s(t_2)=\mathrm{𝚜𝚝}(w)(t_3t_2).$$ $`(3.18)`$ Although reflection has been used to determine relation (3.18) and a linear-like interpretation involving reflection seems difficult to express, there is a simple nonreflection analogue model for this behavior. Suppose that a NSPPM disturbance is transmitted from a position $`F_1,`$ to a position $`F_2.`$ Let $`F_1`$ and $`F_2`$ have no NSP-world relative motion. Suppose that a NSPPM disturbance is transmitted from $`F_1`$ to $`F_2`$ with a constant velocity $`v`$ with the duration of the transmission $`t^{\prime \prime }t^{},`$ where the path of motion is considered as linear. The disturbance continues linearly after it passes point $`F_2`$ but has increased during its travel through the monadic cluster at $`F_2`$ to the velocity $`v+\mathrm{𝚜𝚝}(w).`$ The disturbance then travels linearly for the same duration $`t^{\prime \prime }t^{}.`$ The linear difference in the two distances traveled is $`w(t^{\prime \prime }t^{}).`$ Such results in the NSP-world should be construed only as behavior mimicked by the analogue NSPPM model. Equations (3.10) and (3.15) show that in the NSP-world NSPPM disturbances propagate. Except for the effects of material objects, it is assumed that in the N-world the path of motion displayed by a NSPPM disturbance is linear. This includes the path of motion within an N-world light-clock. We continue this derivation based upon what, at present, appears to be additional parameters, a private NSP-world time and an NSP-world rule. Of course, the idea of the N-world light-clock is being used as a fixed means of identifying the different effects the NSPPM is having upon these two distinct worlds. A question yet to be answered is how can we compensate for differences in these two time measurements, the NSP-world private time measurement of which we can have no knowledge and N-world light-clocks. The weighted mean value theorem for integrals in nonstandard form, when applied to equations (3.15) and (3.17), states that there are two NSP-world times $`t_a,t_b^{}[a,b]`$ such that $`t_1t_at_2t_bt_3`$ and $$\mathrm{𝚜𝚝}(w)=\mathrm{𝚜𝚝}(^{}v(t_a))_{t_1}^{t_2}\frac{1}{x}dx=\mathrm{𝚜𝚝}(^{}v_1(t_b))_{t_2}^{t_3}\frac{1}{x}dx.$$ $`(3.19)`$ \[See note 1 part c.\] Now suppose that within the local N-world an $`F_1F_2,`$ $`F_2F_1`$ light-clock styled measurement for the velocity of light using a fixed instrumentation yields equal quantities. (Why this is the case is established in Section 6.) Model this by (\) $`\mathrm{𝚜𝚝}(^{}v(t_a))=\mathrm{𝚜𝚝}(^{}v_1(t_b))=c`$ in the NSPPM. I point out that there are many nonconstant \-continuous functions that satisfy property (\). For example, certain standard nonconstant linear functions and nonlinear modifications of them. Property (\) yields $$_{t_1}^{t_2}\frac{1}{x}𝑑x=_{t_2}^{t_3}\frac{1}{x}𝑑x.$$ $`(3.20)`$ And solving (3.20) yields $$\mathrm{ln}\left(\frac{t_2}{t_1}\right)=\mathrm{ln}\left(\frac{t_3}{t_2}\right).$$ $`(3.21)`$ From this one has $$t_2=\sqrt{t_1t_3}.$$ $`(3.22)`$ Expression (3.22) is Prokhovnik’s equation (6.3.3) in reference . However, the interpretation of this result and the others that follow cannot, for the NSP-world, be those as proposed by Prokhovnik. The times $`t_1,t_2,t_3,`$ are standard NSPPM times. Further, it is not logically acceptable when considering how to measure such time in the NSP-world or N-world to consider just any mode of measurement. The mode of light velocity measurement must be carried out within the confines of the language used to obtain this derivation. Using this language, a method for time calculation that is permissible in the N-world is the light-clock method. Any other described method for time calculation should not include significant terms from other sources. Time as expressed in this derivation is not a mystical absolute something or other. It is a measured quantity based entirely upon some mode of measurement. They are two major difficulties with most derivations for expressions used in the Special Theory. One is the above mentioned absolute time concept. The other is the ad hoc nonderived N-world relative velocity. In this case, no consideration is given as to how such a relative velocity is to be measured so that from both $`F_1`$ and $`F_2`$ the same result would be obtained. It is possible to achieve such a measurement method because of the logical existence of the NSPPM. In a physical-like sense, the “times” can be considered as the numerical values recorded by single device stationary in the NSPPM. It is conceptual time in that, when events occur, then such numerical event-times “exist.” It is the not yet identified NSPPM properties that yield the unusual behavior indicated by (3.22). One can use light-clocks and a counter that indicates, from some starting count, the number of times the light pulse has traversed back and forth between the mirror and source of our light-clock. Suppose that $`F_1`$ and $`F_2`$ can coincide. When they do coincide, the $`F_2`$ light-clock counter number that appears conceptually first after that moment can be considered to coincide with the counter number for the $`F_1`$ light-clock. After $`F_2`$ is perceived to no longer coincide with $`F_1,`$ a light pulse is transmitted from $`F_1`$ towards $`F_2`$ in an assumed linear manner. The “next” $`F_1`$ counter number after this event is $`\tau _{11}.`$ We assume that the relative velocity of $`F_2`$ with respect to $`F_1`$ may have altered the light-clock counter numbers, compared to the count at $`F_1,`$ for a light-clock riding with $`F_2`$. The length $`L`$ used to define a light-clock is measured by the NSP-world ruler and would not be altered. Maybe the light velocity $`c`$, as produced by the standard part operator, is altered by N-world relative velocity. Further, these two N-world light-clocks are only located at the two positions $`F_1,F_2,`$ and this light pulse is represented by a NSPPM disturbance. The light pulse is reflected back to $`F_1`$ by a mirror similar to the light-clock itself. The first counter number on the $`F_2`$ light-clock to appear, intuitively, “after” this reflection is approximated by $`\tau _{21}.`$ The $`F_1`$ counter number first perceived after the arrival of the returning light pulse is $`\tau _{31}.`$ From a linear viewpoint, at the moment of reflection, denoted by $`\tau _{21},`$ the pulse has traveled an operational linear light-clock distance of $`(\tau _{21}\tau _{11})L.`$ After reflection, under our assumptions and nonfavored position concept, a NSPPM disturbance would trace out the same operational linear light-clock distance measured by $`(\tau _{31}\tau _{21})L.`$ Thus the operational light-clock distance from $`F_1`$ to $`F_2`$ would be at the moment of operational reflection, under our linear assumptions, 1/2 the sum of these two distances or $`S_1=(1/2)(\tau _{31}\tau _{11})L.`$ Now we can also determine the appropriate operational relation between these light-clock counter numbers for $`S_1=(\tau _{21}\tau _{11})L.`$ Hence, $`\tau _{31}=2\tau _{21}\tau _{11},`$ and $`\tau _{21}`$ operationally behaves like an Einstein measure. After, measured by light-clock counts, the pulse has been received back to $`F_1,`$ a second light pulse (denoted by a second subscript of 2) is immediately sent to $`F_2.`$ Although $`\tau _{31}\tau _{12}`$, it is assumed that $`\tau _{31}=\tau _{12}`$ \[See note 2.5\]. The same analysis with new light-clock count numbers yields a different operational distance $`S_2=(1/2)(\tau _{32}\tau _{12})L`$ and $`\tau _{32}=2\tau _{22}\tau _{12}.`$ One can determine the operational light-clock time intervals by considering $`\tau _{22}\tau _{21}=(1/2)((\tau _{32}\tau _{31})+(\tau _{12}\tau _{11}))`$ and the operational linear light-clock distance difference $`S_2S_1=(1/2)((\tau _{32}\tau _{31})(\tau _{12}\tau _{11}))L.`$ Since we can only actually measure numerical quantities as discrete or terminating numbers, it would be empirically sound to write the N-world time intervals for these scenarios as $`t_1=\tau _{12}\tau _{11},t_3=(\tau _{32}\tau _{31}).`$ This yields the operational Einstein measure expressions in (6.3.4) of as $`\tau _{22}\tau _{21}=t_E`$ and operational light-length $`r_E=S_2S_1,`$ using our specific light-clock approach. This allows us to define, operationally, the N-world relative velocity as $`v_E=r_E/t_E.`$ \[In this section, the $`t_1,t_3`$ are not the same Einstein measures, in form, as described in . But, in section 4, 5, 6 these operational measures are used along with infinitesimal light-clock counts to obtain the exact Einstein measure forms for the time measure. This is: the $`t_1`$ is a specific starting count and the $`t_3`$ is $`t_1`$ plus an appropriate lapsed time.\] Can we theoretically turn the above approximate operational approach for discrete N-world light-clock time into a time continuum? Light-clocks can be considered from the NSP-world viewpoint. In such a case, the actual NSP-world length used to form the light-clock might be considered as a nonzero infinitesimal. Thus, at least, the numbers $`\tau _{32},\tau _{21},\tau _{31},\tau _{22}`$ are infinite hyperreal numbers, various differences would be finite and, after taking the standard part operator, all of the N-world times and lengths such as $`t_E,r_E,S_1,S_2`$ should be exact and not approximate in character. These concepts will be fully analyzed in section 6. Indeed, as previously indicated, for all of this to hold the velocity $`c`$ cannot be measured by any means. As indicated in section 6, the actual numerical quantity $`c`$ as it appears in (3.22) is the standard part of pure NSP-world quantities. Within the N-world, one obtains an “apparent” constancy for the velocity of light since, for this derivation, it must be measured by means of a to-and-fro light-clock styled procedure with a fixed instrumentation. As yet, we have not discussed relations between N-world light-clock measurements and N-world physical laws. It should be self-evident that the assumed linearity of the light paths in the N-world can be modeled by the concept of projective geometry. Relative to the paths of motion of a light path in the NSP-world, the NSPPM disturbances, the N-world path behaves as if it were a projection upon a plane. Prokhovnik analyzes such projective behavior and comes to the conclusions that in two or more dimensions the N-world light paths would follow the rules of hyperbolic geometry. In Prokhovnik, the equations (3.22) and the statements establishing the relations between the operational or exact Einstein measures $`t_E,r_E`$ and $`v_E`$ lead to the Einstein expression relating the light-clock determined relative velocities for three linear positions having three NSP-world relative and uniform velocities $`w_1,w_2,w_3.`$ In the appendix, in terms of light-clock determined Einstein measures and based upon the projection idea, the basic Special Theory coordinate transformation are correctly obtained. Thus, all of the NSP-world times have been removed from the results and even the propagation differences with respect to light-clock measurements. Just use light-clocks in the N-world to measure all these quantities in the required manner and the entire Special Theory is forthcoming. I mention that it can be shown that $`w`$ and $`c`$ may be measured by probes that are not N-world electromagnetic in character. Thus $`w`$ need not be obtained in the same manner as is $`v_E`$ except that N-world light-clocks would be used for N-world time measurements. For this reason, $`\mathrm{𝚜𝚝}(w)=\omega `$ is not directly related to the so-called textual expansion of the space within our universe. The NSPPM is not to be taken as a nonstandard translation of the Maxwell EMF equations. 4. The Time Continuum. With respect to models that use the classical continuum approach (i.e. variables are assumed to vary over such things as an interval of real numbers) does the mathematics perfectly measure quantities within nature – quantities that cannot be perfectly measured by a human being? Or is the mathematics only approximate in some sense? Many would believe that if “nature” is no better than the human being, then classical mathematics is incorrect as a perfect measure of natural system behavior. However, this is often contradicted in the limit. That is when individuals refine their measurements, as best as it can done at the present epoch, then the discrete human measurements seem to approach the classical as a limit. Continued exploration of this question is a philosophical problem that will not be discussed in this paper, but it is interesting to model those finite things that can, apparently, be accomplished by the human being, transfer these processes to the NSP-world and see what happens. For what follows, when the term “finite” (i.e. limited) hyperreal number is used, since it is usually near to a nonzero real number, it will usually refer to the ordinary nonstandard notion of finite except that the infinitesimals have been removed. This allows for the existence of finite multiplicative inverses. First, suppose that $`t_E=\mathrm{𝚜𝚝}(t_{Ea}),r_E=\mathrm{𝚜𝚝}(r_{Ea}),S_1=\mathrm{𝚜𝚝}(S_{1a}),S_2=\mathrm{𝚜𝚝}(S_{2a})`$ and each is a nonnegative real number. Thus $`t_{Ea},r_{Ea},S_{1a},S_{2a}`$ are all nonnegative finite hyperreal numbers. Let $`L=1/10^\omega >0,\omega \mathrm{𝙸𝙽}_{\mathrm{}}^+.`$ By transfer and the result that $`S_{1a},S_{2a},`$ are considered finite (i.e. near standard), then $`S_{1a}(1/2)L(\tau _{31}\tau _{11})L(\tau _{21}\tau _{11})(1/2)(\tau _{31}\tau _{11}),(\tau _{21}\tau _{11})`$ cannot be finite. Thus, by Theorem 11.1.1 , it can be assumed that there exist $`\eta ,\gamma \mathrm{𝙸𝙽}_{\mathrm{}}^+`$ such that $`(1/2)(\tau _{31}\tau _{11})=\eta ,(\tau _{21}\tau _{11})=\gamma .`$ This implies that each $`\tau `$ corresponds to an infinite light-clock count and that $$\tau _{31}=2\eta +\tau _{11},\tau _{21}=\gamma +\tau _{11}.$$ $`(4.1)`$ In like manner, it follows that $$\tau _{32}=2\lambda +\tau _{12},\tau _{22}=\delta +\tau _{12},\lambda ,\delta \mathrm{𝙸𝙽}_{\mathrm{}}^+.$$ $`(4.2)`$ Observe that the second of the double subscripts being 2 indicates the light-clock counts for the second light transmission. Now for $`t_{Ea}`$ to be finite requires that the corresponding nonnegative $`t_{1a},t_{3a}`$ be finite. Since a different mode of conceptual time might be used in the NSP-world, then there is a need for a number $`u=L/c`$ that adjusts NSP-world conceptual time to the light-clock count numbers. \[See note 18.\] By transfer of the case where these are real number counts, this yields that $`t_{3a}u(\tau _{32}\tau _{31})=2u(\lambda \eta )+u(\tau _{12}\tau _{11})2u(\lambda \eta )+t_{1a}`$ and $`t_{Ea}u(\tau _{22}\tau _{21})u(\delta \gamma )+t_{1a}.`$ Hence for all of this to hold in the NSP-world $`u(\delta \gamma )`$ must be finite or that there exists some $`r\mathrm{𝙸𝚁}^+`$ such that $`u(\delta \gamma )\mu (r).`$ Let $`\tau _{12}=\alpha ,\tau _{11}=\beta .`$ Then $`t_{Ea}u(\delta \gamma )+u(\alpha \beta )`$ implies that $`u(\alpha \beta )`$ is also finite. The requirement that these infinite numbers exist in such a manner that the standard part of their products with $`L`$ \[resp. $`u`$\] exists and satisfies the continuum requirements of classical mathematics is satisfied by Theorem 11.1.1 , where in that theorem $`10^\omega =1/L`$ \[resp. $`1/u`$\]. \[See note 2.\] It is obvious that the nonnegative numbers needed to satisfy this theorem are nonnegative infinite numbers since the results are to be nonnegative and finite. Theorem 11.1.1 allows for the appropriate $`\lambda ,\eta ,\delta ,\gamma `$ to satisfy a bounding property in that we know two such numbers exist such that $`\lambda ,\eta <1/L^2,\delta ,\gamma <1/u^2.`$ \[Note: It is important to realize that due to this correspondence to a continuum of real numbers that the entire analysis as it appears in section 3 is now consistent with a mode of measurement. Also the time concept is replaced in this analysis with a “count” concept. This count concept will be interpreted in section 8 as a count per some unit of time measure.\] Also note that the concepts are somewhat simplified if it is assumed that $`\tau _{12}=\tau _{31}.`$ In this case, substitution into 4.1 yields that $`t_{1a}2u\eta `$ and $`t_{3a}2u\lambda .`$ Consequently, $`t_{Ea}=(1/2)(t_{1a}+t_{3a})u(\lambda +\eta ).`$ This predicts what is to be expected, that, in this case, the value of $`t_E`$ from the NSP-world viewpoint is not related to the first “synchronizing” light pulse sent. 5. Standard Light-clocks and c. I mention that the use of subparticles or the concept of the NSPPM are not necessary for the derivation in section 3 to hold. One can substitute for the NSPPM the term “NS-substratum” or the like and for the term “monadic cluster” of possible subparticles just the concept of a “monadic neighborhood.” It is not necessary that one assume that the NS-substratum is composed of subparticles or any identifiable entity, only that NSPPM transmission of such radiation behaves in the simplistic manner stated. It is illustrative to show by a diagram of simple light-clock counts how this analysis actually demonstrates the two different modes of propagation, the NSP-world mode and the different mode when viewed from the N-world. In general, $`L`$ is always fixed and for the following analysis and, for this particular scenario, inf. light-clock $`c`$ may change. This process of using N-world light-clocks to approximate the relative velocity should only be done once due to the necessity of “indexing” the light-clocks when $`F_1`$ and $`F_2`$ coincide. In the following diagram, the numbers represent actual light-clock count numbers as perceived in the N-world. The first column are those recorded at $`F_1,`$ the second column those required at $`F_2.`$ The arrows and the numbers above them represent our $`F_1`$ comprehension of what happens when the transmission is considered to take place in the N-world. The Einstein measures are only for the $`F_1`$ position. $$\begin{array}{ccccc}F_1& & \mathrm{N}\mathrm{world}& & F_2\\ \tau _{11}=20& & & & & \\ & & \stackrel{20}{}& & \\ \tau =40& & & & \tau _{21}=40\\ & & \stackrel{20}{}& & \\ \tau _{31}=60& & & & & \\ & & & & & & \\ \tau _{12}=80& & & & & \\ & & \stackrel{30}{}& & \\ & & & & \tau _{22}=110\\ & & \stackrel{30}{}& & \\ \tau _{32}=140& & & & & \end{array}$$ Certainly, the above diagram satisfies the required light-clock count equations. The only light-clock counts that actually are perceivable are those at $`F_1.`$ And, for the transformation equations, the scenario is altered. When the Special Theory transformation equations are obtained, two distinct N-world observers are used and a third N-world distinct fundamental position. All light-clock counts made at each of these three positions are entered into the appropriate expressions for the Einstein measures as obtained for each individual position. 6. Infinitesimal Light-clock Analysis. In the originally presented Einstein derivation, time and length are taken as absolute time and length. It was previously pointed out that this assumpt yields logical error. The scientific community extrapolated the language used in the derivation, a language stated only in terms of light propagation behavior, without logical reason, to the “concept” of Newtonian absolute time and length. Can the actual meaning of the “time” and “length” expressed in the Lorentz transformation be determined? In what follows, a measure by light-clock counts is used to analyze the classical transformation as derived in the Appendix-A and, essentially, such “counts” will replace conceptional time. \[See note 1.5\] The superscripts indicate the counts associated with the light-clocks, the Einstein measures, and the like, at the positions $`F_1,F_2.`$ The $`1`$ being the light-clock measures at $`F_1`$ for a light pulse event from $`P,`$ the $`2`$ for the light-clock measures at the $`F_2`$ for the same light pulse event from $`P,`$ and the 3 for the light-clock measures and its corresponding Einstein measures at $`F_1`$ for the velocity of $`F_2`$ relative to $`F_1.`$ The NSP-world measured angle, assuming linear projection due to the constancy of the velocities, from $`F_1`$ to the light pulse event from $`P`$ is $`\theta ,`$ and that from $`F_2`$ to $`P`$ is an exterior angle $`\varphi .`$ The expressions for our proposes are $`x_E^{(1)}=v_E^{(1)}t_E^{(1)}\mathrm{cos}\theta ,x_E^{(2)}=v_E^{(2)}t_E^{(2)}\mathrm{cos}\varphi .`$ \[Note: The negative is required since $`\pi /2\varphi \pi `$ and use of the customary coordinate systems.\] In all that follows, $`i`$ varies from 1 to 3. We investigate what happens when the standard model is now embedded back again into the non-infinitesimal finite NSP-world. All of the “coordinate” transformation equations are in the Appendix and they actually only involve $`\omega _i/c.`$ These equations are interpreted in the NSP-world. But as far as the light-clock counts are concerned, their appropriate differences are only infinitely near to a standard number. The appropriate expressions are altered to take this into account. For simplicity in notation, it is again assumed that “immediate” in the light-clock count process means $`\tau _{12}^{(i)}=\tau _{31}^{(i)}.`$ \[See note 3.\] Consequently, $`t_{1a}^{(i)}2u\eta ^{(i)},t_{3a}^{(i)}2u\lambda ^{(i)},\eta ^{(i)},\lambda ^{(i)}\mathrm{𝙸𝙽}_{\mathrm{}}^+.`$ Then $$t_{Ea}^{(i)}u(\lambda ^{(i)}+\eta ^{(i)}),\lambda ^{(i)},\eta ^{(i)}\mathrm{𝙸𝙽}_{\mathrm{}}^+.$$ $`(6.1)`$ Now from our definition $`r_E^{(i)}L(\lambda ^{(i)}\eta ^{(i)}),(\lambda ^{(i)}\eta ^{(i)})\mathrm{𝙸𝙽}_{\mathrm{}}^+.`$ Hence, since all of the numbers to which st is applied are nonnegative and finite and $`\mathrm{𝚜𝚝}(v_{Ea}^{(i)})\mathrm{𝚜𝚝}(t_{Ea}^{(i)})=\mathrm{𝚜𝚝}(r_{Ea}^{(i)}),`$ it follows that $$v_{Ea}^{(i)}L\frac{(\lambda ^{(i)}\eta ^{(i)})}{u(\lambda ^{(i)}+\eta ^{(i)})}.$$ $`(6.2)`$ Now consider a set of two 4-tuples $$(\mathrm{𝚜𝚝}(x_{Ea}^{(1)}),\mathrm{𝚜𝚝}(y_{Ea}^{(1)}),\mathrm{𝚜𝚝}(z_{Ea}^{(1)}),\mathrm{𝚜𝚝}(t_{Ea}^{(1})),$$ $$(\mathrm{𝚜𝚝}(x_{Ea}^{(2)}),\mathrm{𝚜𝚝}(y_{Ea}^{(2)}),\mathrm{𝚜𝚝}(z_{Ea}^{(2)}),\mathrm{𝚜𝚝}(t_{Ea}^{(2)})),$$ where they are viewed as Cartesian coordinates in the NSP-world. First, we have $`\mathrm{𝚜𝚝}(x_{Ea}^{(1)})=\mathrm{𝚜𝚝}(v_{Ea}^{(1)})\mathrm{𝚜𝚝}((t_{Ea}^{(1)})\mathrm{𝚜𝚝}(^{}\mathrm{cos}\theta ),\mathrm{𝚜𝚝}(x_{Ea}^{(2)})=\mathrm{𝚜𝚝}(v_{Ea}^{(2)})\mathrm{𝚜𝚝}(t_{Ea}^{(2)})\mathrm{𝚜𝚝}(^{}\mathrm{cos}\varphi ).`$ Now suppose the local constancy of $`c`$. The N-world Lorentz transformation expressions are $$\mathrm{𝚜𝚝}(t_{Ea}^{(1)})=\beta _3(\mathrm{𝚜𝚝}(t_{Ea}^{(2)})+\mathrm{𝚜𝚝}(v_{Ea}^{(3)})\mathrm{𝚜𝚝}(x_{Ea}^{(2)})/c^2),$$ $$\mathrm{𝚜𝚝}(x_{Ea}^{(1)})=\beta _3(\mathrm{𝚜𝚝}(x_{Ea}^{(2)})+\mathrm{𝚜𝚝}(v_{Ea}^{(3)})\mathrm{𝚜𝚝}(t_{Ea}^{(2)})),$$ where $`\beta _3=\mathrm{𝚜𝚝}((1(v_{Ea}^{(3)})^2/c^2)^{1/2}).`$ Since $`L(\lambda ^{(i)}\eta ^{(i)})cu(\lambda ^{(i)}\eta ^{(i)})`$, the finite character of $`L(\lambda ^{(i)}\eta ^{(i)}),u(\lambda ^{(i)}\eta ^{(i)})`$ yields that $`c=\mathrm{𝚜𝚝}(L/u)`$ \[See note 8\]. When transferred to the NSP-world with light-clock counts, substitution yields $$t_{Ea}^{(1)}u(\lambda ^{(1)}+\eta ^{(1)})\beta [u(\lambda ^{(2)}+\eta ^{(2)})u(\lambda ^{(2)}+\eta ^{(2)})K^{(3)}K^{(2)}{}_{}{}^{}\mathrm{cos}\varphi ],$$ $`(6.3)`$ where $`K^{(i)}=(\lambda ^{(i)}\eta ^{(i)})/(\lambda ^{(i)}+\eta ^{(i)}),\beta =(1(K^{(3)})^2)^{1/2}.`$ For the “distance” transformation, we have $$x_{Ea}^{(1)}L(\lambda ^{(1)}\eta ^{(1)})^{}\mathrm{cos}\theta $$ $$\beta (L(\lambda ^{(2)}\eta ^{(2)})^{}\mathrm{cos}\varphi +\frac{L(\lambda ^{(3)}\eta ^{(3)})}{u(\lambda ^{(3)}+\eta ^{(3)})}u(\lambda ^{(2)}+\eta ^{(2)})).$$ $`(6.4)`$ Assume in the NSP-world that $`\theta \pi /2,\varphi \pi .`$ Consequently, substituting into 6.4 yields $$L(\lambda ^{(2)}\eta ^{(2)})\frac{L(\lambda ^{(3)}\eta ^{(3)})}{u(\lambda ^{(3)}+\eta ^{(3)})}u(\lambda ^{(2)}+\eta ^{(2)}).$$ $`(6.5)`$ Applying the finite property for these numbers, and, for this scenario, taking into account the different modes of the corresponding light-clock measures, yields $$\frac{L(\lambda ^{(3)}\eta ^{(3)})}{u(\lambda ^{(3)}+\eta ^{(3)})}\frac{L(\eta ^{(2)}\lambda ^{(2)})}{u(\lambda ^{(2)}+\eta ^{(2)})}v_{Ea}^{(3)}v_{Ea}^{(2)}.$$ $`(6.6)`$ Hence, $`\mathrm{𝚜𝚝}(v_{Ea}^{(3)})=\mathrm{𝚜𝚝}(v_{Ea}^{(2)}).`$ \[Due to the coordinate-system selected, these are directed velocities.\] This predicts that, in the N-world, the light-clock determined relative velocity of $`F_2`$ as measured from the $`F_1`$ and $`F_1`$ as measured from the $`F_2`$ positions would be the same if these special infinitesimal light-clocks are used. If noninfinitesimal N-world light-clocks are used, then the values will be approximately the same and equal in the limit. Expression 6.4 relates the light-clock counts relative to the measure of the to-and-fro paths of light transmission. By not substituting for $`x_{Ea}^{(2)},`$ it is easily seen that $`x_{Ea}^{(2)}LG,`$ where $`G`$ is an expression written entirely in terms of various light-clock count numbers. This implies that the so-called 4-tuples $`(\mathrm{𝚜𝚝}(x_{Ea}^{(1)}),\mathrm{𝚜𝚝}(y_{Ea}^{(1)}),`$ $`\mathrm{𝚜𝚝}(z_{Ea}^{(1)}),`$ $`\mathrm{𝚜𝚝}(t_{Ea}^{(1})),`$ $`(\mathrm{𝚜𝚝}(x_{Ea}^{(2)}),`$$`\mathrm{𝚜𝚝}(y_{Ea}^{(2)}),`$ $`\mathrm{𝚜𝚝}(z_{Ea}^{(2)}),\mathrm{𝚜𝚝}(t_{Ea}^{(2)}))`$ are not the absolute Cartesian type coordinates determined by Euclidean geometry and used to model Galilean dynamics. These coordinates are dynamically determined by the behavior of electromagnetic radiation within the N-world. Indeed, in , the analysis within the (outside of the monadic clusters) that leads to Prokhovnik’s conclusions is only relative to electromagnetic propagation and is done by pure number Galilean dynamics. Recall that the monadic cluster analysis is also done by Galilean dynamics. In general, when it is claimed that “length contracts” with respect to relative velocity the “proof” is stated as follows: $`x^{}=\mathrm{𝚜𝚝}(\beta )(x+vt);\overline{x}^{}=\mathrm{𝚜𝚝}(\beta )(\overline{x}+\overline{v}\overline{t}).`$ Then these two expressions are subtracted. Supposedly, this yields $`\overline{x}^{}x^{}=\mathrm{𝚜𝚝}(\beta )(\overline{x}x)`$ since its assumed that $`\overline{v}\overline{t}=vt.`$ For defined coordinates $`\overline{x}_E^{(i)},x_E^{(i)},i=1,2,`$ a more complete expression would be $$\overline{x}_E^{(1)}x_E^{(1)}=\mathrm{𝚜𝚝}(\beta )((\overline{x}_E^{(2)}x_E^{(2)})+(\overline{v}_E^{(3)}\overline{t}_E^{(2)}v_E^{(3)}t_E^{(2)})).$$ $`(6.7)`$ In this particular analysis, it has been assumed that all NSP-world relative velocities $`\omega _i,\overline{\omega _i}0.`$ To obtain the classical length contraction expression, let $`\omega _i=\overline{\omega _i},i=1,2,3.`$ Now this implies that $`\overline{\theta }=\theta ,\overline{\varphi }=\varphi `$ as they appear in the velocity figure on page 52 and that $$\overline{x}_E^{(1)}x_E^{(1)}=\mathrm{𝚜𝚝}(\beta )(\overline{x}_E^{(2)}x_E^{(2)}).$$ $`(6.8)`$ The difficulty with this expression has been its interpretation. Many modern treatments of Special Relativity argue that (6.8) has no physical meaning. But in these arguments it is assumed that $`\overline{x}_E^{(1)}x_E^{(1)}`$ means “length” in the Cartesian coordinate sense as related to Galilean dynamics. As pointed out, such a physical meaning is not the case. Expression (6.8) is a relationship between light-clock counts and, in general, displays properties of electromagnetic propagation within the N-world. Is there a difference between the right and left-hand sides of 6.8 when viewed entirely from the NSP-world. First, express 6.8 as $`\overline{x}_E^{(1)}x_E^{(1)}=\mathrm{𝚜𝚝}(\beta )\overline{x}_E^{(2)}\mathrm{𝚜𝚝}(\beta )x_E^{(2)}.`$ In terms of operational light-clock counts, this expression becomes $$L(\overline{\lambda }^{(1)}{}_{}{}^{}\mathrm{cos}\theta \overline{\eta }^{(1)}{}_{}{}^{}\mathrm{cos}\theta )L(\lambda ^{(1)}{}_{}{}^{}\mathrm{cos}\theta \eta ^{(1)}{}_{}{}^{}\mathrm{cos}\theta )$$ $`(6.9)`$ $$L(\overline{\lambda }^{(2)}\beta \|^{}\mathrm{cos}\varphi \|\overline{\eta }^{(2)}\beta \|^{}\mathrm{cos}\varphi \|)L(\lambda ^{(2)}\beta \|^{}\mathrm{cos}\varphi \|\eta ^{(2)}\beta \|^{}\mathrm{cos}\varphi \|),$$ where finite $`\beta =(1(K^{(3)})^2)^{1/2}`$ and $`\|\|`$ is used so that the Einstein velocities are not directed numbers and the Einstein distances are comparable. Also as long as $`\theta ,\varphi `$ satisfy the velocity figure on page 45, then (6.9) is independent of the specific angles chosen in the N-world since in the N-world expression (6.8) no angles appear relating the relative velocities. That is, the velocities are not vector quantities in the N-world, but scalars. Assuming the nontrivial case that $`\theta \pi /2,\varphi \pi ,`$ we have from Theorem 11.1.1 that there exist $`\overline{\mathrm{\Lambda }}^{(i)},\overline{N}^{(i)},\mathrm{\Lambda }^{(i)},N^{(i)}\mathrm{𝙸𝙽}_{\mathrm{}},i=1,2`$ such that $`{}_{}{}^{}\mathrm{cos}\theta \overline{\mathrm{\Lambda }}^{(1)}/\overline{\lambda }^{(1)}\overline{N}^{(1)}/\overline{\eta }^{(1)}`$ $`\mathrm{\Lambda }^{(1)}/\lambda ^{(1)}N^{(1)}/\eta ^{(1)},`$ $`\beta \|^{}\mathrm{cos}\varphi \|\overline{\mathrm{\Lambda }}^{(2)}/\overline{\lambda }^{(2)}\overline{N}^{(2)}/\overline{\eta }^{(2)}`$ $`\mathrm{\Lambda }^{(2)}/\lambda ^{(2)}N^{(1)}/\eta ^{(2)}.`$ Consequently, using the finite character of these quotients and the finite character of $`L(\overline{\lambda }^{(i)}),L(\overline{\eta }^{(i)}),L(\lambda ^{(i)}),L(\eta ^{(i)}),i=1,2,`$ the general three body NSP-world view 6.9 is $$L(\overline{\mathrm{\Lambda }}^{(1)}\overline{N}^{(1)})L(\mathrm{\Lambda }^{(1)}N^{(1)})=L\mathrm{\Gamma }^{(1)}$$ $$L\mathrm{\Gamma }_1^{(2)}=L(\overline{\mathrm{\Lambda }}^{(2)}\overline{N}^{(2)})L(\mathrm{\Lambda }^{(2)}N^{(2)}).$$ $`(6.10)`$ The obvious interpretation of 6.10 from the simple NSP-world light propagation viewpoint is displayed by taking the standard part of expression 6.10. $$\mathrm{𝚜𝚝}(L(\overline{\mathrm{\Lambda }}^{(1)}\overline{N}^{(1)}))\mathrm{𝚜𝚝}(L(\mathrm{\Lambda }^{(1)}N^{(1)}))=\mathrm{𝚜𝚝}(L\mathrm{\Gamma }^{(1)})=$$ $$\mathrm{𝚜𝚝}(L\mathrm{\Gamma }_1^{(2)})=\mathrm{𝚜𝚝}(L(\overline{\mathrm{\Lambda }}^{(1)}\overline{N}^{(1)}))\mathrm{𝚜𝚝}(L(\mathrm{\Lambda }^{(1)}N^{(1)})).$$ $`(6.11)`$ This is the general view as to the equality of the standard NSP-world distance traveled by a light pulse moving to-and-fro within a light-clock as used to measure at $`F_1`$ and $`F_2,`$ as viewed from the NSPPM only, the occurrence of the light pulse event from $`P`$. In order to interpret 6.9 for the N-world and a single NSP-world relative velocity, you consider additionally that $`\omega _1=\omega _2=\omega _3.`$ Hence, $`\theta =\pi /3`$ and correspondingly $`\varphi =2\pi /3.`$ In this case, $`\beta `$ is unaltered and since $`\mathrm{cos}\pi /3,\mathrm{cos}2\pi /3`$ are nonzero and finite, 6.9 now yields $$\mathrm{𝚜𝚝}(L(\overline{\lambda }^{(1)}\overline{\eta }^{(1)}))\mathrm{𝚜𝚝}(L(\lambda ^{(1)}\eta ^{(1)}))=$$ $$\mathrm{𝚜𝚝}(\beta )(\mathrm{𝚜𝚝}(L(\overline{\lambda }_1^{(2)}\overline{\eta }_1^{(2)}))\mathrm{𝚜𝚝}(L(\lambda _1^{(2)}\eta _1^{(2)})))$$ $$(\mathrm{𝚜𝚝}(L\overline{\lambda }^{(1)})\mathrm{𝚜𝚝}(L\overline{\eta }^{(1)}))(\mathrm{𝚜𝚝}(L\lambda ^{(1)})\mathrm{𝚜𝚝}(L\eta ^{(1)}))=$$ $$\mathrm{𝚜𝚝}(\beta )((\mathrm{𝚜𝚝}(L\overline{\lambda }_1^{(2)})\mathrm{𝚜𝚝}(L\overline{\eta }_1^{(2)}))(\mathrm{𝚜𝚝}(L\lambda _1^{(2)})\mathrm{𝚜𝚝}(L\eta _1^{(2)}))).$$ $`(6.12)`$ Or $$\mathrm{𝚜𝚝}(L(\overline{\lambda }^{(1)}\overline{\eta }^{(1)})L(\lambda ^{(1)}\eta ^{(1)}))=$$ $$\mathrm{𝚜𝚝}(L[(\overline{\lambda }^{(1)}\overline{\eta }^{(1)})(\lambda ^{(1)}\eta ^{(1)})])=$$ $$\mathrm{𝚜𝚝}(L\mathrm{\Pi }^{(1)})=\mathrm{𝚜𝚝}(\beta )\mathrm{𝚜𝚝}(L\mathrm{\Pi }_1^{(2)})=\mathrm{𝚜𝚝}(\beta L\mathrm{\Pi }_1^{(2)})=$$ $`(6.13)`$ $$\mathrm{𝚜𝚝}(L[(\overline{\lambda }^{(1)}\overline{\eta }^{(1)})(\lambda ^{(1)}\eta ^{(1)})])=$$ $$\mathrm{𝚜𝚝}(\beta L[(\overline{\lambda }_1^{(2)}\overline{\eta }_1^{(2)})(\lambda _1^{(2)}\eta _1^{(2)})]).$$ In order to obtain the so-called “time dilation” expressions, follow the same procedure as above. Notice, however, that (6.3) leads to a contradiction unless $$u((\overline{\lambda }^{(1)}+\overline{\eta }^{(1)})(\lambda ^{(1)}+\eta ^{(1)}))\beta u((\overline{\lambda }^{(2)}+\overline{\eta }^{(2)})(\lambda ^{(2)}+\eta ^{(2)})).$$ $`(6.14)`$ It is interesting, but not surprising, that this procedure yields (6.14) without hypothesizing a relation between the $`\omega _i,i=1,2,3`$ and implies that the timing infinitesimal light-clocks are the fundamental constitutes for the analysis. In the NSP-world, 6.14 can be re-expressed as $$u((\overline{\lambda }^{(1)}+\overline{\eta }^{(1)})(\lambda ^{(1)}+\eta ^{(1)}))u(\overline{\lambda }_2^{(2)}\lambda _2^{(2)}).$$ $`(6.15)`$ Or $$\mathrm{𝚜𝚝}(u((\overline{\lambda }^{(1)}+\overline{\eta }^{(1)}))=\mathrm{𝚜𝚝}(u\mathrm{\Pi }_2^{(1)})=$$ $$\mathrm{𝚜𝚝}(u\mathrm{\Pi }_3^{(2)})=\mathrm{𝚜𝚝}(u(\overline{\lambda }_2^{(2)}\lambda _2^{(2)})).$$ $`(6.16)`$ \[See note 4.\] From the N-world, the expression becomes, taking the standard part operator, $$\mathrm{𝚜𝚝}(u(\overline{\lambda }^{(1)}+\overline{\eta }^{(1)}))\mathrm{𝚜𝚝}(u(\lambda ^{(1)}+\eta ^{(1)}))=$$ $$\mathrm{𝚜𝚝}(\beta )(\mathrm{𝚜𝚝}(u(\overline{\lambda }^{(2)}+\overline{\eta }^{(2)}))\mathrm{𝚜𝚝}(u(\lambda ^{(2)}+\eta ^{(2)}))).$$ $`(6.17)`$ Or $$\mathrm{𝚜𝚝}(u\mathrm{\Pi }_2^{(1)})=\mathrm{𝚜𝚝}(\beta )\mathrm{𝚜𝚝}(u\mathrm{\Pi }_4^{(2)})=\mathrm{𝚜𝚝}(\beta u\mathrm{\Pi }_4^{(2)})=$$ $$\mathrm{𝚜𝚝}(u((\overline{\lambda }^{(1)}+\overline{\eta }^{(1)})(\lambda ^{(1)}+\eta ^{(1)})))=\mathrm{𝚜𝚝}(\beta u[(\overline{\lambda }^{(2)}+\overline{\eta }^{(2)})(\lambda ^{(2)}+\eta ^{(2)})]).$$ $`(6.18)`$ Note that using the standard part operator in the above expressions, yields continuum time and space coordinates to which the calculus can now be applied. However, the time and space measurements are not to be made with respect to an universal (absolute) clock or ruler. The measurements are relative to electromagnetic propagation. The Einstein time and length are not the NSPPM time and length, but rather they are concepts that incorporate a mode of measurement into electromagnetic field theory. This mode of measurement follows from the one wave property used for Special Theory scenarios, the property that, in the N-world, the propagation of a photon do not take on the velocity of its source. It is this that helps clarify properties of the NSPPM. Expressions such as (6.13), (6.18) will be interpreted in the next sections of this paper. 7. An Interpretation. In each of the expressions $`(6.i),i=10,\mathrm{},18`$ the infinitesimal numbers $`L,u`$ are unaltered. If this is the case, then the light-clock counts would appear to be altered. As shown in Note , alteration of $`c`$ can be represented as alterations that yield infinite counts. Thus, in one case, you have a specific infinitesimal $`L`$ and for the other infinitesimal light-clocks a different light-clock $`c`$ is used. But, $`L/u=c.`$ Consequently the only alteration that takes place in N-world expressions $`(6.i),i=12,13,17,18`$ is the infiniteimal light-clocks that need to be employed. This is exactly what (6.13) and (6.18) state if you consider it written as say, $`(\beta L)`$ rather than $`L(\beta )`$. Although these are external expressions and cannot be “formally” transferred back to the N-world, the methods of infinitesimal modeling require the concepts of “constant” and “not constant” to be preserved. These N-world expressions can be re-described in terms of N-world approximations. Simply substitute $``$ for $`=,`$ a nonzero real $`d`$ \[resp. $`\mu `$\] for $`L`$ \[resp. $`u`$\] and real natural numbers for each light-clock count in equations $`(6.i)`$,$`i=12,17.`$ Then for a particular $`d`$ \[resp. $`\mu `$\] any change in the light-clock measured relative velocity $`v_E`$ would dictate a change in the the light-clocks used. Hence, the N-world need not be concerned with the idea that “length” contracts but rather it is the required light-clocks change. It is the required change in infiniteimal light-clocks that lead to real physical changes in behavior as such behavior is compared to a standard behavior. But, in many cases, the use of light-clocks is not intended to be a literal use of such instruments. For certain scenarios, light-clocks are to be considered as analog models that incorporate electromagnetic energy properties. \[See note 18, first paragraph.\] The analysis given in the section 3 is done to discover a general property for the transmission of electromagnetic radiation. It is clear that property (\) does not require that the measured velocity of light be a universal constant. All that is needed is that for the two NSP-world times $`t_a,t_b`$ that $`\mathrm{𝚜𝚝}(\mathrm{}(t_a))=\mathrm{𝚜𝚝}(\mathrm{}_1(t_b)).`$ This means that all that is required for the most basic aspects of the Special Theory to hold is that at two NSP-world times in the $`F_1F_2,F_2F_1`$ reflection process $`\mathrm{𝚜𝚝}(\mathrm{}(t_a))=\mathrm{𝚜𝚝}(\mathrm{}_1(t_b))`$, $`t_a`$ a time during the transmission prior to reflection and $`t_b`$ after reflection. If $`\mathrm{},\mathrm{}_1`$ are nonstandard extensions of standard functions $`v,v_1`$ continuous on $`[a,b]`$, then given any $`ϵ\mathrm{𝙸𝚁}^+`$ there is a $`\delta `$ such that for each $`t,t^{}[a,b]`$ such that $`\|tt^{}\|<\delta `$ it follows that $`\|v(t)v(t^{})\|<ϵ/3`$ and $`\|v_1(t)v_1(t^{})\|<ϵ/3.`$ Letting $`t_3t_1<\delta ,`$ then $`\|t_at_b\|<\delta .`$ Since $`{}_{}{}^{}v(t_a)=\mathrm{}(t_a)^{}v_1(t_b)=\mathrm{}_1(t_b),`$ \-transfer implies $`\|^{}v(t_2)^{}v_1(t_2)\|<ϵ.`$ \[ See note 5.\] Since $`t_2`$ is a standard number, $`\|v(t_2)v_1(t_2)\|<ϵ`$ implies that $`v(t_2)=v_1(t_2).`$ Hence, in this case, the two functions $`\mathrm{},\mathrm{}_1`$ do not differentiate between the velocity $`c`$ at $`t_2.`$ But $`t_2`$ can be considered an arbitrary (i.e. NSPPM) time such that $`t_1<t_2<t_3.`$ This does not require $`c`$ to be the same for all cosmic times only that $`v(t)=v_1(t),t_1<t<t_3.`$ The restriction that $`\mathrm{},\mathrm{}_1`$ are extended standard functions appears necessary for our derivation. Also, this analysis is not related to what $`\mathrm{}`$ may be for a stationary laboratory. In the case of stationary $`F_1,F_2,`$ then the integrals are zero in equation (19) of section 3. The easiest thing to do is to simply postulate that $`\mathrm{𝚜𝚝}(^{}v(t_a))`$ is a universal constant. This does not make such an assumption correct. One of the properties that will allow the Einstein velocity transformation expression to be derived is the equilinear property. This property is weaker than the $`c`$ = constant property for light propagation. Suppose that you have within the NSP-world three observers $`F_1,F_2,F_3`$ that are linearly related. Further, suppose that $`w_1`$ is the NSP-world velocity of $`F_2`$ relative to $`F_1`$ and $`w_2`$ is the NSP-world velocity of $`F_3`$ relative to $`F_2.`$ It is assumed that for this nonmonadic cluster situation, that Galilean dynamics also apply and that $`\mathrm{𝚜𝚝}(w_1)+\mathrm{𝚜𝚝}(w_2)=\mathrm{𝚜𝚝}(w_3).`$ Using the description for light propagation as given in section 3, let $`t_1`$ be the cosmic time when a light pulse leaves $`F_1`$, $`t_2`$ when it “passes” $`F_2,`$ and $`t_3`$ the cosmic time when it arrives at $`F_3.`$ From equation (3.15), it follows that $$\mathrm{𝚜𝚝}(w_1)=\mathrm{𝚜𝚝}(^{}v_1(t_{1a}))\mathrm{𝚜𝚝}\left({}_{}{}^{}_{t_1}^{t_2}\frac{1}{x}dx\right)+[\mathrm{𝚜𝚝}(w_2)=]\mathrm{𝚜𝚝}(^{}v_2(t_{2a}))\mathrm{𝚜𝚝}\left({}_{}{}^{}_{t_2}^{t_3}\frac{1}{x}dx\right)=$$ $$\mathrm{𝚜𝚝}(w_3)=\mathrm{𝚜𝚝}(^{}v_3(t_{3a}))\mathrm{𝚜𝚝}\left({}_{}{}^{}_{t_1}^{t_3}\frac{1}{x}dx\right).$$ $`(7.1)`$ If $`\mathrm{𝚜𝚝}(^{}v_1(t_{1a}))=\mathrm{𝚜𝚝}(^{}v_2(t_{2a}))=\mathrm{𝚜𝚝}(^{}v_3(t_{3a})),`$ then we say that the velocity functions $`{}_{}{}^{}v_{1}^{},^{}v_2,^{}v_3`$ are equilinear. The constancy of $`c`$ implies equilinear, but not conversely. In either case, functions such as $`{}_{}{}^{}v_{1}^{}`$ and $`{}_{}{}^{}v_{2}^{}`$ need not be the same within a stationary laboratory after interaction. Experimentation indicates that electromagnetic propagation does “appear” to behave in the N-world in such a way that it does not accquire the velocity of the source. The light-clock analysis is consistent with the following speculation. Depending upon the scenario, the uniform velocity yields an effect via interactions with the subparticle field (the NSPPM) that uses a photon particle behavioral model. This is termed the (emis) effect. Recall that a “light-clock” can be considered as an analog model for the most basic of the electromagnetic properties. On the other hand, only those experimental methods that replicate or are equivalent to the methods of Einstein measure would be relative to the Special Theory. This is one of the basic logical errors in theory application. The experimental language must be related to the language of the derivation. The concept of the light-clock, linear paths and the like are all intended to imply NSPPM interactions. Any explanation for experimentally verified Special Theory effects should be stated in such a language and none other. I also point out that there are no paradoxes in this derivation for you cannot simply “change your mind” with respect to the NSPPM. For example, an observer is either in motion or not in motion, and not both with respect to the NSPPM. 8. A Speculation and Ambiguous Interpretations Suppose that the correct principles of infinitesimal modeling were known prior to the M-M (i.e. Michelson-Morley) experiment. Scientists would know that the (mathematical) NSPPM is not an N-world entity. They would know that they could have very little knowledge as to the refined workings of this NSP-world NSPPM since $``$ is not an $`=.`$ They would have been forced to accept the statement of Max Planck that “Nature does not allow herself to be exhaustively expressed in human thought.”\[The Mechanics of Deformable Bodies, Vol. II, Introduction to Theoretical Physics, Macmillian, N.Y. (1932),p. 2.\] Further suppose, that human comprehension was advanced enough so that all scientific experimentation always included a theory of measurement. The M-M experiment would then have been performed to learn, if possible, more about this NSP-world NSPPM. When a null finding was obtained then a derivation such as that in section 3 might have been forthcoming. Then the following two expressions would have emerged from the derivation. The Einstein method for measurement - the “radar” method - is used (see A3, p. 52) to determining the relative velocity of the moving light-clock. Using Appendix-A equations (A14), let $`P`$ correspond to $`F_2`$. Then $`\theta =0,\varphi =\pi /2.`$ Since, $`x^{(2)}=0`$ from page 54, then $`F_2`$ is the $`s`$-point Hence, $`t_E^2=t^{(2)}`$. The superscript and subscript $`s`$ represents local measurements about the $`s`$-point, using various devices, for laboratory standards (i.e. standard behavior) and using infinitesimal light-clocks or approximating devices such as atomic-clocks. \[Due to their construction atomic clocks are effected by relativistic motion and gravitational fields approximately as the infinitesimal light-clock’s counts are effected.\] Superscript or subscript $`m`$ indicates local measurements, using the same devices, for an entity considered at the $`m`$-point in motion relative to the $`s`$-point, where Einstein time and distance via the radar method as registered at $`s`$ are used to investigate $`m`$-point behavior. For example, $`m`$-point time is measured at the $`s`$-point via infinitesimal light clock and the radar method and this represents time at the $`m`$-point. To determine how physical behavior is being altered, the $`m`$ and $`s`$-measurements are compared. Many claim that you can replace each $`s`$ with $`m`$, and $`m`$ with $`s`$ in what follows. This may lead to various controversies which are elimianted in part 3. A specific interpretation of $$\mathrm{𝚜𝚝}(\beta )^1(\overline{t}^{(s)}t^{(s)})=\overline{t}_E^{(m)}t_E^{(m)}$$ $`(8.1)`$ or the corresponding $$\mathrm{𝚜𝚝}(\beta )^1(\overline{x}^{(s)}x^{(s)})=\overline{x}_E^{(m)}x_E^{(m)}$$ $`(8.2)`$ seems necessary. However, (8.2) is unnecessary since $`v_E(\mathrm{𝚜𝚝}(\beta )^1(\overline{t}^{(s)}t^{(s)}))=v_E(\overline{t}_E^{(m)}t_E^{(m)})`$ yields (8.2), which can be used when convienient. Thus, only the infinitesimal light-clock “time” alterations are significant. Actual length as measured via the radar method is not altered. It is the clock counts that are altered. If, in (8.2), which is employed for convenience, $`\overline{x}^{(s)}x^{(s)}=U^s`$ (note that $`x^{(s)}=v_Et^{(s)}`$ etc.) is interpreted as “any” standard unit for length measurement at the $`s`$-point and $`\overline{x}_E^{(m)}x_E^{(m)})=U^m`$ the same “standard” unit for length measurement in a system moving with respect to the NSPPM (without regard to direction), then for equality to take place the unit of measure $`U^m`$ may seem to be altered in the moving system. Of course, it would have been immediately realized that the error in this last statement is that $`U^s`$ is “any” unit of measure. Once again, the error in these two statements is the term “any.” (This problem is removed by application of $`(14)_a`$ or $`(14)_b`$ p. 60.) If, in (8.2), which is employed for convenience, $`\overline{x}^{(s)}x^{(s)}=U^s`$ (Note that $`x^{(s)}=v_Et^{(s)}.`$) is interpreted as “any” standard unit for length measurement at the $`s`$-point and $`\overline{x}_E^{(m)}x_E^{(m)})=U^m`$ the same “standard” unit for length measurement in a system moving with respect to the NSPPM (without regard to direction), then for equality to take place the unit of measure $`U^m`$ may seem to be altered in the moving system. Of course, it would have been immediately realized that the error in this last statement is that $`U^s`$ is “any” unit of measure. Once again, the error in these two statements is the term “any.” (This problem is removed by application of $`(14)_a`$ or $`(14)_b`$ p. 60.) Consider experiements such as the M-M, Kennedy-Thorndike and many others. When viewed from the wave state, the interferometer measurement technique is determined completely by a light-clock type process – the number of light waves in the linear path. We need to use $`L_{sc}^m,`$ a scenario associated light unit, for $`U^m`$ and use a $`L_{sc}^s`$ for $`U^s.`$ It appears for this particular scenario, that $`L_{sc}^s`$ may be considered the private unit of length in the NSP-world, such as $`L,`$ used to measure NSP-world light-path length. The “wavelength” $`\lambda `$ of any light source must also be measured in the same light units. Let $`\lambda =N^sL_{sc}^s.`$ Taking into consideration a unit conversion factor $`k`$ between the unknown NSP-world private units, such that $`\mathrm{𝚜𝚝}(kL_{sc}^s)=U^s`$, the number of light waves in $`s`$-laboratory would be $`A^s\mathrm{𝚜𝚝}(kL_{sc}^s)/N^s\mathrm{𝚜𝚝}(kL_{sc}^s)=A^s/N^s,`$ where $`A^s`$ is a pure number such that $`A^s\mathrm{𝚜𝚝}(kL_{sc}^s)`$ is the “path-length” using the units in the $`s`$-system. In the moving system, assuming that this simple aspect of light propagation holds in the NSP-world and the N-world which we did to obtain the derivation in section 3, it is claimed that substitution yields $`\mathrm{𝚜𝚝}(\beta ^1A^skL_{sc}^s)/\mathrm{𝚜𝚝}(\beta ^1N^skL_{sc}^s)=A^s\mathrm{𝚜𝚝}(\beta ^1kL_{sc}^s)/N^s\mathrm{𝚜𝚝}(\beta ^1kL_{sc}^s)=A^s/N^s.`$ Thus there would be no difference in the number of light waves in any case where the experimental set up involved the sum of light paths each of which corresponds to the to-and-fro process \[1: 24\]. Further, the same conclusions would be reached using (8.2). not relevant to a Sagnac type of experiment. However, this does not mean that a similar derivation involving a polygonal propagation path cannot be obtained. \[Indeed, this may be a consequence of a result to be derived in article 3. However, see note 8 part 4, p. 80.\] Where is the logical error in the above argument? The error is the object upon which the $`\mathrm{𝚜𝚝}(\beta )^1`$ operates. Specifically (6.13) states that $$\mathrm{𝚜𝚝}(\beta )^1(A^skL_{sc}^s)\stackrel{\left(\mathrm{emis}\right)}{}\beta ^1(L\mathrm{\Pi }^{(s)})=(\beta ^1L)\mathrm{\Pi }^{(s)}\mathrm{and}$$ $`(8.3)`$ $$\mathrm{𝚜𝚝}(\beta )^1(N^skL_{sc}^s)\stackrel{\left(\mathrm{emis}\right)}{}\beta ^1(L\mathrm{\Pi }_1^{(s)})=(\beta ^1L)\mathrm{\Pi }_1^{(s)}.$$ $`(8.4)`$ It is now rather obvious that the two (emis) aspects of the M-M experiment nullify each other. Also for no finite $`w`$ can $`\beta 0.`$ There is a great difference between the propagation properties in the NSP-world and the N-world. For example, the classical Doppler effect is an N-world effect relative to linear propagation. Rather than indicating that the NSPPM is not present, the M-M results indicate indirectly that the NSP-world NSPPM exists. Apparently, the well-known Ives-Stillwell, and all similar, experiments used in an attempt to verify such things as the relativistic redshift are of such a nature that they eliminate other effects that motion is assumed to have upon the scenario associated electromagnetic propagation. What was shown is that the frequency $`\nu `$ of the canal rays vary with respect to a representation for $`v_E`$ measured from electromagnetic theory in the form $`\nu _m=\mathrm{𝚜𝚝}(\beta )^1\nu _s.`$ First, we must investigate what the so-called time dilation statement (8.2) means. What it means is exemplified by (6.14) and how the human mind comprehends the measure of “time.” In the scenario associated (8.2) expression, for the right and left-sides to be comprehensible, the expression should be conceived of as a measure that originates with infinitesimal light-clock behavior. It is the experience with a specific unit and the number of them that “passes” that yields the intuitive concept of “observer time.” On the other hand, for some purposes or as some authors assume, (8.2) might be viewed as a change in a time unit $`T^s`$ rather than in an infinitesimal light-clock. Both of these interpretations can be incorporated into a frequency statement. First, relative to the frequency of light-clock counts, for a fixed stationary unit of time $`T^s,`$ (8.2) reads $$\mathrm{𝚜𝚝}(\beta )^1C_{sc}^s/T^sC_{sc}^m/T^s\mathrm{𝚜𝚝}(\beta )^1C_{sc}^sC_{sc}^m.$$ $`(8.5)`$ But according to (6.18), the $`C_{sc}^s`$ and $`C_{sc}^m`$ correspond to infinitesimal light-clocks measures and nothing more than that. Indeed, (8.5) has nothing to do with the concept of absolute “time” only with the different infinitesimal light-clocks that need to be used due to relative motion. This requirement may be due to (emis). Indeed, the “length contraction” expression (8.1) and the “time dilation” expression (8.2) have nothing to do with either absolute length or absolute time. These two expressions are both saying the same thing from two different viewpoints. There is an alteration due to the (emis). \[Note that the second $``$ in (8.5) depends upon the $`T^s`$ chosen.\] On the other hand, for a relativistic redshift type experiment, the usual interpretation is that $`\nu _sp/T^s`$ and $`\nu _mp/T^m.`$ This leads to $`p/T^m\mathrm{𝚜𝚝}(\beta )^1p/T^sT^m\mathrm{𝚜𝚝}(\beta )T^s.`$ Assuming that all frequency alterations due to (emis) have been eliminated then this is interpreted to mean that “time” is slower in the moving excited hydrogen atom than in the “stationary” laboratory. When compared to (8.5), there is the ambiguous interpretation in that the $`p`$ is considered the same for both sides (i.e. the concept of the frequency is not altered by NSPPM motion). It is consistent with all that has come before that the Ives-Stillwell result be written as $`\nu _sp/T^s`$ and that $`\nu _mq/T^s,`$ where “time” as a general notion is not altered. This leads to the expression $$\mathrm{𝚜𝚝}(\beta )^1pq[=\mathrm{in}\mathrm{the}\mathrm{limit}].$$ $`(8.6)`$ Expression (8.6) does not correspond to a concept of “time” but rather to the concept of alterations in emitted frequency due to (emis). One, therefore, has an ambiguous interpretation that in an Ives-Stillwell scenario the number that represents the frequency of light emitted from an atomic unit moving with velocity $`\omega `$ with respect to the NSPPM is altered due to (emis). This (emis) alteration depends upon $`K^{(3)}.`$ It is critical that the two different infinitesimal light-clock interpretations be understood. One interpretation is relative to electromagnetic propagation theory. In this case, the light-clock concept is taken in its most literal form. The second interpretation is relative to an infinitesimal light-clock as an analogue model. This means that the cause need not be related to propagation but is more probably due to how individual constituents interact with the NSPPM. The exact nature of this interaction and a non-ambiguous approach needs further investigation based upon constituent models since the analogue model specifically denies that there is some type of absolute time dilation but, rather, signifies the existences of other possible causes. \[In article 3, the $`\nu _m=\mathrm{𝚜𝚝}(\beta )^1\nu _s`$ is formally and non-ambiguously derived from a special line-element, a universal functional requirement and Schrödinger’s equation.\] In our analysis it has been assumed that $`F_1`$ is stationary in the NSP-world NSPPM. It is clear, however, that under our assumption that the scalar velocities in the NSP-world are additive with respect to linear motion, then if $`F_1`$ has a velocity $`\omega `$ with respect to the NSPPM and $`F_2`$ has the velocity $`\omega ^{},`$ then it follows that the light-clock counts for $`F_1`$ require the use of a different light-clock with respect to a stationary $`F_0`$ due to the (emis) and the light-clocks for $`F_2`$ have been similarly changed with respect to a stationary $`F_0`$ due to (emis). Consequently, a light-clock related expressed by $`K^{(3)}`$ is the result of the combination, so to speak, of these two (emis) influences. The relative NSPPM velocity $`\omega _2`$ of $`F_1`$ with respect to $`F_2`$ which yields the difference between these influences is that which would satisfies the additive rule for three linear positions. As previously stated, within the NSP-world relative to electromagnetic propagation, observer scalar velocities are either additive or related as discussed above. Within the N-world, this last statement need not be so. Velocities of individual entities are modeled by either vectors or, at the least, by signed numbers. Once the N-world expression is developed, then it can be modified in accordance with the usual (emis) alterations, in which case the velocity statements are N-world Einstein measures. For example, deriving the so-called relativistic Dopplertarian effect, the combination of the classical and the relativistic redshift, by means of a NSPPM argument such as appears in where it is assumed that the light propagation laws with respect to the photon concept in the NSP-world are the same as those in the N-world, is in logical error. Deriving the classical Doppler effect expression then, when physically justified, making the wave number alteration in accordance with the (emis) would be the correct logic needed to obtain the relativistic Dopplertarian effect. \[See note 6.\] Although I will not, as yet, re-interpreted Special Relativity results with respect to this purely electromagnetic interpretation, it is interesting to note the following two re-interpretations. The so-called variation of “mass” was, in truth, originally derived for imponderable matter (i.e. elementary matter.) This would lead one to believe that the so-called rest mass and its alteration, if experimentally verified, is really a manifestation of the electromagnetic nature of such elementary matter. Once again the so-called mass alteration can be associated with an (emis) concept. The $`\mu `$-meson decay rate may also show the same type of alteration as appears to be the case in an Ives-Stillwell experiment. It does not take a great stretch of the imagination to again attribute the apparent alteration in this rate to an (emis) process. This would lead to the possibility that such decay is controlled by electromagnetic properties. Indeed, in order to conserve various things, $`\mu `$-meson decay is said to lead to the generation of the neutrino and antineutrino. \[After this paper was completed, a method was discovered that establishes that predicted mass and decay time alterations are (emis) effects. The derivations are found in article 3.\] I note that such things as neutrinos and antineutrinos need not exist. Indeed, the nonconservation of certain quantities for such a scenario leads to the conclusion that subparticles exist within the NSP-world and carry off the “missing” quantities. Thus the invention of such objects may definitely be considered as only a bookkeeping technique. As pointed out, all such experimental verification of the properly interpreted transformation equations can be considered as indirect evidence that the NSP-world NSPPM exists. But none of these results should be extended beyond the experimental scenarios concerned. Furthermore, I conjecture that no matter how the human mind attempts to explain the (emis) in terms of a human language, it will always be necessary to postulate some interaction process with the NSPPM without being able to specifically describe this interaction in terms of more fundamental concepts. Finally, the MA-model specifically states that the Special Theory is a local theory and should not be extended, without careful consideration, beyond a local time interval $`[a,b].`$ 9. Reciprocal Relations As is common to many mathematical models, not all relations generated by the mathematics need to correspond to physical reality. This is the modern approach to the length contradiction controversy . Since this is a mathematical model, there is a theory of correspondence between the physical language and the mathematical structure. This correspondence should be retained throughout any derivation. This is a NSPPM theory and what is stationary or what is not stationary with respect to the NSPPM must be maintained throughout any correspondence. This applies to such reciprocal relations as $$\mathrm{𝚜𝚝}(\beta )^1(\overline{t}_E^{(m)}t_E^{(m)})=\overline{t}^{(s)}t^{(s)}$$ $`(9.1)`$ and $$\mathrm{𝚜𝚝}(\beta )^1(\overline{x}_E^{(m)}x_E^{(m)})=\overline{x}^{(s)}x^{(s)}$$ $`(9.2)`$ Statement (8.1) and (9.1) \[resp. (8.2) and (9.2)\] both hold from the NSPPM viewpoint only when $`v_E=0`$ since it is not the question of the N-world viewpoint of relative velocity but rather the viewpoint that $`F_1`$ is fixed and $`F_2`$ is not fixed in the NSPPM or $`\omega \omega ^{}.`$ The physical concept of the $`(s)`$ and $`(m)`$ must be maintained throughout the physical correspondence. Which expression would hold for a particular scenario depends upon laboratory confirmation. This is a scenario associated theory. All of the laboratory scenarios discussed in this paper use infinitesimalized (9.1) and (9.2) as derived from line-elements and the “view” or comparison is always made relative to the $`(s)`$. Other authors, such as Dingle and Builder , have, in a absolute sense, excepted one of these sets of equations, without derivation, rather the other set. I have not taken this stance in this paper. One of the basic controversies associated with the Special Theory is whether (8.2) or (8.1) \[resp. (9.1) or (9.2)\] actually have physical meaning. The notion is that either “length” is a fundamental concept and “time” is defined in terms of it, or “time” is a fundamental concept and length is defined in terms of it. Ives, and many others assumed that the fundamental notion is length contraction and not time dilation. The modern approach is the opposite of this. Length contraction in the N-world has no physical meaning, but time dilation does . We know that time is often defined in terms of length and velocities. But, the length or time being considered here is Einstein length or Einstein time. This is never mentioned when this problem is being considered. As discussed at the end of section 3, Einstein length is actually defined in terms of infinitesimal light-clocks or in terms of the Einstein velocity and Einstein time. As shown after equation (8.2) is considered, it is only infinitesimal light-clock “time” that is altered and length altertions is but a technical artefact. The changes in the infinitesimal light-clock counts yields an analogue model for physical changes that cause Special Theory effects. \[See note 7.\] $`\{`$Remark: Karl Popper notwithstanding, it is not the sole purpose of mathematical models to predict natural system behavior. The major purpose is to maintain logical rigor and, hopefully, when applicable to discover new properties for natural systems. I have used in this speculation a correspondence theory that takes the stance that any verifiable Special Theory effect is electromagnetic in character rather than a problem in measure. However, whether such effects are simply effects relative to the propagation of electromagnetic information or whether they are effects relative to the constituents involved cannot be directly obtain from the Special Theory. All mathematically stated effects involve the Einstein measure of relative velocity, $`v_E`$ – a propagation related measure. The measure of an effect should also be done in accordance with electromagnetic theory. As demonstrated, the Special Theory should not be unnecessarily applied to the behavior of all nature systems since it is related to electromagnetic interaction; unless, of course, all natural systems are electromagnetic in character. Without strong justification, the assumption that one theory does apply to all scenarios is one of the greatest errors in mathematically modeling. But, if laboratory experiments verify that alterations are taking place in measured quantities and these variations are approximated in accordance with the Special Theory, then this would indicate that either the alterations are related to electromagnetic propagation properties or the constituents have an appropriate electromagnetic character.$`\}`$ NOTES (a) Equation (3.9) is obtained as follows: since $`t[a,b],`$ $`t`$ finite and not infinitesimal. Thus division by $`t`$ preserves $`.`$ Hence, $$\left[t\left(\frac{{}_{}{}^{}s(t+dt)s(t)}{dt}\right)s(t)\right]/t^2\frac{\mathrm{}(t)}{t}.$$ $`(1)`$ Since $`t`$ is an arbitrary standard number and $`dt`$ is assume to be an arbitrary and appropriate nonzero infinitesimal and the function $`s(t)/t`$ is differentiable, the standard part of the left-side equals the standard part of the right-side. Thus $$\frac{d(s(t)/t)}{dt}=\frac{v(t)}{t},$$ $`(2)`$ for each $`t[a,b].`$ By \-transfer, equation (3.9) holds for each $`t^{}[a,b].`$ (b) Equation (3.10) is then obtained by use of the \-integral and the fundamental theorem of integral calculus \-transferred to the NSP-world. It is useful to view the definite integral over a standard interval say $`[t_1,t]`$ as an operator, at least, defined on the set $`C([t_1,t],\mathrm{𝙸𝚁})`$ of all continuous real valued functions defined on $`[t_1,t].`$ Thus, in general, the fundamental theorem of integral calculus can be viewed as the statement that $`(f^{},f(t)f(t_1))_{t_1}^t.`$ Hence $`{}_{}{}^{}(f^{},f(t)f(t_1))^{}_{t_1}^t(^{}f^{},^{}(f(t)f(t_1)))^{}_{t_1}^t(^{}f^{},f(t)f(t_1))^{}_{t_1}^t.`$ (c) To obtain the expressions in (3.19), consider $`f(x)=1/x.`$ Then $`{}_{}{}^{}f`$ is limited and S-continuous on $`{}_{}{}^{}[a,b].`$ Hence $`(^{}f,\mathrm{ln}t_2\mathrm{ln}t_1)^{}_{t_1}^{t_2}.`$ Hence $`\mathrm{𝚜𝚝}((^{}f,\mathrm{ln}t_2\mathrm{ln}t_1))=(f,\mathrm{ln}t_2\mathrm{ln}t_1)_{t_1}^{t_2}.`$ Further (3.19) can be interpreted as an interaction property. \[1.5\] Infinitesimal light-clocks are based upon the QED model as to how electrons are kept in a range of distances in a hydrogen atom proton. The back-and-forth exchanges of photons between a proton and electron replaces “reflection” and the average distance between the proton and electron is infinitesimalized to the $`L`$. In this case, the proton and electron are also infinitesimalized. The large number of such interchanges over a second, in the model, is motivation for the use of the members of $`\mathrm{𝙸𝙽}_{\mathrm{}}^+`$ as count numbers. The basic theorem that allows for the entire concept of infinitesimal light-clocks and the analysis that appears in this monograph has not been stated. As taken from “The Theory of Ultralogics,”the theorem, for this application, is: Theorem 11.1.1 Let $`10^\omega \mathrm{𝙸𝙽}_{\mathrm{}}.`$ Then for each $`r\mathrm{𝙸𝚁}`$ there exists an $`x\{2m/10^\omega (2m^{}𝐙)(\|2m\|<\lambda 10^\omega )\},`$ for any $`\lambda \mathrm{𝙸𝙽}_{\mathrm{}},`$ such that $`xr`$ (i.e. $`x\mu (r).)`$ Theorem 11.1.1 holds for other members of $`\mathrm{𝙸𝙽}_{\mathrm{}}.`$ Let $`L=1/10^\omega `$ where $`\omega `$ is any hyperreal infinite natural number (i.e. $`\omega \mathrm{𝙸𝙽}_{\mathrm{}}).`$ Hence, by this theorem, for any positive real number $`r`$ there exists some $`m\mathrm{𝙸𝙽}_{\mathrm{}}`$ such that $`2\mathrm{𝚜𝚝}(m/10^\omega )=r.`$ I point out that for this nonzero case it is necessary that $`m\mathrm{𝙸𝙽}_{\mathrm{}}`$ for if $`m\mathrm{𝙸𝙽},`$ then $`\mathrm{𝚜𝚝}(m/10^\omega )=0.`$ Since $`c=\mathrm{𝚜𝚝}(L/u),`$ then $`2\mathrm{𝚜𝚝}(um)=2\mathrm{𝚜𝚝}((L/c)m)=t=r/c`$ as required. Thus, the infinitesimal light-clock determined length $`r`$ and interval of time $`t`$ are determined by the difference in infinitesimal light-clock counts $`2m=(\lambda \eta )`$. Note that our approach allows the calculus to model this behavior by simply assuming that the standard functions are differentiable etc. \[2.5\] (4 JUN 2000) Equating these counts here and elsewhere is done so that the “light pulse” is considered to have a “single instantaneous effect” from a global viewpoint and as such is not a signal in that globally it contains no information. Thus additional analysis is needed before one can state that the Special Theory applies to informational transmissions. It’s obvious from section 7 that the actual value for $`c`$ may depend upon the physical application of this theory. At this point and on, the subscripts on the $`\tau `$ have a different meaning than previously indicated. The subscripts denote process numbers while the superscript denotes the position numbers. For example, $`\tau _{12}^2`$ means the light-clock count number when the second light pulse leaves $`F_2`$ and $`\tau _{31}^2`$ would mean the light-clock count number when the first light pulse returns to position $`F_2.`$ The additional piece of each subscript denoted by the $`a`$ on this and the following pages indicates, what I thought was obvious from the lines that follow their introduction, that these are approximating numbers that are infinitesimally near to standard NSP-world number obtained by taking the standard part. Note that such infinite hyperreal numbers as $`\mathrm{\Pi }_3^{(2)}`$ (here and elsewhere) denote the difference between two infinitesimal light-clock counts and since we are excluding the finite number infinitesimally near to 0, these numbers must be infinite hyperreal. Infinitesimal light-clocks can be assumed to measure this number by use of a differential counter. BUT it is always to be conceived of as an infinitesimal light-clock “interval” (increment, difference, etc.) It is important to recall this when the various line-elements in the next article are considered. This result is obtained as follows: since $`t_at_2t_b,`$ it follows that $`\|t_at_2\|<\delta ,\|t_bt_2\|<\delta .`$ Hence by \-transfer, $`\|^{}v(t_2)^{}v(t_a)\|<ϵ/3,\|^{}v_1(t_b)^{}v_1(t_2)\|<ϵ/3.`$ Since we assume arbitrary $`ϵ/3`$ is a standard positive number, then $`{}_{}{}^{}v(t_a)=\mathrm{}(t_a)^{}v_1(t_b)=\mathrm{}_1(t_b)\|^{}v(t_a)^{}v_1(t_b)\|<ϵ/3.`$ Hence $`\|^{}v(t_2)^{}v_1(t_2)\|<ϵ.`$ In this article, I mention that all previous derivations for the complete Dopplertarian effect (the N-world and the transverse) are in logical error. Although there are various reasons for a redshift not just the Dopplertarian, the electromagnetic redshift based solely upon properties of the NSPPM can be derived as follows: (i) let $`\nu ^s`$ denote the “standard” laboratory frequency for radiation emitted from an atomic system. This is usually determined by the observer. The NSP-world alteration in emitted frequency at an atomic structure due to (emis) is $`\gamma \nu ^s=\nu ^{\mathrm{radiation}},`$ where $`\gamma =\sqrt{1v_E^2/c^2}`$ and $`v_E`$ is the Einstein measure of the relative velocity using light-clocks only. (ii) Assuming that an observer is observing this emitted radiation in a direct line with the propagation and the atomic structure is receding with velocity $`v`$ from the observer, the frequency of the electromagnetic propagation, within the N-world, is altered compared to the observers standards. This alteration is $`\nu ^{\mathrm{radiation}}(1/(1+v/c))=\nu ^{\mathrm{received}}.`$ Consequently, this yields the total alteration as $`\gamma \nu ^s(1/(1+v/c))=\nu ^{\mathrm{received}}.`$ Note that $`v`$ is measured in the N-world and can be considered a directed velocity. Usually, if due to the fact that we are dealing with electromagnetic radiation, we consider $`v`$ the Einstein measure of linear velocity (i.e. $`v=v_E`$), then the total Dopplertarian effect for $`v0`$ can be written as $$\nu ^s\left(\frac{1v_E/c}{1+v_E/c}\right)^{1/2}=\nu ^{\mathrm{received}}.$$ $`(3)`$ It should always be remembered that there are other reasons, such as the gravitational redshift and others yet to be analyzed, that can mask this total Dopplertarian redshift. A question that has been asked relative to the new derivation that yields Special Theory resilts is why in the N-world do we have the apparent nonballistic effects associated with electromagnetic radiation? In the derivation, the opposite was assumed for the NSP-world monadic clusters. The constancy of the measure, by light-clocks and the like, of the $`F_1F_2,F_2F_1`$ velocity of electromagnetic radiation was modeled by letting $`\mathrm{𝚜𝚝}(t_a)=\mathrm{𝚜𝚝}(t_b).`$ As mentioned in the section on the Special Theory, the Einstein velocity measure transformation expression can be obtained prior to embedding the world into a hyperbolic velocity space. It is obtained by considering three in-line standard positions $`F_1,F_2,F_3`$ that have the NSP-world velocities $`w_1`$ for $`F_2`$ relative to $`F_1`$, $`w_2`$ for $`F_3`$ relative to $`F_2`$ and the simple composition $`w_3=w_1+w_2`$ for $`F_3`$ relative to $`F_1`$. Then simple substitution in this expression yields $$v_E^{(3)}=(v_E^{(1)}+v_E^{(2)})/\left(1+\frac{v_E^{(1)}v_E^{(2)}}{c^2}\right).$$ $`(4)`$ This relation is telling us something about the required behavior in the N-world of electromagnetic radiation. To see that within the N-world we need to assume for electromagnetic radiation effects the nonballistic property, simply let $`v_E^{(2)}=c`$ or $`v_E^{(2)}c.`$ Then $`v_E^{(3)}=c,`$ or $`c.`$ Of course, the reason we do not have a contradiction is that we have two distinctly different views of the behavior of electromagnetic radiation, the NSP-world view and the N-world view. Further, note how, for consistency, the velocity of electromagnetic radiation is to be measured. It is measured by the Einstein method, or equivalent, relative to a to-and-fro path and measures of “time” and “distance” by means of a (infinitesimal) light-clock counts. Since one has the NSPPM, then letting $`F_1`$ be fixed in that medium, assuming that “absolute” physical standards are measured at $`F_1,`$ equation (4) indicates why, in comparison, physical behavior varies at $`F_2`$ and $`F_3`$. The hyperbolic velocity space properties are the cause for such behavior differences. I am convinced that the dual character of the Special theory derivation requires individual reflection in order to be understood fully. In the NSP-world, electromagnetic radiation behaves in one respect, at least, like a particle in that it satisfies the ballistic nature of particle motion. The reason that equation (3) is derivable is due to the definition of Einstein time. But Einstein time, as measured by electromagnetic pulses, models the nonballistic or one and only one wave-like property in that a wave front does not partake of the velocity of the source. This is the reason why I wrote that a NSPPM disturbance would trace the same operational linear light-clock distance. The measuring light-clocks are in the N-world in this case. $`F_1`$ is modeled as fixed in the NSPPM and $`F_2`$ has an NSP-world relative velocity. The instant the light pulse is reflected back to $`F_1`$ it does not, from the N-world viewpoint, partake of the N-world relative velocity and therefore traces out the exact same apparent N-world linear path. The position $`F_2`$ acts like a virtual position having no other N-world effect upon the light pulse except a reversal of direction. This expression implies that the “$`c`$” that appears here and elsewhere is to be measured by infinitesimal light-clocks. As noted $`uL/c,`$ but infinitesimal light-clock construction yields that $`u=L/c.`$ For a fixed $`L`$, from the NSPPM viewpoint, $`u`$ is fixed. Notice that $`t^{(i)}u(2\eta ^{(i)})=u(\gamma ^{(i)}),\gamma ^{(i)}\mathrm{𝙸𝙽}_{\mathrm{}}^+.`$ REFERENCES 1 Dingle, H., The Special Theory of Relativity, Methuen’s Monographs on Physical Subjects, Methuen & Co. LTD, London, 1950. 2 Herrmann, R. A., Some Applications of Nonstandard Analysis to Undergraduate Mathematics – Infinitesimal Modeling, Instructional Development Project, Mathematics Dept. U. S. Naval Academy, Annapolis, MD 21402-5002, 1991. http://arxiv.org/abs/math/0312432 3 Herrmann, R. A., Fractals and ultrasmooth microeffects. J. Math. Physics., 30(1989), 805–808. 4 Herrmann, R. A., Physics is legislated by a cosmogony. Speculat. Sci. Technol., 11(1988), 17–24. 5 Herrmann, R. A., Rigorous infinitesimal modelling, Math. Japonica 26(4)(1981), 461–465. 6 Lawden, D. F., An Introduction to Tensor Calculus, Relativity and Cosmology, John Wiley & Son, New York, 1982. 7 Prokhovnik, S. J., The Logic of Special Relativity, Cambridge University Press, Cambridge, 1967. 8 Stroyan, K. D. and Luxemburg, W. A. J. 1976. Introduction to the Theory of Infinitesimals, Academic Press, New York. 9 Herrmann, R. A., The Theory of Ultralogics, (1992) http://arxiv.org/abs/math.GM/9903081 http://arxiv.org/abs/math.GM/9903082 10 Davis M., Applied Nonstandard Analysis, John Wiley & Sons, New York, 1977. Appendix-A 1. The Need for Hyperbolic Geometry In this appendix, it is shown that from equations (3.21) and (3.22) the Lorentz transformation are derivable. All of the properties for the Special Theory are based upon “light” propagation. In Article 2, the concern is with two positions $`F_1,F_2`$ in the NSPPM within the NSP-world and how the proposed NSPPM influences such behavior. Prior to applications to the N-world, with the necessity for the N-world Einstein measures, the NSPPM exhibits infinitesimal behavior and special NSPPM non-classical global behavior. The behavior at specific moments of NSPPM time for global positions and classical uniform velocities are investigated. The following is a classical description for photon behavior. Only NSPPM relative velocities (speeds) are being considered. Below is a global diagram for four points that began as the corners of a square, where $`u`$ and $`\omega `$ denote uniform relative velocities between point locations and no other point velocities are considered. The meanings for the symbolized entities are discussed below. $`F_1t`$ $``$ $`\omega F_2(t(p_1))`$ $`F_1^{}t^{}`$ $`u`$ $`F_1^{}t(p_2)`$ $`\omega F_2^{}(t(p_2))`$ $`u`$ $`u`$ Consider the following sequence of (conceptual) NSPPM time-ordered events. First, the N-world position points $`F_1,F_2,F_1^{},F_2^{}`$ are stationary with respect to each other and form the corners of a very small rhombus, say the side-length is the average distance $`d`$ between the electron and proton within an hydrogen atom. The sides are $`\overline{F_1,F_2},\overline{F_2F_2^{}},\overline{F_2^{},F_1^{}},\overline{F_1^{}F_1}`$. At the NSPPM time $`t_g`$, the almost coinciding $`F_2,F_2^{}`$ uniformally recede from the almost coinciding $`F_1,F_1^{}`$ with constant velocity $`\omega `$. At a time $`t>t_g,`$ where the distance between the two groups is significantly greater than $`d`$, one process occurs simultaneously. The point $`F_1^{}`$ separates from $`F_1`$ with relative velocity $`u`$ and $`F_2^{}`$ separates from $`F_2`$ with a relative velocity $`u`$. \[Using NSP-world processes, such simultaneity is possible relative to a non-photon transmission of information (Herrmann, 1999).\] At any time $`t,`$ the elongating line segments $`\overline{F_1F_1^{}}`$ and $`\overline{F_2F_2^{}}`$ are parallel and they are not parallel to the parallel elongating line segments $`\overline{F_1F_2}`$ and $`\overline{F_1^{}F_2^{}}.`$ At NSPPM time $`t,`$ a photon $`p_1`$ is emitted from $`F_1`$ towards $`F_2`$ and passes through $`F_2`$ and continues on. As $`F_1^{}`$ recedes from $`F_1`$, at $`t^{}>t,`$ a photon $`p_2`$ is emitted from $`F_1^{}`$ towards $`F_2^{}`$. The original classical photon-particle property that within a monadic cluster photons prorogate with velocity $`\omega +c`$ is extended to this global environment. \[Again there are NSP-world processes that can ensure that the emitted photons acquire this prorogation velocity (Herrmann, 1999).\] Also, this classical photon-particle property is applied to $`u`$. Thus, photon $`p_2`$ is assumed to take on an additional velocity component $`u`$. Photon, $`p_1`$, passes through $`F_2`$ at the NSPPM time $`t(p_1)`$. Then $`p_2`$ is received at point $`F_2^{}`$ at time $`t(p_2)`$. Classically, $`t(p_1^{})>t(p_1).`$ From a viewpoint relative to elongating $`\overline{F_1F_1^{}}`$, the distance between the two photon-paths of motion measured parallel to elongating $`\overline{F_1F_1^{}}`$ is $`u(t(p_2)t).`$ On the other hand, from the viewpoint of elongating $`\overline{F_1F_1^{}}`$, the distance between photon-paths, if they were parallel, is $`u(t^{}t).`$ By the relativity principle, from the viewpoint of $`F_1^{}`$, the first equation in (3.19) should apply. Integrating, where $`\mathrm{𝚜𝚝}(^{}v(t_a)))=c`$, one obtains $`u(t(p_2)t))=ue^{\omega /c}(t^{}t).`$ \[Note: No reflection is required for this restricted application of (3.19).\] This result is not the classical expression $`u(t(p_2)t)).`$ For better comprehension, use infinitesimal light-clocks to measure NSPPM time. Then using the same NSPPM process that yields information instantaneously throughout the standard portion of the NSPPM, all clocks used to determine these times can be set at zero when they indicate the time $`t`$. This yields that the two expressions for the distance are $`ut(p_2)`$ and $`ue^{\omega /c}t^{}.`$ However, the classical expression $`ut(p_2)`$ has the time $`t(p_2)`$ dependent upon both $`\omega `$ and, after the $`t^{}`$ moment, upon $`u.`$ But, for the relativistic expression, the $`t^{}`$ is neither dependent upon the $`u`$ velocity after $`t^{}`$ nor the $`\omega `$ and the factor $`e^{\omega /c}`$ has only one variable $`\omega .`$ What property does this NSPPM behavior have that differentiates it from the classical? Consider the two velocities $`u`$ and $`ue^{\omega /c}.`$ These two velocities only correspond when $`\omega =0.`$ Hence, if we draw a velocity diagram, one would conclude that, in this case, the velocities are trivially “parallel.” Using Lobatchewskian’s horocycle construction, Kulczycki (1961) shows that for “parallel geometric” lines in hyperbolic space, the distance between each pair of such lines increases (or decreases) by a factor $`e^{x/k},`$ as one moves an ordinary distance $`x`$ along the lines and $`k`$ is some constant related to the $`x`$ unit of measurement. Phrasing this in terms of velocities, where $`x=\omega `$ and $`k=c`$, then, for this case, the velocities, as represented in the NSPPM by standard real numbers, appear to satisfy the properties for an hyperbolic velocity-space. Such velocity behavior would lead to this non-classical NSPPM behavior. When simple classical physics is applied to this simple Euclidian configuration within the NSPPM, then there is a transformation $`\mathrm{\Phi }:\mathrm{NSPPM}`$ N-world, which is characterized by hyperbolic velocity-space properties. This is also the case for relative velocity and collinear points, which are exponentially related to the Einstein measure of relative velocity in the N-world. In what follows, this same example is used but generalized slightly by letting $`F_1`$ and $`F_2`$ coincide. 2. The Lorentz Transformations Previously, we obtained the expression that $`t_2=\sqrt{t_1t_3}.`$ The Einstein measures are defined formally as $$\{\begin{array}{cc}t_E=(1/2)(t_3+t_1)\hfill & \\ r_E=(1/2)c(t_3t_1)\hfill & \\ v_E=r_E/t_E,\mathrm{where}\mathrm{defined}.\hfill & \end{array}$$ $`(A1)`$ Notice that when $`r_E=0,`$ then $`v_E=0`$ and $`t_E=t_3=t_1=t_2`$ is not Einstein measure. The Einstein time $`t_E`$ is obtained by considering the “flight-time” that would result from using one and only one wave-like property not part of the NSPPM but within the N-world. This property is that the $`c`$ is not altered by the velocity of the source. This Einstein approach assumes that the light pulse path-length from $`F_1`$ to $`F_2`$ equals that from $`F_2`$ back to $`F_1.`$ Thus, the Einstein flight-time used for the distance $`r_E`$ is $`(t_3t_1)/2`$. The $`t_E,`$ the Einstein time corresponding to an infinitesimal light-clock at $`F_2,`$ satisfies $`t_3t_E=t_Et_1.`$ From $`(A1),`$ we have that $$t_3=(1+v_E/c)t_E\mathrm{and}t_1=(1v_E/c)t_E,$$ $`(A2)`$ and, hence, $`t_2=(\sqrt{1v_E^2/c^2})t_E.`$ Since $`e^{\omega /c}=\sqrt{t_3/t_1},`$ this yields $$e^{\omega /c}=\left(\frac{1+v_E/c}{1v_E/c}\right)^{(1/2)}.$$ $`(A3)`$ Although it would not be difficult to present all that comes next in terms of the nonstandard notions, it is not necessary since all of the functions being consider are continuous and standard functions. The effect the NSPPM has upon the N-world are standard effects produced by application of the standard part operator “st.” From the previous diagram, let $`F_1`$ and $`F_2`$ coincide and not separate. Call this location $`P`$ and consider the diagram below. This is a three position classical NSPPM light-path and relative velocity diagram used for the infinitesimal light-clock analysis in section 6 of Article 2. This diagram is not a vector composition diagram but rather represents linear light-paths with respect to Einstein measures for relative velocities. It is also a relative velocity diagram to which hyperbolic “geometry” is applied. $`P`$ $`\omega _1\omega _2`$ $`\|`$ $`\|`$ $`\|`$ $`\|`$ $`\|n`$ $`\|`$ $`\mathrm{}\mathrm{-}\mathrm{-}`$$`p_1`$$`\mathrm{-}\mathrm{-}\mathrm{}`$$`\|`$$`\mathrm{}\mathrm{-}\mathrm{-}`$$`p_2`$$`\mathrm{-}\mathrm{-}\mathrm{}`$ $`\varphi `$ $`\omega _1`$ $`\theta `$ $``$ $`\omega _2`$ $`F_1`$ $`\omega _3`$ $`F_2`$ Since Einstein measures are to be associated with this diagram, then this diagram should be obtained relative to infinitesimal light-clock counts and processes in the NSPPM. The three locations $`F_1,F_2,P`$ are assumed, at first, to coincide. When this occurs, the infinitesimal light-clock counts coincide. The positions $`F_2,F_1`$ recede from each other with velocity $`\omega _3`$. The positions $`F_2,F_1`$ recede from each other with velocity $`\omega _3`$. The object denoted by location $`P`$ recedes from the $`F_1,F_2`$ locations with uniform NSPPM velocities, in standard form, of $`\omega _1,\omega _2`$, respectively. Further, consider the special case where both are observing the pulse sent from $`P`$ at the exact some $`P`$-time. This produces the internal angle $`\theta `$ and exterior angle $`\varphi `$ for this velocity triangle. The segments marked $`p_1`$ and $`p_2`$ are the projections of the velocity representations (not vectors) $`F_1P`$ and $`F_2P`$ onto the velocity representation $`F_1F_2`$. The $`n`$ is the usual normal for this projection. We note that $`p_1+p_2=\omega _3.`$ We apply hyperbolic trigonometry in accordance with , where we need to consider a particular $`k`$. We do this by scaling the velocities in terms of light units and let $`k=c`$. From \[2, p. 143\] $$\{\begin{array}{cc}\mathrm{tanh}(p_1/c)=(\mathrm{tanh}(\omega _1/c))\mathrm{cos}\theta \hfill & \\ \mathrm{tanh}(p_2/c)=(\mathrm{tanh}(\omega _2/c))\mathrm{cos}\varphi \hfill & \end{array},$$ $`(A4)`$ and also $$\mathrm{sinh}(n/c)=(\mathrm{sinh}(\omega _1/c))\mathrm{sin}\theta =(\mathrm{sinh}(\omega _2/c))\mathrm{sin}\varphi .$$ $`(A5)`$ Now, eliminating $`\theta `$ from $`(A4)`$ and $`(A5)`$ yields \[1, p. 146\] $$\mathrm{cosh}(\omega _1/c)=(\mathrm{cosh}(p_1/c))\mathrm{cosh}(n/c).$$ $`(A6)`$ Combining (A4), (A5) and (A6) leads the hyperbolic cosine law \[2, p. 167\]. $$\mathrm{cosh}(\omega _1/c)=(\mathrm{cosh}(\omega _2/c))\mathrm{cosh}(\omega _3/c)+(\mathrm{sinh}(\omega _2/c))(\mathrm{sinh}(\omega _3/c))\mathrm{cos}\varphi .$$ $`(A7)`$ From $`(A3),`$ where each $`v_i`$ is the Einstein relative velocity, we have that $$e^{\omega _i/c}=\left(\frac{1+v_i/c}{1v_i/c}\right)^{(1/2)},i=1,2,3.$$ $`(A3)^{}`$ From the basic hyperbolic definitions, we obtain from $`(A3)^{}`$ $$\{\begin{array}{cc}\mathrm{tanh}(\omega _i/c)=v_i/c\hfill & \\ \mathrm{cosh}(\omega _i/c)=(1v_i^2/c^2)^{1/2}=\beta _i\hfill & \\ \mathrm{sinh}(\omega _i/c)=\beta _iv_i/c\hfill & \end{array}.$$ $`(A8)`$ Our final hyperbolic requirement is to use $$\mathrm{tanh}(\omega _3/c)=\mathrm{tanh}(p_1/c+p_2/c)=\frac{\mathrm{tanh}(p_1/c)+\mathrm{tanh}(p_2/c)}{1+(\mathrm{tanh}(p_1/c))\mathrm{tanh}(p_2/c)}.$$ $`(A9)`$ Now into $`(A9),`$ substitute $`(A4)`$ and then substitute the first case from $`(A8).`$ One obtains $$v_1\mathrm{cos}\theta =\frac{v_3v_2\mathrm{cos}\varphi }{1\alpha },\alpha =\frac{v_3v_2\mathrm{cos}\varphi }{c^2}.$$ $`(A10)`$ Substituting into $`(A7)`$ the second and third cases from $`(A8)`$ yields $$\beta _1=\beta _2\beta _3(1\alpha ),\beta _i=(1v_i^2/c^2)^{1/2}.$$ $`(A11)`$ From equations $`(A11),(A5)`$ and the last case in $`(A8)`$ is obtained $$v_1\mathrm{sin}\theta =\frac{v_2\mathrm{sin}\varphi }{\beta _3(1\alpha )}.$$ $`(A12)`$ For the specific physical behavior being displayed, the photons received from $`P`$ at $`F_1`$ and $`F_2`$ are “reflected back” at the NSPPM $`P`$-time $`t^r.`$ We then apply to this three point scenario our previous results. \[Note: For comprehension, it may be necessary to apply certain relative velocity viewpoints such as from $`F_1`$ the point $`P`$ is receding from $`F_1`$ and $`F_2`$ is receding from $`P`$. In this case, the NSPPM times when the photons are sent from $`F_1`$ and $`F_2`$ are related. Of course, as usual there is assumed to be no time delay between the receiving and the sending of a “reflected” photon.\] In this case, let $`t^{(1)},r^{(1)},v_1`$ be the Einstein measures at $`F_1`$ for this $`P`$-event, and $`t^{(2)},r^{(2)},v_2`$ be the Einstein measures at $`F_2`$. Since $`t^r=\beta _1^1t^{(1)},t^r=\beta _2^1t^{(2)}`$ (p. 52), then $$\frac{t^{(1)}}{\beta _1}=\frac{t^{(2)}}{\beta _2}\mathrm{and}r^{(1)}=v_1t^{(1)},r^{(2)}=v_2t^{(2)}.$$ $`(A13)`$ Suppose that we have the four coordinates, three rectangular, for this $`P`$ event as measured from $`F_1=(x^{(1)},y^{(1)},z^{(1)},t^{(1)})`$ and from $`F_2=(x^{(2)},y^{(2)},z^{(2)},t^{(2)})`$ in a three point plane. It is important to recall that the $`x,y,z`$ are related to Einstein measures of distance. Further, we take the $`x`$-axis as that of $`F_1F_2`$. The $`v_3`$ is the Einstein measure of the $`F_2`$ velocity as measured by an inf. light-clock at $`F_1`$. To correspond to the customary coordinate system employed \[1, p. 32\], this gives $$\{\begin{array}{cc}x^{(1)}=v_1t^{(1)}\mathrm{cos}\theta ,y^{(1)}=v_1t^{(1)}\mathrm{sin}\theta ,z^{(1)}=0\hfill & \\ x^{(2)}=v_2t^{(2)}\mathrm{cos}\varphi ,y^{(2)}=v_2t^{(2)}\mathrm{sin}\varphi ,z^{(2)}=0\hfill & \end{array}.$$ $`(A14)`$ It follows from $`(A10),\mathrm{},(A14)`$ that $$t^{(1)}=\beta _3(t^{(2)}v_3x^{(2)}/c^2),x^{(1)}=\beta _3(x^{(2)}v_3t^{(2)}),y^{(1)}=y^{(2)},z^{(1)}=z^{(2)}.$$ $`(A15)`$ Hence, for this special case $`\omega _1,\omega _2,\theta ,\varphi `$ are eliminated and the Lorentz Transformations are established. If $`PF_1,PF_2`$, then the fact that $`x^{(1)},x^{(2)}`$ are not the measures for a physical ruler but are measures for a distance related to Einstein measures, which are defined by the properties of the propagation of electromagnetic radiation and infinitesimal light-clock counts, shows that the notion of actual natural world “length” contraction is false. For logical consistency, Einstein measures as determined by the light-clock counts are necessary. This analysis is relative to a “second” pulse when light-clock counts are considered. The positions $`F_1`$ and $`F_2`$ continue to coincide during the first pulse light-clock count determinations. Infinitesimal light-clock counts allow us to consider a real interval as an interval for “time” measure as well as to apply infinitesimal analysis. This is significant when the line-element method in Article 3 is applied to determine alterations in physical behavior. All of the coordinates being considered must be as they would be understood from the Einstein measure viewpoint. The interpretations must always be considered from this viewpoint as well. Finally, the model theoretic error of generalization is eliminated by predicting alterations in clock behavior rather than by the error of inappropriate generalization. REFERENCES Dingle, H. The Special Theory of Relativity, Methuem & Co., London, 1950. Herrmann, R. A. “The NSP-world and action-at-a-distance,” In “Instantaneous Action at a Distance in Modern Physics: ‘Pro’ and ‘Contra,’ ” Edited by Chubykalo, A., N. V. Pope and R. Smirnov-Rueda, (In CONTEMPORARY FUNDAMENTAL PHYSICS - V. V. Dvoeglazov (Ed.)) Nova Science Books and Journals, New York, 1999, pp. 223-235. Also see http://arxiv.org/abs/math/9903082 pp. 118-119. Kulczycki, S. Non-Euclidean Geometry, Pergamon Press, New York, 1961.
warning/0005/hep-th0005217.html	ar5iv	text	# Localized (Super)Gravity and Cosmological Constant ## I Introduction In the Brane World scenario the Standard Model gauge and matter fields are assumed to be localized on $`p9`$ spatial dimensional branes (or an intersection thereof), while gravity lives in a larger (10 or 11) dimensional bulk of space-time . The volume of dimensions transverse to the branes is then assumed to be finite. This is automatically achieved if the transverse dimensions are compact. However, as was pointed out in , if one considers certain warped compactifications, then the volume of the transverse dimensions can be finite even if the latter are non-compact. Another way of phrasing this result is that gravity is localized on a lower dimensional subspace (that is, the brane)The idea to use warped compactifications to achieve localization of gravity was originally presented in . In a concrete realization of this idea was given in the context of one extra dimension. In this was generalized to cases with more extra dimensions. In possible string theory embeddings of such a brane world scenario were discussed. For a general discussion of localization of various fields in such backgrounds, see .. One motivation for considering such unconventional compactifications is the moduli problem. In particular, the extra dimensions in such scenarios are non-compact while their volume is finite and fixed in terms of other parameters in the theory such as those in the scalar potential. That is, the scalars descending from the components of the higher dimensional metric corresponding to the extra dimensions are actually massive, and their expectation values are fixed. Also, scenarios with localized gravity could possibly have implications for the cosmological constant problem. In particular, in this context together with the localized graviton zero mode one has a continuum of massive bulk graviton modes. One can then ask whether this can somehow help avoid the usual four-dimensional field theory arguments why having vanishing (or small) cosmological constant in the absence of supersymmetry generically seems to require gross fine-tuning. Thus, on one hand, a priori it is not completely obvious whether the four-dimensional effective field theory arguments do always apply in this context. On the other hand, the fact that the volume of extra dimensions is finite seems to suggest that they should, at least at low enough energies. One of the purposes of this paper is to get additional insights into these issues. In this paper we consider localization of gravity arising in $`(D1)`$-dimensional domain wall solutions in the system of $`D`$-dimensional Einstein gravity coupled to a single real scalar field with a generic scalar potential. In particular, we discuss conditions on the scalar potential such that the corresponding domain wall solutions are non-singular (in the sense that singularities do not arise at finite values of the coordinate transverse to the domain wall). The usual kink type of solutions are non-singular as they interpolate between two adjacent local AdS minima of the scalar potential. Such solutions always have vanishing $`(D1)`$-dimensional cosmological constant. On the other hand, we point out that there exist other non-singular solutions (subject to the aforementioned non-singularity conditions on the scalar potential) which do not interpolate between AdS minima. In fact, such solutions exist even for potentials which have no minima at all and are unbounded below. Domain walls of this type have infinite tension. The $`(D1)`$-dimensional cosmological constant in such solutions can be vanishing or negative. However, positive cosmological constant is allowed for singular solutions. In fact, we discuss non-trivial conditions for physically allowed singularities, which arise from the requirement that truncating the space at the singularities be consistent. In particular, as we point out in the following, these conditions are not satisfied in some of the recently discussed “self-tuning” solutions. One possible implication of the existence of non-singular solutions with infinite tension is that if the corresponding potentials can be obtained from, say, $`D=5`$ $`𝒩=2`$ gauged supergravity, then one would obtain a supersymmetric domain wall with localized supergravity. In particular, a priori it is unclear whether such potentials avoid recent “no-go” theorems, which seem to imply that potentials with more then one AdS minima cannot be obtained in this context thus ruling out supersymmetric kink type of solutions. If not, then it would be important to show that the aforementioned “no-go” theorems also extend to infinite tension domain walls. In this paper we also study the issue of stability of domain walls which localize gravity. In particular, a priori it is not obvious why, say, infinite tension domain walls, arising in theories with potentials which have no minima at all and are unbounded below, should be stable. Indeed, the $`D`$-dimensional theory in such cases appears to be sick, and could even have tachyons. Nonetheless, we argue that (non-singular) $`(D1)`$-dimensional solutions are (classically) stable even for non-vanishing $`(D1)`$-dimensional cosmological constant (in particular, in the latter case the corresponding domain walls are not BPS saturated). The basic reason for stability is that in the case of smooth domain walls (that is, those without any ad hoc $`\delta `$-function “brane” sources which break $`D`$-dimensional diffeomorphism invariance explicitly), localization of gravity is a Higgs mechanism for the graviton field in the process of which the scalar field is eaten (or, more precisely, its corresponding modes are) by the would-be massless $`(D1)`$-dimensional graviphoton which acquires non-zero mass. In particular, we point out the importance of the full $`D`$-dimensional diffeomorphism invariance for the self-consistency of this Higgs mechanism. In the light of the above discussion it is appropriate to mention that discontinuous domain walls arising in the presence of $`\delta `$-function “brane” sources cannot be though of as limits of smooth solutions. One way to understand this is to note that such “brane” sources break translational invariance explicitly, while in the case of smooth domain walls this breakdown is spontaneous. In particular, there is an explicit discontinuity between these two different setups unless the $`(D1)`$-dimensional Planck scale is vanishing. At the end of the paper we discuss some aspects of the cosmological constant problem in the context of a brane world realized via a domain wall which localizes gravity. In particular, we point out that higher derivative such as higher curvature terms do not seem to be under control unless we consider the following type of domain walls. Consider a potential with no local minima such that it approaches finite negative values for infinite values of the scalar field. In this case we argue that the higher derivative terms are under control (as long as the AdS curvature at infinity is small enough compared with the cut-off scale for the higher derivative terms). Then provided that quantum corrections do not modify the aforementioned behavior of the scalar potential (that is, if the scalar potential still approaches constant negative values at infinity), one might hope to solve the cosmological constant problem in this context as the $`(D1)`$-dimensional cosmological constant in such solutions is ultimately zero (albeit it is not completely clear how to control the local corrections which could modify the behavior of the scalar potential at infinity). Even though the above setup might appear promising in the context of the cosmological constant problem, we point out certain difficulties with such a scenario, which, in fact, seem to be generic to theories with finite volume non-compact extra dimensions. In particular, even though higher curvature terms appear to be under control as far as preserving the desirable properties of the corresponding warped background is concerned, they generically delocalize gravity. That is, once we include such terms, generically there is no longer a normalizable graviton zero mode, and the usual $`(D1)`$-dimensional Newton’s law is no longer reproduced. We also point out some possibilities for avoiding these problems, which will be discussed in more detail elsewhere. ## II Setup In this section we discuss the setup within which we will discuss solutions with localized gravity. Thus, consider a single real scalar field $`\varphi `$ coupled to gravity with the following actionHere we focus on the case with one scalar field for the sake of simplicity. In particular, in this case we can absorb a (non-singular) metric $`Z(\varphi )`$ in the $`(\varphi )^2`$ term by a non-linear field redefinition. This cannot generically be done in the case of multiple scalar fields $`\varphi ^i`$, where one must therefore also consider the metric $`Z_{ij}(\varphi )`$.: $$S=M_P^{D2}d^Dx\sqrt{G}\left[R\frac{4}{D2}(\varphi )^2V(\varphi )\right],$$ (1) where $`M_P`$ is the $`D`$-dimensional (reduced) Planck mass. The equations of motion read: $`{\displaystyle \frac{8}{D2}}^2\varphi V_\varphi =0,`$ (2) $`R_{MN}{\displaystyle \frac{1}{2}}G_{MN}R={\displaystyle \frac{4}{D2}}\left[_M\varphi _N\varphi {\displaystyle \frac{1}{2}}G_{MN}(\varphi )^2\right]{\displaystyle \frac{1}{2}}G_{MN}V.`$ (3) The subscript $`\varphi `$ in $`V_\varphi `$ denotes derivative w.r.t. $`\varphi `$. In the following we will be interested in solutions to the above equations of motion with the warped metric of the following form: $$ds_D^2=\mathrm{exp}(2A)ds_{D1}^2+dy^2,$$ (4) where $`yx^D`$, the warp factor $`A`$, which is a function of $`y`$, is independent of the coordinates $`x^\mu `$, $`\mu =1,\mathrm{},D1`$, and the $`(D1)`$-dimensional interval is given by $$ds_{D1}^2=\stackrel{~}{g}_{\mu \nu }dx^\mu dx^\nu ,$$ (5) with the $`(D1)`$-dimensional metric $`\stackrel{~}{g}_{\mu \nu }`$ independent of $`y`$. With the above ansätz, we have: $`R_{\mu \nu }=\stackrel{~}{R}_{\mu \nu }\mathrm{exp}(2A)\left[A^{\prime \prime }+(D1)(A^{})^2\right]\stackrel{~}{g}_{\mu \nu },`$ (6) $`R_{DD}=(D1)\left[A^{\prime \prime }+(A^{})^2\right],`$ (7) $`R_{\mu D}=0,`$ (8) where prime denotes derivative w.r.t. $`y`$. Also, $`\stackrel{~}{R}_{\mu \nu }`$ is the $`(D1)`$-dimensional Ricci tensor corresponding to the metric $`\stackrel{~}{g}_{\mu \nu }`$. In the following we will be interested in solutions where $`\varphi `$ depends non-trivially on $`y`$. From the above equations it then follows that $`\varphi `$ is independent of $`x^\mu `$. The equations of motion for $`\varphi `$ and $`A`$ then become: $`{\displaystyle \frac{8}{D2}}\left[\varphi ^{\prime \prime }+(D1)A^{}\varphi ^{}\right]V_\varphi =0,`$ (9) $`(D1)(D2)(A^{})^2{\displaystyle \frac{4}{D2}}(\varphi ^{})^2+V{\displaystyle \frac{D1}{D3}}\stackrel{~}{\mathrm{\Lambda }}\mathrm{exp}(2A)=0,`$ (10) $`(D2)A^{\prime \prime }+{\displaystyle \frac{4}{D2}}(\varphi ^{})^2+{\displaystyle \frac{1}{D3}}\stackrel{~}{\mathrm{\Lambda }}\mathrm{exp}(2A)=0.`$ (11) The first equation is the dilaton equation of motion, the second equation is that for $`R_{DD}`$, and the third equation is a linear combination of the latter and the equation for $`R_{\mu \nu }`$. In fact, the equation for $`R_{\mu \nu }`$ implies that $`\stackrel{~}{\mathrm{\Lambda }}`$ is independent of $`x^\mu `$ and $`y`$. In fact, it is nothing but the cosmological constant of the $`(D1)`$-dimensional manifold, which is therefore an Einstein manifold, described by the metric $`\stackrel{~}{g}_{\mu \nu }`$. Our normalization of $`\stackrel{~}{\mathrm{\Lambda }}`$ is such that the $`(D1)`$-dimensional metric $`\stackrel{~}{g}_{\mu \nu }`$ satisfies Einstein’s equations $$\stackrel{~}{R}_{\mu \nu }\frac{1}{2}\stackrel{~}{g}_{\mu \nu }\stackrel{~}{R}=\frac{1}{2}\stackrel{~}{g}_{\mu \nu }\stackrel{~}{\mathrm{\Lambda }},$$ (12) so that the $`(D1)`$-dimensional Ricci scalar is given by $$\stackrel{~}{R}=\frac{D1}{D3}\stackrel{~}{\mathrm{\Lambda }}.$$ (13) Note that we have only two fields $`\varphi `$ and $`A`$, yet we have three equations. However, only two of these equations are independent. This can be seen as follows. Using the second equation one can express $`\varphi ^{}`$ ($`A^{}`$) via $`A^{}`$ ($`\varphi ^{}`$) and $`V`$. One can then compute $`\varphi ^{\prime \prime }`$ ($`A^{\prime \prime }`$) and plug it in the first (third) equation. This equation can then be seen to be automatically satisfied as long as the third (first) equation is satisfied. As usual, this is a consequence of Bianchi identities. ## III Solutions with $`(D1)`$-dimensional Poincaré Invariance In this section we will discuss solutions of the aforementioned equations with $`\stackrel{~}{\mathrm{\Lambda }}=0`$. In this case the equations of motion read: $`(D1)(D2)(A^{})^2{\displaystyle \frac{4}{D2}}(\varphi ^{})^2+V=0,`$ (14) $`(D2)A^{\prime \prime }+{\displaystyle \frac{4}{D2}}(\varphi ^{})^2=0.`$ (15) To study possible solutions to these equations, it is convenient to rewrite the scalar potential $`V`$ as follows: $$V=W_\varphi ^2\gamma ^2W^2,$$ (16) where $$\gamma =\frac{2\sqrt{D1}}{D2},$$ (17) and $`W`$ is some function of $`\varphi `$. Note that this is always possible (at least for a large class of scalar potentials) if we view (16) as a differential equation for $`W`$. Also note that even if $`V`$ is a simple function of $`\varphi `$, generically $`W`$ is a complicated function, which is due to the fact that the differential equation (16) is highly non-linear. The advantage of using $`W`$ instead of $`V`$ is that the equations of motion for $`\varphi `$ and $`A`$ can now be written as the first order differential equations. Thus, the above two equations are satisfied if $`\varphi `$ and $`A`$ satisfy $`\varphi ^{}=\alpha W_\varphi ,`$ (18) $`A^{}=\beta W,`$ (19) where $$\alpha ϵ\frac{\sqrt{D2}}{2},\beta ϵ\frac{2}{(D2)^{3/2}}.$$ (20) Here $`ϵ=\pm 1`$. Here one can ask whether the equations of motion imply that $`\varphi `$ and $`A`$ must satisfy the first order differential equations for some $`W`$. The answer to this question is positive. To see this, it is convenient to first rewrite the equations of motion by treating $`A`$ as a function of $`\varphi `$ (instead of $`y`$), while still treating $`\varphi `$ as a function of $`y`$. Thus, we have $`{\displaystyle \frac{4}{D2}}(\varphi ^{})^2\left[1{\displaystyle \frac{1}{4}}(D1)(D2)^2(A_\varphi )^2\right]=V,`$ (21) $`A_{\varphi \varphi }(\varphi ^{})^2={\displaystyle \frac{D2}{8}}A_\varphi V_\varphi {\displaystyle \frac{1}{D2}}V.`$ (22) Next, let $$V\left[1\frac{1}{4}(D1)(D2)^2(A_\varphi )^2\right]^1h^2,$$ (23) where $`h`$ is a function of $`\varphi `$. Then $`\varphi `$ satisfies the first order differential equation $$\varphi ^{}=\alpha h.$$ (24) Now define $`W`$ via $$V=h^2\gamma ^2W^2.$$ (25) It then follows that $`A`$ satisfies the first order differential equation $$A^{}=\eta \|\beta \|W,$$ (26) where $`\eta =\pm 1`$. Requiring that the equation for $`A^{\prime \prime }`$ (or, equivalently, that for $`A_{\varphi \varphi }`$) is satisfied then implies that $$h=ϵ\eta W_\varphi ,$$ (27) which is what we wished to show. Note that without loss of generality above we have chosen $`\eta =ϵ`$. ### A Non-singularity Conditions In this subsection we would like to discuss the conditions on $`W`$ such that the corresponding solutions do not blow up at finite values of $`y`$. More precisely, in this section we will focus on solutions such that $`\varphi `$ is non-singular<sup>§</sup><sup>§</sup>§We will refer to the corresponding domain walls as non-singular. However, as we will point out in the following, some of such solutions are actually singular in the sense that the $`D`$-dimensional Ricci scalar $`R`$ blows up, but the singularities are located at $`y=\pm \mathrm{}`$. at finite $`y`$. (We will discuss singular solutions in section V.) To begin with note that if $`V`$ is non-singular, which we will assume in the following, then $`W`$ and $`W_\varphi `$ should (generically) be non-singular as well. This then guarantees that solutions are continuous for finite values of $`\varphi `$. However, a priori it is still possible that $`\varphi `$ blows up at finite values of $`y`$. Note that the question of non-singularity of solutions can be studied without any reference to the coupling to gravity. Indeed, the equation we would like to study is $$\varphi ^{}=\alpha W_\varphi .$$ (28) This equation arises in a non-gravitational theory described by the action $$𝒮=d^Dx\left[\frac{4}{D2}(\varphi )^2𝒱(\varphi )\right].$$ (29) We can rewrite the potential $`𝒱`$ as $$𝒱=W_\varphi ^2$$ (30) if we view it as a differential equation for $`W`$. Note that this is possible as long as $`𝒱`$ is non-negative. Now let us study solutions which depend on $`y`$ only. Then the equation of motion for $`\varphi `$ is given by (28), which is nothing but the BPS equation. As should be clear from our discussion, one can write this equation for any non-negative potential $`𝒱`$. In particular, the theory need not be supersymmetric. However, a priori there is no guarantee that solutions of this equation, that is, BPS solutions, satisfying certain physical requirements exist for a given $`𝒱`$. This is precisely the question we would like to study here. It is convenient to divide the discussion of this question according to whether $`W_\varphi `$ vanishes for some value(s) of $`\varphi `$ or its non-vanishing for all real $`\varphi `$. We will first discuss the latter case. Non-vanishing $`W_\varphi `$ Suppose $`W_\varphi `$ is non-vanishing for all real $`\varphi `$. In a (globally) supersymmetric setup this would be equivalent to the statement that the theory does not have a supersymmetric vacuum. More generally, we have two possibilities. The potential $`𝒱`$ has extrema if and only if $`W_{\varphi \varphi }`$ vanishes for some values of $`\varphi `$. Some of these extrema may correspond to (local) minima (if $`𝒱_{\varphi \varphi }=2W_\varphi W_{\varphi \varphi \varphi }>0`$ at the corresponding values of $`\varphi `$). On the other hand, there might be cases where the potential $`𝒱`$ has no extrema at all - this happens if not only $`W_\varphi `$ but also $`W_{\varphi \varphi }`$ is non-vanishing. In this case we have a runaway type of potential. At any rate, for our discussion here it will have little relevance whether the potential has local minima or not. As we will point out, non-singular BPS solutions exist in either case provided that a certain simple criterion is satisfied. It might be unfamiliar, and even sound surprising, that potentials with, say, runaway behavior can have BPS solutions. An example of this type was originally discussed in . Before we give our general discussion here, we would like to give an even simpler example. Thus consider a theory with $$W=\rho \varphi ,$$ (31) where $`\rho `$ is a parameter. Note that this is a theory with a constant potential $`𝒱=\rho ^2`$. In the globally supersymmetric context supersymmetry is completely broken in $`D`$ dimensions. However, there are BPS solutions with only half of the supersymmetries broken. These are given by $$\varphi (y)=\alpha \rho y.$$ (32) Thus, even though supersymmetry is broken if we preserve $`D`$-dimensional Lorentz invariance, we can preserve half of the supersymmetries at the cost of breaking Lorentz invariance, more concretely, to that of a $`(D1)`$-dimensional theory. If we view this (delocalized) solution as a domain wall, then its tension is infinite. In fact, so is the corresponding central charge in the superalgebra. Thus, for instance, in the context of $`𝒩=1`$ supersymmetry in $`D=4`$ (in which case we really have one complex scalar filed in the corresponding chiral supermultiplet) there is a central extension of the $`𝒩=1`$ superalgebra with an infinite central charge: $$\{Q_\alpha ,Q_\beta \}=\mathrm{\Sigma }_{\alpha \beta }𝒵,$$ (33) where $`\mathrm{\Sigma }_{\alpha \beta }`$ is proportional to the area tensor in the plane perpendicular to the $`y`$ direction, and $$𝒵=2[W(y=L)W(y=L)]=(\mathrm{const}.)\times L\mathrm{}.$$ (34) Thus, solutions of this type can be viewed as domain walls with infinite tension/central charge. As we will see in the following, gravitational generalizations of such domain walls, if they are non-singular, automatically localize gravity. Next, let us discuss the general condition for such domain walls to be non-singular. That is, we would like to find the condition under which $`\varphi `$ does not blow up at finite values of $`y`$. For this to be the case, it is necessary and sufficient that the function $$F(\varphi )\frac{d\varphi }{W_\varphi }$$ (35) is unbounded at $`\varphi \pm \mathrm{}`$. Indeed, let $`\varphi (y_0)=\varphi _0`$ be finite for some finite point $`y_0`$. Then $$_{\varphi _0}^{\mathrm{}}\frac{d\varphi }{W_\varphi }=\alpha (y_{\mathrm{}}y_0),$$ (36) where $`y_{\mathrm{}}`$ is the point where $`\varphi `$ becomes $`\mathrm{}`$. If the above condition is satisfied, then the above integral is divergent, and $`y_{\mathrm{}}`$ is either $`+\mathrm{}`$ or $`\mathrm{}`$ (depending on the form of $`W`$ and sign of $`\alpha `$). Similar considerations apply to $`y_{\mathrm{}}`$ where $`\varphi `$ becomes $`\mathrm{}`$. The above condition can be reformulated directly in terms of $`W`$: $`W`$ should notMore precisely, this is correct up to usual “logarithmic” factors (that is, $`\mathrm{log}(\varphi )`$, $`\mathrm{log}(\mathrm{log}(\varphi ))`$, etc., or, more generally, the appropriate products thereof). Thus, for instance, the non-singularity condition on (35) is satisfied for $`W=\xi \varphi ^2\mathrm{log}(\varphi )`$. grow faster than $`\varphi ^2`$ for $`\varphi \pm \mathrm{}`$ . It is clear that there exist infinite choices of $`W`$ which satisfy this condition. Let us discuss general properties of such domain walls assuming that the above non-singularity condition is satisfied. Note that $`W_\varphi `$ never changes sign as $`W_\varphi `$ is non-vanishing. This then implies that $`\varphi `$ monotonically increases/decreases from $`\mathrm{}/+\mathrm{}`$ to $`+\mathrm{}/\mathrm{}`$. That is, non-singular domain walls of this type have infinite tension (and, equivalently, infinite central charge in the supersymmetric context). Here the following remark is in order. Consider the cases where $`W_\varphi `$ is non-vanishing at finite $`\varphi `$ but goes to zero for $`\varphi +\mathrm{}`$ and/or $`\varphi \mathrm{}`$. If $`W`$ blows up at large $`\varphi `$ (but more slowly than $`\varphi `$, so that $`W_\varphi `$ vanishes there), then the tension of the domain wall is infinite. However, if $`W`$ approaches finite values at $`\varphi \pm \mathrm{}`$ (which implies that $`W_\varphi `$ vanishes at $`\varphi \pm \mathrm{}`$), then the tension of the domain wall is finite - it is proportional to $`W(+\mathrm{})W(\mathrm{})`$, which is finite in this case. The situation here is similar to that where $`W_\varphi `$ vanishes at finite values of $`\varphi `$. $`W_\varphi `$ Vanishing at One Point Let us now discuss the case where $`W_\varphi `$ vanishes for only one value of $`\varphi =\varphi _0`$. This corresponds to a global minimum of the potential. (More precisely, there might also be a runaway branch with vanishing $`𝒱`$ in the large $`\|\varphi \|`$ limit.) In such cases we can also have non-singular domain walls. First, let us study the $`y`$ dependence of $`\varphi `$ near $`\varphi _0`$. We have $$W_\varphi =W_{\varphi \varphi }(\varphi _0)(\varphi \varphi _0)+\frac{1}{2}W_{\varphi \varphi \varphi }(\varphi _0)(\varphi \varphi _0)^2+\mathrm{}.$$ (37) Let us first assume that $`W_{\varphi \varphi }(\varphi _0)0`$. Then the leading $`y`$ dependence of $`\varphi `$ near $`\varphi _0`$ is given by $$\varphi (y)\varphi _0C\mathrm{exp}(ay),$$ (38) where $`C`$ is the integration constant, and $`a\alpha W_{\varphi \varphi }(\varphi _0)`$. For definiteness let us assume that $`a>0`$. Then $`\varphi `$ approaches $`\varphi _0`$ at $`y\mathrm{}`$. In the complete solution $`\varphi `$ will then monotonically increase/decrease to $`+\mathrm{}/\mathrm{}`$ at $`y+\mathrm{}`$ provided that the solution is non-singular. This is the case provided that $`W`$ does not grow faster than $`\varphi ^2`$ at $`\varphi +\mathrm{}/\mathrm{}`$. If this condition is satisfied only for, say, $`\varphi +\mathrm{}`$, then we only have non-singular solutions with $`\varphi \varphi _0`$ at $`y\mathrm{}`$ and $`\varphi +\mathrm{}`$ at $`y\pm \mathrm{}`$. If it is also satisfied for $`\varphi \mathrm{}`$, we then also have the corresponding solutions with $`\varphi \varphi _0`$ at $`y\mathrm{}`$ and $`\varphi \mathrm{}`$ at $`y\pm \mathrm{}`$. Cases with $`W_{\varphi \varphi }(\varphi _0)=0`$ can be discussed similarly. In fact, let us assume that the lowest (non-trivial) derivative of $`W`$ which is non-vanishing at $`\varphi _0`$ is the $`k`$-th derivative ($`k3`$). Then the leading $`y`$ dependence of $`\varphi `$ near $`\varphi _0`$ is given by $$(\varphi (y)\varphi _0)^{k2}\frac{a}{y},$$ (39) where $`a(k1)!/(k2)\alpha W^{(k)}(\varphi _0)`$. Once again, we have domain walls interpolating between $`\varphi =\varphi _0`$ at $`y\mathrm{}`$ and $`\varphi =\pm \mathrm{}`$ at $`y\pm \mathrm{}`$ provided that the corresponding non-singularity conditions are satisfied. Let us consider a simple example of such domain walls. Let $$W=\frac{1}{2}\zeta \varphi ^2.$$ (40) Then we have a single minimum at $`\varphi =0`$. Note that the non-singularity conditions are satisfied for this choice of $`W`$. The domain wall solutions are given by $$\varphi (y)=C\mathrm{exp}(\alpha \zeta y),$$ (41) where $`C`$ is an arbitrary integration constant. $`W_\varphi `$ Vanishing at Multiple Points If $`W_\varphi `$ vanishes at more then one points, call them $`\varphi _a`$, then we have familiar domain walls (with finite tension) interpolating between adjacent vacua. However, we can also have domain walls of the aforementioned type, which interpolate between $`\varphi =\mathrm{min}(\varphi _a)/\varphi =\mathrm{max}(\varphi _a)`$ and $`\varphi =\mathrm{}/\varphi =+\mathrm{}`$ provided that $`W`$ does not grow faster then $`\varphi ^2`$ at $`\varphi \mathrm{}/+\mathrm{}`$. Before we end this subsection, let us illustrate some of the points discussed above on the example of the $`\varphi ^4`$ theory. More concretely, consider the potential $$𝒱=V_0+m^2\varphi ^2+\lambda \varphi ^4,$$ (42) where $`\lambda >0`$. First consider the case $`m^2>0`$. In this case we must assume that $`V_00`$ so that $`𝒱`$ is non-negative. Then the theory has one minimum, which corresponds to $`W_\varphi =0`$ for $`V_0=0`$, but for $`V_0>0`$ $`W_\varphi `$ is non-vanishing. In neither case, however, do we have non-singular domain walls as the non-singularity conditions on $`W_\varphi `$ at $`\varphi \pm \mathrm{}`$ are not satisfied. Next, let us consider the case $`m^2<0`$. Let $$m^22\lambda v^2.$$ (43) Then we can rewrite the potential as $$𝒱=(V_0\lambda v^4)+\lambda (\varphi ^2v^2)^2,$$ (44) where we must assume that $`V_0\lambda v^4`$ so that the potential is non-negative at the minima $`\varphi =\pm v`$. These correspond to $`W_\varphi =0`$ for $`V_0=\lambda v^4`$, but for $`V_0>\lambda v^4`$ $`W_\varphi `$ is non-vanishing. It then follows that for $`V_0>\lambda v^4`$ there are no non-singular domain walls. For $`V_0=\lambda v^4`$, however, we have the well-known kink solutions interpolating between the two vacua $`\varphi =\pm v`$. Note that in this case we do not have non-singular domain walls interpolating between $`\varphi =\pm v`$ and $`\varphi \pm \mathrm{}`$ (the signs are correlated) as the non-singularity conditions are not satisfied. Thus, for the $`\varphi ^4`$ theory we have reproduced a known result that non-singular domain wall solutions exist (for $`m^2<0`$) only if we “fine-tune” $`V_0`$ in terms of $`m^2`$ and $`\lambda `$ as follows: $$V_0=\frac{m^4}{4\lambda }.$$ (45) In this particular case this is equivalent to the requirement that the (two degenerate) ground state(s) of the theory have vanishing vacuum energy. This is because in this case non-singularity conditions are not satisfied for non-singular infinite tension domain walls to exist. However, as should be clear from our previous discussions, such domain walls can exist in theories with ground state(s) with non-vanishing vacuum energy or even theories with no ground state at all. Normalizable Modes Before we end this subsection, we would like to discuss the condition for the existence of normalizable modes localized on a domain wall (in the non-gravitational context discussed in this subsection). Thus, consider small fluctuations $`\phi (x^\mu ,y)`$ around the domain wall background. The linearized action for $`\phi `$ reads (up to surface terms) $`𝒮[\phi ]`$ $`𝒮[\varphi +\phi ]𝒮[\varphi ]`$ (47) $`={\displaystyle \frac{1}{\alpha ^2}}{\displaystyle d^Dx\left\{_\mu \phi ^\mu \phi +(\phi ^{})^2+\alpha ^2\left[W_{\varphi \varphi }^2+W_\varphi W_{\varphi \varphi \varphi }\right]\phi ^2\right\}}.`$ Here $`\varphi `$ satisfies (28). Let us define $$\phi \varphi ^{}\omega .$$ (48) The action for $`\omega `$ reads: $$𝒮[\omega ]=d^{D1}x𝑑yW_\varphi ^2\left[_\mu \omega ^\mu \omega +(\omega ^{})^2\right].$$ (49) The zero mode is given by the configurations where $`\omega `$ is independent of $`y`$ $$\omega =\omega (x^\mu ).$$ (50) To have such a normalizable zero mode, the integral $$T2_{\mathrm{}}^+\mathrm{}𝑑yW_\varphi ^2=\frac{2}{\alpha }W\|_{y=\mathrm{}}^{y=+\mathrm{}},$$ (51) which is nothing but the domain wall tension, must be finite. Thus, domain walls with infinite tension do not have normalizable zero modes in this context. On the other hand, domain walls with finite tension do have normalizable zero modesAs we will see in the following, once we consider the coupling to gravity, this conclusion is modified.. Here the following remark is in order. The action (49) is actually exact, that is, the $`\omega `$ zero mode is a free field. This is due to the fact that it describes the Goldstone mode of the broken translational invariance in the $`y`$ direction. Thus, let $`\varphi (y)`$ be a solution describing a non-singular domain wall. Consider the following fluctuations around this solution: $`\varphi (y+\omega (x^\mu ))`$. Then we have $`\varphi (y+\omega (x^\mu ))=\varphi (y)+\varphi ^{}(y)\omega (x^\mu )+𝒪(\omega ^2)`$, so that $`\omega (x^\mu )`$ is precisely the linearized zero mode described above. However, if we plug $`\varphi (y+\omega (x^\mu ))`$ into the action $`𝒮[\varphi ]`$, after shifting the integration variable $`y`$ (note that integration over $`y`$ is from $`\mathrm{}`$ to $`+\mathrm{}`$ for non-singular domain walls), we obtain the action $`𝒮[\omega ]`$ given in (49) for arbitrary $`\omega (x^\mu )`$. Also, the equation of motion for $`\omega `$, which follows from the action (49), is given by $$W_\varphi ^2_\mu ^\mu \omega +\left[W_\varphi ^2\omega ^{}\right]^{}=0.$$ (52) Let $$\omega (x^\mu ,y)=\stackrel{~}{\omega }(x^\mu )\sigma (y),$$ (53) where $`\stackrel{~}{\omega }`$ satisfies the $`(D1)`$-dimensional Klein-Gordon equation: $$_\mu ^\mu \stackrel{~}{\omega }=m^2\stackrel{~}{\omega }.$$ (54) Then we have the following equation for $`\sigma `$: $$W_\varphi ^2m^2\sigma +\left[W_\varphi ^2\sigma ^{}\right]^{}=0.$$ (55) In particular, a zero mode ($`m^2=0`$) must satisfy $$\sigma ^{}=C/W_\varphi ^2,$$ (56) where $`C`$ is a constant. Above we claimed that the zero mode corresponds to taking $`C=0`$. To see that $`C0`$ does not give rise to normalizable zero modes, note that in this case $`\sigma `$ is a monotonous function of $`y`$ as can be seen from (56). Thus, it either goes to a constant at, say, $`y+\mathrm{}`$, or is unbounded. In either case, if the domain wall tension is infinite, such a solution is not normalizable. In fact, it is not difficult to show that such zero modes are not normalizable even if the domain wall tension is finite. Normalizability of other (that is, non-zero) modes can be discussed in a similar fashion. For this it is sometimes convenient to treat $`\sigma `$ as function of $`\varphi `$. We then have $$\left[\sigma _\varphi W_\varphi ^3\right]_\varphi +W_\varphi m^2\sigma =0,$$ (57) and $$\phi ^2𝑑y\phi ^2=\frac{1}{\alpha }\stackrel{~}{\omega }^2𝑑\varphi W_\varphi \sigma ^2.$$ (58) Thus, for illustrative purposes let us consider an example where $`W=\frac{1}{2}\zeta \varphi ^2`$. As we discussed above, in this case the domain wall solution is given by $$\varphi (y)=C\mathrm{exp}(\alpha \zeta y),$$ (59) where $`C`$ is an integration constant. The equation for $`\sigma `$ now reads: $$\left[\varphi ^3\sigma _\varphi \right]_\varphi +\frac{m^2}{\alpha ^2\zeta ^2}\varphi \sigma =0.$$ (60) The general solution to this equation is given by: $$\sigma =D_+\varphi ^{\rho _+}+D_{}\varphi ^\rho _{},$$ (61) where $$\rho _\pm =1\pm \sqrt{1\frac{m^2}{\alpha ^2\zeta ^2}},$$ (62) and $`D_\pm `$ are integration constants. Since the domain wall interpolates between $`\varphi =0`$ and $`\varphi =\pm \mathrm{}`$ (depending on the sign of $`C`$), it then follows that $`\phi ^2`$ is infinite for all allowed $`m^2`$, so that there are no normalizable modes. Note that the tension of this domain wall is infinite. ### B What about the Warp Factor? In this subsection we would like to discuss what happens to the warp factor once we consider the gravitational generalizations of the above domain walls. (Note that the discussion of the behavior of $`\varphi `$ is unchanged.) To answer this question, let us look at the equation for $`A`$: $$A^{}=\beta W.$$ (63) To begin with let us discuss the behavior of $`A`$ as a function of $`\varphi `$ near the points where $`W_\varphi `$ vanishes. Let us refer to such a point as $`\varphi _0`$. First, let us focus on the case where $`W_{\varphi \varphi }(\varphi _0)0`$. In this case we have $$A^{}=\beta \left[W(\varphi _0)+\frac{1}{2}W_{\varphi \varphi }(\varphi (y)\varphi _0)^2+\mathrm{}\right].$$ (64) Since in this case $`\varphi (y)\varphi _0C\mathrm{exp}(ay)`$ vanishes exponentially fast for the corresponding limit $`y+\mathrm{}`$ or $`y\mathrm{}`$, it then follows that all the contributions to $`A`$ except for that corresponding to the $`W(\varphi _0)`$ term vanish exponentially as well. That is, $`A(y)`$ goes to a constant if $`W(\varphi _0)`$ vanishes, while in the case $`W(\varphi _0)0`$ the leading behavior of $`A`$ as a function of $`y`$ is given by: $$A(y)\beta W(\varphi _0)y.$$ (65) This implies that for $`W(\varphi _0)=0`$ the factor $`\mathrm{exp}(2A)`$ in the metric goes to a constant in the corresponding limit for $`y`$, so that the volume of the compactification in this case is actually infinite. However, if $`W(\varphi _0)`$ is non-zero and has a correct sign, then $`\mathrm{exp}(2A)`$ vanishes exponentially in the corresponding limit for $`y`$. In fact, with appropriately chosen $`W`$ such that the domain wall interpolates between two adjacent vacua we can have the correct behavior at both $`y\mathrm{}`$ and $`y+\mathrm{}`$ so that the compactification volume is finite. More concretely, let $`W_\varphi `$ vanish at two points, call them $`\varphi _1`$ and $`\varphi _2`$, and be non-vanishing for $`\varphi _1<\varphi <\varphi _2`$. Then if $`W(\varphi _1)`$ and $`W(\varphi _2)`$ are non-zero with $`\mathrm{sign}(W(\varphi _1))=\mathrm{sign}(W(\varphi _2))=\mathrm{sign}(\beta )`$ (that is, $`W`$ must change sign between $`\varphi _1`$ and $`\varphi _2`$), $`A(y)`$ goes to $`\mathrm{}`$ at both $`y\pm \mathrm{}`$. Let us illustrate this with a simple example. Consider the theory with $$W=\xi \left[\varphi \frac{1}{3}\varphi ^3\right].$$ (66) The solution for $`\varphi `$ is given by ($`y_0`$ is the integration constant corresponding to the center of the domain wall) $$\varphi (y)=\mathrm{tanh}(\alpha \xi (yy_0)),$$ (67) while the solution for $`A`$ reads ($`C`$ is an integration constant) $$A(y)=\frac{2\beta }{3\alpha }\left[\mathrm{ln}(\mathrm{cosh}(\alpha \xi (yy_0)))\frac{1}{4}\frac{1}{\mathrm{cosh}^2(\alpha \xi (yy_0))}\right]+C.$$ (68) We therefore have $$\mathrm{exp}(2A)=\mathrm{exp}(2C)\left[\mathrm{cosh}(\alpha \xi (yy_0))\right]^{4\beta /3\alpha }\mathrm{exp}\left[\frac{\beta }{3\alpha }\frac{1}{\mathrm{cosh}^2(\alpha \xi (yy_0))}\right].$$ (69) Since $`\alpha `$ and $`\beta `$ have opposite signs, it then follows that $`\mathrm{exp}(2A)`$ falls off exponentially at large $`y`$, and the compactification volume is indeed finite. The cases where not only $`W_\varphi `$ but also higher derivatives vanish at $`\varphi _0`$ can be discussed similarly. Suppose the lowest non-trivial derivative of $`W`$ which is non-vanishing is the $`k`$-th derivative ($`k3`$). Then we have $$W=W(\varphi _0)+\frac{1}{k!}W^{(k)}(\varphi _0)(\varphi \varphi _0)^k+\mathrm{}.$$ (70) Since in this case $`(\varphi (y)\varphi _0)^{k2}a/y`$, then it follows that in the corresponding limit for $`y`$ the contributions to $`A`$ from all the terms except the $`W(\varphi _0)`$ term vanish. For instance, the first subleading term proportional to $`W^{(k)}(\varphi _0)`$ goes to zero as $`1/y^{2/(k2)}`$. Thus, as in the previous case, here we have a finite volume compactification only if $`W(\varphi _0)0`$. Let us now discuss the behavior of $`A`$ in the cases where in the corresponding limit for $`y`$ $`\varphi `$ goes to $`\pm \mathrm{}`$. (This is the case of domain walls with infinite tension.) Let us assume that for large $`\varphi `$ $`W`$ has the following behavior: $$W\xi \varphi ^\nu ,$$ (71) where $`\nu 2`$ so that the non-singularity conditions are satisfied. Moreover, first let us consider the case where $`\nu >0`$ as the analog of the $`\nu <0`$ case requires a separate treatment. Let us discuss the $`\nu =2`$ and $`\nu <2`$ cases separately. For $`\nu =2`$ we have $`W_\varphi 2\xi \varphi `$, and $$\varphi (y)C\mathrm{exp}(2\alpha \xi y),$$ (72) where C is the integration constant. In this case we therefore have the following leading behavior for $`A`$: $$A(y)\frac{\beta }{4\alpha }C^2\mathrm{exp}(4\alpha \xi y)\frac{\beta }{4\alpha }\varphi ^2,$$ (73) so that $`A(y)`$ goes to $`\mathrm{}`$ exponentially fast (recall that $`\alpha `$ and $`\beta `$ have opposite signs), and $`\mathrm{exp}(2A)`$ goes to zero very fast, more concretely, as an exponential of an exponential. Next, let us discuss $`0<\nu <2`$ cases. We then have $`W_\varphi \nu \xi \varphi ^{\nu 1}`$, and $$\varphi (y)\left[\alpha \nu (2\nu )\xi \right]^{\frac{1}{2\nu }}y^{\frac{1}{2\nu }}.$$ (74) In this case we therefore have the following leading behavior for $`A`$: $$A(y)\frac{\beta }{2\nu \alpha }\left[\alpha \nu (2\nu )\xi \right]^{\frac{2}{2\nu }}y^{\frac{2}{2\nu }}\frac{\beta }{2\nu \alpha }\varphi ^2,$$ (75) so for $`0<\nu <2`$ $`A(y)`$ goes to $`\mathrm{}`$ as a power of $`y`$. Let $$u\left\|\frac{\beta }{\nu \alpha }\left[\alpha \nu (2\nu )\xi \right]^{\frac{2}{2\nu }}\right\|.$$ (76) Then we have $$\mathrm{exp}(2A)\mathrm{exp}(u\|y\|^{\frac{2}{2\nu }}).$$ (77) Thus, $`\mathrm{exp}(2A)`$ rapidly decays to zero. In particular, for $`\nu =1`$ we have a Gaussian fall-off for $`\mathrm{exp}(2A)`$, while $`\varphi `$ is linear in $`y`$ for large $`y`$. So, the compactification volume is finite as long as $`0<\nu 2`$. Thus, as we see, such infinite tension domain walls localize gravity even more efficiently than usual domain walls interpolating between adjacent minima. In the latter case the warp factor $`\mathrm{exp}(2A)`$ decays as $$\mathrm{exp}(2A)\mathrm{exp}(\stackrel{~}{u}\|y\|),$$ (78) where $`\stackrel{~}{u}\|2\beta W(\varphi _0)\|`$. On the other hand, for $`0<\nu <2`$ the exponent of $`\|y\|`$ in (77) is larger than 1, while for $`\nu =2`$ $`\mathrm{exp}(2A)`$ decays even more rapidly, namely, as an exponential of an exponential. Next, let us assume that for large $`\varphi `$ $`W`$ has the following behavior: $$W\zeta +\xi \varphi ^\sigma .$$ (79) where $`\sigma >0`$. Here we have written the leading constant term $`\zeta `$ which $`W`$ approaches in the infinite $`\varphi `$ limit, and also the next-to-leading term which vanishes in this limit. We then have $`W_\varphi \sigma \xi \varphi ^{\sigma 1}`$, and $$\varphi (y)\left[\alpha \sigma (2+\sigma )\xi \right]^{\frac{1}{2+\sigma }}y^{\frac{1}{2+\sigma }}.$$ (80) In this case we therefore have the following leading behavior for $`A`$: $$A(y)\beta \zeta y\frac{\beta }{2\sigma \alpha }\left[\alpha \sigma (2+\sigma )\xi y\right]^{\frac{2}{2+\sigma }},$$ (81) so that for $`\zeta 0`$ the leading behavior is $$A(y)\beta \zeta y,$$ (82) while for $`\zeta =0`$ we have $$A(y)\frac{\beta }{2\sigma \alpha }\varphi ^2.$$ (83) Since the signs of $`\alpha `$ and $`\beta `$ are opposite, in the $`\zeta =0`$ case $`A`$ goes to $`+\mathrm{}`$ (instead of $`\mathrm{}`$), so that $`\mathrm{exp}(2A)`$ blows up, and the compactification volume is infinite. However, for non-zero $`\zeta `$ with the appropriate sign $`A`$ goes to $`\mathrm{}`$, so that $`\mathrm{exp}(2A)`$ decays exponentially, and the compactification volume is finite (provided that $`\mathrm{exp}(2A)`$ goes to zero fast enough in the opposite limit for $`y`$ as well). Here the situation is analogous to what we have found for the behavior of $`A`$ near the points where $`W_\varphi `$ vanishes. In particular, to have decaying $`\mathrm{exp}(2A)`$ we need a constant term $`\zeta `$ in $`W`$ in the corresponding limit. Before we end this subsection, let us discuss the analog of the $`\nu =0`$ case above. Thus, let us assume that for large $`\varphi `$ $`W`$ has the following behavior: $$W\xi \mathrm{ln}(\varphi ).$$ (84) In this case we have $`W_\varphi \xi /\varphi `$, and $$\varphi (y)(2\alpha \xi )^{\frac{1}{2}}y^{\frac{1}{2}}.$$ (85) For definiteness let as assume that $`\alpha \xi >0`$, so that we are dealing with the $`y+\mathrm{}`$ limit here. The warp factor then has the following leading behavior: $$A(y)\frac{1}{2}\beta \xi y\mathrm{ln}(y).$$ (86) Note that $`\beta \xi <0`$ (which follows from the fact that $`\alpha \xi >0`$ and $`\alpha \beta <0`$), which implies that in the $`y+\mathrm{}`$ limit $`A(y)`$ goes to $`\mathrm{}`$, and the compactification volume is finite (once again, provided that $`\mathrm{exp}(2A)`$ goes to zero fast enough in the $`y\mathrm{}`$ limit). ## IV Killing Spinors and Supersymmetry In this section we would like to discuss the aforementioned domain walls in the supersymmetric context, where their BPS property is expected to persist to the quantum level. In particular, such domain walls, as usual, are expected to be stable. The domain walls we study in this paper preserve one half of the original supersymmetries in $`D`$ dimensions. More precisely, the domain walls preserve one half of the supersymmetries of a supersymmetric minimum of the scalar potential $`V`$ if such minima are present. However, one of the key points is that domain walls have the same number of supersymmetries (which act trivially on the domain wall) even if $`V`$ has no supersymmetric minima at all. In this sense such domain walls are similar to those in the previous case. To see the aforementioned properties of domain walls explicitly, let us study the corresponding Killing spinor equations, which (up to equivalent representations) read: $`\left[\mathrm{\Gamma }^M_M\varphi \alpha W_\varphi \right]\epsilon =0,`$ (87) $`\left[𝒟_M{\displaystyle \frac{1}{2}}\beta W\mathrm{\Gamma }_M\right]\epsilon =0.`$ (88) Here $`\epsilon `$ is a Killing spinor, $`\mathrm{\Gamma }_M`$ are $`D`$-dimensional Dirac gamma matrices satisfying $$\{\mathrm{\Gamma }_M,\mathrm{\Gamma }_N\}=2G_{MN},$$ (89) and $`𝒟_M`$ is the covariant derivative $$𝒟_M_M+\frac{1}{4}\mathrm{\Gamma }_{AB}\omega _{}^{AB}{}_{M}{}^{}.$$ (90) The spin connection $`\omega _{}^{AB}{}_{M}{}^{}`$ is defined via the vielbeins $`e_{}^{A}{}_{M}{}^{}`$ in the usual way (here the capital Latin indices $`A,B,\mathrm{}=1,\mathrm{},D`$ are lowered and raised with the $`D`$-dimensional Minkowski metric $`\eta _{AB}`$ and its inverse, while the capital Latin indices $`M,N,\mathrm{}=1,\mathrm{},D`$ are lowered and raised with the $`D`$-dimensional metric $`G_{MN}`$ and its inverse). Furthermore, $$\mathrm{\Gamma }_{AB}\frac{1}{2}[\mathrm{\Gamma }_A,\mathrm{\Gamma }_B],$$ (91) where $`\mathrm{\Gamma }_A`$ are the constant Dirac gamma matrices satisfying $$\{\mathrm{\Gamma }_A,\mathrm{\Gamma }_B\}=2\eta _{AB}.$$ (92) Finally, $`W`$, which determines the scalar potential via (16), is interpreted as the superpotential in this context. Note that (87) comes from the requirement that the variation of the superpartner $`\lambda `$ of the scalar field $`\varphi `$ vanish. On the other hand, (88) arises from the requirement that the gravitino $`\psi _M`$ have vanishing variations under the corresponding supersymmetry transformations. Next, we would like to study the above Killing spinor equations in the warped backgrounds of the form (4) with the flat $`(D1)`$-dimensional metric $`\stackrel{~}{g}_{\mu \nu }=\eta _{\mu \nu }`$ (and the scalar field depending only on $`y`$): $`\left[\mathrm{\Gamma }_D\varphi ^{}\alpha W_\varphi \right]\epsilon =0,`$ (93) $`\epsilon ^{}{\displaystyle \frac{1}{2}}\beta W\mathrm{\Gamma }_D\epsilon =0,`$ (94) $`_\mu \epsilon +{\displaystyle \frac{1}{2}}\mathrm{exp}(A)\stackrel{~}{\mathrm{\Gamma }}_\mu \left[A^{}\mathrm{\Gamma }_D\beta W\right]\epsilon =0.`$ (95) Here $`\stackrel{~}{\mathrm{\Gamma }}_\mu `$ are the constant $`(D1)`$-dimensional Dirac gamma matrices satisfying $$\{\stackrel{~}{\mathrm{\Gamma }}_\mu ,\stackrel{~}{\mathrm{\Gamma }}_\nu \}=2\stackrel{~}{g}_{\mu \nu }=2\eta _{\mu \nu }.$$ (96) Also, note that $`\mathrm{\Gamma }_D`$, which is the $`D`$-dimensional Dirac gamma matrix $`\mathrm{\Gamma }_M`$ with $`M=D`$ (that is, the Dirac gamma matrix corresponding to the $`x^D=y`$ direction) is constant in this background. Note that (93), which is an algebraic equation for $`\epsilon `$, comes from (87), while the two differential equations (94) and (95) come from (88). These last two equations correspond to the coupling to gravity, while (93) appears already in a non-gravitational system with the scalar potential $`𝒱=W_\varphi ^2`$. Since on the solution we have $`\varphi ^{}=\alpha W_\varphi `$ and $`A^{}=\beta W`$, the Killing spinor satisfying (93), (94) and (95) has positive helicity w.r.t. $`\mathrm{\Gamma }_D`$, and is given by $$\epsilon _+=\mathrm{exp}(A/2)\epsilon _+^{(0)},$$ (97) where $`\epsilon _+^{(0)}`$ is a constant spinor of positive helicity. Actually, the above conclusion that we only have Killing spinors of positive helicity is only correct if $`W`$ has non-trivial $`\varphi `$ dependence. Thus, consider the case where $`W`$ is constant. It is clear that (93) is then satisfied for both positive as well as negative helicity spinors (as long as $`\varphi `$ is constant). As to (94) and (95), the general solution to this set of equations for constant $`W`$ is given by $$\epsilon =\left[\beta W\mathrm{exp}(A/2)x^\mu \stackrel{~}{\mathrm{\Gamma }}_\mu +\mathrm{exp}(A/2)\right]\epsilon _{}^{(0)}+\mathrm{exp}(A/2)\epsilon _+^{(0)},$$ (98) where $`\epsilon _{}^{(0)}`$ and $`\epsilon _+^{(0)}`$ are constant negative respectively positive helicity spinors. Thus, for constant $`W`$, that is, in the AdS space (for constant $`W`$ we have constant $`V=\gamma ^2W^2<0`$) we have Killing spinors of both positive and negative helicity. However, if we consider a domain wall for a non-constant superpotential $`W`$, then we have only the positive helicity Killing spinor (97). Thus, such domain walls are BPS objects preserving one half of the supersymmetries. The BPS property of such a domain wall implies that the modes on the wall have boson-fermion degeneracy corresponding to the unbroken supersymmetries<sup>\</sup><sup>\</sup>\Actually, the state corresponding to the domain wall itself does not have a fermionic superpartner as the unbroken supersymmetry acts trivially on it.. Thus, for instance, we expect to find a normalizable gravitino zero mode localized on the wall. In fact, this zero mode is given by the variation of $`\psi _\mu `$ corresponding to the supersymmetry transformation parameter $`ϵ`$ for which the variations of $`\lambda `$ and $`\psi _D`$ vanish. Such a spinor is given by $$\epsilon =\mathrm{exp}(A/2)\chi _+,$$ (99) where $`\chi _+=\chi _+(x^\mu )`$ is an arbitrary spinor of positive helicity (w.r.t. $`\mathrm{\Gamma }_D`$) which depends on $`x^\mu `$ but is independent of $`y`$. The corresponding gravitino variation is given by $$\delta \psi _\mu =\mathrm{exp}(A/2)_\mu \chi _+,$$ (100) and it is not difficult to check that $`\delta \psi _\mu `$ defined this way has precisely the correct normalization (in terms of powers of the warp factor) for it to be the superpartner of the four dimensional graviton. In fact, both the Einstein and the Rarita-Schwinger terms in the $`(D1)`$-dimensional action are proportional to $$M_P^{D2}_{\mathrm{}}^+\mathrm{}𝑑y\mathrm{exp}\left[(D3)A\right]\stackrel{~}{M}_P^{D3},$$ (101) where $`M_P`$ is the $`D`$-dimensional Planck scale, while $`\stackrel{~}{M}_P`$ is the $`(D1)`$-dimensional Planck scale. Thus, supersymmetric BPS domain walls discussed in this paper localize supergravity. Before we end this section, we would like to make the following comment. As we have already mentioned, in the supersymmetric context BPS domain walls are expected to be stable. This is the case not only for the domain walls interpolating between two supersymmetric AdS minima, but also for the infinite tension domain walls (or, more precisely, their gravitational counterparts). That is, the $`D`$-dimensional scalar potential $`V`$ might have no minima at all, yet the corresponding BPS domain walls are supersymmetric and stable. In particular, in this context there cannot be any tachyonic modes present, which is guaranteed by the unbroken supersymmetries. Finally, the following remark is in order. The above discussion applies only to the cases where we can sypersymmetrize the system of the scalar field $`\varphi `$ coupled to gravity. A priori such a supersymmetrization might not exist for a given potential $`V(\varphi )`$. In particular, it might not be possible to couple the aforementioned bosonic system to fermions in a fashion consistent with the corresponding superalgebra. However, the question of whether a given potential can be obtained from a supergravity theory is outside of the scope of this paper. ## V Solutions with $`\stackrel{~}{\mathrm{\Lambda }}0`$ In this section we will discuss solutions with non-zero $`(D1)`$-dimensional cosmological constant $`\stackrel{~}{\mathrm{\Lambda }}`$. Thus, let us look at the equations for $`A`$ and $`\varphi `$, namely, (10) and (11). In particular, let us rewrite (11) as follows: $$\frac{4}{D2}(\varphi ^{})^2\mathrm{exp}(2A)+\frac{1}{D3}\stackrel{~}{\mathrm{\Lambda }}=(D2)A^{\prime \prime }\mathrm{exp}(2A).$$ (102) To begin with, let us discuss non-singular domain walls. To have a non-singular domain wall with finite $`\stackrel{~}{M}_P`$, it is necessary that $`\mathrm{exp}(A)`$ asymptotically goes to zero. It then follows that $`A^{\prime \prime }\mathrm{exp}(2A)`$ (as well as, say, $`(A^{})^2\mathrm{exp}(2A)`$) asymptotically goes to zero as well. Then the r.h.s. of (102) asymptotically goes to zero at $`y\pm \mathrm{}`$, while the l.h.s. is strictly positive definite for $`\stackrel{~}{\mathrm{\Lambda }}>0`$. This then implies that $`\stackrel{~}{\mathrm{\Lambda }}`$ cannot be positive for non-singular domain walls that localize gravity. ### A Non-singular Domain Walls with $`\stackrel{~}{\mathrm{\Lambda }}<0`$ Here we can ask whether there exist non-singular domain walls with negative cosmological constant. The answer to this question is affirmative. To construct such a domain wall, pick $`A(y)`$ such that $`\stackrel{~}{M}_P`$ is finite. That is, at $`y\pm \mathrm{}`$ $`\mathrm{exp}[(D3)A]`$ goes to zero faster than $`1/y`$. Moreover, let $`A(y)`$ be such that $`A^{\prime \prime }`$ is non-positive everywhere. This then guarantees that the r.h.s. of the equation (which follows from (11)) $$\frac{4}{D2}(\varphi ^{})^2=\frac{1}{D3}\stackrel{~}{\mathrm{\Lambda }}\mathrm{exp}(2A)(D2)A^{\prime \prime }$$ (103) is strictly positive. We can then solve for $`\varphi (y)`$, and express $`y`$ in terms of $`\varphi `$ as $`\varphi (y)`$ is a monotonous function of $`y`$. Using (10) we can express the potential $`V`$ as a function of $`y`$, and, subsequently, as a function of $`\varphi `$. This way we can construct an infinite class of non-singular domain walls with $`\stackrel{~}{\mathrm{\Lambda }}<0`$ which localize gravity<sup>††</sup><sup>††</sup>††A similar procedure was described in . However, in the second order equations of motion for $`\varphi `$ and $`A`$ were rewritten as the first order equations. Unlike the $`\stackrel{~}{\mathrm{\Lambda }}=0`$ case, for non-zero (in particular, negative) $`\stackrel{~}{\mathrm{\Lambda }}`$ in this process one loses a class of solutions. Thus, it is not difficult to show that, for any choice of $`V`$, the first order equations of do not have non-singular solutions with $`\stackrel{~}{\mathrm{\Lambda }}<0`$ which localize gravity. More concretely, in such solutions $`A^{}`$ must change sign for some finite $`y`$ as $`A`$ goes to $`\mathrm{}`$ for $`y\pm \mathrm{}`$. However, solutions of this type are lost in the process of taking the square root while deriving a first order equation for $`A^{}`$ from (10) - see below.. Here we would like to study some properties of such domain walls. In particular, let us understand the asymptotic behavior of $`\varphi `$ and $`A`$ at large $`y`$. From (103) we obtain the following leading behavior for $`\varphi ^{}`$ at large $`y`$: $$\frac{4}{D2}(\varphi ^{})^2\frac{1}{D3}\stackrel{~}{\mathrm{\Lambda }}\mathrm{exp}(2A).$$ (104) On the other hand, from (10) we have: $$V\frac{D2}{D3}\stackrel{~}{\mathrm{\Lambda }}\mathrm{exp}(2A).$$ (105) That is, at large $`y`$ the potential $`V`$ behaves as $$V4(\varphi ^{})^2,$$ (106) so that it is unbounded below there. This, in particular, implies that such solutions with negative $`\stackrel{~}{\mathrm{\Lambda }}`$ cannot asymptotically approach AdS space (that is, they cannot be of the kink type). Let us define $`\chi (\varphi )`$ at large $`\varphi `$ via $$V(\varphi )4[\chi (\varphi )]^2.$$ (107) Then at large $`y`$ we have $$\varphi ^{}\chi (\varphi ).$$ (108) Note that, for the domain wall to be non-singular, $`\chi (\varphi )`$ must not (up to usual “logarithmic” factors - see below) grow faster than $`\varphi `$. Thus, let $`\chi (\varphi )\xi \varphi ^\nu `$, where $`0<\nu <1`$. Then we have $$\varphi (y)[(1\nu )\xi y]^{\frac{1}{1\nu }},$$ (109) and $$A(y)\frac{\nu }{1\nu }\mathrm{ln}(\xi y).$$ (110) Note that to have finite $`\stackrel{~}{M}_P`$ we must assume that $`\mathrm{exp}[(D3)A]`$ goes to zero at large $`y`$ faster than $`1/y`$. This then implies the following restriction on $`\nu `$: $$\nu >\frac{1}{D2}.$$ (111) Thus, if this condition is not satisfied, then the corresponding domain walls with negative $`\stackrel{~}{\mathrm{\Lambda }}`$ do not localize gravity. This is to be contrasted with the $`\stackrel{~}{\mathrm{\Lambda }}=0`$ case where such potentials do give rise to walls which localize gravity. Also, note that to have finite volume in the $`y`$ direction, which is given by $$\mathrm{Vol}=𝑑y\mathrm{exp}[(D1)A],$$ (112) we must assume that $`\mathrm{exp}[(D1)A]`$ goes to zero at large $`y`$ faster than $`1/y`$. This is satisfied provided that $`\nu >1/D`$. Thus, for $`1/D<\nu <1/(D2)`$ the compactification volume is finite, yet gravity is not localized. Next, assume that $`\chi (\varphi )\varphi \rho (\mathrm{ln}\|\varphi \|)`$. For the domain wall to be non-singular, the function $`\rho (u)`$ must be such that the integral $$\frac{du}{\rho (u)}$$ (113) is unbounded at large $`u`$. The leading behavior for $`A`$ is then obtained from the following equation: $$A^{}=\rho (A).$$ (114) Thus, for instance, if $`\rho (u)\zeta u^\sigma `$, where $`0\sigma <1`$, then $$A[(1\sigma )\zeta y]^{\frac{1}{1\sigma }}.$$ (115) If $`\rho (u)\zeta u`$, then $$AC\mathrm{exp}(\zeta y).$$ (116) Here $`C`$ is a negative integration constant, and $`\zeta y`$ is assumed to be positive. First Order Equations Sometimes it is convenient to rewrite the second order equations of motion for $`\varphi `$ and $`A`$ as the first order equations. However, as we have already mentioned, for non-zero (in particular, negative) $`\stackrel{~}{\mathrm{\Lambda }}`$ some extra care is needed not to lose some of the solutions. Thus, consider the case of non-singular domain walls with $`\stackrel{~}{\mathrm{\Lambda }}<0`$ such that $`A`$ goes to $`\mathrm{}`$ at $`y\pm \mathrm{}`$. Then it follows that at some finite $`y`$, call it $`y_0`$, $`A^{}`$ must change sign. It is then not difficult to show that the first order equations for $`\varphi `$ and $`A`$ are given by: $`\varphi ^{}=\alpha h,`$ (117) $`A^{}=\mathrm{sign}(yy_0)\beta \mathrm{\Omega }W,`$ (118) where $`W_\varphi =\mathrm{sign}(yy_0)\mathrm{\Omega }h,`$ (119) $`\mathrm{\Omega }=\sqrt{1+{\displaystyle \frac{D1}{D3}}{\displaystyle \frac{\stackrel{~}{\mathrm{\Lambda }}}{\gamma ^2W^2}}\mathrm{exp}(2A)},`$ (120) $`V=h^2\gamma ^2W^2.`$ (121) Here we note that $`W`$ does not vanish anywhere, and has an extremum at $`y=y_0`$. Also, $`\varphi ^{}`$ does not vanish at $`y=y_0`$, where by assumption $`A^{\prime \prime }`$ is non-positive. ### B Singular Domain Walls By a singular domain wall we mean a solution where $`\varphi `$ blows up at some finite value(s) of $`y`$. We can then have domain walls that are singular on both sides (they interpolate between $`y_{}`$ and $`y_+`$, $`y_\pm `$ being the points where $`\varphi `$ blows up), or domain walls that are singular on one side only (that is, those interpolating between $`\mathrm{}`$ and $`y_+`$, or $`y_{}`$ and $`+\mathrm{}`$). To begin with, note that if $`\varphi `$ blows up at some finite $`y`$, call it $`y_{}`$, then $`A`$ is bounded above at $`y=y_{}`$. Indeed, let us rewrite (102) as follows: $$\mathrm{exp}(2A)\left[\frac{4}{D2}(\varphi ^{})^2+(D2)A^{\prime \prime }\right]=\frac{1}{D3}\stackrel{~}{\mathrm{\Lambda }}.$$ (122) If $`A`$ goes to $`+\mathrm{}`$ at $`y=y_{}`$, then $`A^{\prime \prime }`$ also goes to $`+\mathrm{}`$, and the above equation cannot be satisfied. Thus, $`A`$ goes to either $`\mathrm{}`$ or a constant at $`y=y_{}`$. In either case $`A^{\prime \prime }`$ goes to $`\mathrm{}`$ there (which follows from (122)). This then implies that for domain walls singular on both sides the compactification volume is automatically finite if we cut off the space in the $`y`$ direction at the singularity<sup>‡‡</sup><sup>‡‡</sup>‡‡Here one might suspect that not all singularities are physical. The purpose of the next subsection is to determine the conditions for physically allowed singularities.. Similarly, for domain walls singular on one side only the compactification volume is finite as long as on the non-singular side $`\mathrm{exp}(2A)`$ goes to zero fast enough. We now wish to show that if the scalar potential $`V`$ is bounded above at the singularity, then $`A`$ necessarily goes to $`\mathrm{}`$ there (that is, it cannot go to a constant). To see this, consider the sum of (10) and (11): $$(D2)\left[A^{\prime \prime }+(D1)(A^{})^2\right]+V\frac{D2}{D3}\stackrel{~}{\mathrm{\Lambda }}\mathrm{exp}(2A)=0.$$ (123) Now suppose $`V`$ is bounded above at $`y=y_{}`$. Let us assume that $`A`$ goes to a constant at $`y=y_{}`$. It then follows from (123) that $`(A^{})^2`$ goes to $`\mathrm{}`$. In fact, $`(A^{})^2`$ goes to $`\mathrm{}`$ at least as fast as $`A^{\prime \prime }`$. However, from (122) it follows that $`A^{\prime \prime }`$ goes to $`\mathrm{}`$ as fast as $`(\varphi ^{})^2`$. We therefore conclude that $`A^{}`$ blows up at least as fast as $`\varphi ^{}`$. This, however, implies that $`A`$ cannot go to a constant at $`y=y_{}`$ as $`\varphi `$ blows up there. Here we note that if we relax the requirement that $`V`$ be bounded above at singularities, then $`A`$ could a priori go to a constant at a singularity. However, as we will point out in the next subsection, for singularities to be physical an additional non-trivial condition must be satisfied. As we will see in a moment, singular domain walls with the property that $`A`$ goes to a constant at at least one singularity do not satisfy this consistency condition. ### C Physically Allowed Singularities Let us substitute the domain wall ansätz into the action $`S`$ given in (1). We then obtain (here $`(\stackrel{~}{}\varphi )^2\stackrel{~}{g}^{\mu \nu }\stackrel{~}{}_\mu \varphi \stackrel{~}{}_\nu \varphi `$): $$S/M_P^{D2}=d^Dx\sqrt{\stackrel{~}{g}}\mathrm{exp}[(D3)A]\left[\stackrel{~}{R}\frac{1}{\alpha ^2}(\stackrel{~}{}\varphi )^2\right]d^{D1}x\sqrt{\stackrel{~}{g}}E[A,\varphi ],$$ (124) where the energy functional $`E[A,\varphi ]`$ is given by $$E[A,\varphi ]=𝑑y\mathrm{exp}[(D1)A]\left[(D1)[2A^{\prime \prime }+D(A^{})^2]+\frac{1}{\alpha ^2}(\varphi ^{})^2+V(\varphi )\right].$$ (125) Here the integral over $`y`$ is taken from $`y_{}`$ to $`y_+`$, where $`y_\pm `$ correspond to the edges of the domain wall ($`y_\pm `$ can be finite or infinite). This energy functional can be rewritten as $`E[A,\varphi ]=`$ $`{\displaystyle 𝑑y\mathrm{exp}[(D1)A]\left[\frac{1}{\alpha ^2}(\varphi ^{})^2(D1)(D2)(A^{})^2+V(\varphi )\right]}+`$ (127) $`2(D1)\left[A^{}\mathrm{exp}[(D1)A]\right]\|_y_{}^{y_+}.`$ The boundary term will become important in the following, so we will keep it. Here we note that (9) and a linear combination of (10) and (11) with $`\stackrel{~}{\mathrm{\Lambda }}=0`$ are the Euler-Lagrange equations for the energy functional (127) . Similarly, for non-zero $`\stackrel{~}{\mathrm{\Lambda }}`$ the same equations are the Euler-Lagrange equations for the functional $$\widehat{E}[A,\varphi ]E[A,\varphi ]𝑑y\frac{D1}{D3}\stackrel{~}{\mathrm{\Lambda }}\mathrm{exp}[(D3)A].$$ (128) Using the equations of motion for $`\varphi `$ and $`A`$, it is not difficult to show that on the solution we have $`E[A,\varphi ]=\stackrel{~}{\mathrm{\Lambda }}{\displaystyle 𝑑y\mathrm{exp}[(D3)A]}+2\left[A^{}\mathrm{exp}[(D1)A]\right]\|_y_{}^{y_+}.`$ (129) This then implies that on the solution, where $`\stackrel{~}{}_\mu \varphi =0`$, and $`\stackrel{~}{g}_{\mu \nu }`$ is independent of $`y`$, we have $$S=\stackrel{~}{M}_P^{D3}d^{D1}x\sqrt{\stackrel{~}{g}}\left[\stackrel{~}{R}\stackrel{~}{\mathrm{\Lambda }}\right]2M_P^{D2}\left[A^{}\mathrm{exp}[(D1)A]\right]\|_y_{}^{y_+}d^{D1}x\sqrt{\stackrel{~}{g}},$$ (130) where $$\stackrel{~}{M}_P^{D3}M_P^{D2}_y_{}^{y_+}𝑑y\mathrm{exp}[(D3)A].$$ (131) This coincides with the $`(D1)`$-dimensional action with the Planck scale $`\stackrel{~}{M}_P`$ and the cosmological constant $`\stackrel{~}{\mathrm{\Lambda }}`$ $$\stackrel{~}{S}=\stackrel{~}{M}_P^{D3}d^{D1}x\sqrt{\stackrel{~}{g}}\left[\stackrel{~}{R}\stackrel{~}{\mathrm{\Lambda }}\right]$$ (132) up to the surface term $$\frac{2}{D1}M_P^{D2}𝒜^{}\|_y_{}^{y_+}d^{D1}x\sqrt{\stackrel{~}{g}},$$ (133) where $$𝒜\mathrm{exp}[(D1)A].$$ (134) Thus, to have a consistent $`(D1)`$-dimensional interpretation, we must require that the aforementioned boundary contribution vanish. That is, we have an additional consistency condition: $$𝒜^{}\|_y_{}^{y_+}=0.$$ (135) Note that for non-singular domain walls with finite $`\stackrel{~}{M}_P`$ this condition is automatically satisfied - indeed, for such walls $`𝒜`$ and, therefore, $`𝒜^{}`$ asymptotically go to zero at $`y_\pm =\pm \mathrm{}`$. However, for singular domain walls this condition is non-trivial. In fact, in the case of singular domain walls $`𝒜^{}(y_\pm )`$ need not even be finite. For the condition (135) to be meaningful, we must then require that $$𝒜^{}(y_\pm )\text{are finite, and}𝒜^{}(y_+)𝒜^{}(y_{})=0.$$ (136) Note that if $`𝒜`$ goes to zero at singularities, then $`𝒜^{}`$ is non-positive at $`y=y_+`$, and non-negative at $`y=y_{}`$, so that the condition (135) is satisfied if and only if $$𝒜^{}(y_\pm )=0.$$ (137) As we will show in a moment, singularities where $`𝒜`$ is finite do not satisfy the condition (136). This then implies that at physically allowed singularities $`𝒜`$ as well $`𝒜^{}`$ must necessarily vanish. Potentials Bounded Above Let us first consider the case where $`V`$ is bounded above at singularities. In this case $`𝒜`$ goes to zero at $`y_\pm `$, and, therefore, $`𝒜^{}`$ must also go to zero there. If $`\stackrel{~}{\mathrm{\Lambda }}0`$ these conditions imply a certain restriction on the scalar potential. Thus, multiplying both sides of (123) by $`\mathrm{exp}[(D1)A]`$, we obtain $$V=\frac{D2}{D1}\frac{𝒜^{\prime \prime }}{𝒜}+\frac{D2}{D3}\frac{\stackrel{~}{\mathrm{\Lambda }}}{𝒜^{\frac{2}{D1}}}.$$ (138) Now, both $`𝒜`$ and $`𝒜^{}`$ must vanish at $`y_\pm `$. Then it follows that $`𝒜^{\prime \prime }`$ is non-negative at $`y_\pm `$. This then implies that for $`\stackrel{~}{\mathrm{\Lambda }}0`$ at the singularities the scalar potential must be unbounded below<sup>\</sup><sup>\</sup>\This is a much stronger condition than that recently conjectured in to be necessary for a singularity to be physical. In particular, it was argued in that (for solutions with $`\stackrel{~}{\mathrm{\Lambda }}=0`$) a singularity is physical only if $`V`$ is bounded above there.: $$V\mathrm{}\text{at}y=y_\pm \text{for}\stackrel{~}{\mathrm{\Lambda }}0.$$ (139) Note that this condition is necessary but might not always be sufficient for (137) to be satisfied. Thus, generally the full condition (137) should be used to check whether a given singularity is physically allowed. Here we note that for $`\stackrel{~}{\mathrm{\Lambda }}>0`$ the corresponding condition on $`V`$ is more detail dependent. In particular, as we will see in the following, there exist (discontinuous) domain walls with physically acceptable singularities and $`\stackrel{~}{\mathrm{\Lambda }}>0`$ such that $`V`$ is not unbounded below at singularities. Potentials Unbounded Above Let us now consider potentials unbounded above at the singularities. Let us see if we can satisfy the condition (136). First, the above discussion for singularities where $`A\mathrm{}`$ applies to this case as well. In particular, solutions with $`\stackrel{~}{\mathrm{\Lambda }}0`$ with singularities where $`A\mathrm{}`$ and $`V`$ is unbounded above are not physically allowed. Next, let us assume that $`A`$ goes to a constant at, say, $`y_+`$. Recall that $`A^{\prime \prime }`$ must go to $`\mathrm{}`$. In particular, $`(\varphi ^{})^2`$ goes to $`\mathrm{}`$ as fast as $`A^{\prime \prime }`$. Let $`A^{}N`$, and $`z1/(y_+y)`$. As $`yy_+`$ we have $`z+\mathrm{}`$. Now, $`N^{}=z^2N_z`$, where the subscript $`z`$ denotes derivative w.r.t. $`z`$. On the other hand, $`\varphi ^{}=z^2\varphi _z`$. This then implies that at $`z+\mathrm{}`$ we have (see (122)) $$\varphi _z\eta \frac{D2}{2}\frac{1}{z}\sqrt{N_z},$$ (140) where $`\eta =\pm 1`$. For $`\varphi `$ to blow up at $`z+\mathrm{}`$ (that is, at $`yy_+`$) it is then necessary that $`N_z`$ be a non-decreasing<sup>\†</sup><sup>\†</sup>\†More precisely, this is correct up to the usual “logarithmic” factors. For instance, let $`N_z1/\mathrm{ln}^2(z)`$. Then $`\varphi _z1/z\mathrm{ln}(z)`$, and $`\varphi \mathrm{ln}(\mathrm{ln}(z))`$, so that $`\varphi `$ blows up at $`z+\mathrm{}`$. However, such logarithmic factors do not affect the following discussion. In particular, in this case $`Nz/\mathrm{log}^2(z)`$, which blows up at large $`z`$. function of $`z`$. But then it follows that $`N`$ blows up at $`z+\mathrm{}`$ at least as fast as $`z`$. That is, $`A^{}`$ blows up at $`yy_+`$ at least as fast as $`1/(y_+y)`$. Thus, the condition (136) cannot be satisfied for such domain walls (recall that $`𝒜^{}=(D1)A^{}\mathrm{exp}[(D1)A]`$, so to have finite $`𝒜^{}(y_+)`$ we must have finite $`A^{}(y_+)`$ as $`A`$ goes to a constant there by assumption). Similarly, if $`A`$ goes to a constant at $`y=y_{}`$, (136) cannot be satisfied there. We therefore conclude that singularities where $`A`$ goes to a constant are unphysical (regardless of the value of $`\stackrel{~}{\mathrm{\Lambda }}`$). Physical Interpretation From the above discussion it follows that for singularities to be physical they must satisfy the condition (137). Here we would like to discuss a physical interpretation of this condition. Thus, we would like to be able to cut off the space in the $`y`$ direction at singularities. That is, singularities correspond to boundaries of the space. On the other hand, the condition (137) can be rewritten as (note that $`𝒜=\sqrt{G}`$) $$(\sqrt{G})^{}\|_{y_\pm }=0.$$ (141) That is, at a boundary the normal derivative of the volume density must vanish. This is just as well for otherwise there would be no reason why we should not continue the space in the $`y`$ direction beyond such a boundary. Note that if we have a domain wall which is singular on both sides, then topologically the space is an interval, which is the same as $`𝐒^1/𝐙_2`$. Similarly, if the domain wall is singular on one side only, then topologically the space is $`𝐑/𝐙_2`$. However, the actual compactification, say, in the former case does not correspond to that on $`𝐒^1/𝐙_2`$ in the sense that it cannot be viewed as an orbifold of an $`𝐒^1`$ compactification. In particular, in this case we do not have the usual Kaluza-Klein modes (or, more precisely, their combinations invariant under the corresponding $`𝐙_2`$ orbifold action) arising in the $`𝐒^1`$ compactifications. ### D Discontinuous Domain Walls In the previous subsections we discussed smooth domain walls. Here we would like to consider discontinuous domain walls. More concretely, both $`\varphi `$ and $`A`$ could be continuous, but need not have continuous derivatives. In particular, to have a non-singular domain wall that localizes gravity we must make sure that $`A`$ goes to $`\mathrm{}`$ at $`y\pm \mathrm{}`$. Also, in the case of singular domain walls $`A`$ must go to $`\mathrm{}`$ at singularities. This implies that $`A^{}`$ must change sign at finite $`y`$. This can happen continuously as in the previous subsections. However, a priori we also have the following possibility<sup>\‡</sup><sup>\‡</sup>\‡Here we should point out that one can also consider discontinuous domain walls where $`A^{}`$ does not change sign at the discontinuity. Furthermore, the following discussion can be straightforwardly generalized to include multiple discontinuities.. Let $`A^{}>0`$ for $`y<y_0`$, and $`A^{}<0`$ for $`y>y_0`$. At $`y=y_0`$ $`A^{}`$ is discontinuous. Then $`A^{\prime \prime }`$ behaves as a $`\delta `$-function at $`y=y_0`$. To be able to satisfy the equations of motion for $`A`$ and $`\varphi `$ we then must add a source term to the action $`S`$ in (1): $$S_{\mathrm{total}}=S+S_{\mathrm{source}},$$ (142) where $$S_{\mathrm{source}}=M_P^{D2}_\mathrm{\Sigma }d^{D1}x\sqrt{g}f(\varphi ).$$ (143) Here $`\mathrm{\Sigma }`$ is the $`y=y_0`$ hypersurface, $$g_{\mu \nu }\delta _{\mu }^{}{}_{}{}^{M}\delta _{\nu }^{}{}_{}{}^{N}G_{MN},$$ (144) and $`f(\varphi )`$ describes the coupling of $`\varphi `$ to the space-time defect (often referred to as a “brane”) corresponding to the hypersurface $`\mathrm{\Sigma }`$. The above source term (generically) leads to discontinuities in $`A^{}`$ and $`\varphi ^{}`$ at $`y=y_0`$, and the corresponding jump conditions must be imposed on the solution. The equations of motion now read: $`{\displaystyle \frac{8}{D2}}^2\varphi V_\varphi {\displaystyle \frac{\sqrt{g}}{\sqrt{G}}}f_\varphi \delta (yy_0)=0,`$ (145) $`R_{MN}{\displaystyle \frac{1}{2}}G_{MN}R={\displaystyle \frac{4}{D2}}\left[_M\varphi _N\varphi {\displaystyle \frac{1}{2}}G_{MN}(\varphi )^2\right]{\displaystyle \frac{1}{2}}G_{MN}V`$ (146) $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{\sqrt{g}}{\sqrt{G}}}\delta _{M}^{}{}_{}{}^{\mu }\delta _{N}^{}{}_{}{}^{\nu }g_{\mu \nu }f\delta (yy_0).`$ (147) As before, we are interested in solutions with the metric of the following form: $$ds_D^2=\mathrm{exp}(2A)\stackrel{~}{g}_{\mu \nu }dx^\mu dx^\nu +dy^2,$$ (148) where the warp factor $`A`$, which is a function of $`y`$, is independent of the coordinates $`x^\mu `$, and the $`(D1)`$-dimensional metric $`\stackrel{~}{g}_{\mu \nu }`$ is independent of $`y`$. Moreover, $`\varphi `$ is independent of $`x^\mu `$, but can depend non-trivially on $`y`$. With this ansätz, we have the following equations of motion for $`\varphi `$ and $`A`$: $`{\displaystyle \frac{8}{D2}}\left[\varphi ^{\prime \prime }+(D1)A^{}\varphi ^{}\right]V_\varphi f_\varphi \delta (yy_0)=0,`$ (149) $`(D1)(D2)(A^{})^2{\displaystyle \frac{4}{D2}}(\varphi ^{})^2+V{\displaystyle \frac{D1}{D3}}\stackrel{~}{\mathrm{\Lambda }}\mathrm{exp}(2A)=0,`$ (150) $`(D2)A^{\prime \prime }+{\displaystyle \frac{4}{D2}}(\varphi ^{})^2+{\displaystyle \frac{1}{D3}}\stackrel{~}{\mathrm{\Lambda }}\mathrm{exp}(2A)+{\displaystyle \frac{1}{2}}f\delta (yy_0)=0.`$ (151) The jump conditions read: $`{\displaystyle \frac{8}{D2}}\left[\varphi ^{}(y_0+)\varphi ^{}(y_0)\right]f_\varphi (\varphi _0)=0,`$ (152) $`(D2)\left[A^{}(y_0+)A^{}(y_0)\right]+{\displaystyle \frac{1}{2}}f(\varphi _0)=0,`$ (153) $`\varphi (y_0+)=\varphi (y_0)\varphi _0,`$ (154) $`A(y_0+)=A(y_0)A_0.`$ (155) From our previous discussions it should be clear that if we demand that $`A`$ goes to $`\mathrm{}`$ at $`y=y_\pm `$, then $`A^{}(y_0)0`$ and $`A^{}(y_0+)0`$, so that from the jump condition on $`A^{}`$ we have $$f(\varphi _0)0.$$ (156) Since $`M_P^{D2}f(\varphi _0)`$ is nothing but the tension of the aforementioned space-time defect (“brane”), it then follows that gravity can be localized in this context only if this tension is non-negativeNote that this need not be the case if, say, $`A^{}`$ does not change sign at the discontinuity.. It is now clear how to obtain a solution of this type with finite compactification volume. Start with two smooth domain walls, call them (1) and (2), which are solutions of the equations of motion without the source term. Let these domain walls be such that $`A^{(1)}\mathrm{}`$ as $`yy_{}^{(1)}`$, and $`A^{(2)}\mathrm{}`$ as $`yy_+^{(2)}`$. Moreover, let $`y_+^{(2)}>y_{}^{(1)}`$, and assume that there is a point $`y_0`$, $`y_{}^{(1)}<y_0<y_+^{(2)}`$, such that $`\varphi ^{(1)}(y_0)=\varphi ^{(2)}(y_0)\varphi _0`$, and $`A^{(1)}(y_0)=A^{(2)}(y_0)A_0`$ with some finite $`\varphi _0`$ and $`A_0`$. Now consider the domain wall given by $`\varphi (y)=\varphi ^{(1)}(y),y_{}^{(1)}<yy_0,`$ (157) $`\varphi (y)=\varphi ^{(2)}(y),y_0<y<y_+^{(2)},`$ (158) $`A(y)=A^{(1)}(y),y_{}^{(1)}<yy_0,`$ (159) $`A(y)=A^{(2)}(y),y_0<y<y_+^{(2)}.`$ (160) This discontinuous domain wall will satisfy the equations of motion with the source term as long as the function $`f(\varphi )`$ is such that the following jump conditions are satisfied: $`{\displaystyle \frac{8}{D2}}\left[(\varphi ^{(2)})^{}(y_0)(\varphi ^{(1)})^{}(y_0)\right]f_\varphi (\varphi _0)=0,`$ (161) $`(D2)\left[(A^{(2)})^{}(y_0)(A^{(1)})^{}(y_0)\right]+{\displaystyle \frac{1}{2}}f(\varphi _0)=0.`$ (162) Note that from our discussion at the beginning of this section it follows that discontinuous domain walls with finite $`\stackrel{~}{M}_P`$ which are non-singular or singular on one side only cannot have positive cosmological constant. On the other hand, discontinuous domain walls with finite $`\stackrel{~}{M}_P`$ which are singular on both sides can a priori have either vanishing or non-vanishing cosmological constant of either sign. Here we can repeat the analysis of the previous subsection for discontinuous domain walls. It is then not difficult to show that $`S_{\mathrm{total}}`$ on the solution has the same form as (130), so that we must impose the consistency condition (135). The latter is automatically satisfied for non-singular domain walls, but for singular domain walls it is non-trivial. Here we would like to give a simple example which does not satisfy this condition. Examples with Unphysical Singularities Thus, consider the case where the scalar potential vanishes: $$V(\varphi )0.$$ (163) Let us look for domain wall solutions with zero cosmological constant. Let us rewrite the scalar potential in terms of $`W`$: $$V=W_\varphi ^2\gamma ^2W^2.$$ (164) Without loss of generality we can choose $$W=\xi \mathrm{exp}(\gamma \varphi ),$$ (165) where $`\xi `$ is an arbitrary constant. The smooth domain walls with $`\stackrel{~}{\mathrm{\Lambda }}=0`$ are then the solutions to the equations $`\varphi ^{}=\alpha W_\varphi =\alpha \xi \gamma \mathrm{exp}(\gamma \varphi ),`$ (166) $`A^{}=\beta W=\beta \xi \mathrm{exp}(\gamma \varphi ).`$ (167) These solutions are given by: $`\varphi ={\displaystyle \frac{1}{\gamma }}\mathrm{ln}[\alpha \xi \gamma ^2(y_{}y)],`$ (168) $`A=C{\displaystyle \frac{\beta }{\alpha \gamma ^2}}\mathrm{ln}[\alpha \xi \gamma ^2(y_{}y)],`$ (169) where $`C`$ is an integration constant, and $`y_{}`$ corresponds to the location of the singularity. As expected, at the singularity $`A\mathrm{}`$ (recall that $`\alpha `$ and $`\beta `$ have opposite signs). Note, however, that at the other edge $`A+\mathrm{}`$, so that these smooth domain walls do not localize gravity. Now take two such domain walls satisfying the aforementioned properties, and construct a discontinuous domain wall with the edges corresponding to singularities (that is, such a domain wall is singular on both sides). This domain wall then has finite compactification volume. However, it does not satisfy the consistency condition (135), and therefore the singularities are unphysical. Indeed, we have (for definiteness we assume that $`\alpha \xi >0`$, so that $`y<y_{}`$, that is, $`y_{}`$ corresponds to $`y_+`$) $$𝒜=\mathrm{exp}[(D1)A]=\mathrm{exp}[(D1)C]\alpha \xi \gamma ^2(y_{}y),$$ (170) so that $$𝒜^{}=\mathrm{exp}[(D1)C]\alpha \xi \gamma ^2.$$ (171) This then implies that the consistency condition (136) is not satisfiedThis observation was also made by Gary Shiu. for these solutions<sup>\∥</sup><sup>\∥</sup>\∥Here we note that such discontinuous domain walls were recently discussed in .. In fact, it is not difficult to show that the same conclusion applies to discontinuous singular domain walls with non-vanishing $`\stackrel{~}{\mathrm{\Lambda }}`$ obtained by starting with $`V(\varphi )0`$ and adding a source term<sup>\</sup><sup>\</sup>\*Such domain walls were recently studied in .. ### E Smooth vs. Discontinuous Domain Walls The purpose of this subsection is to point out that a discontinuous domain wall of the aforementioned type cannot be thought of as a limit of a smooth domain wall. The physical reason for this is actually quite simple. Thus, note that in the case of a smooth domain wall with finite $`\stackrel{~}{M}_P`$ the localization widths for the scalar field $`\varphi `$ and gravity are related to each other - essentially they are both determined by the fall-off of the warp factor $`\mathrm{exp}(2A)`$. Thus, if we take the localization width of the scalar field to zero with the intension of reproducing the source term in (142), the localization width for gravity will also go to zero. Then for finite $`D`$-dimensional Planck scale $`M_P`$ the $`(D1)`$-dimensional Planck scale $`\stackrel{~}{M}_P`$ goes to zero as well. This indicates that we cannot approximate a discontinuous domain wall by a limit of a smooth domain wall. To understand this point in a bit more detail, let us consider the original Randall-Sundrum scenario , where the scalar potential is a negative constant $$V=\xi ^2.$$ (172) Let us consider solutions with vanishing cosmological constant $`\stackrel{~}{\mathrm{\Lambda }}=0`$. We have the following smooth solutions ($`\varphi _0`$ is a constant): $`\varphi =\varphi _0,`$ (173) $`A={\displaystyle \frac{\beta \xi }{\gamma }}y+\mathrm{const}.,`$ (174) which correspond to the $`D`$-dimensional AdS space-time. Now let us add the source term with $`f(\varphi )`$ independent of $`\varphi `$: $$f(\varphi )f_0.$$ (175) We then find the following discontinuous domain walls with finite $`\stackrel{~}{M}_P`$ (for definiteness we are assuming $`\beta \xi >0`$): $`\varphi (y)=\varphi _0,`$ (176) $`A(y)=A_0\beta \xi \|yy_0\|,`$ (177) with the following jump condition: $$f_0=4(D2)\beta \xi ,$$ (178) which implies that the “brane” tension must be fine-tuned to the $`D`$-dimensional cosmological constant ($`\xi =\sqrt{V}`$) for such a solution to exist. Note that this solution is precisely the Randall-Sundrum solution up to the extra spectator scalar field $`\varphi `$. We can now ask whether we can approximate this solution by a limit of a smooth solution. The starting point, as it should become clear in a moment, is not very important here as long as it correctly reproduces the features of the above solution such as the fact that away from the “brane” we have AdS vacua. Thus, for illustrative purposes, let us consider the kink solution discussed in subsection B of section III. Now, $`W`$ is given by $$W=\frac{3\xi }{2\gamma }\left[\varphi \frac{1}{3}\varphi ^3\right],$$ (179) so that in the AdS minima with $`W_\varphi =0`$ ($`\varphi =\pm 1`$) we have $`W=\pm \xi /\gamma `$, and $`V=\xi ^2`$. The solution for $`\varphi `$ and $`A`$ is given by $`\varphi (y)=\mathrm{tanh}\left[{\displaystyle \frac{3\alpha \xi }{2\gamma }}(yy_0)\right],`$ (180) $`A={\displaystyle \frac{2\beta }{3\alpha }}\left[\mathrm{ln}\left[\mathrm{cosh}\left[{\displaystyle \frac{3\alpha \xi }{2\gamma }}(yy_0)\right]\right]{\displaystyle \frac{1}{4}}{\displaystyle \frac{1}{\mathrm{cosh}^2\left[\frac{3\alpha \xi }{2\gamma }(yy_0)\right]}}\right]+C.`$ (181) It is now clear that the Randall-Sundrum solution for $`A`$ is obtained in the limit $`\xi \mathrm{}`$. (In this limit $`\varphi (y)=\mathrm{sign}(yy_0)`$.) However, in this limit the $`D`$-dimensional cosmological constant $`V=\xi ^2\mathrm{}`$, and the $`(D1)`$-dimensional Planck scale $$\stackrel{~}{M}_P^{D3}=\frac{2}{D3}\frac{1}{\beta \xi }M_P^{D2}0$$ (182) for finite $`D`$-dimensional Planck scale $`M_P`$. Note that even if we attempt to take $`M_P`$ to infinity to have $`\stackrel{~}{M}_P`$ finite, the effective $`D`$-dimensional cosmological constant $`M_P^{D2}V`$ would still go to $`\mathrm{}`$. This indicates that discontinuous solutions cannot be thought of as limits of smooth solutions. Examples of Discontinuous Solutions with $`\stackrel{~}{\mathrm{\Lambda }}0`$ Before we end this section, we would like to discuss discontinuous solutions with non-zero cosmological constant. In particular, let us consider the Randall-Sundrum setup. As before this is achieved by setting $`\varphi (y)\varphi _0`$, and $`V=\xi ^2`$. A restricted set of solutions with $`\stackrel{~}{\mathrm{\Lambda }}0`$ in this context were considered in . Here we would like to discuss the most general solutions of this type<sup>\††</sup><sup>\††</sup>\††Discontinuous domain walls with non-zero $`\stackrel{~}{\mathrm{\Lambda }}`$ were recently studied in a more general context in .. Thus, first consider the equations of motion for $`A`$ without the source term: $`(D1)(D2)(A^{})^2=\xi ^2+{\displaystyle \frac{D1}{D3}}\stackrel{~}{\mathrm{\Lambda }}\mathrm{exp}(2A),`$ (183) $`(D2)A^{\prime \prime }={\displaystyle \frac{1}{D3}}\stackrel{~}{\mathrm{\Lambda }}\mathrm{exp}(2A).`$ (184) Let $$\mathrm{exp}(A)U.$$ (185) Then we have: $`(D1)(D2)(U^{})^2=\xi ^2U^2+{\displaystyle \frac{D1}{D3}}\stackrel{~}{\mathrm{\Lambda }},`$ (186) $`(D1)(D2)U^{\prime \prime }=\xi ^2U.`$ (187) The general solution to this system of equations is given by $$U=C\left[\mathrm{exp}(ay)\frac{\stackrel{~}{\mathrm{\Lambda }}}{4\xi ^2C^2}\mathrm{exp}(ay)\right],$$ (188) where $$a\frac{\xi }{\sqrt{(D1)(D2)}},$$ (189) and $`C`$ is an integration constant. Note that if $`\stackrel{~}{\mathrm{\Lambda }}<0`$, then we must choose $`C>0`$, and the solution is non-singular. Moreover, $`U`$ blows up at both $`y\pm \mathrm{}`$, so even if we add the source term, the corresponding solution cannot localize gravity. On the other hand, for $`\stackrel{~}{\mathrm{\Lambda }}>0`$ $`U`$ vanishes for some finite $`y`$, so that the domain wall is singular. Note that without loss of generality we can assume that $`C>0`$. Then we can obtain a singular domain wall with non-vanishing $`\stackrel{~}{\mathrm{\Lambda }}`$ and finite $`\stackrel{~}{M}_P`$ by adding the source term and gluing together the appropriate parts of two smooth domain walls. Thus, consider the solution (for definiteness we assume $`a>0`$, that is, $`\xi >0`$) $$U(y)=C\left[\mathrm{exp}(a\|y\|)\frac{\stackrel{~}{\mathrm{\Lambda }}}{4\xi ^2C^2}\mathrm{exp}(a\|y\|)\right],$$ (190) where $$C>\frac{\sqrt{\stackrel{~}{\mathrm{\Lambda }}}}{2\xi }.$$ (191) The singularities (which are physical as they satisfy (137)) are located at $$y_\pm =\pm \frac{1}{a}\mathrm{ln}\left[\frac{2\xi C}{\sqrt{\stackrel{~}{\mathrm{\Lambda }}}}\right].$$ (192) The discontinuity is located at $`y_0=0`$, and the corresponding jump condition reads: $$f_0=4(D2)a\mathrm{coth}(a\|y_\pm \|)>0.$$ (193) Here we note that the solutions discussed in correspond to taking $`C\stackrel{~}{\mathrm{\Lambda }}/4\xi ^2C=1`$. However, generally $`C`$ is an independent integration constant. In particular, a priori $`\stackrel{~}{\mathrm{\Lambda }}`$ is not determined by specifying $`\xi `$ and $`f_0`$. However, the ratio of the effective $`(D1)`$-dimensional cosmological constant $`\stackrel{~}{M}_P^{D3}\stackrel{~}{\mathrm{\Lambda }}`$ to $`\stackrel{~}{M}_P^{D1}`$, that is, $`\stackrel{~}{\mathrm{\Lambda }}/\stackrel{~}{M}_P^2`$, depends on $`\stackrel{~}{\mathrm{\Lambda }}`$ and $`C`$ via the ratio $`\stackrel{~}{\mathrm{\Lambda }}/C^2`$, which is fixed via $`\xi `$ and $`f_0`$. Next, consider the same $`U(y)`$ as in (190) with $$C<\frac{\sqrt{\stackrel{~}{\mathrm{\Lambda }}}}{2\xi }.$$ (194) It then describes a non-singular domain wall which does not localize gravity. In fact, in this case the jump condition reads: $$f_0=4(D2)a\mathrm{coth}(\zeta )<0,$$ (195) where $$\zeta \frac{1}{a}\mathrm{ln}\left[\frac{\sqrt{\stackrel{~}{\mathrm{\Lambda }}}}{2\xi C}\right].$$ (196) Note that such a domain wall would require a negative tension “brane”. ### F Summary In this subsection we would like to briefly summarize the results of this section. (i) Non-singular smooth domain walls with $`A\mathrm{}`$ at $`y\pm \mathrm{}`$ (that is, finite compactification volume and $`\stackrel{~}{M}_P`$) cannot have $`\stackrel{~}{\mathrm{\Lambda }}>0`$. On the other hand, non-singular smooth domain walls approaching AdS vacua on at least one side can only have $`\stackrel{~}{\mathrm{\Lambda }}=0`$. (ii) Singular smooth domain walls can a priori have vanishing or non-vanishing $`\stackrel{~}{\mathrm{\Lambda }}`$ of either sign. Singularities with $`A\mathrm{}`$ are physical only if the scalar potential is unbounded below there. However, singularities where $`A\mathrm{const}.`$, which can only arise if the potential is unbounded above there, are unphysical. In particular, it is not consistent to cut off the space at such singularities in the $`y`$ direction. (iv) The aforementioned conclusions also apply to discontinuous domain walls (with a $`\delta `$-function source term). However, such a domain wall cannot be thought of as a limit of a smooth domain wall with non-zero $`\stackrel{~}{M}_P`$. ## VI Normalizable Modes In this section we would like to discuss normalizable modes living on the domain walls discussed in the previous sections. In particular, let us focus on non-singular domain walls. To study normalizable modes, it is convenient to consider small fluctuations around the background metric $`G_{MN}`$ and the scalar field $`\varphi `$, which we will denote by $`h_{MN}`$ and $`\phi `$, respectively. Using the diffeomorphism invariance of the theory we can choose the following gauge: $$h_{\mu D}=h_{DD}=0.$$ (197) Instead of the remaining metric fluctuations $`h_{\mu \nu }`$, it will be convenient to work with $`\stackrel{~}{h}_{\mu \nu }`$ defined via $$h_{\mu \nu }=\mathrm{exp}(2A)\stackrel{~}{h}_{\mu \nu }.$$ (198) The quadratic action for $`\stackrel{~}{h}_{\mu \nu }`$ and $`\phi `$ is given by $`S[\stackrel{~}{h}_{\mu \nu },\phi ]/M_P^{D2}=`$ (199) $`{\displaystyle d^Dx\mathrm{exp}[(D3)A]\left\{\left(\sqrt{\widehat{g}}\left[\widehat{R}\stackrel{~}{\mathrm{\Lambda }}\right]\right)^{(2)}\sqrt{\stackrel{~}{g}}\frac{1}{\alpha ^2}(\stackrel{~}{}\phi )^2\right\}}`$ (200) $`{\displaystyle d^Dx\mathrm{exp}[(D1)A]\sqrt{\stackrel{~}{g}}\left\{\frac{1}{\alpha ^2}(\phi ^{})^2+\frac{1}{2}V_{\varphi \varphi }\phi ^2\frac{1}{\alpha ^2}\varphi ^{}\stackrel{~}{h}^{}\phi +\frac{1}{4}\left[(\stackrel{~}{h}_{\mu \nu }^{})^2(\stackrel{~}{h}^{})^2\right]\right\}}.`$ (201) Here we are using the following notations. First, $`\widehat{g}det(\widehat{g}_{\mu \nu })`$, where $`\widehat{g}_{\mu \nu }\stackrel{~}{g}_{\mu \nu }+\stackrel{~}{h}_{\mu \nu }`$. Also, $`\widehat{R}`$ is the $`(D1)`$-dimensional Ricci scalar constructed from the metric $`\widehat{g}_{\mu \nu }`$. The superscript “$`(2)`$” in the term $$\left(\sqrt{\widehat{g}}\left[\widehat{R}\stackrel{~}{\mathrm{\Lambda }}\right]\right)^{(2)}$$ (202) indicates that only the terms quadratic in $`\stackrel{~}{h}_{\mu \nu }`$ should be kept in the corresponding expression. Finally, $`\stackrel{~}{h}\stackrel{~}{h}_\mu ^\mu `$, and $`(\stackrel{~}{h}_{\mu \nu }^{})^2(\stackrel{~}{h}_{\mu \nu })^{}(\stackrel{~}{h}^{\mu \nu })^{}`$, where the indices are lowered and raised with the metric $`\stackrel{~}{g}_{\mu \nu }`$ and its inverse. We also note that the boundary contributions discussed in subsection C of section V vanish in the case of non-singular domain walls. As in the case of non-gravitational domain walls, the above action has a much simpler form if instead of $`\phi `$ we use the field $`\omega `$ defined via $$\phi \varphi ^{}\omega .$$ (203) Thus, the quadratic action for $`\stackrel{~}{h}_{\mu \nu }`$ and $`\omega `$ is given by: $`S[\stackrel{~}{h}_{\mu \nu },\omega ]/M_P^{D2}=`$ (204) $`{\displaystyle d^Dx\mathrm{exp}[(D3)A]\left\{\left(\sqrt{\widehat{g}}\left[\widehat{R}\stackrel{~}{\mathrm{\Lambda }}\right]\right)^{(2)}\sqrt{\stackrel{~}{g}}\frac{1}{\alpha ^2}(\varphi ^{})^2\left[(\stackrel{~}{}\omega )^2\frac{D1}{D2}\frac{\stackrel{~}{\mathrm{\Lambda }}}{D3}\omega ^2\right]\right\}}`$ (205) $`{\displaystyle d^Dx\mathrm{exp}[(D1)A]\sqrt{\stackrel{~}{g}}\left\{\frac{1}{\alpha ^2}(\varphi ^{})^2(\omega ^{})^2+\frac{1}{4}\left[(\stackrel{~}{h}_{\mu \nu }^{})^2\frac{D2}{D1}\left(\stackrel{~}{h}^{}+2\gamma ^2(\varphi ^{})^2\omega \right)^2\right]\right\}}.`$ (206) Here $$\stackrel{~}{h}_{\mu \nu }^{}\stackrel{~}{h}_{\mu \nu }\frac{1}{D1}\stackrel{~}{h}\stackrel{~}{g}_{\mu \nu }$$ (207) is the traceless part of $`\stackrel{~}{h}_{\mu \nu }`$. Next, we would like to study normalizable modes. Here we are assuming that the warp factor $`\mathrm{exp}(2A)`$ goes to zero at $`y\pm \mathrm{}`$ fast enough, so that the compactifications volume as well as $`\stackrel{~}{M}_P`$ are finite. Then normalizable modes correspond to configurations with $`\omega ^{}=0`$. It is not difficult to show that the equation of motion for such $`\omega `$ simplifies as follows: $$\stackrel{~}{h}^{}=2\mathrm{exp}(2A)\stackrel{~}{}^2\omega +2(D1)A^{\prime \prime }\omega .$$ (208) Note that this equation does not explicitly contain $`\stackrel{~}{\mathrm{\Lambda }}`$. Thus, the equation of motion (208) for $`\omega =\omega (x^\mu )`$ seems to imply that we have normalizable modes satisfying the $`(D1)`$-dimensional Klein-Gordon equation $$\stackrel{~}{}^2\omega =m^2\omega $$ (209) for any $`m^2`$ (including massive and tachyonic). At first this might appear to imply that the domain wall is unstable. This is, however, not the case due to the following. Note that $`\omega `$ (or, equivalently, $`\phi `$) mixes with $`\stackrel{~}{h}^{}`$, that is, the $`y`$ derivative of the trace part of $`\stackrel{~}{h}_{\mu \nu }`$. Moreover, as we will see in a moment, there is a residual diffeomorphism invariance which respects the gauge choice (197). These facts have important implications which we would like to discuss next. ### A Residual Diffeomorphism Invariance The gauge choice (197) does not fix all the gauge freedom in the system. Thus, there is a residual diffeomorphism invariance which preserves (197)<sup>\‡‡</sup><sup>\‡‡</sup>\‡‡This residual gauge invariance was discussed in the context of the Randall-Sundrum model in , and also in .. For our purposes here it will suffice to consider infinitesimal residual diffeomorphisms. Thus, we have $$\delta h_{MN}=_M\xi _N+_N\xi _M.$$ (210) In particular, $$\delta h_{DD}=2\xi _D^{},$$ (211) and the gauge condition $`h_{DD}=0`$ is preserved if and only if $$\xi _D^{}=0.$$ (212) Furthermore, we have $$\delta h_{\mu D}=\stackrel{~}{}_\mu \xi _D+\xi _\mu ^{}2A^{}\xi _\mu ,$$ (213) where on the r.h.s. we have kept only the terms independent of $`h_{MN}`$, which is consistent with the linearized approximation. Requiring that the gauge condition $`h_{\mu D}=0`$ is preserved, we obtain $$\xi _\mu =\mathrm{exp}(2A)\stackrel{~}{\xi }_\mu ,$$ (214) where $`\stackrel{~}{\xi }_\mu `$ satisfies the following first order differential equation $$\stackrel{~}{\xi }_\mu ^{}=\mathrm{exp}(2A)\stackrel{~}{}_\mu \xi _D.$$ (215) Finally, for the rest of the $`h_{MN}`$ components we obtain $$\delta h_{\mu \nu }=\stackrel{~}{}_\mu \xi _\nu +\stackrel{~}{}_\nu \xi _\mu +2A^{}\mathrm{exp}(2A)\xi _D\stackrel{~}{g}_{\mu \nu },$$ (216) or, equivalently, $$\delta \stackrel{~}{h}_{\mu \nu }=\stackrel{~}{}_\mu \stackrel{~}{\xi }_\nu +\stackrel{~}{}_\nu \stackrel{~}{\xi }_\mu +2A^{}\xi _D\stackrel{~}{g}_{\mu \nu }.$$ (217) In particular, we have $$\delta \stackrel{~}{h}=2\stackrel{~}{}^\mu \stackrel{~}{\xi }_\mu +2(D1)A^{}\xi _D,$$ (218) from which it follows that $$\delta \stackrel{~}{h}^{}=2\mathrm{exp}(2A)\stackrel{~}{}^2\xi _D+2(D1)A^{\prime \prime }\xi _D.$$ (219) The last result implies that for any choice of $`\omega `$ in (208) (where $`\omega `$ is assumed to be independent of $`y`$), the corresponding $`\stackrel{~}{h}`$ is a pure gauge, and the corresponding gauge parameter is given by $$\xi _D=\omega .$$ (220) That is, $`\omega `$ is not a propagating degree of freedom, but a gauge parameter, and, up to a gauge transformation, $`\stackrel{~}{h}`$ is independent of $`y`$. This result has a simple physical interpretation to which we now turn. ### B Localization of Gravity as a Higgs Mechanism The above analyses suggest that localization of gravity is actually a Higgs mechanism for the graviton. In particular, in the $`D`$-dimensional language we have the fields $`h_{MN}`$ and $`\phi `$ with $`D(D3)/2`$ and 1 physical degrees of freedom, respectively. In the $`(D1)`$-dimensional language together with the scalar $`\phi `$ we have the fields $`h_{\mu \nu }`$, $`h_{\mu D}`$ and $`h_{DD}`$ corresponding to the $`(D1)`$-dimensional graviton, a graviphoton and a scalar, respectively. If we turn off the coupling to gravity, there is a $`\phi `$ zero mode (which might or might not be normalizable in this context) corresponding to the Goldstone mode of the broken translational invariance in the $`y`$ direction. However, once we gauge diffeomorphisms, that is, once we consider the coupling to gravity, the $`\phi `$ zero mode is eaten by the graviphoton, which now becomes massive. The scalar $`h_{DD}`$ is also massive - the volume of the compactification, which is finite, is not a free parameter but is completely fixed in terms of the parameters in the scalar potential $`V`$. This is analogous to what happens, say, in the Abelian Higgs model, where a $`U(1)`$ gauge field is coupled to a single complex scalar field. Once the latter acquires an expectation value, the angular part is eaten by the gauge field, which therefore becomes massive, while the radial part acquires a mass due to a non-trivial scalar potential. However, as we will see in a moment, the Higgs mechanism that takes place in the localization of gravity has certain novel features, which is due to important differences between gravity and an ordinary gauge theory. ### C Zero and Massive Modes To analyze normalizable modes, consider the standard factorized form of the corresponding field, which is $`\stackrel{~}{h}_{\mu \nu }`$ in this case: $$\stackrel{~}{h}_{\mu \nu }(x^\lambda ,y)=\zeta _{\mu \nu }(x^\lambda )\mathrm{\Sigma }(y).$$ (221) Using the fact that $`\omega `$ in (208) is a gauge parameter, without loss of generality we can set it to zero, which implies $`\stackrel{~}{h}^{}=0`$, and, consequently, either $`\mathrm{\Sigma }^{}=0`$, or $`\zeta _\mu ^\mu =0`$. In the former case we have a zero mode $$\stackrel{~}{h}_{\mu \nu }(x^\lambda ,y)=\mathrm{const}.\times \zeta _{\mu \nu }(x^\lambda ),$$ (222) where $`\zeta _{\mu \nu }`$ satisfies the equation of motion for the massless $`(D1)`$-dimensional graviton, which is nothing but the Euler-Lagrange equation obtained by varying the quadratic $`(D1)`$-dimensional action $$\stackrel{~}{S}[\stackrel{~}{h}_{\mu \nu }]\stackrel{~}{M}_P^{D3}d^{D1}x\left(\sqrt{\widehat{g}}\left[\widehat{R}\stackrel{~}{\mathrm{\Lambda }}\right]\right)^{(2)}$$ (223) w.r.t. $`\stackrel{~}{h}_{\mu \nu }`$ (or, equivalently, $`\zeta _{\mu \nu }`$). In this case the zero mode is quadratically normalizable, that is, the $`(D1)`$-dimensional Planck mass $`\stackrel{~}{M}_P`$ is finite. Next, if $`\mathrm{\Sigma }^{}0`$, then using (204) it is not difficult to show that $`\mathrm{\Sigma }(y)`$ must satisfy the equation $$\left(\mathrm{exp}[(D1)A]\mathrm{\Sigma }^{}\right)^{}+m^2\mathrm{exp}[(D3)A]\mathrm{\Sigma }=0,$$ (224) where $`m^2`$ is a constant. As to $`\zeta _{\mu \nu }`$, it then satisfies the equation of motion of a $`(D1)`$-dimensional symmetric traceless (recall that $`\stackrel{~}{h}=0`$ in this case) two-index tensor field with mass squared $`m^20`$ propagating in the background metric $`\stackrel{~}{g}_{\mu \nu }`$. Using (224) the action for $`\stackrel{~}{h}_{\mu \nu }`$ can be written as $$\stackrel{~}{S}[\stackrel{~}{h}_{\mu \nu }]M_P^{D2}d^{D1}x𝑑y\mathrm{exp}[(D3)A]\left\{\left(\sqrt{\widehat{g}}\left[\widehat{R}\stackrel{~}{\mathrm{\Lambda }}\right]\right)^{(2)}\frac{1}{4}m^2(\stackrel{~}{h}_{\mu \nu }^{})^2\right\},$$ (225) and the norm of $`\stackrel{~}{h}_{\mu \nu }`$ is proportional to $$\stackrel{~}{h}_{\mu \nu }^2𝑑y\mathrm{exp}[(D3)A]\mathrm{\Sigma }^2.$$ (226) We would now like to show that massive modes with $`m^2>0`$ are plane-wave normalizable (in contrast to the zero mode, which is quadratically normalizable), while tachyonic modes with $`m^2<0`$ are not normalizable. To see this, let us make the coordinate transformation $`yz`$ so that the metric takes the form: $$ds_D^2=\mathrm{exp}(2A)\left(\stackrel{~}{g}_{\mu \nu }dx^\mu dx^\nu +dz^2\right).$$ (227) That is, $$dy=\mathrm{exp}(A)dz,$$ (228) where we have chosen the overall sign so that $`z\pm \mathrm{}`$ as $`y\pm \mathrm{}`$. In this new coordinate system (224) reads: $$\mathrm{\Sigma }_{zz}+(D2)A_z\mathrm{\Sigma }_z+m^2\mathrm{\Sigma }=0,$$ (229) where the subscript $`z`$ denotes derivative w.r.t. $`z`$. Let $$\mathrm{\Sigma }\mathrm{exp}\left[\frac{1}{2}(D2)A\right]\widehat{\mathrm{\Sigma }}.$$ (230) The equation for $`\widehat{\mathrm{\Sigma }}`$ reads: $$\widehat{\mathrm{\Sigma }}_{zz}+\left[m^2\frac{1}{2}(D2)A_{zz}\frac{1}{4}(D2)^2(A_z)^2\right]\widehat{\mathrm{\Sigma }}=0.$$ (231) Let us assume that $`m^20`$. Then for large $`z`$ the $`m^2`$ term is dominant in the square brackets in (231) for $`A_z=A^{}\mathrm{exp}(A)0,`$ (232) $`A_{zz}=\left[A^{\prime \prime }+(A^{})^2\right]\mathrm{exp}(2A)0`$ (233) for large $`y`$ (that is, large $`z`$). This then implies that for $`m^2>0`$ $`\widehat{\mathrm{\Sigma }}`$ has the following asymptotic behavior at large $`z`$: $$\widehat{\mathrm{\Sigma }}(z)C_1\mathrm{sin}(mz)+C_2\mathrm{cos}(mz),$$ (234) where $`C_1,C_2`$ are some constants. On the other hand, in terms of $`\widehat{\mathrm{\Sigma }}`$ the norm of $`\stackrel{~}{h}_{\mu \nu }`$ (226) is given by $$𝑑y\mathrm{exp}[(D3)A]\mathrm{\Sigma }^2=𝑑z\widehat{\mathrm{\Sigma }}^2.$$ (235) This implies that massive $`\stackrel{~}{h}_{\mu \nu }`$ modes are plane-wave normalizable (but not quadratically normalizable). Thus, we have a continuum of massive bulk modes just as in the Randall-Sundrum model. What about the tachyonic modes? It is not difficult to see that these are not normalizable. Indeed, (231) can be written as $$(QQ^{})\widehat{\mathrm{\Sigma }}=m^2\widehat{\mathrm{\Sigma }},$$ (236) where $`Q_z+{\displaystyle \frac{1}{2}}(D2)A_z,`$ (237) $`Q^{}=_z+{\displaystyle \frac{1}{2}}(D2)A_z.`$ (238) Evidently (236) does not have normalizable solutions with $`m^2<0`$ as it is nothing but Schrödinger’s equation of a supersymmetric quantum mechanics with the Hamiltonian $`QQ^{}`$ . Before we end this section, let us make the following remark. The existence of the quadratically normalizable zero mode as well as plane-wave normalizable massive bulk modes in the above setup is similar to what happens in the original Randall-Sundrum model. However, localization of gravity in the case of smooth domain walls arises via a spontaneous breakdown of diffeomorphism invariance, while in the Randall-Sundrum model as well as the cases where we have discontinuous domain walls diffeomorphism invariance is explicitly broken by the “brane”. That is, in the former case we have Higgs mechanism, while in the latter case we do not, and, in particular, there is no massless graviphoton to begin with. This is precisely the underlying reason why discontinuous domain walls cannot be thought of as limits of smooth domain walls. ## VII The Upshot In the previous sections we have studied possible types of domain wall solutions of (1) which localize gravity. In this section we would like to discuss some aspects of the cosmological constant problem within the context of a brane world realized via such a domain wall. In particular, for the reasons which should be clear from our previous discussions, we will focus on smooth domain walls. To begin with, let us consider non-singular domain walls with infinite tension. In particular, suppose the corresponding scalar potential is such that it does not have local minima and is unbounded below. Then it is clear that no solution that does not localize gravity is consistent. However, as we saw in the previous sections, subject to non-singularity conditions on the scalar potential, there are (classically) consistent domain wall solutions which localize gravity. This implies that the theory will ultimately<sup>†\</sup><sup>†\</sup>†\Here we should point out that a priori there might also exist consistent localized solutions of codimension $`r>1`$ if we have $`r`$ scalar fields. In such a case the corresponding theory might also be able to choose such backgrounds. end up choosing such a background. In such cases we therefore have spontaneous localization of gravity via the Higgs mechanism discussed in subsection B of the previous section. Had it been the case that domain wall solutions with infinite tension could only have vanishing $`(D1)`$-dimensional cosmological constant, the aforementioned property a priori might have been considered as a possible solution to the cosmological constant problem. However, as we saw in the previous sections, infinite tension domain walls can have negative cosmological constant. Could we then conclude that such domain walls solve “half” of the cosmological constant problem as they do not allow $`\stackrel{~}{\mathrm{\Lambda }}>0`$? Here we would like to point out that such a conclusion might be premature. The reason is twofold. First, to have a non-singular domain wall, the scalar potential must satisfy certain conditions discussed in the previous sections. Quantum corrections can a priori modify the potential in such a way that these conditions are no longer satisfied. For singular domain walls, on the other hand, non-zero cosmological constant of either sign is a priori allowed. Nonetheless, one could in principle argue that (at least in some sense) the subset of potentials that satisfy non-singularity conditions has non-zero measure, so this particular objection might be avoidable. There is, however, another point here, which could potentially be more problematic. Quantum corrections modify not only the scalar potential, but generically also introduce higher derivative, in particular, higher curvature terms, which become important in such backgrounds. Thus, for instance, the $`D`$-dimensional Ricci scalar is given by: $$R=\stackrel{~}{R}\mathrm{exp}(2A)(D1)\left[2A^{\prime \prime }+D(A^{})^2\right],$$ (239) where $`\stackrel{~}{R}`$ is the $`(D1)`$-dimensional Ricci scalar. Even if $`\stackrel{~}{R}=0`$ on the solution (that is, if the $`(D1)`$-dimensional cosmological constant is vanishing), the $`D`$-dimensional Ricci scalar blows up at large $`y`$ unless $`A^{}`$ goes to a constant there<sup>††</sup><sup>††</sup>††Note that if on the solution $`\stackrel{~}{R}0`$, then $`R`$ ultimately blows up at large $`y`$ as long as $`A`$ goes to $`\mathrm{}`$ there.. In the case of infinite tension domain walls $`A^{}`$ diverges at large $`y`$, and so does the Ricci scalar $`R`$, so that higher curvature (or, more generally, higher derivative) terms cannot be ignored. What is worse, in this case it is not even clear if there exists a controlled approximation scheme for taking into account such terms. The underlying reason for this is that such domain walls are actually singular but the singularities are located at $`y\pm \mathrm{}`$ (that is, at infinite distances from the “core” of the domain wall). In particular, a priori there does not seem to be a reason to believe that any conclusions about allowed values of the $`(D1)`$-dimensional cosmological constant would not be modified by higher curvature terms. In fact, as we will point out in the following, in this context it is not even clear whether such domain walls do indeed localize gravity. Note that the above discussion also applies to singular domain walls with singularities located at finite values of $`y`$. The above discussion suggests that we consider solutions where the $`D`$-dimensional Ricci scalar $`R`$ goes to constant values at large $`y`$. Such solutions should have vanishing $`(D1)`$-dimensional cosmological constant, and $`A^{}`$ should go to constant values at large $`y`$. In particular, if the cut-off scale for higher derivative terms in the action is of order $`M_{}`$, then for such solutions we have a controlled approximation scheme for including higher curvature terms as long as $$\|R\|(A^{})^2M_{}^2.$$ (240) In fact, from our previous discussions it should be clear that such solutions must necessarily interpolate between two adjacent local AdS minima of the scalar potential. However, if such minima are located at finite values of $`\varphi `$, then a priori there exist other solutions with non-zero cosmological constant<sup>†‡</sup><sup>†‡</sup>†‡If the scalar potential is bounded below, then before taking into account higher derivative terms such solutions are necessarily singular and have positive cosmological constant. Otherwise one can a priori have singular or non-singular solutions with negative cosmological constant. However, we should mention that these conclusions might be modified by higher curvature terms.. These solutions have the property that $`\varphi `$ “overshoots” the values corresponding to the local AdS minima. To avoid this, one could consider potentials which have no local minima but approach constant negative values at large $`\varphi `$. ### A Runaway Potentials Thus, we are lead to consider the following setup. Let the scalar potential, such that it has no local minima, be of the runaway type, that is, it approaches constant values at large $`\varphi `$. More precisely, let these values be negative, and satisfy the condition $$\|V(\varphi \pm \mathrm{})\|M_{}^2,$$ (241) with $`M_{}`$ defined as above. Then the corresponding domain walls necessarily have vanishing cosmological constant. In fact, it is not difficult to see that this statement holds even if we take into account higher curvature (or, more generally, higher derivative) terms as these are now under control<sup>†§</sup><sup>†§</sup>†§More precisely, this is correct subject to the convergence of the $`\|R\|/M_{}^2`$ expansion.. We therefore conclude that for such domain walls (at least in some sense - see below) the higher curvature corrections do not change the conclusion that the cosmological constant must vanish. There is, however, the following issue here. A priori there does not seem to be a reason to believe that quantum corrections will not modify the behavior of the scalar potential at large $`\varphi `$. In particular, if instead of constant values the potential goes to $`+\mathrm{}`$ or $`\mathrm{}`$ at large $`\varphi `$, then the conclusion that the cosmological constant must vanish would be modified<sup>†¶</sup><sup>†¶</sup>†¶Actually, if the potential is of the aforementioned runaway type on at least one side, it is enough to ensure vanishing of the cosmological constant.. Note that supersymmetry locally broken in the “core” of the domain wall does not seem to help here - since the compactification volume is actually finite, supersymmetry breaking in the bulk is non-vanishing as it is only suppressed by the volume factor. At any rate, if the runaway behavior of the scalar potential is not lifted by quantum corrections<sup>†∥</sup><sup>†∥</sup>†∥Here we should mention that typically it is not easy to lift runaway behavior perturbatively., then such domain walls could a priori lead to a solution to the cosmological constant problem. Here we would like to discuss an explicit example of such a domain wall. Since the domain wall has vanishing $`\stackrel{~}{\mathrm{\Lambda }}`$, it is convenient to express the scalar potential in terms of $`W`$ via (16). Thus, consider the example where $$W=\xi \mathrm{tanh}(\varphi ).$$ (242) The scalar potential reads: $$V=\xi ^2\left[\frac{1}{\mathrm{cosh}^4(\varphi )}\gamma ^2\mathrm{tanh}^2(\varphi )\right],$$ (243) and it approaches a constant value $`V\gamma ^2\xi ^2`$ at $`\varphi \pm \mathrm{}`$. The solution for $`\varphi `$ and $`A`$ is given by $`2\varphi +\mathrm{sinh}(2\varphi )=4\alpha \xi (yy_0),`$ (244) $`A=C+{\displaystyle \frac{\beta }{4\alpha }}\mathrm{cosh}(2\varphi ),`$ (245) where $`y_0`$ and $`C`$ are integration constants. Note that $`\|\varphi \|`$ grows logarithmically for large $`\|y\|`$, whereas $`A`$ linearly goes to $`\mathrm{}`$. In particular, higher derivative terms are either vanishing at large $`y`$, or subleading provided that $`\xi ^2M_{}^2`$. ### B Delocalization of Gravity The discussion in the previous subsection might appear promising in the context of the cosmological context problem. There is, however, an additional issue here, which we would like to discuss next. The point is that higher curvature terms appear to delocalize gravity. Indeed, let us, say, consider adding to (1) a generic term of the form $$\zeta d^Dx\sqrt{G}R^k,$$ (246) where $`\zeta `$ has dimension of $`(\mathrm{mass})^{D2k}`$. We have $$R^k=\stackrel{~}{R}^k\mathrm{exp}(2kA)+\mathrm{},$$ (247) where the ellipses stand for terms containing lower powers of $`\stackrel{~}{R}`$. The contribution of the $`\stackrel{~}{R}^k`$ term to the action then reads: $$\zeta d^{D1}x𝑑y\mathrm{exp}[(D12k)A]\sqrt{\stackrel{~}{g}}\stackrel{~}{R}^k.$$ (248) Assuming that $`A`$ goes to $`\mathrm{}`$ at large $`y`$, for large enough $`k`$ the corresponding $`y`$ integral $$𝑑y\mathrm{exp}[(D12k)A]$$ (249) diverges, which implies that gravity is no longer localized if we include such a higher curvature term. In fact, for $`D=5`$ delocalization of gravity generically occurs already at the four-derivative level, that is, once we include $`R^2`$, $`R_{MN}R^{MN}`$ and $`R_{MNST}R^{MNST}`$ terms. In particular, it is not difficult to show that once we include higher curvature terms there is no normalizable graviton zero mode, and we no longer have $`(D1)`$-dimensional gravity localized on the wall, in particular, $`(D1)`$-dimensional Newton’s law is no longer reproduced<sup>†\</sup><sup>†\</sup>†\*Here we could ask how is gravity modified. However, generically inclusion of four-derivative terms leads to negative-norm states in the Hilbert space of the theory thus violating unitarity.. As we have already mentioned, in the case of runaway potentials considered in the previous subsection we have control over higher derivative terms as far as the background is concerned. However, what we see here is that we do not have control over the fluctuations around such a background. The reason for this is that at large $`y`$ we approach the horizon of the AdS space where gravity is strongly coupled. It is unclear whether there is any “non-perturbative” approach to this problem, but in any case the requirement that one reproduce the $`(D1)`$-dimensional Einstein gravity at low energies seems reasonable<sup>†††</sup><sup>†††</sup>†††Here we would like to point out a possible way out of this difficulty. If the warp factor $`A`$ goes to $`+\mathrm{}`$ instead of $`\mathrm{}`$ at large $`y`$, then in, say, $`D=5`$ terms with more than four derivatives are now normalizable. On the other hand, one then must arrange for the lower-derivative terms to be normalizable as well, which a priori does not seem impossible. Such a framework is interesting as here the volume of the fifth dimension is actually infinite. As was pointed out in , in the supersymmetric context this could have important implications for the cosmological constant problem. In particular, supersymmetry locally broken on a “brane” (e.g., the “core” of a smooth domain wall) is not transmitted to the bulk as supersymmetry breaking in the bulk is suppressed by the compactification volume, which is infinite. For a recent discussion of infinite volume extra dimension scenarios, see .. Before we end this subsection, we would like to make the following remark. If we forget about higher curvature terms beyond four derivatives, then for $`D=5`$ there is an exception to the aforementioned observation that gravity is delocalized. In particular, consider adding the Gauss-Bonnet term to the 5-dimensional action $$\zeta d^5x\sqrt{G}\left[R^24R_{MN}R^{MN}+R_{MNST}R^{MNST}\right].$$ (250) The corresponding term containing four derivatives w.r.t. the 4-dimensional coordinates $`x^\mu `$ then reads: $$\zeta d^4x𝑑y\sqrt{\stackrel{~}{g}}\left[\stackrel{~}{R}^24\stackrel{~}{R}_{\mu \nu }\stackrel{~}{R}^{\mu \nu }+\stackrel{~}{R}_{\mu \nu \sigma \tau }\stackrel{~}{R}^{\mu \nu \sigma \tau }\right].$$ (251) Note that the integral over $`y`$ is divergent - there is no warp factor in this expression. However, this term is harmless as long as the four-dimensional space-time is topologically trivial. Indeed, the Gauss-Bonnet term in four dimensions is a total derivative. Note that if we considered adding the Gauss-Bonnet term in $`D>5`$, then the analog of the aforementioned term would be normalizable as the $`y`$ integral is now non-trivially weighted: $$𝑑y\mathrm{exp}[(D5)A].$$ (252) Note, however, that if we expand around the $`(D1)`$-dimensional Minkowski space, the term quadratic in $`\stackrel{~}{h}_{\mu \nu }`$ is a total derivative (albeit the full Gauss-Bonnet term in $`D1>4`$ dimensions no longer is). This implies that the $`(D1)`$-dimensional graviton zero mode is unaffected even if $`D>5`$. Discontinuous domain wall solutions in the presence of the Gauss-Bonnet term were recently studied in . Smooth domain walls in this context appear to have some interesting properties which will be discussed elsewhere <sup>†‡‡</sup><sup>†‡‡</sup>†‡‡Here we should point out the following issue arising in this context. As we have already mentioned, if we expand the ($`D`$-dimensional) Gauss-Bonnet term around a flat background, the term quadratic in $`h_{MN}`$ is a total derivative. In particular, it does not lead to violation of unitarity unlike a generic four-derivative action. However, in a non-trivial warped background the linearized theory has a non-trivial propagator. Nonetheless, we expect that if gravity is localized, then the corresponding Higgs mechanism ensures absence of propagating negative-norm states.. ### C Comments Before we end our discussion, we would like to make a few comments. First, the aforementioned difficulties with the higher derivative terms might be under better control if we consider a supersymmetric setup. However, attempts to construct non-singular smooth supersymmetric domain walls in the context of, say, $`D=5`$ $`𝒩=2`$ gauged supergravity have not been successful so far (see, e.g., <sup>‡\</sup><sup>‡\</sup>‡\In a recent paper a construction of such a domain wall was reported. However, as was pointed out in , the corresponding solution only exists if one ad hoc freezes one of the scalar fields whose dynamics is such that the domain wall is actually singular.). In particular, typically one searches for solutions interpolating between two local AdS minima of the scalar potential. However, as was argued in , on general grounds potentials with more then one AdS minima are not expected to exist in $`D=5`$ $`𝒩=2`$ gauged supergravity<sup>‡†</sup><sup>‡†</sup>‡†Note, however, that such domain walls have been constructed within the framework of $`D=4`$ $`𝒩=1`$ supergravity - see for a review.. A possible way around this difficulty could be to consider potentials with, say, no local minima at all. As we discussed in the previous sections, such potentials admit non-singular solutions subject to certain conditions. Understanding whether such potentials can arise in $`D=5`$ $`𝒩=2`$ gauged supergravity (that is, if they avoid the general “no-go” arguments of ) is beyond the scope of this paper. At any rate, it would be interesting if one could construct a supersymmetric infinite tension domain wall, or else prove that such solutions do not exist in $`D=5`$ $`𝒩=2`$ gauged supergravity<sup>‡‡</sup><sup>‡‡</sup>‡‡For a recent analysis of general $`D=5`$ $`𝒩=2`$ gauged supergravity models, see .. The second comment concerns the “self-tuning” approach of to the cosmological constant problem. In it was argued that fine-tuning in these solutions is hidden in singularities. In it was argued that in such solutions one either has a naked singularity or fine-tuning. Here we would like to point out that, as we have already mentioned, the singularities in these solutions are unphysical. In particular, it does not seem consistent to truncate the space in the $`y`$ direction at such singularities. At any rate, the self-tuning property would only imply that the four-dimensional cosmological constant is not affected by the quantum effects on the “brane”. However, once bulk contributions are included, one seems to lose control over the cosmological constant. In fact, a setup where pure brane contributions to the cosmological constant vanish at least in the conformal/large $`N`$ limit has been known , but as usual there too one loses control over the cosmological constant once the coupling to gravity is included. ###### Acknowledgements. I would like to thank Gregory Gabadadze, Martin Roček, Gary Shiu, Stefan Vandoren and Peter van Nieuwenhuizen for discussions, and especially Gia Dvali for valuable discussions and collaboration at some stages of this work. This work was supported in part by the National Science Foundation. I would also like to thank Albert and Ribena Yu for financial support.
warning/0005/math0005278.html	ar5iv	text	# Narrow operators and rich subspaces of Banach spaces with the Daugavet property ## 1. Introduction Following we say that a Banach space $`X`$ has the Daugavet property if for every operator $`T:XX`$ of rank $`1`$ the Daugavet equation (1.1) $$\mathrm{Id}+T=1+T$$ is fulfilled. It is known that then every weakly compact operator, even every strong Radon-Nikodým operator, and every operator not fixing a copy of $`\mathrm{}_1`$ satisfies (1.1) as well (, ). Incidentally, this shows that our definition of the Daugavet property is equivalent to the ones which have been proposed in and . Classical results due to Daugavet , Lozanovskii , and Foiaş, Singer and Pełczyński state that $`C(K)`$, $`L_1(\mu )`$ and $`L_{\mathrm{}}(\mu )`$ have the Daugavet property provided that $`K`$ is perfect and $`\mu `$ is non-atomic. Recently, corresponding results in the non-commutative setting were obtained by Oikhberg . The papers and study Banach spaces with the Daugavet property from a structural point of view; for example it is shown that such a space never embeds into a space having an unconditional basis, and it contains (many) subspaces isomorphic to $`\mathrm{}_1`$. Also, hereditary properties of the Daugavet property are established there. Returning to the classical spaces $`C(K)`$ and $`L_1(\mu )`$ we mention that a different approach to (1.1) on these spaces was launched earlier in and . These papers study a duality between certain operators, called narrow operators, and certain subspaces, called rich subspaces, of such spaces. (For the definitions, which differ in the two cases, see Section 2.) One of the key features of this approach is that the concept of a narrow operator on $`C(K)`$ or $`L_1(\mu )`$, which makes sense for operators from these spaces into an arbitrary range space, only depends on the values $`Tx`$, but not on the images $`Tx`$ themselves. The idea of the present paper is to introduce narrow operators and rich subspaces in general. In Section 2 we propose a formalism in order to deal with those properties of an operator which depend only on the norms of the images of elements. We define corresponding equivalence classes and their formal sums and differences, which is reminiscent of certain procedures in the theory of operator ideals. Then, in Section 3 we introduce and study narrow operators on Banach spaces with the Daugavet property. We show, in particular, that strong Radon-Nikodým operators are narrow and that narrow operators mapping $`X`$ to itself satisfy (1.1). In Section 4 we prove that operators not fixing a copy of $`\mathrm{}_1`$ are narrow, thus extending a result from . To do so we need an extension of a theorem due to Rosenthal characterising separable Banach spaces that fail to contain isomorphic copies of $`\mathrm{}_1`$ (Theorem 4.3), which seems to be of independent interest. Section 5 deals with rich subspaces. As in the classical case of $`C(K)`$ or $`L_1(\mu )`$, a closed subspace $`YX`$ is called rich if the quotient map $`q:XX/Y`$ is narrow. One of the main results here is that the Daugavet property passes to rich subspaces, which leads to new hereditary properties. We also study a related class of subspaces which we term wealthy. What looks like a quibble of words is another main result from Section 5: a subspace is rich if and only if it is wealthy. In fact, we also need to deal with a slightly more general class of operators called strong Daugavet operators. It turns out that there are strong Daugavet operators which are not narrow; an example to this effect is presented in Section 6. As for notation, we denote the closed unit ball of a Banach space by $`B(X)`$ and its unit sphere by $`S(X)`$. The slice of $`B(X)`$ determined by a functional $`x^{}S(X^{})`$ and $`\epsilon >0`$ is the set $$S(x^{},\epsilon )=\{xB(X):x^{}(x)1\epsilon \}.$$ We shall repeatedly make use of the following characterisation of the Daugavet property in terms of slices or weakly open sets from \[13, Lemma 2.2\] and \[21, Lemma 2.2\] respectively. ###### Lemma 1.1. The following assertions are equivalent: * $`X`$ has the Daugavet property. * For every $`xS(X)`$, $`\epsilon >0`$ and every slice $`S`$ of $`B(X)`$ there exists some $`yS`$ such that $`x+y>2\epsilon `$. * For every $`xS(X)`$, $`\epsilon >0`$ and every nonvoid relatively weakly open subset $`U`$ of $`B(X)`$ there exists some $`yU`$ such that $`x+y>2\epsilon `$. Actually, this lemma characterises Daugavet pairs $`(Y,X)`$, meaning a Banach space $`X`$ and a subspace $`YX`$ such that $$J+T=1+T$$ for every operator from $`Y`$ into $`X`$ of rank $`1`$; here $`J`$ denotes the canonical embedding map. The only modification to be made in the formulation of Lemma 1.1 is that $`S`$ and $`U`$ refer to slices and subsets of $`B(Y)`$. Finally we mention that all the Banach spaces in this paper are tacitly assumed to be real. Acknowledgement. We are grateful to the referee whose detailed suggestions have led to a substantial improvement of the exposition of this paper. ## 2. The semigroup $`𝒪𝒫(X)`$ Throughout the paper the symbol $`X`$ will be used for a fixed Banach space, the symbols $`T`$, $`T_i`$ etc. for bounded linear operators, acting from $`X`$ to some other Banach space (not necessarily the same one). ###### Definition 2.1. We say that two operators $`T_1`$ and $`T_2`$ are equivalent (in symbols $`T_1T_2`$) if $`T_1x=T_2x`$ for every $`xX`$. A class $``$ of operators is said to be admissible if for every $`T`$ all the members of the equivalence class of $`T`$ also belong to $``$. In other words, the operators $`T_1`$ and $`T_2`$ are equivalent if there is an isometry $`U:T_1(X)T_2(X)`$ such that $`T_2=UT_1`$. For example, the classes of finite-rank operators, compact operators, weakly compact operators, operators bounded from below are admissible; surjections, isomorphisms, projections are examples of non-admissible operator classes. ###### Definition 2.2. We say that $`T_1T_2`$ if $`T_1xT_2x`$ for every $`xX`$. A class $``$ of operators forms an order ideal if for every $`T`$ every operator $`T_1T`$ also belongs to $``$. In other words, $`T_1T_2`$ if there is a bounded operator $`U:T_2(X)T_1(X)`$ of norm $`1`$ such that $`T_1=UT_2`$. Order ideals are clearly admissible. The classes of finite-rank operators, compact operators, weakly compact operators are order ideals. ###### Definition 2.3. A sequence $`(T_n)`$ of operators is said to be $``$convergent to an operator $`T`$ if $`T_nxTx`$ uniformly on $`B(X)`$. In terms of $``$convergence we define the notions of a $``$closed set of operators, $``$closure, etc. in a natural way. Of course, the $``$limit of a sequence is not unique, but it is unique up to equivalence of operators. For example, the class $`(X)`$ of finite-rank operators on an infinite-dimensional space $`X`$ is not $``$closed: its $``$closure contains all compact operators. Indeed, let $`T:XY`$ be compact. Then, for the canonical isometry $`U`$ from $`Y`$ into $`C(B(Y^{}))`$, $`T_1:=UT`$ is compact, too, and by definition $`T_1T`$. Since $`C(B(Y^{}))`$ has the approximation property, $`T_1`$ can be approximated by finite-rank operators in the above sense. In fact, the $``$closure of $`(X)`$ coincides with the class $`𝒞(X)`$ of all compact operators since $`𝒞(X)`$ is $``$closed. To see this suppose that $`(T_n)`$ is a $``$convergent sequence of compact operators on $`X`$ with limit $`T`$. Let $`(x_n)`$ be a bounded sequence in $`X`$; using a diagonal procedure one can find a subsequence $`(x_n^{})`$ such that $`(T_kx_n^{})_n^{}`$ is convergent for each $`k`$. But $`T_kxTx`$ uniformly on bounded sets as $`k\mathrm{}`$; hence $`(Tx_n^{})`$ is a Cauchy sequence and thus convergent. ###### Definition 2.4. Let $`𝒩`$ be a collection of subsets in $`X`$. We define a class of operators $`𝒩^{}`$ as follows: $`T𝒩^{}`$ if for every $`A𝒩`$, $`T`$ is unbounded from below on $`A`$; i.e., $$\epsilon >0xA:Tx\epsilon .$$ Evidently, $`𝒩^{}`$ is a $``$closed order ideal, and it is homogeneous in the sense that $`\lambda T𝒩^{}`$ whenever $`\lambda `$ and $`T𝒩^{}`$. For example, if $`𝒩=\{S(X)\}`$, then $`𝒩^{}=𝒰(X)`$, the class of operators that are unbounded from below which is defined by $$T𝒰(X)inf\{Tx:x=1\}=0.$$ A significant example for us is the class of all $`C`$-narrow operators on the space $`C(K)`$. This class was introduced in as the class of those operators $`T:C(K)Y`$ which are unbounded from below on the unit sphere of each subspace $`J_F:=\{fC(K):f_F=0\}`$, where $`F`$ is a proper closed subset of $`K`$. To put it another way, if $`𝒩`$ denotes the collection of these unit spheres, then the class of $`C`$-narrow operators is just $`𝒩^{}`$. Another important example is the class of all $`L_1`$-narrow operators on the space $`L_1=L_1(\mathrm{\Omega },\mathrm{\Sigma },\mu )`$. An operator $`T:L_1Y`$ is called $`L_1`$-narrow if for each $`B\mathrm{\Sigma }`$ and each $`\epsilon >0`$ there is a function vanishing off $`B`$ and taking only the values $`1`$ and $`1`$ on $`B`$ such that $`Tf\epsilon `$. In other words, if $`𝒩`$ now denotes the collection of these sets of functions, then the class of $`L_1`$-narrow operators is again $`𝒩^{}`$. $`L_1`$-narrow operators were formally introduced in , but the complement of this class was studied earlier by Ghoussoub and Rosenthal who called non-$`L_1`$-narrow operators norm-sign-preserving. An operator is not $`L_1`$-narrow if and only if it is not a sign-embedding on any $`L_1(B)`$-subspace (, , ). We caution the reader that in and only the term “narrow” is used. In this paper we prefer to speak of $`C`$\- and $`L_1`$-narrow operators in order not to mix up these notions with our concept of a narrow operator in Section 3; cf., however, Theorem 3.7. We now define $`𝒪𝒫(X)`$ as the class of all operators on $`X`$ with the convention that equivalent operators will be identified. Hence $`𝒪𝒫(X)`$ is actually a collection of equivalence classes, and in fact it is a set. Namely, for an operator $`T`$ on $`X`$ its equivalent class can be identified with the seminorm $`xTx`$, and the collection of seminorms on $`X`$ is clearly a set. Thus, admissible families of operators can be identified with subsets of $`𝒪𝒫(X)`$, and it makes sense to write $`T𝒪𝒫(X)`$ or $`𝒪𝒫(X)`$. We now introduce addition and subtraction on $`𝒪𝒫(X)`$. If $`T_1:XY_1`$ and $`T_2:XY_2`$ are two operators, define $$T_1\stackrel{~}{+}T_2:XY_1_1Y_2,x(T_1x,T_2x);$$ i.e., $$(T_1\stackrel{~}{+}T_2)x=T_1x+T_2x.$$ ###### Definition 2.5. If $`_1,_2𝒪𝒫(X)`$ are non-empty, then their $``$sum is defined by $`_1\stackrel{~}{+}_2=\{T_1\stackrel{~}{+}T_2:T_1_1,T_2_2\}`$. Their $``$difference is defined by $`_2\stackrel{~}{}_1=\{T𝒪𝒫(X):T\stackrel{~}{+}T_1_2`$ whenever $`T_1_1\}`$. The operation $`\stackrel{~}{+}`$ is a commutative and associative operation on $`𝒪𝒫(X)`$, and we have $`0_2\stackrel{~}{}_1`$ if and only if $`_1_2`$. Let us give some examples. ###### Example 2.6. Let $`K`$ be a compact Hausdorff space and let $`𝒰(C(K))`$ denote the class of operators equivalent to some multiplication operator $`M_h:fhf`$ on $`C(K)`$ which is unbounded from below; i.e., where $`h`$ has a zero. Then $`𝒰(C(K))\stackrel{~}{}𝒰(C(K))`$ consists exactly of the $`C`$-narrow operators described above. ###### Proof. Let $`T:C(K)Y`$ be $`C`$-narrow. If $`h`$ has a zero, we have to show that, given $`\epsilon >0`$, there is some $`fS(C(K))`$ such that both $`Tf\epsilon `$ and $`hf_{\mathrm{}}\epsilon `$. Now, if $`F=\{\|h\|\epsilon \}`$, which is a proper subset of $`K`$, and $`fS(J_F)`$ such that $`Tf\epsilon `$, then $`hf_{\mathrm{}}\epsilon `$ as well. Conversely, if a closed proper subset $`FK`$ is given, pick some $`hS(C(K))`$ such that $`h=1`$ on $`F`$, $`h=0`$ off a neighbourhood $`V`$ of $`F`$. If $`f_{\mathrm{}}1`$, $`Tf\epsilon `$ and $`hf_{\mathrm{}}\epsilon `$, then in particular $`\|f\|\epsilon `$ on $`F`$. Hence it is possible to replace $`f`$ by a function $`gS(J_F)`$ such that $`Tg2\epsilon `$, which proves that $`T`$ is $`C`$-narrow. ∎ For our next example recall that an operator on $`X`$ is a left semi-Fredholm operator if its kernel is finite-dimensional and its range is closed, and it is strictly singular if it is unbounded from below on (the unit sphere of) each infinite-dimensional subspace of $`X`$. ###### Example 2.7. The class $`𝒰(X)\stackrel{~}{}(X)`$ consists of all operators that are not left semi-Fredholm operators; $`𝒰(X)\stackrel{~}{}(𝒰(X)\stackrel{~}{}(X))`$ consists of all strictly singular operators. ###### Proof. Let us denote $`𝒢(X)=𝒰(X)\stackrel{~}{}(X)`$ and $`(X)=𝒰(X)\stackrel{~}{}𝒢(X)`$. If $`T`$ is a left semi-Fredholm operator, then, since $`\mathrm{ker}T`$ is complemented by a finite-codimensional subspace $`YX`$, $`T_Y`$ is bounded from below, because $`T`$ acts as an isomorphism from $`Y`$ onto $`T(X)`$. On the other hand, if $`T_Y`$ is bounded from below on some finite-codimensional subspace $`YX`$, then $`T(Y)`$ and, consequently, $`T(X)`$ must be closed, and $`\mathrm{ker}T`$ is finite-dimensional, since otherwise $`Y\mathrm{ker}T\{0\}`$. This shows that $`T`$ is not a left semi-Fredholm operator if and only if * $`T_Y`$ is not bounded from below on any finite-codimensional subspace $`YX`$. Now, if $`T`$ satisfies (a), $`F(X)`$ and $`Y=\mathrm{ker}F`$, then $`T\stackrel{~}{+}F𝒰(X)`$, i.e., $`T𝒢(X)`$. Conversely, if $`T𝒢(X)`$, $`YX`$ is finite-codimensional and $`q:XX/Y`$ is the quotient map, then, since $`T\stackrel{~}{+}q𝒰(X)`$, $`T`$ satisfies (a). Thus, we have shown the announced characterisation of $`𝒢(X)`$ and, moreover, we have shown that (a) provides another characterisation of $`𝒢(X)`$. It follows from (a) that $`T𝒢(X)`$ if and only if * for every $`\epsilon >0`$ there exists an infinite-dimensional subspace $`ZX`$ such that $`T_Z\epsilon `$; see \[14, Prop. 2.c.4\]. From (b) it is clear that every strictly singular operator belongs to $`(X)`$. On the other hand, if $`S`$ is not strictly singular and is bounded from below on some infinite-dimensional subspace $`Z`$, then we have for the quotient map $`q:XX/Z`$ that $`S\stackrel{~}{+}q`$ is bounded from below. Since $`q`$ obviously satisfies (b), this shows that $`S(X)`$. ∎ Let us list some elementary properties of the operation $`\stackrel{~}{}`$ that follow directly from the definition. ###### Proposition 2.8. Suppose that $`_1,_2𝒪𝒫(X)`$ contain the zero operator. * $`_2\stackrel{~}{}_1`$ is an order ideal or $``$closed whenever $`_2`$ is. * If $`_1`$ and $`_2`$ are order ideals, then $`_2\stackrel{~}{}_1`$ is homogeneous whenever $`_2`$ is. Of particular relevance are subsets of $`𝒪𝒫(X)`$ that are semigroups with respect to the operation $`\stackrel{~}{+}`$. ###### Proposition 2.9. Suppose that $`_1,_2𝒪𝒫(X)`$ contain the zero operator. * $`_1`$ is a subsemigroup of $`𝒪𝒫(X)`$ if and only if $`_1\stackrel{~}{}_1_1`$, in which case $`_1\stackrel{~}{}_1=_1`$. * Let $`_1`$ be a subsemigroup of $`𝒪𝒫(X)`$, and let $`_1_2`$. Then $`_2\stackrel{~}{}(_2\stackrel{~}{}_1)`$ is again a subsemigroup. * $`_2\stackrel{~}{}_2`$ is always a subsemigroup of $`𝒪𝒫(X)`$. ###### Proof. (a) is clear from the definition. For (b) we note first that (2.1) $$_2\stackrel{~}{}(_2\stackrel{~}{}(_2_1))=_2_1.$$ Indeed, by definition of $`\stackrel{~}{}`$ we have (2.2) $$_2\stackrel{~}{}(_2\stackrel{~}{}_1)_1,$$ whence $$_2\stackrel{~}{}(_2\stackrel{~}{}(_2\stackrel{~}{}_1))_2\stackrel{~}{}_1.$$ On the other hand, an application of (2.2) with $`_1`$ replaced by $`_2\stackrel{~}{}_1`$ gives “$``$” in (2.1). Now, by elementary arithmetic involving $`\stackrel{~}{+}`$ and $`\stackrel{~}{}`$ we have, writing $`𝒟=_2\stackrel{~}{}_1`$ for short, $`(_2\stackrel{~}{}𝒟)\stackrel{~}{}(_2\stackrel{~}{}𝒟)`$ $`=_2\stackrel{~}{}(𝒟\stackrel{~}{+}(_2\stackrel{~}{}𝒟))`$ $`=_2\stackrel{~}{}((_2\stackrel{~}{}𝒟)\stackrel{~}{+}𝒟)`$ $`=(_2\stackrel{~}{}(_2\stackrel{~}{}𝒟))\stackrel{~}{}𝒟`$ $`=𝒟\stackrel{~}{}𝒟\text{(by (}\text{2.1}\text{))}`$ $`=(_2\stackrel{~}{}_1)\stackrel{~}{}𝒟`$ $`=_2\stackrel{~}{}(_1\stackrel{~}{+}𝒟).`$ Because $`_1`$ is a semigroup, one can easily deduce that $`_1\stackrel{~}{+}𝒟𝒟`$; indeed, $`_1\stackrel{~}{+}𝒟`$ $`=(_2\stackrel{~}{}_1)\stackrel{~}{+}_1`$ $`=(_2\stackrel{~}{}(_1\stackrel{~}{+}_1))\stackrel{~}{+}_1`$ $`=((_2\stackrel{~}{}_1)\stackrel{~}{}_1)\stackrel{~}{+}_1`$ $`_2\stackrel{~}{}_1.`$ Therefore $$(_2\stackrel{~}{}𝒟)\stackrel{~}{}(_2\stackrel{~}{}𝒟)_2\stackrel{~}{}𝒟,$$ completing the proof that $`_2\stackrel{~}{}(_2\stackrel{~}{}_1)`$ is a semigroup. Finally, (c) is the special case $`_1=\{0\}`$ of (b). ∎ The following definition is important for our abstract semigroup approach. ###### Definition 2.10. Let $`𝒪𝒫(X)`$, and let $`_1`$ be a subsemigroup of $`𝒪𝒫(X)`$. $`_1`$ is called a maximal subsemigroup of $``$ if every subsemigroup $`_2`$ which includes $`_1`$ coincides with $`_1`$. We call the intersection of all maximal subsemigroups of $``$ the central part of $``$ and denote it by $`cp()`$. Here is a characterisation of the central part of $``$. ###### Theorem 2.11. Let $`𝒪𝒫(X)`$ have the following properties: $`0`$ and every element of $``$ is contained in a subsemigroup of $``$ (this happens for example if $``$ is homogeneous). Then $`cp()=\stackrel{~}{}`$. ###### Proof. Let $`_1`$ be a maximal subsemigroup of $``$. Put $`_2=\stackrel{~}{}`$. We have proved above in Proposition 2.9(c) that $`_2`$ is a subsemigroup, so $`_2\stackrel{~}{+}_1`$ is a subsemigroup, too. By definition of $`_2`$ we have $`_2\stackrel{~}{+}_1`$. So the maximality of $`_1`$ implies that $`_1_2`$. This proves the inclusion $`cp()\stackrel{~}{}`$. Let us now prove the inverse inclusion. Let $`Tcp()(\stackrel{~}{})`$. Then there is some $`T_1`$ such that $`T_1\stackrel{~}{+}T`$ does not belong to $``$. Consider the maximal subsemigroup $`_3`$ of $``$ which contains $`T_1`$. Then $`_3`$ cannot contain $`T`$, so $`cp()`$ cannot contain $`T`$ either. ∎ For every operator $`T`$ and $`\epsilon >0`$ we define the tube $$U_{T,\epsilon }=\{xX:Tx<\epsilon \}.$$ Let $`𝒪𝒫(X)`$. Put $$_{}=\{U_{T,\epsilon }S(X):T,\epsilon >0\}$$ Then $`(_{})^{}=𝒰(X)\stackrel{~}{}`$. ###### Proposition 2.12. Let $`𝒪𝒫(X)`$ and let $`𝒩`$ be a collection of subsets in $`X`$. Then $`𝒩^{}\stackrel{~}{}=𝒩_{\mathcal{1}}^{}{}_{}{}^{}`$, where $`𝒩_\mathcal{1}`$ consists of all intersections of the form $`U_{T,\epsilon }A`$, $`T`$, $`A𝒩`$, $`\epsilon >0`$. In particular, if $`𝒩^{}\stackrel{~}{}`$ is non-empty, then all the intersections $`U_{T,\epsilon }A`$ are non-empty and $`𝒩^{}`$. ###### Proof. Let $`T_1𝒩^{}\stackrel{~}{}`$. Then for every $`T`$ we have $`T_1\stackrel{~}{+}T𝒩^{}`$. This means that for every $`A𝒩`$ and $`\epsilon >0`$ there is an element $`xA`$ such that $`(T_1\stackrel{~}{+}T)x<\epsilon `$. This in turn implies that $`xAU_{T,\epsilon }`$ and $`T_1x<\epsilon `$. So $`T_1𝒩_{\mathcal{1}}^{}{}_{}{}^{}`$. Now let $`T_1𝒩_{\mathcal{1}}^{}{}_{}{}^{}`$. Then for every $`T`$, every $`A𝒩`$ and $`\epsilon >0`$ there is an element $`xAU_{T,\epsilon /2}`$ such that $`T_1x<\epsilon /2`$. But by the definition of tubes, $`Tx<\epsilon /2`$. So $`(T_1\stackrel{~}{+}T)x<\epsilon `$ and $`T_1𝒩^{}\stackrel{~}{}`$. ∎ ## 3. Narrow operators In this section we define the class of narrow operators on a Banach space with the Daugavet property. But first we need to introduce a closely related class of operators. ###### Definition 3.1. An operator $`T𝒪𝒫(X)`$ is said to be a strong Daugavet operator if for every two elements $`x,yS(X)`$ and for every $`\epsilon >0`$ there is an element $`z(y+U_{T,\epsilon })S(X)`$ such that $`z+x>2\epsilon `$. We denote the class of all strong Daugavet operators on $`X`$ by $`𝒮𝒟(X)`$. It follows from Lemma 1.1 that a finite-rank operator on a space with the Daugavet property is a strong Daugavet operator, and conversely, if every rank-$`1`$ operator is strongly Daugavet, then $`X`$ has the Daugavet property. There is an obvious connection between strong Daugavet operators and the Daugavet equation. ###### Lemma 3.2. If $`T:XX`$ is a strong Daugavet operator, then $`T`$ satisfies the Daugavet equation (1.1). ###### Proof. We assume without loss of generality that $`T=1`$. Given $`\epsilon >0`$ pick $`yS(X)`$ such that $`Ty1\epsilon `$. If $`x=Ty/Ty`$ and $`z`$ is chosen according to Definition 3.1, then $$2\epsilon <z+xz+Ty+\epsilon z+Tz+2\epsilon ,$$ hence $$z+Tz23\epsilon ,$$ which proves the lemma. ∎ We now relate the strong Daugavet property to a collection of subsets of $`X`$. ###### Definition 3.3. For every ordered pair of elements $`(x,y)`$ of $`S(X)`$ and every $`\epsilon >0`$ let us define a set $`D(x,y,\epsilon )`$ by $$zD(x,y,\epsilon )z+x+y>2\epsilon \&z+y<1+\epsilon .$$ By $`𝒟(X)`$ we denote the collection of all sets $`D(x,y,\epsilon )`$, where $`x,yS(X)`$ and $`\epsilon >0`$. ###### Proposition 3.4. $`𝒮𝒟(X)=𝒟(X)^{}`$. ###### Proof. $`T𝒟(X)^{}`$ if and only if for every pair $`x,yS(X)`$ and $`\epsilon >0`$ there is an element $`zD(x,y,\epsilon )`$ such that $`Tz<\epsilon `$. This in turn is equivalent to the following condition: for every pair $`x,yS(X)`$ and $`\epsilon >0`$ there is an element $`v`$ such that $`v<1+\epsilon `$, $`x+v>2\epsilon `$ and $`v`$ belongs to the tube $`y+U_{T,\epsilon }`$ (just put $`v=z+y`$). Evidently, the last equation coincides with the strong Daugavet property of the operator $`T`$. ∎ Let us consider an example. ###### Theorem 3.5. For a compact Hausdorff space $`K`$, the class $`𝒮𝒟(C(K))`$ of strong Daugavet operators coincides with the class of $`C`$-narrow operators. ###### Proof. The fact that every $`C`$-narrow operator is a strong Daugavet operator has been proved in a slightly different form in \[12, Th. 3.2\]. Consider the converse implication. Let $`T𝒮𝒟(C(K))`$. Fix a closed subset $`FK`$ and $`0<\epsilon <1/4`$. According to the definition it is sufficient to prove that there is a function $`fS(C(K))`$ for which the restriction to $`F`$ is less than $`2\epsilon `$ and $`Tf<2\epsilon `$ (cf. Example 2.6). Let us fix a neighbourhood $`U`$ of $`F`$ and an open set $`VK`$, $`VU=\mathrm{}`$. Select inductively functions $`x_n,y_nS(C(K))`$ and $`f_n,g_nC(K)`$ as follows. All the $`y_n`$ are supported on $`U`$, and the $`x_n`$ are non-negative functions supported on $`V`$. Given $`x_n`$ and $`y_n`$ pick $`f_nD(x_n,y_n,\epsilon )`$ with $`Tf_n<\epsilon `$, and let $`g_n=f_1+f_2+\mathrm{}+f_n`$. Then choose $`y_{n+1}S(C(K))`$ subject to the above support condition such that $`sup_{tF}\|g_n(t)\|y_{n+1}`$ coincides on $`F`$ with $`g_n`$, and let $`x_{n+1}`$ be a non-negative continuous function supported on the subset of $`V`$ where $`g_n`$ attains its supremum on $`V`$ up to $`\epsilon `$, i.e., on the set $`\{tV:g_n(t)>sup_{sV}g_n(s)\epsilon \}`$, etc. (There is no initial restriction in the choice of $`y_1`$ and $`x_1`$ apart from the support and positivity conditions.) We first claim that $$g_n_F:=\underset{tF}{sup}\|g_n(t)\|3+n\epsilon .$$ This is certainly true for $`n=1`$ since $`f_1+y_1<1+\epsilon `$. Now induction yields for $`tF`$ $`\|g_{n+1}(t)\|`$ $`=\|g_n(t)+f_{n+1}(t)\|`$ $`=\left\|g_n_Fy_{n+1}(t)+f_{n+1}(t)\right\|`$ $`=\left\|y_{n+1}(t)+f_{n+1}(t)+(g_n_F1)y_{n+1}(t)\right\|`$ $`\|y_{n+1}(t)+f_{n+1}(t)\|+g_n_F1`$ $`1+\epsilon +2+n\epsilon =3+(n+1)\epsilon .`$ (We have tacitly assumed that $`g_n_F1`$ since the induction step is clear otherwise, because $`f_{n+1}2+\epsilon `$.) Next, we have that $$\underset{tV}{sup}g_n(t)>n(12\epsilon ).$$ Indeed, the functions $`x_1`$ and $`y_1`$ are disjointly supported; hence by the definition of $`D(x_1,y_1,\epsilon )`$ there is a point in the support of $`x_1`$ at which $`f_1=g_1`$ is bigger than $`1\epsilon `$. Thus, the above inequality holds for $`n=1`$. To perform the induction step we use the same argument to find a point $`t_0`$ in the support of $`x_{n+1}`$ at which $`f_{n+1}`$ exceeds $`1\epsilon `$. At this point $`t_0`$ the function $`g_n`$ attains its supremum on $`V`$ up to $`\epsilon `$. So $`\underset{tV}{sup}g_{n+1}(t)`$ $`g_{n+1}(t_0)=g_n(t_0)+f_{n+1}(t_0)>\underset{tV}{sup}g_n(t)+12\epsilon `$ $`>n(12\epsilon )+12\epsilon =(n+1)(12\epsilon ).`$ Therefore $`g_{n+1}>(n+1)(12\epsilon )`$, and on the other hand we have $`Tg_n_{k=1}^nTf_n<n\epsilon `$. So for $`n`$ big enough the function $`f=g_n/g_n`$ will satisfy the desired conditions. ∎ Actually, a somewhat smaller class of operators turns out to be crucial. ###### Definition 3.6. Let $`X`$ be a space with the Daugavet property. Define the class of narrow operators by $`𝒩𝒜(X)=𝒮𝒟(X)\stackrel{~}{}X^{}`$. In other words, an operator $`T`$ is said to be a narrow operator if, for every $`x^{}X^{}`$, $`T\stackrel{~}{+}x^{}`$ is a strong Daugavet operator. Incidentally, it follows from the defining property of a narrow operator that a Banach space on which at least one narrow operator is defined must automatically have the Daugavet property. Proposition 3.4 and Proposition 2.8 imply that $`𝒩𝒜(X)`$ is a $``$closed homogeneous order ideal, and hence so is $`cp(𝒩𝒜(X))`$. We now show that we won’t get anything new on $`C(K)`$ if $`K`$ is perfect. ###### Theorem 3.7. For a perfect compact Hausdorff space $`K`$, the classes of $`C`$-narrow operators and of narrow operators coincide on $`C(K)`$. ###### Proof. Since a narrow operator is a strong Daugavet operator and a strong Daugavet operator is $`C`$-narrow (Theorem 3.5), it is left to prove that a $`C`$-narrow operator $`T`$ on $`C(K)`$ is narrow if $`K`$ is perfect. Let $`x^{}C(K)^{}`$ be a functional, represented by a regular Borel measure $`\mu `$; we have to show that $`T\stackrel{~}{+}x^{}`$ is a strong Daugavet operator. Thus, let $`f,gS(C(K))`$ and $`\epsilon >0`$. Let $`\epsilon ^{}=\epsilon /(4+T)`$, and consider the open set $`U=\{t:\|f(t)\|>1\epsilon ^{}\}`$. Pick an open non-empty subset $`VU`$ with the property that $`fg`$ is almost constant on $`V`$ in that for some number $`c[2,2]`$ $$\|f(t)g(t)c\|\epsilon ^{}\text{ for }tV,$$ and $`\|\mu \|(V)\epsilon ^{}`$; the latter is possible since $`K`$ has no isolated points. Since $`T`$ is $`C`$-narrow, there is some $`\phi S(C(K))`$ vanishing off $`V`$ such that $`T\phi \epsilon ^{}`$; in fact, $`\phi `$ can (and will) be chosen positive \[12, Lemma 1.4\]. Let $`h=\phi f+(1\phi )g`$. Then $`h1`$, and $$f+h\underset{tV}{sup}\|f(t)+h(t)\|22\epsilon ^{}2\epsilon ;$$ furthermore $`\|x^{}(g)x^{}(h)\|`$ $`=\|x^{}(\phi (fg))\|\|\mu (V)\|\phi fg2\epsilon ^{}`$ $`T(g)T(h)`$ $`=T(\phi (fg))T\phi (fgc)+T\phi \|c\|`$ $`T\epsilon ^{}+2\epsilon ^{}`$ so that $$(T\stackrel{~}{+}x^{})(gh)\epsilon ,$$ which proves that $`T\stackrel{~}{+}x^{}`$ is a strong Daugavet operator. ∎ We shall show below (Section 6) that in general, on a space with the Daugavet property narrow and strong Daugavet operators are not the same. We don’t know if in general $`𝒩𝒜(X)`$ is a subsemigroup of $`𝒪𝒫(X)`$ as it will be shown to be the case for $`X=C(K)`$ (Theorem 4.8), but we will show that its central part $`cp(𝒩𝒜(X))`$ is always large. It contains, in particular, all strong Radon-Nikodým operators and all operators which do not fix copies of $`\mathrm{}_1`$. Hence all the operators which are majorized by linear combinations of strong Radon-Nikodým operators and operators not fixing copies of $`\mathrm{}_1`$, as well as $``$-limits of sequences of such operators belong to $`cp(𝒩𝒜(X))`$. We now formulate a number of lemmas. Eventually, Proposition 3.11 will present a geometric description of what distinguishes a narrow operator from a strong Daugavet operator. ###### Lemma 3.8. Let $`T𝒩𝒜(X)`$. Then for every $`x,yS(X)`$, $`\epsilon >0`$ and every slice $`S=S(x^{},\alpha )`$ of the unit ball of $`X`$ containing $`y`$ there is an element $`vS`$ such that $`x+v>2\epsilon `$ and $`T(yv)<\epsilon `$. ###### Proof. Fix some $`0<\delta <\epsilon `$ and find an element $`y_1`$ in norm-interior of $`S`$ such that $`yy_1<\delta `$. By Proposition 2.12, for every $`0<\delta _1<\epsilon `$ there is an element $`uU_{x^{},\delta _1}D(x,y_1/y_1,\delta _1)`$ such that $`Tu<\delta _1`$. If $`\delta _1`$ is small enough, then $`v:=(y_1+y_1u)/(y_1+y_1u)S`$. So, if in turn $`\delta `$ is small enough, then $`v`$ satisfies our requirements. ∎ ###### Lemma 3.9. For every $`\tau >0`$ and every pair of positive numbers $`a,b`$ there is a $`\delta >0`$ such that if $`v,xS(X)`$ and $`x+v>2\delta `$, then $`ax+bv>a+b\tau `$. ###### Proof. Select $`\delta <\tau /\mathrm{max}(a,b)`$. There are two symmetric cases: $`ab`$ or $`ba`$. Consider for example the first of them. If we assume that our statement is not true, then we obtain $`2\delta `$ $`<x+v=(1a/b)x+1/b(ax+bv)`$ $`1a/b+1/b(a+b\tau )=2\tau /b,`$ a contradiction. ∎ ###### Lemma 3.10. Let $`T𝒩𝒜(X)`$. * Let $`S_1,\mathrm{},S_n`$ be a finite collection of slices and $`UB(X)`$ be a convex combination of these slices, i.e., there are $`\lambda _k0`$, $`k=1,\mathrm{},n`$, $`_{k=1}^n\lambda _k=1`$, such that $`\lambda _1S_1+\mathrm{}+\lambda _nS_n=U`$. Then for every $`\epsilon >0`$, every $`x_1S(X)`$ and every $`wU`$ there exists an element $`uU`$ such that $`u+x_1>2\epsilon `$ and $`T(wu)<\epsilon `$. * The same conclusion is true if $`U`$ is a relatively weakly open set. ###### Proof. (a) First of all let us fix elements $`y_kS_k`$ such that $`\lambda _1y_1+\mathrm{}+\lambda _ny_n=w`$. Applying repeatedly Lemma 3.8 with sufficiently small $`\epsilon _j`$ to $`S_j`$, $`y_jS_j`$ and $$x_j=\left(x_1+\underset{k=1}{\overset{j1}{}}\lambda _kv_k\right)/x_1+\underset{k=1}{\overset{j1}{}}\lambda _kv_k,$$ we may select elements $`v_kS_k`$ with $`T(y_kv_k)<\epsilon `$, $`k=1,\mathrm{},n`$, in such a way that for every $`j=1,\mathrm{},n`$ $$x_1+\underset{k=1}{\overset{j}{}}\lambda _kv_k>1+\underset{k=1}{\overset{j}{}}\lambda _k(1\epsilon )$$ (to get the last inequality, we need to apply the previous lemma at each step). The element $`u=\lambda _1v_1+\lambda _2v_2+\mathrm{}+\lambda _nv_n`$ will be as required. (b) This follows from (a) since given $`uU`$ there is a convex combination $`V`$ of slices such that $`uVU`$; see \[8, Lemma II.1\] or . ∎ ###### Proposition 3.11. An operator $`T`$ on a Banach space $`X`$ with the Daugavet property is narrow if and only if for every $`x,yS(X)`$, $`\epsilon >0`$ and every slice $`S`$ of the unit ball of $`X`$ containing $`y`$ there is an element $`vS`$ such that $`x+v>2\epsilon `$ and $`T(yv)<\epsilon `$. ###### Proof. It only remains to show that the above condition is sufficient for $`T`$ to be narrow. We first note that an operator satisfying that condition will also satisfy the conclusion of Lemma 3.10; see the proof of that lemma. Now, if $`x_0^{}X^{}`$, $`x,yS(X)`$ and $`\epsilon >0`$, consider the relatively weakly open set $$U:=\{zB(X):\|x_0^{}(zy)\|<\epsilon /2\}.$$ By Lemma 3.10 there exists some $`wU`$ such that $`w+x>2\epsilon /2`$ and $`T(wy)<\epsilon /2`$; note that $`yU`$. By definition this means that $`T\stackrel{~}{+}x_0^{}`$ is a strong Daugavet operator; i.e., $`T`$ is narrow. ∎ Let $`T`$ be a strong Radon-Nikodým operator on a space $`X`$ with the Daugavet property; this means that the closure of $`T(B(X))`$ is a set with the Radon-Nikodým property, cf. for this notion. We shall show that such an operator is narrow. For $`\epsilon >0`$, consider the subset $`A(T,\epsilon )`$ of $`B(X)`$ defined by $`yA(T,\epsilon )`$ if there exists a convex combination $`U`$ of slices of the unit ball such that $`yU`$ and $`Uy+U_{T,\epsilon }`$. ###### Lemma 3.12. The set $`A(T,\epsilon )`$ introduced above is a convex dense subset of $`B(X)`$. ###### Proof. The convexity is evident. To prove the density we need to show, by the Hahn-Banach theorem, that for every $`x^{}S(X^{})`$ and every $`0<\delta <\epsilon `$ there is an element $`yA(T,\epsilon )`$ such that $`x^{}(y)>1\delta `$ (in other words, $`yS=S(x^{},\delta )`$). Let us fix an element $`xB(X)`$ with $`x^{}(x)>1\delta /2`$ and consider the operator $`T_1=x^{}\stackrel{~}{+}T`$. Consider further $`\overline{T_1(B(X))}`$ and a $`\delta /2`$-neighbourhood $`W`$ of $`T_1x`$ in $`\overline{T_1(B(X))}`$. By the Radon-Nikodým property of the set $`\overline{T_1(B(X))}`$ there is a convex combination $`W_1`$ of slices of $`\overline{T_1(B(X))}`$ in $`W`$. The preimages in $`B(X)`$ of these slices of $`\overline{T_1(B(X))}`$ are slices in $`B(X)`$. The corresponding convex combination $`U`$ of these slices in $`B(X)`$ lies in the preimage of $`W`$ in $`B(X)`$, so this convex combination is contained in $`(x+U_{T_1,\delta /2})B(X)`$. Fix an element $`yU`$. By our construction $`yU(x+U_{T_1,\delta /2})B(X)S`$. On the other hand, $$Ux+U_{T_1,\delta /2}y+U_{T_1,\delta }y+U_{T,\delta }y+U_{T,\epsilon },$$ so $`yA(T,\epsilon )`$. ∎ The following result is a generalisation of \[13, Th. 2.3\]. It can be understood as a transfer theorem: in Definition 3.6 one can pass from one-dimensional operators to a much wider class of operators. Let us denote the class of strong Radon-Nikodým operators on $`X`$ by $`𝒮𝒩(X)`$. ###### Theorem 3.13. Let $`X`$ be a space with the Daugavet property, $`T`$ be narrow and $`T_1`$ be a strong Radon-Nikodým operator on $`X`$. Then $`T\stackrel{~}{+}T_1`$ is narrow; that is, we have $`𝒩𝒜(X)\stackrel{~}{+}𝒮𝒩(X)=𝒩𝒜(X)`$. In particular every strong Radon-Nikodým operator $`T_1`$ on $`X`$ is a narrow operator. ###### Proof. Let us fix $`\epsilon >0`$, $`x,yS(X)`$ and $`y_1A(T_1,\epsilon )`$ such that $`yy_1<\epsilon `$. According to the definition of $`A(T_1,\epsilon )`$ there exists a convex combination $`U`$ of slices of the unit ball such that $`y_1U`$ and $`Uy_1+U_{T_1,\epsilon }`$. By Lemma 3.10 there is an element $`zU`$ such that $`z+x>2\epsilon `$ and $`T(y_1z)<\epsilon `$. But the inclusion $`zy_1+U_{T_1,\epsilon }`$ means that $`T_1(y_1z)<\epsilon `$. So $$(T\stackrel{~}{+}T_1)(yz)<\epsilon T\stackrel{~}{+}T_1+(T\stackrel{~}{+}T_1)(y_1z)<\epsilon T\stackrel{~}{+}T_1+2\epsilon .$$ Because $`\epsilon `$ is arbitrarily small, the last inequality shows that $`T\stackrel{~}{+}T_1`$ satisfies the definition of a strong Daugavet operator. Now let $`x^{}X^{}`$ and consider $`T_2=T_1\stackrel{~}{+}x^{}`$. This is a strong Radon-Nikodým operator, too. So $`(T\stackrel{~}{+}T_1)\stackrel{~}{+}x^{}=T\stackrel{~}{+}T_2`$ is a strong Daugavet operator by what we have just proved; by definition, this says that $`T\stackrel{~}{+}T_1`$ is narrow. ∎ ###### Corollary 3.14. Let $`X`$ be a Banach space with the Daugavet property. * $`𝒩𝒜(X)\stackrel{~}{+}X^{}=𝒩𝒜(X)`$. * $`cp(𝒩𝒜(X))=𝒮𝒟(X)\stackrel{~}{}𝒩𝒜(X)`$. * $`𝒮𝒩(X)cp(𝒩𝒜(X))`$. ###### Proof. (a) follows from the previous theorem, because every finite-rank operator is a strong Radon-Nikodým operator. For (b) use Theorem 2.11 and note that $`𝒮𝒟(X)\stackrel{~}{}𝒩𝒜(X)`$ $`=𝒮𝒟(X)\stackrel{~}{}(𝒩𝒜(X)\stackrel{~}{+}X^{})`$ $`=(𝒮𝒟(X)\stackrel{~}{}X^{})\stackrel{~}{}𝒩𝒜(X)`$ $`=𝒩𝒜(X)\stackrel{~}{}𝒩𝒜(X).`$ (c) is a restatement of Theorem 3.13. ∎ ## 4. Operators which do not fix copies of $`\mathrm{}_1`$ It is proved in that an operator $`T:XX`$ on a space with the Daugavet property which does not fix a copy of $`\mathrm{}_1`$ satisfies the Daugavet equation. Recall that $`T𝒪𝒫(X)`$ does not fix a copy of $`\mathrm{}_1`$ if there is no subspace $`EX`$ isomorphic to $`\mathrm{}_1`$ on which the restriction $`T:ET(E)`$ is an isomorphism. By Rosenthal’s $`\mathrm{}_1`$-theorem, this is equivalent to saying that for every bounded sequence $`(x_n)X`$, the sequence of images $`(Tx_n)`$ admits a weak Cauchy subsequence. We shall investigate the class of operators not fixing a copy of $`\mathrm{}_1`$ in the present context. We will use the following theorem, due to H.P. Rosenthal : ###### Theorem 4.1. Let $`X`$ be a separable Banach space without $`\mathrm{}_1`$-subspaces. If $`AX`$ is bounded and $`x^{}X^{}`$ is a weak$`^{}`$ limit point of $`A`$, then there is a sequence in $`A`$ which converges to $`x^{}`$ in the weak$`^{}`$ topology of $`X^{}`$. In fact, we shall need a generalization of this result and first provide a lemma. ###### Lemma 4.2. Let $`X`$ be a Banach space without subspaces isomorphic to $`\mathrm{}_1`$, and let $`\{x_{n,m}\}_{n,m}X`$ be a bounded double sequence. Let $`x^{}X^{}`$ be a $`\sigma (X^{},X^{})`$-limit point of every column $`\{x_{n,m}\}_n`$ of $`\{x_{n,m}\}_{n,m}`$. Then there are strictly increasing sequences $`(n(k))`$, $`(m(k))`$ of indices such that $`x_{n(k),m(k)}x^{}`$ in $`\sigma (X^{},X^{})`$. ###### Proof. Consider an auxiliary space $`Y=X\times `$ and an auxiliary matrix $`\{y_{n,m}\}_{n,m}Y`$, $`y_{n,m}=(x_{n,m},1/n+1/m)`$. Since $`Y`$ contains no copies of $`\mathrm{}_1`$ either and since $`(x^{},0)`$ is a $`\sigma (Y^{},Y^{})`$-limit point of $`\{y_{n,m}\}_{n,m}`$, there is, according to Theorem 4.1, a sequence of the form $`(y_{n(k),m(k)})`$ which converges to $`(x^{},0)`$ in $`\sigma (Y^{},Y^{})`$. This means in particular that $`x_{n(k),m(k)}x^{}`$ in $`\sigma (X^{},X^{})`$ and $`1/n+1/m0`$. So $`(n(k))`$ and $`(m(k))`$ both tend to $`\mathrm{}`$, which, after passing to a subsequence, provides the desired sequence. ∎ The next result is a direct generalisation of Theorem 4.1. ###### Theorem 4.3. Let $`X`$ be a separable Banach space without $`\mathrm{}_1`$-subspaces, $`(\mathrm{\Gamma },)`$ be a directed set, and let $`F:\mathrm{\Gamma }X`$ be a bounded function. Then for every $`\sigma (X^{},X^{})`$-limit point $`x^{}`$ of the function $`F`$ there is a strictly increasing sequence $`\gamma (1)\gamma (2)\mathrm{}`$ in $`\mathrm{\Gamma }`$ such that $`\left(F(\gamma (n))\right)`$ converges to $`x^{}`$ in $`\sigma (X^{},X^{})`$. ###### Proof. Using inductively Theorem 4.1 we can select a doubly indexed sequence $`\{\gamma _{n,m}\}_{n,m}`$ in $`\mathrm{\Gamma }`$ with the following properties: 1. for every $`m`$, $`x^{}X^{}`$ is a $`\sigma (X^{},X^{})`$-limit point of every column $`\{F(\gamma _{n,m})\}_n`$; 2. for every $`m,n,k,l`$, if $`\mathrm{max}\{k,l\}<m`$, then $`\gamma _{k,l}\gamma _{n,m}`$. Applying Lemma 4.2 and passing to a subsequence if necessary, we obtain strictly increasing sequences $`(n(k))`$, $`(m(k))`$ such that $`\mathrm{max}_{k<j}\{n(k),m(k)\}<m(j)`$ and $`(F(\gamma _{n(k),m(k)})`$ converges to $`x^{}`$ in $`\sigma (X^{},X^{})`$. To finish the proof put $`\gamma (k)=\gamma _{n(k),m(k)}`$. ∎ We wish to prove that an operator not fixing a copy of $`\mathrm{}_1`$ is narrow (Theorem 4.13 below). To cover the case of non-separable spaces as well we first show that the Daugavet property is separably determined. The next lemma prepares this result. ###### Lemma 4.4. Let $`X`$ be a Banach space with the Daugavet property. Then for any $`\epsilon >0`$ and $`x,yS(X)`$, there exists a finite-dimensional subspace $`Y=Y(x,y,\epsilon )`$ of $`X`$ with $`x,yY`$ such that for every slice $`S(x^{},\epsilon /2)`$ containing $`y`$ there is some $`y_1S(Y)S(x^{},\epsilon )`$ such that $`y_1+x>2\epsilon `$. ###### Proof. Assume there exist $`\epsilon >0`$ and $`x,yS(X)`$ such that for every finite-dimensional subspace $`YX`$ there is a slice $`S(x_Y^{},\epsilon /2)`$ containing $`y`$ with $`y_1+x2\epsilon `$ for all $`y_1S(Y)S(x_Y^{},\epsilon )`$. Take a weak cluster point $`x^{}`$ of the net $`(x_Y^{})`$ and let $`x_0^{}=x^{}/x^{}`$. We have $`x^{}(y)1\epsilon /2`$ since $`yS(x_Y^{},\epsilon /2)`$ and therefore $`x^{}1\epsilon /2`$. Now if $`y_1S(x_0^{},\epsilon /2)`$, then $`x^{}(y_1)x^{}(1\epsilon /2)>1\epsilon `$ and therefore $`x_{Y_1}^{}(y_1)>1\epsilon `$ for some $`Y_1`$ that contains $`y_1`$. So by assumption $`y_1+x2\epsilon `$, which contradicts Lemma 1.1 when applied to the slice $`S(x_0^{},\epsilon /2)`$. ∎ ###### Theorem 4.5. A Banach space $`X`$ has the Daugavet property if and only if for every separable subspace $`YX`$ there is a separable subspace $`ZX`$ which contains $`Y`$ and has the Daugavet property. ###### Proof. Suppose $`X`$ has the Daugavet property. Let $`(v_n)`$ be a dense sequence in $`Y`$. We select a sequence $`V_1V_2\mathrm{}`$ of finite-dimensional subspaces of $`X`$ by the following inductive procedure. Put $`V_1=linv_1`$. Suppose $`V_n`$ has already been constructed. Fix a $`2^n`$-net $`(x_k^n,y_k^n)`$, $`k=1,\mathrm{},N_n`$, in $`S(V_n)\times S(V_n)`$ provided with the sum norm, select by Lemma 4.4 finite-dimensional subspaces $`Y_k=Y(x_k^n,y_k^n,\epsilon )`$, $`k=1,\mathrm{},N_n`$, for $`\epsilon =2^n`$ and define $`V_{n+1}=lin(\{v_{n+1}\}Y_1\mathrm{}Y_{N_n})`$. If $`Z`$ is defined to be the closure of the union of all the $`V_n`$, then $`YZ`$ and $`Z`$ has the Daugavet property by Lemma 1.1. Conversely, let $`xS(X)`$, $`\epsilon >0`$ and let $`SB(X)`$ be a slice. Fix a point $`zS`$. If $`Z`$ is a separable subspace with the Daugavet property containing $`x`$ and $`z`$, then by Lemma 1.1 there exists some $`ySZ`$ such that $`y+x>2\epsilon `$. Again by Lemma 1.1 this shows that $`X`$ has the Daugavet property. ∎ We shall need the operator version of this theorem, which is based on the following lemma. The proofs of Lemma 4.6 and Theorem 4.7 are virtually the same as those of Lemma 4.4 and Theorem 4.5 (one uses Proposition 3.11). ###### Lemma 4.6. Let $`X`$ be a Banach space with the Daugavet property and let $`T`$ be a narrow operator on $`X`$. Then for any $`\epsilon >0`$ and $`x,yS(X)`$, there exists a finite-dimensional subspace $`Y=Y(x,y,\epsilon )`$ of $`X`$ with $`x,yY`$ such that for every slice $`S(x^{},\epsilon /2)`$ containing $`y`$ there is some $`y_1S(Y)S(x^{},\epsilon )`$ with $`Ty_1Ty<\epsilon `$ such that $`y_1+x>2\epsilon `$. ###### Theorem 4.7. An operator $`T`$ on a Banach space $`X`$ is narrow if and only if for every separable subspace $`Y`$ of $`X`$ there is a separable subspace $`ZX`$ containing $`Y`$ such that the restriction of $`T`$ to $`Z`$ is narrow. This theorem leads to an important structural result on narrow operators on $`C(K)`$. ###### Theorem 4.8. If $`K`$ is a perfect compact Hausdorff space, then $`𝒩𝒜(C(K))`$ is a subsemigroup of $`𝒪𝒫(C(K))`$. ###### Proof. First, let $`K`$ be a perfect compact metric space. It follows from Theorem 3.7 and \[12, Th. 1.8\] that the set of all narrow operators on $`C(K)`$ is stable under the operation $`\stackrel{~}{+}`$, i.e., it is a semigroup. (In fact, only deals with $`K=[0,1]`$, but the arguments work as well for a metric $`K`$.) We shall now reduce the general case to the metric one. Let now $`K`$ be a perfect compact Hausdorff space, and let $`T_1`$ and $`T_2`$ be two narrow operators on $`C(K)`$; we shall verify that $`T_1\stackrel{~}{+}T_2`$ is narrow, using Theorem 4.7 above. Thus, let $`Y`$ be a separable subspace of $`C(K)`$. We shall first argue that there is a separable space $`Z_1`$ containing $`Y`$ such that $`T_1_{Z_1}`$ and $`T_2_{Z_1}`$ are strong Daugavet operators. Let $`A`$ be a countable dense subset of $`S(Y)`$. For every pair $`(x,y)`$ in $`A\times A`$ and every $`\epsilon =1/k`$ there is some $`z_1`$ (resp. $`z_2`$) according to the definition of the strong Daugavet property of $`T_1`$ (resp. $`T_2`$). The countable collection of these $`z`$’s and $`Y`$ span a closed separable subspace $`X_1`$. Repeating this procedure starting from $`X_1`$ yields some closed separable subspace $`X_2`$, etc. The closed linear span $`Z_1`$ of $`X_1,X_2,X_3,\mathrm{}`$ then has the desired property. Now by \[22, Lemma 2.4\] there is a separable space $`Z_2Z_1`$ isometric to some space $`C(M_2)`$ for a perfect compact metric space $`M_2`$. By the same token as above, we can extend $`Z_2`$ to a separable space $`Z_3`$ so that $`T_1`$ and $`T_2`$ are strong Daugavet operators on $`Z_3`$, and we can extend $`Z_3`$ to a separable space $`Z_4`$ isometric to some space $`C(M_4)`$ for a perfect compact metric space $`M_4`$, etc. Let $`Z`$ be the closed linear span of $`Z_1,Z_2,Z_3,\mathrm{}`$. Then $`T_1_Z`$ and $`T_2_Z`$ are strong Daugavet operators, and $`Z`$ is isometric to some space $`C(M)`$ for a perfect compact metric space $`M`$. By what we already know, $`(T_1\stackrel{~}{+}T_2)_Z`$ is a narrow operator on $`ZC(M)`$; recall that the classes of narrow and strong Daugavet operators coincide on $`C(M)`$. Finally, Theorem 4.7 implies that $`T_1\stackrel{~}{+}T_2`$ is narrow on $`C(K)`$, which proves the theorem. ∎ Next we introduce a topology related to an order ideal of operators. ###### Definition 4.9. Let $`𝒪𝒫(X)`$ be an order ideal of operators, closed under the operation $`\stackrel{~}{+}`$. Then the system of tubes $`U_{T,\epsilon }`$, $`T`$, $`\epsilon >0`$, defines a base of neighbourhoods of $`0`$ for some locally convex topology on $`X`$. We denote this topology by $`\sigma (X,)`$. If $`=(X)`$, the class of all finite-rank operators, then $`\sigma (X,)`$ coincides with the weak topology; if $`=𝒪𝒫(X)`$, then $`\sigma (X,)`$ coincides with the norm topology. For classes which are in between one gets topologies which are between the weak and the norm topology. If $`𝒩`$ is a collection of subsets in $`X`$ such that $`𝒩^{}`$ is closed under the operation $`\stackrel{~}{+}`$, then $`\sigma (X,𝒩^{})`$ is the strongest locally convex topology on $`X`$ continuous with respect to the norm, in which the zero vector belongs to the closure of every element of $`𝒩`$. ###### Definition 4.10. A locally convex topology $`\tau `$ on $`X`$ is said to be a Daugavet topology if for every two elements $`x,yS(X)`$, for every $`\epsilon >0`$ and every $`\tau `$-neighbourhood $`U`$ of $`y`$ there is an element $`zUS(X)`$ such that $`z+x>2\epsilon `$. Of course, $`\sigma (X,)`$ is a Daugavet topology if and only if every operator $`T`$ is a strong Daugavet operator. ###### Lemma 4.11. Let $`X`$ be a Banach space with the Daugavet property, $`T`$ a narrow operator, $`A=\{a_1,\mathrm{},a_n\}S(X)`$, $`\epsilon >0`$ and $`yS(X)`$. Then for every $`\sigma (X,cp(𝒩𝒜(X)))`$-neighbourhood $`W`$ of $`y`$ there is an element $`wWS(X)`$ such that $`T(wy)<\epsilon `$ and $`w+a>2\epsilon `$ for every $`aA`$. ###### Proof. We shall argue by induction on $`n`$. First of all consider $`n=1`$. Every $`\sigma (X,cp(𝒩𝒜(X)))`$-neighbourhood of $`y`$ can be represented as $`W=y+U_{R,\delta }`$, where $`Rcp(𝒩𝒜(X))`$. Since $`T_1=R\stackrel{~}{+}T`$ is a strong Daugavet operator by definition of the central part, there is an element $`wS(X)`$ such that $`w+a_1>2\epsilon `$ and $`T_1(wy)<\mathrm{min}(\delta ,\epsilon )`$. The last inequality means, in particular, that $`T(wy)<\epsilon `$ and $`wW`$. Now suppose our assertion is true for $`n`$, let us prove it for $`n+1`$. Let $`A=\{a_1,\mathrm{},a_n,a_{n+1}\}S(X)`$, and let us assume that an element $`w_1WS(X)`$ such that $`T(w_1y)<\epsilon /2`$ and $`w_1+a_k>2\epsilon `$, $`k=1,\mathrm{},n`$, has already been selected. Then there is a weak neighbourhood $`U`$ of $`w_1`$ such that the inequalities $`u+a_k>2\epsilon `$, $`k=1,\mathrm{},n`$, hold for every $`uU`$. The intersection $`UW`$ is a $`\sigma (X,cp(𝒩𝒜(X)))`$-neighbourhood of $`w_1`$, so according to our inductive assumption for $`n=1`$, there is an element $`wS(X)UW`$ such that $`w+a_{n+1}>2\epsilon `$ and $`T(ww_1)<\epsilon /2`$. This element $`w`$ satisfies all the requirements. ∎ Using an $`\epsilon `$-net of the unit ball of the finite-dimensional subspace $`Z`$ below one can easily deduce the following corollary. ###### Proposition 4.12. Let $`X`$ be a Banach space with the Daugavet property, $`T`$ be a narrow operator and $`ZX`$ be a finite-dimensional subspace. Then for every $`\epsilon >0`$, every $`yS(X)`$ and every $`\sigma (X,cp(𝒩𝒜(X)))`$-neighbourhood $`W`$ of $`y`$ there is an element $`wWS(X)`$ such that $`T(wy)<\epsilon `$ and $`z+w>(1\epsilon )(z+w)`$ for every $`zZ`$. ###### Theorem 4.13. Let $`X`$ be a Banach space with the Daugavet property and let $`T`$ be an operator on $`X`$ which does not fix a copy of $`\mathrm{}_1`$. Then $`Tcp(𝒩𝒜(X))`$, so in particular $`T`$ is a narrow operator. ###### Proof. Lemma 1(xii) of implies that every operator which does not fix a copy of $`\mathrm{}_1`$ can be factored through a space without $`\mathrm{}_1`$-subspaces. So every operator which does not fix a copy of $`\mathrm{}_1`$ can be majorized by an operator which maps into a space without $`\mathrm{}_1`$-subspaces. Since the class of narrow operators is an order ideal, it is enough to prove our theorem for $`T:XY`$, where $`Y`$ has no $`\mathrm{}_1`$-subspaces. Also, by Theorem 4.7 we may assume that $`X`$ and $`Y`$ are separable. Let us fix a narrow operator $`R`$, $`\epsilon >0`$ and $`x,yS(X)`$. Let us introduce a directed set $`(\mathrm{\Gamma },)`$ as follows: the elements of $`\mathrm{\Gamma }`$ are finite sequences in $`S(X)`$ of the form $`\gamma =(x_1,\mathrm{},x_n)`$, $`n`$, with $`x_1=x`$. The (strict) ordering is defined by $$(x_1,\mathrm{},x_n)(y_1,\mathrm{},y_m)n<m\&\{x_1,\mathrm{},x_n\}\{y_1,\mathrm{},y_{m1}\}$$ and of course $`\gamma _1\gamma _2`$ if $`\gamma _1\gamma _2`$ or $`\gamma _1=\gamma _2`$. Now define a bounded function $`F:\mathrm{\Gamma }Y\times \times `$ by $$F(\gamma )=(Tx_n,\alpha (\gamma ),R(yx_n)),$$ where $$\alpha (\gamma )=sup\{a>0:z+x_n>a(z+x_n)zlin\{x_1,x_2,\mathrm{},x_{n1}\}\}.$$ Due to Proposition 4.12, for every weak neighbourhood $`U`$ of $`y`$ in $`B(X)`$, every $`\epsilon >0`$ and every finite collection $`\{v_1,\mathrm{},v_n\}X`$ there is some $`v_{n+1}U`$ for which $`\alpha \left((v_1,\mathrm{},v_{n+1})\right)>1\epsilon `$ and $`R(yv_{n+1})<\epsilon `$. This means that $`(Ty,1,0)`$ is a weak limit point of the function $`F`$. So, by Theorem 4.3 there is a strictly $``$-increasing sequence $`(\gamma _j)=\left((x_1,\mathrm{},x_{n(j)})\right)`$ for which $`(Tx_{n(j)})`$ tends weakly to $`Ty`$, $`(R(yx_{n(j)}))`$ tends to 0 and $`(\alpha (\gamma _j))`$ tends to $`1`$. Passing to a subsequence we can select points $`x_{n(j)}`$ in such a way that the sequence $`\{x,x_{n(1)},x_{n(2)},\mathrm{}\}`$ is $`\epsilon `$-equivalent to the canonical basis of $`\mathrm{}_1`$. According to Mazur’s theorem, there is a sequence $`z_nconv\{x_{n(j)}\}_{j>n}`$ such that $`TyTz_n0`$. Evidently $`z_n+x>2\epsilon `$ and $`(R\stackrel{~}{+}T)(yz_n)0`$, which means that $`R\stackrel{~}{+}T𝒮𝒟(X)`$ and thus proves the theorem by Corollary 3.14(b). ∎ There are other applications of Theorem 4.3 which are not related to the Daugavet property. As an example let us prove the following theorem which was earlier established by E. Behrends under the more restrictive condition of separability of $`X^{}`$. ###### Theorem 4.14. Let $`X`$ be a Banach space without $`\mathrm{}_1`$-subspaces and $`A_nX`$ be bounded subsets with $`0\overline{\mathrm{conv}}A_n`$ for each $`n`$. Then there exists a sequence $`(a_n)`$ in $`X`$ with $`a_nA_n`$ for every $`n`$ such that $`0\overline{\mathrm{conv}}\{a_1,a_2,\mathrm{}\}`$. ###### Proof. In each $`A_n`$ there is a separable subset whose closed convex hull contains $`0`$. So, passing to the linear span of these separable subsets we may assume that $`X`$ is separable. Introduce a directed set $`(\mathrm{\Gamma },)`$ as follows: the elements of $`\mathrm{\Gamma }`$ are of the form $$\gamma =(n,m,\{a_k\}_{k=n}^m,\{\lambda _k\}_{k=n}^m),$$ where $`n,m`$, $`n<m`$, $`a_kA_k`$, $`\lambda _k>0`$, $`_{k=n}^m\lambda _k=1`$. Define $``$ as follows: let $`\gamma _1=(n_1,m_1,\{a_k\}_{k=n_1}^{m_1},\{\lambda _k\}_{k=n_1}^{m_1})`$, $`\gamma _2=(n_2,m_2,\{b_k\}_{k=n_2}^{m_2},\{\mu _k\}_{k=n_2}^{m_2})`$; then $`\gamma _1\gamma _2`$ if $`m_1<n_2`$. Define $`F:\mathrm{\Gamma }X`$ by the formula $`F(\gamma )=_{k=n}^m\lambda _ka_k`$. Now, $`0`$ is a weak limit point of $`F`$; see the proof of \[2, Th. 4.3\]. So, by Theorem 4.3 there is a sequence of elements $$\gamma _j=(n_j,m_j,\{a_k\}_{k=n_j}^{m_j},\{\lambda _k\}_{k=n_j}^{m_j})$$ such that $`n_1<m_1<n_2<m_2<n_3<\mathrm{}`$ and $`_{k=n_j}^{m_j}\lambda _ka_k`$ tends weakly to zero. To finish the proof one just needs to apply Mazur’s theorem. ∎ ## 5. Rich subspaces In a subspace $`Y`$ of $`L_1`$ is called rich if the quotient map $`q:L_1L_1/Y`$ is $`L_1`$-narrow, and likewise a subspace $`Y`$ of $`C(K)`$ is called rich in if the quotient map $`q:C(K)C(K)/Y`$ is $`C`$-narrow. We are now in a position to discuss rich subspaces in general. ###### Definition 5.1. Let $`X`$ be a Banach space with the Daugavet property. A subspace $`Y`$ is said to be almost rich if the quotient map $`q:XX/Y`$ is a strong Daugavet operator. A subspace $`Y`$ is said to be rich if the quotient map $`q:XX/Y`$ is a narrow operator. By Theorem 3.7 the new definition comprises the old one for subspaces of $`C(K)`$. The necessity to distinguish rich and almost rich subspaces will become apparent later when we show that the following theorem does not extend to almost rich subspaces; see Theorem 6.4. ###### Theorem 5.2. A rich subspace $`Y`$ of a Banach space $`X`$ with the Daugavet property has the Daugavet property itself. Moreover, $`(Y,X)`$ is a Daugavet pair. ###### Proof. Consider elements $`xS(X)`$, $`yS(Y)`$, a slice $`S=S(x^{},\epsilon )`$ and $`yS`$. According to our assumption the quotient map $`q:XX/Y`$ is a narrow operator. So there is an element $`uS`$ such that $`u+x>2\epsilon `$ and $`q(yu)=q(u)<\epsilon `$. The last condition means that the distance from $`u`$ to $`Y`$ is smaller than $`\epsilon `$, so there is an element $`vY`$ with $`vu<\epsilon `$. The norm of $`v`$ is close to $`1`$, viz. $`12\epsilon <v<1+\epsilon `$. Put $`w=v/v`$. For this $`w`$ we have $`wu<3\epsilon `$, so $`wS(x^{},4\epsilon )`$ and $`w+x>24\epsilon `$. ∎ This theorem leads to new hereditary properties for the Daugavet property. ###### Proposition 5.3. Suppose $`Y`$ is a subspace of a Banach space $`X`$ with the Daugavet property. * If the quotient space $`X/Y`$ has the Radon-Nikodým property, then $`Y`$ is rich. * If the quotient space $`X/Y`$ contains no copy of $`\mathrm{}_1`$, then $`Y`$ is rich in $`X`$. * If $`(X/Y)^{}`$ has the Radon-Nikodým property, then $`Y`$ is rich. In either case $`Y`$ has the Daugavet property itself. ###### Proof. (a) follows from Theorem 3.13, (b) from Theorem 4.13, and (c) follows from (b). ∎ That $`Y`$ has the Daugavet property under assumption (a) has been proved earlier in . ###### Remark 5.4. If the quotient map $`q:XX/Y`$ belongs to $`cp(𝒩𝒜(X))`$, then the restriction to $`Y`$ of every narrow operator on $`X`$ is a narrow operator itself. If $`Y`$ is a rich subspace of a space $`X`$ having the Daugavet property, then the restriction to $`Y`$ of every operator $`Tcp(𝒩𝒜(X))`$ is a narrow operator. ###### Definition 5.5. We say that a subspace $`Y`$ of a space $`X`$ with the Daugavet property is wealthy if $`Y`$ and every subspace of $`X`$ containing $`Y`$ have the Daugavet property. It is plain that if $`Y`$ is an (almost) rich subspace of a space $`X`$ with the Daugavet property, then every bigger subspace is (almost) rich, too. Thus, if $`Y`$ is rich, then it is wealthy. We now investigate the converse implication. ###### Lemma 5.6. The following conditions for a subspace $`Y`$ of a Banach space $`X`$ with the Daugavet property are equivalent: * $`Y`$ is wealthy. * Every finite-codimensional subspace of $`Y`$ is wealthy. * For every pair $`x,yS(X)`$, the linear span of $`Y`$, $`x`$ and $`y`$ has the Daugavet property. * For every $`x,yS(X)`$, for every $`\epsilon >0`$ and for every slice $`S`$ of $`S(X)`$ which contains $`y`$ there is an element $`vlin(\{x,y\}Y)S`$ such that $`x+v>2\epsilon `$. ###### Proof. Due to Proposition 5.3 every finite-codimensional subspace of a space with the Daugavet property has the Daugavet property itself (see also \[13, Th. 2.14\]); this is the reason for the equivalence of (i) and (ii). The implication (i) $``$ (iii) follows immediately from the definition of a wealthy subspace; (iii) $``$ (iv) and (iv) $``$ (i) are consequences of Lemma 1.1. ∎ Let us say that a pair of elements $`x,yS(X)`$ is $`\epsilon `$-fine if there is a slice $`S`$ of $`S(X)`$ which contains $`y`$ and the diameter of $`Slin\{x,y\}`$ is less than $`\epsilon `$. ###### Lemma 5.7. Let $`Y`$ be a wealthy subspace of a Banach space $`X`$ with the Daugavet property and let a pair $`x,yS(X)`$ be $`\epsilon `$-fine. Then $`Y`$ intersects $`D(x,y,2\epsilon )`$. ###### Proof. First of all let us fix a slice $`S=S(x^{},\epsilon _1)`$ from the definition of an $`\epsilon `$-fine pair and fix a $`\delta >0`$ such that the set $`W=\{wlin\{x,y\}:w<1+\delta `$, $`x^{}(w)>1\epsilon _1\}`$ still has diameter less than $`\epsilon `$. Now let us find a finite-codimensional subspace $`EY`$ such that 1. $`x^{}=0`$ on $`E`$, 2. if $`eE`$ and $`wlin\{x,y\}`$, then $`w<(1+\delta )e+w`$; the last condition can be satisfied by a variant of the Mazur argument leading to the basic sequence selection principle; see \[10, Lemma 6.3.1\]. According to our assumptions $`lin(\{x,y\}E)`$ has the Daugavet property. So there is an element $`vlin(\{x,y\}E)S`$ such that $`x+v>2\epsilon `$. Let us represent $`v`$ in the form $`v=e+w`$, where $`eE`$, $`wlin\{x,y\}`$. By choice of $`E`$ this means that $`w<1+\delta `$ and $`x^{}(w)=x^{}(v)>1\epsilon _1`$. Thus, $`wW`$ and $`yw<\epsilon `$. Finally we have that the element $`e`$ belongs to $`ED(x,y,2\epsilon )`$, which concludes the proof. ∎ Let us recall the following result , which can also be deduced from our Proposition 4.12: ###### Lemma 5.8. Let $`X`$ be a Banach space with the Daugavet property and $`ZX`$ be a finite-dimensional subspace. Then, for every $`\epsilon >0`$ in every slice of the unit sphere of $`X`$ there is an element $`x`$ such that $$z+x>(1\epsilon )(z+x)zZ.$$ We now present two easy lemmas. ###### Lemma 5.9. A subspace $`Y`$ of a Banach space with the Daugavet property which is almost rich together with all of its $`1`$-codimensional subspaces is rich. ###### Proof. Let $`q:XX/Y`$ be the quotient map and let $`x^{}S(X^{})`$; further let $`Y_1=Y\mathrm{ker}x^{}`$ and let $`q_1:XX/Y_1`$ be the corresponding quotient map. Then $`Y_1=Y`$ or $`Y_1`$ is $`1`$-codimensional in $`Y`$. Now, in either case we have $`q(x)+\|x^{}(x)\|2q_1(x)`$ for all $`xX`$. Since $`q_1`$ is a strong Daugavet operator by assumption, so is $`q\stackrel{~}{+}x^{}`$, and $`q`$ is narrow. ∎ ###### Lemma 5.10. A subspace $`Y`$ of a Banach space $`X`$ with the Daugavet property is almost rich if and only if $`Y`$ intersects all the elements of $`𝒟(X)`$. ###### Proof. If $`Y`$ intersects all the elements of $`𝒟(X)`$, then the quotient map $`q:XX/Y`$ is unbounded from below on every element of $`𝒟(X)`$. So the quotient map belongs to $`𝒟(X)^{}`$ which coincides with the class of strong Daugavet operators by Proposition 3.4. Now consider the converse statement. If $`Y`$ is almost rich, then for every $`\epsilon >0`$ the map $`q`$ is unbounded from below on every set of the form $`D(x,y,\epsilon /2)`$. This means that there is an element $`zY`$ for which $`dist(z,D(x,y,\epsilon /2))<\epsilon /2`$. In this case $`z`$ belongs to $`D(x,y,\epsilon )`$, so the intersection of this set with $`Y`$ is non-empty. ∎ The following is the key result for establishing that wealthy subspaces are rich. ###### Lemma 5.11. Every wealthy subspace $`Y`$ of a Banach space $`X`$ having the Daugavet property is almost rich. ###### Proof. According to Lemma 5.10 we need to prove that for every positive $`\epsilon <1/10`$ and every pair $`x,yS(X)`$ the subspace $`Y`$ intersects $`D(x,y,\epsilon )`$. To do this, according to Lemma 5.7, it is enough to show that for every $`\epsilon >0`$ and every pair $`x,yS(X)`$ there is an $`\epsilon `$-fine pair $`x_1,y_1S(X)`$ which approximates $`(x,y)`$ well; i.e., $`xx_1+yy_1<\epsilon `$. Let us fix a positive $`\delta <\epsilon ^2/8`$ and select an element $`zS(X)`$ in such a way that for every $`wlin\{x,y\}`$ and for every $`t>0`$ $$w+tz(1\delta )(w+\|t\|)$$ (we use Lemma 5.8). Put $`x_1=x+\epsilon z`$, $`y_1=y`$. To show that $`(x_1,y)`$ is an $`\epsilon `$-fine pair it is sufficient to demonstrate that, for every $`vlin\{x_1,y\}`$ with $`v\epsilon `$, $`\mathrm{max}\{y+v,yv\}>1`$. To do this let us argue ad absurdum. Take some $`v=ay+b(x+\epsilon z)`$ with $`v\epsilon `$ and assume that $`\mathrm{max}\{y+v,yv\}=1`$. Then $`1`$ $`=\mathrm{max}\{y+ay+b(x+\epsilon z),yayb(x+\epsilon z)\}`$ $`(1\delta )(\mathrm{max}\{y+ay+bx+,yaybx\}+\|b\|\epsilon )`$ $`(1\delta )(1+\|b\|\epsilon ).`$ So $`\|b\|\epsilon /4`$. But in this case $`\|a\|>\epsilon /2`$ and $$\mathrm{max}\{y+v,yv\}>\mathrm{max}\{y+ay,yay\}\epsilon /3>1+\epsilon /6,$$ which provides a contradiction. ∎ ###### Theorem 5.12. The following properties of a subspace $`Y`$ of a Daugavet space $`X`$ are equivalent: * $`Y`$ is wealthy. * $`Y`$ is rich. * Every finite-codimensional subspace of $`Y`$ is rich. ###### Proof. It is clear that (iii) $``$ (ii) $``$ (i), see the remark following Definition 5.5. Now suppose (i). Every $`1`$-codimensional subspace of $`Y`$ is wealthy by Lemma 5.6 and is hence almost rich by Lemma 5.11. An appeal to Lemma 5.9 completes the proof. ∎ ## 6. Operators on $`L_1`$ In this section we shall study strong Daugavet and narrow operators on $`L_1`$. We first introduce a technical definition. Let $`(\mathrm{\Omega },\mathrm{\Sigma },\mu )`$ be an atomless probability space. A function $`fL_1=L_1(\mu )`$ is said to be a balanced $`\epsilon `$-peak on $`A\mathrm{\Sigma }`$ if $`f1`$, $`suppfA`$, $`_\mathrm{\Omega }f𝑑\mu =0`$ and $`\mu \{t:f(t)=1\}>\mu (A)\epsilon `$. The collection of all balanced $`\epsilon `$-peaks on $`A`$ will be denoted by $`P(A,\epsilon )`$. ###### Theorem 6.1. $`𝒩𝒜(L_1)=\{P(A,\epsilon ):A\mathrm{\Sigma },\epsilon >0\}^{}`$. ###### Proof. Let $`T𝒩𝒜(L_1)`$, $`\delta ,\epsilon >0`$, and $`A\mathrm{\Sigma }`$. Consider a slice in $`L_1`$ of the form $$S=\{fB(L_1):_Af𝑑\mu >1\delta \}.$$ Applying Lemma 3.8 to this slice, the elements $`x=\chi _A/\mu (A)`$, $`y=\chi _A/\mu (A)`$ and $`\delta `$ we get a function $`vS`$ such that (6.1) $$v\chi _A/\mu (A)>2\delta ,T(v\chi _A/\mu (A))<\delta .$$ Denote by $`B`$ the set $`\{tA:v(t)>0\}`$. The condition $`vS`$ implies that $`v\chi _Bv<\delta `$, so $$v\chi _B\chi _A/\mu (A)>22\delta .$$ Next, introduce $`C=\{tA:v(t)>1/\mu (A)\}`$. By the last inequality $$v\chi _C\chi _A/\mu (A)>22\delta ,v\chi _Cv<3\delta $$ and (6.2) $$\mu (C)<\delta \mu (A);$$ to see this observe that $`22\delta `$ $`<\chi _Bv{\displaystyle \frac{\chi _A}{\mu (A)}}{\displaystyle _C}\left(\chi _Bv{\displaystyle \frac{1}{\mu (A)}}\right)𝑑\mu +{\displaystyle \frac{1}{\mu (A)}}(\mu (A)\mu (C))`$ $`22{\displaystyle \frac{\mu (C)}{\mu (A)}}.`$ Put $`f=(\mu (A)/\beta )\chi _Cv\chi _A`$ with $`\beta =_Cv𝑑\mu `$ so that $`_\mathrm{\Omega }f𝑑\mu =0`$. Since $`_Av𝑑\mu >1\delta `$ we have from $`v\chi _Cv<3\delta `$ that $`\beta 14\delta `$. By (6.1) we conclude that $$Tf=\mu (A)T\left(\frac{\chi _Cv}{\beta }\frac{\chi _A}{\mu (A)}\right)\mu (A)\left(T\frac{\chi _Cv}{\beta }v+\delta \right)$$ and $$\frac{\chi _Cv}{\beta }v\frac{\chi _Cvv}{\beta }+\frac{v}{\beta }v\frac{3\delta }{\beta }+\left(\frac{1}{\beta }1\right)\frac{7\delta }{14\delta },$$ and if $`\delta `$ is small enough, by (6.2) $`fP(A,\epsilon )`$. This proves the inclusion $`𝒩𝒜(L_1)\{P(A,\epsilon ):A\mathrm{\Sigma },\epsilon >0\}^{}`$. To prove the opposite inclusion we use Proposition 3.11. Let us fix $`T\{P(A,\epsilon ):A\mathrm{\Sigma },\epsilon >0\}^{}`$. Let $`x,yS(L_1)`$, $`y^{}S(L_{\mathrm{}})`$ and $`\epsilon >0`$ be such that $`y^{},y>1\epsilon `$. Without loss of generality we may assume that there is a partition $`A_1,\mathrm{},A_n`$ of $`\mathrm{\Omega }`$ such that the restrictions of $`x`$, $`y`$ and $`y^{}`$ on $`A_k`$ are constants, say $`a_k`$, $`b_k`$ and $`c_k`$ respectively. By our assumption $`T`$ is unbounded from below on each of the $`P(A_k,\delta )`$ for every $`\delta >0`$, $`k=1,\mathrm{},n`$. Let us fix functions $`f_kP(A_k,\delta )`$ such that $`Tf_k<\delta `$, $`k=1,\mathrm{},n`$, and put $$v=\underset{k=1}{\overset{n}{}}b_k(\chi _{A_k}+f_k).$$ By definition of balanced $`\delta `$-peaks $`y^{},v>1\epsilon `$, $`v=1`$, and $`T(yv)`$ and $`\mu (suppv)`$ become arbitrarily small when $`\delta `$ is small enough. Thus $`\delta `$ can be chosen so that $`v`$ fulfills the conditions $`T(yv)<\epsilon `$ and $`x+v>2\epsilon `$. ∎ The characterisation of narrow operators on $`L_1`$ proved above looks similar to the definition of $`L_1`$-narrow operators. It is easy to prove that every $`L_1`$-narrow operator is narrow. We don’t know whether the classes of narrow operators and $`L_1`$-narrow operators on $`L_1`$ coincide. The aim of the remainder of this section is to construct an example of a strong Daugavet operator on $`L_1`$ which is not narrow. In fact, we shall define a subspace $`YL_1[0,1]`$ so that the quotient map $`q:L_1L_1/Y`$ is a strong Daugavet operator, but $`Y`$ fails the Daugavet property. By Theorem 5.2, $`q`$ cannot be narrow. Likewise, $`Y`$ is almost rich, but not rich. Let $`I_{n,k}=[\frac{k1}{2^n},\frac{k}{2^n})`$ for $`n_0`$ and $`k=1,2,\mathrm{},2^n`$. Fix $`N`$. We define $`g_{0,1}`$ $`=`$ $`(2^N1)\chi _{I_{N,1}}\chi _{I_{0,1}I_{N,1}}`$ $`g_{1,1}`$ $`=`$ $`(2^{N^2N}1)\chi _{I_{N^2,1}}\chi _{I_{N,1}I_{N^2,1}}`$ $`g_{1,k}`$ $`=`$ $`g_{1,1}(t\frac{k1}{2^n}),k=2,\mathrm{},2^N,`$ $`\mathrm{}`$ $`g_{n,1}`$ $`=`$ $`(2^{N^{n+1}N^n}1)\chi _{I_{N^{n+1},1}}\chi _{I_{N^n,1}I_{N^{n+1},1}}`$ $`g_{n,k}`$ $`=`$ $`g_{n,1}(t\frac{k1}{2^{N^n}}),k=2,\mathrm{},2^{N^n}.`$ Denote by $`P_n`$ the “peak set” of the $`n`$’th generation, i.e., $$P_n=\{t[0,1]:\underset{k=1}{\overset{2^{N^n}}{}}g_{n,k}(t)>0\},$$ and $`P=_nP_n`$. Clearly $`\|P_n\|=2^{N^n}/2^{N^{n+1}}=\left(1/2^{N1}\right)^{N^n}`$ and $`\|P\|1/(2^N1)`$. Notice also that $`_0^1g_{n,k}(t)𝑑t=0`$ for all $`n`$ and $`k`$. First we formulate a lemma. All the norms appearing below are $`L_1`$-norms. ###### Lemma 6.2. Let $$g=\underset{n=0}{\overset{M}{}}\underset{k=1}{\overset{2^{N^n}}{}}a_{n,k}g_{n,k}.$$ Then $$g\chi _{[0,1]P}3g\chi _P.$$ ###### Proof. Denote $$g^{\prime \prime }=\underset{suppg_{n,k}P}{}a_{n,k}g_{n,k},g^{}=gg^{\prime \prime }.$$ Since $`g^{}`$ and $`g`$ coincide off $`P`$, we clearly have (6.3) $$g^{}\chi _{[0,1]P}=g\chi _{[0,1]P}.$$ We also have that (6.4) $$g^{}\chi _Pg\chi _P.$$ Indeed, we can write $`P`$ as a countable union of disjoint (half-open) intervals; denote by $`I`$ any one of these. Then $`g^{}`$ is constant on $`I`$, and $`_0^1g^{\prime \prime }(t)𝑑t=0`$. Hence $$g^{}\chi _I=\left\|_0^1g^{}(t)\chi _I(t)𝑑t\right\|=\left\|_0^1(g^{}(t)\chi _I(t)+g^{\prime \prime }(t)\chi _I(t))𝑑t\right\|g\chi _I.$$ Summing up over all $`I`$ gives the result. Next, we claim that (6.5) $$g^{}\chi _{[0,1]P}3g^{}\chi _P.$$ To see this, we label the intervals $`I`$ from the previous paragraph as follows. For every $`l`$ write $`B_0=P_0`$ and $`B_l=P_l_{i=1}^{l1}P_i`$. Each $`B_l`$ can be written as $`_{dD_l}I_{N^{l+1},d}`$ where $`D_l`$ is some subset of $`\{1,\mathrm{},2^{N^{l+1}}\}`$ with cardinality $`<2^{N^l}`$. Let us write $`g^{}=_{n=0}^M_{k=1}^{2^{N^n}}b_{n,k}g_{n,k}`$. We then have the estimates $$_0^1\|g^{}(t)\chi _{B_0}(t)\|𝑑t=\|b_{0,1}\|\frac{2^N1}{2^N}$$ and $`{\displaystyle _0^1}\|g^{}(t)\chi _{B_l}(t)\|𝑑t`$ $`=`$ $`{\displaystyle \underset{dD_l}{}}{\displaystyle _{I_{N^{l+1},d}}}\|b_{0,1}{\displaystyle \underset{n=1}{\overset{l1}{}}}{\displaystyle \underset{k=1}{\overset{2^{N^n}}{}}}b_{n,k}\chi _{suppg_{n,k}}`$ $`+b_{l,(d1)/(2^{N1})^{N^l}+1}\left(2^{N^{l+1}N^l}1\right)\|dt`$ $``$ $`{\displaystyle \underset{k=1}{\overset{2^{N^l}}{}}}\left({\displaystyle \frac{1}{2^{N^l}}}{\displaystyle \frac{1}{2^{N^{l+1}}}}\right)\|b_{l,k}\|`$ $`{\displaystyle \frac{1}{(2^{N1})^{N^l}}}\|b_{0,1}\|{\displaystyle \frac{1}{(2^{N1})^{N^l}}}{\displaystyle \underset{n=1}{\overset{l1}{}}}{\displaystyle \underset{k=1}{\overset{2^{N^l}}{}}}\|b_{n,k}\|.`$ Summing up over all $`l`$ gives us $`{\displaystyle _P}\|g^{}(t)\|𝑑t`$ $``$ $`\|b_{0,1}\|\left({\displaystyle \frac{2^N1}{2^N}}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2^{N1})^{N^m}}}\right)`$ $`+{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{1}{2^{N^l}}}{\displaystyle \frac{1}{2^{N^{l+1}}}}{\displaystyle \underset{m=l+1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2^{N1})^{N^m}}}\right){\displaystyle \underset{k=1}{\overset{2^{N^l}}{}}}\|b_{l,k}\|`$ $``$ $`{\displaystyle \frac{1}{2}}\|b_{0,1}\|+{\displaystyle \frac{1}{2}}{\displaystyle \underset{l=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{2^{N^l}}}{\displaystyle \underset{k=1}{\overset{2^{N^l}}{}}}\|b_{l,k}\|.`$ On the other hand, by the triangle inequality $$_0^1\|g^{}(t)\|𝑑t2\left(\|b_{0,1}\|+\underset{l=1}{\overset{\mathrm{}}{}}\frac{1}{2^{N^l}}\underset{k=1}{\overset{2^{N^l}}{}}\|b_{l,k}\|\right),$$ hence the claim follows. The lemma now results from (6.3)–(6.5) ∎ ###### Theorem 6.3. Let $`Y_NL_1[0,1]`$ be the closed subspace generated by the system $`\{g_{n,k}\}`$ and the constants. Then the quotient map $`q_N:L_1L_1/Y_N`$ is a strong Daugavet operator for all $`N`$, but $`Y_N`$ fails the Daugavet property if $`N4`$. ###### Proof. Let us fix $`x,yS(L_1)`$ and $`\epsilon >0`$. Without loss of generality we may assume that $`x=_{k=1}^{2^{N^n}}a_{n,k}\chi _{I_{n,k}}`$ for a big enough $`n`$ to be chosen later. Put $`h=_{k=1}^{2^{N^n}}a_{n,k}g_{n,k}`$. Then $$x+h=\underset{k=1}{\overset{2^{N^n}}{}}2^{N^{n+1}N^n}\chi _{N^{n+1},d_{n,k}}a_{n,k}$$ with $`d_{n,k}=1+(k1)(2^{N1})^{N^n}`$. So $$x+h=\underset{k=1}{\overset{2^{N^n}}{}}\frac{\|a_{n,k}\|}{2^{N^n}}=x=1,$$ and $`supp(x+h)P_n`$. Since $`\|P_n\|0`$ we can pick $`n`$ big enough to satisfy $`x+h+y>2\epsilon `$. This shows that $`q_N`$ is a strong Daugavet operator. To show that $`Y_N`$ fails the Daugavet property if $`N4`$, take $`g^{}=\chi _{[0,1]P}Y_N^{}`$ and $`\epsilon =2\|P\|`$. Since $`\mathrm{𝟏}S(Y_N)`$, we get $$g^{}g^{}(\mathrm{𝟏})=1\epsilon /2>1\epsilon .$$ Thus, $`S(g^{},\epsilon )B(Y_N)\mathrm{}`$. We show that there is no $`f`$ in this slice such that $`f\mathrm{𝟏}>2\epsilon `$. Suppose, on the contrary, that there is such an $`f`$. Without loss of generality we can assume that $$f=a_0\mathrm{𝟏}+g$$ where $`g`$ is as in Lemma 6.2. It follows from our conditions that (6.6) $$f\chi _P=_P\|f(t)\|𝑑t=fg^{}(\|f\|)1g^{}(f)<\epsilon .$$ Hence, $$1_0^1f(t)𝑑t=_Pf(t)𝑑t+g^{}(f)>12\epsilon ,$$ and since $`_0^1f(t)𝑑t=a_0`$, we get (6.7) $$12\epsilon <a_01.$$ By (6.6) and (6.7), (6.8) $$g\chi _P\epsilon +\|P\|<2\epsilon ,$$ thus (6.7) and (6.8) yield $$g\chi _{[0,1]P}g2\epsilon =fa_0\mathrm{𝟏}2\epsilon f\mathrm{𝟏}4\epsilon >25\epsilon .$$ But now Lemma 6.2 and (6.8) imply $$25\epsilon <g\chi _{[0,1]P}3g\chi _P<6\epsilon ,$$ which yields $`\epsilon >2/11`$, i.e., $`\|P\|>1/11`$, which is false for $`N4`$. ∎ Theorems 6.3 and 5.2 immediately yield the following result. ###### Theorem 6.4. There is an almost rich subspace of $`L_1[0,1]`$ which fails the Daugavet property and hence fails to be rich. Thus, on $`L_1[0,1]`$ the class of strong Daugavet operators does not coincide with the class of narrow operators. ## 7. Questions We finish this paper with some questions which have remained open. We intend to deal with these problems in a future publication. 1. Does the class of narrow operators on a Banach space $`X`$ form a subsemigroup of $`𝒪𝒫(X)`$? \[Added in the final version: We have recently constructed a counterexample on the space $`X=C([0,1],L_1[0,1])`$.\] 2. Is every narrow operator on $`L_1`$ also $`L_1`$-narrow? 3. Is the sum of two $`L_1`$-narrow operators from $`L_1`$ to $`L_1`$ again $`L_1`$-narrow? This question is clearly related to the previous ones; we remark that the proof in \[17, p. 69\] which purportedly shows this to be true appears to have a gap. 4. If $`X`$ has the Daugavet property, does $`X`$ have a subspace isomorphic to $`\mathrm{}_2`$? 5. If $`T`$ is an operator on a space $`X`$ with the Daugavet property which does not fix a copy of $`\mathrm{}_2`$, is $`T`$ then narrow? We remark that the answer is affirmative in the case $`X=C[0,1]`$ by our Theorem 4.13 and a result due to Bourgain .
warning/0005/cond-mat0005519.html	ar5iv	text	# Absence of lattice strain anomalies at the electronic topological transition in zinc at high pressure ## I Introduction Zinc and cadmium are unique among the hexagonal close-packed (hcp) transition metals in that the axial ratio ($`c/a=1.856`$ for zinc and 1.886 for cadmium) is far from the ideal value defined by hard sphere packing ($`c/a=\sqrt{8/3}=1.633`$). Upon compression, the axial ratio decreases towards the ideal value. Lynch and Drickamer first observed that the decrease in $`c/a`$ with increasing pressure was not smooth; subsequent experiments yielded inconsistent results on the nature of this anomaly . Takemura confirmed the anomaly using a methanol-ethanol-water mixture (MEWM) as a pressure medium in diamond anvil cell experiments: he observed the $`a`$-axis expanding over a small range of compression, yielding a rapid decrease of the axial ratio $`c/a`$. Ab initio computations found similar behavior and provided an explanation for the anomaly by means of changes in the Fermi surface topology under compression. Takemura recently repeated his experiments but using helium as pressure medium which is more nearly hydrostatic than MEWM, but found that both axes compressed monotonically with no anomaly in $`c/a`$, contrary to his earlier experiments and theory. The most recent experimental results call previous theoretical studies into question. All previous theoretical studies show the anomaly in the axial ratio whether the local density approximation (LDA) or the generalized gradient approximation (GGA) to the exchange correlation potential is used. The anomaly has been connected to changes in the electronic structure : Fast et al. observe one electronic topological transition (ETT) at the high symmetry point K on the Brillouin zone boundary forming a ellipsoidal piece (needle) in the Fermi surface. Novikov et al. see at least one additional ETT at approximately the same compression, also at K, where disconnected pieces form a three-leg structure of the Fermi surface along the K-M directions upon compression. Depending on $`c/a`$ with compression Novikov et al. propose one more ETT at L (butterfly) reconciling previous contradictory results from first principles calculations. All previous computations were performed with typical Brillouin zone sampling (no more than 1000 k-points in the irreducible wedge of the Brillouin zone). In an attempt to understand the discrepancy of the hydrostatic experiments and previous computational results we calculate the equation of state, lattice constants, and electronic structure of zinc over a wide compression range from first principles paying particular attention to the convergence of the calculations with respect to reciprocal space integration. In section II we introduce the method used and elaborate the computational details of our first principles calculations. Section III focuses on our results for the equation of state, lattice constants, and electronic structure. We compare our results to experiments at ambient conditions and high pressure and to previous theoretical work. Discussion and conclusions follow. ## II Method We investigate the energetics of hcp zinc using the full-potential linearized-augmented plane-wave method (LAPW) with GGA. Core states are treated self-consistently using the full Dirac equation for the spherical part of the potential, while valence states are treated in a semirelativistic approximation neglecting spin-orbit coupling. $`3s`$, $`3p`$, $`3d`$, $`4s`$, and $`4p`$ states are treated as valence electrons. The muffin-tin radius $`R_{MT}`$ is 2.0 Bohr over the whole compression range considered. We perform calculations at three sets of Brillouin zone sampling, 24x24x24, 32x32x32, and 48x48x48 special k-points , yielding 732, 1632, and 5208 k-points in the irreducible wedge of the Brillouin zone for the hcp lattice, respectively. The lowest k-point sampling is comparable to the previous GGA study while the latter two are much denser than any previously published results. The size of the basis is set by R<sub>MT</sub>K$`{}_{max}{}^{}=9.0`$, where K<sub>max</sub> is the largest reciprocal space vector. We use Fermi broadening with an electronic temperature of 5 mRy. For the densest k-point mesh we also perform calculations without electronic broadening for a selected subset of volumes and do not see any significant change in our results: equilibrium axial ratios remain within $`\pm `$ 0.005, the uncertainty of our results. We carry out total energy calculations over a wide range of unit cell volumes. At each volume we perform calculations for several different values of the axial ratio and find the equilibrium structure by fitting the results to a quadratic. The equation of state is obtained by describing the energy-volume curve with a third order expansion in Eulerian finite strain. We consider unit cell volumes ranging from 60-110 Bohr<sup>3</sup> for 24x24x24 k-points and focus on the range in which the anomaly in $`c/a`$ occurs (90-102.5 Bohr<sup>3</sup>) for 32x32x32 and 48x48x48 k-point meshes. ## III results ### Equation of State A comparison of the pressure-volume relation between our results and static experiments (Fig. 1 and Table I) shows good agreement at low pressure. At higher pressure theory differs significantly from the results of the MEWM diamond anvil cell experiments ; this is consistent with previous theoretical results . To investigate whether non-hydrostaticity may be responsible for the discrepancy we also compare to the results of shock wave experiments where hydrostaticity is readily achieved (Fig. 1 and Table I). The Hugoniot is reduced to a 0 K isotherm by solving the Rankin-Hugoniot equation . We estimate the thermal pressure ($`P_{th}=\gamma E_{th}/V`$) along the Hugoniot, with $`\gamma `$ the Grüneisen parameter and $`E_{th}`$ the thermal energy . We approximate the thermal energy by the Dulong Petit law ($`C_V^{lat}=3R`$); the electronic contribution to the thermal pressure is negligible (the temperature along the Hugoniot is less than 2000 K). We assume the Grüneisen parameter is proportional to compression ($`\gamma =\gamma _0V/V_0`$) with $`\gamma _0`$ its zero pressure value evaluated from the thermodynamic definition ($`\gamma =\alpha K_T/C_V\rho `$), where the thermal expansivity $`\alpha `$, the isothermal bulk modulus $`K_T`$, specific heat $`C_V`$ and density $`\rho `$ at zero pressure are taken from the literature . The reduced Hugoniot agrees with our GGA results much better than the static experiments; differences in volume are less than 1.5%. The large discrepancy between the static and shock wave experiments indicates that the MEWM experiments may be biased by non-hydrostaticity. ### Lattice Constants Total energy as a function of axial ratio for the 24x24x24 k-point mesh shows an unusually large scatter about the quadratic fit in $`c/a`$ (Fig. 2). With increasing number of k-points the scatter decreases and the minimum becomes better defined. In contrast to the previous GGA results we do not see multiple minima in $`c/a`$ for any volume and find the axial ratio reliably resolved to within $`\pm 0.005`$, within the symbol size in Fig. 1. The curvature of energy as a function of $`c/a`$ varies considerably for the different k-point meshes, showing that elastic constants will also be strongly dependent on k-point sampling, as the shear elastic constant ($`C_S`$) is related to this strain . The development of the axial ratio $`c/a`$ with compression differs for the three sets of computations considerably (Fig. 1). For 24x24x24 k-points we see an anomaly similar to that in the MEWM experiments: after an initial linear decrease in the axial ratio (102.5-95 Bohr<sup>3</sup>) the slope in $`c/a`$ steepens (95-90 Bohr<sup>3</sup>) before decreasing again at higher pressures. The dependence of $`c/a`$ on compression for k-point meshes of 32x32x32 and 48x48x48 is much smoother; the anomaly in $`c/a`$ has disappeared. The difference between experiment and theory is less than 4% in $`c/a`$ which is typical of all electron calculations . At higher compression (V$`<`$ 70 Bohr<sup>3</sup>) the theoretical value smoothly approaches 1.61; the MEWM experiments converge to 1.59. The nature of the anomaly is revealed by considering the lattice constants separately (Fig. 3). The $`c`$-axis compresses monotonically with decreasing volume in all computations and experiments considered. Theory overpredicts $`c`$ by less than 2%, and there is little difference in $`c`$ for the two denser k-point meshes. An expansion of the $`a`$-axis for the 24x24x24 k-point calculations and the MEWM experiments cause the anomaly in $`c/a`$ (Fig. 3). For the two denser k-point meshes $`a`$ compresses monotonically; for volumes smaller than $`V=95`$ Bohr<sup>3</sup> $`a`$ follows a linear trend with the same slope as the helium experiments. For volumes greater than V=95 Bohr<sup>3</sup>, $`a`$ is less compressible than it is at higher pressure. For the the larger k-point samplings the calculations underestimate $`a`$ by less than 1%, while with 24x24x24 k-point the maximum difference is approximately 1.5%. To illustrate this point further we evaluate the linear compressibility for the two axes $`k_x=(1/x)(x/P)`$ (with $`x=a,c`$) for our results and the static experiments using central differences (Fig. 4). Our results for 24x24x24 k-points show an anomaly in $`k_a`$ similar in character and magnitude to that found in the MEWM experiments. For denser k-point sampling the anomaly in $`k_a`$ is shifted towards lower pressure and its magnitude decreases with increasing number of k-points. For $`k_c`$ an anomaly exists as well for both the MEWM experiments and the calculations with the smallest k-point mesh, it is however less pronounced than for $`k_a`$ and is absent from the results for the two denser k-point meshes. ### Electronic Structure The band structure of zinc under compression changes considerably from $`V=102.5`$ Bohr<sup>3</sup> to $`V=95`$ Bohr<sup>3</sup> (Fig. 5). The electronic structure is in excellent agreement with the previous GGA results . The major change in band structure occurs at the high symmetry point K on the Brillouin zone boundary where three bands (K<sub>7</sub>,K<sub>8</sub>,K<sub>9</sub>) cross the Fermi energy under compression, changing the topology of the Fermi surface. From the band structure we see the needle around K and also the connection of the three-leg piece at K. Focusing on the development of the band structure at K we consider the eigenvalues of the K<sub>7</sub>, K<sub>8</sub>, and K<sub>9</sub> states (Fig. 6). For the 24x24x24 k-point calculations these bands show a quadratic volume dependence and cross the Fermi energy at V=97.5 Bohr<sup>3</sup> (K<sub>7</sub> and K<sub>8</sub>) and V=97 Bohr<sup>3</sup> (K<sub>9</sub>). For the two denser k-point meshes the eigenvalues depend linearly on volume and the crossing points are indistinguishable for 32x32x32 and 48x48x48 k-points. The crossings occur at slightly higher volume than for the 24x24x24 k-point mesh (V=98 and 97.5 Bohr<sup>3</sup> for K<sub>7</sub>, K<sub>8</sub>, and K<sub>9</sub>, respectively), but the difference is small compared with the effect of k-point sampling on the lattice parameters. ## IV Discussion The ETT discussed in the last section has previously been invoked as an explanation for the anomaly in $`a`$-axis compressibilities . In contrast to these studies we find that the occurrence of the ETT is independent of the calculated anomaly in $`a`$-axis compressibility, as the location of abnormal compression of $`a`$ shifts with increasing k-point sampling towards higher volumes (Figs. 3 and 4) while the ETT always occurs at approximately the same volume (Fig. 6). The anomaly in $`a`$-axis compression seen in previous calculations appears to be a consequence of insufficient k-point sampling. The results presented here for $`a`$ (Fig. 3) and linear compressibility $`k_a`$ (Fig. 4) suggest that even for the densest k-point mesh we use (48x48x48) the lattice parameters are not converged. The discrepancy between the MEWM and helium experiments can be attributed to freezing of the MEWM pressure medium which is known to occur at about 10 GPa . Freezing substantially increases the non-hydrostatic component of stress as recognized previously in high pressure experiments on forsterite (Mg<sub>2</sub>SiO<sub>4</sub>) . At room temperature helium also freezes within the pressure range of the experiment (11.5 GPa) but remains soft enough to maintain hydrostaticity . Recent neutron inelastic-scattering experiments under compression show no softening or anomaly in the phonon frequency, supporting the monotonic compression of both axes as seen in our dense k-point calculations and the helium experiments. In retrospect it occurs as a fortuitous (or unfortunate) coincidence that for typical computational parameters comparable behavior in linear compressibilities is found in first principles electronic structure calculations and for experiments with non-hydrostatic conditions, despite the fundamentally different underlying physical problem. Following the same notion the observation of anomalies in the axial ratio under compression for cadmium in experiments and theory using both LDA and GGA may also be an artifact due to non-hydrostatic conditions in the experiments and insufficient convergence with respect to computational parameters as well. The ETT, however, might have important effects on higher order physical properties such as elasticity. For tantalum a similar change in electronic structure as for zinc has been found under compression which has little effect on the equation of state but appears in the elastic constants . ## V Conclusions Using the first-principles LAPW method with GGA we calculate the equation of state, structural parameters, and electronic structure of zinc over a wide compression range. We perform calculations for three different k-point samplings of the first Brillouin zone (24x24x24, 32x32x32, and 48x48x48 k-points) and find lattice parameters, in particular the $`a`$-axis, strongly dependent on the number of k-points, while little or no effect can be seen on the equation of state and band structure. For lattice constants we find that a previously observed anomaly in $`a`$-axis compressibility shifts to lower pressure and decreases in amplitude as we increase k-point sampling from 24x24x24 to 48x48x48. This anomaly is not coupled to a change in electronic band structure as has been proposed before; we observe the ETT occurring at approximately the same volume for all sets of computational parameters. The disappearance of the anomaly in lattice constants for our results is in agreement with recent static experiments using helium as a pressure medium. The remaining anomaly in $`a`$-axis compressibility indicates that structural parameters are not fully converged even for the prohibitively large k-point sampling we perform. ## Acknowledgement We greatly appreciate helpful discussion with W. Holzapfel and K. Takemura at different stages of this project and thank H. Krakauer and D. Singh for use of their LAPW code. This work was supported by the National Science Foundation under grants EAR-9614790 and EAR-9980553 (LS), and in part by DOE ASCI/ASAP subcontract B341492 to Caltech DOE W-7405-ENG-48 (REC). Computations were performed on the SGI Origin 2000 at the Department of Geological Sciences at the University of Michigan and the Cray SV1 at the Geophysical Laboratory, support by NSF grant EAR-9975753 and by the W. M. Keck Foundation.
warning/0005/hep-th0005164.html	ar5iv	text	# 1 Introduction ## 1 Introduction In recent developments in string theory, the notion of duality has played a crucial role in the understanding of its nonperturbative properties. Some of the dualities are clearly understood in terms of twelve dimensional F-theory . The “definition” of F-theory is as follows. We consider F-theory compactified on a Calabi-Yau $`(n+1)`$-fold $`X_{n+1}`$ which admits an elliptic fibration $`\pi :X_{n+1}B_n`$, where $`B_n`$ is an $`n`$ complex dimensional manifold. Then, this theory is dual to type IIB string theory compactified on $`B_n`$ with 7-branes wrapped on divisors in $`B_n`$. Via this type IIB description, various string vacua can be analyzed in the context of F-theory. The power of F-theory would be well illustrated by the duality between F-theory and heterotic string. For instance, F-theory compactified on an elliptically fibered $`K3`$ surface is dual to heterotic string compactified on a 2-torus. In this duality, the Narain moduli space on the heterotic side, which contains the kähler and complex moduli of the 2-torus as well as the Wilson line data, coincides with the complex structure moduli space of the $`K3`$ on the F-theory side. The origin of gauge symmetry in 8 dimensions is quite different in these dual theories. In heterotic theory, gauge symmetry comes from the breaking of the $`E_8\times E_8`$ or $`SO(32)`$ gauge group to its subgroups. On the other hand, F-theory acquires unbroken gauge symmetry when the $`K3`$ surface develops the corresponding $`ADE`$ type singularity. The definition of F-theory implies that the 8 dimensional theory given above is represented also by type IIB string theory on $`B_1=𝐏^1`$ with 24 7-branes located at points on the $`𝐏^1`$. The condition that the $`K3`$ surface in the F-theory description develops a singularity, namely the condition that a gauge symmetry enhancement occurs, corresponds to the situation where in type IIB theory several 7-branes coalesce at the same position. The perturbative spectra of heterotic string are identified with string junctions stretched between 7-branes. The 7-brane configurations which give rise to $`ADE`$ type gauge groups have been obtained in . By modding out $`E_8\times E_8`$ heterotic string compactified on a circle by some involution, we can construct another heterotic theory which is now referred to as the CHL string . It has been recently found by Bershadsky, Pantev and Sadov that the CHL string in 8 dimensions also has a dual description in terms of F-theory on a $`K3`$ . The point is that incorporating $`B`$ field background in the F-theory part is equivalent to the CHL involution. Specifically, we switch on NSNS and/or RR 2-form fluxes along $`𝐏^1`$ so that the $`SL(2,𝐙)`$ symmetry of type IIB string is broken to its subgroup $`\mathrm{\Gamma }_0(2)`$ . Then, the complex moduli space of the $`K3`$ is confined to its subspace, which is identical with the moduli space for the CHL string. It should be noted that $`Sp`$ gauge groups are realized in the CHL string in 8 dimensions and its F-theory counterpart. This is a remarkable feature absent in the ordinary vacua where no flux is turned on. <sup>1</sup><sup>1</sup>1 It has also been found in that $`Sp`$ gauge groups appear when several $`D7`$-branes meet an $`O7^+`$-plane. Our aim in this paper is to analyze string junctions in the presence of non-vanishing $`B`$ fields, and to clarify how $`Sp`$ gauge groups are realized in the type IIB framework. The paper is organized as follows. In section 2, we give a brief review of the CHL string in 8 dimensions. In section 3, the F-theory dual description of the CHL string is reviewed. The CHL involution is realized in the F-theory context by the presence of quantized $`B`$ fluxes. We consider type IIB string theory in section 4, and translate the geometric data of F-theory into string junctions. The last section is devoted to conclusions and discussions. ## 2 Review of The CHL String Here we give a review of the eight-dimensional CHL string theory and its possible gauge symmetry enhancements. In the following argument we take the left-moving sector to be supersymmetric, and we use the bosonic realization for the right-moving sector. Let us begin with the ordinary toroidal compactification. In the $`T^d`$-compactification of the $`E_8\times E_8`$ heterotic string we have $`d`$ periodic coordinates $`X^i(z,\overline{z})=X^i(z)+\stackrel{~}{X}^i(\overline{z})`$ and sixteen antiholomorphic bosons $`X^I(\overline{z})`$. Their mode expansions take the following form<sup>2</sup><sup>2</sup>2 Here and throughout the section we use the convention $`\alpha ^{}=2`$. $`X^i(z)`$ $`=`$ $`x^iip^i\mathrm{ln}z+i{\displaystyle \underset{n}{}}{\displaystyle \frac{\alpha _n^i}{nz^n}},`$ $`\stackrel{~}{X}^i(\overline{z})`$ $`=`$ $`\stackrel{~}{x}^ii\stackrel{~}{p}^i\mathrm{ln}\overline{z}+i{\displaystyle \underset{n}{}}{\displaystyle \frac{\stackrel{~}{\alpha }_n^i}{n\overline{z}^n}},`$ (2.1) $`\stackrel{~}{X}^I(\overline{z})`$ $`=`$ $`\stackrel{~}{x}^Ii\stackrel{~}{p}^I\mathrm{ln}\overline{z}+i{\displaystyle \underset{n}{}}{\displaystyle \frac{\stackrel{~}{\alpha }_n^I}{n\overline{z}^n}},`$ and the canonical quantization leads to the following commutation relations: $$\begin{array}{ccccc}[x^i,p_j]& =& [\stackrel{~}{x}^i,\stackrel{~}{p}_j]& =& i\delta _j^i,\hfill \\ [\alpha _m^i,\alpha _n^j]& =& [\stackrel{~}{\alpha }_m^i,\stackrel{~}{\alpha }_n^j]& =& mg^{ij}\delta _{m+n,0},\hfill \end{array}\begin{array}{ccc}[\stackrel{~}{x}^I,\stackrel{~}{x}^J]& =& i\delta ^{IJ},\hfill \\ [\stackrel{~}{\alpha }_m^I,\stackrel{~}{\alpha }_n^J]& =& m\delta ^{IJ}\delta _{m+n,0}.\hfill \end{array}$$ (2.2) The momenta $`(p_i,\stackrel{~}{p}_i,\stackrel{~}{p}_I)`$ are quantized in the following way: $`\stackrel{~}{p}_I`$ $`=`$ $`q_I+A_{iI}w^i,`$ $`\stackrel{~}{p}_i`$ $`=`$ $`n_i+(g_{ij}B_{ij}+A_i^IA_j^I){\displaystyle \frac{w^j}{2}}+A_{iI}q^I,`$ (2.3) $`p_i`$ $`=`$ $`n_i+(g_{ij}B_{ij}+A_i^IA_j^I){\displaystyle \frac{w^j}{2}}+A_{iI}q^I,`$ $$n_i,w^i𝐙,(q^I)\stackrel{}{q}=(\stackrel{}{q}_1,\stackrel{}{q}_2)\mathrm{\Gamma }_8\mathrm{\Gamma }_8.$$ (2.4) Here the (half-)integers $`(n_i,w^i,q^I)`$ correspond respectively to the momenta, the winding numbers and the gauge charges, and specify a point in the Narain lattice $`\mathrm{\Gamma }^{d,d+16}`$. The background fields $`g_{ij},B_{ij},A_{iI}`$ parameterize the following moduli space of vacua: $$O(d,d+16,\mathrm{\Gamma }^{d,d+16})\backslash O(d,d+16)/(O(d)\times O(d+16)).$$ In fact, if we define the norm of the momentum vector by $`p^2p_ip^i\stackrel{~}{p}_i\stackrel{~}{p}^i\stackrel{~}{p}^I\stackrel{~}{p}^I`$, it is indeed independent of the moduli: $$p^2p_ip^i\stackrel{~}{p}_i\stackrel{~}{p}^i\stackrel{~}{p}^I\stackrel{~}{p}^I=2n_iw^iq^Iq^I2𝐙.$$ A choice of a vacuum uniquely determines the decomposition of $`V𝐑^{d,d+16}`$ into two subspaces $`V_\pm `$, each having positive and negative-definite norm. The gauge symmetry enhancements occur precisely when some lattice vectors of squared norm $`2`$ are contained in $`V_{}`$. Thus the analysis of possible gauge enhancements is reduced to a purely mathematical problem. The CHL string is defined as the asymmetric $`𝐙_2`$ orbifold of the $`E_8\times E_8`$ heterotic string theory compactified onto torus. The $`𝐙_2`$ symmetry to be considered is the exchange of two $`E_8`$’s accompanied by a half-period shift. The orbifolding gives rise to two sectors. In the untwisted sector we have the same quantization of the momenta (2.3), (2.4) as in the usual toroidal compactification, and we only take the $`𝐙_2`$-symmetric states as physical states. In the twisted sector we have $$(\stackrel{}{n},\stackrel{}{w})(𝐙^d,𝐙^d+\stackrel{}{v}),$$ where $`\stackrel{}{v}`$ denotes the half-period shift vector. Namely if $`\stackrel{}{v}`$ is taken along the first direction we have $`w^1𝐙+\frac{1}{2}`$. The quantization of the sixteen antiholomorphic bosons gives $$\stackrel{}{q}=(\stackrel{}{q}_1,\stackrel{}{q}_1)(2\stackrel{}{q}_1\mathrm{\Gamma }_8),\stackrel{~}{\alpha }_m^I=()^{2m}\stackrel{~}{\alpha }_m^{I\pm 8}(m𝐙/2).$$ In order for the orbifolding to be possible, the Wilson line must be invariant under the exchange of two $`E_8`$’s, namely $`A_i^I=A_i^{I\pm 8}`$. Hence the moduli space is reduced to the following form: $$O(d,d+8,\mathrm{\Gamma }_{(d)})\backslash O(d,d+8)/(O(d)\times O(d+8)),\mathrm{\Gamma }_{(d)}\mathrm{\Gamma }_8\mathrm{\Gamma }^{d1,d1}(2)\mathrm{\Gamma }^{1,1}.$$ (2.5) $$\mathrm{\Gamma }^{d,d}(2)\left\{\stackrel{}{x}=(\xi _i,\eta ^i)𝐙^{2d}\right\},\stackrel{}{x}^2=4\xi _i\eta ^i.$$ Here the lattice $`\mathrm{\Gamma }_{(d)}`$ is the Narain lattice of the CHL string , and any automorphism of this lattice is a symmetry of the CHL string theory. The way to associate a point of $`\mathrm{\Gamma }_{(d)}`$ to a state with charges $`(n_i,w^i,\stackrel{}{q})`$ is given as follows: $$𝐱(n_i,2w^i,\stackrel{}{Q}\stackrel{}{q}_1+\stackrel{}{q}_2)\mathrm{\Gamma }_{(d)},𝐱^24n_iw^i\stackrel{}{Q}^2.$$ As in the case of ordinary toroidal compactification, a choice of vacuum determines the decomposition of any vector $`𝐱`$ into left and right-moving projection components, $`𝐱=𝐱_++𝐱_{}`$. We define a classification of lattice points $`\left\{𝐱\right\}`$ as follows: type 1. $`𝐱2(\mathrm{\Gamma }^{1,1}\mathrm{\Gamma }_8)\mathrm{\Gamma }^{d1,d1}(2)`$. type 2. $`𝐱`$ is not in the above subset, and $`𝐱^2`$ is divisible by 4. type 3. $`𝐱^2`$ is not divisible by 4. In the table below we summarize the classification of string states according to the above definition. Here the half-period shift is chosen in the first direction. \| (A) \| 2$`w^1`$ even (untwisted sector) \| $`Q2\mathrm{\Gamma }_8`$ \| $`n_1`$ even \| type 1 \| \| --- \| --- \| --- \| --- \| --- \| \| (B) \| \| \| $`n_1`$ odd \| type 2 \| \| (C) \| \| $`Q/2\mathrm{\Gamma }_8`$ \| $`Q^2/2`$ even \| type 2 \| \| (D) \| \| \| $`Q^2/2`$ odd \| type 3 \| \| (E) \| $`2w^1`$ odd (twisted sector) \| \| $`n_1+Q^2/2`$ even \| type 2 \| \| (F) \| \| \| $`n_1+Q^2/2`$ odd \| type 3 \| A careful evaluation of the mass-shell condition yields the following mass formula: the mass of the lightest state among those corresponding to a given vector $`𝐱`$ of type $`i`$ is given by $$m^2=\frac{𝐱_+^2}{2}=\frac{𝐱_{}^2}{2}N_i,(N_i=2,0,1\text{for }i=1,2,3).$$ (2.6) Here $`N_i`$ are obtained as the sum of the zero-point energy and the term $`(\stackrel{}{q}_1\stackrel{}{q}_2)^2`$. Remarkably the mass formula depends only on the type of $`𝐱`$ and not on any of the additional informations. Hence the massless charged vector bosons arise either from the $`𝐱`$ of type 1 and squared norm $`4`$, or of type 3 and squared norm $`2`$. Let us focus on the $`T^2`$ compactification, and take the half-period along the first direction. To study the possible enhancements of gauge symmetry, it is the best way to consider first the charged vector bosons with $`𝐱^2=4`$, corresponding to long roots. From the previous argument it follows that all the long roots $`𝐱`$ arise from the type 1 of the three subsets. Further analysis leads us to the following lemma: (1) Any two long roots are orthogonal. (2) The average of any two long roots belongs to $`\mathrm{\Gamma }_{(2)}`$. These can be proved straightforwardly once we realize that $`n_2,w^2`$ must be odd for any long roots. From these lemma it follows that, given the set of long roots $`𝐱_{j(j=1,\mathrm{},2+n)}`$, the lattice vectors $$\{\pm 𝐱_j,(\pm 𝐱_i\pm ^{}𝐱_j)/2\}$$ form the root system of $`Sp(2+n)`$. Hence the CHL string can realize the $`Sp`$ type gauge symmetry in 8 dimensions. The set of long roots for $`Sp(2+n)`$ can be explicitly constructed in the following way. Take an $`A_n`$ sublattice of $`\mathrm{\Gamma }_8`$ and choose its basis $`e_{i(i=1,\mathrm{}n)}`$ so that they satisfy $`e_ie_j=1+\delta _{ij}`$. Then the long roots are expressed as $$\begin{array}{ccccccc}& & n_1\hfill & \hfill w^1& \hfill n_2& \hfill w^2& \hfill \stackrel{}{Q}\\ 𝐱_1\hfill & =& (0,\hfill & \hfill 0,& \hfill 1,& \hfill 1,& \hfill 0),\\ 𝐱_2\hfill & =& (2,\hfill & \hfill 1,& \hfill 1,& \hfill 1,& \hfill 0),\\ 𝐱_{2+j}\hfill & =& (0,\hfill & \hfill 1,& \hfill 1,& \hfill 1,& \hfill 2e_j).\end{array}$$ (2.7) Generically the enhanced gauge symmetry is given by the product of the $`Sp`$ type level one algebra and the $`ADE`$ type level two algebra. For example, using the formula $`A_n\times E_{8n}E_8`$ we find that the following enhancements are possible: $$Sp(2+n)\times E_{8n}(n=0,\mathrm{},8),$$ (2.8) where we have used the notation $`E_5=SO(10),E_4=SU(5),E_3=SU(3)\times SU(2),E_2=SU(2)\times U(1)`$ and $`E_1=SU(2)`$. ## 3 F-theory Dual of The CHL String In this section we present the F-theory compactification to 8 dimensions which precisely yields a dual description of the CHL string explained in the previous section, following . Since the ordinary Narain moduli space is restricted to its subspace due to the CHL involution, we have to find a similar mechanism which freezes some of the complex moduli of $`K3`$ so as to obtain the same moduli space. This is systematically done by turning on 2-form background fields $`B`$ taking values in $`H^2(𝐏^1,𝐙/2)`$. Let us consider the case where 2-forms $`B_{NS}`$ and $`B_R`$ take the values $$\left(\begin{array}{c}B_{NS}\\ B_R\end{array}\right)=\left(\begin{array}{c}1/2\\ 0\end{array}\right).$$ (3.1) If $`(B_{NS},B_R)`$ is non-vanishing, the $`SL(2,𝐙)`$ monodromy group around 7-branes on $`𝐏^1`$ must be reduced to a smaller group which keeps $`(B_{NS},B_R)`$ invariant modulo $`H^2(𝐏^1,𝐙)`$. In the present case (3.1), the reduced monodromy group is $`\mathrm{\Gamma }_0(2)`$, which is represented by the matrices of the form $$\mathrm{\Gamma }_0(2)=\{\left(\begin{array}{cc}a& b\\ c& d\end{array}\right);a,b,d𝐙,c2𝐙,adbc=1\}.$$ (3.2) The elliptic $`K3`$ surface with a section can be expressed by the familiar Weierstrass equation, which in general has the entire monodromy group $`SL(2,𝐙)`$, not $`\mathrm{\Gamma }_0(2)`$. In order to get the elliptic $`K3`$ with $`\mathrm{\Gamma }_0(2)`$ monodromy, which we denote by $`X`$, the Weierstrass equation must take a special form. To find the form of the equation, note that the elliptic fibration of the $`K3`$ must have one more section which intersects the fiber at the point corresponding to a half-period of the fiber torus. We thus conclude that $`X`$ is of the form $$X:y^2=(xa_4(z))(x^2+a_4(z)x+b_8(z)),$$ (3.3) where $`z`$ is the coordinate on the $`𝐏^1`$ base, and $`a_4(z),b_8(z)`$ are polynomials of degree 4 and 8, respectively. The Weierstrass form (3.3) indeed has a section $`s_1:(x=a_4(z),y=0)`$, in addition to the ordinary one $`s_0:(x=\mathrm{},y=\mathrm{})`$. If we denote by $`\alpha `$ and $`\beta `$ the two independent 1-cycles of the fiber torus which correspond to NSNS and RR directions, then the additional section $`s_1`$ corresponds to the half-period $`(1/2)\alpha +0\beta `$, in agreement with the quantized flux (3.1). The discriminant $`\mathrm{\Delta }`$ of $`X`$ is given by $$\mathrm{\Delta }=(4b_8a_4^2)(b_8+2a_4^2)^2.$$ (3.4) This expression shows that $`X`$ always has eight $`A_1`$ type singularities, and eight $`I_1`$ type singular fibers. All the eight $`A_1`$ singularities are situated at the section $`s_1`$. The $`K3`$ surface $`X`$ given by (3.3) with $`\mathrm{\Gamma }_0(2)`$ monodromy can be constructed from another $`K3`$ surface $`Y`$, through the following two steps. First, we mod out $`Y`$ by an involution $`\sigma `$, to make the “intermediate” $`K3`$ $`\stackrel{~}{X}=Y/\sigma `$. The involution $`\sigma `$ is none other than a geometric realization of the CHL involution. Second, we apply a birational transformation to $`\stackrel{~}{X}`$, yielding the desired $`K3`$ surface $`X`$. Let us start with the construction of $`Y`$. We take the $`K3`$ surface $`Y`$ to be a double cover of the Hirzebruch surface $`𝐅_0=𝐏^1\times 𝐏^1`$, which branches along the zero locus of a section of the line bundle $`𝒪_{𝐅_0}(4,4)`$ over $`𝐅_0`$. Denoting the inhomogeneous coordinates on the two $`𝐏^1`$ factors in $`𝐅_0`$ by $`s`$ and $`t`$, we can represent $`Y`$ as the hypersurface equation $$Y:y^2=\underset{i,j=0}{\overset{4}{}}a_{i,j}s^it^j.$$ (3.5) Here, $`a_{i,j}`$ parametrize the complex structure moduli of $`Y`$. By fixing either one of $`s`$ and $`t`$, one can see that $`Y`$ admits two elliptic fibrations $`p:Y𝐏^1`$ and $`q:Y𝐏^1`$ whose bases are the two $`𝐏^1`$ factors in $`𝐅_0`$ parametrized by $`s`$ and $`t`$, respectively. Let us turn to the first step, the application of the involution $`\sigma `$ to $`Y`$. The involution $`\sigma `$ is defined by the $`𝐙_2`$ action on $`Y`$, $$\sigma :(y,s,t)(y,s,t).$$ (3.6) This means that $`\sigma `$ inverses the orientation of one of the $`𝐏^1`$ and the elliptic fiber over it simultaneously. Since $`Y`$ should be invariant under the $`𝐙_2`$ action (3.6) in order to be modded out by $`\sigma `$, $`Y`$ has to take the form $$y^2=\underset{i=0}{\overset{4}{}}\underset{j=0}{\overset{2}{}}a_{i,2j}s^it^{2j}.$$ (3.7) The hypersurface (3.7) is embedded in $`𝐂^3`$ parametrized by $`(y,s,t)`$ which is reduced to $`𝐂^2/𝐙_2\times 𝐂`$ by the involution $`\sigma `$. The orbifold $`𝐂^2/𝐙_2`$ is described by the three well-defined coordinates $`(u=y^2,v=t^2,w=yt)`$ and the relation among them $$uv=w^2.$$ (3.8) By combining (3.7) with (3.8), we can express $`\stackrel{~}{X}=Y/\sigma `$ by the hypersurface equation parametrized by $`(w,s,v)`$, $`\stackrel{~}{X}:w^2`$ $`=`$ $`vP(s,v),`$ $`P(s,v)`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{4}{}}}{\displaystyle \underset{j=0}{\overset{2}{}}}a_{i,2j}s^iv^j.`$ (3.9) Note that, similarly to $`Y`$, $`\stackrel{~}{X}`$ can also be interpreted as a double cover of $`𝐅_0`$, which is now coordinatized by $`s`$ and $`v`$. The only difference between $`Y`$ and $`\stackrel{~}{X}`$ is that the branch divisor defining $`\stackrel{~}{X}`$ consists of three components $`\stackrel{~}{C}_1:\{v=\mathrm{}\}`$, $`\stackrel{~}{C}_2:\{v=0\}`$ and $`\stackrel{~}{C}_3:\{P(s,v)=0\}`$, whereas the branch divisor for a generic $`Y`$ consists of a single component. In other words, the involution $`\sigma `$ has the effect of restricting the complex moduli space of $`Y`$, so that the branch divisor defining $`Y`$ splits to $`\stackrel{~}{C}_1`$, $`\stackrel{~}{C}_2`$ and $`\stackrel{~}{C}_3`$, which are identified with the zero loci of sections of line bundles $`q^{}𝒪_{𝐏^1}(1)`$, $`q^{}𝒪_{𝐏^1}(1)`$ and $`𝒪_{𝐅_0}(4,2)`$. The second step is to convert $`\stackrel{~}{X}`$ into $`X`$ by a birational transformation. To do this, we focus on the fact that $`\stackrel{~}{X}`$ has eight $`A_1`$ type singularities, four of which are on $`\stackrel{~}{C}_1`$, and the other four on $`\stackrel{~}{C}_2`$. The four $`A_1`$ singularities on $`\stackrel{~}{C}_1`$ comes from the four points where $`\stackrel{~}{C}_3`$ intersects $`\stackrel{~}{C}_1`$. The same is true for the singularities on $`\stackrel{~}{C}_2`$. By blowing up $`𝐅_0`$ at the four points on $`\stackrel{~}{C}_1`$, we obtain a surface birational to $`𝐅_0`$, whose generic $`𝐏^1`$ fiber becomes reducible at the four points on the $`𝐏^1`$ base where $`\stackrel{~}{C}_3`$ meets $`\stackrel{~}{C}_1`$. Each of the four reducible fibers consists of two $`𝐏^1`$’s, one of which is generated by the blowing up procedure. Then, we blow down the remaining four $`𝐏^1`$’s in the reducible fibers. The resulting surface is $`𝐅_4`$, on which there exist three divisors $`C_i(i=1,2,3)`$ which inherit from $`\stackrel{~}{C}_i`$. From the above birational transformation, it is easy to see that $`C_2`$ intersects $`C_3`$ at eight points, while $`C_1`$ does not intersect $`C_3`$ at all. Finally, we identify $`X`$ with the double cover of $`𝐅_4`$, whose branch divisor consists of $`C_i(i=1,2,3)`$. The two sections $`s_0`$ and $`s_1`$ of $`X`$ originate from $`C_1`$ and $`C_2`$, and $`X`$ indeed has eight $`A_1`$ singularities on $`s_1`$ which appear from the eight points where $`C_2`$ and $`C_3`$ intersect. It is also important to note that the elliptic fibration $`\pi `$ for $`X`$ traces back to the elliptic fibration $`p`$ for $`Y`$. The $`K3`$ surface $`X`$ constructed in this way can be represented by the equation (3.3), which has monodromy group $`\mathrm{\Gamma }_0(2)`$. Then we switch on the background flux (3.1), to ensure that the monodromy group of $`X`$ is not enlarged to $`SL(2,𝐙)`$. To determine the unbroken gauge symmetry in F-theory on $`X`$, we need to examine the relation between the singularity structures of $`Y`$ and $`X`$. In what follows, we shall not make an exhaustive analysis, but illustrate the main idea by just one example. The case to be considered in this paper is an $`A_{2n1}`$ type singularity in $`Y`$, which is located on a fixed elliptic fiber $`\{t=0\}`$ with respect to the fibration $`q`$. By appropriately adjusting the coefficients $`a_{i,2j}`$ in (3.7), we can let $`Y`$ develop an $`A_{2n1}(n=1,\mathrm{},8)`$ singularity at the point $`\{s=t=0\}`$. Explicitly, the $`K3`$ surface $`Y`$ with such a singularity is given by $`A_1`$ $`:`$ $`y^2=t^4+t^2+s^2,`$ $`A_3`$ $`:`$ $`y^2=t^4+st^2+s^2,`$ $`A_5`$ $`:`$ $`y^2=st^4+t^4+2st^2+s^2,`$ $`A_7`$ $`:`$ $`y^2=s^2t^4+t^4+2st^2+s^2,`$ $`A_9`$ $`:`$ $`y^2=s^3t^4+t^4+2st^2+s^2,`$ $`A_{11}`$ $`:`$ $`y^2=s^4t^4+t^4+2st^2+s^2,`$ $`A_{13}`$ $`:`$ $`y^2=s^4t^4+s^3t^4+2s^2t^4+t^4+s^4t^2+2s^3t^2+2st^2+s^2,`$ $`A_{15}`$ $`:`$ $`y^2=s^3t^4+t^4+s^4t^2+2st^2+s^2.`$ (3.10) If we resolve the above $`A_{2n1}`$ singularity, there appear, from the singularity, $`2n1`$ $`𝐏^1`$’s which intersect one another according to the $`A_{2n1}`$ Dynkin diagram. Let us denote these $`2n1`$ components by $`S_i(i=1,\mathrm{},2n1)`$. Under the involution $`\sigma `$, the middle component $`S_n`$ remains unchanged, whereas the other $`2n2`$ components are exchanged in the manner $$S_iS_{2ni}(i=1,\mathrm{},n1).$$ (3.11) This suggests that $`\sigma `$ acts on $`S_i`$’s as an outer automorphism of the $`A_{2n1}`$ algebra, giving rise to the new algebra $`C_n`$. Therefore we expect that $`Sp(n)`$ gauge symmetry enhancement takes place in F-theory on $`X`$, when $`Y`$ has an $`A_{2n1}`$ type singularity on $`\{t=0\}`$ (or on $`\{t=\mathrm{}\}`$). The next problem is then to determine what type of singularity $`X`$ should develop in order for $`Sp(n)`$ gauge symmetry to be produced. This is easily done by just translating the singularity data (3) for $`Y`$ into the one for $`X`$. It turns out that $`\stackrel{~}{X}`$ develops an $`A_3`$ (resp. $`D_{n+2}(n=2,\mathrm{},8)`$) singularity when $`Y`$ has a singularity of the type $`A_1`$ (resp. $`A_{2n1}`$). The explicit form of $`\stackrel{~}{X}`$ can be determined by using (3.7), (3) and (3) to be $`A_3`$ $`:`$ $`w^2=v(v^2+v+s^2),`$ $`D_4`$ $`:`$ $`w^2=v(v^2+sv+s^2),`$ $`D_5`$ $`:`$ $`w^2=v(sv^2+v^2+2sv+s^2),`$ $`D_6`$ $`:`$ $`w^2=v(s^2v^2+v^2+2sv+s^2),`$ $`D_7`$ $`:`$ $`w^2=v(s^3v^2+v^2+2sv+s^2),`$ $`D_8`$ $`:`$ $`w^2=v(s^4v^2+v^2+2sv+s^2),`$ $`D_9`$ $`:`$ $`w^2=v(s^4v^2+s^3v^2+2s^2v^2+v^2+s^4v+2s^3v+2sv+s^2),`$ $`D_{10}`$ $`:`$ $`w^2=v(s^3v^2+v^2+s^4v+2sv+s^2).`$ (3.12) In the present situation, the singularity structures of $`X`$ and $`\stackrel{~}{X}`$ are the same. We are thus led to the following dictionary relating the type of singularity in $`X`$ to gauge symmetry: $$\begin{array}{cc}\mathrm{singularity}\mathrm{in}X& \mathrm{gauge}\mathrm{symmetry}\\ & \\ A_3& Sp(1)\\ D_{n+2}(n=2,\mathrm{},8)& Sp(n)\end{array}$$ (3.13) One might wonder why $`A_3`$ (resp. $`D_{n+2}`$) singularity does not give rise to $`SU(4)`$ (resp. $`SO(2n+4)`$), but to $`Sp`$ gauge symmetry. This is due to the presence of quantized fluxes along the base of $`X`$. Recall that $`X`$ necessarily has eight $`A_1`$ singularities and eight $`I_1`$ fibers. The $`A_1`$ singularities cannot be deformed, since deforming them would force $`X`$ to have the full monodromy group $`SL(2,𝐙)`$, and therefore conflict with the presence of background fluxes. The $`A_3`$ singularity in $`X`$ is made from the collision of two $`A_1`$ singularities. Similarly, the $`D_{n+2}`$ singularity in $`X`$ comes from the collision of two $`A_1`$ singularities and $`n`$ $`I_1`$ fibers. If the $`A_1`$ singularities could be deformed, then $`SU(4)`$ or $`SO(2n+4)`$ gauge symmetry would be generated as it should be. However, since deforming the $`A_1`$ singularities is in fact forbidden, it is possible that $`Sp`$ gauge groups are realized though it seems strange at first sight. In the next section, we will find further evidence that $`Sp(n)(n=1,\mathrm{},8)`$ gauge symmetry is actually generated, by explicitly constructing the string junctions which realize the simple roots for $`Sp(n)`$ Lie algebra. ## 4 String Junction In this section we consider the type IIB string theory compactified on $`B_1=𝐏^1`$. There are 24 7-branes located at points on the $`𝐏^1`$, and we turn on the background $`B`$ field as explained in the previous section. The nonzero $`B`$ fields do not break the supersymmetry, because the supersymmetry transformation law of fermions contain only their field strengths $`H=dB`$, which simply vanish in this case. We will focus on the vicinity of the 7-branes realizing the $`Sp`$ type gauge symmetry. Indeed, such an enhancement of the gauge symmetry is possible in the CHL string theory, and we may as well expect the same symmetry also in the F-theory by constructing the $`K3`$ surface $`X`$ from the other $`K3`$ surface $`Y`$ by an orbifolding. However, it is not clear whether we can realize $`Sp(n)`$ gauge symmetry by the above construction of $`X`$, since it is only $`X`$ and not $`Y`$ which is of physical significance. This is true even when we consider the M-theory and membranes by further compactification on $`S^1`$. In the following we explain that $`Sp(n)(n2)`$ gauge symmetry can be realized by an $`SO(2n+4)`$-type fiber singularity, by using string junctions and taking account of the background $`B`$ field appropriately. The case of $`Sp(1)SU(2)`$ will be discussed at the end of this section. We follow the notation of and in particular denote by $`X_{[p,q]}`$ the $`[p,q]`$ 7-brane or the 7-brane on which $`[p,q]`$-string can end. We also denote the 7-branes $`X_{[1,0]},X_{[1,1]},X_{[1,1]}`$ and $`X_{[0,1]}`$ as $`𝐀,𝐁,𝐂`$ and $`𝐃`$, respectively, for notational simplicity. The monodromy of a $`[p,q]`$ 7-brane is given by the formula $$K_{[p,q]}=\left(\begin{array}{cc}1+pq& p^2\\ q^2& 1pq\end{array}\right),$$ (4.1) and the 7-branes $`𝐀,𝐁,𝐂`$ and $`𝐃`$ have the following monodromy matrices $$K_𝐀=\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right),K_𝐁=\left(\begin{array}{cc}0& 1\\ 1& 2\end{array}\right),K_𝐂=\left(\begin{array}{cc}2& 1\\ 1& 0\end{array}\right),K_𝐃=\left(\begin{array}{cc}1& 0\\ 1& 1\end{array}\right).$$ (4.2) The $`D_{n+2}(n2)`$ singularity of $`X`$ (3.12), which is necessary for realizing the $`Sp(n)`$ gauge symmetry, corresponds to the collection of 7-branes $`𝐀^{n+2}\mathrm{𝐁𝐂}`$ in the type IIB picture. The $`D_{n+2}`$ singularity can be deformed to two $`A_1`$ singularities and $`n`$ $`I_1`$ fibers. However, further resolution is not allowed due to the presence of the background $`B`$ field. In order to see this in terms of 7-brane configurations, we notice the equivalence between $`𝐀^{n+2}\mathrm{𝐁𝐂}`$ and $`𝐀^n𝐃^2𝐂^2`$. This can be seen as follows. First, by moving a 7-brane across the branch cut of another 7-brane we obtain the following formulae $$X_{[p,q]}X_{[r,s]}=X_{[r,s]+(psqr)[p,q]}X_{[p,q]}=X_{[r,s]}X_{[p,q]+(psqr)[r,s]},$$ (4.3) which give, in particular, $`\mathrm{𝐀𝐁}=\mathrm{𝐃𝐀},\mathrm{𝐀𝐃}=\mathrm{𝐃𝐂}`$ and $`\mathrm{𝐂𝐀}=\mathrm{𝐃𝐂}`$. Using them we obtain $$𝐀^{n+2}\mathrm{𝐁𝐂}=𝐀^{n+1}\mathrm{𝐃𝐀𝐂}=𝐀^n\mathrm{𝐃𝐂𝐀𝐂}=𝐀^n𝐃^2𝐂^2.$$ (4.4) From the new expression for the collection of 7-branes we identify the two $`A_1`$ singularities with $`𝐃^2`$ and $`𝐂^2`$, and $`I_1`$ fibers with $`𝐀`$’s. Indeed, $`K_𝐀,K_{𝐃}^{}{}_{}{}^{2}`$ and $`K_{𝐂}^{}{}_{}{}^{2}`$ are elements of $`\mathrm{\Gamma }_0(2)`$ defined in (3.2) although $`K_𝐃`$ and $`K_𝐂`$ are not. The VEV of the $`n+2`$ complex scalars, which are in the multiplets containing the gauge fields for the Cartan subalgebra of the gauge symmetry, represent the positions of the $`n`$ $`𝐀`$’s and relative positions of the two $`𝐃`$’s and two $`𝐂`$’s on the $`𝐏^1`$. Each of these scalars originates from the open string whose two endpoints are on a single 7-brane. The massless BPS states are represented by string junctions which connect the 7-branes collapsing to a single point on $`𝐏^1`$. If there is no background $`B`$ field, the string junctions connecting some of the 7-branes $`𝐀^{n+2}\mathrm{𝐁𝐂}`$ form the root lattice of $`SO(2n+4)`$ . We show in fig. $`1`$ the junctions which correspond to the simple roots<sup>3</sup><sup>3</sup>3 In this paper we use same symbol for a simple root and a junction corresponding to it. $`\alpha _i`$ for $`SO(2n+4)`$. Now we analyze how the junctions are transformed by the chain of rearrangements of 7-branes that brings $`𝐀^{n+2}\mathrm{𝐁𝐂}`$ into $`𝐀^n𝐃^2𝐂^2`$. By the Hanany-Witten effect a $`[p,q]`$-string and a $`[r,s]`$ 7-brane passing through each other create $`(psqr)`$ $`[r,s]`$-strings between the two. Hence the string junction which have $`n_1`$ endpoints on $`X_{[p,q]}`$ and $`n_2`$ on $`X_{[r,s]}`$ are transformed according to the following identity $$n_1[p,q]+n_2[r,s]=n_2\left\{[r,s]+(psqr)[p,q]\right\}+\left\{n_1(psqr)n_2\right\}[p,q]$$ in the process of the rearrangement $`X_{[p,q]}X_{[r,s]}=X_{[r,s]+(psqr)[p,q]}X_{[p,q]}`$. Using this formula we find that the junction corresponding to the $`\alpha _i`$ in fig. 1 is transformed to that in fig. 2. Once we turn on the $`B`$ field, the gauge fields corresponding to the junction $`\alpha _j(njn+2)`$ disappear from the massless spectrum. This is because the subgroup $`SU(2)`$ of the $`SO(2n+4)`$ which is generated by $`\alpha _j`$ and $`\alpha _j`$ contains the fluctuation modes that separate the two $`𝐃`$’s or two $`𝐂`$’s. In other words, since the $`A_1`$ singularities cannot be deformed, the scalars corresponding to their deformation parameters cannot have nonzero VEV and must not remain massless. On the other hand, junctions which have the same number of endpoints on each of the two $`𝐃`$’s (and each of the two $`𝐂`$’s) remain massless, since they correspond to the fluctuations that move $`𝐃^2`$ (and $`𝐂^2`$) altogether without breaking them into constituents. By taking these rules into account, we conclude that the massless spectrum representing the simple roots for $`Sp(n)`$ is generated by $`\stackrel{~}{\alpha }_i\alpha _i(1in1)`$ and $$\stackrel{~}{\alpha }_n\alpha _n+(\alpha _n+\alpha _{n+1}+\alpha _{n+2}),$$ (4.5) as depicted in fig. 3. In particular we can check that $`\stackrel{~}{\alpha }_{n1}\stackrel{~}{\alpha }_n=2`$ and that the root $`\stackrel{~}{\alpha }_n`$ has the norm $`(\stackrel{~}{\alpha }_n)^2=4`$, namely it is a long root. Note that the inner product of the roots is given by minus the intersection number of the corresponding string junctions, i.e. $`𝐉_{\alpha _i}𝐉_{\alpha _j}=\alpha _i\alpha _j`$ where $`𝐉_{\alpha _i}`$ is the junction corresponding to the root $`\alpha _i`$. It is also known that, by compactifying further onto $`S^1`$ and going to the M-theory picture, the intersection number of junctions is identified with that of membranes wrapping on the corresponding two-cycles in $`K3`$. Furthermore we expect the BPS condition to be unchanged when we turn on the background $`B`$ field .<sup>4</sup><sup>4</sup>4 The junctions corresponding to the frozen moduli are also expected to be BPS, but they do not yield massless states according to the argument in the previous paragraph. Thus we can realize $`Sp(n)`$ gauge group in type IIB string theory by 7-branes and the non-vanishing background $`B`$ field. Here we notice that, according to the arguments in , the BPS junction $`𝐉`$ have to satisfy $$𝐉𝐉=2g2+b2,$$ (4.6) where $`g`$ and $`b`$ are, the genus and the number of the boundary, of the corresponding two-cycles in $`K3`$. Thus it seems impossible to have BPS junctions with self-intersection number $`4`$ corresponding to the long roots for $`Sp(n)`$. However, this is not true because every junction corresponding to the long roots consists of two junctions orthogonal to each other. For example, $`\stackrel{~}{\alpha }_n`$ is expressed as the sum of $`\alpha _n`$ and $`(\alpha _n+\alpha _{n+1}+\alpha _{n+2})`$, which would satisfy the BPS condition if the background $`B`$ field were turned off. Moreover the two junctions are equivalent as cycles of the corresponding $`K3`$, since under the presence of the $`B`$ field the $`A_1`$ singularities cannot be deformed. Thus every junction corresponding to a long root is a bound state of two junctions which cannot be separated due to the background $`B`$ field. Note that when we turn off the $`B`$ field, it becomes a marginal bound state and can be separated into two BPS junctions. The 7-brane configuration for $`Sp(1)`$ has an expression slightly different from the ones for the sequence $`Sp(n)(n2)`$. The configuration is given by two pairs of $`𝐃`$ branes, namely $`𝐃^2𝐃^2`$. As pointed out before, this configuration would give $`SU(4)`$ gauge symmetry if both pairs of $`𝐃`$ branes could be separated. However, we have $`Sp(1)`$ symmetry here, since the rule concerning massless spectrum allows only one string junction $`\alpha `$ consisting of two coincident D-strings, which emanate from one of the two $`𝐃^2`$’s and get absorbed into the other as depicted in fig. 4. We identify this junction with the simple root for $`Sp(1)`$. The $`O7^+`$-plane has also been used to produce $`Sp(n)`$ gauge group, and a certain dual description of CHL string theory contains it . In this description, a modification of the BPS condition of string junction is required . In fact this is same as our result $`\mathrm{𝐉𝐉}=4`$ for the long root. The $`O7^+`$-plane might be realized in our picture by eliminating the $`B`$ field by a singular gauge transformation and localize it to a single point on $`𝐏^1`$. By identifying the $`O7^+`$-plane with the point onto which the nonzero $`B`$ field is localized, we might be able to prove the equivalence between the F-theory with quantized flux and the description with the $`O7^+`$-plane. ## 5 Conclusions and Discussions In this paper we have clarified how the $`Sp`$ type gauge symmetry can be realized in type IIB string theory on $`𝐏^1`$ with quantized $`B`$ flux. In the presence of the $`B`$ flux the monodromy around the 7-branes are restricted to the subgroup $`\mathrm{\Gamma }_0(2)`$ of $`SL(2,𝐙)`$, and due to this condition some 7-branes are forced to be bound together. The analysis of the string junctions on the collapsed 7-branes shows that, by excluding the junctions corresponding to the deformation that would enlarge the monodromy group, we can obtain $`Sp(n)`$ gauge group out of the $`D_{n+2}`$ singular fiber. Some junctions in the root system of $`Sp(n)`$ have the self-intersection $`4`$ and at first sight they cannot be BPS saturated. However, a careful look into the corresponding cycles in $`K3`$ leads us to realize that the long roots always correspond to two degenerate cycles, so that they are BPS saturated. It would be a very interesting problem to see the full correspondence between the CHL string and type IIB string theory with nonzero $`B`$ field. To see this, note that the gauge groups of the type (2.8) is possible in the CHL string. However the singular fiber of the type $`E_6`$ or $`E_8`$ is impossible in the type IIB string with the $`B`$ field, since the monodromies around the corresponding collections of 7-branes are not in $`\mathrm{\Gamma }_0(2)`$. In the realization of such gauge groups as well as $`Sp(9),Sp(10)`$ has been given in a very complicated way. Some enhancements require that the monodromy group become further smaller to the group $`\mathrm{\Gamma }(2)`$. It would be interesting to study how such gauge groups are realized by 7-branes and string junctions in the type IIB string theory. On the other hand, according to our construction of $`Sp`$ gauge symmetry the gauge groups of the form $`Sp\times Sp`$ are apparently possible, which is in contradiction with the result in the CHL string and the description with the $`O7^+`$-plane. In the CHL string we can easily see that one cannot realize the product of $`Sp`$’s, since by the lemma (2) of section 2 the embedding of the root lattice of $`Sp\times Sp`$ in $`\mathrm{\Gamma }_{(2)}`$ would result in the embedding of the larger $`Sp`$ gauge group. Translating this fact into our framework we must conclude that, if we realize two $`Sp`$ singularity at two different points on $`𝐏^1`$ we obtain a large $`Sp`$ group, with some junctions linking two distant fiber singularities becoming massless. At present we have no idea for resolving this discrepancy. ## Acknowledgements We would greatly like to thank Y. Yoshida for collaboration at the early stage of this work. We also thank T. Takayanagi for useful discussions. J. H. and K. H. thank K. Ohta and S. Sugimoto for helpful comments. The work of J. H., K. H. and S. T. was supported in part by the Japan Society for the Promotion of Science under the Postdoctoral Research Programs No. 11-09295, No. 12-02721 and No. 11-08864, respectively.
warning/0005/cond-mat0005431.html	ar5iv	text	# Vortices Observed and to be Observed ## 1 INTRODUCTION Phases of superfluid <sup>3</sup>He provided the first example of unconventional superfluidity, which triggered the search for the identification of the unconventional superconductivity and superfluidity first in the systems with heavy fermions, then in high-temperature superconductors, and most recently in a layered superconductor Sr<sub>2</sub>RuO<sub>4</sub> and in a multi-component Bose condensates. Unconventional superfluidity of <sup>3</sup>He-A manifests itself in exotic types of the topological defects and in unusual behavior at low temperature caused by the gap nodes. <sup>3</sup>He-A is the first example of superfluids with the gap nodes in the energy spectrum for quasiparticles. Though existence of nodes in <sup>3</sup>He-A was proven only recently in experiments with the AB interface where the $`T^4`$ behavior has been verified ,$`\mathrm{\backslash }sevenrm^\text{?}`$ these nodes influence the investigation of low-$`T`$ properties of unconventional superconductors. In particular, the observed low-$`T`$ and low-$`B`$ scaling for the specific heat in a mixed state of high-$`T_c`$ superconductors related to the gap nodes and vortices $`\mathrm{\backslash }sevenrm^\text{?}`$ has its history started in 1981, when it was found that gap nodes in <sup>3</sup>He-A lead to finite nonanalytic density of states at $`T=0`$ in the presence of the order parameter texture .$`\mathrm{\backslash }sevenrm^\text{?}`$ The intermediate steps were the Ref. ,$`\mathrm{\backslash }sevenrm^{\text{?}\text{}\text{?}}`$ where the superfluid velocity field has been used instead of a texture, and Ref. $`\mathrm{\backslash }sevenrm^\text{?}`$ where this was extended for the superflow around the vortex in $`d`$-wave superconductor. Another example of influence of <sup>3</sup>He-A is the observation in high-$`T_c`$ superconductors of the fractional flux attached to the tricrystal line, which is the junction of three grain boundaries.$`\mathrm{\backslash }sevenrm^\text{?}`$ This was inspired by the discussion of the flux quantization with the half of conventional magnitude in a heavy fermionic superconductors ,$`\mathrm{\backslash }sevenrm^\text{?}`$ which in turn has been inspired by <sup>3</sup>He-A where the possibility of a vortex with half-integer circulation has been shown .$`\mathrm{\backslash }sevenrm^\text{?}`$ The two examples above are the two sides of the same phenomenon: the nontrivial topology. The gap nodes in momentum space and vortices – the nodes in position space – are topological objects, which can be continuously transformed to each other by transformation in the extended momentum+position space, as also was first discussed in <sup>3</sup>He-A .$`\mathrm{\backslash }sevenrm^\text{?}`$ The same topology is also in the basis of the Standard Model of electroweak and strong interactions in particle physics $`\mathrm{\backslash }sevenrm^\text{?}`$ (see Review $`\mathrm{\backslash }sevenrm^\text{?}`$). As a result, in the vicinity of the gap nodes the <sup>3</sup>He-A liquid as well as the high-temperature superconductors and other similar systems acquire some attributes of relativistic quantum field theory such as Lorentz invariance, gauge invariance, chirality and chiral fermions, etc. The conceptual similarity between these condensed matter systems and the quantum vacuum makes them an ideal laboratory for simulating many effects in high energy physics and cosmology. Now it is believed that the tetragonal layered superconductor Sr<sub>2</sub>RuO<sub>4</sub> has the order parameter of the same symmetry class as <sup>3</sup>He-A $`\mathrm{\backslash }sevenrm^{\text{?}\text{}\text{?}}`$ and thus shares many unusual properties of the latter. In particular, it is assumed that the superconducting state of Sr<sub>2</sub>RuO<sub>4</sub> is chiral, i.e. reflection and time reversal symmetries are spontaneously broken. There are other systems which share some of unusual properties of <sup>3</sup>He-A: the multicomponent atomic Bose-condensates ;$`\mathrm{\backslash }sevenrm^{\text{?}\text{}\text{?}}`$ neutron superfluids; quark matter with colour superconductivity … It is worthwhile to consider where we are now: what is observed from these exotic properties and what is to be observed. ## 2 CONTINUOUS VORTICITY According to the Landau picture of superfluidity, the superfluid flow is potential: its velocity $`𝐯_s`$ is curl-free: $`\times 𝐯_s`$. Later Onsager $`\mathrm{\backslash }sevenrm^\text{?}`$ and Feynman $`\mathrm{\backslash }sevenrm^\text{?}`$ found that this statement must be generalized: $`\times 𝐯_s0`$ at singular lines, the quantized vortices, around which the phase of the order parameter winds by $`2\pi N`$. Discovery of superfluid <sup>3</sup>He-A made further weakening of the rule: The nonsingular vorticity can be produced by the regular texture of the order parameter .$`\mathrm{\backslash }sevenrm^\text{?}`$ The order parameter – the wave function of Cooper pair with spin $`S=1`$ and orbital momentum $`L=1`$ – is determined by two vectors $$\mathrm{\Delta }(𝐤)=\mathrm{\Delta }_0(\sigma \widehat{𝐝})\left(𝐤(\widehat{𝐞}_1+i\widehat{𝐞}_2)\right).$$ (1) The real unit vector $`\widehat{𝐝}`$ enters the spin part of the order parameter, where $`\sigma `$ are the Pauli matrices. The orbital part, which shows the momentum $`𝐤`$ dependence, is described by the complex vector $`\widehat{𝐞}_1+i\widehat{𝐞}_2`$, whith mutually orthogonal real and imaginary parts. The order parameter is intrinsically complex, i.e. its phase cannot be eliminated by a gauge transformation. This violates the time reversal symmetry and leads to the spontaneous angular momentum along the unit vector $`\widehat{𝐥}=\widehat{𝐞}^{(1)}\times \widehat{𝐞}^{(2)}`$. The superfluid velocity of the chiral condensate is $`𝐯_s=\frac{\mathrm{}}{2m}\widehat{e}_i^{(1)}\widehat{e}_i^{(2)}`$ (where $`2m`$ is the mass of the Cooper pair) with continuous vorticity satisfying the Mermin-Ho relation: $`\times 𝐯_s=\frac{\mathrm{}}{4m}e_{ijk}\widehat{l}_i\widehat{l}_j\times \widehat{l}_k`$. Quantization of vorticity is provided by topology of textures: The texture which is homogeneous at large distances $`\mathrm{\backslash }sevenrm^\text{?}`$ has $`4\pi `$ winding (i.e. $`N=2`$) if the $`\widehat{𝐥}`$ field covers all $`4\pi `$ orientations in the soft core of continuous vortex. If the order parameter is a spinor (the “half of vector”), as it occurs in spin-1/2 Bose condensates $`\mathrm{\backslash }sevenrm^\text{?}`$ and in the Standard Model of electroweak interactions, the continuous vortices and correspondingly continuous cosmic strings $`\mathrm{\backslash }sevenrm^\text{?}`$ have twice less winding number, $`N=1`$. The NMR measurements already in 1983 provided the indication for existence of the continuous $`4\pi `$ lines in rotating <sup>3</sup>He-A (see Review $`\mathrm{\backslash }sevenrm^\text{?}`$). However, the first proof that the isolated continuous vortex contains $`N=2`$ circulation quanta has been obtained only recently .$`\mathrm{\backslash }sevenrm^\text{?}`$ Lines with continuous vorticity has been used for “cosmological” experiments. Investigation of the dynamics of these vortices presented the first demonstration of the condensed matter analog of the axial anomaly in relativistic field theory, which is believed responsible for the present excess of matter over antimatter .$`\mathrm{\backslash }sevenrm^\text{?}`$ Another cosmological phenomenon is related to nucleation of these vortices investigated in .$`\mathrm{\backslash }sevenrm^\text{?}`$ Vortex formation marks the helical instability of the normal/superfluid counterflow in <sup>3</sup>He-A (the counterflow is the flow of the normal component with respect to the superfluid one). The instability is described by the same equations and actually represents the same physics (see Review $`\mathrm{\backslash }sevenrm^\text{?}`$) as the helical instability of the bath of right-handed electrons towards formation of the helical hypermagnetic field discussed in .$`\mathrm{\backslash }sevenrm^{\text{?}\text{}\text{?}}`$ This experiment thus supported the Joyce-Shaposhnikov scenario of formation of primordial cosmological magnetic field. Another example of continuous vorticity observed in <sup>3</sup>He-A is the vortex sheet .$`\mathrm{\backslash }sevenrm^\text{?}`$ This is the chain of the so called Mermin-Ho continuous vortices with $`N=1`$ trapped by the soliton. Each vortex is the kink in the soliton structure. The kink can live only within the soliton, precisely as the Bloch line within the Bloch wall in magnets. ## 3 FRACTIONAL VORTICES IN SUPERFLUIDS AND SUPERCONDUCTORS In <sup>3</sup>He-A the fractional vorticity is still to be observed. The discrete symmetry, which supports the half-quantum vortex comes from the identification of points $`\widehat{𝐝}`$, $`\widehat{𝐞}_1+i\widehat{𝐞}_2`$ and $`\widehat{𝐝}`$, $`(\widehat{𝐞}_1+i\widehat{𝐞}_2)`$. In a half-quantum vortex both vectors change sign after circling around the vortex while the order parameter in Eq.(1) returns to its initial value $$\widehat{𝐝}=\widehat{𝐱}\mathrm{cos}\frac{\varphi }{2}+\widehat{𝐲}\mathrm{sin}\frac{\varphi }{2},\widehat{𝐞}_1+i\widehat{𝐞}_2=(\widehat{𝐱}+i\widehat{𝐲})e^{i\varphi /2}.$$ (2) Here $`\varphi `$ is the azimuthal angle of the cylindrical coordinate system; the magnetic field is applied along $`z`$ to keep the $`\widehat{𝐝}`$-vector in $`xy`$ plane. $`\widehat{𝐝}`$ is quantization axis for spin. Since it rotates by $`\pi `$ around the vortex, a “person” living in <sup>3</sup>He-A, who moves around the vortex, insensibly finds its spin reversed with respect to the fixed environment. This is an analog of Alice string in particle physics .$`\mathrm{\backslash }sevenrm^\text{?}`$ A particle which moves around an Alice string continuously flips its charge or parity or enters the “shadow” world .$`\mathrm{\backslash }sevenrm^\text{?}`$ In superconductors the crystalline structure must be taken into account. In simplest representation, which preserves the tetragonal symmetry, the order parameter in Sr<sub>2</sub>RuO<sub>4</sub>: $$\mathrm{\Delta }(𝐤)=\mathrm{\Delta }_0(\widehat{𝐝}\sigma )\left(\mathrm{sin}𝐤𝐚+i\mathrm{sin}𝐤𝐛\right)e^{i\theta },$$ (3) where $`\theta `$ is the phase of the order parameter; $`𝐚`$ and $`𝐛`$ are the elementary vectors of the crystal lattice within the layer. Vortices with fractional $`N`$ can be obtained in two ways. If $`\widehat{𝐝}`$-field is flexible enough, there is analog of $`N=1/2`$ vortex in Eq.(2), in which $`\widehat{𝐝}\widehat{𝐝}`$ and $`\theta \theta +\pi `$ around the vortex .$`\mathrm{\backslash }sevenrm^\text{?}`$ Another is the Möbius strip geometry produced by twisting the crystal axes $`𝐚`$ and $`𝐛`$ around the loop .$`\mathrm{\backslash }sevenrm^\text{?}`$ The closed wire traps the fractional flux, if it is twisted by an angle $`\pi /2`$ before gluing the ends. Since the local orientation of the crystal lattice continuously changes by $`\pi /2`$ around the loop, $`𝐚𝐛`$ and $`𝐛𝐚`$, the order parameter is multiplied by $`i`$ after encircling the loop. The single-valuedness of the order parameter requires that this must be compensated by a change of its phase $`\theta `$ by $`\pi /2`$. As a result the phase winding around the loop is $`\pi /2`$, i.e. $`N=1/4`$. This, however, does not mean that the loop of the chiral $`p`$-wave superconductor traps 1/4 of the magnetic flux $`\mathrm{\Phi }_0`$ of conventional Abrikosov vortex. Because of the breaking of time reversal symmetry in chiral crystalline superconductors, persistent electric current arises not only due to the phase coherence but also due to deformations of the crystal :$`\mathrm{\backslash }sevenrm^\text{?}`$ $$𝐣=\rho _s\left(𝐯_s\frac{e}{mc}𝐀\right)+Ka_ib_i,𝐯_s=\frac{\mathrm{}}{2m}\theta .$$ (4) Magnetic flux trapped by the loop is obtained from the condition of no current in the wire, $`𝐣=0`$ in Eq.(4). Thus, the trapped flux depends on the parameter $`K`$ in the deformation current. In the limit case of $`K=0`$ the flux is $`\mathrm{\Phi }_0/4`$ (or $`\mathrm{\Phi }_0/6`$ if the underlying crystal lattice has hexagonal symmetry). In the nonchiral $`d`$-wave superconductor in layered cuprate oxides the order parameter can be represented by: $$\mathrm{\Delta }(𝐤)=\mathrm{\Delta }_0\left(\mathrm{sin}^2𝐤𝐚\mathrm{sin}^2𝐤𝐛\right)e^{i\theta }.$$ (5) The same twisted vortex loop with $`𝐚𝐛`$ and $`𝐛𝐚`$, leads to the change of sign of the order parameter, which must be compensated by a change of its phase $`\theta `$ by $`\pi `$. This corresponds to the $`N=1/2`$ of circulation quantum. The fractional flux trapped by the loop is now exactly $`\mathrm{\Phi }_0/2`$, since the parameter $`K`$ in Eq.(4) is exactly zero in nonchiral superconductors. The same reasoning gives rise to the $`\mathrm{\Phi }_0/2`$ flux attached to the tricrystal line :$`\mathrm{\backslash }sevenrm^\text{?}`$ around this line one has $`𝐚𝐛`$ and $`𝐛𝐚`$. Observation of the fractional flux different from $`\mathrm{\Phi }_0/2`$ would indicate breaking of the time reversal symmetry .$`\mathrm{\backslash }sevenrm^{\text{?}\text{}\text{?}\text{}\text{?}}`$ Let us also mention that the vortex sheet with fractional elementary vortices has been predicted to exist in chiral superconductors ,$`\mathrm{\backslash }sevenrm^\text{?}`$ even before the experimental identification of the vortex sheet in <sup>3</sup>He-A. The object which traps vorticity is the domain wall separating domains with the opposite orientations of the $`\widehat{𝐥}`$-vector .$`\mathrm{\backslash }sevenrm^\text{?}`$ The kink – the Bloch line – in the domain wall is the vortex with the fractional winding number $`N=1/2`$. When there are many fractional vortices (kinks) trapped, they form the vortex sheet, which as suggested in $`\mathrm{\backslash }sevenrm^\text{?}`$ can be responsible for the flux flow dynamics in the low-T phase of the heavy fermionic superconductor UPt<sub>3</sub>. In conclusion of this Section let us mention an exotic combined object in <sup>3</sup>He-A which still have not been observed – vortex terminated by hedgehog .$`\mathrm{\backslash }sevenrm^{\text{?}\text{}\text{?}}`$ Four equally “charged” half-quantum vortices, i.e. each with $`N=+1/2`$ winding number, can meet each other at one point – the hedgehog in the $`\widehat{𝐥}`$ field (see Review $`\mathrm{\backslash }sevenrm^\text{?}`$). Such combined object, which reminds the Dirac magnetic monopole with one or several physical Dirac strings, is called a nexus in relativistic theories .$`\mathrm{\backslash }sevenrm^\text{?}`$ In electroweak theory the monopole-antimonopole pair connected by $`Z`$ string is called a dumbell .$`\mathrm{\backslash }sevenrm^\text{?}`$ These objects can exist in chiral superconductors too, where the nexus has a magnetic charge emanating radially from the hedgehog: this charge is compensated by the flux supplied to the hedgehog by $`N=1/2`$ Abrikosov vortices. Nexus provides a natural trap for massive magnetic monopole – the ‘t Hooft-Polyakov magnetic monopole. If such a monopole enters the chiral superconductor, it is bound to the hedgehog by two $`N=1`$ or four $`N=1/2`$ vortices .$`\mathrm{\backslash }sevenrm^\text{?}`$ ## 4 VORTICES IN ATOMIC BOSE-CONDENSATES In a vector Bose-condensate with the hyperfine spin $`F=1`$ an order parameter consists of 3 complex amplitudes of spin projections $`M=(+1,0,1)`$. They can be organized to form the complex vector $`𝐚`$: $$\mathrm{\Psi }_\nu =\left(\begin{array}{c}\mathrm{\Psi }_{+1}\\ \mathrm{\Psi }_0\\ \mathrm{\Psi }_1\end{array}\right)=\left(\begin{array}{c}\frac{a_x+ia_y}{\sqrt{2}}\\ a_z\\ \frac{a_xia_y}{\sqrt{2}}\end{array}\right).$$ (6) There are two symmetrically distinct phases of the $`F=1`$ Bose-condensates: (i) The chiral state, which mimics the orbital part of the <sup>3</sup>He-A order parameter in Eq.(1), $`𝐚=f(\widehat{𝐞}_1+i\widehat{𝐞}_2)`$, $`\widehat{𝐥}=\widehat{𝐞}_1\times \widehat{𝐞}_2)`$, occurs when the scattering length $`a_2`$ in the scattering channel of two atoms with the total spin 2 is less than that with the total spin zero, $`a_2<a_0`$ .$`\mathrm{\backslash }sevenrm^{\text{?}\text{}\text{?}}`$ (ii) The nonchiral state, which mimics the spin part of the <sup>3</sup>He-A order parameter in Eq.(1), $`𝐚=f\widehat{𝐝}e^{i\theta }`$, occurs for $`a_2>a_0`$. Each of the two states shares some properties of superfluid <sup>3</sup>He-A. In particular, the chiral state (i) has continuous $`4\pi `$ vortices ,$`\mathrm{\backslash }sevenrm^{\text{?}\text{}\text{?}}`$ which are called skyrmions there. An optical method to create half of the full skyrmion – the Mermin-Ho continuous vortices with $`2\pi `$ winding – in the $`F=1`$ Bose-condensates has been recently discussed in Ref. .$`\mathrm{\backslash }sevenrm^\text{?}`$ The nonchiral state (ii) may contain $`N=1/2`$ vortex – the combination of $`\pi `$-vortex in the phase $`\theta `$ and $`\pi `$-disclination in the vector $`\widehat{𝐝}`$ as in Eq.(2):$`\mathrm{\backslash }sevenrm^\text{?}`$ $$𝐚=f\left(\widehat{𝐱}\mathrm{cos}\frac{\varphi }{2}+\widehat{𝐲}\mathrm{sin}\frac{\varphi }{2}\right)e^{i\varphi /2}.$$ (7) In spin projection representation the order parameter far from the core is $$\mathrm{\Psi }_\nu =fe^{i\varphi /2}\left(\begin{array}{c}e^{i\varphi /2}\\ 0\\ e^{i\varphi /2}\end{array}\right)=f\left(\begin{array}{c}e^{i\varphi }\\ 0\\ 1\end{array}\right).$$ (8) This means that the $`N=1/2`$ vortex can be represented as a vortex in the $`M=+1`$ component, while the $`M=1`$ component is vortex-free. The same two cases can occur in a mixture of two condensates. The skyrmion and $`N=1/2`$-vortex can be written in the same form: $$\left(\begin{array}{c}\mathrm{\Psi }_{}\\ \mathrm{\Psi }_{}\end{array}\right)=f\left(\begin{array}{c}e^{i\varphi }\mathrm{cos}\frac{\beta (r)}{2}\\ \mathrm{sin}\frac{\beta (r)}{2}\end{array}\right),\widehat{𝐥}=(\mathrm{sin}\beta \mathrm{cos}\varphi ,\mathrm{sin}\beta \mathrm{sin}\varphi ,\mathrm{cos}\beta ).$$ (9) For skyrmion one has $`\beta (\mathrm{})=0`$ and $`\beta (0)=\pi `$. This represents the $`N=1`$ vortex in condensate $`\|>`$, whose soft core is filled by the condensate $`\|>`$ as was observed in Ref..$`\mathrm{\backslash }sevenrm^\text{?}`$ Note that this is the full skyrmion, i.e. the $`\widehat{𝐥}`$ vector sweeps the whole sphere, however it has $`N=1`$ winding instead of $`N=2`$ in a vector Bose-condensate .$`\mathrm{\backslash }sevenrm^\text{?}`$ The $`N=1/2`$-vortex is obtained if $`\beta (\mathrm{})=\pi /2`$: it is the $`N=1`$ vortex in the $`\|>`$ component, while the $`\|>`$ component is vortex-free. ## 5 FERMIONS ON EXOTIC VORTICES Interesting phenomena to be observed are related to behavior of fermionic quasiparticles in topologically nontrivial environments provided by fractional vortices and monopoles. In the presence of a monopole the quantum statistics of particles can change, e.g. isospin degrees of freedom are trasformed to spin degrees .$`\mathrm{\backslash }sevenrm^\text{?}`$ There are also the fermion zero modes: bound states of fermions at monopole or vortex with exactly zero energy. The low-energy fermions bound to the vortex core play the main role in the thermodynamics and dynamics of the vortex state in superconductors and Fermi-superfluids. The spectrum of low-energy bound states in the core of the $`N=\pm 1`$ vortex in $`s`$-wave superconductor was obtained in Ref.$`\mathrm{\backslash }sevenrm^\text{?}`$ $$E_n=\omega _0\left(n+\frac{1}{2}\right).$$ (10) This spectrum is two-fold degenerate due to spin degrees of freedom. The integral quantum number $`n`$ is related to the angular momentum of the fermions $`n=NL_z`$. The level spacing is small compared to the energy gap of the quasiparticles outside the core, $`\omega _0\mathrm{\Delta }^2/E_F\mathrm{\Delta }`$, and the discrete nature of the spectrum becomes revealed only at $`T\omega _0`$. There are no quasiparticle states with energies below $`\omega _0/2`$. The latter, however, can appear if the core structure is deformed ,$`\mathrm{\backslash }sevenrm^\text{?}`$ or in some regions of the magnetic field due to Zeeman splitting :$`\mathrm{\backslash }sevenrm^\text{?}`$ one or several energy levels cross zero as a function of momentum $`p_z`$, and one obtains the 1D Fermi liquid(s) living in the vortex core, which can exhibit all the exotic properties of 1D Fermi systems including Peierls instability $`\mathrm{\backslash }sevenrm^{\text{?}\text{}\text{?}}`$ and Luttinger liquid physics .$`\mathrm{\backslash }sevenrm^\text{?}`$ $`N=1`$ vortex in <sup>3</sup>He-A has also equidistant energy levels but the Maslov index in Bohr-Sommerfeld quantization is $`0`$ instead of $`1/2`$ in Eq.(10) :$`\mathrm{\backslash }sevenrm^\text{?}`$ $$E_n=\omega _0n.$$ (11) The spectrum now contains the state with $`n=0`$, which has exactly zero energy. This $`E=0`$ level is doubly degenerate due to spin. Such fermion zero mode is robust to any deformations, which preserve the spin degeneracy .$`\mathrm{\backslash }sevenrm^\text{?}`$ The most exotic situation will occur for the $`N=1/2`$ vortex in layered chiral superconductors. As we discussed above, the $`N=1/2`$ vortex is equivalent to the $`N=1`$ vortex in only one spin component. That is why each layer of superconductor contains only one energy level with $`E=0`$. Since the $`E=0`$ level can be either filled or empty, there is a fractional entropy $`(1/2)ln2`$ per layer per vortex. The factor (1/2) appears because in superconductors the particle excitation coincides with its antiparticle (hole), i.e. the quasiparticle is a Majorana fermion .$`\mathrm{\backslash }sevenrm^\text{?}`$ Majorana fermions at $`E=0`$ level lead to the non-Abelian statistics of $`N=1/2`$ vortices: the interchange of two point vortices becomes identical operation (up to an overall phase) only on being repeated four times .$`\mathrm{\backslash }sevenrm^\text{?}`$ This can be used for quantum computing .$`\mathrm{\backslash }sevenrm^\text{?}`$ Also the spin of the vortex in a chiral superconductor can be fractional, as well as the electric charge per layer per vortex ,$`\mathrm{\backslash }sevenrm^\text{?}`$ but this is still not conclusive. The problem with the fractional charge, spin and statistics, related to the topological defects in chiral superconductors is still open. This is related also to the problem of the quantization of physical parameters in 2D systems, such as Hall conductivity and spin Hall conductivity .$`\mathrm{\backslash }sevenrm^{\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}}`$ ## ACKNOWLEDGMENTS I thank B. Anderson for illuminating discussions. This work was supported in part by Russian Foundation for Fundamental Research and ESF.
warning/0005/astro-ph0005119.html	ar5iv	text	# Particle Heating by Nonlinear Alfvénic Turbulence in ADAFs ## 1. Introduction The Advection-Dominated Accretion Flows (ADAFs) are a class of hot, optically thin accretion solutions Ichimaru 1977; Rees et al. 1982; Narayan & Yi 1995; Abramowitz et al. 1995; see e.g., Narayan, Mahadevan, & Quataert 1998 for review which describe quite well the spectral characteristics of a number of low luminosity accreting black hole systems, e.g., black hole binaries (Narayan, Barret, & McClintock 1997; Hameury et al. 1997) and low luminosity galactic nuclei e.g., Sgr A: Narayan, Yi, & Mahadevan 1995; Manmoto et al. 1997; Mahadevan 1998; Narayan et al. 1998; NGC 4258: Lasota et al. 1996; Gammie, Narayan, & Blandford 1999; M87 and other ellipticals: Reynolds et al. 1996; Mahadevan 1997; Di Matteo et al. 1999. In ADAFs, the protons and electrons are thermally decoupled (Coulomb collisions are rare), so that all the energy generated by turbulent viscous stresses is stored as thermal energy of the protons and ultimately advected beyond the black hole horizon, while the electrons remain relatively cool and radiate much less energy, hence the sub-Eddington luminosities. In such hot accretion flows, however, the particle mean free path is often comparable to the size of the system, and thus collective plasma effects are likely to be significant (Rees et al. 1982). In the absence of collisions, strong MHD turbulence is required for the angular momentum transport and energy dissipation in accretion flows. The way MHD turbulence dissipates in hot accretion flows turns out to be crucial for the relevance of ADAF models. Let’s denote $`P_e`$ and $`P_p`$ the amounts of energy that heat the electrons and protons, respectively. In traditional ADAF models, the branching ratio parameter $`\delta =P_e/P_p`$ is commonly assumed to be less than or similar to $`10^2`$. The predicted spectra are only weakly sensitive to the exact value of $`\delta `$ so long as $`\delta 10^2`$, while for larger values of $`\delta `$ the changes in spectrum characteristics are drastic because the electrons become hot and radiate a large fraction of the internal energy of the flow via the synchrotron emission so that the flow is no longer advection-dominated. The above constraint on the value of $`\delta `$ has been recently relaxed in the modified ADAF model with wind outflows (Blandford & Begelman 1999). Because of the lower accretion rates (due to the mass loss) at the inner regions, where most of the radiation is produced, the density of the inflowing gas and the magnetic field strength decrease (assuming constant $`\beta `$) and result in a much lower radiative efficiency of the flow. It was demonstrated by Quataert & Narayan (1999) that the values of $`\delta `$ as high as $`\delta 0.3`$ are still plausible in the advection-dominated flows with strong winds. The question of turbulent particle heating and non-thermal coupling of electrons and protons in hot accretion flows has been addressed in several works (Quataert 1998; Quataert & Gruzinov 1999; Bisnovatyi-Kogan & Lovelace 1997; Begelman & Chiueh 1988; Narayan & Yi 1995). To proceed further, we need to understand what type of MHD turbulence is most likely to be realized in ADAFs. It is commonly believed that MHD turbulence in accretion flows is generated by the magneto-rotational instability (Balbus & Hawley 1991), which generates the large-scale magnetic field as well as produces random gas motions and strong fluctuating magnetic fields in the nonlinear regime. Hence, a variety of plasma waves may be generated. Only low-frequency waves may be efficiently excited by large scale motions and carry a significant fraction of the turbulent energy. Quataert (1998) shows that among the low-frequency MHD waves, fast and slow magnetosonic waves are strongly damped in the two-temperature $`T_p10^{12}\text{ K},T_e10^9\text{ K}`$ ($`T_p`$ and $`T_e`$ are the proton and electron temperatures) ADAF plasma by collisionless dissipation (known as Landau damping and its magnetic analog the transit time damping) while Alfvén waves are still very weakly damped in the incompressible approximation. As argued by Quataert (1998), the “impedance mismatch” will likely inhibit excitation of heavily damped modes and thus the MHD turbulence in ADAFs will be (predominantly) Alfvénic. Additionally, it is natural to believe that the turbulent fluctuating magnetic fields are of order of the mean, large scale magnetic field, $`\stackrel{~}{B}B_0`$, and the Balbus-Hawley instability is nonlinearly saturated. A detailed study particle heating processes by Alfvénic turbulence in ADAFs neglecting the effects of compressibility has been performed by Quataert (1998) and Quataert & Gruzinov (1999). They assumed that the energy which is injected into the system on large scales comparable to the size of the accretion flow, $`LR`$, will cascade in $`k`$-space to small scales, typically the proton Larmor scale, $`l\rho _p=v_{\mathrm{th}}^{}{}_{p}{}^{}/\mathrm{\Omega }_p`$ (with $`v_{\mathrm{th}}`$ being the thermal particle velocity and $`\mathrm{\Omega }=eB_0/mc`$ being the cyclotron frequency) or even smaller, where it dissipates. Since the Alfvén wave cascade is likely anisotropic (Goldreich & Sridhar 1995), the $`k`$ components are related as $`k_{}k_{}^{2/3}L^{1/3}`$. Thus, when the cascade reaches $`k_{}`$ such that $`k_{}\rho _p1`$, the wave energy dissipates only via Landau damping (Cherenkov resonance) and may preferentially heat either the electrons or the protons, depending on the value of plasma $`\beta `$ ($`\beta =8\pi p/B_0^2`$ is the ratio of gas to magnetic pressure).<sup>1</sup><sup>1</sup>1The plasma physics $`\beta `$ used in this paper is different from that used in ADAF models and is related to it as $`\beta _{\mathrm{adaf}}=\beta /(\beta +1)`$. The cyclotron damping which may be very efficient to heat protons is unimportant in this case, because $`k_{}k_{}`$ and $`\omega =k_{}v_A\mathrm{\Omega }_p`$ (where $`v_A^2=B_0^2/4\pi m_pn`$ is the Alfvén speed and $`n`$ is the particle density), so that the cyclotron resonance ($`\omega k_{}v\mathrm{\Omega }_p=0`$) is satisfied for a negligibly small population of very energetic particles from the tail of the Maxwellian distribution function. Quataert & Gruzinov (1999) also argue that only a fraction of Alfvénic energy may be dissipated on the scale $`k_{}\rho _p^1`$ and the rest of it passes this ‘barrier’. On scales $`k_{}\rho _p^1`$ Alfvén waves may be converted to whistler waves and cascade even further where dissipation on electrons may be dominant. Because the details of how Alfvén waves are converted into whistler waves on scales $`k_{}\rho _p^1`$ are not known, they estimate that the transition from the electron dominated heating regime to the proton dominated one occurs somewhere inbetween $`\beta 5`$ and $`\beta 10^2`$. The above analysis is accurate for low-amplitude turbulence, $`\stackrel{~}{B}/B_01`$, where one can neglect the effect of finite magnetic field pressure, $`\stackrel{~}{B}^2/8\pi `$, (i.e., incompressible plasma) and use the linear theory. In a high-amplitude turbulence, $`\stackrel{~}{B}/B_01`$, which is likely to be present in accretion flows, finite magnetic pressure of waves nonlinearly couples Alfvénic energy to ion-acoustic-like (i.e., density) quasi-mode perturbations, which are, in general, dissipative. In this paper, we show that: 1. The Alfvénic parallel cascade towards higher $`k_{}`$ exists in a compressible MHD turbulence as a result of nonlinear wave steepening (modulational instability), in addition to the perpendicular (incompressible) cascade in $`k_{}`$. This parallel cascade is nonlinearly dissipative because it proceeds via the excitation of ion-acoustic (i.e., compressible) quasi-modes which are damped via the Landau mechanism. Even for relatively small amplitudes of turbulence, $`\stackrel{~}{B}/B_00.2`$, and for $`\beta 3`$, much or almost all magnetic energy of Alfvén waves dissipates before it reaches the perpendicular ‘dissipation barrier’, $`k_{}\rho _p1`$. The perpendicular cascade, thus, turns out to be energetically unimportant. Dissipation in the nonlinear parallel cascade preferentially heats the protons and yields $`\delta \mathrm{few}\times 10^2`$. 2. If $`\beta 3`$, the turbulence bifurcates to another ‘phase’ where the Afvénic dissipation in the parallel cascade is weak and other dissipation mechanisms (at small scales) are required to maintain a steady-state. Since the parallel cascade proceeds up to $`k_{}4(\stackrel{~}{B}/B_0)(\mathrm{\Omega }_p/v_A)`$, the Alfvénic turbulent energy may be efficiently dissipated due the cyclotron damping on protons, if $`\omega =k_{}v_A\mathrm{\Omega }_p`$, i.e., if $`\stackrel{~}{B}/B_00.3`$. Only protons may be heated in this regime. For lower amplitudes, the perpendicular Goldreich-Sridhar cascade is likely to dominate and energy dissipates, as discussed by Quataert (1998) and Quataert & Gruzinov (1999). The three-wave cascade (Ng & Bhattacharjee 1996), if it dominated the above four-wave cascade, turns out to be insignificant for electron energetics. We also consider other ‘higher-order’ nonlinear effects, such as particle trapping, and their relevance to particle heating in ADAFs. The paper is organized as follows. In §2 we discuss the noisy-KNLS model of strong MHD turbulence. In §3 we apply the model to the conditions in ADAFs. We compare the model at hand with other competing turbulent processes in §4. Section §5 is the conclusion. ## 2. The model of nonlinear Alfvén wave turbulence ### 2.1. General considerations It is known that in an incompressible plasma, $`𝐯=0`$, the Reynolds and magnetic stresses (i.e., the fluid and magnetic nonlinearities) in a finite-amplitude Alfvén wave cancel each other exactly. Such a wave behaves as if it is linear. Hence, only the Goldreich-Sridhar Alfvénic perpendicular, $`k_{}`$, cascade occurs in a turbulent regime. The mutual cancellation of Reynolds and magnetic stresses in the wave breaks down when plasma is compressible. Indeed, the magnetic field pressure, $`\stackrel{~}{B}^2/8\pi `$, associated with the wave field exerts an additional stress onto an ionized gas and change its local density, $`n`$. Variations in $`n`$ affect, in turn, the local wave phase speed, $`v_A`$, and introduce a positive nonlinear feedback into the wave dynamics (Cohen & Kulsrud 1971). This nonlinear phase speed–amplitude coupling in nonlinear Alfvén waves results in the modulational instability which is ultimately responsible for the wave-front steepening, parallel, $`k_{}`$, Alfvénic cascade, and formation of collisionless shocks. The compressional nonlinearity may also be viewed as an effective parametric coupling of a finite-amplitude Alfvén wave to ion-acoustic (i.e., density) wave-like perturbations in a medium.<sup>2</sup><sup>2</sup>2 Such ion-acoustic modes are not plasma eigenmodes. They are driven by the Alfvén wave and, hence, propagate with the Alfvén speed, which is, in general, different from the sound speed in plasma. These ion-acoustic quasi-modes are compressional and have a longitudinal component of the electric field (with respect to the ambient magnetic field). Hence, they efficiently dissipate the energy via Landau damping. A quantitative theory of nonlinear Alfvén waves is now well established see a short review by Medvedev 1999 and references therein. Note that the dissipation of the finite-amplitude Alfvén waves is intrinsically nonlinear process and leads to a number of unusual properties of such waves (Medvedev 1999). ### 2.2. The model There is only one (to our knowledge) model of strong nonlinear Alfvén wave turbulence in a compressible plasma which self-consistently includes the effects of wave nonlinearity and particle kinetics, e.g., Landau damping (Medvedev & Diamond 1997). It is formulated as a stochastic modification of the dynamic equation of Alfvén waves, — the kinetic nonlinear Schrödinger equation (KNLS). This model of turbulence, referred to as the “noisy-KNLS,” is thus structure-based, meaning that it describes the turbulence as a collection of strongly interacting coherent structures, such as shocks, solitons, cnoidal waves, rotational discontinuities, etc., generated by an external random noise source.<sup>3</sup><sup>3</sup>3In this respect, the model is analogous to the noisy-Burgers model of hydrodynamic turbulence which describes the turbulence as a collection of randomly interacting collisional shocks. The noisy-KNLS equation is $$\frac{b}{\tau }+v_A\frac{}{x}\left(N_1b\stackrel{~}{\|b\|^2}+N_2b\widehat{}\left[\stackrel{~}{\|b\|^2}\right]\right)+i\frac{v_A^2}{2\mathrm{\Omega }_p}\frac{^2b}{x^2}=\stackrel{~}{f},$$ (1) where $`\widehat{}`$ is the integral operator, referred to as the Hilbert operator (or the Hilbert integral transform), acting on the wave field: $$\widehat{}\left[\stackrel{~}{\|b\|^2}\right]=\frac{1}{\pi }_{\mathrm{}}^{\mathrm{}}\frac{𝒫}{(x^{}x)}\stackrel{~}{\|b(x^{})\|^2}dx^{}.$$ (2) Here also $`b=(\stackrel{~}{B}_y+i\stackrel{~}{B}_z)/B_0`$ is the normalized complex wave magnetic field, $`\stackrel{~}{\|b\|^2}=\|b\|^2\|b\|^2`$ is the fluctuating component of the magnetic pressure, $`\mathrm{}`$ means appropriate averaging over the spatial domain, $`𝒫`$ means principal value integration, $`\stackrel{~}{f}`$ is the random (turbulent) noise, and the coefficients $`N_1`$ and $`N_2`$ are complicated functions of parameters of a plasma and are discussed below. The first nonlinear term in equation (1) describes the nonlinear (ponderomotive) coupling of finite-amplitude Alfvén waves to ion-acoustic modes. Indeed, the first two terms are similar to a continuity equation for the wave magnetic field, $`_\tau b+_x(ub)=0`$, where the “hydrodynamic” velocity, $`u`$, is generated by the wave pressure, $`uv_A\stackrel{~}{\|b\|^2}`$. The process of nonlinear steepening stops when it is balanced by linear Alfvén wave dispersion which occurs at very small scales comparable to proton’s Larmor radius, $`k_{}\rho _p^1`$, as represented by the last, linear in $`b`$, term. The second nonlinear term in equation (1) describes Landau (i.e., collisionless) damping of the induced density perturbations, not the magnetic fluctuations. This collective kinetic effect is represented by the integral Hilbert operator. We emphasize that the damping is nonlocal which can be interpreted as the effect of particle “memory” which is associated with particle transit through the wave packet. Another interesting and important feature is the absence of any particular scale beyond which the damping dominates. Indeed, the Fourier representation of the Hilbert operator $`\widehat{}[f]`$ acting on a function $`f`$ is $`(ik_{}/\|k_{}\|)f_ki\mathrm{sgn}(k_{})f_k`$, i.e., independent of the magnitude of $`k_{}`$. Thus, Landau damping of nonlinear Alfvén waves occurs at all scales; it is scale invariant. We should comment that both nonlinearities result in the wave steepening and formation of sharp fronts (shocks) which is, practically, equivalent to the parallel cascade of wave energy to smaller scales. Unlike the conventional, incoherent cascade paradigm, however, this cascade proceeds through random interactions of coherent wave structures (such as shocks, nonlinear waves, etc.). Hence, collective wave–particle interactions, such as the nonlinear Landau damping and particle trapping also contribute. We now briefly discuss basic assumptions used in the noisy-KNLS model. First, equation (1) is the envelope equation. That is, Alfvén waves themselves (carrier waves) are assumed to be linear and, thus, obey the dispersion $`\omega =k_{}v_A`$ while the amplitude of these waves may vary in space and time and is described by equation (1). The time variable $`\tau `$ is the “slow” time of the large-scale dynamics of the wave envelope, $`\tau =(B_0/\stackrel{~}{B})^2t_A`$, where $`t_A=\omega ^1`$ is the typical “Alfvénic” time. It has been assumed in derivation that $`(\stackrel{~}{B}/B_0)^2`$ is small but finite. Comparison of the predictions of this model with observations of nonlinear Alfvén waves in the solar wind by spacecrafts shows that the KNLS theory is “robust” and supports its validity even on the margin of applicability, $`(\stackrel{~}{B}/B_0)1`$. Second, the noisy-KNLS also assumes that waves are propagating in one direction, only. The existence of counter-propagating waves allows for another type of parametric wave–wave interactions, — the decay instability, due to which an Alfvén wave may decay into an acoustic wave and a counter-propagating Alfvén wave, — which greatly complicates any analytical treatment of the problem. The decay instability, however, may be greatly suppressed in hot, two-temperature ADAF conditions. An acoustic (compressible) wave that is to be generated is heavily damped by collisionless dissipation. The same argument that the “impedance mismatch” inhibits resonant excitation of strongly damped modes suggests here (but does not prove, though) that the parametric resonance discussed above will be inefficient. Third, the KNLS-based theory considers planar waves, only. That is, no wave packet modulations are allowed in the plane perpendicular to the direction of wave propagation; the wave fronts are always flat. The dynamics of waves is, thus, purely one-dimensional. This strong limitation of the dynamics and evolution of wave structures in the Alfvénic turbulence has, probably, very little effect (if any) on the processes of particle heating. Indeed, on large scales, particles move along field lines, only. Thus, heating of particles may occur (in Cherenkov resonance) via increase of their parallel velocity, only. The perpendicular wave evolution, therefore, hardly affects particle energy, unless $`k_{}\rho _p^1`$. Note, the effective one-dimensionality of the problem does not mean that waves may propagate along the ambient magnetic field, $`𝐁_\mathrm{𝟎}`$, only. It has been shown that waves propagating at angle are still described by equation (1) if one formally writes the wave field as $`b=(\stackrel{~}{B}_y+B_0\mathrm{sin}\mathrm{\Theta }+i\stackrel{~}{B}_z)/B_0`$, where $`\mathrm{\Theta }`$ is the angle<sup>4</sup><sup>4</sup>4 Strictly speaking, the angle $`\mathrm{\Theta }`$ cannot be pushed close to $`90^{}`$ because dispersion becomes different for left-hand- and right-hand-polarized waves (dispersion on electrons vs. dispersion on protons). This sets the limit on the angle $`(\pi /2\mathrm{\Theta })\sqrt{m_e/m_p}`$. between $`𝐤`$ and $`𝐁_\mathrm{𝟎}`$. Fourth, the nonlinear coupling coefficients, $`N_1`$ and $`N_2`$, presented in the Appendix are calculated by solving the particle kinetic (Vlasov) equation (Spangler 1989, 1990) and agree with those calculated using the guiding center formalism (Mjølhus & Wyller 1988). Both methods, however, used perturbative techniques and do not include nonlinear effects such as particle trapping. Maxwellian distribution was assumed for bulk particles. Gruzinov & Quataert (1999) showed that in a strong Goldreich-Sridhar turbulence the time-asymptotic distribution function of protons may be different. It is not clear whether this situation occurs in ADAF because (i) the infall (accretion) time may be shorter than the velocity diffusion time required to establish a new distribution and (ii) the emergent distribution should, in general, be anisotropic (because the turbulence is anisotropic) and thus may be unstable against, for instance, fire-hose instability which will isotropize and thermalize the particles. Since Landau damping is a resonant process, it cares about the local slope of a distribution function, but its overall shape is irrelevant. Therefore, using the Maxwellian distribution is reasonable. Finally, we do not specify the origin of random forcing for the moment. In general, involvement of a particular mechanism would force re-ordering in equation (1) and changes in $`N_1`$ and $`N_2`$. This complication is automatically eliminated in the renormalization-group approach, the key idea of which is to use bare nonlinear couplings to compute the noise-renormalized (i.e., turbulent) quantities (Medvedev & Diamond 1997). Going ahead, the reason why the bifurcation point (see next section) seems to be determined by the bare $`N`$’s and not the turbulent ones is simple: both nonlinearities are cubic and renormalize in the same way. ### 2.3. The structure of Alfvénic turbulence We are interested in a steady-state turbulence in which the energy input due to random forcing is balanced by Landau dissipation. Since this dissipation is independent of $`\|k\|`$, it occurs at all scales; thus, no inertial interval exists. An analytical study of the model for arbitrary noise is an extremely laborious task. That is why it has been analyzed only in the simplest case of $`\delta `$-correlated in space and time, zero-mean, white in $`k`$ noise (Medvedev & Diamond 1997). A one-loop renormalization group technique has been utilized. The noise-dependent renormalized turbulent transport coefficients has been obtained. Quite importantly, they are mostly determined by properties of the noise source on the largest scales (i.e., the system size), — the so called infrared divergences. The large-$`k`$ tail of the noise spectrum seems to be relatively insignificant to control global properties of turbulence. Thus, the assumption of white noise is likely to be unproblematic in the context of turbulence in hot accretion flows, because the spectrum of turbulence is not known, anyway. The existence of two distinct phases (or regimes) of compressible Alfvénic turbulence constitutes the most important and interesting prediction of the theory. The system will bifurcate (rather than smoothly transit) from one phase to another, as plasma parameters smoothly change. The bifurcation point is determined by the ratio of the nonlinearity-to-damping coefficients: $$\frac{\|N_1\|}{\|N_2\|}\|_{\mathrm{bif}}1.3.$$ (3) This result has been obtained rigorously from the analysis of the self-consistency of the infrared cutoff scalings of the solutions for turbulent transport coefficients. The physical interpretation, however, is very simple and clear. If $`\|N_1\|/\|N_2\|<1.3`$, the collisionless dissipation always dominates the nonlinear parallel energy cascade (i.e., the wave steepening). Most of the energy injected on a large scale dissipates before it reaches the dispersion scale, and all the energy is dissipated via Landau damping in the Alfvénic cascade. Thus, the steady-state, so-called “hydrodynamic” ($`\omega ,k0`$) turbulence with no steep fronts is predicted. In the opposite case $`\|N_1\|/\|N_2\|>1.3`$, no mathematically rigorous prediction can be made, since no solution to renormalized equations exists. Speculatively, however, it is clear that the collisionless damping is weak compared to the nonlinear steepening. Only a small fraction of the injected energy dissipates, while most of it cascades along $`k_{}`$ to the smallest scales set by the linear dispersion. Since the source continuously injects energy into the system, it will be accumulated at the smallest scales. Hence, a non-steady-state, small-scale turbulence of Alfvénic shocklets is expected. Stationarity may be enforced by including other dissipation mechanisms into the model. The most efficient one is the cyclotron damping on protons because it operates at nearly the same scales (i.e., the proton Larmor radius) where the wave energy is accumulated. Another mechanism which dominates for linear and weakly nonlinear waves ($`\stackrel{~}{B}/B_01`$) is the perpendicular cascade (Goldreich & Sridhar 1995; Ng & Bhattacharjee 1996). The energy may cascade through $`k_{}\rho _p`$ to smaller scales where Alfvén waves are probably converted into whistler waves. How energy is dissipated in this case is unclear (Quataert & Gruzinov 1999). The discussion of these issues and quantitative estimates of the rates of the perpendicular vs. parallel cascade and cyclotron damping are given in §4. ## 3. Application to hot accretion flows ### 3.1. Preliminary notes In this section, we primarily focus on the ADAF solutions which are the only known thermally stable, self-consistent models of hot, two-temperature accretion flows. The self-similar scalings of plasma parameters relevant for our problem are Narayan & Yi 1995, see also Quataert 1998: $`T_p`$ $``$ $`2\times 10^{12}{\displaystyle \frac{\beta }{\beta +1}}r^1\text{ K},`$ (4a) $`T_e`$ $``$ $`10^910^{10}\text{ K},`$ (4b) $`\beta `$ $`=`$ $`const.,`$ (4c) $`B_0`$ $``$ $`10^9\alpha ^{1/2}(1+\beta )^{1/2}m^{1/2}\dot{m}^{1/2}r^{5/4}\text{ G},`$ (4d) $`{\displaystyle \frac{\rho _p}{R}}`$ $``$ $`6\times 10^9\alpha ^{1/2}\beta ^{1/2}m^{1/2}\dot{m}^{1/2}r^{1/4},`$ (4e) $`{\displaystyle \frac{v_A}{c}}`$ $``$ $`0.9(1+\beta )^{1/2}r^{1/2},`$ (4f) $`{\displaystyle \frac{v_r}{c}}`$ $``$ $`0.37\alpha r^{1/2},`$ (4g) where $`\alpha `$ is the Shakura-Sunyaev turbulent viscosity parameter, $`m=M/M_{\mathrm{}}`$ is the mass of the central object in solar mass units, $`\dot{m}=\dot{M}/\dot{M}_{\mathrm{Edd}}`$ is the accretion rate in Eddington units ($`\dot{M}_{\mathrm{Edd}}=1.39\times 10^{18}m\text{ g s}^1`$), and $`r=R/R_S`$ is the local radius in Schwarzschild units ($`R_S=2.95\times 10^5m\text{ cm}`$). Here $`B_0`$ is the strength of the large-scale magnetic field, $`\beta `$ is assumed to be independent of $`R`$ and takes, likely, sub-equipartition values ($`\beta 110`$), $`\rho _p/R`$ is the proton Larmor scale in terms of the local scale of the accretion flow, and $`v_r/c`$ and $`v_A/c`$ are the radial inflow velocity and the non-relativistic Alfvén speed,<sup>5</sup><sup>5</sup>5In the relativistic case, the displacement current cannot be neglected in Maxwell’s equations. The relativistic Alfvén velocity then reads $`v_A^{\mathrm{rel}}=v_A/\sqrt{1+v_A^2/c^2}`$. We neglect this relativistic effect in our paper because it may be significant very close to the inner edge of the flow, near the last stable orbit, $`r3`$, only. respectively, in units of the speed of light. The electron temperature is nearly constant throughout the accretion flow and starts to decrease at large radii, $`r10^210^3`$. The efficient cooling of relativistic electrons prevents higher temperatures at smaller radii. Thus the proton-to-electron temperature ratio ranges from $`T_p/T_e10^3`$ deeply inside the flow at $`R\text{few}\times R_S`$ to $`T_p/T_e1`$ in the outer regions, $`R/R_s10^3`$. The MHD turbulence in hot, collisionless accretion flows is likely to originate as a result of the magneto-rotational instability (Balbus & Hawley 1991). In the linear phase, the Balbus-Hawley instability predominantly amplifies the toroidal component of the large-scale magnetic field, so that the fluctuations are weak, $`\stackrel{~}{B}B_0`$. When the instability reaches nonlinear amplitudes, it becomes large scale, $`k_{BH}R^1`$, so that long wavelength MHD waves are excited and the MHD turbulence results, as seen from numerical simulations see, e.g., a review by Balbus & Hawley 1998 and references therein. This turbulence is generated on large scales and, hence, carries most of the gravitational potential energy of the accreting gas. Since the ADAF gas/plasma is two-temperature, with the protons much hotter than the electrons, all compressional MHD modes are heavily damped by the collisionless (Landau and transit time) damping and thus are hard to excite, as suggested by Quataert (1998). Thus, only noncompressive (Alfvén) MHD modes may be efficiently excited and reach nonlinear amplitudes. The nonlinear saturation of the Halbus-Hawley instability means that the local fluctuating magnetic fields in the accretion flow become comparable to the averaged global magnetic field and prevent its further growth; thus in the absence of dissipation<sup>6</sup><sup>6</sup>6This condition does not place any constraint on the value of $`B_0`$. $`\stackrel{~}{B}/B_01`$. At such amplitudes, noncompressibility of Alfvén waves is no longer a good approximation and the nonlinear theory of strong MHD turbulence discussed in the previous section should be used instead. In this case, Landau damping will decrease the amplitude of turbulence to a level when dissipation balances energy input. Given the noise properties, the amplitude of turbulence may be estimated from Medvedev & Diamond (1997). ### 3.2. Dissipation-dominated parallel cascade As we discussed in the previous section (§2), there is a regime of strong nonlinear Alfvénic turbulence in which all energy that is injected into a system on large scales dissipates in the cascade and does not reach the proton Larmor radius scales. The bifurcation point, which is given by Eq. (3), is a function of $`T_p/T_e`$ and $`\beta `$, only (via $`N_1`$ and $`N_2`$, see Appendix). It is independent of the fluctuation level $`b\stackrel{~}{B}/B_0`$ because steepening and damping are both cubic in $`b`$ and the wave amplitude cancels out. (The rate of dissipation, however, does depend on the amplitude of turbulence, as discussed below in §4.1.) Using the expressions for $`N_1`$ and $`N_2`$ given in Appendix and Eq. (3), we now plot the borderline between the two regimes, as a function of plasma parameters. The $`T_p/T_e`$-$`\beta `$-diagram of state of the nonlinear Alfvénic turbulence is shown in Fig. Particle Heating by Nonlinear Alfvénic Turbulence in ADAFs. The shaded region corresponds to $`\|N_1/N_2\|<1.3`$ and the unshaded region — to $`\|N_1/N_2\|>1.3`$. The bold-faced part of the boundary curve represents the range of the temperature ratio relevant for ADAFs: $`T_p/T_e`$ decreases from $`10^3`$ in the inner regions of the accretion flow at radii $`R10R_S`$ to $`1`$ in the outer regions at $`R10^210^3R_S`$. As is seen from the figure, the boundary between the two phases of turbulence is insensitive to the temperature ratio for the ADAF conditions and set only by plasma $`\beta `$. If $`\beta >\beta _{\mathrm{crit}}2.6`$, the Alfvénic turbulence in an ADAF is in the regime in which the parallel cascade is strongly dissipative. If additionally the parallel cascade is dominant (i.e., faster than any other cascade/dissipation process), then it dissipates all injected energy. Alternatively, for $`\beta <\beta _{\mathrm{crit}}2.6`$, Alfvén wave energy does not dissipate in the parallel cascade (even if it is dynamically dominant) and reaches $`k_{}\rho _p^1`$ where it may be dissipated by other mechanisms, or cascade beyond this point to even smaller scales. ### 3.3. Electron vs. proton heating in the parallel cascade We can also investigate the electron-to-proton heating ratio, $`\delta `$, using the equations for $`N_1`$ and $`N_2`$ as follows. The electron and proton contributions to the wave dissipation are proportional to the imaginary part of the plasma dispersion function, Eq. (A17). The factor in Eq. (A17) with integral over the distribution function has been explicitly incorporated into the wave equation, Eq. (1), as the Hilbert operator. The numerical factor was left in the coefficients, Eqs. (A6) – (A16), and is proportional to $`e^{X_e^2}`$ and $`e^{X_p^2}`$ for the electrons and protons, respectively. By vanishing by hands all those terms which are multiplied by $`e^{X_e^2}`$, we may easily determine how much the electrons contribute to the overall dissipation of waves. Calculaing the ratio $`(1N_2^{})/N_2^{}=\delta `$, where $`N_2^{}`$ is the nonlinear dissipation coupling coefficient without the contribution of the electrons to dissipation, we obtain the heating ratio. Note, that in $`N_2^{}`$ the electrons still contribute to the linear and nonlinear dispersion of waves, as represented by the real part of the plasma distribution function. A contour plot showing the levels of constant heating ratio, $`\delta =.005,\mathrm{\hspace{0.17em}.01},\mathrm{},\mathrm{\hspace{0.17em}.035}`$, is superimposed onto the Alfvénic dominated dissipation region in Fig. Particle Heating by Nonlinear Alfvénic Turbulence in ADAFs. Fig. Particle Heating by Nonlinear Alfvénic Turbulence in ADAFs also depicts the heating ratio as a function of the temperature ratio $`T_p/T_e`$ for two values of $`\beta `$$`\beta =3`$ and $`40`$. Clearly, $`\delta `$ depends on $`\beta `$ only weakly and ranges with temperature from $`.025`$ to $`.005`$ for the typical ADAF conditions. Since, the spectral ADAF models weakly depend on $`\delta `$ so long as $`\delta \text{few}\times 10^2`$, we set a conservative upper limit on the electron heating in the dissipative parallel cascade of nonlinear Alfvénic turbulence: $$\delta _{}\frac{P_e}{P_p}\|_{\mathrm{casc}}.025,\text{ if }\beta >\beta _{\mathrm{crit}}2.6.$$ (5) ## 4. Other competing processes In the previous section (§3.3) we demonstrated that the energy which heats the electrons in the compressible Alfvénic dissipative parallel cascade in typical ADAF conditions constitues a few percents of the total energy which goes into the protons. If this parallel cascade would be the only one which operates, then equation (5) gives a solid prediction. In fact, there are other competing processes which transfer Alfvénic energy in $`k`$-space and thus may affect the energetics. We consider them in order. In addition to the compressible parallel cascade, there is an incompressible perpendicular cascade due to direct interactions of linear Alfvén wave packets. Recently two different theories of this process have been proposed. Goldreich & Sridhar (1995) (hereafter GS) showed that 4-wave interactions lead to a critically balanced anisotropic cascade in which $`k_{}`$ and $`k_{}`$ are uniquely related. A nice paper by Ng & Bhattacharjee (1996) (hereafter NB) demonstrates that 3-wave interactions of $`k_{}=0`$ perturbations are not empty and result in a more conventional Iroshnikov-Kraichnan-type cascade. Although there are some indications from both theoretical arguments and numerical simulations (Goldreich & Sridhar 1997; Cho & Vishniac 2000) that the GS cascade dominates over the NB cascade, the debate is still not completely resolved. For this reason, we consider both possibilities below. Briefly, we obtained that the GS cascade puts some limitations on our predictions of §3.3. The NB cascade, however, does not alter our conclusions about the electron vs proton heating, despite that the total wave energy in the perpendicular direction may be substantial. ### 4.1. Parallel vs. GS cascade The critically balanced GS cascade (Goldreich & Sridhar 1995) is an anisotropic, (almost) perpendicular cascade in which $`k_{}k_{}k_{}^{3/2}`$. This cascade involves linear Alfvén waves and, thus, is independent of the amplitude of the turbulence. No damping occurs in this process. The MHD (Alfvén) energy is transferred from the scales where it is injected, $`k_iR^1`$, to the scales where it is dissipated or converted into other waves, $`k_f\rho _p^1`$. We have to compare it with the competing compressible parallel cascade during which much of injected energy is dissipated. The rate of the energy transfer depends on the turbulence level via the compressibility. Since these two cascade processes are independent (at least in the KNLS approximation), we may compare relative contribution of each of them to the electron heating as a function of the amplitude of MHD turbulence in ADAFs. The characteristic rate (i.e., the inverse $`e`$-folding time) of the GS cascade is the Alfvénic frequency, $`\gamma _{\mathrm{GS}}\omega _A`$. The characteristic dissipation rate in the parallel cascade and the nonlinear cascade rate itself are comparable and estimated from Eq. (1) to be $`\gamma _{\mathrm{NL}}\omega _A\|N_2\|\stackrel{~}{\|b\|^2}`$. Note, although the rate of the GS cascade, $`\gamma _{\mathrm{GS}}k_{}`$, increases with $`k_{}`$, so does the energy dissipation rate. Thus, the total energy dissipated in the cascade is determined by the width of the “inertial range,” $`R^1k\rho _p^1`$. The energy dissipation rate is, by definition, $$\frac{dϵ}{dt}=\frac{dϵ}{dk_{}}\frac{dk_{}}{dt}=\gamma _{\mathrm{NL}}ϵ.$$ (6) Using the definition for the total energy cascade rate, $`dk_{}/dt=(\gamma _{\mathrm{GS}}+\gamma _{\mathrm{NL}})k_{}`$, and integrating the resultant equation, we obtain $$\frac{\mathrm{ln}ϵ/ϵ_0}{\mathrm{ln}k_m/k_0}=\frac{\|N_2\|\stackrel{~}{\|b\|^2}}{\|N_2\|\stackrel{~}{\|b\|^2}+1},$$ (7) where $`ϵ_0`$ and $`k_0`$ are initial values. The largest scale is set by the size of the accretion flow, $`k_0R^1`$. The smallest scale is estimated from the scaling, $`k_{}k_{}^{2/3}R^{1/3}`$, of the GS cascade with $`k_m\rho _p^1`$.<sup>7</sup><sup>7</sup>7 Note, here we set a lower limit on $`k_m`$. Since both cascades proceed independently, $`k_m`$ may only be larger (i.e., closer to $`\rho _p^1`$) due to the parallel cascade. Thus, $`k_m/k_0(R/\rho _p)^{2/3}`$. Quataert & Gruzinov (1999) argue that much of the energy dissipated at $`k_{}\rho _p^1`$ may heat electrons, depending on the value of $`\beta `$ and numerical constants estimated from numerical simulations which are not very well constrained. We assume here the worst possible case in which all the wave energy reaching the Larmor scale heats electrons. For the electron heating to dominate by the compressible cascade, the amount of energy that reaches the Larmor scale must be smaller than the fraction of energy that is dissipated on electrons in the compressible cascade: $$ϵ(k_m)/ϵ_0<\delta _{},$$ (8) where $`\delta `$ is given by Eq. (5). Combining all together, this inequality may now be re-written as follows: $$\stackrel{~}{\|b\|^2}>\frac{1}{\|N_2\|}\frac{\frac{3}{2}\mathrm{ln}\delta _{}/\mathrm{ln}(R/\rho _p)}{1+\frac{3}{2}\mathrm{ln}\delta _{}/\mathrm{ln}(R/\rho _p)}.$$ (9) The nonlinear coupling to dissipation, $`N_2`$, depends on plasma parameters. The dependence of $`\|N_2\|`$ vs. $`\beta `$ for two values of the proton-to-electron temperature ratio: $`T_p/T_e=1,10^3`$ is shown in Fig. Particle Heating by Nonlinear Alfvénic Turbulence in ADAFs. Clearly, it weakly depends on $`T_p/T_e`$ and is roughly $`\sqrt{\beta }`$. To estimate the lower bound on $`\|b\|`$, we take the lowest value of $`\|N_2\|1`$ for $`\beta 3`$, i.e., near the bifurcation value. The normalized amplitude may be written in a general case as follows (see §2.2): $`\stackrel{~}{\|b\|^2}`$ $`=`$ $`(\stackrel{~}{B}+B_0k_{}/k)^2/B_0^2(\stackrel{~}{B}+B_0k_{}/k)^2/B_0^2`$ (10) $`=`$ $`(\stackrel{~}{B}/B_0)^2(\stackrel{~}{B}/B_0)^2+2(\stackrel{~}{B}/B_0)k_{}/k`$ $``$ $`2(\stackrel{~}{B}/B_0)`$ for $`\stackrel{~}{B}B_0`$ and an oblique and nearly perpendicular angle of propagation, $`k_{}k`$. We take \[see Eqs. (5) and (4)\] $`\delta 2.5\times 10^2`$ and $`\rho _p/R6\times 10^9`$. Then the conservative limit on the fluctuation amplitude of MHD turbulence at which the electron heating is dominated by the compressible cascade is readily estimated as $$\frac{\stackrel{~}{B}}{B_0}0.2.$$ (11) Thus, for $`\stackrel{~}{B}/B_00.2`$, the electron heating due to the parallel, $`\delta _{}`$, and perpendicular, $`\delta _{}`$, cascades are of comparable value and the total electron heating is $`\delta =\delta _{}+\delta _{}0.05`$. Hence, for the levels of turbulence which are believed to exist in ADAFs, $`\stackrel{~}{B}/B_01`$, the heating of electrons is dominated by the compressible cascade by at least an order of magnitude and yields $`\delta 0.025`$, provided $`\beta >\beta _{\mathrm{crit}}2.6`$. ### 4.2. Parallel vs. NB cascade In a similar way we consider competition of the NB and parallel cascades. The NB cascade (Ng & Bhattacharjee 1996) proceeds in $`k_{}`$ only, hence it is totally independent of $`k_{}`$. There is no dissipation during the NB cascade. The Alfvén wave energy may be damped only when it reaches the $`k_{}\rho _p^1`$ scale and, as discussed in §4.1, will predominantly heat the electrons. To estimate the energy lost via the NB cascade, it is sufficient to compare the rate, $`\gamma _{\mathrm{NB}}(\rho _p^1)`$, of the NB cascade at $`k_{}\rho _p^1`$ to the maximum damping rate in the $`k_{}`$ cascade, $`\gamma _{\mathrm{NL},\mathrm{max}}`$. The latter is $`\gamma _{\mathrm{NL},\mathrm{max}}\omega _{A,max}\|N_2\|\stackrel{~}{\|b\|^2}v_A\rho _p^1\|N_2\|\stackrel{~}{\|b\|^2}`$. The former is estimated from Ng & Bhattacharjee (1996) as follows. The energy transfer rate is $$\dot{ϵ}v^2/(N\tau )\text{ constant},$$ (12) where $`v^2E_kk`$ is the energy per unit mass, $`N(\delta v/v)`$ is the number of interactions, $`\tau (\kappa _{}v_A)^1`$ is the interaction time, $`\delta v`$ is the perturbation, and $`\kappa _{}`$is a longitudinal scale of interacting Alfvén wave packets. Therefore, $`\gamma _{\mathrm{NB}}`$ $``$ $`(N\tau )^1\dot{ϵ}/v^2\dot{ϵ}(E_kk)^1`$ (13) $``$ $`(E_{k0}k_0)^1\left({\displaystyle \frac{k}{k_0}}\right)^{1/2}\left({\displaystyle \frac{\delta v_0}{v_0}}\right)^2\kappa _0v_A,`$ where we used the NB scaling $`E_kk^{3/2}`$ and the subscript “0” denotes quantities at the outer scale of the turbulence. Assuming strong turbulence, $`\delta v_0v_0`$, we obtain an upper limit on $`\gamma _{\mathrm{NB}}`$. Assuming also that on the outer scale $`k_0\kappa _0\kappa _0\kappa _0R^1`$, we have $$\gamma _{\mathrm{NB}}(\rho _p^1)v_A(\kappa _0k)^{1/2}\|_{k\rho _p^1}v_A(R\rho _p)^{1/2}.$$ (14) Thus, the ratio is $$\frac{\gamma _{\mathrm{NB}}(\rho _p^1)}{\gamma _{\mathrm{NL},\mathrm{max}}}\left(\|N_2\|\stackrel{~}{\|b\|^2}\right)^1\left(\frac{\rho _p}{R}\right)^{1/2}10^4.$$ (15) Thus, even if all the energy cascaded in $`k_{}`$ heats the electrons, it constitutes about $`10^4`$ of the total energy dissipated on the protons (unless $`\stackrel{~}{B}/B_01`$), i.e., is completely negligible compared to equation (5). Note that the above result means only that the NB cascade is “slower” than the parallel cascade, while the total wave energy in $`k_{}`$ may be non-negligible. ### 4.3. Proton-cyclotron damping vs. perpendicular cascade The nonlinear interaction of finite-amplitude wave-packets results in generation of high-$`k`$ harmonics, as discussed in §2.2. Such a cascade transfers wave energy along the local magnetic field to high $`k_{}`$, where the cyclotron resonance, $`\omega k_{}v\mathrm{\Omega }_p=0`$, may be satisfied and the cyclotron damping on (bulk) protons becomes very efficient. This effect was completely missing from previous studies because both the GS and NB cascades are (nearly) perpendicular; $`k_{}`$ is always small, so that only a very small number of particles from the tail of a particle distribution function are in the resonance. The proton-cyclotron damping is known to heat protons only, no energy goes to electrons, since $`\mathrm{\Omega }_p\mathrm{\Omega }_e`$ and electrons are off resonance. It is, however, difficult to rigorously estimate the $`\stackrel{~}{\|b\|^2}`$, because for the cyclotron damping to dominate, it must be faster than the typical GS or NB cascading time, i.e., $`\gamma _c\omega _A`$. Such a situation may occur only when $`\omega _A\mathrm{\Omega }_p`$, so that (i) Alfvén waves are heavily damped and (ii) they may convert to proton cyclotron waves, which greatly complicates a rigorous analysis. One can, however, roughly estimate the fluctuation level at which the proton-cyclotron damping is so strong that most of the wave energy dissipates on protons and does not convert into whistlers. The parallel cascade proceeds until it hits a scale where the nonlinear steepening is balanced by the wave dispersion \[last term in Eq. (1)\]. Thus, the maximum parallel wave-number to which wave energy cascades in the compressible cascade is readily estimated from Eq. (1) to be $`k_{}^{}{}_{\mathrm{max}}{}^{}\stackrel{~}{\|b\|^2}(2\mathrm{\Omega }_p/v_A\|N_{1,2}\|)`$, where $`\|N_{1,2}\|`$ is the maximum of $`\|N_1\|`$ and $`\|N_2\|`$. The proton-cyclotron damping is strong when $`\omega _A=k_{}v_A\mathrm{\Omega }_p`$. Thus, if $`k_{}^{}{}_{\mathrm{max}}{}^{}\mathrm{\Omega }_p/v_A`$, the MHD turbulence heats protons only. This occurs when $$\frac{\stackrel{~}{B}}{B_0}\frac{1}{4\|N_{1,2}\|},$$ (16) where we again took $`\stackrel{~}{\|b\|^2}2(\stackrel{~}{B}/B_0)`$. Here $`\|N_{1,2}\|1`$ for $`\beta 3`$ and increases with $`\beta `$ and weakly with $`T_p/T_e`$ as shown in Fig. Particle Heating by Nonlinear Alfvénic Turbulence in ADAFs. For a range of $`\beta `$’s, $`1\beta 2.6`$, i.e., above the bifurcation threshold (dissipation in the compressible cascade is weak), but below equipartition, the conservative condition on $`\stackrel{~}{B}`$ reads: $$\frac{\stackrel{~}{B}}{B_0}0.3.$$ (17) ### 4.4. Is nonlinear trapping important? The parallel (coherent) cascade of the wave energy occurs due to the formation and nonlinear evolution of coherent wave structures, such as shocks, Alfvénic discontinuities, solitary wave-packets, etc., as discussed in §2.2. The resonant particles, i.e., those which take the energy from waves and result in damping, turn out to be trapped in the wave potential created by the wave magnetic field. Those resonant particles which move a little slower than the wave will be reflected from the rear part of the wave potential, gain energy, and start to move faster than the wave. The particles moving faster than the wave will be reflected from the front part of the wave potential and, correspondingly, give energy to the wave and decelerate. Upon a half bounce time, $`\tau _\mathrm{b}/2`$, the particles which have been reflected from the rear of the potential cross though it and reach its front part, where they are reflected back and give their energy to the wave. The slow particles do the opposite at the rear part. Such a process may repeat over and over again, provided the wave amplitude is maintained constant by an external source (e.g., the Balbus-Hawley instability in ADAFs), to prevent particles from escape. Clearly, the number of slow ($`v<v_A`$) and fast ($`v>v_A`$) particles will oscillate with time, so that the wave damping rate will change in time too and may become negative (i.e., the wave amplitude will grow) during some periods of time. The bounce period of a particle depends on its energy, because a wave trapping potential is, in general, anharmonic. Due to the difference in the bounce periods, groups of particles with different energies gain a phase shift which grows with time. Upon many cycles (bounces), the particles near the resonance, $`vv_A`$, become completely randomized in phase and form an equilibrium, plateau distribution. At this moment, the wave collisionless dissipation quenches.<sup>8</sup><sup>8</sup>8This process does not affect the cyclotron damping on protons, which operates at the cyclotron, and not Cherenkov, resonance. The process described above is referred to as the nonlinear Landau damping see Medvedev et al. 1998 for more details. Numerical simulations of the process show that it takes $`10^210^3\tau _\mathrm{b}`$ or even longer for particles to phase mix and quench collisionless damping. We note here that the nonlinear Landau damping is significant for large amplitude waves $`\stackrel{~}{B}/B_01`$, for which the number of trapped particles may be large. We now estimate whether particle trapping is important for the physics of ADAFs and, in particular, whether the Alfvén wave damping during the parallel (compressible) cascade in ADAFs quenches. First, the kinetic energy of a particle in a potential is of order its potential energy, so that the typical velocity of a particle of mass $`m_j`$ ($`j=e,p`$) is $`v_j(\stackrel{~}{B}^2/8\pi m_jn)^{1/2}v_{A}^{}{}_{j}{}^{}(\stackrel{~}{B}/B_0)`$. The typical scale of the potential, $`\lambda `$, is set by the scale at which the energy is pumped into the system, i.e., the scale of the Balbus-Hawley instability, $`\lambda k_{BH}^1R`$. Thus the characteristic oscillation period of a particle (the bounce time) at a given radial position in the accretion flow is $`\tau _\mathrm{b}\lambda /v_j`$, which yields $$\tau _{\mathrm{b}}^{}{}_{p}{}^{}R/(v_A\stackrel{~}{B}/B_0),\tau _{\mathrm{b}}^{}{}_{e}{}^{}=(m_e/m_p)^{1/2}\tau _{\mathrm{b}}^{}{}_{p}{}^{}$$ (18) for protons and electrons, respectively. Second, accreting gas flows into the central object and, thus, continuously replenished with a fresh material on the outer edge of the ADAF. The typical infall time at a given radius $`R`$ is $$\tau _\mathrm{r}R/v_r,$$ (19) where $`v_r`$ is the radial velocity of a gas given by Eqs. (4). Estimating the number of bounces of a particle of species $`j`$ in an Alfvén wave at a radial position $`R`$ as $`𝒩_j(R)\tau _\mathrm{r}/\tau _{\mathrm{b}}^{}{}_{j}{}^{}`$, we obtain $$𝒩_p0.4\alpha (1+\beta )^{1/2}(\stackrel{~}{B}/B_0),𝒩_e43𝒩_p.$$ (20) Note that $`𝒩`$ is independent of radius. The above equations show that the protons in the accretion flow are always in the linear regime of damping (no trapping, $`𝒩_p1`$). The electrons may experience up to a few tens of bounces within the infall time, which is probably not enough to quench damping on electrons at all, but still may lower its value a little. We thus conclude that the damping process in the compressible parallel cascade is hardly affected by nonlinear particle trapping and the value of $`\delta \mathrm{few}\times 10^2`$ represents a conservative upper limit on the the electron to proton heating ratio in hot, advection-dominated accretion flows with strong magnetic field turbulence, $`\stackrel{~}{B}B_0`$. ## 5. Conclusions The particle heating is one of the key issues for the advection-dominated accretion flow models. All the models assume that protons receive most of the energy released due to viscous (turbulent) heating of the accretion gas, while electrons remain relatively cold; hence the low luminosities of ADAFs. In the collisionless plasma of the ADAFs, particle heating is determined solely by excitation and dissipation of plasma collective motions (waves). The instability of Balbus & Hawley which is believed to operate and produce magnetic fields in such differentially rotating flows may naturally result in strong MHD turbulence with highly fluctuating magnetic fields, $`\stackrel{~}{B}B_0`$, during the nonlinear stages of its evolution. A detailed analysis of particle heating by such turbulence is presented in this paper. We show that in the most natural case of the ADAF parameters, $`\beta 3`$ and $`\stackrel{~}{B}/B_00.2`$, dissipation of the magnetic energy occurs predominantly in the parallel, compressible cascade (the effect absent for linear, low-amplitude Alfvén waves). Most of the dissipated energy heats protons, while electrons receive only $`\delta \mathrm{few}\%`$ of the energy, which in agreement with the empirical values inferred from the spectral fits for various, low-luminosity accreting systems. If, alternatively, the magnetic field in ADAFs is close to equipartition, $`\beta 3`$, the energy of MHD turbulence is dissipated only on protons via the cyclotron damping mechanism, provided the amplitude of the turbulence is somewhat higher, $`\stackrel{~}{B}/B_0\mathrm{few}\times 10^1`$. The electron heating is negligible in this case. At lower amplitudes, however, the nonlinear effects are weak and the results of the linear analysis for the dissipation of the linear Alfvén waves, addressed by other authors, hold. The author is grateful to Ramesh Narayan for his interest in this work and helpful discussions, to the referee Amitava Bhattacharjee for many insightful comments and suggestions which helped to greatly improve the manuscript and stimulated further work, and to Eliot Quataert for discussions. This work was supported by NASA grant NAG 5-2837 and NSF grant PHY 9507695. ## Appendix A Exact formulae for the nonlinear coefficients $`N_1`$ and $`N_2`$ The nonlinear coefficients in Eq. (1) has been calculated from the full kinetic treatment by solving the Vlasov equation for arbitrary wave profile and determining the wave-induced plasma density perturbation (Spangler 1989, 1990). They are $$N_1=M_1m_1,N_2=M_2m_2,$$ (A1) where $`M_1`$ and $`M_2`$ are the isotropic pressure contributions and $`m_1`$ and $`m_2`$ are the contributions due to pressure anisotropy (i.e., the temperatures and, hence, pressures may be different along the ambient magnetic field and perpendicular to it, $`p_{}p_{}`$): $`M_1`$ $`=`$ $`{\displaystyle \frac{1}{4}}\left(f_7+{\displaystyle \frac{f_5\left(f_1f_2\pi f_3f_4\right)+\pi f_6\left(f_3f_2+f_1f_4\right)}{\left(f_2^2+\pi f_4^2\right)}}\right),`$ (A2) $`M_2`$ $`=`$ $`{\displaystyle \frac{\sqrt{\pi }}{4}}\left(f_6+{\displaystyle \frac{f_6\left(f_1f_2\pi f_3f_4\right)f_5\left(f_3f_2+f_1f_4\right)}{\left(f_2^2+\pi f_4^2\right)}}\right),`$ (A3) $`m_1`$ $`=`$ $`{\displaystyle \frac{1}{4X_p^2}}\left(g_1+{\displaystyle \frac{g_4\left(f_1f_2\pi f_3f_4\right)+\pi g_3\left(f_3f_2+f_1f_4\right)}{2\left(f_2^2+\pi f_4^2\right)}}\right),`$ (A4) $`m_2`$ $`=`$ $`{\displaystyle \frac{\sqrt{\pi }}{4X_p^2}}\left(g_2+{\displaystyle \frac{g_3\left(f_1f_2\pi f_3f_4\right)g_4\left(f_3f_2+f_1f_4\right)}{2\left(f_2^2+\pi f_4^2\right)}}\right),`$ (A5) where $`\sqrt{\pi }`$ appears in numerators of $`M_2`$ and $`m_2`$ because the Hilbert operator carries an extra ($`1/\pi `$) and additional (1/2) in all four coefficients absorbs it from the wave equation, as compared to Spangler (1989, 1990). The coefficients $`f_1,f_2,\mathrm{},f_7,g_1,\mathrm{},g_4`$ are defined as $`f_1`$ $`=`$ $`X_pZ_R(X_p)+X_eZ_R(X_e),`$ (A6) $`f_2`$ $`=`$ $`[1+X_pZ_R(X_p)]+T_r[1+X_eZ_R(X_e)],`$ (A7) $`f_3`$ $`=`$ $`X_pe^{X_p^2}X_ee^{X_e^2},`$ (A8) $`f_4`$ $`=`$ $`X_pe^{X_p^2}+T_rX_ee^{X_e^2},`$ (A9) $`f_5`$ $`=`$ $`1+X_pZ_R(X_p),`$ (A10) $`f_6`$ $`=`$ $`X_pe^{X_p^2},`$ (A11) $`f_7`$ $`=`$ $`X_pZ_R(X_p),`$ (A12) $`g_1`$ $`=`$ $`[X_p^2+X_p^3Z_R(X_p)X_pZ_R(X_p)]+(1/T_r)[X_e^2+X_e^3Z_R(X_e)X_eZ_R(X_e)],`$ (A13) $`g_2`$ $`=`$ $`(X_p^3X_p)e^{X_p^2}+(1/T_r)(X_e^3X_e)e^{X_e^2},`$ (A14) $`g_3`$ $`=`$ $`(2X_p^3X_p)e^{X_p^2}(2X_e^3X_e)e^{X_e^2},`$ (A15) $`g_4`$ $`=`$ $`[2X_p^2+2X_p^3Z_R(X_p)X_pZ_R(X_p)][2X_e^2+2X_e^3Z_R(X_e)X_eZ_R(X_e)].`$ (A16) Here $`Z_R(X)`$ is the real part of the plasma dispersion function for argument $`X`$: $$Z(X)=\frac{1}{\sqrt{\pi }}_{\mathrm{}}^{\mathrm{}}\frac{\mathrm{exp}(\xi ^2)}{\xi X}d\xi =2ie^{X^2}_{\mathrm{}}^{iX}e^{\xi ^2}d\xi .$$ (A17) The imaginary part of $`Z(X)`$ appears explicitly in Eq. (1) as the Hilbert integral operator. Physically, $`X=X_R+iX_I`$ is the ratio of wave phase velocity to thermal velocity. For Alfvén waves, $`X_p`$ and $`X_e`$ become: $$X_p=\left(1+\frac{T_e}{T_p}\right)^{1/2}\frac{1}{\sqrt{\beta }},X_e=\left(\frac{T_p}{T_e}+1\right)^{1/2}\sqrt{\frac{m_e}{m_p}}\frac{1}{\sqrt{\beta }}$$ (A18) At last, $`T_r=T_p/T_e`$ is the proton-to-electron temperature ratio and $`m_e/m_p1/1836`$ is the electron-to-proton mass ratio. Figs. Particle Heating by Nonlinear Alfvénic Turbulence in ADAFs,Particle Heating by Nonlinear Alfvénic Turbulence in ADAFs Fig. Particle Heating by Nonlinear Alfvénic Turbulence in ADAFs
warning/0005/cond-mat0005260.html	ar5iv	text	# Entropic Lattice Boltzmann Methods ## 1 Introduction Though lattice models have been used to study equilibrium systems since the 1920’s, their application to hydrodynamic systems is a much more recent phenomenon. Lattice-gas models of hydrodynamics were first advanced in the late 1980’s, and the related lattice Boltzmann method was developed in the early 1990’s. This paper describes a particularly interesting and useful category of lattice Boltzmann models, which we call Entropic Lattice Boltzmann models, and provides a number of mathematical tools for constructing and studying them. This introductory section describes and contrasts lattice methods for hydrodynamics, traces their development, and endeavors to place the current work in historical context. ### 1.1 Lattice Gases The first isotropic lattice-gas models of hydrodynamics were introduced in the late 1980’s . Such models consist of discrete particles moving and colliding on a lattice, conserving mass and momentum as they do. If the lattice has sufficient symmetry, it can be shown that the density and hydrodynamic velocity of the particles satisfy the Navier-Stokes equations in the appropriate scaling limit. This method exploits an interesting hypothesis of kinetic theory: The Navier-Stokes equations are the dynamic renormalization group fixed point for the hydrodynamic behavior of a system of particles whose collisions conserve mass and momentum. We refer to this assertion as a hypothesis, because it is notoriously difficult to demonstrate in any rigorous fashion. Nevertheless, it is very compelling because it explains why a wide range of real fluids with dramatically different molecular properties – such as air, water, honey and oil – can all be described by the Navier-Stokes equations. A lattice gas can then be understood as a “minimalist” construction of such a set of interacting particles. Viewed in this way, it is perhaps less surprising that they satisfy the Navier-Stokes equations in the scaling limit. For a time, lattice-gas models were actively investigated as an alternative methodology for computational fluid dynamics (CFD) . Unlike all prior CFD methodologies, they do not begin with the Navier-Stokes equations; rather, these equations are an emergent property of the particulate model. Consequently, the use of such models to simulate the Navier-Stokes equations tends to parallel theoretical and experimental studies of natural fluids. One derives a Boltzmann equation for the lattice-gas particles, and applies the Chapman-Enskog analysis to determine the form of the hydrodynamic equations and the transport coefficients. In fact, numerical experimentalists often circumvent this theory by simply measuring lattice-gas transport coefficients in a controlled setting in advance of using them to study a particular flow problem, just as do laboratory experimentalists. One often overlooked advantage of lattice-gas models is their unconditional stability. By insisting that lattice-gas collisions obey a detailed-balance condition <sup>1</sup><sup>1</sup>1This is certainly possible for an ideal single-phase viscous fluid. For complex fluids with finite-range interaction potentials, it is an outstanding problem., we are ensured of the validity of Boltzmann’s $`H`$ theorem, the fluctuation-dissipation theorem, Onsager reciprocity, and a host of other critically important properties with macroscopic consequences. By contrast, when the microscopic origins of the Navier-Stokes equations are cavalierly ignored, and they are “chopped up” into finite-difference schemes, these important properties can be lost. The discretized evolution equations need no longer possess a Lyapunov functional, and the notion of underlying fluctuations may lose meaning altogether. As the first practioneers of finite-difference methods found in the 1940’s and 1950’s, the result can be the development and growth of high-wavenumber numerical instabilities, and indeed these have plagued essentially all CFD methodologies in all of the decades since. Such instabilities are entirely unphysical because they indicate the absence of a Lyapunov functional, analogous to Boltzmann’s $`H`$. Numerical analyists have responded to this problem with methods for mitigating these anomalies – such as upwind differencing – but from a physical point of view it would have been much better if the original discretization process had retained more of the underlying kinetics, so that these problems had not occurred in the first place. Lattice gases represent an important first step in this direction. As was shown shortly after their first applications to hydrodynamics, lattice-gas models for single-phase ideal fluids can be constructed with an $`H`$ theorem that rigorously precludes any kind of numerical instability. More glibly stated, lattice gases avoid numerical instabilities in precisely the same way that Nature herself does so. ### 1.2 Lattice Boltzmann Methods In spite of these appealing features, the presence of intrinsic kinetic fluctuations makes lattice-gas models less than ideal as a CFD methodology. Accurate values for the velocity field at selected locations, or even for bulk coefficients such as drag and lift, have an intrinsic statistical error that can be reduced only by extensive averaging. On the other hand, the presence of such fluctuations in a simple hydrodynamic model make lattice gases an ideal tool for those studying the statistical physics of fluids, molecular hydrodynamics, and complex-fluid hydrodynamics; not surprisingly, these remain the method’s principal application areas . Consequently, many CFD researchers who appreciated the emergent nature of lattice-gas hydrodynamics but wanted to eliminate (or at least control) the level of fluctuations, turned their attention to the direct simulation of the Boltzmann equation for lattice gases. Such simulations are called lattice Boltzmann models . These models evolve the real-valued single-particle distribution function, rather than the discrete particles themselves, and in this manner eliminate the kinetic fluctuations. Early attempts along these lines restricted attention to Boltzmann equations for actual lattice gases. It was soon realized, however, that these quickly become unwieldy as the number of possible particle velocities increases. In order for lattice Boltzmann models to become practical tools, it was necessary to develop simplified collision operators that did not necessarily correspond to an underlying lattice-gas model. The most successful collision operators of this type are those of the form due to Bhatnager, Gross and Krook , and these have given rise to the so-called lattice BGK models . ### 1.3 Lattice BGK Models Lattice BGK collision operators allow the user to specify the form of the equilibrium distribution function to which the fluid should relax. For lattice gases obeying a detailed balance <sup>2</sup><sup>2</sup>2Actually, a weaker condition called semi-detailed balance suffices for this purpose. condition, this is known to be a Fermi-Dirac distribution. Having abandoned lattice-gas collision operators, however, it seemed unnecessary to continue to use lattice-gas equilibria, and practitioners exploited the freedom of choosing lattice BGK equilibria to achieve certain desiderata, such as Galilean invariance, and the correct form of the compressible Navier-Stokes equations. The alert reader will have noticed, however, that the evolution to lattice-BGK methods has jettisoned the last vestiges of kinetic-level physics that might have been left over from the original lattice-gas models. The move to a Boltzmann description of the lattice gas eliminated kinetic fluctuations, but at least it retained an $`H`$-theorem; more precisely, its global equilibrium still extremized a Lyapunov functional of the dynamics. The move to lattice-BGK operators and the arbitrary legislation of the equilibrium distribution function, however, completely abandoned even the concept of detailed balance and the $`H`$-theorem. Without a Lyapunov functional, lattice-BGK methods became susceptible to a wide variety of numerical instabilities which are ill-understood and remain the principal obstacle to the wider application of the technique to the present day. ### 1.4 Motivation for Entropic Lattice Boltzmann Methods In this paper, we argue that the most natural way to eliminate numerical instabilities from lattice Boltzmann models is to return to the method its kinetic underpinnings, including the notion of a Lyapunov functional. We present a general program for the construction of “entropically stabilized” lattice Boltzmann models, and illustrate its application to several example problems. The crux of the idea is to encourage the model builder to specify an appropriate Lyapunov functional for the model, rather than try to blindly dictate an equilibrium state. Of course, specifying a Lyapunov functional determines the equilibrium distribution, but it also governs the approach to this equilibrium. It can therefore be used to control the stability properties of the model. It should be emphasized that the presence of a Lyapunov functional guarantees the nonlinear stability of the model, which is a much stronger condition than linear stability. As we shall show, collision operators that are constructed in this way are similar in form to the lattice BGK collision operators, except that their relaxation parameter may depend on the current state. As a result, the transport coefficients may have a certain minimum value in models of this type but, as we shall see, this value may actually be zero. In fact, the ultimate limitation to the application of these algorithms for very small transport coefficients will come from considerations of accuracy, rather than stability. ### 1.5 Structure of Paper The plan of this paper is as follows: In Section 2 we present the general program of construction of entropic lattice Boltzmann models. This includes the introduction of a conservation representation of the lattice Boltzmann distribution function, by means of which the collision process is most naturally described. This representation also provides an interesting geometric interpretation, as it allows us to describe a collision process as a mathematical map from a certain polytope into (and perhaps onto) itself. We then show a way of constructing Lyapunov functionals on this polytope, contrast these with Boltzmann’s $`H`$, and use them to construct collision operators for absolutely stable models. Section 3 carries out this program for a very simple lattice Boltzmann model of diffusion in one dimension. The model is useful from a pedagogical point of view, since an entropically stabilized collision operator can be written for it in closed form. Moreover, it is possible to determine the transport coefficient of this model analytically. The result is a fully explicit, perfectly conservative, absolutely stable method of integrating the diffusion equation. We present numerical simulations for various values of the diffusivity in order to probe the limitations of the technique. In the course of presenting this model, we show that a linear algebraic procedure, called the Fourier-Motzkin algorithm, is very useful for constructing and visualizing the conservation representation. Section 4 applies the method to a simple five-velocity model of fluid dynamics in one dimension, first considered by Renda et al. in 1997 . Here the geometric picture involves a four dimensional polytope, and the Fourier-Motzkin algorithm is shown to be very useful in describing it. Indeed, Appendix A gives the reader a “tour” of this polytope. Finally, Section 5 applies the method to the two-dimensional FHP model and the three-dimensional FCHC model of fluid dynamics. For these examples, we are able to provide appropriate conservation representations. Though the latter example has too many degrees of freedom to allow geometric visualization of the collision dynamics, it is nevertheless possible to develop and present a Lyapunov functional and an entropic model. We use these more complicated models to demonstrate that entropic lattice Boltzmann models are not computationally prohibitive, even when the number of velocities involved is large. We describe the computational aspects of such models in some detail, and end with a discussion of Galilean invariance. ## 2 Entropic Lattice Boltzmann Models In this section, we present the general program of construction of entropic lattice Boltzmann models. We note that the lattice Boltzmann distribution function admits a variety of representations, and we show how to transform between these. In particular, we introduce a conservation representation, by means of which the collision process has a very natural description. We also examine the geometric aspects of this representation, and use these insights to show how to construct lattice Boltzmann collision operators possessing a Lyapunov functional. ### 2.1 Representation Lattice Boltzmann models are constructed on a $`D`$ dimensional regular lattice $``$, and evolve in discrete time intervals $`\mathrm{\Delta }t`$. We imagine that a population of particles exists at each site $`𝐱`$, and that these particles can have one of a finite number $`n`$ of velocities, $`𝐜_j/\mathrm{\Delta }t`$ where $`j\{1,\mathrm{},n\}`$. The displacement vectors $`𝐜_j`$ are required to be combinations of lattice vectors with integer coefficients, so that if $`𝐱`$, then $`𝐱+𝐜_j`$. We note that there is nothing precluding some of the $`𝐜_j`$ from being the same, or from being zero; indeed, this latter choice is often made to incorporate so-called “rest particles” into the model. Mathematically, the state of the system of particles is represented by $`n`$ real numbers at each site $`𝐱`$ and time step $`t`$. These are denoted by $`N_j(𝐱,t)\mathrm{}`$, and can be thought of as the expected number of particles with mass $`m_j`$ and velocity $`𝐜_j/\mathrm{\Delta }t`$ at site $`𝐱`$ and time $`t`$. As such, we expect them to be nonnegative, $`N_j(𝐱,t)0`$. Taken together, these quantities constitute the single-particle distribution function of the system. This is all the information that is retained in a Boltzmann description of interacting particles. In what follows, we shall think of the set of single-particle distribution functions as a vector space. That is, we shall regard the $`n`$ quantities $`N_j(𝐱,t)`$, for $`i\{1,\mathrm{},n\}`$, as the components of an $`n`$-vector or “ket,” $`\|𝐍(𝐱,t)\mathrm{}^n`$. If no ambiguity is likely to result, we shall often omit the explicit dependence on $`t`$ and write simply $`\|𝐍(𝐱)`$. If discussion is restricted to a single lattice site, we may further abbreviate this as $`\|𝐍`$. In one discrete time step $`\mathrm{\Delta }t`$, the state of the system is modified in a manner that is intended to model collisions between the particles at each site, followed by propagation of the collided particles to new sites. The collisions are modelled by modifications to the distribution function $$\|𝐍(𝐱)\|𝐍^{}.$$ (1) Here and henceforth, we use a prime to denote the postcollision state. It is usually the case that $`\|𝐍^{}(𝐱)`$ depends only on $`\|𝐍(𝐱)`$, though in more sophisticated lattice Boltzmann models it may also depend on the $`\|𝐍`$’s in a local neighborhood of sites about $`𝐱`$. Following convention, we define the collision operator $`C_j(\|𝐍(𝐱))`$ as the difference between the old and new values of the single-particle distribution, $$C_j(\|𝐍(𝐱))=N_j^{}(𝐱)N_j(𝐱).$$ (2) As this is the difference of two kets, it can also be thought of as a ket, $$\|𝐂=\|𝐍^{}\|𝐍.$$ (3) ### 2.2 Conservation Laws The collision process specified in Eq. (1) is usually required to conserve some number of locally defined quantities. Usually, these quantities are additive over the lattice sites and directions. For example, we may require that there be a conserved mass $$M=\underset{𝐱}{}\rho (𝐱,t),$$ (4) where we have defined the mass density $$\rho (𝐱,t)=\underset{j=0}{\overset{n}{}}m_jN_j(𝐱,t)=\underset{j=0}{\overset{n}{}}m_jN_j^{}(𝐱,t).$$ (5) To streamline the notation for this, we can define a covector or “bra” by the prescription $$𝝆\|_j=m_j,$$ (6) and write simply $$\rho =𝝆𝐍=𝝆𝐍^{}.$$ (7) Thus, to the extent that we think of the single-particle distribution function as a vector, to each conserved quantity there corresponds a covector such that the value of the conserved quantity is the contraction of the two. Because collisions are required to preserve this value, the contraction of the covector with the collision operator must vanish, $$𝝆𝐂=0,$$ (8) as can be seen from Eqs. (3) and (7). Lattice Boltzmann models may also conserve momentum, $$𝝅(𝐱,t)=\underset{j=0}{\overset{n}{}}m_j\frac{𝐜_j}{\mathrm{\Delta }t}N_j(𝐱,t),=\underset{j=0}{\overset{n}{}}m_j\frac{𝐜_j}{\mathrm{\Delta }t}N_j^{}(𝐱,t).$$ (9) This is still linear in the $`N_j`$, so that we can define $$𝝅\|_j=m_j\frac{𝐜_j}{\mathrm{\Delta }t},$$ (10) and write $$𝝅(𝐱,t)=𝝅𝐍=𝝅𝐍^{}.$$ (11) Here we must be a bit careful to define the meaning of the quantities involved: Note that $`𝝅\|_j`$ is a covector in its $`j`$ index, but a vector in its spatial index <sup>3</sup><sup>3</sup>3The pedant will note that momentum is more properly thought of as a covector in its spatial index, but we do not bother to distinguish between spatial vectors and covectors in this paper.. Thus, the contraction with $`\|𝐍_j`$ is over the $`j`$ index, and results in the vector $`𝝅`$. Once again, this covector annihilates the collision operator, $$𝝅𝐂=\mathrm{𝟎},$$ (12) where the right-hand side is a null vector. The subsequent propagation process is described mathematically by the prescription $$N_j(𝐱+𝐜_j,t+\mathrm{\Delta }t)=N_j^{}(𝐱,t).$$ (13) Note that this operation must be carried out simultaneously over the entire lattice. This may alter the values of the conserved quantities at each site, but because it is nothing more than a permutation of the values of the distribution function about the lattice, it clearly leaves unaltered the global values of the conserved quantities, $$\begin{array}{ccccc}_𝐱\rho (𝐱,t)& & \text{and}& & _𝐱𝝅(𝐱,t).\end{array}$$ (14) Combining Eqs. (2) and (13), we find the general dynamical equation for a lattice Boltzmann model, $$N_j(𝐱+𝐜_j,t+\mathrm{\Delta }t)N_j(𝐱)=C_j(\|𝐍(𝐱)).$$ (15) For compressible fluids, it is also necessary to pay attention to conservation of energy. Conservation of kinetic energy can be expressed using the bra, $$𝜺\|_j=\frac{m_j}{2}\left\|𝐜_j\right\|^2.$$ (16) Indeed, the right-hand side could be generalized without difficulty to anything that depends only on $`j`$. The problem of including a potential energy function between particles at different sites, however, is much more difficult. If we suppose that there is a potential $`V(\|𝐱_j𝐱_k\|)`$ between particles at $`𝐱_j,𝐱_k`$, then two problems arise: First, the total potential energy will depend on the pair distribution function – that of finding two particles a certain distance apart. This is outside the scope of Boltzmann methods, which retain only the single-particle distribution. We can avoid this problem by making the mean-field approximation, in which the probability of having one particle at $`𝐱_j`$ and another at $`𝐱_k`$ is simply the product of the two single-particle probabilities, but then the potential energy is quadratic, rather than linear, in the $`N_j`$’s. Second, and perhaps more distressingly, the potential energy is not preserved by the propagation step . These considerations make it very difficult to add potential interaction to lattice Boltzmann models. If we are willing to give up the idea of an exactly conserved energy, and instead consider isothermal systems, then methods of skirting these difficulties have been known and actively investigated for some time now . More recently, a method of incorporating interaction that maintains exact (kinetic plus potential) energy conservation has also been proposed . In any case, this paper will be restricted to systems with kinetic energy only. The possibility of extending our methods to models with nontrivial interaction potentials is discussed briefly in the Conclusions section. ### 2.3 Geometric Viewpoint The requirements of maintaining the conserved quantities and the nonnegativity of the distribution function place very important constraints on the collision process. Much of this paper will be devoted to describing – algebraically and geometrically – the most general set of collisional alterations of the $`N_j(𝐱)`$ that meet these requirements. Suppose that there are $`n_c<n`$ independent conserved quantities. As described above, these define $`n_c`$ linearly independent covectors or bras, whose contraction with the collision operator vanishes. We shall sometimes adopt special names for these; for example, that corresponding to mass conservation shall be called $`𝝆\|`$, that for momentum conservation shall be called $`𝝅\|`$, etc. Generically, however, let us refer to them as $`𝝀_\sigma \|`$, where $`\sigma =1,\mathrm{},n_c`$. These covectors are not necessarily orthogonal in any sense. Note, for example, that the covectors for mass and kinetic energy in Eqs. (6) and (16) are not orthogonal with respect to the Euclidean metric. In fact, there is no reason to insist on any kind of metric in this covector space. The $`n_c`$ covectors are not uniquely defined, since a linear combination of two conserved quantities is also conserved. The $`n_c`$ dimensional subspace defined by the covectors, on the other hand, is uniquely defined, and we shall call it the hydrodynamic subspace <sup>4</sup><sup>4</sup>4Note that we are using the term “hydrodynamic” here to describe degrees of freedom that will result in macroscopic equations of motion, whether or not they are those of a fluid. If mass is the only conserved quantity, a diffusion equation generically results; if momentum and/or energy are conserved as well, a set of fluid equations generically results. We use the term “hydrodynamic” in either case. This terminology is standard in kinetic theory and statistical physics, but often seems strange to hydrodynamicists.. It is always possible to construct $`nn_c`$ more covectors that are linearly independent of each other, and of the $`n_c`$ corresponding to the conserved quantities. For example, this may be done by the Gram-Schmidt procedure. Let us generically call these $`𝜶_\eta \|`$, where $`\eta =1,\mathrm{},nn_c`$; if there are only a few of these and we want to avoid excessive subscripting, we shall use the successive Greek letters $`𝜶\|,𝜷\|,𝜸\|,\mathrm{}`$ for this purpose. Once again, these are not uniquely defined or orthogonal in any sense, but they do span an $`nn_c`$ dimensional subspace which we shall call the kinetic subspace. The union of all $`n`$ of our covectors, namely the $`𝝀_\sigma \|`$’s and the $`𝜶_\eta \|`$’s, thus constitute a basis for the full $`n`$ dimensional covector space. The ket $`\|𝐍`$ whose components are the single-particle distribution function is thus seen to be but one representation of the dynamical variable. As has been recognized for some time now , we are free to make a change of basis in the lattice Boltzmann equation. In fact, one other basis suggests itself naturally. Given a basis of covectors or bras, it is always possible to construct a dual basis of vectors or kets. Indeed, our decomposition of the covector space into hydrodynamic and kinetic subspaces is naturally mirrored by a corresponding decomposition of the vector space. Thus we construct $`\|𝝀_\sigma `$ where $`\alpha =1,\mathrm{},n_c`$ and $`\|𝜶_\eta `$ where $`\beta =1,\mathrm{},nn_c`$, such that $$\begin{array}{cc}𝝀_\sigma 𝝀_\tau =\delta _{\sigma \tau }\hfill & 𝝀_\sigma 𝜶_\eta =0\hfill \\ 𝜶_\eta 𝝀_\sigma =0\hfill & 𝜶_\eta 𝜶_\theta =\delta _{\eta \theta }.\hfill \end{array}$$ (17) Thus, we can expend $`\|𝐍`$ in terms of these hydrodynamic and kinetic basis kets. $$\|𝐍=\underset{\sigma =1}{\overset{n_c}{}}\lambda _\sigma \|𝝀_\sigma +\underset{\eta =1}{\overset{nn_c}{}}\alpha _\eta \|𝜶_\eta .$$ (18) We shall call the set of coefficients $`\lambda _\sigma `$ and $`\alpha _\eta `$ the conservation representation. The explicit construction of this representation is best illustrated by example, and we give several of these in the following sections. The advantage of the conservation representation for describing collisions is clear: When $`\|𝐍`$ is expanded in terms of hydrodynamic and kinetic kets, collisions may change only the coefficients of the latter. From Eqs. (17) and (18), we see that the coefficients of the hydrodynamic kets, namely the $`\lambda _\sigma `$’s, are the conserved quantities themselves. These are precisely what must remain unchanged by a collision. So, in the conservation representation, the collision process is simply an alteration of the $`nn_c`$ coefficients of the kinetic kets, namely the $`\alpha _\eta `$’s. Thus, this representation effectively reduces the dimensionality of the space needed to describe the collision process to $`nn_c`$. Again, we shall illustrate the construction of and transformations between these alternative representations of the single-particle distribution function for several examples of entropic lattice Boltzmann models. ### 2.4 Nonnegativity While the conservation representation of $`\|𝐍`$ is most natural for describing collisions, the original representation is more natural for describing the constraint of nonnegativity of the distribution function components. In the original representation, we had a restriction to $`N_j0`$ for $`j=1,\mathrm{},n`$. In the conservation representation, these $`n`$ inequalities transform to a corresponding set of linear inequalities on the hydrodynamic and kinetic parameters, $`\lambda _\sigma `$ and $`\alpha _\eta `$. As is well known, such a set of linear inequalities define a convex polytope in the parameter space. This construction, in the full $`n`$ dimensional space of hydrodynamic and kinetic parameters, shall be called the master polytope. Specification of the hydrodynamic parameters then defines the cross section of the master polytope that bounds the kinetic parameters, $`\alpha _\eta `$; we shall call these cross sections the kinetic polytopes, and it is clear that they must also be convex. We shall construct the master and kinetic polytopes for simple entropic lattice Boltzmann models later in this paper. We shall see that they become very difficult to visualize when the number of particle velocities becomes large. This difficulty raises the question of whether or not there is a general method to describe the shape of polytopes defined by linear equalities in this way. It turns out that such a method is well known in linear programming and optimization theory, and is called Fourier-Motzkin elimination . It is constructive in nature, and works for any set of inequalities in any number of unknowns. If the inequalities cannot be simultaneously satisfied, the method will indicate that. Otherwise, it will yield an ordered sequence of inequalities for each variable, the bounds of which depend only on the previously bounded variables in the sequence. If it were desired to perform a multiple integral over the polytopic region, for example, the Fourier-Motzkin method would provide a systematic procedure for setting up the limits of integration. ### 2.5 Collisions Once we are able to construct and visualize these polytopes, it is straightforward to describe the constraints that conservation imposes on the collision process: Because the propagation process is nothing more than a permutation of the values of the $`N_j(𝐱)`$ on the lattice, it is clear that it will maintain nonnegativity. That is, if the $`N_j(𝐱)`$ were positive prior to the propagation, they will be so afterward. This means that the post-propagation state of each site will lie within the allowed master polytope. Given such a post-propagation state $`\|𝐍`$, we transform to the conservation represenation. The coefficients $`\lambda _\sigma `$ of the hydrodynamic kets are the conserved quantities and must remain unchanged by the collision. Geometrically, these determine a cross section of the master polytope within which the state must remain. This cross section is the kinetic polytope of the pre-collision state. The essential point is that the post-collision state must also lie within this kinetic polytope to preserve nonnegativity. To the extent that the postcollision state is determined by the precollision state, this means that: ###### Collision Property 1 The collision process is a map from the kinetic polytope into itself. We note that this requirement is not without some controversy. It may be argued that lattice Boltzmann algorithms are ficticious kinetic models from which realistic hydrodynamics are emergent. Since the details of the kinetics are ficticious anyway, why not also dispense with the requirement that the single-particle distribution function be positive? As long as the conserved densities of positive-definite quantities, such as mass and kinetic energy, are positive, why should we care if the underlying lattice Boltzmann distribution function is likewise? There are two reasons: The first is that, even if the system is initialized with nonnegative physical densities, the propagation step may give rise to negative physical densities if negative values of the distribution are allowed. To see this, imagine a postcollision state in which all of the neighbors of site $`𝐱`$ have a single negative distribution function component in the direction heading toward $`𝐱`$. Even if all these neighboring sites had positive mass density, site $`𝐱`$ will have a negative mass density after one propagation step. The second reason for demanding nonnegativity of the distribution function is, to some extent, a matter of taste. We like to think that the kinetic underpinnings of the lattice Boltzmann algorithm are more than just a mathematical trick to yield a desired set of hydrodynamic equations. Though there may be no physical system with such a collision operator (not to mention dynamics constrained to a lattice), we feel that the more properties of real kinetics that can be maintained, the more useful the algorithm is likely to be. This is particularly important for complex fluids, for which the form of the hydrodynamic equations is often unknown, and for which we must often appeal to some kinetic level of description. Nevertheless, we shall revisit this question in Subsection 5.5. Collision Property 1 is simple in statement and motivation, but in fact it weeds out many putative lattice Boltzmann collision operators, including those most commonly used in computational fluid dynamics research today – namely, the overrelaxed lattice BGK operators. For sufficiently large overrelaxation parameter (collision frequency), such operators are well known to give rise to negative values of the $`N_j(𝐱)`$. This is usually symptomatic of the onset of a numerical instability in the lattice Boltzmann algorithm. We shall discuss such instabilities in more detail later in this paper. For the present, we emphasize that we are not claiming that Collision Property 1 will eliminate such instabilities. The property does mandate a lower bound of zero on the $`N_j(𝐱)`$, and hence it restricts the manner in which such instabilities might grow and saturate. Nevertheless, it is generally still possible for collision operators that obey Collision Property 1 to exhibit numerical instability if they lack a Lyapunov functional. ### 2.6 Reversibility Collision Property 1 is the minimum requirement that we impose on our lattice Boltzmann collision operators. It is possible to define more stringent requirements for them, and we shall continue to do exactly that to satisfy various desiderata. For example, we might want to demand that our lattice Boltzmann algorithm be reversible. A reversible algorithm could be run backwards in time from any final condition by alternately applying the inverse propagation operator, $$N_j(𝐱𝐜_j)N_j(𝐱),$$ (19) followed by an inverse collision operator to recover an initial condition at an earlier time. For such an inverse collision operator to exist, the map from the allowed polytope to itself must be one-to-one. Since we have already demanded that the map be into, this means that it must also be onto, hence: ###### Collision Property 2 A collision process is reversible if it is a map from the allowed polytope onto itself. ### 2.7 Imposition of a Lyapunov Functional The criteria that we have set out thusfar assure that the conservation laws – and hence the First Law of Thermodynamics – will be obeyed. To ensure stability and thermodynamic consistency, however, it is necessary to also incorporate the requirements of the Second Law. We suppose that our system has a function $`H`$ of the state variables $`N_j`$ that is additive over sites $`𝐱`$, $$H=\underset{𝐱}{}h\left(𝐍(𝐱)\right),$$ (20) and additive over directions $`j`$, $$h\left(𝐍(𝐱)\right)=\underset{j=1}{\overset{n}{}}\theta _j\left(N_j\left(𝐱\right)\right),$$ (21) where the functions $`\theta _j`$ are defined on the nonnegative real numbers. One might wonder if Eq. (21) could be generalized, but in fact it has been shown to be a necessary condition for an $`H`$-theorem . It is clear from this construction that $`H`$ is preserved under the propagation operation of Eq. (13). If we require that our collisions never decrease the contribution to $`H`$ from each site – that is, that $`h\left(𝐍(𝐱)\right)`$ can only be increased by a collision – then $`H`$ is a Lyapunov functional for our system, and the existence and stability of an equilibrium state is guaranteed. We note that $`h`$ is a function of the $`N_j`$ for all $`N_j0`$, and can therefore be thought of as a scalar function on the master polytope. Specification of the coefficients $`\lambda _\sigma `$ of the hydrodynamic kets then determine the cross section of allowed collision outcomes, or the kinetic polytope, parametrized by the coefficients of the kinetic kets $`\alpha _\eta `$. Thus, for a given incoming state, $`h`$ can be thought of as a scalar function on the kinetic polytope. We denote this by $`h(\alpha _\eta )`$, and we demand that the collision process increase this function: ###### Collision Property 3 To ensure the existence and stability of an equilibrium state, a collision at site $`𝐱`$ must not decrease the restriction of the function $`h`$ to the kinetic polytope. From a more geometric point of view, note that the (hyper)surfaces of constant $`h(\alpha _\eta )`$ provide a codimension-one foliation of the kinetic polytope. These (hyper)surfaces degenerate to a point where $`h`$ reaches a maximum. An incoming state lies on such a (hyper)surface. A legitimate collision is required to map such an incoming state to an outgoing one that lies inside, or at least on, this (hyper)surface. Clearly, the point with maximal $`h`$ must be mapped to itself. Indeed, we can subsume Collision Property 1 into Collision Property 3 by constructing the function $`h`$ so that it takes a minimum value on the boundary of the polytope, and increases to a single maximum somewhere inside. The master polytope is defined as the region for which all of the $`N_j(𝐱)`$’s are nonnegative. The boundary of the master polytope is therefore the place where at least one of the $`N_j(𝐱)`$’s vanishes. It follows that $`_j^nN_j(𝐱)`$ is constant – in fact, it is zero – on the polytope boundary. More generally, if $`\zeta _j(x)`$ for $`j=1,\mathrm{},n`$ are such that $$\zeta _j(0)=0,$$ (22) then $`_j^n\zeta _j\left(N_j(𝐱)\right)`$ also goes to zero on the polytope boundary. To ensure that there is only one maximum inside the convex (master or kinetic) polytope, we also require that the $`\zeta _j`$ be strictly increasing, $$\zeta _j^{}(x)>0,$$ (23) for nonnegative $`x`$. Indeed, from Eqs. (22) and (23) it follows that $$\zeta _j(x)0,$$ (24) for nonnegative $`x`$. Thus, the simplest choice would be to make $`h\left(𝐍(𝐱)\right)`$ a function of $`_j^n\zeta _j\left(N_j(𝐱)\right)`$. To make this consistent with Eq. (21), however, we see that the functions $`\theta _j`$ should be chosen to be the logarithms of the functions $`\zeta _j`$, so $$h\left(𝐍(𝐱)\right)=\underset{j=1}{\overset{n}{}}\mathrm{ln}\left[\zeta _j\left(N_j(𝐱)\right)\right].$$ (25) This goes to negative infinity on the boundaries of the master (and hence kinetic) polytopes, and it has a unique maximum in the interior. ###### Collision Property 4 A valid functional form for $`h`$ is given by Eq. (25), where the functions $`\zeta _j`$ obey Eq. (22), Eq. (23), and Eq. (24). This may well be the most controversial of the four collision properties that we have presented. Indeed, it is not necessary to assure that $`H`$ have a single maximum within the polytope. We shall discuss an $`H`$ function that violates Collision Property 4 in Subsection 5.5. The equilibrium distribution is the point within the kinetic polytope where $`h`$ has its maximum value. If we denote the kinetic parameters by $`\alpha _j`$, then we can find this point by demanding that the gradient of $`h`$ vanish, $$0=\frac{h}{\alpha _\eta }=\underset{j=1}{\overset{n}{}}\frac{\zeta _j^{}\left(N_j\right)}{\zeta _j\left(N_j\right)}\frac{N_j}{\alpha _\eta }=\underset{j=1}{\overset{n}{}}\frac{\zeta _j^{}\left(N_j\right)}{\zeta _j\left(N_j\right)}\|𝜶_\eta _j.$$ (26) If we denote the equilibrium distribution function by $`N_j^{\text{eq}}`$, it follows that the covector with components $`\zeta _j^{}\left(N_j^{\text{eq}}\right)/\zeta _j\left(N_j^{\text{eq}}\right)`$ lies in the hydrodynamic subspace, $$\frac{\zeta _j^{}\left(N_j^{\text{eq}}\right)}{\zeta _j\left(N_j^{\text{eq}}\right)}=\underset{\sigma =1}{\overset{n_c}{}}Q_\sigma 𝝀_\sigma \|_j.$$ (27) The $`n_c`$ coefficients $`Q_\sigma `$ must be chosen so that the values of the conserved quantities are as desired, $$𝝀_\tau 𝐍^{\text{eq}}=\lambda _\tau $$ (28) for $`\tau =1,\mathrm{},n_c`$. Eqs. (27) and (28) can be regarded as $`n+n_c`$ equations for the $`n+n_c`$ unknowns, $`N_j^{\text{eq}}`$ (for $`j=1,\mathrm{},n`$) and $`Q_\sigma `$ (for $`\sigma =1,\mathrm{},n_c`$). In general, these equations are nonlinear and it is not always possible to write the $`N_j^{\text{eq}}`$ in closed form. The equilibrium distribution is something over which the model builder would like to retain as much control as possible, since it is often used to tailor the form of the resultant hydrodynamic equations. For example, a judicious choice of the equilibrium distribution function is required to obtain Galilean invariant hydrodynamic equations for lattice Boltzmann models of fluids. Customarily, lattice Boltzmann model builders have simply dictated the form of the equilibrium distribution function. To the extent that this is necessary, however, it is more in keeping with the philosophy espoused in this paper to (try to) dictate the form of the Lyapunov functional, rather than that of the equilibrium distribution. The challenge is to find a set of functions $`\zeta _j`$, subject to the requirements of Eqs. (22), (23) and (24), such that Eq. (27) yields the desired equilibrium distribution. We shall return to this problem in Subsection 5.5. While this formulation of an entropic principle for lattice Boltzmann models seems reasonable, note that it is rather different from the usual one encountered in kinetic theory. The usual choice there would be that of Boltzmann, $`h=_{j=1}^nN_j\mathrm{ln}N_j`$. A moment’s examination of Eq. (25) indicates that this would correspond to the choice $`\zeta _j(x)=x^x`$, but this function decreases for $`x(0,1/e)`$ and hence violates Eq. (23). Thus, while it is possible to construct lattice Boltzmann models with Maxwell-Boltzmann equilibria , the Lyapunov functionals of the models described in this paper need to be rather different from those commonly used in kinetic theory. With this caveat in mind, we shall henceforth abuse terminology by taking “$`H`$ function” and “Lyapunov functional” to be synonymous. ### 2.8 Collision Operator Finally, we turn our attention to the construction of a collision operator that is in keeping with the collision properties described above. Obviously, there are many ways to describe maps from the kinetic polytopes into (and perhaps onto) themselves, that do not decrease an $`H`$ function. For the purposes of this paper, however, we restrict our attention to the BGK form of collision operator in which the outgoing state is a linear combination of the incoming state and the equilbrium, $$\|𝐂=\frac{1}{\tau }\left(\|𝐍^{\text{eq}}\|𝐍\right),$$ (29) where $`\tau `$ is called the relaxation time. Usually $`\tau `$ is taken to be constant, but more generally it may depend on the conserved quantities, and most generally on all of the $`N_j`$. If we combine Eqs. (29) and (3), we get the BGK equation, $$\|𝐍^{}=\|𝐍+\frac{1}{\tau }\left(\|𝐍^{\text{eq}}\|𝐍\right).$$ (30) From a geometric point of view, this equation tells us to draw a line in the kinetic polytope from the position of the incoming state $`\|𝐍`$ through the equilibrium state $`\|𝐍^{\text{eq}}`$. The final state is a weighted combination of these two states and hence lies on this line. The incoming state is weighted by $`11/\tau `$, and the equilibrium state is weighted by $`1/\tau `$. Thus, for $`\tau 1`$ the outgoing state lies somewhere on the segment between the incoming state and the equilibrium. Since both of these states are in the kinetic polytope, and since this polytope is convex, the outgoing state must lie within it as well. Thus Collision Property 1 is satisfied. Moreover, since the restriction of $`h`$ to this segment increases as one approaches the equilibrium, Collision Property 3 is also satisfied. The instabilities associated with lattice BGK operators arise because practitioners try to overrelax them. A large class of lattice BGK models for the Navier-Stokes equations have shear viscosity $`\nu (\tau 1/2)`$. In an effort to achieve lower viscosity (and hence higher Reynolds number), practitioners set $`\tau `$ to values between $`1/2`$ and unity. In this situation, the outgoing state “overshoots” the equilibrium and lies on the other side of the polytope, opposite the equilibrium state. For sufficiently small $`\tau `$, the outgoing state will lie on a contour of $`h`$ that is lower than that of the incoming state, thereby violating Collision Property 3. Still smaller values of $`\tau `$ might cause the outgoing state to lie outside the kinetic polytope altogether, thereby violating Collision Property 1. In either case, numerical instability is likely to result. ### 2.9 Entropically Stabilized BGK Operators A potential solution to this problem is suggested by our geometrical viewpoint. For $`\tau =1`$ the outgoing state is the equilibrium, for which $`h`$ is maximal. As $`\tau `$ is decreased from unity, the final value of $`h`$ decreases from its maximal value. At some value of $`\tau `$ less than unity, the outgoing value of $`h`$ will be equal to the incoming one. In order to respect Collision Property 3, we must not make $`\tau `$ any lower than this value, given by the solution to the equation $$h\left(𝐍^{}\right)=h\left(𝐍\right),$$ (31) where $`\|𝐍^{}`$ is given by Eq. (30). This may be regarded as a nonlinear algebraic equation for the single scalar unknown $`\tau `$. Indeed, this limitation on $`\tau `$ was suggested independently by Karlin et al. , and by Chen and Teixeira . It is straightforward to find the solution to this equation numerically, since it is easily bounded: We know that the desired solution has an upper bound of unity. The lower bound will be that for which the solution leaves the kinetic polytope; this happens when one of the $`N_j^{}`$’s vanish, or equivalently when $`\tau `$ is equal to the largest value of $`1N_j^{\text{eq}}/N_j`$, for $`j\{1,\mathrm{},n\}`$, that lies between zero and one. Given these two bounds on $`\tau `$, the regula falsi algorithm will reliably locate the desired solution. Call this solution $`\tau _{}(𝐍)`$. This will generally be a function of the incoming state. A useful way to parametrize $`\tau `$ is then to write $$\tau (𝐍)=\frac{\tau _{}(𝐍)}{\kappa },$$ (32) where $`0<\kappa <1`$ is a constant parameter. The case $`\kappa 0`$ corresponds to $`\tau \mathrm{}`$ so that the collision operator vanishes; the case $`\kappa 1`$ corresponds to $`\tau \tau _{}`$, which is the largest value possible that respects Collision Property 3. Thus, the entropically stabilized version of the lattice BGK equation is $$\|𝐍^{}=\|𝐍+\frac{\kappa }{\tau _{}(𝐍)}\left(\|𝐍^{\text{eq}}\|𝐍\right),$$ (33) where $`\tau _{}`$ is the solution to Eq. (31), and $`0\kappa 1`$. Of course, making $`\tau `$ a nontrivial function of the incoming state will impact the hydrodynamic equations derived from the model. The challenge to the designer of entropic lattice Boltzmann models is then to choose the $`\zeta _j`$ very judiciously, so that Eq. (33) yields the desired hydrodynamic equations, while stability is guaranteed by keeping $`0<\kappa <1`$. In constructing a lattice Boltzmann model in this fashion, as opposed to using the simpler prescription of specifying $`\tau `$, it may be argued that we are relinquishing some control over the transport coefficients. After all, if $`\nu (\tau 1/2)`$, it appears that we can specify arbitrarily small viscosity by reducing $`\tau `$. In fact, this is not the case since uncontrolled instabilities are known to set in for $`\tau `$ well above $`1/2`$. The objective of entropic lattice Boltzmann models is to allow the user some ability to overrelax the collision process, without sacrificing stability. The price that one may have to pay for this stable overrelaxation is living with a bounded transport coefficient. Thus, entropic lattice Boltzmann models for fluid flow may be restricted to minimum values of viscosity. As we shall show, however, for certain entropically stabilized lattice Boltzmann models, this minimum value can actually be zero. This allows for fully explicit, perfectly conservative, absolutely stable algorithms at arbitrarily small transport coefficient. Finally, we note that this particular prescription is only one way of creating a stable lattice Boltzmann algorithm. In fact, any map obeying Collision Properties 1 and 3 will work. More general mappings of polytopes to themselves can and should be considered, and we leave this to future study. ## 3 One-Dimensional Diffusion Model Having discussed entropic lattice Boltzmann models in general terms, we now apply the formalism to mass diffusion in one dimension. In elementary books on numerical analysis, it is demonstrated that the fully explicit finite-difference approximation to the one-dimensional diffusion equation is stable only if the Courant condition , $`\mathrm{\Delta }t(\mathrm{\Delta }x)^2/(2D)`$, is satisfied; note that this places an upper bound on the transport coefficient. It is also shown that this condition may be removed by adopting an implicit differencing scheme, such as that of Crank and Nicolson , or an alternating-direction implicit scheme, such as that of Peaceman and Rachford . The DuFort-Frankel algorithm is fully explicit and unconditionally stable, but it achieves this by a differencing scheme that involves three time steps, even though the diffusion equation is first-order in time. For the problem of achieving high Reynolds number, one would like the transport coefficient to be as small as possible. We shall show that the entropic lattice Boltzmann algorithm provides a fully explicit, perfectly conservative, two-time-step algorithm that is absolutely stable for arbitrarily small transport coefficient. While this result seems significant in and of itself, part of our purpose is pedagogical. In the course of our development of this model, we shall discuss optimal conservation representations and present the Fourier-Motzkin method for visualizing the master and kinetic polytopes. Nongeneric features of the example are noted, in preparation for more sophisticated examples in subsequent sections. ### 3.1 Description of Model We suppose that we have a regular one-dimensional lattice $``$ whose sites $`x`$ are occupied by particles whose velocities may take on one of only $`n=3`$ discrete values, namely $`1`$, $`0`$ and $`+1`$. We abuse notation slightly by letting $`j`$ take its values from the set of symbols $`\{,0,+\}`$, so that we may write the single-particle distribution as $`N_j(x)`$. Thus, the state of a given site is captured by the ket, $$\|𝐍=\left(\begin{array}{c}N_{}\hfill \\ N_0\hfill \\ N_+\hfill \end{array}\right),$$ (34) where we have suppressed the dependence on the coordinate $`x`$ for simplicity. Diffusion conserves mass, so we suppose that the mass per unit lattice site (mass density in lattice units), $$\rho =N_{}+N_0+N_+=𝝆𝐍,$$ (35) is conserved by the collisions. Here we have introduced the hydrodynamic “bra,” $`𝝆\|`$ $`=`$ $`\left(\begin{array}{ccc}\hfill +1& \hfill +1& \hfill +1\end{array}\right).`$ (37) This is the one and only hydrodynamic degree of freedom in this example. Because there are a total of three degrees of freedom, the other two must be kinetic in nature. To span these kinetic degrees of freedom and thereby make the bra basis complete, we introduce the linearly independent bras, $`𝜶\|`$ $`=`$ $`\left(\begin{array}{ccc}\hfill 1& \hfill +2& \hfill 1\end{array}\right)`$ (40) $`𝜷\|`$ $`=`$ $`\left(\begin{array}{ccc}\hfill +1& \hfill 0& \hfill 1\end{array}\right).`$ (42) Next, we form a matrix of the three bras and invert it to get the dual basis of kets, $$\left(\begin{array}{ccc}\hfill \|𝝆& \hfill \|𝜶& \hfill \|𝜷\end{array}\right)=\left(\begin{array}{c}\hfill 𝝆\|\\ \hfill 𝜶\|\\ \hfill 𝜷\|\end{array}\right)^1=\left(\begin{array}{ccc}\hfill +1& \hfill +1& \hfill +1\\ \hfill 1& \hfill +2& \hfill 1\\ \hfill +1& \hfill 0& \hfill 1\end{array}\right)^1=\frac{1}{6}\left(\begin{array}{ccc}\hfill +2& \hfill 1& \hfill +3\\ \hfill +2& \hfill +2& \hfill 0\\ \hfill +2& \hfill 1& \hfill 3\end{array}\right).$$ (43) From Eq. (43) we identify the linearly independent basis kets $$\begin{array}{ccccc}\|𝝆=\frac{1}{3}\left(\begin{array}{c}\hfill +1\\ \hfill +1\\ \hfill +1\end{array}\right)& & \|𝜶=\frac{1}{6}\left(\begin{array}{c}\hfill 1\\ \hfill +2\\ \hfill 1\end{array}\right)& & \|𝜷=\frac{1}{2}\left(\begin{array}{c}\hfill +1\\ \hfill 0\\ \hfill 1\end{array}\right).\end{array}$$ (44) The first of these is a hydrodynamic basis ket, while the last two are kinetic basis kets. One nongeneric feature should be noted: In this example, we were able to choose $`𝜶\|`$ and $`𝜷\|`$ so that each ket is proportional to the transpose of a corresponding bra. Such a choice is convenient but unnecessary. The only real requirement in choosing the kinetic bras is that they be linearly independent of each other and of the hydrodynamic bras. The next example will illustrate a situation in which there is no obvious correspondence between the individual bras and kets. We can expand the state ket in this basis as follows $$\|𝐍=\rho \|𝝆+\alpha \|𝜶+\beta \|𝜷=\frac{\rho }{3}\left(\begin{array}{c}1\frac{\overline{\alpha }}{2}+\frac{3\overline{\beta }}{2}\\ 1+\overline{\alpha }\\ 1\frac{\overline{\alpha }}{2}\frac{3\overline{\beta }}{2}\end{array}\right),$$ (45) where we have defined $`\overline{\alpha }\alpha /\rho `$ and $`\overline{\beta }\beta /\rho `$. The coefficients $`\rho `$, $`\alpha `$ and $`\beta `$ (or equivalently $`\rho `$, $`\overline{\alpha }`$ and $`\overline{\beta }`$) constitute the conservation representation. The inverse of this transformation is seen to be $`\rho `$ $`=`$ $`𝝆𝐍=+N_{}+N_0+N_+`$ $`\alpha `$ $`=`$ $`𝜶𝐍=N_{}+2N_0N_+`$ (46) $`\beta `$ $`=`$ $`𝜷𝐍=+N_{}N_+.`$ ### 3.2 Nonnegativity The nonnegativity of the components of the single-particle distribution, $`\|𝐍`$, places inequality constraints on the parameters $`\rho `$, $`\alpha `$ and $`\beta `$. For example, it is clear that $`\rho 0`$. Referring to Eq. (45), we see that we must also demand $`0`$ $``$ $`1{\displaystyle \frac{\overline{\alpha }}{2}}+{\displaystyle \frac{3\overline{\beta }}{2}}`$ $`0`$ $``$ $`1+\overline{\alpha }`$ (47) $`0`$ $``$ $`1{\displaystyle \frac{\overline{\alpha }}{2}}{\displaystyle \frac{3\overline{\beta }}{2}}.`$ It is easy to see that all three of these inequalities may be subsumed by the single statement $$1\overline{\alpha }23\left\|\overline{\beta }\right\|.$$ (48) This restricts the kinetic parameters to a triangular region in the $`(\overline{\alpha },\overline{\beta })`$ plane, as is shown in Fig. 1. Note that the effect of varying $`\rho `$ at constant $`\overline{\alpha }`$ and $`\overline{\beta }`$ is to simply scale the components of the distribution function. The triangular region bounding the parameters $`\overline{\alpha }`$ and $`\overline{\beta }`$ is then independent of $`\rho `$. That is, the bounds on the barred kinetic parameters do not depend on the hydrodynamic parameter, so there is really no need for distinction between the master and kinetic polytopes. It should be noted that this is another nongeneric feature of the present model. While it is always possible to scale out by a single nonnegative-definite hydrodynamic parameter, as we have done with $`\rho `$ in this case, more sophisticated models will have several hydrodynamic parameters. In such cases, the shape of the region bounding the kinetic parameters will depend on the remaining hydrodynamic parameters. We shall see this more clearly in the example of Section 4. ### 3.3 Optimality of Representation In this example, there is clearly a great deal of latitude in the choice of the kinetic bras. The only requirement that we have placed on these is that they be linearly independent of the hydrodynamic bra and of each other. This raises the question of whether there is some optimal choice that might be made for these. Of course, this depends entirely on what is meant by “optimal” in this context. It has been noted that the original representation of $`\|𝐍`$ is more natural for expressing the constraint of nonnegativity, while the conservation representation is more natural for expressing the collision process. For this reason, any computer implementation of entropic lattice Boltzmann methods will require frequent transformations between the two representations. This transformation is precisely what we have worked out in Eqs. (45) and (46) above. Thus, one natural figure of merit is the number of arithmetic operations required to perform such transformations. To investigate this question, rather than requiring that the kinetic bras be given by Eqs. (40) and (42), we leave them in the general form $`𝜶\|`$ $`=`$ $`\left(\begin{array}{ccc}\hfill \alpha _1& \hfill \alpha _2& \hfill \alpha _3\end{array}\right)`$ (50) $`𝜷\|`$ $`=`$ $`\left(\begin{array}{ccc}\hfill \beta _1& \hfill \beta _2& \hfill \beta _3\end{array}\right).`$ (52) When we invert to get the kets, the results are $$\begin{array}{ccccc}\|𝝆=\frac{1}{\mathrm{\Delta }}\left(\begin{array}{c}\hfill \alpha _2\beta _3\alpha _3\beta _2\\ \hfill \alpha _3\beta _1\alpha _1\beta _3\\ \hfill \alpha _1\beta _2\alpha _2\beta _1\end{array}\right)& & \|𝜶=\frac{1}{\mathrm{\Delta }}\left(\begin{array}{c}\hfill \beta _2\beta _3\\ \hfill \beta _3\beta _1\\ \hfill \beta _1\beta _2\end{array}\right)& & \|𝜷=\frac{1}{\mathrm{\Delta }}\left(\begin{array}{c}\hfill \alpha _3\alpha _2\\ \hfill \alpha _1\alpha _3\\ \hfill \alpha _2\alpha _1\end{array}\right),\end{array}$$ (53) where we have defined the determinant $$\mathrm{\Delta }\alpha _1\left(\beta _2\beta _3\right)+\alpha _2\left(\beta _3\beta _1\right)+\alpha _3\left(\beta _1\beta _2\right).$$ (54) One way to simplify the transformation process would be to find a representation in which as many components as possible of the kinetic bras and kets vanish, without making the transformation singular. This means that we want to choose the $`\alpha _j`$’s and $`\beta _j`$’s such that $`\mathrm{\Delta }0`$, while making vanish as many as possible of the following twelve quantities: $`\alpha _1,`$ $`\alpha _2,`$ $`\alpha _3,`$ $`\beta _1,`$ $`\beta _2,`$ $`\beta _3,`$ $`\alpha _3\alpha _2,`$ $`\alpha _1\alpha _3,`$ $`\alpha _2\alpha _1,`$ $`\beta _2\beta _3,`$ $`\beta _3\beta _1,`$ $`\beta _1\beta _2.`$ In this example, it is straightforward to see that there are a number of ways to make six of these quantities equal to zero. For example, we could choose $`\alpha _10,`$ $`\alpha _2=0,`$ $`\alpha _3=0,`$ $`\beta _1=0,`$ $`\beta _20,`$ $`\beta _3=0,`$ (55) resulting in $`\mathrm{\Delta }=\alpha _1\beta _20`$ and $`\alpha _3\alpha _2=0,`$ $`\alpha _1\alpha _30,`$ $`\alpha _2\alpha _10,`$ $`\beta _2\beta _30,`$ $`\beta _3\beta _1=0,`$ $`\beta _1\beta _20.`$ (56) This choice also has the added virtue of making the hydrodynamic ket equal to $$\|𝝆=\left(\begin{array}{c}0\\ 0\\ 1\end{array}\right),$$ (57) which also has two vanishing components. The computation of the $`\|𝐍`$ components from the parameters $`\rho `$, $`\overline{\alpha }`$ and $`\overline{\beta }`$ is thus reduced to $`N_{}`$ $`=`$ $`\rho \overline{\alpha }`$ $`N_0`$ $`=`$ $`\rho \overline{\beta }`$ $`N_+`$ $`=`$ $`\rho \left(1\overline{\alpha }\overline{\beta }\right),`$ involving a total of two multiplications and two additions/subtractions; this may be contrasted with Eqs. (45) which involve six multiplications and five additions/subtractions. Likewise, using the choice of Eqs. (55), the computation of the parameters $`\rho `$, $`\overline{\alpha }`$ and $`\overline{\beta }`$ from the $`𝐍`$ is reduced to $`\rho `$ $`=`$ $`N_{}+N_0+N_+`$ $`\overline{\alpha }`$ $`=`$ $`N_{}/\rho `$ $`\overline{\beta }`$ $`=`$ $`N_0/\rho ,`$ involving two additions/subtractions and two divisions; this may be contrasted to Eqs. (46) which involve five additions/subtractions, one multiplication, and two divisions. Clearly, more sophisticated figures of merit could be devised to optimize the transformations used in lattice Boltzmann computations. As we have noted, vanishing components of the kinetic bras and kets eliminate the addition/subtraction of terms. Likewise, components equal to $`\pm 1`$ do require an addition/subtraction, but not a multiplication, and this fact could be taken into account in a more refined figure of merit. The computation of the collision outcome is the principal “inner loop” of a lattice Boltzmann computation, insofar as it must be performed at each site of a spatial grid at each time step. Such considerations may be especially important for lattice Boltzmann models with large numbers of velocities. ### 3.4 Fourier-Motzkin Elimination For this simple example, we had no difficulty visualizing the master and kinetic polytopes. In preparation for the succeeding sections where the task will be substantially more difficult, however, we take this opportunity to introduce the Fourier-Motzkin elimination method for this purpose. The algorithm consists of the following sequence of steps: 1. We rewrite the set of inequalities, Eq. (47) in matrix format, $$\left(\begin{array}{ccc}\hfill \frac{1}{2}& \hfill +\frac{3}{2}& \hfill +1\\ \hfill +1& \hfill 0& \hfill +1\\ \hfill \frac{1}{2}& \hfill \frac{3}{2}& \hfill +1\end{array}\right)\left(\begin{array}{c}\overline{\alpha }\\ \overline{\beta }\\ 1\end{array}\right)0$$ (58) We have adopted the convention of including constant terms in an extra column of the matrix, using the device of appending $`1`$ to the column vector of unknowns. In general, there will be $`m`$ inequalities for $`n`$ unknowns, and the matrix will be of size $`m\times (n+1)`$. 2. We scale each inequality by a positive factor so that the pivot <sup>5</sup><sup>5</sup>5As in discussions of Gaussian elimination, the term “pivot” refers to the first nonzero entry in a row. is either $`1`$, $`0`$ or $`+1`$. (Recall that scaling by a positive factor preserves the sense of an inequality.) We then reorder the inequalities, sorting by their (scaled) pivots so that the zero pivots are last. Beginning with Eq. (58), this yields $$\left(\begin{array}{ccc}\hfill 1& \hfill +3& \hfill +2\\ \hfill 1& \hfill 3& \hfill +2\\ \hfill +1& \hfill 0& \hfill +1\end{array}\right)\left(\begin{array}{c}\overline{\alpha }\\ \overline{\beta }\\ 1\end{array}\right)0.$$ (59) Note that there were no zero pivots in this first step. We did, however, reorder the inequalities so that those with pivot $`1`$ are together, and preceed that with pivot $`+1`$. 3. We now add all pairs of inequalities, such that the first member of the pair has pivot $`1`$ and the second member of the pair has pivot $`+1`$, and we append the new zero-pivot inequalities thus obtained to the system. If there are $`m_{}`$ inequalities with pivot $`1`$, and $`m_+`$ inequalities with pivot $`+1`$, this results in the addition of $`m_{}m_+`$ new inequalities to the system. Since there were $`mm_{}m_+`$ zero-pivot inequalities to begin with, the new system will have a total of $`m_0=mm_{}m_++m_{}m_+`$ zero-pivot inequalities. For our above example, $`m=3`$, $`m_{}=2`$ and $`m_+=1`$, so we get $`m_0=2`$ zero-pivot inequalities in the system, which can now be written $$\left(\begin{array}{ccc}\hfill 1& \hfill +3& \hfill +2\\ \hfill 1& \hfill 3& \hfill +2\\ \hfill +1& \hfill 0& \hfill +1\\ \hfill 0& \hfill +3& \hfill +3\\ \hfill 0& \hfill 3& \hfill +3\end{array}\right)\left(\begin{array}{c}\overline{\alpha }\\ \overline{\beta }\\ 1\end{array}\right)0.$$ (60) 4. Finally, we recurse by returning to step 2 for the $`m_0\times n`$ submatrix obtained by taking only the zero-pivot inequalities and neglecting the first unknown (since the zero-pivot inequalities do not involve it anyway). We continue in this fashion until all of the first $`n`$ elements of the rows of the matrix are zero. For our example, the recursion step asks us to return to step 2 for the submatrix indicated below $$\left(\text{ }\text{-1+3+2-1-3+2+10+10+3+30-3+3}\text{ }\right)\left(\begin{array}{c}\overline{\alpha }\\ \overline{\beta }\\ 1\end{array}\right)0.$$ (61) As before, we begin by normalizing the rows and sorting them by pivot, $$\left(\text{ }\text{-1+3+2-1-3+2+10+10-1+10+1+1}\text{ }\right)\left(\begin{array}{c}\overline{\alpha }\\ \overline{\beta }\\ 1\end{array}\right)0.$$ (62) In this example there is now one row of the submatrix with pivot $`1`$ and one with pivot $`+1`$, so we append the sum of these to the system to get $$\left(\text{ }\text{-1+3+2-1-3+2+10+10-1+10+1+100+2}\text{ }\right)\left(\begin{array}{c}\overline{\alpha }\\ \overline{\beta }\\ 1\end{array}\right)0.$$ (63) At this point, the first $`n`$ elements of the last row of the matrix are zero, so we stop the recursion and examine the last row. It states a true inequality, namely $`+20`$, so we can conclude that the original inequalities are not mutually exclusive; that is, a nonnull polytope of solutions will exist. If the $`n+1`$ element of the final row had been negative, we would have concluded that no such solution was possible. Once we have stopped the recursion and established the consistency of the inequalities, we can get the final set of inequalities by “back substituting” up the matrix. The last two nontrivial inequalities tell us that $`1\overline{\beta }+1`$, or $$\left\|\overline{\beta }\right\|1.$$ (64) The first three equations then yield $`1\overline{\alpha }\mathrm{min}(23\overline{\beta },2+3\overline{\beta })`$ which is equivalent to Eq. (48), describing the triangular region of Fig. 1. More generally, each variable will have a lower bound that is the maximum of all the lower bounds determined by the inequalities with pivot $`+1`$, and an upper bound that is the minimum of all the upper bounds determined by the inequalities with pivot $`1`$. The arguments of the maximum and minimum functions will depend only on those variables that have already been bounded. If we were using the technique to find limits of integration for a multiple integral, the inequalities at the bottom of the Fourier-Motzkin-eliminated matrix would correspond to the outermost integrals. We shall revisit this technique in the next section. ### 3.5 Collision Operator To illustrate the construction of an entropically stable collision operator for this model, we adopt the simplest possible $`H`$ function by taking $`\zeta _j(x)=x`$. Since there is only one hydrodynamic degree of freedom, Eqs. (27) reduce to $$N_{}^{\text{eq}}=N_0^{\text{eq}}=N_+^{\text{eq}}=\frac{1}{Q}.$$ (65) Eq. (28) then tells us that $`Q=3/\rho `$, so $$𝐍^{\text{eq}}=\frac{\rho }{3}\left(\begin{array}{c}1\\ 1\\ 1\end{array}\right).$$ (66) Comparing this with Eq. (45), we see that, within the kinetic polytope, the equilibrium point is the origin, $`\alpha ^{\text{eq}}=\beta ^{\text{eq}}=0`$. A contour plot of the $`H`$ function on the kinetic polytope is presented in Fig. 2. The lattice BGK equation in the conservation representation, Eq. (30), is then equivalent to the prescription $$\rho ^{}=\rho $$ (67) and $$\left(\begin{array}{c}\alpha ^{}\hfill \\ \beta ^{}\hfill \end{array}\right)=\left(\begin{array}{c}\alpha \hfill \\ \beta \hfill \end{array}\right)+\frac{1}{\tau }\left[\left(\begin{array}{c}\alpha ^{\text{eq}}\hfill \\ \beta ^{\text{eq}}\hfill \end{array}\right)\left(\begin{array}{c}\alpha \hfill \\ \beta \hfill \end{array}\right)\right]=z\left(\begin{array}{c}\alpha \hfill \\ \beta \hfill \end{array}\right),$$ (68) where $`z11/\tau `$. These can also be written as the single equation $$\left(\begin{array}{c}\rho ^{}\hfill \\ \alpha ^{}\hfill \\ \beta ^{}\hfill \end{array}\right)=\left(\begin{array}{ccc}1& 0& 0\\ 0& z& 0\\ 0& 0& z\end{array}\right)\left(\begin{array}{c}\rho \hfill \\ \alpha \hfill \\ \beta \hfill \end{array}\right).$$ (69) As promised, the collision is dramatically simplified in this representation. The hydrodynamic parameter is unchanged; only the two kinetic parameters are altered by the collision. Transformation back to the original representation yields $$\|𝐍^{}=\frac{\rho }{3}\left(\begin{array}{c}1\frac{z\overline{\alpha }}{2}+\frac{3z\overline{\beta }}{2}\\ 1+z\overline{\alpha }\\ 1\frac{z\overline{\alpha }}{2}\frac{3z\overline{\beta }}{2}\end{array}\right),$$ (70) or equivalently $$\|𝐍^{}=\frac{1}{3}\left(\begin{array}{ccc}1+2z& 1z& 1z\\ 1z& 1+2z& 1z\\ 1z& 1z& 1+2z\end{array}\right)\|𝐍.$$ (71) ### 3.6 Entropic Stabilization The condition for marginal entropic stabilization, Eq. (31), for this model is then $$\mathrm{ln}N_+^{}+\mathrm{ln}N_0^{}+\mathrm{ln}N_{}^{}=\mathrm{ln}N_++\mathrm{ln}N_0+\mathrm{ln}N_{}.$$ (72) This needs to be solved for the limiting value $`z=z_{}`$, where $`z_{}=11/\tau _{}`$. Using Eqs. (45) and (70), this becomes $$\left(1\frac{z_{}\overline{\alpha }}{2}+\frac{3z_{}\overline{\beta }}{2}\right)\left(1+z_{}\overline{\alpha }\right)\left(1\frac{z_{}\overline{\alpha }}{2}\frac{3z_{}\overline{\beta }}{2}\right)=\left(1\frac{\overline{\alpha }}{2}+\frac{3\overline{\beta }}{2}\right)\left(1+\overline{\alpha }\right)\left(1\frac{\overline{\alpha }}{2}\frac{3\overline{\beta }}{2}\right).$$ (73) This appears to be a cubic equation for $`z_{}`$, but in fact $`z_{}=1`$ (corresponding to $`\tau _{}\mathrm{}`$) is clearly a root. It is an uninteresting root since it corresponds to no collision. Once it is removed, we are left with a quadratic for $`z_{}`$, so the relevant solution is easily written in closed form, $$z_{}=\frac{3\overline{\alpha }^2\overline{\alpha }^3+9\overline{\beta }^2+9\overline{\alpha }\overline{\beta }^2\sqrt{3(\overline{\alpha }^2+\overline{\alpha }^3+3\overline{\beta }^29\overline{\alpha }\overline{\beta }^2)(3\overline{\alpha }^2\overline{\alpha }^3+9\overline{\beta }^2+9\overline{\alpha }\overline{\beta }^2)}}{2(\overline{\alpha }^39\overline{\alpha }\overline{\beta }^2)}.$$ (74) A contour plot of the function $`z_{}(\overline{\alpha },\overline{\beta })`$ on the kinetic polytope is displayed in Fig. 3. There it can be seen that $`z_{}`$ has minimum value of $`2`$ at the midpoints of the sides of the triangle, has maximum value of $`1/2`$ at the vertices, and approaches $`1`$ at the equilibrium $`\overline{\alpha }=\overline{\beta }=0`$. The entropically stabilized collision operator is then given by Eq. (70), with $`\tau \tau _{}/\kappa `$, or equivalently $$z1\kappa (1z_{}).$$ (75) It is useful to discuss three possible choices for $`\kappa `$: * $`\kappa 0`$ means that $`z1`$ and $`\tau \mathrm{}`$. This is the uninteresting limit in which the collision operator vanishes and hydrodynamic behavior is lost. * $`\kappa =1/2`$ means that $`z=(1+z_{})/2`$ and $`\tau =2/(1z_{})`$. This is interesting because near equilibrium, where $`z_{}1`$, it coincides with the prescription $`\tau =1`$. This is the limit, mentioned in Subsection 2.9, wherein the outgoing state is the local equilibrium. This is the smallest value of $`\tau `$ for which the standard lattice BGK algorithm is guaranteed to be stable. (In fairness, the standard algorithm is likely to be stable for smaller $`\tau `$; it’s just that this is not guaranteed.) * $`\kappa 1`$ means that $`zz_{}`$ and $`\tau \tau _{}`$. This is the largest value of $`\tau `$ for which the entropic lattice BGK algorithm is guaranteed to be stable. ### 3.7 Attainable Transport Coefficient The complicated form of Eq. (74) may lead one to believe that the hydrodynamic equation (in this case, a diffusion equation) would be very difficult to derive for the entropically stabilized collision operator. In fact, this is not the case at all. The Chapman-Enskog analysis that yields the hydrodynamic equations requires only that one linearize the collision operator about the local equilibrium. Since the local equilibrium has $`\alpha =\beta =0`$, it is clear from Eqs. (68) and (69), that only the value of $`z`$ at equilibrium will enter. Since $`z_{}1`$ at the equilibrium, Eq. (75) indicates that the transport coefficient will be precisely that obtained by the standard lattice BGK algorithm with $`\tau =1/(2\kappa )`$. The Chapman-Enskog analysis of the standard lattice BGK algorithm is an elementary exercise . The result for the diffusivity in natural lattice units is $$D=\frac{1}{3}\left(\tau \frac{1}{2}\right).$$ (76) Hence, by the above argument, the result for the diffusivity of the entropic lattice BGK algorithm is $$D=\frac{1}{6}\left(\frac{1}{\kappa }1\right).$$ (77) From this, we can clearly see the benefit of the entropic algorithm. The standard lattice BGK algorithm is not guaranteed stable for $`\tau <1`$ or $`\kappa >1/2`$, as discussed above. From Eq. (76) or (77), respectively, we see that this corresponds to a minimum diffusivity of $`1/6`$. By contrast, the entropic lattice BGK algorithm loses stability only when $`\kappa =1`$, corresponding to $`D=0`$. It is remarkable that we have found a stable, conservative, explicit algorithm that seems to allow the transport coefficient to become arbitrarily small. After all, the analysis leading to the shear viscosity of a fluid model is not very different from that leading to the diffusivity above, and arbitrarily small shear viscosity would allow for arbitrarily large Reynolds number. In the field of computational fluid dynamics, however, results that seem too good to be true usually are just that. For one thing, stability does not imply accuracy. A perfectly stable algorithm is not terribly useful if it does not converge to the correct answer. Flows at higher Reynolds number involve ever smaller eddies, and at some point the lattice spacing becomes insufficient to resolve these. While it is comforting to have an algorithm that does not lose stability in this situation, it is probably a mistake to attach much physical significance to its results. Another problem is that the time required for the system to come to equilibrium goes to infinity as $`\kappa 1`$. In fact, when $`\kappa =1`$ the entropic lattice BGK collision operator is its own inverse: If incoming state $`\|𝐍`$ yields outgoing state $`\|𝐍^{}`$, then incoming state $`\|𝐍^{}`$ would yield outgoing state $`\|𝐍`$. This means that if the entire lattice were initialized slightly away from equilibrium with no spatial gradients whatsoever, it would simply thrash back and forth between two states, without ever converging to the desired equilibrium. When the time required to achieve local equilibrium exceeds time scales of interest, hydrodynamic behavior breaks down. To illustrate this problem, we simulated the entropically stabilized diffusion model with the initial density profile $$\rho =\rho _0+\rho _1\mathrm{sin}\left(\frac{2\pi \mathrm{}x}{L}\right)$$ (78) on a periodic lattice. We used lattice size $`L=32`$, and wavenumber $`\mathrm{}=3`$ in all the simulations reported below. We initialized the state of each site in a local equilibrium $`N_{}=N_0=N_+=\rho /3`$, without the Chapman-Enskog correction due to gradient. Now Eq. (78) is a solution of the diffusion equation, if $`\rho _1`$ decays in time as $`\mathrm{exp}(\gamma t)`$, where $$\gamma =D\left(\frac{2\pi \mathrm{}}{L}\right)^2.$$ (79) So we fit our results to the functional form of Eq. (78), measured $`\rho _1(t)`$, and did a least-squares fit of its logarithm to a linear function of time in order to determine $`\gamma `$. We then used Eq. (79) to get the diffusion coefficient, and we compared this to the theoretical value provided by Eq. (77) for several different values of $`\kappa `$, approaching unity. Fig. 4 shows the measured value of $`\rho _1(t)`$ for $`\kappa =0.9`$ and $`\kappa =0.99`$. As can be seen, there is an initial transient, due to the inadequacy of the form used for the local equilibrium $`N_{}=N_0=N_+=\rho /3`$ in the presence of a spatial gradient. The left-hand side of Fig. 5 shows the measured value of $`\rho _1(t)`$ for $`\kappa =0.999`$, and the right-hand side is an enlargement of the transient region. Fig. 6 shows the same things for $`\kappa =0.9999`$. Note that the transient period lengthens considerably as $`\kappa `$ nears unity. We could have substantially reduced these transient periods had we used the Chapman-Enskog correction to the local equilibrium; instead, we simply waited for the transient to die away before measuring the decay constant $`\gamma `$. The results for the diffusivity are displayed in Table 1. Note that we find reasonable agreement until we get to the cases $`\kappa =0.9999`$ and $`\kappa =0.99999`$. To see what is going wrong for these cases, the upper-left-hand corner of Fig. 7 shows the measured value of $`\rho _1(t)`$ for $`\kappa =0.99999`$, and the upper-right-hand corner shows the same plot with a reduced range for the ordinate. In addition to the oddly shaped envelope of the initial transient period, similar to that seen in Fig. 6, and enlarged in the lower part of this figure, we note that the plot of $`\rho _1(t)`$ in the upper right seems to be increasing in thickness, even after the initial transient. Upon further enlargement, shown on the left-hand side of Fig. 8, we see that this thickness is in fact a subtle transient of extremely long duration. The right-hand side of the figure shows this same transient at the longest time for which the simulation in Fig. 7 was run. Though reduced in magnitude (the scales of the ordinates of both graphs in Fig. 8 are equal), the transient is still present, even though most of the signal itself has diffused away. We believe that this transient is causing the anomaly in the measured diffusivity. A much longer simulation would be required to test this assertion. To better understand this loss of separation between the kinetic and hydrodynamic time scales, note that the time required for the signal to decay away is $$\tau _{\text{signal}}\frac{1}{\gamma }\frac{1}{D}\frac{\kappa }{1\kappa }.$$ (80) The time required for the transient to decay to $`1/e`$ of its initial value, on the other hand, obeys $$\kappa ^{\tau _{\text{transient}}}=\frac{1}{e},$$ (81) whence $$\tau _{\text{transient}}\frac{1}{\mathrm{ln}\kappa }.$$ (82) For $`\kappa =1ϵ`$, where $`ϵ`$ is small, the characteristic times in Eqs. (80) and (82) both go like $`1/ϵ`$. Thus, the transients begin to linger for hydrodynamic time scales, and it becomes impossible to separate the kinetic behavior from the hydrodynamic behavior. When the time required for the decay of transients is comparable to hydrodynamic time scales of interest, the usefulness of the simulation is questionable. Thus, while Table 1 indicates good agreement with theory, and the transport coefficient really does tend to zero without loss of stability, one must be exceedingly careful in the preparation of the initial condition to exploit this feature of entropic lattice Boltzmann models. In particular, transient behavior could be dramatically reduced by including the Chapman-Enskog correction to first-order in the gradient. Indeed, if this problem turns out to be the only obstacle to stable, conservative, explicit algorithms with arbitrarily small transport coefficients, higher-order gradient corrections may also be worth investigating. ## 4 One-Dimensional Compressible Fluid Model In this section, we apply the entropic lattice Boltzmann method to a simple five-velocity model of fluid dynamics in one dimension, first considered by Renda et al. in 1997 . We shall find that the geometric picture is much richer than that for the diffusion model. The master polytope for this model is four-dimensional, and we shall show that the Fourier-Motzkin algorithm is very useful for describing it. ### 4.1 Description of Model As a second example, we consider a lattice Boltzmann model for a one-dimensional compressible fluid with conserved mass, momentum and energy, first studied by Renda et al. . The velocity space of this model consists of five discrete values of velocity, namely $`0`$, $`\pm 1`$, and $`\pm 2`$. The single-particle distribution at site $`x`$ with velocity $`j\{2,1,0,+1,+2\}`$ is denoted by $`N_j`$. As in the last example, the state vector $`\|𝐍`$ at a given site $`x`$ will be denoted by a ket, $$\|𝐍=\left(\begin{array}{c}N_2\hfill \\ N_1\hfill \\ N_0\hfill \\ N_{+1}\hfill \\ N_{+2}\hfill \end{array}\right),$$ (83) where we have suppressed the dependence on $`x`$ for simplicity. A compressible fluid conserves mass, momentum and kinetic energy, so we suppose that the corresponding densities $`\rho `$ $`=`$ $`{\displaystyle \underset{j=2}{\overset{+2}{}}}N_j=𝝆𝐍`$ (84) $`\pi `$ $`=`$ $`{\displaystyle \underset{j=2}{\overset{+2}{}}}jN_j=𝝅𝐍`$ (85) $`\epsilon `$ $`=`$ $`{\displaystyle \underset{j=2}{\overset{+2}{}}}{\displaystyle \frac{j^2}{2}}N_j=𝜺𝐍,`$ (86) are conserved by the collisions. Here we have introduced the bras $`\begin{array}{ccc}𝝆\|_j& =& 1\\ & & \\ 𝝅\|_j& =& j\\ & & \\ 𝜺\|_j& =& j^2/2,\end{array}`$ (92) or explicitly $`𝝆\|`$ $``$ $`\left(\begin{array}{ccccc}\hfill +1& \hfill +1& \hfill +1& \hfill +1& \hfill +1\end{array}\right)`$ (94) $`𝝅\|`$ $``$ $`\left(\begin{array}{ccccc}\hfill 2& \hfill 1& \hfill 0& \hfill +1& \hfill +2\end{array}\right)`$ (96) $`𝜺\|`$ $``$ $`{\displaystyle \frac{1}{2}}\left(\begin{array}{ccccc}\hfill +4& \hfill +1& \hfill 0& \hfill +1& \hfill +4\end{array}\right).`$ (98) These are the three hydrodynamic degrees of freedom in this example. Because there are a total of five degrees of freedom, the other two must be kinetic in nature. To span these kinetic degrees of freedom and thereby make the bra basis complete, we introduce the linearly independent bras, $`𝜶\|`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\begin{array}{ccccc}\hfill 0& \hfill +1& \hfill 0& \hfill +1& \hfill 0\end{array}\right)`$ (100) $`𝜷\|`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\begin{array}{ccccc}\hfill 0& \hfill 1& \hfill 0& \hfill +1& \hfill 0\end{array}\right).`$ (102) Next, we form a matrix of the five bras and invert it to get the dual basis of kets, $`\left(\begin{array}{ccccc}\|𝝆& \|𝝅& \|𝜺& \|𝜶& \|𝜷\end{array}\right)=`$ (120) $`\left(\begin{array}{c}𝝆\|\\ 𝝅\|\\ 𝜺\|\\ 𝜶\|\\ 𝜷\|\end{array}\right)^1=\left[{\displaystyle \frac{1}{2}}\left(\begin{array}{ccccc}\hfill +2& \hfill +2& \hfill +2& \hfill +2& \hfill +2\\ \hfill 4& \hfill 2& \hfill 0& \hfill +2& \hfill +4\\ \hfill +4& \hfill +1& \hfill 0& \hfill +1& \hfill +4\\ \hfill 0& \hfill +1& \hfill 0& \hfill +1& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0& \hfill +1& \hfill 0\end{array}\right)\right]^1={\displaystyle \frac{1}{4}}\left(\begin{array}{ccccc}\hfill 0& \hfill 1& \hfill +1& \hfill 1& \hfill +2\\ \hfill 0& \hfill 0& \hfill 0& \hfill +4& \hfill 4\\ \hfill +4& \hfill 0& \hfill 2& \hfill 6& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill +4& \hfill +4\\ \hfill 0& \hfill +1& \hfill +1& \hfill 1& \hfill 2\end{array}\right).`$ From Eq. (120) we identify the linearly independent basis kets, consisting of the hydrodynamic basis kets $$\begin{array}{ccc}\|𝝆=\left(\begin{array}{c}\hfill 0\\ \hfill 0\\ \hfill +1\\ \hfill 0\\ \hfill 0\end{array}\right),\hfill & \|𝝅=\frac{1}{4}\left(\begin{array}{c}\hfill 1\\ \hfill 0\\ \hfill 0\\ \hfill 0\\ \hfill +1\end{array}\right),\hfill & \|𝜺=\frac{1}{4}\left(\begin{array}{c}\hfill +1\\ \hfill 0\\ \hfill 2\\ \hfill 0\\ \hfill +1\end{array}\right),\hfill \end{array}$$ (121) and the kinetic basis kets $$\begin{array}{cc}\|𝜶=\frac{1}{4}\left(\begin{array}{c}\hfill 1\\ \hfill +4\\ \hfill 6\\ \hfill +4\\ \hfill 1\end{array}\right),\hfill & \|𝜷=\frac{1}{2}\left(\begin{array}{c}\hfill +1\\ \hfill 2\\ \hfill 0\\ \hfill +2\\ \hfill 1\end{array}\right).\hfill \end{array}$$ (122) Unlike the previous example, there is no obvious correspondence between the individual bra and ket basis elements. This is because there is no natural notion of a metric in this space. Nevertheless, as noted in Section 2, it is always possible to construct an entire ket basis from an entire bra basis, and that is what we have done here. We can then expand the state ket in this basis as follows $$\|𝐍=\rho \|𝝆+\pi \|𝝅+\epsilon \|𝜺+\alpha \|𝜶+\beta \|𝜷=\rho \left(\begin{array}{c}\frac{1}{4}\left(\overline{\epsilon }\overline{\pi }\overline{\alpha }+2\overline{\beta }\right)\\ \overline{\alpha }\overline{\beta }\\ 1\frac{1}{2}\left(\overline{\epsilon }+3\overline{\alpha }\right)\\ \overline{\alpha }+\overline{\beta }\\ \frac{1}{4}\left(\overline{\epsilon }+\overline{\pi }\overline{\alpha }2\overline{\beta }\right)\end{array}\right),$$ (123) where we have continued to use an overbar to denote conserved quantities per unit mass; for example, $`\overline{\pi }\pi /\rho `$ is the hydrodynamic velocity. As with the last example, this representation of the distribution function $`\|𝐍`$ has the virtue of making the conserved quantities manifest. In this case, we have the inverse transformation $`\rho `$ $`=`$ $`𝝆𝐍=+N_2+N_1+N_0+N_{+1}+N_{+2}`$ $`\pi `$ $`=`$ $`𝝅𝐍=2N_2N_1+N_{+1}+2N_{+2}`$ $`\epsilon `$ $`=`$ $`𝜺𝐍=+2N_2+N_1/2+N_{+1}/2+2N_{+2}`$ (124) $`\alpha `$ $`=`$ $`𝜶𝐍=+N_1/2+N_{+1}/2`$ $`\beta `$ $`=`$ $`𝜷𝐍=N_1/2+N_{+1}/2.`$ ### 4.2 Nonnegativity We now consider the shape of the master and kinetic polytopes for the compressible fluid model described in the last subsection. Here we shall find a much richer structure than the simple triangular region that we found for the diffusive example in the last section above. Referring to Eq. (123), we see that that nonnegativity of the components of the single-particle distribution function is guaranteed by the set of inequalities $`\rho 0`$ and $`0`$ $``$ $`\overline{\epsilon }\overline{\pi }\overline{\alpha }+2\overline{\beta }`$ $`0`$ $``$ $`\overline{\alpha }\overline{\beta }`$ $`0`$ $``$ $`2\overline{\epsilon }3\overline{\alpha }`$ (125) $`0`$ $``$ $`\overline{\alpha }+\overline{\beta }`$ $`0`$ $``$ $`\overline{\epsilon }+\overline{\pi }\overline{\alpha }2\overline{\beta }.`$ These define the master polytope in the four dimensional $`(\overline{\pi },\overline{\epsilon },\overline{\alpha },\overline{\beta })`$ space. If we fix the hydrodynamic parameters, $`\overline{\pi }`$ and $`\overline{\epsilon }`$, the above still constitute a set of linear inequalities specifying the corresponding kinetic polytope in the two dimensional $`(\overline{\alpha },\overline{\beta })`$ space. Geometrically, these are the intersections of the master polytope with the (hyper)planes of fixed hydrodynamic parameters $`\overline{\pi }`$ and $`\overline{\epsilon }`$. Thus, though the mass density $`\rho `$ still scales out of the distribution function, the shape of the kinetic polytopes in the $`(\overline{\alpha },\overline{\beta })`$ plane will depend on the hydrodynamic velocity $`\overline{\pi }`$ and the kinetic energy per unit mass $`\overline{\epsilon }`$. It is not particularly easy to visualize the shape of the kinetic polytope in the $`(\overline{\alpha },\overline{\beta })`$ plane for given hydrodynamic parameters $`\overline{\pi }`$ and $`\overline{\epsilon }`$. In fact, the task is tedious enough that we have relegated it to Appendix A, which takes the reader on a detailed tour of the four-dimensional master polytope and its kinetic polytope cross sections. To summarize the conclusions of that Appendix, the projection of the master polytope on the $`(\overline{\pi },\overline{\epsilon })`$ plane is illustrated in Fig. 9. Values of $`(\overline{\pi },\overline{\epsilon })`$ outside the shaded regions are not possible. The shape of the kinetic polytope is then a triangle when $`(\overline{\pi },\overline{\epsilon })`$ is in the most lightly shaded regions of Fig. 9, a quadrilateral when $`(\overline{\pi },\overline{\epsilon })`$ is in the intermediately shaded regions of Fig. 9, and a pentagon when $`(\overline{\pi },\overline{\epsilon })`$ is in the most darkly shaded region of Fig. 9. Illustrations and detailed descriptions of the kinetic polytopes for particular values of $`(\overline{\pi },\overline{\epsilon })`$ in each of the seven distinct regions with $`\overline{\pi }0`$ are provided in the seven figures of Appendix A; the chart below Fig. 9 indicates which figure corresponds to each of its regions. Finally, since the inequalities are invariant under $`(\overline{\pi },\overline{\beta })(\overline{\pi },\overline{\beta })`$, the kinetic polytopes for $`\overline{\pi }0`$ are obtained from those for $`\overline{\pi }0`$ by simply reflecting in $`\overline{\beta }`$. ### 4.3 Fourier-Motzkin Analysis for the 1D Compressible Fluid Model To compute the master polytope for the compressible fluid model, we cast the inequalities of Eqs. (125) into matrix format as follows, $$\left(\begin{array}{ccccc}\hfill 1& \hfill +2& \hfill 1& \hfill +1& \hfill 0\\ \hfill +1& \hfill 1& \hfill 0& \hfill 0& \hfill 0\\ \hfill 3& \hfill 0& \hfill 0& \hfill 1& \hfill +2\\ \hfill +1& \hfill +1& \hfill 0& \hfill 0& \hfill 0\\ \hfill 1& \hfill 2& \hfill +1& \hfill +1& \hfill 0\end{array}\right)\left(\begin{array}{c}\overline{\alpha }\\ \overline{\beta }\\ \overline{\pi }\\ \overline{\epsilon }\\ 1\end{array}\right)0.$$ (126) Here we have put the hydrodynamic parameters $`\overline{\pi }`$ and $`\overline{\epsilon }`$ after the kinetic parameters $`\overline{\alpha }`$ and $`\overline{\beta }`$ in the column vector of unknowns, for reasons that will become apparent in a moment. Performing the Fourier-Motzkin algorithm on this set of inequalities, we finally arrive at the 23 inequalities, $$\left(\begin{array}{ccccc}\hfill 1& \hfill 2& \hfill +1& \hfill +1& \hfill 0\\ \hfill 1& \hfill +2& \hfill 1& \hfill +1& \hfill 0\\ \hfill 1& \hfill 0& \hfill 0& \hfill \frac{1}{3}& \hfill +\frac{2}{3}\\ \hfill +1& \hfill 1& \hfill 0& \hfill 0& \hfill 0\\ \hfill +1& \hfill +1& \hfill 0& \hfill 0& \hfill 0\\ \hfill 0& \hfill 1& \hfill +1& \hfill +1& \hfill 0\\ \hfill 0& \hfill 1& \hfill +\frac{1}{3}& \hfill +\frac{1}{3}& \hfill 0\\ \hfill 0& \hfill 1& \hfill 0& \hfill \frac{1}{3}& \hfill +\frac{2}{3}\\ \hfill 0& \hfill +1& \hfill 1& \hfill +1& \hfill 0\\ \hfill 0& \hfill +1& \hfill \frac{1}{3}& \hfill +\frac{1}{3}& \hfill 0\\ \hfill 0& \hfill +1& \hfill 0& \hfill \frac{1}{3}& \hfill +\frac{2}{3}\\ \hfill 0& \hfill 0& \hfill 1& \hfill +2& \hfill 0\\ \hfill 0& \hfill 0& \hfill 1& \hfill +\frac{2}{3}& \hfill +\frac{2}{3}\\ \hfill 0& \hfill 0& \hfill 1& \hfill 0& \hfill +2\\ \hfill 0& \hfill 0& \hfill +1& \hfill +2& \hfill 0\\ \hfill 0& \hfill 0& \hfill +1& \hfill +\frac{2}{3}& \hfill +\frac{2}{3}\\ \hfill 0& \hfill 0& \hfill +1& \hfill 0& \hfill +2\\ \hfill 0& \hfill 0& \hfill 0& \hfill 1& \hfill +2\\ \hfill 0& \hfill 0& \hfill 0& \hfill +1& \hfill +4\\ \hfill 0& \hfill 0& \hfill 0& \hfill +1& \hfill +1\\ \hfill 0& \hfill 0& \hfill 0& \hfill +1& \hfill +\frac{1}{4}\\ \hfill 0& \hfill 0& \hfill 0& \hfill +1& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0& \hfill 0& \hfill +1\end{array}\right)\left(\begin{array}{c}\overline{\alpha }\\ \overline{\beta }\\ \overline{\pi }\\ \overline{\epsilon }\\ 1\end{array}\right)0.$$ (127) The last of these is a consistency condition which is seen to be satisfied. The five inequalities prior to that yield $$\mathrm{max}(0,\frac{1}{4},1,4)\overline{\epsilon }2,$$ or $$0\overline{\epsilon }2,$$ (128) which explains the extent of Fig. 9 along the ordinate. The six inequalities prior to that yield $$\mathrm{max}(2,\frac{2}{3}\frac{2}{3}\overline{\epsilon },2\overline{\epsilon })\overline{\pi }\mathrm{min}(+2,+\frac{2}{3}+\frac{2}{3}\overline{\epsilon },+2\overline{\epsilon }),$$ or $$\left\|\overline{\pi }\right\|\mathrm{min}(2,\frac{2}{3}+\frac{2}{3}\overline{\epsilon },2\overline{\epsilon }),$$ (129) which explains the extent of Fig. 9 along the abscissa. Thus, we have bounded the projection of the master polytope in the hydrodynamic parameter space. Given hydrodynamic parameters in the allowed region thus determined, the first 11 inequalities of Eq. (127) then specify the kinetic polytopes. (This is why we put the hydrodynamic parameters after the kinetic ones in the column vector of unknowns.) Continuing our back substitution, the bounds on $`\overline{\beta }`$ are $$\mathrm{max}(\frac{2\overline{\epsilon }}{3},\overline{\epsilon }+\overline{\pi },\frac{\overline{\epsilon }\overline{\pi }}{3})\overline{\beta }\mathrm{min}(+\frac{2\overline{\epsilon }}{3},+\overline{\epsilon }+\overline{\pi },+\frac{\overline{\epsilon }+\overline{\pi }}{3}).$$ (130) Finally, the bounds on $`\overline{\alpha }`$ are $$\mathrm{max}(\overline{\beta },\overline{\beta })\overline{\alpha }\mathrm{min}(+2\overline{\beta }\overline{\pi }+\overline{\epsilon },2\overline{\beta }+\overline{\pi }+\overline{\epsilon }).$$ (131) Eqs. (130) and (131) specify the shape of the kinetic polytope in the $`(\overline{\alpha },\overline{\beta })`$ plane, given the hydrodynamic parameters $`\overline{\pi }`$ and $`\overline{\epsilon }`$. In fact, they provide a more succinct way of determining and expressing the bounds of both the master polytope and the kinetic polytopes than the analysis presented in Appendix A. ## 5 Two and Three-Dimensional Fluid Dynamics In this section, we apply our methodology to much richer models, namely those most commonly used for two and three-dimensional Navier-Stokes hydrodynamics. Thought the Fourier-Motzkin method cannot easily be applied to the latter, we are able to present a valid collision representation for both examples. We indicate how a computer simulation for such models might be implemented. ### 5.1 FHP Model The diffusion and compressible fluid models that we considered above had three and five particle velocities, respectively. Such models are of academic and pedagogical interest, but not particularly useful as simulational tools. The smallest model that has been used for serious computational fluid dynamics calculations in two dimensions is called the FHP model. This is a six-velocity model on a triangular grid which has been shown to yield isotropic, incompressible Navier-Stokes flow in the appropriate scaling limit . The six velocities are given by $$𝐜_j=\widehat{𝐱}\mathrm{cos}\left(\frac{2\pi j}{6}\right)+\widehat{𝐲}\mathrm{sin}\left(\frac{2\pi j}{6}\right),$$ (132) for $`0j5`$. The model conserves mass with density $$\rho =\underset{j=0}{\overset{5}{}}N_j$$ (133) and momentum with density $$𝝅=\underset{j=0}{\overset{5}{}}N_j𝐜_j,$$ (134) so the hydrodynamic bras are $`𝝆\|`$ $`=`$ $`\left(\begin{array}{cccccc}+1& +1& +1& +1& +1& +1\end{array}\right)`$ (136) $`𝝅_x\|`$ $`=`$ $`\left(\begin{array}{cccccc}+1& +\frac{1}{2}& \frac{1}{2}& 1& \frac{1}{2}& +\frac{1}{2}\end{array}\right)`$ (138) $`𝝅_y\|`$ $`=`$ $`\left(\begin{array}{cccccc}0& +\frac{\sqrt{3}}{2}& +\frac{\sqrt{3}}{2}& 0& \frac{\sqrt{3}}{2}& \frac{\sqrt{3}}{2}\end{array}\right)`$ (140) and a reasonable choice for the kinetic bras might be $`𝜶\|`$ $`=`$ $`\left(\begin{array}{cccccc}+1& 1& +1& 1& +1& 1\end{array}\right)`$ (142) $`𝜷\|`$ $`=`$ $`\left(\begin{array}{cccccc}1& +\frac{1}{2}& +\frac{1}{2}& 1& +\frac{1}{2}& +\frac{1}{2}\end{array}\right)`$ (144) $`𝜸\|`$ $`=`$ $`\left(\begin{array}{cccccc}0& +\frac{\sqrt{3}}{2}& \frac{\sqrt{3}}{2}& 0& +\frac{\sqrt{3}}{2}& \frac{\sqrt{3}}{2}\end{array}\right).`$ (146) The kets can then be found as with the preceeding examples; the hydrodynamic kets are $$\begin{array}{ccc}\|𝝆=\frac{1}{6}\left(\begin{array}{c}\hfill +1\\ \hfill +1\\ \hfill +1\\ \hfill +1\\ \hfill +1\\ \hfill +1\end{array}\right),\hfill & \|𝝅_x=\frac{1}{6}\left(\begin{array}{c}\hfill +2\\ \hfill +1\\ \hfill 1\\ \hfill 2\\ \hfill 1\\ \hfill +1\end{array}\right),\hfill & \|𝝅_y=\frac{\sqrt{3}}{6}\left(\begin{array}{c}\hfill 0\\ \hfill +1\\ \hfill +1\\ \hfill 0\\ \hfill 1\\ \hfill 1\end{array}\right),\hfill \end{array}$$ (148) and the kinetic kets are $$\begin{array}{ccc}\|𝜶=\frac{1}{6}\left(\begin{array}{c}\hfill +1\\ \hfill 1\\ \hfill +1\\ \hfill 1\\ \hfill +1\\ \hfill 1\end{array}\right),\hfill & \|𝜷=\frac{1}{6}\left(\begin{array}{c}\hfill 2\\ \hfill +1\\ \hfill +1\\ \hfill 2\\ \hfill +1\\ \hfill +1\end{array}\right),\hfill & \|𝜸=\frac{\sqrt{3}}{6}\left(\begin{array}{c}\hfill 0\\ \hfill +1\\ \hfill 1\\ \hfill 0\\ \hfill +1\\ \hfill 1\end{array}\right).\hfill \end{array}$$ (149) The general representation of the state of the FHP model is then $`\|𝐍`$ $`=`$ $`\rho \|𝝆+\pi _x\|𝝅_x+\pi _y\|𝝅_y+\alpha \|𝜶+\beta \|𝜷+\gamma \|𝜸`$ (156) $`=`$ $`{\displaystyle \frac{\rho }{6}}\left(\begin{array}{c}1+2\overline{\pi }_x+\overline{\alpha }2\overline{\beta }\\ 1+\overline{\pi }_x+\sqrt{3}\overline{\pi }_y\overline{\alpha }+\overline{\beta }+\sqrt{3}\overline{\gamma }\\ 1\overline{\pi }_x+\sqrt{3}\overline{\pi }_y+\overline{\alpha }+\overline{\beta }\sqrt{3}\overline{\gamma }\\ 12\overline{\pi }_x\overline{\alpha }2\overline{\beta }\\ 1\overline{\pi }_x\sqrt{3}\overline{\pi }_y+\overline{\alpha }+\overline{\beta }+\sqrt{3}\overline{\gamma }\\ 1+\overline{\pi }_x\sqrt{3}\overline{\pi }_y\overline{\alpha }+\overline{\beta }\sqrt{3}\overline{\gamma }\end{array}\right).`$ In addition, the factors of $`\sqrt{3}`$ could be absorbed into $`\overline{\pi }_y`$ and $`\overline{\gamma }`$ to further simplify the final expression. ### 5.2 FCHC Model Early attempts to generalize the FHP model to obtain a three-dimensional lattice model of the Navier-Stokes equations were thwarted by the observation that no regular three-dimensional lattice would yield the required isotropy – as the triangular lattice did in two dimensions. The problem was solved by noting that a lattice with the required isotropy existed in four dimensions, and its three-dimensional projection could be used to construct an isotropic model . The four dimensional lattice is called the Face-Centered Hypercubic (FCHC) lattice. It is a self-dual lattice with coordination number 24. It can be defined as all sites $`(i,j,k,l)`$ on a Cartesian lattice such that $`i+j+k+l`$ is even. The 24 lattice vectors – or, equivalently, the 24 lattice sites that neighbor the origin $`(0,0,0,0)`$ – are enumerated in Table 2. To project these lattice vectors to three-dimensional space, we simply ignore the ficticious $`l`$ coordinate. When we do this, we note that 6 of the resulting three-dimensional lattice vectors will have two preimages from the original set of 24 FCHC lattice vectors. Rather than maintain these as separate three-dimensional vectors (as is done in lattice-gas studies), it is possible to combine them in lattice Boltzmann studies, simply weighting contributions from those directions by a factor of two. In fact, this weighting factor can be thought of as a direction-dependent particle mass $`m_j`$, so that the three-dimensional projection of the FCHC model is as presented in Table 3. The model is required to conserve mass and three components of momentum, so that the four hydrodynamic bras are given by $`𝝆\|_j`$ $`=`$ $`m_j`$ (157) $`𝝅_x\|_j`$ $`=`$ $`m_jc_{j}^{}{}_{x}{}^{}`$ (158) $`𝝅_y\|_j`$ $`=`$ $`m_jc_{j}^{}{}_{y}{}^{}`$ (159) $`𝝅_z\|_j`$ $`=`$ $`m_jc_{j}^{}{}_{z}{}^{}.`$ (160) These are shown in the first four rows of Table 4 in Appendix B, followed by a choice of twenty linearly independent kinetic bras. The corresponding kets are shown in Table 5 in Appendix B, and these can be used to construct a general analytic form for the state ket, $`\|𝐍`$. ### 5.3 Attainable Transport Coefficient As with the one-dimensional diffusion model, we expect that the entropic lattice Boltzmann model will allow us to reduce the transport coefficient of the model. For fluids, this would mean smaller shear viscosity and hence higher Reynolds number. To see if this is possible, we consider the FHP model in the incompressible limit. In the incompressible limit, the Mach number scales like the Knudsen number, so the Chapman-Enskog method involves only the equilibrium distribution at zero momentum. Taking $`\zeta _j(x)=x`$, it is straightforward to extremize $`h`$ to find the equilibrium distribution. This may be done perturbatively in Mach (equivalently, Knudsen) number, and the result to tenth order is $`\overline{\alpha }^{\text{eq}}`$ $`=`$ $`2\pi _x\left(\pi _x^23\pi _y^2\right)\left(1\pi _x^2\pi _y^2\right)^2\left(1+2\pi _x^2+2\pi _y^2\right)+𝒪\left(\pi ^{10}\right)`$ (161) $`\overline{\beta }^{\text{eq}}`$ $`=`$ $`\pi _y^2\pi _x^2\left(\pi _x^46\pi _x^2\pi _y^2+\pi _y^4\right)+`$ (162) $`2\left(\pi _x^65\pi _x^4\pi _y^25\pi _x^2\pi _y^4+\pi _y^6\right)+\left(\pi _x^2+\pi _y^2\right)^2\left(3\pi _x^4+6\pi _x^2\pi _y^25\pi _y^4\right)+𝒪\left(\pi ^{10}\right)`$ $`\overline{\gamma }^{\text{eq}}`$ $`=`$ $`2\pi _x\pi _y\left[12\left(\pi _x^2\pi _y^2\right)+4\left(\pi _x^4\pi _y^4\right)2\left(\pi _x^2+\pi _y^2\right)^2\left(3\pi _x^2+\pi _y^2\right)\right]+𝒪\left(\pi ^{10}\right).`$ (163) Thus the equilibrium values of $`\overline{\beta }`$ and $`\overline{\gamma }`$ are of second order, while that of $`\overline{\alpha }`$ is actually third order. For the computation of viscosity in the incompressible limit, the Chapman-Enskog analysis requires the linearized collision operator only to zeroth order in the Mach number. Thus, for the purposes of working out the attainable viscosity, we may simply set the equilibrium values of all three kinetic parameters to zero. It follows that the outgoing kinetic parameters are equal to the incoming ones multiplied by $`z=11/\tau `$. We can now construct an entropic lattice Boltzmann model. Using Eq. (156) for guidance, we see that the analog of Eq. (73) for $`z_{}`$ is $`\left(1+2\overline{\pi }_x+\overline{\alpha }2\overline{\beta }\right)\left(1+\overline{\pi }_x+\sqrt{3}\overline{\pi }_y\overline{\alpha }+\overline{\beta }+\sqrt{3}\overline{\gamma }\right)\left(1\overline{\pi }_x+\sqrt{3}\overline{\pi }_y+\overline{\alpha }+\overline{\beta }\sqrt{3}\overline{\gamma }\right)`$ $`\left(12\overline{\pi }_x\overline{\alpha }2\overline{\beta }\right)\left(1\overline{\pi }_x\sqrt{3}\overline{\pi }_y+\overline{\alpha }+\overline{\beta }+\sqrt{3}\overline{\gamma }\right)\left(1+\overline{\pi }_x\sqrt{3}\overline{\pi }_y\overline{\alpha }+\overline{\beta }\sqrt{3}\overline{\gamma }\right)=`$ $`\left(1+2\overline{\pi }_x+z_{}\overline{\alpha }2z_{}\overline{\beta }\right)\left(1+\overline{\pi }_x+\sqrt{3}\overline{\pi }_yz_{}\overline{\alpha }+z_{}\overline{\beta }+\sqrt{3}z_{}\overline{\gamma }\right)\left(1\overline{\pi }_x+\sqrt{3}\overline{\pi }_y+z_{}\overline{\alpha }+z_{}\overline{\beta }\sqrt{3}z_{}\overline{\gamma }\right)`$ $`\left(12\overline{\pi }_xz_{}\overline{\alpha }2z_{}\overline{\beta }\right)\left(1\overline{\pi }_x\sqrt{3}\overline{\pi }_y+z_{}\overline{\alpha }+z_{}\overline{\beta }+\sqrt{3}z_{}\overline{\gamma }\right)\left(1+\overline{\pi }_x\sqrt{3}\overline{\pi }_yz_{}\overline{\alpha }+z_{}\overline{\beta }\sqrt{3}z_{}\overline{\gamma }\right),`$ (164) where it bears repeating that, just for the purposes of this computation of viscosity, we have neglected corrections to the equilibrium of order Mach number squared <sup>6</sup><sup>6</sup>6Note that these terms of order Mach number squared are necessary to derive the remainder of the Navier-Stokes equations.. This is a sixth-order equation for $`z_{}`$; since $`z_{}=1`$ is clearly a root, it can be reduced to fifth order. In spite of the fact that fifth order equations generally do not admit solutions in radicals, this one does happen to do so. The result, however, is complicated and uninspiring, and so we omit it here. The interested reader can simply type the above equation into, e.g., $`\text{Mathematica}^\mathrm{\copyright }`$, and use the Solve utility to see the result. The important point is that the relevant root for $`z_{}`$ approaches $`1`$ as $`\overline{\alpha },\overline{\beta },\overline{\gamma }0`$. The rest of the analysis is exactly as it was for the one-dimensional diffusion model. We set $`z=1\kappa (1z_{})`$, where $`0<\kappa <1`$. Since the usual lattice Boltzmann version of the FHP model has the shear viscosity proportional to $`\tau 1/2`$, the entropic version will have it proportional to $`1/\kappa 1`$ which goes to zero as $`\kappa `$ approaches unity. Moreover, absolute stability is maintained as this limit is approached, even though the algorithm is fully explicit and perfectly conservative. As with the diffusion model, we are prevented from achieving arbitrarily small transport coefficient only by considerations of accuracy (resolution of the smallest eddies), and of the time scale for the approach to equilibrium. ### 5.4 Computer Implementation We pause to consider the computer implementation of the entropic lattice Boltzmann algorithm for the FHP fluid. The propagation step can be handled in precisely the same fashion that it is for conventional lattice Boltzmann simulations. The principal difference lies in the implementation of the collision operator, and that is what we describe here. Unlike the analysis in the last subsection, a complete computer implementation needs to solve for the equilibrium exactly – not just to within the Mach number squared. Given the six quantities, $`N_j`$ for $`j=0,\mathrm{},5`$, entering a site, we first contract them with the bras, Eqs. (136) through (146), to obtain the conservation representation, namely $`\rho `$, $`\pi _x`$, $`\pi _y`$, $`\alpha `$, $`\beta `$ and $`\gamma `$. Given $`\pi _x`$ and $`\pi _y`$, it is then necessary to solve for the equilibrium values of $`\overline{\alpha }`$, $`\overline{\beta }`$ and $`\overline{\gamma }`$. We can use Eqs. (161) through (163) to get excellent approximations for these, and then use a Newton-Raphson iteration to get the correct values to machine precision. Even with an excellent initial guess, there is some potential for loss of convergence in any Newton-Raphson iteration, and so it might be worthwhile to investigate better ways of solving for these values. Assume that we can do this, and denote the equilibrium values with an “eq” superscript. Given the equilibrium values, we set $`\overline{\alpha }^{}`$ $`=`$ $`\overline{\alpha }+{\displaystyle \frac{1}{\tau }}\left(\overline{\alpha }^{\text{eq}}\overline{\alpha }\right)`$ (165) $`\overline{\beta }^{}`$ $`=`$ $`\overline{\beta }+{\displaystyle \frac{1}{\tau }}\left(\overline{\beta }^{\text{eq}}\overline{\beta }\right)`$ (166) $`\overline{\gamma }^{}`$ $`=`$ $`\overline{\gamma }+{\displaystyle \frac{1}{\tau }}\left(\overline{\gamma }^{\text{eq}}\overline{\gamma }\right),`$ (167) and find the value of $`\tau `$ that solves the nonlinear equation $`h(𝐍^{})=h(𝐍)`$. This equation will look much like Eq. (164), except that $`\overline{\alpha }^{\text{eq}}`$, $`\overline{\beta }^{\text{eq}}`$ and $`\overline{\gamma }^{\text{eq}}`$ will enter, corresponding to corrections of order Mach number squared. Since this nonlinear equation has only one scalar unknown, namely $`\tau `$, it can be solved using an iteration method that is guaranteed to converge, such as regula falsi. We know that $`\tau =1`$ is one bound on the solution; the other bound can be taken as the point at which the line of outgoing states intersects the boundary of the kinetic polytope. This means finding the largest value of $`\tau `$ between zero and one for which $`N_j^{}=0`$ for some $`j\{0,\mathrm{},5\}`$. Using these two bounds, we iterate the regula falsi algorithm to obtain the solution $`\tau _{}`$. We then set $`\tau =\tau _{}/\kappa `$, and use Eqs. (165) through (167) to get the outgoing values of the kinetic parameters. We finish the collision procedure by using Eq. (156) to get the outgoing state in the original representation. ### 5.5 Galilean Invariance As noted in the Introduction, practitioners of lattice BGK models have long known that it is necessary to tailor the equilibrium distribution function in a certain way in order to maintain Galilean invariance. Recall that the zeroth moment (with respect to the velocity vector) of the distribution function is determined by the mass density. Likewise, the first moment is determined by the momentum density. For a lattice BGK model of an ideal gas , the trace of the second moment is determined by the kinetic energy density. To obtain the correct, Galilean-invariant form of the compressible Navier-Stokes equations, it is also necessary to mandate the rest of the second moment, the third moment, and the trace of the fourth moment. While it is possible to construct equilibrium distributions with these moments for various lattices, no lattice model has ever been found for which these equilibria are guaranteed to be nonnegative. This means that the equilibrium toward which we would like to relax in order to maintain Galilean invariance may itself lie outside of the master polytope! There is thus no way that such an equilibrium could be the maximum of an entropy function that goes to $`\mathrm{}`$ on the polytope boundary, such as do the ones that we have been considering. In such a situation, we may ask whether it is possible to jettison the requirement of nonnegativity, while still demanding absolute stability. In fact, our formalism suggests an interesting way of doing just that. If we no longer demand that the function $`h`$ go to negative infinity on the polytope boundary, then we need no longer demand Eq. (22), though Eq. (23) is still useful to ensure that the equilibrium is unique, and Eq. (24) is useful to keep $`h`$ real-valued in light of Eq. (25). Thus, we are free to make the choice $`\zeta _j(x)`$ is $`\mathrm{exp}[x^2/(2g_j)]`$, where the $`g_j`$’s are positive constants. Indeed, this choice is similar to one made by Karlin et al. . Eq. (27) is then the set of $`n`$ equations, $$N_j^{\text{eq}}=g_j\underset{\sigma =1}{\overset{n_c}{}}Q_\sigma \|𝝀_\sigma _j,$$ (168) and Eqs. (28) give $`n_c`$ additional requirements. One may then try to solve these $`n+n_c`$ equations for the $`n+n_c`$ unknowns, $`Q_\sigma `$ and $`g_j`$, subject to the requirement that $`g_j>0`$. In fact, not all of these unknowns are independent; the scaling $`g_jg_j/x`$ and $`Q_\sigma Q_\sigma x`$, where $`x`$ is arbitrary, leaves Eq. (168) invariant, so there are more requirements than equations, and hence not all equilibria are derivable from an $`h`$ function of this form. Nevertheless, many interesting equilibria are derivable in this way, and this approach leads immediately to the construction of an entropically stabilized algorithm for such equilibria. Investigation of this approach is work in progress . Its successful resolution may yield an absolutely stable entropic lattice Boltzmann algorithm for compressible flow, albeit one that allows for negative $`N_j`$. ## 6 Conclusions In summary, we have presented a general methodology for constructing lattice Boltzmann models of fluid dynamics that conserve a given set of moments of the distribution function. The method guarantees the nonnegativity of the distribution function, and allows for convenient visualization of the dynamics by construction of the polytope of allowed states using Fourier-Motzkin elimination. We have shown how such models can be endowed with a Lyapunov functional, resulting in unconditional numerical stability, and we presented various choices for the $`H`$ function for this purpose. We also described the computer implementation of such models in some detail. For lattice Boltzmann models of diffusion and of incompressible fluid flow, we were able to present fully explicit, perfectly conservative, absolutely stable algorithms that appear to work for arbitrarily small transport coefficient. Indeed, we showed that the limitations on attainable transport coefficients for these models arise from considerations of accuracy, rather than stability. There are numerous avenues for extension of the current work. We close by presenting a list of these. We hope that the richness and promise of this style of model will inspire others to take up some of these questions. 1. We have not discussed the question of boundary conditions in this paper. For a viscous fluid, for example, collisions at a solid boundary wall must conserve mass, but not momentum or energy. The latter two quantities can enter or leave the domain through the wall, resulting in drag, lift, surface heating, etc. This results in very different polytope structure at the boundaries than in the interior. In addition, there is nothing preventing entropy from entering or leaving through the wall, so the entire question of how to construct $`h`$ functions at the boundary sites needs to be revisited. 2. While we touched on the question of optimality of conservation representations, clearly much more work could be done along these lines. In addition to making it computationally easier to transform back and forth between the original and conservation representations, one might like to try to optimize the form of the entropically stabilized collision operator. 3. We mentioned reversible models, but did not discuss or study them in any detail. Our entropically stabilized models are reversible in the limit as $`\kappa 1`$, but this is precisely the limit in which they begin to fail to simulate the parabolic equations of interest, namely the diffusion and Navier-Stokes equations. It would be interesting to know if such reversible models are capable of simuating time-symmetric partial differential equations, such as Euler’s equations of inviscid fluid dynamics. 4. The nature of the anomalies encountered as $`\kappa 1`$ remains to be completely elucidated. For example, the lower right plot in Fig. 7 appears to indicate that the approach to equilibrium is extremely slow for this model. Does it eventually reach an exponentially decaying equilibrium state for which the diffusion coefficient matches that of the theory? Does roundoff error play a role in this regime? These open questions will require substantially more numerical simulation than that presented here. 5. Given that the anomalies as $`\kappa 1`$ are due to the slowness of the approach to the equilibrium state, much effort could be expended in trying to start the simulation in as close to a global equilibrium state as possible. This means taking spatial gradient corrections into account. At a minimum, the Chapman-Enskog corrections – accurate to first-order in the spatial gradient – could be incorporated with little effort. 6. The question of how to deal with interaction potentials in this formalism could be addressed. At the level of mean-field theory, this would require the introduction of conserved quantities that are quadratic in the distribution function. The allowed regions would then no longer be polytopes, convex, or even simply connected. Moreover, multiple local equilibria within the allowed region may be expected. For near-equilibrium dynamics, this might give rise to a Ginzburg-Landau model with multiple free-energy minima. This is thus a difficult extension of the current theory, but one for which the rewards would be great. 7. One could try to construct an adaptive mesh refinement (AMR) version of this algorithm for which the nonnegativity of the physical densities is the refinement criterion. Using such an approach, it might be possible to relax the criterion of nonnegativity of the distribution function. That is, we could carry negative distribution function values, but insist that all physical densities be positive. Should the propagation step threaten to create a negative density, we would try to refine the lattice locally until that is no longer the case. 8. Finally, it would be interesting to consider the potential utility of this algorithm for the construction of “eddy-viscosity” subgrid models of turbulence. As noted at several points in this paper, while entropic lattice Boltzmann algorithms may be stable for arbitrarily small transport coefficient, they may begin to lose accuracy. The smallest eddy sizes in three dimensional Navier-Stokes turbulence, for example, scale as $`\text{Re}^{3/4}`$, where the Reynolds number Re goes inversely with viscosity. Thus, at sufficiently small viscosity, the eddies will be smaller than any fixed grid spacing. If one’s goal is an ab initio simulation of the Navier-Stokes equations, this circumstance is grounds for rejecting the entropic lattice Boltzmann solution as unphysical. On the other hand, stable, coarse-grained, turbulent flow must obey the same conservation laws as the underlying Navier-Stokes equations. Since the entropic lattice Boltzmann collision operators are the most general that respect these conservation laws, it is interesting to speculate as to whether some physical significance might yet be attached to their solutions, even when the smallest eddies are not resolved by the lattice. ## Acknowlegements This work was presented by BMB at the Seventh International Conference on the Discrete Simulation of Fluids in Oxford, England in July, 1998. During the late stages of writing this paper, we became aware that two other groups were independently working on requirements of the form of Eq. (31). As the works were conducted independently, there are numerous difference in style, formulation and emphasis, and we believe that the resulting papers, including this one, are nicely complementary to one another. It is a pleasure to acknowledge helpful conversations with Frank Alexander, Bastien Chopard, Nicos Martys, Sauro Succi and Julia Yeomans. BMB was supported in part by an IPA assignment agreement with the Air Force Research Laboratory, Hanscom AFB, Massachusetts, and in part by the United States Air Force Office of Scientific Research under grant number F49620-95-1-0285. BMB and PVC would like to thank Schlumberger Research and the EPSRC’s Collaborative Computing Project No. 5 for making possible BMB’s visit to Schlumberger Cambridge Research and the University of Oxford, where this work was begun in early 1997. PVC is grateful to Wolfson College and the Department of Theoretical Physics, University of Oxford, for a Visiting Fellowship (1996 – 1999). ## Appendix A A Tour of the Master Polytope for the Compressible-Fluid Model The five inequalities of Eq. (125) can be recast as $$0\|\overline{\beta }\|\overline{\alpha }\mathrm{min}(\frac{2\overline{\epsilon }}{3},\overline{\epsilon }\|2\overline{\beta }\overline{\pi }\|).$$ (169) This implies $$\overline{\epsilon }\overline{\alpha }+\|2\overline{\beta }\overline{\pi }\|\overline{\alpha }\|\overline{\beta }\|0$$ (170) and $$\overline{\epsilon }23\overline{\alpha }23\|\overline{\beta }\|2,$$ (171) and hence $`\overline{\alpha }0`$ and $`\overline{\epsilon }[0,2]`$. Moreover, we note that all of the inequalities are invariant under the transformation $`\overline{\pi }`$ $``$ $`\overline{\pi }`$ $`\overline{\beta }`$ $``$ $`\overline{\beta },`$ so that we may restrict our attention to $`\overline{\pi }0`$ without loss of generality. Next note that the second argument of the $`\mathrm{min}`$ function in Eq. (169) will be smaller than the first if $`0\overline{\epsilon }<(2\overline{\epsilon })/3`$, or $`0\overline{\epsilon }<1/2`$. In this case, for sufficiently small $`\overline{\pi }`$, the boundary of the polytope in the $`(\overline{\alpha },\overline{\beta })`$ plane will be a quadrilateral, such as that in Fig. 10, whose lower bound in $`\overline{\alpha }`$ (shown in black) is $`\|\overline{\beta }\|`$ and whose upper bound in $`\overline{\alpha }`$ (shown in gray) is $`\overline{\epsilon }\|2\overline{\beta }\overline{\pi }\|`$. The allowed region of the $`(\overline{\alpha },\overline{\beta })`$ plane is the shaded area in between these bounds. The bottom vertex <sup>7</sup><sup>7</sup>7Throughout this appendix, we denote vertices in parentheses, so as not to confuse them with other variables. $`(a)`$ of this quadrilateral is at the origin, $`(\overline{\alpha },\overline{\beta })=(0,0)`$, and the upper vertex $`(c)`$ is at $`(\overline{\alpha },\overline{\beta })=(\overline{\epsilon },\overline{\pi }/2)`$. The right vertex $`(b)`$ is then at $`\overline{\beta }=\overline{\epsilon }(2\overline{\beta }\overline{\pi })`$ or $`(\overline{\alpha },\overline{\beta })=((\overline{\epsilon }+\overline{\pi })/3,(\overline{\epsilon }+\overline{\pi })/3)`$, and the left vertex $`(d)`$ is then at $`\overline{\beta }=\overline{\epsilon }+(2\overline{\beta }\overline{\pi })`$ or $`(\overline{\alpha },\overline{\beta })=((\overline{\epsilon }\overline{\pi })/3,(\overline{\epsilon }+\overline{\pi })/3)`$. These results are summarized in the table included in Fig. 10. We now ask for what range of $`\overline{\epsilon }`$ and $`\overline{\pi }`$ the above-described quadrilateral boundary in the $`(\overline{\alpha },\overline{\beta })`$ plane is valid. We note that vertices $`(a)`$ and $`(b)`$ will be degenerate when $`\overline{\epsilon }=\overline{\pi }`$, and that vertices $`(a)`$ and $`(d)`$ will be degenerate when $`\overline{\epsilon }=+\overline{\pi }`$. Also, it is easy to see that one or the other of these two degeneracies will occur before any degeneracy involving vertex $`(c)`$. It follows that the above-described quadrilateral boundary in the $`(\overline{\alpha },\overline{\beta })`$ plane is valid only for $`\|\overline{\pi }\|<\overline{\epsilon }<1/2`$. This region of the $`(\overline{\pi },\overline{\epsilon })`$ plane is the shaded triangle $`AIH`$ in Fig. 9. Beginning in triangle $`AIH`$, boundary $`AI`$ is encountered when $`\overline{\epsilon }=\overline{\pi }`$ so that vertices $`(a)`$ and $`(d)`$ are degenerate. If we cross this boundary the allowed region in the $`(\overline{\alpha },\overline{\beta })`$ plane is no longer a quadrilateral, but rather becomes a triangle, such as that in Fig. 11, whose upper bound in $`\overline{\alpha }`$ (shown in gray) is $`\overline{\epsilon }\|2\overline{\beta }\overline{\pi }\|`$ and whose lower bound in $`\overline{\alpha }`$ (shown in black) is $`\overline{\beta }`$. The bottom vertex $`(e)`$ of this triangle is at $`\overline{\alpha }=\overline{\beta }=\overline{\epsilon }+(2\overline{\beta }\overline{\pi })`$ or $`(\overline{\alpha },\overline{\beta })=(\overline{\pi }\overline{\epsilon },\overline{\pi }\overline{\epsilon })`$; this vertex replaces the degenerate vertices $`(a)`$ and $`(d)`$ of the above-described quadrilateral. The expressions for the coordinates of the triangle’s upper vertex $`(c)`$ and its right vertex $`(b)`$ are identical to those of the corresponding vertices of the quadrilateral. These results are summarized in the table included in Fig. 11. Again, we ask for what range of $`\overline{\epsilon }`$ and $`\overline{\pi }`$ the above-described triangular boundary in the $`(\overline{\alpha },\overline{\beta })`$ plane is valid. We note that vertices $`(b)`$ and $`(c)`$ will be degenerate when $`(\overline{\epsilon }+\overline{\pi })/3=\overline{\epsilon }`$ or $`\overline{\epsilon }=\overline{\pi }/2`$; for smaller values of $`\overline{\epsilon }`$ the set of allowed points in the $`(\overline{\alpha },\overline{\beta })`$ plane is empty. We also note that, just as for the quadrilateral of Fig. 10, the first argument of the $`\mathrm{min}`$ function in Eq. (169) will become the determining factor when $`\overline{\epsilon }>1/2`$, also invalidating the above argument. Thus, the allowed region in the $`(\overline{\alpha },\overline{\beta })`$ plane will be a triangle with vertices described in Fig. 11 only if $`(\overline{\epsilon },\overline{\pi })`$ is in the shaded triangle $`ABI`$ in Fig. 9. Since the region $`AEH`$ is the image of $`ABI`$ under the map $`\overline{\pi }\overline{\pi }`$, it follows that the allowed region of the $`(\overline{\alpha },\overline{\beta })`$ plane is also triangular for $`(\overline{\epsilon },\overline{\pi })\mathrm{}AEH`$, with vertices identical to those in Fig. 11 with $`\overline{\beta }\overline{\beta }`$. We have now described all of Fig. 9 below the line $`\overline{\epsilon }=1/2`$ (that is, below line $`EB`$). We next consider the situation for $`\overline{\epsilon }>1/2`$. This means that we have to start taking into account the inequality $$\overline{\alpha }\frac{2\overline{\epsilon }}{3}$$ (172) of Eq. (169). This is a constant upper bound on $`\overline{\alpha }`$ which will become less than the $`\overline{\alpha }`$ coordinate of vertex $`(c)`$ in Figs. 10 and 11 when $`\overline{\epsilon }>1/2`$. This results in a horizontal truncation of the quadrilateral of Fig. 10 so that it becomes a pentagon, and of the triangle of Fig. 11 so that it becomes a quadrilateral. These are shown in Figs. 12 and 13, respectively. The top vertices $`(f)`$ and $`(g)`$ in both of these figures are at the upper bound $`\overline{\alpha }=(2\overline{\epsilon })/3`$, and $`(2\overline{\epsilon })/3=\overline{\epsilon }(2\overline{\beta }v)`$ or $`\overline{\beta }=(3v\pm 4\overline{\epsilon }2)/6`$, respectively; these vertices replace vertex $`(c)`$ in Figs. 10 and 11. The expressions for the coordinates of vertices $`(a)`$, $`(b)`$, $`(d)`$ and $`(e)`$ are identical to those derived previously. The vertex coordinates are given in the tables in their corresponding figures. Yet again, we ask for what range of $`\overline{\epsilon }`$ and $`\overline{\pi }`$ the boundaries pictured in Figs. 12 and 13 are valid. As $`\overline{\epsilon }`$ increases, the upper bound on $`\overline{\alpha }`$ decreases until vertices $`(b)`$ and $`(f)`$ coincide. This happens when $`(\overline{\epsilon }+\overline{\pi })/3=(2\overline{\epsilon })/3`$, or $`\overline{\epsilon }=1\overline{\pi }/2`$. This is line $`FB`$ in Fig. 9. Above this line, the pentagonal region of Fig. 12 degenerates to a quadrilateral, and the quadrilateral region of Fig. 13 degenerates to a triangle. In both cases, vertices $`(b)`$ and $`(f)`$ are replaced by a new vertex $`(h)`$ whose coordinates are $`\overline{\alpha }=\overline{\beta }=(2\overline{\epsilon })/3`$. These situations are shown in Figs. 14 and 15, along with corresponding tables of vertex coordinates. The quadrilateral of Fig. 14 will degenerate into the triangle of Fig. 15 when vertices $`(a)`$ and $`(d)`$ merge and are replaced by vertex $`(e)`$. As before, this degeneracy happens when $`\overline{\epsilon }=\overline{\pi }`$, and this is the boundary line $`CJ`$ in the illustration of the $`(\overline{\epsilon },\overline{\pi })`$ plane shown in Fig. 9. The triangle of Fig. 15, in turn, degenerates into the empty set when $`(2\overline{\epsilon })/3=\overline{\pi }\overline{\epsilon }`$, or $`\overline{\epsilon }=3\overline{\pi }/21`$. Referring to Fig. 9, this is the boundary line $`BC`$; thus, we see that the quadrilaterals of Fig. 14 are obtained when $`(\overline{\epsilon },\overline{\pi })\mathrm{}CFJ`$, and the triangles of Fig. 15 are obtained when $`(\overline{\epsilon },\overline{\pi })\mathrm{}BCJ`$. Finally, there is one other way that the quadrilaterals of Fig. 14 can degenerate. Vertices $`(d)`$ and $`(g)`$ will coincide if $`(2\overline{\epsilon })/3=(\overline{\epsilon }\overline{\pi })/3`$, or $`\overline{\epsilon }=1+\overline{\pi }/2`$. For values of $`\overline{\epsilon }`$ greater than this, vertices $`(d)`$ and $`(g)`$ are replaced by vertex $`(i)`$ with coordinates given by $`\overline{\alpha }=\overline{\beta }=(2\overline{\epsilon })/3`$. The resulting isoceles triangular region is shown in Fig. 16, along with a corresponding table of vertex coordinates. This triangluar region degenerates to the empty set only when $`\overline{\epsilon }=2`$. Referring to Fig. 9, we see that the triangles of Fig. 16 are obtained when $`(\overline{\epsilon },\overline{\pi })\mathrm{}CFL`$; in fact, since these regions are symmetric in $`\overline{\beta }`$, they are also symmetric in $`\overline{\pi }`$, so their description is the same for all $`(\overline{\epsilon },\overline{\pi })\mathrm{}CDF`$. ## Appendix B Basis Bras and Kets for FCHC Model In this appendix, we present one possible choice for the bras and kets of the FCHC model for three-dimensional fluid dynamics. The basis bras are shown in Table 4, and the basis kets are shown in Table 5.

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